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# Kinematics of Metal-Poor Stars in the Galaxy. II. Proper Motions for a Large Non-Kinematically Selected Sample
## 1 Introduction
Studies of the kinematics of various stellar populations in the Galaxy, in particular the thick disk and the nascent halo, have long been limited by the availability of large samples of stars with measurements of velocities, distances, and metallicities. Such a database is required in order to constrain plausible scenarios for the formation and evolution of the Milky Way and other large spiral galaxies like it. Current issues which might be addressed with such data include: (a) the rotational character of the thick disk and halo (see, e.g., Beers & Sommer-Larsen 1995, hereafter BSL, and references therein), (b) the existence and the observed lower limit on the metal abundance of stars in the so-called metal-weak thick disk (MWTD) (Morrison, Flynn, & Freeman 1990, hereafter MFF; BSL; Carney et al. 1996; Chiba & Yoshii 1998; Martin & Morrison 1998) (c) a dual-component (flattened plus spherical) halo population in the Galaxy (Hartwick 1987; Sommer-Larsen & Zhen 1990; Norris 1994; Kinman, Suntzeff, & Kraft 1994; Sommer-Larsen et al. 1997), (d) quantitative estimates of the local density of thick disk and halo stars (Yoshii 1982; Preston, Shectman, & Beers 1991; Morrison 1993), (e) tests for the existence of a putative “counter-rotating” halo component beyond 4–5 kpc from the plane (Majewski 1992; Wilhelm 1995; Carney et al. 1996; Wilhelm et al. 1996; Zinn 1996; Carney 1999), (f) measures of the local halo velocity ellipsoid for comparison with the derived ellipsoid for more distant halo stars (Sommer-Larsen et al. 1997), and (g) derivation of a reliable RR Lyrae absolute magnitude estimate based on statistical parallax analyses (Layden et al. 1996; Fernley et al. 1998; Popowski & Gould 1998). As more extensive searches are carried out for evidence of past (and present) mergers of smaller galaxies with the Milky Way (e.g., Preston, Beers, & Shectman 1994; Harding et al. 1998; Helmi & White 1999; Majewski 1999), it is of equal importance to obtain secure knowledge of the “global” Galactic kinematic properties, so that deviations from the expected behavior can be reliably assessed.
To confidently address the above issues (and many others) stars chosen as kinematic tracers should be identified in a manner which does not depend on a kinematic selection criterion itself. Although it may be possible to statistically correct for an input selection bias of this nature (Bahcall & Casertano 1986; Ryan & Norris 1991a; Carney et al. 1994; Carney 1999), one is left with lingering doubt concerning the derived kinematic parameters based on post-facto modifications of the results. It is similarly important that the tracer stars cover a wide range of metallicities and distances (over both northern and southern Galactic hemispheres), so that correlations of kinematics as a function of these parameters can be investigated.
Although hints of the impact of a kinematic selection criterion are evident in the work of Yoshii & Saio (1979), the first large database suitable for exploration of many of these issues was that of Norris (1986), which included some 400 spectroscopically and/or photometrically selected stars with abundances $`[\mathrm{Fe}/\mathrm{H}]0.6`$, and with available radial velocities and distance estimates. In Paper I of this series BSL extended the Norris catalog by inclusion of some 900 stars identified from the HK object objective-prism survey of Beers & colleagues (Beers, Preston, & Shectman 1985; Beers, Preston, & Shectman 1992a; Beers et al. 1992b), as well as some 600 additional stars from other smaller samples which appeared in the literature subsequent to the publication of the Norris catalog, obtaining a total sample of 1936 stars.
In this paper we present a revision of the BSL catalog of Paper I, based on additional observational information which has recently become available. In addition to revisions of metallicities for the published HK survey stars, new photometric distance estimates have been made for the entire BSL catalog based on an internally self-consistent methodology, and in some cases, new photometry. Radial velocities have been updated based on recently-obtained high-resolution data for a number of stars. Much more accurate positional information for the stars in our revised catalog has been obtained by comparison with astrometric positions available for many of the brighter stars in our sample, plus improved information from automated scans of wide-field photographic plates, such as compiled in USNO-A V2.0 (Monet et al. 1998), NPM1 (Klemola, Hanson, & Jones 1993), and SPM 2.0 (Platais et al. 1998). We have also added 182 RR Lyrae variables from the recent work of Layden (1994), Layden et al. (1996), and Fernley et al. (1998), so that useful comparisons of the kinematics of these stars with the non-variables (which presumably sample the same Galactic phase space) can be carried out.
The major difference between the present catalog and the BSL catalog is the addition of proper motions, from a variety of sources, for over half of the stars in our catalog. In §2 we discuss the assemblage of the present catalog. Revisions of radial velocities and abundances, and in particular, distance estimates, are discussed in §3. In §4 we present the new proper motion information, and discuss the averaging we have carried out in order to minimize statistical errors. In §5 we discuss the observed characteristics of the stars in the revised catalog, and derive estimates of space motions and orbital parameters for the subset of stars with complete kinematical information . In the accompanying analysis paper (Chiba & Beers 2000; Paper III) we use this wealth of new information to consider many of the questions put forth above.
## 2 Additions and Subtractions from the Beers & Sommer-Larsen (1995) Catalog
The BSL catalog of Paper I contained 1936 stars with available abundances, distances, and radial velocities. As it was our intent to use this catalog as the starting point in searches for stars with measured proper motions, the first requirement was to obtain updated positions for the stars of as high an accuracy as possible. A small number of fainter stars in the BSL catalog with crude positions (and no finding charts) were eliminated entirely. The remaining stars were then compared with the positions listed in the SIMBAD database. A first search was done on the star name, followed by a search on the reported position in the BSL catalog, with the requirement that the magnitude reported in SIMBAD be commensurate with that of the target star. Many of the brighter stars from the BSL catalog now have astrometric positions from Hipparcos, which are accurate at the milli-arcsecond (mas) level. For the stars of intermediate and fainter magnitudes, searches were conducted on the BSL positions by comparison with the USNO-A 2.0 catalog, again with the requirement that the magnitudes be commensurate. The typical accuracy of stellar positions for USNO-A 2.0 is at the 0$`\stackrel{}{\mathrm{.}}`$25 (one-sigma) level (Deutsch 1999). Further searches were conducted within the catalogs we use below to find proper motions – NPM1, SPM 2.0, STARNET, and ACT – which resulted in positional accuracies on the order of 0$`\stackrel{}{\mathrm{.}}`$2, 30 mas, 0$`\stackrel{}{\mathrm{.}}`$3, and 25 mas, respectively. For a small number of stars no improved positional information could be found, hence we retained the positions in the original BSL catalog.
The Hipparcos, Tycho, NPM1, and SPM 2.0 catalogs included a number of “targeted” RR Lyrae stars which were not included in the BSL catalog, some of which have radial velocity and abundance estimates available in the literature. We thus added a total of 182 RR Lyrae stars to our revised catalog, including many with metallicities greater than the nominal cutoff of the BSL catalog. A small number of non-variables which were not included in the BSL catalog, but which have available radial velocities and abundance estimates, were also added. As a result of the revised abundances, described in more detail below, a small number of HK survey stars which were originally assigned metallicities $`[\mathrm{Fe}/\mathrm{H}]0.6`$ presently have values above this limit. These stars are noted in the revised catalog.
Column (1) of Table 1A lists the star names. We have endeavored to follow the IAU-recommended nomenclature, but in a few cases had to abbreviate the star name in order to save space. In any case, the positional identification is usually unambiguous. All adopted positions were updated to the ICRS (2000.0) system. Columns (2) and (3) of Table 1A list the right ascension and declination of the adopted positions, respectively. The source of the position is indicated in column (4).
## 3 Revisions of Abundance Estimates, Velocities, and Distances
### 3.1 Revisions of Abundance Estimates
Roughly 45 % of the stars in the revised catalog were originally identified during the course of the HK objective-prism survey. In the BSL catalog the abundance estimates for these stars were obtained from a calibration of the strength of the CaII K line as a function of $`(BV)_o`$ color described in Beers et al. (1990). This calibration is now known to suffer from several deficiencies, the most worrisome of which is the fact that assigned abundances in the range $`[\mathrm{Fe}/\mathrm{H}]>1.5`$ are less than optimal as a result of the saturation of the CaII K line, in particular at cooler temperatures. Stars in the metallicity regime $`2.0[\mathrm{Fe}/\mathrm{H}]1.0`$ represent the transition between the halo and thick disk populations, and in particular, may be members of the MWTD component of the Galaxy, hence it is important that the abundances be correctly estimated. At the low end of the metallicity scale ($`[\mathrm{Fe}/\mathrm{H}]<2.5`$), the lack of available calibrators in the Beers et al. (1990) treatment resulted in stellar abundance estimates which were, on the whole, somewhat lower than has proven to be justified on the basis of recent high-resolution spectroscopy, in particular for the hotter stars near the main-sequence turnoff of an old halo population.
The re-calibration of metallicity estimates based on medium resolution spectroscopy described by Beers et al. (1999) avoids, to a great extent, the above difficulties. This re-calibration makes use of an additional metallicity estimator based on the Auto-Correlation Function of metallic lines in a stellar spectrum (originally described by Ratnatunga & Freeman 1989), which provides a superior estimate of abundance to CaII K in the regime $`[\mathrm{Fe}/\mathrm{H}]>1.0`$. The combination of the CaII K and ACF approaches implemented in Beers et al. (1999) results in metallicity estimates over the interval $`4.0[\mathrm{Fe}/\mathrm{H}]0.0`$ with typical errors on the order of 0.1-0.2 dex, and no significant systematic offsets.
Figure 1 is a comparison of the abundance determinations of the HK survey stars in the revised catalog based on the new and old calibration. Note that the broad sloping region about the one-to-one line, which is dominated by stars with $`0.3(BV)_o0.5`$, indicates that the revised abundances for the hotter HK survey stars are generally higher at the low metallicity end of the scale, and lower at the high metallicity end of the scale, as compared to the old calibration. The abundance estimates for the cooler stars with $`0.5<(BV)_o1.2`$ have changed relatively less.
As part of the re-calibration effort, Table 7 of Beers et al. (1999) lists averaged high-resolution abundance determinations from the literature for 551 stars. For stars which also appear in the revised catalog, we adopted these averaged abundance determinations.
The catalog of Cayrel de Strobel et al. (1997) provides detailed information on abundances for stars with determinations based on high-resolution spectroscopy. Suitably averaged abundances for stars which did not appear in Beers et al. (1999) were adopted for the revised catalog listing. A number of stars from MFF have been shown by Ryan & Lambert (1995) to suffer from mis-estimated abundances (generally too low) arising from a faulty calibration of DDO photometry. For MFF stars with high-resolution abundances obtained by Ryan & Lambert, we simply adopt their values. For other MFF stars which are included in the revised catalog, the abundances should be considered uncertain.
Metallicity estimates for the revised catalog are reported in column (5) of Table 1A. The source of the metallicity estimate is indicated in column (6). Uncertain abundances are indicated with a ’:’ appended to the reported \[Fe/H\] .
### 3.2 Revision of Radial Velocities
Radial velocities of improved accuracy have been reported for a number of stars in our revised catalog based on high-resolution spectroscopic follow-up, in particular for the HK survey stars (McWilliam et al. 1995; Norris, Ryan, & Beers 1996). The adopted radial velocities are reported in column (7) of Table 1A. Errors on the radial velocity measurements are reported in column (8), and the source of the radial velocity estimate is provided in column (9). For the HK survey stars with only medium-resolution spectroscopy available, we adopt a conservative error estimate of 10 $`\mathrm{km}\mathrm{s}^1`$ (based on previous comparisons). For non-variable stars where the original authors did not report an error on the velocity determination, we have adopted an error of 10 $`\mathrm{km}\mathrm{s}^1`$. For RR Lyrae variables without reported radial velocity errors, we assign an error of 30 $`\mathrm{km}\mathrm{s}^1`$, which accounts for the uncertainty in the systemic velocity arising from limited phase coverage in spectra for a given variable (Smith 1995).
### 3.3 Revision of Distance Estimates
We have determined new photometric estimates of stellar distances for stars in our catalog based on $`M_V`$ vs. $`(BV)_o`$ relations for non-variable stars, as described below, and an adopted $`M_V`$ vs. $`[\mathrm{Fe}/\mathrm{H}]`$ relation for RR Lyrae variables.
#### 3.3.1 Apparent Magnitudes and Colors
Our primary source of $`V`$-band apparent magnitudes and $`BV`$ colors for non-variable stars is the photometry obtained during the course of the HK survey, if available (with typical accuracy in $`V`$ magnitudes and $`BV`$ colors in the range 0.005 to 0.01, respectively). Additional photometric information was taken from the SIMBAD database, the Hipparcos Catalogue, and the calibrated photographic photometry of SPM 2.0, and GSC 1.2.
There are several hundred HK survey stars in our catalog which presently lack measured photometry. For these stars we have estimated the $`V`$ magnitudes and $`BV`$ colors in the following manner. An estimate of the $`B`$ magnitude is obtained from GSC 1.2, where available. An approximate de-reddened $`(BV)_o`$ color is obtained from available medium resolution spectroscopy, based on the calibration of the Balmer line index $`HP2`$ described by Beers et al. (1999). This is then reddened using the listed estimated to reddening in the direction toward the star, and the resulting $`BV`$ is obtained. This color is then used to convert the GSC $`B`$ magnitude estimate to an approximate $`V`$ magnitude. Given the multiple approximations in these procedures, these estimates should be regarded with appropriate caution – these stars are indicated in Table 1A with parentheses around the $`V`$ and $`BV`$ magnitude estimates. In a few cases, apparent magnitudes for HK survey stars that were not found in the GSC are estimated from the intensity of the spectrum on the original prism plate, and are indicated as such with brackets around the $`V`$ and $`BV`$ magnitudes in the table.
For a number of stars, $`V`$ magnitudes and $`by`$ colors are available from Schuster et al. (1996), Schuster et al. (2000), or Anthony-Twarog et al. (2000). For these stars, we have listed an approximate $`BV`$ color obtained by adopting the transformation $`BV=1.35(by)`$, and indicate these cases with parentheses about the approximate $`BV`$ color in . For RR Lyrae variables, the intensity-averaged $`V`$ magnitudes, where available, are taken from either Layden (1994), Layden et al. (1996), or Fernley et al. (1998). The adopted photometric quantities are reported in columns (10) and (11) of Table 1A; the source of the photometry is given in column (12); a reddening estimate (as described below) is listed in column (13).
#### 3.3.2 Estimates of Reddening
Our primary source of total reddening in the directions of the sample stars is the new map of Schlegel, Finkbeiner, & Davis (1998, hereafter SFD) based on COBE and IRAS measurements of dust emission in the Galaxy. The SFD map is superior to the widely used map of Burstein & Heiles (1982, hereafter BH) in its spatial resolution and its zero point (see also Gould & Popowski 1998 for arguments in favor of the SFD map). Note, however, that Arce & Goodman (1999) caution that the SFD map overestimates the reddening values when the color excess $`E_{BV}`$ is more than about 0.15 mags, where SFD’s calibration from dust column density to reddening may no longer be accurate. For the small number of stars where this large reddening applies, we have adopted estimates from the BH map, where available. The few program stars in such highly reddened directions are generally located at small distances from the Sun, so that the effects of correction for reddening on distance estimates are small.
We adopt a $`V`$-band absorption $`A_V=3.1E_{BV}`$ and assume that the dust layer has a scale height $`h=125`$ pc. The reddening to a given star at distance $`D`$ is reduced compared to the total reddening by a factor $`\mathrm{exp}[|D\mathrm{sin}b|/h]`$, where $`b`$ is the galactic latitude. An iteration procedure is employed to obtain consistent values of $`E_{BV}`$ and $`D`$. The final adopted reddening is listed in column (13) of Table 1A.
#### 3.3.3 Classification of Stellar Types
Column (14) of Table 1A lists the classification of each stellar type. We follow the coding of BSL – D: main-sequence dwarf star; A: main-sequence A-type star; TO: main-sequence turnoff star; SG: subgiant star; G: giant star; AGB: asymptotic giant branch star; FHB: field horizontal-branch star; V: variable stars other than the RR Lyrae type. For the RR Lyraes, we use a code ’RRV’ to distinguish them from other types of variable stars. Most of the classifications are taken from the BSL catalog, with the exception of a few stars where we have adopted the classifications from Norris, Bessel, & Pickles (1995) (HD 13359: G; HD 17233: SG; HD 219221: D; HD 217808: G; CPD-62 394: G) and Anthony-Twarog & Twarog (1994) (HD 161770: TO; BD-01 1792: G; HD 17072: AGB; HD 41667: G). These revised classifications provide photometric distance estimates which are more consistent with the Hipparcos parallaxes, as described below. Revised classifications are also provided for a number of the HK survey stars, based on more recent observations.
There still remains some uncertainty in the classification of a number of stars with relatively blue colors \[$`0.3(BV)_o0.5`$\]. If the corresponding stars have $`(UB)_o`$ color data from the HK survey, the stars are classified as FHB if $`(UB)_o>0.1`$, or TO if $`(UB)_o0.1`$ (Wilhelm, Beers, & Gray 1999). This crude discrimination scheme is adequate for stars with abundances $`[\mathrm{Fe}/\mathrm{H}]<1.5`$. For more metal-rich stars, we also consider the alternative classification as possible, and assign a classification of FHB/TO.
#### 3.3.4 Templates for Distance Estimation
RR Lyrae distance scales have been reconsidered extensively subsequent to the release of the Hipparcos Catalogue, based on a variety of methods including main-sequence fitting, the Baade-Wesselink method, and statistical parallax. Combining all of the results, Chaboyer (1999) proposed the following $`M_V`$ vs. $`[\mathrm{Fe}/\mathrm{H}]`$ relation:
$$M_V(RR)=0.23([\mathrm{Fe}/\mathrm{H}]+1.6)+0.56,$$
(1)
where $`M_V(RR)`$ denotes the absolute magnitude of RR Lyrae stars. We adopt this relation to estimate the distances of our RR Lyrae sample.
For non-variable stars, we first calibrate fiducial points in the $`M_V`$ vs $`(BV)_o`$ relation, using published data from various Galactic globular clusters and the Hyades and Pleiades open clusters to cover large ranges of colors and metal abundances. The globular clusters we use include M92 ($`[\mathrm{Fe}/\mathrm{H}]=2.24`$), M15 ($`2.15`$), M3 ($`1.66`$), M13 ($`1.65`$), NGC 6752 ($`1.54`$), NGC 362 ($`1.27`$), and 47 Tuc ($`0.71`$) (Sandage 1970; Buonanno, Corsi, & Fusi Pecci 1989; Durrell & Harris 1993; Penny & Dickens 1986; Harris 1982; Hesser et al. 1987), where metal abundances and extinctions are taken from the compilation of Chaboyer, Demarque, & Sarajedini (1996). The absolute magnitude of the horizontal branch in each cluster is scaled according to Eqn. (1). For the Hyades cluster ($`[\mathrm{Fe}/\mathrm{H}]=+0.12`$), we adopt the work of Perryman et al. (1998), which utilized the Hipparcos parallaxes to determine the distances to the cluster stars. Using their Figure 21, we read off the fiducial $`M_V`$ vs $`(BV)_o`$ relation for the stars which are not classified as double or multiple in the Hipparcos Catalogue. For the Pleiades cluster ($`[\mathrm{Fe}/\mathrm{H}]=+0.11`$, essentially identical to the Hyades abundance), we adopt the $`M_V`$ vs $`(BV)_o`$ relation published by Mermilliod (1981).
Next, polynomial fits to the $`M_V`$ vs $`(BV)_o`$ relation are made for each type of star; the fitting formula and the coefficients of the fits are provided in Table 2. For G, SG, and AGB stars with abundances $`[\mathrm{Fe}/\mathrm{H}]<2.5`$, we use the tabulated relation for $`[\mathrm{Fe}/\mathrm{H}]=2.5`$. For D stars we employ an interpolation/extrapolation procedure among the tabulated relations. Figure 2a demonstrates the adopted $`M_V`$ vs $`(BV)_o`$ relations of G, SG, and AGB stars at metal abundances appropriate for M92$`+`$M15 ($`[\mathrm{Fe}/\mathrm{H}]=2.20`$), M13 ($`1.65`$) and 47 Tuc ($`0.71`$), whereas Figure 2b shows those of D stars for M92$`+`$M15, 47 Tuc, and Hyades. Unfortunately, there remains some additional uncertainty for TO stars, due to their positions at the “curved” part of the adopted $`M_V`$ vs $`(BV)_o`$ relation. We employ the following approximate method. For a given $`(BV)_o`$ color, $`M_V`$’s assuming both a D and SG classification are first estimated. If $`M_V(\mathrm{D})`$ is brighter than $`M_V(\mathrm{SG})`$, $`M_V(\mathrm{D})`$ is adopted as $`M_V(\mathrm{TO})`$. If this is not the case, the intermediate $`M_V`$ which locates between $`M_V(\mathrm{D})`$ and $`M_V(\mathrm{SG})`$ is calibrated using the formula given in Table 2. This formula is designed to reproduce the turnoff point from neighboring points of the isochrones used in Beers et al. (1999).
In order to estimate distances for the hotter main-sequence gravity A-type stars (i.e., those with colors $`0.10(BV)_o0.34`$) we apply a small vertical shift to the $`M_V`$ vs $`(BV)_o`$ relation for the Pleiades so that it matches that of the Hyades at $`(BV)_o=0.34`$. To account for the expected decline in luminosity with decreasing abundance, the adjusted Pleiades line is corrected by applying shifts of the same size as those derived for the cooler stars on the main sequence of the metal-poor globulars as compared to the Hyades. We note that this is only one of several choices we could have made, and it may not be the optimal one, but it should serve for the present purpose.
Taking all of the above information into account, distances are then estimated by means of an iteration procedure which includes modification of the interstellar reddening. These distances, $`D_{pho}`$, are indicated in column (2) of Table 1B. In some cases there was no reliable photometry or classification available, and, occasionally, our iteration scheme failed to converge. In these circumstances we made use of the distances provided in the original BSL catalog (who, in such cases, adopted the distance estimates of the authors who provided the abundance and radial velocity information), listed as $`D_{BSL}`$ in column (5) of Table 1B. For the stars with ambiguous luminosity classifications, we list two alternative distances in column (2).
The small number of stars with uncertain type classifications (either due to the FHB/TO ambiguity noted above, or for the TO stars that appear too blue for their adopted type designation when compared with an old isochrone) are noted as such in column 12 of Table 1B, and are ultimately excluded from further analyses.
#### 3.3.5 Comparison of Photometric and Astrometric Distance Scales
There are 508 stars in our present catalog which were observed with the Hipparcos astrometric satellite (ESA 1997). To assess the reliability of our photometric distances, $`D_{pho}`$, obtained above, we plot in Figure 3 the Hipparcos trigonometric parallaxes, $`\pi _{HIP}`$, versus our photometric parallaxes, $`\pi _{pho}=1/D_{pho}`$, for 508 stars with available photometric distances. Plotting parallaxes, rather than distances, in Figure 3 has several advantages. All the trigonometric parallaxes, even small or negative values, can be used in an unbiased way (see the discussion on pp. 27–28 of Arenou & Luri 1999). Systematic errors in the data can be detected; an incorrect photometric distance scale would be seen as a slope not equal to unity; a zero-point error in the trigonometric parallaxes as a non-zero intercept. For clarity, we have not plotted error bars in Figure 3; the mean error of $`\pi _{HIP}`$ is $`1.6`$ mas; the mean error of $`\pi _{pho}`$ is $`20\%`$ of $`\pi _{pho}`$.
The general trend of the data in Figure 3 seems to fall slightly below the 45-degree line, indicating that $`\pi _{pho}`$ may be $`20\%`$ too large for the nearer stars in the BSL sample. This largely reflects the fact that these nearby stars have the largest errors in $`\pi _{pho}`$. The bulk of the data, especially the more distant stars, with smaller $`\pi _{pho}`$ errors, falls closer to the 45-degree line. Thus we conclude from Figure 3 that there is not any significant systematic error in our photometric distance scale. It is also worth noting that Figure 3 shows very few outliers, i.e. cases where $`\pi _{pho}`$ and $`\pi _{HIP}`$ badly disagree. The large spread in $`\pi _{HIP}`$ near $`\pi _{pho}=0`$ is due to the large errors in $`\pi _{HIP}`$ (up to 7 mas) for these faint, distant stars.
For the 424 stars with non-negative Hipparcos parallaxes, Table 1B lists in column (3) the Hipparcos distance estimates $`D_{HIP}=1/\pi _{HIP}`$ and in column (4) their relative precisions $`\sigma _{\pi _{HIP}}/\pi _{HIP}`$. Because the photometric distances for the nearest stars in our catalog have lower precision than the trigonometric parallaxes, we chose to adopt the Hipparcos distances for 44 stars with $`\sigma _{\pi _{HIP}}/\pi _{HIP}<0.12`$. The final adopted distance estimates for our program stars are listed as $`D_{adopt}`$ in column (6) of Table 1B.
## 4 Proper Motions
### 4.1 Sources of Proper Motions
#### 4.1.1 The Hipparcos Catalogue
The Hipparcos Catalogue is the primary result of the first space mission dedicated to astrometry (ESA 1997), and it provides high accuracy proper motions for 118,218 stars with $`V<12.5`$ mag covering the entire sky. In the present sample, 508 stars are included in this catalog with average accuracies of 1.61 mas yr<sup>-1</sup> in $`\mu _\alpha ^{}`$ ($`=\mu _\alpha \mathrm{cos}\delta `$) and 1.29 mas yr<sup>-1</sup> in $`\mu _\delta `$. Kinematical analyses for a subset of these Hipparcos stars are reported in Chiba & Yoshii (1998) and Chiba, Yoshii, & Beers (1999).
#### 4.1.2 The SPM Catalog 2.0
The Yale/San Juan Southern Proper Motion (SPM) Catalog 2.0 provides positions, absolute proper motions and photographic $`B,V`$ magnitudes for 321,608 objects with $`5<V<18.5`$, mainly in the declination zone from $`37^{}`$ to $`27^{}`$, except within the Galactic zone of avoidance. In this declination zone an effort was made to measure all HK survey stars from a database provided by T. Beers prior to undertaking measurements of the plates. A total of 505 of our program stars have proper motions extracted from the SPM 2.0 Catalog, using a 10″ search radius centered on the catalogued positions, and satisfying the requirement that the apparent magnitudes be within 1 mag in $`B`$ or $`V`$ of those adopted in Table 1A. The description of the catalog properties can be find in Platais et al. (1998), however, this paper addresses only the separate catalog at the South Galactic Pole. For additional details of the SPM Catalog 2.0, the reader is referred to the World Wide Web at URL http://www.astro.yale.edu/astrom/. The average accuracy of SPM proper motions for the subset of 508 stars is 2.9 mas yr<sup>-1</sup>.
#### 4.1.3 The Lick NPM1 Catalog
The Lick Northern Proper Motion (NPM) program (Hanson 1997) is a photographic survey measuring absolute proper motions, on an inertial system defined by 50,000 faint galaxies, for over 300,000 stars with $`8<B<18`$, covering the northern two-thirds of the sky ($`\delta >23^{}`$). Part I of the NPM program (“NPM1”), outside the disk of the Milky Way ($`|b|>10`$ deg), was completed in 1993; the Lick NPM1 Catalog (Klemola, Hanson, and Jones 1993) contains 148,940 stars, with proper motions accurate to 5 mas/yr, and positions accurate to 0$`\stackrel{}{\mathrm{.}}`$2 or better. Part II of the NPM program (“NPM2”), in the Milky Way sky, will be completed in roughly three years.
We searched for program stars within a $`\pm 10\mathrm{}\times 10\mathrm{}`$ error box, centered on the NPM1 catalog positions (updated to epoch 2000.0 using the NPM1 proper motions). This procedure found matches for 241 stars; 95 % of these comprise a tight core with RMS position difference 0$`\stackrel{}{\mathrm{.}}`$5 in each coordinate, and can be considered unambiguous matches. The remaining 5 % can be considered highly probable matches; agreement between the cataloged and NPM1 apparent magnitudes (better than $`\pm 1`$ mag) confirms these.
NPM1 proper motions for an additional 58 RR Lyrae stars, from the list of Layden et al. (1996), were added to our catalogue, making a total of 299 stars with NPM1 proper motion data.
#### 4.1.4 The STARNET Catalog
The proper motions of the STARNET (STN) catalog were determined using the Astrographic Catalogue (AC) as the first epoch positional source, and the astrometrically upgraded Guide Star Catalog (GSC 1.2) as the second epoch source. The AC has been made available in machine readable form by Nesterov et al. (1990). Its reduction to the FK5 system was performed by S. Röser at the Astronomisches Rechen-Institut (ARI). The astrometrical upgrading of the Guide Star Catalog is described in detail by Morrison et al. (1996). GSC 1.2 is presently available at the Astronomical Data Center. The result of combining both catalogs is the STARNET catalog of proper motions (Röser 1996). It contains about 4.3 million stars with an average density of 100 stars per square degree. The median magnitudes of the stars are $`B=12.0`$ mag in the southern hemisphere and $`V=11.5`$ in the northern hemisphere, respectively, with the faintest stars reaching $`B,V13.5`$. The present rms accuracy of STARNET positions is 0$`\stackrel{}{\mathrm{.}}`$3. The accuracy of the proper motions, which have been determined simply from combining the estimated positional errors at both epochs, are about 5 mas yr<sup>-1</sup>. The catalog is reduced to the FK5 reference system and thus contains some small systematic differences with respect to ICRS (see Vol. 3, Chapter 19 of the Hipparcos and Tycho Catalogues, ESA 1997), in addition to some residual offsets still present in the GSC 1.2 positions. STARNET is as yet unpublished, but is accessible upon request at ARI.
We performed searches for our program stars in STARNET using windows of 40″ diameter centered on the sample stars, and obtained 774 identifications in STARNET, with an average accuracy of the proper motions of 4.7 mas yr<sup>-1</sup> in each component, respectively. Since the STARNET catalogue may contain mis-identifications of stars from first epoch to second, proper motions larger than 200 mas yr<sup>-1</sup> in one component may be suspect. We therefore inspected such proper motions individually. As it turned out all but one of these stars were listed in the Hipparcos Catalogue; the remaining star could be found in the PPM catalog (Röser and Bastian 1991; Bastian et al. 1993). The agreement between the STARNET motions and the other measurements is usually quite good, with deviations of the order of the expected accuracy of the proper motions.
#### 4.1.5 The ACT Reference Catalog
This recently completed catalog by the USNO contains positions and proper motions for 988,758 stars with $`V<1111.5`$ mag covering the entire sky (Urban, Corbin, & Wycoff 1998a). Accurate proper motions were calculated by combining the positions from the Tycho Catalogue (ESA 1997) with those from new reductions of the Astrographic Catalogue, referred to as AC 2000 (Urban et al. 1998b). We performed a search of our sample stars within a $`10^{\prime \prime }`$ radius centered at the positions of the cataloged stars, and obtained unambiguous matches for 525 stars. Average accuracies in proper motions are 3.1 mas yr<sup>-1</sup> and 2.9 mas yr<sup>-1</sup> in $`\mu _\alpha ^{}`$ and $`\mu _\delta `$, respectively.
### 4.2 Comparison of Proper Motions for Stars with Multiple Measures
To assess the quality of the adopted proper motions, Figure 4 shows the difference between the Hipparcos Catalogue and other ground-based measurements for each proper-motion component, for the stars in common between the catalogs. While the Hipparcos errors are smallest in general, the figures show that there is an overall agreement with proper motion measurements from the additional sources, without any clearly systematic differences as compared with Hipparcos. It is worth noting that for several stars, the proper motion accuracy in the SPM 2.0 Catalog is better than that in the Hipparcos Catalogue.
For the stars which have been independently measured in two or more catalogs it is possible, by combining all measurements, to reduce the statistical errors as well as minimize any small remaining systematic errors in the individual catalogs, as was argued in Martin & Morrison (1998). For these stars, we estimate average proper motions, $`<\mu >`$, and their errors, $`<\sigma _\mu >`$, weighted by the inverse variances, as
$`<\mu >`$ $`=`$ $`({\displaystyle \underset{i=1}{\overset{n}{}}}\mu _i/\sigma _{\mu _i}^2)/({\displaystyle \underset{i=1}{\overset{n}{}}}1/\sigma _{\mu _i}^2),`$ (2)
$`<\sigma _\mu >`$ $`=`$ $`({\displaystyle \underset{i=1}{\overset{n}{}}}1/\sigma _{\mu _i}^2)^{1/2},`$ (3)
where $`n`$ denotes the number of measurements. In this manner, the proper motions for 691 stars in the present sample are improved. Our adopted proper motions are listed in columns (7) and (8) of Table 1B; the associated errors are listed in columns (9) and (10). The sources of the proper motions are listed in column (11).
We emphasize the large number of the stars having proper motion data in the present catalog: $`N=1258`$ (1214 with $`[\mathrm{Fe}/\mathrm{H}]0.6`$). This is the largest such dataset among any non-kinematically selected samples, and is substantially larger than a number of recent samples selected on the basis of their high proper motions alone (Ryan & Norris 1991a: $`N=727`$ with $`[\mathrm{Fe}/\mathrm{H}]0.6`$; Carney et al. 1994: $`N=689`$ with $`[\mathrm{Fe}/\mathrm{H}]0.6`$). Furthermore, the quality of the proper motion data has been improved in a meaningful manner, compared with that of earlier studies, due to the continuous improvement of astrometric observations from the ground as well as from space. As a consequence the current database allows us to elucidate a great deal of information concerning the space motions of the metal-poor populations of the Milky Way with a lowest-to-date observational error.
## 5 The Revised Catalog and its Characteristics
### 5.1 The Catalog
Tables 1A and 1B present the final revised catalog for 2106 stars. With a constraint of $`[\mathrm{Fe}/\mathrm{H}]0.6`$, the present catalog contains 2041 stars, including 356 RR Lyrae variables (168 of which are newly added as compared to the BSL catalog of Paper I), 38 variables of other types, and 1647 non-variables. Proper motions are available for 1214 stars with $`[\mathrm{Fe}/\mathrm{H}]0.6`$.
### 5.2 Characteristics of the Revised Catalog
Figure 5a illustrates the distribution of apparent magnitudes for stars in the present catalog. The relatively small number of stars with magnitudes $`V>15`$, which, for the giants of metal-poor populations, are located at distances in excess of 10 kpc from the Sun, clearly points to the need for future work. The distribution of $`(BV)_o`$ colors for our catalog stars is shown in Figure 5b. The dominance of D and MS stars is clear, pointing out the need for additional studies of redder G and AGB stars of the thick disk and halo, as well as the need for the inclusion of blue FHB stars.
Figure 5c shows the distribution of adopted distances for stars in the present catalog. Stars populating the inner few kpc from the Sun are well represented, with a rather sharp decline beyond 3 kpc from the Sun. The distribution of revised stellar abundances shown in Figure 5d underscores the fact that we have included a large number of stars with abundances below the peak metallicity of the halo metallicity distribution function (\[Fe/H\] $`1.6`$, Laird et al. 1988; Ryan & Norris 1991b). Future work which supplements this sample with additional stars in the intermediate abundance range $`1.5[\mathrm{Fe}/\mathrm{H}]0.6`$ is clearly needed.
The dashed histograms in Figures 5a-d indicate the distribution of the above observables for the subset of stars with available proper motions. As is evident from inspection of Figure 5a, the great majority of the brighter stars have measured proper motions, but the completeness falls from about 50% to roughly 30% for the stars in the magnitude range $`12V16`$. Figure 5c shows that the fraction of stars with available proper motions is roughly a constant $``$ 60% out to distances on the order of 4–5 kpc, with no strong dependence on distance.
In Figures 6a and 6b we show the distribution of adopted radial velocities (for all catalog stars with \[Fe/H\] $`0.6`$) and measured proper motions (where available) as a function of the stellar abundances. The largest radial velocities and proper motions are observed near the center of the abundance distribution, but this is presumably a result of the (presently unavoidable) fact that the relatively rare large observed motions will occur with greatest frequency among the stars which are most commonly represented in the catalog.
The referee has raised the question of whether we have “imported” a kinematic bias into the sample of stars with measured proper motions by the addition of information from the sources we have drawn on. This is demonstrably not the case. None of the five catalogs we have used in our assemblage of proper motion information have lower limits on their proper motions (as do, for example, many older catalogs such as those from Giclas and Luyten). The stars in each of these catalogs are chosen in advance of measuring their proper motions, hence a given star is not rejected if its measured proper motion turns out to be small, even if it is consistent with zero. Furthermore, since the proper motions we have included are measured in two independent coordinates, each of which can be individually negative, zero, or positive, errors in the measured proper motions do not result in a net bias away from zero. This bias does exist in many of the older catalogs (where selection was based on the always positive length of a tangential motion vector), but it does not apply to the catalogs we have used for our proper motions.
Since the above biases could lead, in effect, to a distance-related bias in the resulting kinematics, we have performed a simple diagnostic to demonstrate that this is not the case for our sample of stars (other tests are described in Paper III). Figure 6c shows a plot of the estimated tangential velocity of our sample stars as a function of distance from the Sun. Superposed on this figure are lines of constant proper motion, corresponding to 10 mas/yr and 1 mas/yr, respectively. There does not appear to be any obvious changes in the distribution of predicted tangential velocity ($`V_t=4.74\mu D`$, where $`D`$ is the distance in kpc, and $`\mu `$ is in mas/yr) that correlate with increasing distance from the Sun, as might be expected if there were any “hidden” kinematic bias in the sample of stars with measured proper motions.
The spatial distribution of the sample in $`(R,Z)`$ plane is shown in Figure 7, where filled circles and crosses indicate the stars with and without available proper motions, respectively. Below the Galactic plane there are many remote stars with available proper motions, mainly provided by SPM 2.0.
### 5.3 Full Space Motions and Orbital Parameters
We now obtain estimates of the full space motions for the subsample of stars for which proper motions are available, and derive parameters of their orbital motions in a given gravitational potential. The results are summarized in Table 3. Note that two of the stars in this subsample were judged not to have sufficiently accurate distance estimates to derive space motions, V\* KR Vir and BPS CS 29512-0032, so these stars are not included in Table 3. Column (1) lists the star name. Column (2) recalls the adopted metallicity from Table 1A. Columns (3) and (4) list the positions of the stars in the meridional plane $`(R,Z)`$, adopting $`R_{}=8.5`$ kpc as the Galactocentric distance for the Sun. Columns (5)-(7) list the three-dimensional velocities $`U`$, $`V`$, and $`W`$, directed to the Galactic anticenter, the rotational direction, and the north pole, respectively. These velocity components are corrected for the solar motion $`(U_{},V_{},W_{})=(9,12,7)`$ km s<sup>-1</sup> with respect to the local standard of rest (LSR) (Mihalas & Binney 1981). Figure 8 is a plot of the $`U,V,W`$ velocity components for the sample of stars with complete velocity information and $`[\mathrm{Fe}/\mathrm{H}]0.6`$, as a function of \[Fe/H\].
Columns (8) and (9) of Table (3) list the velocity components $`(V_R,V_\varphi )`$ in the cylindrical rest frame $`(R,\varphi )`$, respectively, on the assumption that the rotational speed of the LSR around the Galactic center is $`V_{LSR}=220`$ km s<sup>-1</sup>. Note that for the stars for which accurate distances are not available, columns (2)-(9) contain no data.
To estimate the orbital parameters for these stars, we adopt the analytic Stäckel-type mass model developed by Sommer-Larsen & Zhen (1990), which consists of a flattened, oblate disk and a nearly spherical massive halo. This model reproduces a flat rotation curve beyond $`R=4`$ kpc and the local mass density at $`R_{}`$ in a reasonable manner. Its analytic form has the great advantage of maintaining clarity in further analyses (see, e.g., Chiba & Yoshii 1998). Columns (10) and (11) list the estimated apogalactic distances, $`R_{ap}`$, and the estimated perigalactic distances, $`R_{pr}`$ along the Galactic plane, respectively. Column (12) lists the maximum distance above (or below) the plane, $`Z_{max}`$, explored by each star in the course of its orbital motion. In column (13) we list the characteristic eccentricities of the orbits, defined as $`e=(r_{ap}r_{pr})/(r_{ap}+r_{pr})`$, where $`r_{ap}`$ and $`r_{pr}`$ stand for the apogalactic and perigalactic distances from the Galactic center, respectively. We note that the current mass model of the Galaxy fails to gravitationally bind the 16 stars which lack the data of orbital parameters in columns (10)-(13). These stars probably correspond to those in the “error tail” of overestimated photometric distances, thereby having overestimated transverse velocities.
In the accompanying analysis paper (Chiba & Beers 2000; Paper III), we make use of the kinematic information in Table 3, and the radial velocity and distance information for the stars presently without available proper motions, to investigate a number of questions concerning the nature of the kinematics of the Galaxy.
### 5.4 Future Samples
Exploration of the kinematic properties of the Galaxy is very much a work in progress. New and extensive samples of several thousand additional non-kinematically selected metal-deficient stars with available velocity, distance, and abundance measurements are expected to become available within the next year (Beers et al. 2000; Cayrel et al. 2000; Norris et al. 2000; Rebolo et al. 2000). Roughly 30% of these stars already have available proper motions; completing the search for HK survey stars with proper motions awaits the extension of the SPM catalog to other areas of the southern Galactic hemisphere.
Samples of intermediate distance ($`2d15`$ kpc) FHB and other A-type stars (e.g., Wilhelm et al. 1999), supplemented with observations of more distant metal-poor dwarfs, giants, and FHB and A-type stars from the Hamburg/ESO survey ($`10d25`$ kpc; Christlieb 1999), and from other sources (e.g., the Sloan Digital Sky Survey, Pier 1999), will provide useful extensions of our catalog.
A project to refine systemic radial velocity measurements for the RR Lyraes with poorly determined values is already underway. Additional $`UBV`$ and/or Strömgren photometry is required in order to resolve ambiguities in stellar classifications which remain in the catalog, and to obtain more precise estimates of distances, especially for the stars with available proper motions.
We are grateful to the referee, Bruce Carney, for a careful reading of this manuscript, and for a number of thoughtful comments. TCB acknowledges support for this work from grant AST 95-29454 from the National Science Foundation. MC and YY acknowledge partial support from Grants-in-Aid for Scientific Research (09640328) and COE Research (07CE2002) of the Ministry of Education, Science, Sports and Culture of Japan. MC thanks Hideyuki Saio for his help in distance estimates for the current catalog stars. The SPM program is supported by grants from the NSF to Yale University and the Yale Southern Observatory, Inc. IP thanks T. Girard, V. Kozhurina-Platais, and W. van Altena for their expertise and contribution to the SPM program. RBH thanks the National Science Foundation for its long-term support of the Lick Northern Proper Motion program. Current work on the NPM program is supported by NSF grant AST 95-30632. RBH thanks A. Klemola for his help in providing identifications for the Lick NPM1 Catalog stars. BF thanks to S. Röser and S. Frink for help with the STARNET catalog. SR acknowledges partial support for this work from grant 200068/95-4 CNPq, Brazil, and from the Brazilian Agency FAPESP. This work made use of the SIMBAD database, operated at CDS, Strasbourg, France.
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# On canonically derived families of surfaces of general type over curves
## 1. Preliminaries
### 1.1 Convention
Let $`X`$ be a normal projective variety of dimension $`d`$. We denote by $`\text{Div}(X)`$ the group of Weil divisors on $`X`$. An element $`D\text{Div}(X)`$ is called a $``$-divisor. A $``$-divisor $`D`$ is said to be $``$-Cartier if $`mD`$ is a Cartier divisor for some positive integer $`m`$. For a $``$-Cartier divisor $`D`$ and an irreducible curve $`CX`$, we can define the intersection number $`DC`$ in a natural way. A $``$-Cartier divisor $`D`$ is called nef (namely numerically effective) if $`DC0`$ for any effective curve $`CX`$. A nef divisor $`D`$ is called big if $`D^d>0`$. We say that $`X`$ is $``$-factorial if every Weil divisor on $`X`$ is $``$-Cartier. For a Weil divisor $`D`$ on $`X`$, write $`𝒪_X(D)`$ as the corresponding reflexive sheaf. Denote by $`K_X`$ a canonical divisor of $`X`$, which is a Weil divisor. $`X`$ is called minimal if $`K_X`$ is a nef $``$-Cartier divisor. $`X`$ is said to be of general type if $`\text{kod}(X)=\text{dim}(X)`$. For a positive integer $`m`$, we set $`\omega _X^{[m]}:=𝒪_X(mK_X).`$ We use \[R\] as a nice reference for the definition of canonical, terminal singularities. According to both \[KMM\] and \[K-M\], any given smooth projective 3-fold $`Y`$ of general type has a minimal model $`X`$ which has only $``$-factorial terminal singularities.
### 1.2 Vanishing theorems
Let $`D=a_iD_i`$ be a $``$-divisor on $`X`$ where the $`D_i^{}s`$ are distinct prime divisors and $`a_i`$. We define
the round-down $`\mathrm{}D\mathrm{}:=\mathrm{}a_i\mathrm{}D_i`$, where $`\mathrm{}a_i\mathrm{}`$ is the integral part of $`a_i`$.
the round-up $`\mathrm{}D\mathrm{}:=\mathrm{}D\mathrm{}`$.
the fractional part $`\{D\}:=\mathrm{}D\mathrm{}D\mathrm{}\mathrm{}`$.
Throughout this paper, we will use the Kawamata-Viehweg vanishing theorem (\[Ka1\], \[KMM\] and \[V\]) in the following forms.
###### Theorem 1.1
Let $`X`$ be a smooth complete variety, $`D\text{Div}(X)`$. Assume the following two conditions:
(i) $`D`$ is nef and big;
(ii) the fractional part of $`D`$ has supports with only normal crossings.
Then $`H^i(X,𝒪_X(K_X+\mathrm{}D\mathrm{}))=0`$ for all $`i>0`$.
###### Theorem 1.2
Let $`X`$ be a normal projective variety with only canonical singularities. Let $`D`$ be a $``$-Cartier Weil divisor such that $`D`$ is nef and big. Then $`H^i(X,𝒪_X(K_X+D))=0`$ for all $`i>0`$.
### 1.3 Semi-positivity
Let $`C`$ be a smooth projective curve and $``$ be a vector bundle on $`C`$. We call
$$\mu ():=\frac{\mathrm{deg}()}{\text{rk}()}$$
the slope of $``$. According to \[H-N\], there is the Harder-Narasimhan filtration
$$0=_0_1\mathrm{}_{n1}_n=,$$
where the quotient $`_i/_{i1}`$ is a semistable vector bundle and
$$\mu (_i/_{i1})>\mu (_{i+1}/_i)$$
for all $`i`$. We define
$$\mu _{min}():=\mu (/_{n1}),$$
which is called the minimal slope of $``$.
###### Definition Definition 1.3
The vector bundle $``$ is said to be semi-positive if $`\mu _{min}()0`$.
According to \[Ka2\], \[Ko2\], \[N\] and \[O\], we have the following
###### Fact 1.4
Let $`X`$ be a smooth projective 3-fold and $`f:XC`$ be a proper morphism with connected fibers onto a smooth projective curve $`C`$. Then both $`f_{}\omega _{X/C}^m`$ and $`R^if_{}\omega _{X/C}`$ are semi-positive vector bundles on $`C`$ for all $`m>0`$ and $`i>0`$. In particular,
$$\mathrm{deg}f_{}\omega _{X/C}^m0\text{and}\mathrm{deg}R^if_{}\omega _{X/C}0.$$
### 1.4 Basic formulae
Let $`X`$ be a smooth projective 3-fold and $`f:XC`$ be a fibration onto the smooth projective curve $`C`$. Denote $`b:=g(C)`$. From the spectral sequence
$$E_2^{p,q}:=H^p(C,R^qf_{}\omega _X)E^n:=H^n(X,\omega _X),$$
one obtains the following formulae
$$q(X):=h^1(X,𝒪_X)=b+h^1(C,R^1f_{}\omega _X)$$
$`1.1`$
$$h^2(𝒪_X)=h^1(C,f_{}\omega _X)+h^0(C,R^1f_{}\omega _X).$$
$`1.2`$
## 2. Lemmas
###### Lemma 2.1
Let $`X`$ be a smooth projective 3-fold of general type. $`m2`$ is an integer. Suppose that $`|mK_X|`$ is composed of a pencil of surfaces. Keep the same notations as in the first page of this paper. We have a derived fibration $`f_m:X^{}C`$. Then $`C`$ is either an elliptic curve or $`^1`$.
###### Demonstration Proof
Suppose $`b>0`$. Then $`\varphi _m`$ is a morphism. We have a derived fibration
$$f:=f_m:XC.$$
Let $`_0`$ be the saturated sub-bundle of $`f_{}\omega _X^m`$ which is generated by $`H^0(C,f_{}\omega _X^m)`$. Since $`|mK_X|`$ is composed of a pencil and $`\varphi _m`$ factors through $`f`$, $`_0`$ should be a line bundle on $`C`$. Denote $`:=f_{}\omega _X^m`$. Then we have the following extension
$$0_0_10$$
and the exact sequence
$$f_{}\omega _{X/C}^m_1\omega _C^m0.$$
Note that $`r:=\text{rk}()=h^0(F,mK_F)2`$ because the general fiber $`F`$ is a smooth projective surface of general type. According to Fact 1.4, $`f_{}\omega _{X/C}^m`$ is semi-positive. Therefore $`\mathrm{deg}(_1\omega _C^m)0`$, i.e.
$$\mathrm{deg}_12m(r1)(b1).$$
We have
$`h^1(_0)`$ $`h^0(_1)\mathrm{deg}_1+(r1)(1b)`$ $`2.1`$
$`(2m1)(r1)(b1).`$
Suppose $`h^1(_0)>0`$. Since $`\mathrm{deg}_0>0`$, according to the Clifford’s theorem, we have
$$\mathrm{deg}_02h^0(_0)2h^0(_0)=P_m(X)2.$$
On the other hand, we have
$$h^1(_0)=h^0(_0)\mathrm{deg}(_0)+b1b1.$$
$`2.2`$
Thus, by (2.1) and (2.2), we get
$$b1(2m1)(r1)(b1).$$
The only possibility is $`b=1`$. When $`h^1(_0)=0`$, we also automatically have $`b=1`$ from (2.1). The proof is complete. ∎
###### Lemma 2.2
Let $``$ be a vector bundle of rank $`r`$ on a smooth projective curve $`C`$. Suppose $`\omega _C^1`$ is semi-positive. Then we have $`h^1(C,)r.`$
###### Demonstration Proof
Suppose there are $`r+1`$ independant sections
$$s_1,s_2,\mathrm{},s_{r+1}H^0(C,(\omega _C^1)^{})H^1(C,).$$
Denote $`^{}:=(\omega _C^1)^{}`$. For any point $`xC`$, the stalk $`_x^{}`$ is an $`𝒪_{C,x}`$-module of rank $`r`$. This means that $`s_{1,x}`$, $`s_{2,x}`$, $`\mathrm{}`$, $`s_{r+1,x}`$ are algebraically dependant in $`_x^{}`$. Thus there are $`r+1`$ nontrivial germs
$$f_1,f_2,\mathrm{},f_{r+1}𝒪_{C,x}$$
such that $`_{i=1}^{r+1}f_i(x)s_{i,x}(x)=0`$. Now set
$$s:=\underset{i=1}{\overset{r+1}{}}f_i(x)s_iH^0(C,(\omega _C^1)^{}).$$
$`s`$ is a non-zero section. Otherwise $`s_1`$, $`\mathrm{}`$, $`s_{r+1}`$ are dependant. Because $`s`$ vanishes at $`x`$, $`s`$ defines a line bundle $``$ which has positive degree. So $`\omega _C^1`$ has a quotient bundle with negative degree. This contradicts to the semi-positivity of $`\omega _C^1`$. The proof is completed. ∎
###### Corollary 2.3
Let $`f:XC`$ be a fibration from a smooth projective 3-fold $`X`$ onto a smooth projective curve $`C`$. Let $`F`$ be a general fibre of $`f`$ and set $`b:=g(C)`$. Then $`q(X)b+q(F)`$.
###### Demonstration Proof
This is a direct result from Fact 1.4, (1.1) and Lemma 2.2.
## 3. Proof of the main theorems
Since the behavior of pluricanonical maps is birationally invariant, we may suppose that $`X`$ is a normal projective minimal 3-fold of general type with only $``$-factorial terminal singularities. We make this assumption so as to utilize vanishing theorems. Now suppose that $`|mK_X|`$ is composed of a pencil of surfaces. We can use the same set up as in the first page of this paper. An extra point is that we can take the modification $`\pi _m:X^{}X`$ such that $`\pi _m^{}(mK_X)`$ has supports with only normal crossings. We keep the same notations. Then we get a derived fibration $`f_m:X^{}C`$. Denote by $`F`$ a general fiber of $`f_m`$ and by $`b`$ the genus of $`C`$. Sometimes we simply denote $`f_m`$ by $`f`$ and $`\pi _m`$ by $`\pi `$, respectively.
###### Proposition 3.1
Let $`X`$ be a normal projective minimal 3-fold of general type with only $``$-factorial terminal singularities. Suppose $`|mK_X|`$ is composed of a pencil of surfaces. If $`p_g(X)>0`$, then $`p_g(F)=1`$ under one of the following conditions.
(1) $`m=2`$, $`b=0`$ and $`P_m(X)4`$.
(2) $`m3`$, $`b=0`$ and $`P_m(X)3`$.
(3) $`m=2`$, $`b=1`$ and $`P_m(X)3`$.
(4) $`m3`$, $`b=1`$ and $`P_m(X)2`$.
###### Demonstration Proof
Because $`p_g(X^{})=p_g(X)>0`$, we can choose an effective divisor $`D_1|K_X^{}|`$. We write
$$K_X^{}=\pi ^{}(K_X)+a_iE_i,$$
where $`a_i^+`$, $`E_i`$ is an exceptional prime divisor for all $`i`$. We note that
$$\pi ^{}(K_X)=K_X^{}a_iE_i=D_1a_iE_i$$
is an effective $``$-divisor. So we have
$$\pi ^{}(K_X)=b_iG_i=G_v+G_h,$$
where $`b_i^+`$, $`G_i`$ is a prime divisor on $`X^{}`$ for all $`i`$, the support of $`G_v`$ is contained in fibers of $`f`$ and $`G_h`$ is the horizontal part of $`\pi ^{}(K_X)`$. Both $`G_v`$ and $`G_h`$ are effective $``$-divisors. So we have
$$(m1)\pi ^{}(K_X)=(m1)G_v+(m1)G_h$$
$$m\pi ^{}(K_X)=mG_v+mG_h.$$
Suppose $`M_m`$ is the movable part of $`|mK_X^{}|`$. Then $`M_m_{}mG_v`$ because the support of $`M_m`$ is vertical. Thus we can write
$$mG_v=M_m+G_v^{}=M_m+c_iG_i^{},$$
where $`c_i^+`$, $`G_i^{}`$ is a prime divisor on $`X^{}`$ for all $`i`$. Therefore
$$m\pi ^{}(K_X)=M_m+(G_v^{}+mG_h),$$
where $`G_v^{}+mG_h`$ is an effective $``$-divisor. We can suppose
$$M_m_{\text{lin}}\underset{i=1}{\overset{a}{}}F_i_{\text{num}}aF,$$
where
$$a=\{\begin{array}{cc}P_m(X)1,\hfill & \text{if}b=0,\hfill \\ P_m(X),\hfill & \text{if}b=1.\hfill \end{array}$$
Thus we have
$$\pi ^{}(K_X)_{\text{num}}\frac{a}{m}F+\frac{1}{m}G_v^{}+G_h$$
$$(m1)\pi ^{}(K_X)_{\text{num}}\frac{a(m1)}{m}F+\frac{m1}{m}G_v^{}+(m1)G_h.$$
Denote $`a^{}:=\frac{a(m1)}{m}`$. We can see that $`a^{}>1`$ under one of the assumptions within (1) through (4) of the proposition. So
$$(m1)\pi ^{}(K_X)F\frac{m1}{ma^{}}G_v^{}\frac{m1}{a^{}}G_h_{\text{num}}(m1)(1\frac{1}{a^{}})\pi ^{}(K_X)$$
is nef and big and its fractional part has supports with only normal crossings. According to the Kawamata-Viehweg vanishing theorem, we get
$$H^1(X^{},K_X^{}+\mathrm{}(m1)\pi ^{}(K_X)\frac{m1}{ma^{}}G_v^{}\frac{m1}{a^{}}G_h\mathrm{}F)=0.$$
Set $`G^{\prime \prime }:=\mathrm{}(m1)\pi ^{}(K_X)\frac{m1}{ma^{}}G_v^{}\frac{m1}{a^{}}G_h\mathrm{}`$. Then $`G^{\prime \prime }\mathrm{}(m1)\pi ^{}(K_X)\mathrm{}`$ and so
$$K_X^{}+G^{\prime \prime }K_X^{}+\mathrm{}(m1)\pi ^{}(K_X)\mathrm{}.$$
Therefore we have
$$\text{dim}\mathrm{\Phi }_{|K_X^{}+G^{\prime \prime }|}(F)=0$$
for a general fiber $`F`$. From the exact sequence
$$0𝒪_X^{}(K_X^{}+G^{\prime \prime }F)𝒪_X^{}(K_X^{}+G^{\prime \prime })𝒪_F(K_F+G^{\prime \prime }|_F)0,$$
we get the surjective map
$$H^0(X^{},K_X^{}+G^{\prime \prime })H^0(F,K_F+G^{\prime \prime }|_F).$$
This means
$$|K_X^{}+G^{\prime \prime }||_F=|K_F+G^{\prime \prime }|_F|.$$
Noting that
$$G^{\prime \prime }|_F=\mathrm{}(m1)\pi ^{}(K_X)\frac{m1}{ma^{}}G_v^{}\frac{m1}{a^{}}G_h\mathrm{}|_F=\mathrm{}(m1)(1\frac{1}{a^{}})G_h\mathrm{}|_F$$
is an effective divisor, we have $`h^0(F,K_F+G^{\prime \prime }|_F)2`$ whenever $`p_g(F)2`$. This would lead to $`\text{dim}\varphi _m(F)1`$, which is impossible. So we should have $`p_g(F)=1`$ because $`p_g(F)>0`$ under the assumption $`p_g(X)>0`$. ∎
###### Remark Remark 3.2
The assumption $`p_g(X)>0`$ in Proposition 3.1 is important. If $`p_g(X)=0`$, the above method is invalid because we don’t know whether $`G^{\prime \prime }|_F`$ is effective.
###### Theorem 3.3
Let $`f:XC`$ be a derived family of surfaces from the m-canonical pencil $`|mK_X|`$ of a smooth projective 3-fold $`X`$ of general type. Let $`F`$ be a general fiber of $`f`$ and denote $`b:=g(C)`$. Suppose $`m3`$. Then one of the following occurs:
(A1) $`m=5,6`$, $`p_g(X)=2`$, $`q(X)=0`$, $`b=0`$ and $`p_g(F)=q(F)=1`$.
(A2) $`m=4`$, $`2p_g(X)3`$, $`q(X)=0`$, $`b=0`$ and $`p_g(F)=1`$.
(A3) $`m=3`$, $`2p_g(X)5`$, $`q(X)1`$, $`b=0`$ and $`p_g(F)=1`$.
(B0) $`p_g(X)=1`$, $`q(X)b+1`$ and $`p_g(F)=1`$.
(B1) $`p_g(X)=P_2(X)=\mathrm{}=P_{m1}(X)=1`$, $`P_m(X)=2`$, $`\text{dim}\varphi _{m+1}(X)2`$, $`b=0`$ and $`p_g(F)2`$.
(B2) $`p_g(X)=P_2(X)=\mathrm{}=P_{m1}(X)=1`$, $`P_m(X)=P_{m+1}(X)=2`$, $`\text{dim}\varphi _{m+2}(X)2`$, $`b=0`$ and $`p_g(F)2`$.
(B3) $`p_g(X)=P_2(X)=\mathrm{}=P_{m2}(X)=1`$, $`P_{m1}(X)=P_m(X)=2`$, $`\text{dim}\varphi _{m+1}(X)2`$, $`b=0`$ and $`p_g(F)2`$.
(C0) $`p_g(X)=0`$, $`q(X)2`$ and $`q(F)2`$.
(C1) $`p_g(X)=0`$, $`\text{dim}\varphi _{m+1}(X)2`$ and $`q(F)2`$.
(C2) $`p_g(X)=0`$, $`P_m(X)=2`$, $`\text{dim}\varphi _{m+2}(X)2`$, $`b=1`$ and $`q(F)=2`$.
(C3) $`p_g(X)=0`$, $`P_m(X)=2`$, $`\text{dim}\varphi _{2m}(X)2`$, $`b=0`$ and $`q(F)3`$.
(C4) $`p_g(X)=0`$, $`P_m(X)=2`$, $`P_{2m}(X)=3`$, $`\text{dim}\varphi _{2m+1}(X)=3`$, $`b=0`$ and $`q(F)3`$.
###### Demonstration Proof
We formulate the proof through three steps. Though the proof is slightly longer, it’s a case by case discussion.
Step 1. $`p_g(X)2`$
Suppose $`b=1`$. In this situation, we see that the movable part of $`|3K_X|`$ defines a morphism. Because $`p_g(X)2`$, $`\text{dim}\varphi _1(X)=1`$ and both $`\varphi _1`$ and $`\varphi _3`$ derive the same fibration $`f:XC`$. So the movable part of $`|K_X|`$ is also base point free. Let $`M_1`$ be the movable part of $`|K_X|`$. Then $`M_1_{\text{lin}}F_i`$. Let $`F`$ be a general fiber of $`f`$. Because the singularities on $`X`$ are all isolated, $`F`$ is a smooth projective surface of general type. By Theorem 2.2, we see that $`H^1(X,2K_X)=0`$. Therefore we have
$$|2K_X+F_i||_F=|2K_F|.$$
Because $`p_g(F)>0`$, $`\mathrm{\Phi }_{|2K_F|}`$ is generically finite by Theorem 1 of \[X1\]. This means that $`\varphi _3`$ is generically finite and so is $`\varphi _m`$ for $`m4`$. So we only have to consider the case when $`b=0`$.
Now we have a fibration $`f:X^{}^1`$. Because $`p_g(X)2`$, we have $`P_3(X)4`$. Thus, by Proposition 3.1, we see that $`p_g(F)=1`$. In this situation, Kollár’s technique (the proof of Corollary 4.8 in \[Ko1\]) is still effective. Let $`p_g(X)=k+1`$, $`k1`$. Since $`|K_X^{}|`$ is composed of a pencil, we have $`𝒪(k)f_{}\omega _X^{}`$ on $`^1`$. If $`k5`$, then
$$𝒪(5)𝒪(k)f_{}\omega _X^{}.$$
Thus we have
$$:=𝒪(1)f_{}\omega _{X^{}/^1}^2=𝒪(5)f_{}\omega _X^{}^2f_{}\omega _X^{}^3.$$
The local sections of $`f_{}\omega _X^{}^2`$ give the bicanonical map of the fiber $`F`$ and they extend to global sections of $``$, because $``$ is generated by global sections. On the other hand, $`H^0(^1,)`$ can distinguish different fibers of $`f`$ because $`f_{}\omega _{X^{}/𝒫^1}^2`$ is a sum of line bundles with nonnegative degree on $`^1`$. So $`H^0(^1,)`$ gives a generically finite map on $`X^{}`$ and so does $`H^0(X^{},3K_X^{})`$. This contradicts to our assumption of $`\text{dim}\varphi _3(X)=1`$. Thus we have $`k4`$, i.e. $`p_g(X)5`$. By virtue of this technique, we have
$`k=3,4,`$ $`f_{}\omega _X^{}^4m=3(A3)`$
$`k=2,`$ $`f_{}\omega _X^{}^5m4(A2),(A3)`$
$`k=1,`$ $`f_{}\omega _X^{}^7m6(A1),(A2),(A3).`$
In order to complete the proof for this case, we have to prove $`q(X)=0`$ for (A1), (A2) and that the only possibility of $`F`$ in (A1) is $`p_g(F)=q(F)=1`$. Suppose $`q(X)=1`$ in cases (A1) and (A2). Then $`q(F)=1`$ and $`R^1f_{}\omega _X^{}\omega _^1`$. Because $`f_{}\omega _X^{}`$ is of positive degree, we have $`h^1(^1,f_{}\omega _X^{})=0`$. So by (1.2), $`h^2(𝒪_X^{})=0`$. Then we have $`\chi (𝒪_X^{})2`$. According to Reid’s plurigus formula (\[R\]), we have $`P_2(X)7`$. This means $`𝒪(6)f_{}\omega _X^{}^2`$ and $`f_{}\omega _X^{}^4`$. So $`\varphi _4`$ is generically finite, a contradiction. Finally, with regard to (A1), if $`q(F)=0`$, then $`h^2(𝒪_X^{})=0`$. So $`\chi (𝒪_X^{})1`$. Thus $`P_3(X)6`$ according to Reid. We have $`𝒪(5)f_{}\omega _X^{}^3`$ and $`f_{}\omega _X^{}^5`$. So $`\varphi _5`$ is generically finite, a contradiction.
Step 2. $`p_g(X)=1`$
When $`b=1`$ or $`b=0`$ and $`P_m(X)3`$, we have $`p_g(F)=1`$ according to Proposition 3.1. This leads to (B0). From now on, we can suppose $`b=0`$, $`P_m(X)=2`$ and $`p_g(F)2`$.
We claim that $`P_{m2}(X)=1`$. In fact, if $`P_{m2}(X)>1`$, we must have $`P_{m2}(X)=2`$. So the movable part of $`|(m2)K_X^{}|`$ is a fiber $`F`$ of $`f`$. By Theorem 1.1, we have
$$|K_X^{}+\mathrm{}\pi ^{}(K_X)\mathrm{}+F||_F=|K_F+D|,$$
where $`D:=\mathrm{}\pi ^{}(K_X)\mathrm{}|_F`$ is an effective divisor on $`F`$. So we see that $`\text{dim}\varphi _m(X)2`$, a contradiction.
If $`P_{m1}(X)=2`$, then we can see from the above argument that $`\text{dim}\varphi _{m+1}(X)2`$. This leads to (B3).
If $`P_{m1}(X)=1`$ and we are not in (B1), then $`P_{m+1}(X)=2`$ by virtue of Proposition 3.1. We can easily see that $`\text{dim}\varphi _{m+2}(X)2`$. This leads to (B2).
Step 3. $`p_g(X)=0`$
We can suppose $`b=1`$ and $`q(F)2`$ or $`b=0`$ and $`q(F)3`$. Otherwise we are in case (C0).
Suppose $`b=1`$. If $`q(F)3`$, we can see that $`\text{dim}\varphi _{m+1}(X)2`$. In fact, we can write
$$m\pi ^{}(K_X)_{}aF+E_{}^{(m)},$$
where $`a=P_m(X)2`$ and $`E_{}^{(m)}`$ is an effective $``$-divisor. It is obvious that
$$K_X^{}+\mathrm{}m\pi ^{}(K_X)F\frac{1}{a}E_{}^{(m)}\mathrm{}(m+1)K_X^{}.$$
Since
$$m\pi ^{}(K_X)F\frac{1}{a}E_{}^{(m)}_{\text{num}}m(1\frac{1}{a})\pi ^{}(K_X)$$
is nef and big, we get by the vanishing theorem that
$$H^1(X^{},K_X^{}+\mathrm{}m\pi ^{}(K_X)F\frac{1}{a}E_{}^{(m)}\mathrm{})=0.$$
This means that
$$|K_X^{}+\mathrm{}m\pi ^{}(K_X)\frac{1}{a}E_{}^{(m)}\mathrm{}||_F=|K_F+\mathrm{}m\pi ^{}(K_X)\frac{1}{a}E_{}^{(m)}\mathrm{}|_F|,$$
where
$$\mathrm{}m\pi ^{}(K_X)\frac{1}{a}E_{}^{(m)}\mathrm{}|_F=\mathrm{}(1\frac{1}{a})E_{}^{(m)}\mathrm{}|_F$$
is an effective divisor. According to \[X2\], $`|K_F|`$ gives a generically finite map. So we see that $`\text{dim}\varphi _{m+1}(F)=2`$ and thus $`\text{dim}\varphi _{m+1}(X)2`$. This leads to (C1). If $`q(F)=2`$ and $`P_m(X)3`$, we can still see that $`\text{dim}\varphi _{m+1}(X)2`$. In this situation, $`a3`$. Let $`F_1`$ and $`F_2`$ be two distinct general fibers of $`f`$. Then we see that
$$m\pi ^{}(K_X)F_1F_2\frac{2}{a}E_{}^{(m)}_{\text{num}}m(1\frac{2}{a})\pi ^{}(K_X)$$
is nef and big. So we have the following surjective map
$`H^0(X^{},K_X^{}+\mathrm{}m\pi ^{}(K_X){\displaystyle \frac{2}{a}}E_{}^{(m)}\mathrm{})`$
$`H^0(F_1,K_{F_1}+D_1)H^0(F_2,K_{F_2}+D_2)0,`$
where
$$D_i=\mathrm{}m\pi ^{}(K_X)\frac{2}{a}E_{}^{(m)}\mathrm{}|_{F_i}=\mathrm{}m(1\frac{2}{a})E_{}^{(m)}\mathrm{}|_{F_i}$$
is effective for all $`i`$. This means that
$$|K_X^{}+\mathrm{}m\pi ^{}(K_X)\frac{2}{a}E_{}^{(m)}\mathrm{}|$$
can distinguish two different fibers of $`f`$ and $`\text{dim}\varphi _{m+1}(F_i)1`$. We again see that $`\text{dim}\varphi _{m+1}(X)2`$. This leads to (C1). If $`q(F)=2`$ and $`P_m(X)=2`$, we can use a parallel argument to that in the proof of the case $`b=1`$ of Step 1 to see that $`\text{dim}\varphi _{m+2}(X)2`$. This corresponds to (C2).
Suppose $`b=0`$ and $`q(F)3`$. If $`P_m(X)3`$, we can use the same argument as in the case $`b=1`$ of Step 3 to see that $`\text{dim}\varphi _{m+1}(X)2`$. This leads to (C1). What remains to be studied is the case $`P_m(X)=2`$. This is the most frustrating case. Anyway, it is easy to see that $`\text{dim}\varphi _{2m+1}(X)2`$ in this case. Actually, one only has to consider the system
$$|K_X^{}+\mathrm{}m\pi ^{}(K_X)\mathrm{}+F|.$$
We can see that
$$|K_X^{}+\mathrm{}m\pi ^{}(K_X)\mathrm{}+F||_F=|K_F+\mathrm{}m\pi ^{}(K_X)\mathrm{}|_F|,$$
where $`\mathrm{}m\pi ^{}(K_X)\mathrm{}|_F`$ is effective. This means
$$\text{dim}\varphi _{2m+1}(X)\text{dim}\varphi _{2m+1}(F)=2.$$
Now if $`\text{dim}\varphi _{2m}(X)2`$, we are in (C3). If $`\text{dim}\varphi _{2m}(X)=1`$, we have the following claim which shows that we are in either (C1) or (C4). We note that $`P_{2m}(X)3`$.
Claim. If $`b=0`$, $`q(F)3`$, $`P_m(X)=2`$, $`P_{2m}(X)4`$ and $`\text{dim}\varphi _{2m}(X)=1`$. Then $`\text{dim}\varphi _{m+1}(X)2`$. This leads to (C1).
Since $`\text{dim}\varphi _{2m}(X)=1`$, we can see that both $`\varphi _{2m}`$ and $`\varphi _m`$ derive the same fibration $`f:X^{}^1`$. We can write
$$m\pi ^{}(K_X)_{}F+(D_v+D_h),$$
where $`D_v`$ is a vertical $``$-divisor with respect to the fibration $`f`$ and $`D_h`$ is the horizontal part. The supports of $`D_v`$ and $`D_h`$ are contained in the fixed part of $`|mK_X^{}|`$. $`D_v`$ and $`D_h`$ are both effective $``$-divisors. Similarly, we can write
$$2m\pi ^{}(K_X)_{}\underset{i=1}{\overset{a_2}{}}F_i+(D_v^{}+D_h^{}),$$
where $`a_2P_{2m}(X)13`$, $`D_v^{}`$ is a vertical effective $``$-divisor and $`D_h^{}`$ is a horizontal effective $``$-divisor. Since the support of $`D_h^{}`$ is contained in the fixed part of $`|2mK_X^{}|`$, we can see that $`D_h^{}=2D_h`$. Now we have
$$2m\pi ^{}(K_X)_{\text{num}}a_2F+D_v^{}+2D_h,$$
$$m\pi ^{}(K_X)_{\text{num}}\frac{a_2}{2}F+\frac{1}{2}D_v^{}+D_h.$$
So
$$m\pi ^{}(K_X)F\frac{1}{a_2}D_v^{}\frac{2}{a_2}D_h_{\text{num}}m(1\frac{2}{a_2})\pi ^{}(K_X)$$
is nef and big. This means, according to the vanishing theorem, that
$$H^1(X^{},K_X^{}+\mathrm{}m\pi ^{}(K_X)F\frac{1}{a_2}D_v^{}\frac{2}{a_2}D_h\mathrm{})=0.$$
Denote $`M:=\mathrm{}m\pi ^{}(K_X)\frac{1}{a_2}D_v^{}\frac{2}{a_2}D_h\mathrm{}.`$ Then $`K_X^{}+M(m+1)K_X^{}`$. Then
$$|K_X^{}+M||_F=|K_F+M|_F|,$$
where $`M|_F=\mathrm{}(1\frac{2}{a_2})D_h\mathrm{}|_F`$ is an effective divisor on $`F`$. So
$$\text{dim}\varphi _{m+1}(X)\text{dim}\mathrm{\Phi }_{|K_X^{}+M|}(F)=2.$$
The proof is complete. ∎
###### Theorem 3.4
Let $`f:XC`$ be a derived family of surfaces from the bicanonical pencil $`|2K_X|`$ of a smooth projective 3-fold $`X`$ of general type. Let $`F`$ be a general fiber of $`f`$ and denote $`b:=g(C)`$. Then one of the following occurs.
(A0)’ $`p_g(X)>0`$, $`q(X)b+1`$ and $`p_g(F)=1`$.
(A1)’ $`1p_g(X)2`$, $`2P_2(X)3`$, $`\text{dim}\varphi _3(X)2`$ and $`p_g(F)2`$.
(A2)’ $`p_g(X)=1`$, $`P_2(X)=P_3(X)=2`$, $`\text{dim}\varphi _4(X)2`$, $`b=0`$ and $`p_g(F)2`$.
(B0)’ $`p_g(X)=0`$, $`q(X)2`$ and $`q(F)2`$.
(B1)’ $`p_g(X)=0`$, $`\text{dim}\varphi _3(X)2`$ and $`q(F)2`$.
(B2)’ $`p_g(X)=0`$, $`P_2(X)=2`$, $`\text{dim}\varphi _4(X)2`$, $`b=1`$ and $`q(F)=2`$.
(B3)’ $`p_g(X)=0`$, $`P_2(X)=2`$, $`\text{dim}\varphi _4(X)2`$, $`b=0`$ and $`q(F)3`$.
(B4)’ $`p_g(X)=0`$, $`P_2(X)=2`$, $`P_4(X)=3`$, $`\text{dim}\varphi _5(X)2`$, $`b=0`$ and $`q(F)3`$.
###### Demonstration Proof
The proof is parallel to that of Theorem 3.3 except that we have more cases here. In order to avoid unnecessary redundancy, we only give the proof where it is different from the respective part in the proof of Theorem 3.3.
Step 1. $`p_g(X)2`$.
In this case, we always have $`P_2(X)3`$. When $`b=0`$ and $`P_2(X)4`$ or $`b=1`$, we see from Propositioin 3.1 that $`p_g(F)=1`$. This leads to (A0)’.
So we only have to consider the case with $`b=0`$, $`P_2(X)=3`$ and $`p_g(F)2`$. In this situation, we see that the movable part of $`|2K_X^{}|`$ contains exactly 2 fibers of $`f`$. Since $`P_3(X)4`$, Proposition 3.1 gives $`\text{dim}\varphi _3(X)2`$. This corresponds to (A1)’.
Step 2. $`p_g(X)=1`$.
Excluding the situation (A0)’ while observing Proposition 3.1, we only have to consider the case with $`p_g(F)2`$ and with the following extra properties:
$$b=0,P_2(X)3\text{or}b=1,P_2(X)=2.$$
When $`b=0`$ and $`P_2(X)=3`$ or $`b=1`$ and $`P_2(X)=2`$, we know from Proposition 3.1 that $`\text{dim}\varphi _3(X)2`$. This leads to (A1)’.
When $`b=0`$, $`P_2(X)=2`$, $`p_g(F)2`$ and $`P_3(X)3`$, we see from Proposition 3.1 that $`\text{dim}\varphi _3(X)2`$. This also corresponds to (A1)’. Otherwise we always have $`\text{dim}\varphi _4(X)2`$ because $`P_4(X)3`$. This is just (A2)’.
Step 3. $`p_g(X)=0`$.
The argument in the proof of Theorem 3.3 is still effective in this case. We can see that (C0) through (C4) correspond to (B0)’ through (B4)’, respectively. We omit the proof. ∎
Now we can see that Theorem 1 (i), (iv) and (v), Theorem 2 are direct results from Theorem 3.3 and Theorem 3.4. In order to complete the proof of Theorem 1, we only have to show $`q(X)2`$ whenever $`m11`$ or $`m7`$ and $`p_g(X)>0`$.
###### Proposition 3.5
Let $`X`$ be a minimal projective 3-fold of general type with only $``$-factorial terminal singularities. Suppose $`q(X)3`$, $`P_{k_0}(X)>0`$ and $`P_{k_2}(X)2`$. Then $`\text{dim}\varphi _{k_0+k_2+1}(X)2`$.
###### Demonstration Proof
Choose a 1-dimensional subsystem $`\mathrm{\Lambda }|k_2K_X|`$ while taking a birational modification $`\pi :X^{}X`$ such that the pencil $`\mathrm{\Lambda }`$ defines a morphism $`g:X^{}^1`$. We can even take further modification to $`\pi `$ so that $`\pi ^{}(k_2K_X)`$ has supports with only normal crossings. Taking the Stein factorization of $`g`$, then we get a derived fibration $`p:X^{}C_1`$. We note that this fibration is different from the one which was defined at the first page of this paper. Denote $`b_1:=g(C_1)`$. Let $`M`$ be the movable part of the pencil. We obviously have $`Mk_2K_X^{}`$. We can write $`M_{\text{lin}}_{i=1}^{a_1}F_i`$, where $`a_11`$ and $`F_i`$ is a fiber of $`p`$ for all $`i`$. We also note that $`a_1=1`$ if and only if $`b_1=0`$. A general fiber $`F`$ is a smooth projective surface of general type.
Suppose $`b_1>0`$. Then $`|M|`$ is base point free on $`X`$. Because $`X`$ has only isolated singularities, $`F`$ is smooth. We study the system $`|tK_X+M|`$ where $`t2`$. We know that $`M`$ contains at least two components $`F_1`$ and $`F_2`$. By Theorem 1.2, we see that
$$H^0(X,tK_X+M)H^0(F_1,tK_{F_1})H^0(F_2,tK_{F_2})$$
is surjective. This means that $`\mathrm{\Phi }_{|tK_X+M|}`$ can distinguish $`F_1`$ and $`F_2`$ and the restriction to $`F_i`$ is at least a bicanonical map. We know that $`\text{dim}\mathrm{\Phi }_{|tK_{F_i}|}(F_i)1`$ for all $`t2`$. Noting that the image of $`X`$ through $`\mathrm{\Phi }_{|tK_X+M|}`$ is irreducible, we see that $`\text{dim}\mathrm{\Phi }_{|tK_X+M|}(X)2`$. So $`\text{dim}\varphi _{t+k_2}(X)2`$. Thus $`\text{dim}\varphi _{k_0+k_2+1}(X)2`$.
Suppose $`b=0`$. By Corollary 2.3, we have $`q(F)3`$. In this case, $`M_{\text{lin}}F`$. We have
$$|K_X^{}+\mathrm{}k_0\pi ^{}(K_X)\mathrm{}+F||_F=|K_F+\mathrm{}k_0\pi ^{}(K_X)\mathrm{}|_F|,$$
where $`\mathrm{}k_0\pi ^{}(K_X)\mathrm{}|_F`$ is effective. Thus $`\text{dim}\varphi _{k_0+k_2+1}(X)\text{dim}\varphi _{k_0+k_2+1}(F)=2`$. ∎
###### Lemma 3.6
Let $`X`$ be a smooth projective 3-fold of general type. If $`q(X)3`$, then either
$$P_2(X)>0\text{and}P_4(X)2$$
or
$$P_k(X)2\text{for all}k5.$$
###### Demonstration Proof
This is a byproduct from the proof of both Theorem 6.1, \[Ko1\] and Proposition 4.3, \[Ko1\]. ∎
###### Proposition 3.7
Let $`X`$ be a minimal projective 3-fold of general type with only $``$-factorial terminal singularities. Suppose $`q(X)3`$, $`P_2(X)>0`$ and $`P_4(X)2`$. Then
(1) $`\text{dim}\varphi _7(X)2`$. $`\text{dim}\varphi _m(X)2`$ for all $`m9`$.
(2) If $`\text{dim}\varphi _8(X)=1`$, then $`p_g(X)=0`$, $`P_2(X)=1`$, $`P_4(X)=2`$ and $`P_8(X)=3`$.
###### Demonstration Proof
(1) Let $`k_0=2`$ and $`k_2=4`$. Proposition 3.5 gives $`\text{dim}\varphi _7(X)2.`$ So $`\text{dim}\varphi _{2l+7}(X)2`$ for all $`l^+`$. Let $`k_0=2`$ and $`k_2=7`$. Applying Proposition 3.5 again, we get $`\text{dim}\varphi _{10}(X)2`$. Thus $`\text{dim}_{2l+10}(X)2`$ for all $`l^+`$.
(2) We study $`\varphi _8`$. We have $`P_8(X)3`$. If $`P_8(X)4`$ and $`\text{dim}\varphi _8(X)=1`$, we want to deduce a contradiction. We know that both $`\varphi _4`$ and $`\varphi _8`$ derive the same fibration $`f:X^{}C`$ which was described in the first page of this paper. If $`b>0`$, it is easy to see that $`\text{dim}\varphi _8(X)2`$ by a standard argument which has been used many times in this paper. So we can suppose $`b=0`$. Since $`q(X)3`$, we have $`q(F)3`$. Suppose $`M_4`$, $`M_8`$ are the movable parts of $`|4K_X^{}|`$, $`|8K_X^{}|`$ respectively. Then we have
$$4\pi ^{}(K_X)_{}M_4+E_4,$$
$$8\pi ^{}(K_X)_{}M_8+E_8,$$
where $`E_4`$ and $`E_8`$ are effective $``$-divisors. Let $`E_v`$, $`E_v^{}`$ be the vertical parts of $`E_4`$, $`E_8`$ and $`E_h`$, $`E_h^{}`$ be the horizontal parts of $`E_4`$, $`E_8`$ respectively. Because the support of $`E_h`$ is contained in the fixed part of $`|4K_X^{}|`$ and the support of $`E_h^{}`$ is contained in the fixed part of $`|8K_X^{}|`$, we see that $`E_h^{}=2E_h`$. Now we have
$$8\pi ^{}(K_X)_{\text{num}}a_8F+E_v^{}+2E_h,$$
where $`a_83`$. It follows that
$$4\pi ^{}(K_X)_{\text{num}}\frac{a_8}{2}F+\frac{1}{2}E_v^{}+E_h.$$
Thus
$$4\pi ^{}(K_X)F\frac{1}{a_8}E_v^{}\frac{2}{a_8}E_h_{\text{num}}4(1\frac{2}{a_8})\pi ^{}(K_X)$$
is a nef and big $``$-divisor. Denote
$$G:=\mathrm{}4\pi ^{}(K_X)\frac{1}{a_8}E_v^{}\frac{2}{a_8}E_h\mathrm{}.$$
Then we have $`H^1(X^{},K_X^{}+GF)=0`$. So we see that
$$|K_X^{}+G||_F=|K_F+G|_F|,$$
where
$$G|_F=\mathrm{}4\pi ^{}(K_X)\frac{1}{a_8}E_v^{}\frac{2}{a_8}E_h\mathrm{}|_F=\mathrm{}(1\frac{2}{a_8})E_h\mathrm{}|_F$$
is effective. So $`\text{dim}\varphi _5(X)\text{dim}\mathrm{\Phi }_{|K_F+G|_F|}(F)=2`$. In particular, $`P_5(X)2`$. Now let $`k_0=2`$ and $`k_2=5`$. Applying Proposition 3.5, we see that $`\text{dim}\varphi _8(X)2`$. This contradicts to our assumption. Thus we have seen that $`P_8(X)=3`$ if $`\text{dim}\varphi _8(X)=1`$. It follows immediately that $`P_4(X)=2`$ and $`P_2(X)=1`$. If $`P_g(X)>0`$, it is very easy to see from Proposition 3.5 that $`\text{dim}\varphi _8(X)2`$. So we have completed the proof. ∎
###### Proposition 3.8
Let $`X`$ be a minimal projective 3-fold of general type with only $``$-factorial terminal singularities. Suppose $`q(X)3`$ and $`P_k(X)2`$ for all $`k5`$. Then
(1) $`\text{dim}\varphi _m(X)2`$ for all $`m11`$.
(2) If $`\text{dim}\varphi _{10}(X)=1`$, then $`p_g(X)=P_2(X)=P_3(X)=P_4(X)=0`$.
(3) If $`\text{dim}\varphi _9(X)=1`$, then $`p_g(X)=P_2(X)=P_3(X)=0`$.
(4) If $`\text{dim}\varphi _8(X)=1`$, then $`p_g(X)=P_2(X)=0`$.
(5) If $`\text{dim}\varphi _7(X)=1`$, then $`p_g(X)=0`$.
###### Demonstration Proof
Let $`k_0=5`$ and $`k_2=t5`$. Proposition 3.5 gives $`\text{dim}\varphi _{t+6}(X)2`$ for all $`t5`$. This leads to (1).
If $`p_g(X)>0`$, let $`k_0=1`$ and $`k_2=t5`$. Proposition 3.5 gives $`\text{dim}\varphi _{t+2}(X)2`$ for all $`t5`$.
If $`P_2(X)>0`$, let $`k_0=2`$ and $`k_2=t5`$. Proposition 3.5 gives $`\text{dim}\varphi _{t+3}(X)2`$ for all $`t5`$.
If $`P_3(X)>0`$, let $`k_0=3`$ and $`k_2=t5`$. Proposition 3.5 gives $`\text{dim}\varphi _{t+4}(X)2`$ for all $`t5`$.
If $`P_4(X)>0`$, let $`k_0=4`$ and $`k_2=t5`$. Proposition 3.5 gives $`\text{dim}\varphi _{t+5}(X)2`$ for all $`t5`$.
(2), (3), (4) and (5) follow immediately. ∎
###### Corollary 3.9
Let $`f:XC`$ be a derived family of surfaces from the m-canonical pencil $`|mK_X|`$ of a smooth projective 3-fold $`X`$ of general type. Then
(1) $`q(X)2`$ when $`m11`$.
(2) $`q(X)2`$ when $`m7`$ and $`p_g(X)>0`$.
###### Demonstration Proof
This is obvious from Lemma 3.6, Proposition 3.7 and Proposition 3.8. ∎
## 4. Appendix to Kollár’s method
Given a smooth projective 3-fold $`X`$ of general type, it is uncertain whether $`\text{dim}\varphi _{m+1}(X)\text{dim}\varphi _m(X)`$ for all $`m>0`$. Even if $`P_m(X)>0`$, it is false that $`P_{m+1}(X)>0`$. This makes it difficult to study some stable property of $`\varphi _m`$. That is why Kollár’s bound was bigger. Hereby we would like to study in an alternative way. The bounds are better, however still unsatisfactory.
###### Proposition 4.1
Let $`X`$ be a minimal projective 3-fold of general type with only $``$-factorial terminal singularities. Suppose $`q(X)2`$ and $`P_{k_2}(X)2`$. Then $`\text{dim}\varphi _m(X)2`$ for all $`m4k_2+2`$.
###### Demonstration Proof
Choose a 1-dimensional subsystem $`\mathrm{\Lambda }|k_2K_X|`$ while taking a birational modification $`\pi :X^{}X`$ such that the pencil $`\mathrm{\Lambda }`$ defines a morphism $`g:X^{}^1`$. We can even take further modification to $`\pi `$ so that $`\pi ^{}(k_2K_X)`$ has supports with only normal crossings. Taking the Stein factorization of $`g`$, then we get a derived fibration $`p:X^{}C_2`$. Denote $`b_2:=g(C_2)`$. Let $`M`$ be the movable part of the pencil. We obviously have $`Mk_2K_X^{}`$. We can write $`M_{\text{lin}}_{i=1}^{a_2}F_i`$, where $`a_21`$ and $`F_i`$ is a fiber of $`q`$. A general fiber $`F`$ is a smooth projective surface of general type.
If $`b_2>0`$, then we can see that $`\text{dim}\varphi _{k_2+t}(X)2`$ for all $`t2`$ according to the parallel argument in the proof of Proposition 3.5.
If $`b_2=0`$. we study in an alternative way. We have $`M_{\text{lin}}F`$. Because $`q(X)2`$, we have $`p_g(F)q(F)2`$. According to Theorem 1.1, we have
$$|K_X^{}+\mathrm{}k_2\pi ^{}(K_X)\mathrm{}+F||_F=|K_F+\mathrm{}k_2\pi ^{}(K_X)\mathrm{}|_F|.$$
This means $`\text{dim}\varphi _{2k_2+1}(F)1`$ because $`\mathrm{}k_2\pi ^{}(K_X)\mathrm{}|_F`$ is effective. Suppose $`M_{2k_2+1}`$ is the movable part of $`|(2k_2+1)K_X^{}|`$ and $`M_{2k_2+1}^{}`$ is the movable part of
$$|K_X^{}+\mathrm{}k_2\pi ^{}(K_X)\mathrm{}+F|.$$
Then $`M_{2k_2+1}^{}M_{2k_2+1}`$. Let $`M_0`$ be the movable part of $`|K_F|`$. Then $`h^0(F,M_0)2`$. Considering the following two maps
$$H^0(X^{},K_X^{}+\mathrm{}k_2\pi ^{}(K_X)\mathrm{}+F)\stackrel{𝛼}{}H^0(F,K_F+\mathrm{}k_2\pi ^{}(K_X)\mathrm{}|_F)0$$
$$H^0(X^{},M_{2k_2+1}^{})\stackrel{𝛽}{}H^0(F,M_{2k_2+1}^{}|_F),$$
we know that $`\alpha `$ is surjective and the images of $`\alpha `$ and $`\beta `$ have the same dimension. So
$`h^0(F,M_{2k_2+1}^{}|_F)\text{dim}_{}\text{im}(\beta )=\text{dim}_{}\text{im}(\alpha )`$
$`=h^0(F,K_F+\mathrm{}k_2\pi ^{}(K_X)\mathrm{}|_F).`$
Because
$$M_{2k_2+1}^{}|_FK_F+\mathrm{}k_2\pi ^{}(K_X)\mathrm{}|_F,$$
we see that
$$M_0M_{2k_2+1}^{}|_FM_{2k_2+1}|_F.$$
For all $`t0`$ and two different fibers $`F_1`$, $`F_2`$, we consider the system
$$|K_X^{}+\mathrm{}(t+2k_2+1)\pi ^{}(K_X)\mathrm{}+F_1+F_2|.$$
It is obvious that
$$|K_X^{}+\mathrm{}(t+2k_2+1)\pi ^{}(K_X)\mathrm{}+F_1+F_2||(t+4k_2+2)K_X^{}|.$$
From Theorem 1.1, we have the exact sequence
$`H^0(X^{},K_X^{}+\mathrm{}(t+2k_2+1)\pi ^{}(K_X)\mathrm{}+F_1+F_2)`$
$`H^0(F_1,K_{F_1}+G_1)H^0(F_2,K_{F_2}+G_2)0,`$
where $`G_i=(\mathrm{}(t+2k_2+1)\pi ^{}(K_X)\mathrm{}+F_1+F_2)|_{F_i}`$ for all $`i`$. We can see that
$`K_{F_i}+G_i`$ $`K_{F_i}+\mathrm{}t\pi ^{}(K_X)|_{F_i}\mathrm{}+M_{2k_2+1}|_{F_i}`$
$`K_{F_i}+\mathrm{}t\pi ^{}(K_X)|_{F_i}\mathrm{}+M_0`$
for all $`i`$. Furthermore, one can see that
$$\text{dim}\mathrm{\Phi }_{|K_{F_i}+\mathrm{}t\pi ^{}(K_X)|_{F_i}\mathrm{}+M_0|}(F_i)1.$$
When $`t=0`$, it is obvious. When $`t>0`$, one need to use the vanishing theorem to prove it. Noting that the image of $`X^{}`$ through $`\varphi _{t+4k_2+2}`$ is irreducible and that
$$\text{dim}\varphi _{t+4k_2+2}(F_i)1$$
for all $`i`$, we can see $`\text{dim}\varphi _{t+4k_2+2}(X)2`$. ∎
###### Proposition 4.2
Let $`X`$ be a minimal projective 3-fold of general type with only $``$-factorial terminal singularities. Suppose $`q(X)>0`$ and $`P_{k_2}(X)2`$. Then $`\text{dim}\varphi _m(X)2`$ for all $`m7k_2+3`$.
###### Demonstration Proof
We keep the same set up as in the proof of Proposition 4.1. We only have to study the case when $`b_2=0`$. We have the fibration $`p:X^{}^1`$. We still denote by $`F`$ a general fiber of $`p`$. Since $`q(X)>0`$, we get $`p_g(F)q(F)1`$. If $`p_g(F)2`$, we have seen from the proof of the last proposition that we can get better bounds. The most frustrating case is when $`p_g(F)=q(F)=1`$. Let $`\sigma :FF_0`$ be the contraction onto the minimal model. According to Theorem 3.1 in \[Ci\], we know that $`|2K_{F_0}|`$ is base point free when $`p_g(F)>0`$. So the movable part of $`|2K_F|`$ is just $`\sigma ^{}(2K_{F_0})`$. According to Kollár’s method, we see that
$$|(5k_2+2)K_X^{}||_F|2K_F|.$$
So, if we denote by $`M_{5k_2+2}`$ the movable part of $`|(5k_2+2)K_X^{}|`$, we should have
$$M_{5k_2+2}|_F\sigma ^{}(2K_{F_0}).$$
For all $`t0`$ and two different fibers $`F_1`$, $`F_2`$, we consider the system
$$|K_X^{}+\mathrm{}(t+5k_2+2)\pi ^{}(K_X)\mathrm{}+F_1+F_2|.$$
It is obvious that
$$|K_X^{}+\mathrm{}(t+5k_2+2)\pi ^{}(K_X)\mathrm{}+F_1+F_2||(t+7k_2+3)K_X^{}|.$$
From Theorem 1.1, we have the exact sequence
$`H^0(X^{},K_X^{}+\mathrm{}(t+5k_2+2)\pi ^{}(K_X)\mathrm{}+F_1+F_2)`$
$`H^0(F_1,K_{F_1}+G_1^{})H^0(F_2,K_{F_2}+G_2^{})0,`$
where $`G_i^{}=(\mathrm{}(t+5k_2+2)\pi ^{}(K_X)\mathrm{}+F_1+F_2)|_{F_i}`$ for all $`i`$. We can see that
$`K_{F_i}+G_i^{}`$ $`K_{F_i}+\mathrm{}t\pi ^{}(K_X)|_{F_i}\mathrm{}+M_{5k_2+2}|_{F_i}`$
$`K_{F_i}+\mathrm{}t\pi ^{}(K_X)|_{F_i}\mathrm{}+\sigma ^{}(2K_{F_0})`$
for all $`i`$. Furthermore, one can see that
$$\text{dim}\mathrm{\Phi }_{|K_{F_i}+\mathrm{}t\pi ^{}(K_X)|_{F_i}\mathrm{}+\sigma ^{}(2K_{F_0})|}(F_i)1.$$
When $`t=0`$, it is obvious. When $`t>0`$, one need to use the vanishing theorem to prove it. Noting that the image of $`X^{}`$ through $`\varphi _{t+7k_2+3}`$ is irreducible and that
$$\text{dim}\varphi _{t+7k_2+3}(F_i)1$$
for all $`i`$, we can see $`\text{dim}\varphi _{t+7k_2+3}(X)2`$. The proof is complete. ∎
###### Corollary 4.3
Let $`f:XC`$ be a derived family of surfaces from the m-canonical pencil $`|mK_X|`$ of a smooth projective 3-fold $`X`$ of general type. Then
(1) $`q(X)1`$ whenever $`m82`$.
(2) $`q(X)=0`$ whenever $`m143`$.
###### Demonstration Proof
According to \[F\] and Remark 6.6 in \[Ko1\], we always have $`P_{20}(X)2`$ if $`q(X)>0`$. Let $`k_2=20`$ while applying Proposition 4.1 and Proposition 4.2, we get what we want. ∎
## Acknowledgment
This note was written while I was visiting as a post-doc fellow at the Mathematisches Institut der Universit$`\ddot{\text{a}}`$t G$`\ddot{\text{o}}`$ttingen, Germany. I would like to thank Prof. F. Catanese for useful discussions and helps during my stay at G$`\ddot{\text{o}}`$ttingen.
## References
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# 1 Introduction
## 1 Introduction
Field theories at finite temperature and density have been proposed as the fundamental underlying theory for the description of the physics of the early universe. A proposed scenario for baryogenesis is by the electroweak phase transition . QCD is expected to become deconfined at high temperature. The formation of a quark gluon plasma and the phase transitions of QCD are supposed to be visible in relativistic heavy ion collision and astrophysics . A modern presentation of finite temperature field theory can be found in .
Beyond their phenomenological implications, quantum field theories at finite temperature are very challenging also from the more theoretical point of view. There is a real-time as well as an imaginary-time formalism, the first describing dynamical and the second equilibrium properties . Many fundamental issues and problems are unsolved so far or require a deeper understanding. Quantum field theories are subject to enhanced complexities compared to zero temperature and zero density. This is largely related to the presence of additional length scales, due to the interaction with a heat bath. On the various scales the properties of the theory are considerably different.
The separation of scales is widely believed to be an intrinsic property of the field theory. In QCD the scales are associated to the generation of electric and magnetic screening and plasmon masses. In the framework of perturbation theory, this manifests itself in terms of IR divergences that are “severe”. They are not removable as it is the case at temperature $`T=0`$ by adjusting the renormalization prescription . Various elaborate resummation techniques have been proposed to (at least partially) remove the IR singularities and in addition compute screening masses in perturbation theory. In any case, all the approaches (need to) aim at a clean separation of IR and UV behaviour.
A precondition of all these considerations is renormalizability. Renormalizability is an essential requirement of any local quantum field theory, both at zero and non-zero temperature . It implies that the correlation functions stay finite as the UV-cutoff $`\mathrm{\Lambda }_0`$, say, is removed, $`\mathrm{\Lambda }_0\mathrm{}`$, and that the limit is parametrized by a set of renormalized (relevant) coupling constants. Moreover, it is crucial that renormalization can be achieved in a temperature independent way, which means that the field theory renormalized at zero temperature stays UV finite at every $`T>0`$ as well. This is often taken for granted even for complicated theories, such as gauge theories. Temperature independent renormalizability is indispensable for relating bare and renormalized coupling constants in a $`T`$-independent way. It is thus required when formulating Callan-Symanzik type of equations that govern the $`T`$-dependence of observables, including correlation functions and effective masses. More generally it implies that the static and dynamic properties mediated by the interactions with a heat bath are intrinsic features of the field theory itself.
Various attempts and steps towards proving renormalizability exist . In order to separate off the IR problem from the UV scale, a massive field theory is considered. Both in the real- and in the imaginary-time formalism, the investigations are commonly based on a Feynman diagrammatic approach in momentum space. In the real-time description, it is argued that the part of the propagator which depends on the temperature $`T`$ or the chemical potential $`\mu `$ decays exponentially fast for large momenta, so it should be “innocent” of any UV problem. In the imaginary-time formalism the approach is generically more “cumbersome”, but it is again argued that in the sum over the Matsubara frequencies all $`T`$\- or $`\mu `$-dependent UV divergences cancel out.
Experience obtained by explicit computations to leading orders of perturbation theory confirms that, once IR and UV singularities are properly disentangled, all UV divergences found are $`T`$-independent and are removed by the zero temperature counterterms. However, this is not so for non-zero chemical potential $`\mu `$ (associated to a finite density). A field theory that has been renormalized at $`\mu =0`$ is able to generate $`\mu `$-dependent UV divergences that are not removed by the $`\mu =0`$ counterterms. A simple example is given by a 4-dimensional Yukawa model, with a chemical potential associated to the fermion number. In the framework of the renormalization group, the chemical potential introduces an additional relevant operator, so at least one additional renormalization condition is expected. This also indicates a possible problem for the analytic continuation from the euclidean to the real-time formulation, in agreement with a discussion in the framework of axiomatic quantum field theories at finite temperature, where the problem of proving the existence of correlation functions (even at $`\mu =0`$) in the real-time formalism has been stressed.
The renormalization of field theories at $`T=0`$ is well understood. Strong statements and proofs on the renormalizability of various field theories relevant in physics exist, including several different regularization and renormalization schemes, see e.g. . Unfortunately, this sophistication does not extend to finite $`T`$ so far. Rigorous proofs do not exist, to the best of our knowledge. We would like to point out, however, that recently rigorous bounds, uniform in the temperature, have been established for the perturbative correlation functions of many-fermion models. Here renormalization is necessary to obtain well-behaved bounds on the IR side, when approaching the Fermi surface, whereas the UV regularization is kept fixed. Feldman et al. renormalize the many-fermion models with $`T`$-independent counterterms, as we do.
In this paper we give a mathematical proof that massive $`\phi ^4`$ theory at finite $`T`$, in the imaginary-time formalism, is renormalizable. More precisely, we show, to all orders of the loop expansion, that all correlation functions become UV finite at every finite $`T`$ once the theory has been renormalized at $`T=0`$ by (one of the) usual renormalization prescriptions.
The proof is given in the framework of Wilson’s flow equation. It avoids the analysis of individual Feynman integrals (or Feynman sums), which requires the involved combinatorics encoded in the forest formula for overlapping divergences. Moreover it avoids the formulation and proof of a power counting theorem. Using flow equations, the proof of renormalizability merely amounts to establish appropriate bounds in momentum space on the correlation functions, which are viewed as coefficient functions of the associated generating functional. The proof is by induction on the number of loops.
This paper is organized as follows. In Sect. 2 we introduce our basic notations. This includes the definition of the generating functional $`L^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\phi )`$ of the connected, free propagator amputated Green functions on “momentum scale $`\mathrm{\Lambda }`$”, with $`0\mathrm{\Lambda }\mathrm{\Lambda }_0`$, where $`\mathrm{\Lambda }_0`$ denotes the UV cutoff. The dependence of $`L^{\mathrm{\Lambda },\mathrm{\Lambda }_0}`$ on the scale $`\mathrm{\Lambda }`$ is described by the so-called Wilson flow equation. We recap the basic steps of proving renormalizability of 4-dimensional $`\phi ^4`$ field theories at zero temperature by means of the flow equation. Renormalizability is stated in terms of uniform bounds on the (coefficient functions of the) solution $`L^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\phi )`$ of the flow equation and its derivative with respect to the UV-cutoff $`\mathrm{\Lambda }_0`$, with boundary conditions imposed at $`\mathrm{\Lambda }=0`$ for the relevant couplings and at $`\mathrm{\Lambda }=\mathrm{\Lambda }_0`$ for the irrelevant interactions.
In Sect. 3 we show that the difference $`D^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\underset{¯}{\phi };T)`$ of the generating functionals at temperature $`T>0`$ and $`T=0`$:
$$D^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\underset{¯}{\phi };T)L^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\underset{¯}{\phi };T)L^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\underset{¯}{\phi })$$
(1)
has the properties of an irrelevant operator in the sense of the renormalization group <sup>4</sup><sup>4</sup>4For the definition of the momentum space field variables $`\underset{¯}{\phi }`$ and their position space Fourier transform $`\underset{¯}{\overset{^}{\phi }}`$ we refer to the beginning of sect.3 : Equ. (1) should be understood in the weak sense, i.e. in a formal power series expansion w.r.t. $`\mathrm{}`$ and as an identity for all coefficient functions generated by the generating functionals. For the equation to make sense as it stands the variables $`\underset{¯}{\overset{^}{\phi }}`$ have to be appropriately restricted, for instance to be smooth functions, supported in the interval $`[0,\beta ]`$ in the $`x_0`$-component in position space.. More precisely, $`T`$-independence of the counterterms means that the boundary condition
$$D^{\mathrm{\Lambda }_0,\mathrm{\Lambda }_0}(\underset{¯}{\phi };T)0$$
(2)
holds. From this we derive strong bounds on all scales $`\mathrm{\Lambda }`$ for $`D^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\underset{¯}{\phi };T)`$. Together with the bounds on $`L^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\phi )`$ this proves UV finiteness of massive $`\phi _4^4`$ for every finite $`T`$, that is,
$$\underset{\mathrm{\Lambda }_0\mathrm{},\mathrm{\Lambda }0}{lim}L^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\underset{¯}{\phi };T)$$
(3)
exists, to all orders of the loop expansion. As an immediate consequence, the theory is also made UV finite by imposing normalization conditions on the mass, the wave function constant and on the quartic coupling constant at any fixed temperature $`T_0`$. In Sect. 4 we summarize our central statements and give a short outlook.
## 2 Renormalization of zero temperature $`\phi _4^4`$ theory <br>\- a short reminder
Perturbative renormalizability of euclidean zero temperature $`\phi _4^4`$ theory will be established by analysing the generating functional $`L^{\mathrm{\Lambda },\mathrm{\Lambda }_0}`$ of connected (free propagator) amputated Green functions (CAG). The upper indices $`\mathrm{\Lambda }`$ and $`\mathrm{\Lambda }_0`$ enter through the regularized propagator
$$C^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(p)=\frac{1}{p^2+m^2}\{e^{\frac{p^2+m^2}{\mathrm{\Lambda }_0^2}}e^{\frac{p^2+m^2}{\mathrm{\Lambda }^2}}\}$$
(4)
or its Fourier transform
$$\widehat{C}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(x)=_pC^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(p)e^{ipx},$$
(5)
where we use the shorthand
$$_p:=_{\mathrm{I}\mathrm{R}^4}\frac{d^4p}{(2\pi )^4}.$$
(6)
We assume
$$0\mathrm{\Lambda }\mathrm{\Lambda }_0\mathrm{}$$
(7)
so that the Wilson flow parameter $`\mathrm{\Lambda }`$ takes the role of an infrared (IR) cutoff<sup>5</sup><sup>5</sup>5Such a cutoff is of course not necessary in a massive theory. The IR behaviour is only modified for $`\mathrm{\Lambda }`$ above $`m`$., whereas $`\mathrm{\Lambda }_0`$ is the ultraviolet (UV) regularization. The full propagator is recovered for $`\mathrm{\Lambda }=0`$ and $`\mathrm{\Lambda }_0\mathrm{}`$. We also introduce the convention
$$\widehat{\phi }(x)=_p\phi (p)e^{ipx},\frac{\delta }{\delta \widehat{\phi }(x)}=(2\pi )^4_p\frac{\delta }{\delta \phi (p)}e^{ipx}.$$
(8)
For our purposes the ”fields” $`\widehat{\phi }(x)`$ may be assumed to live in the Schwartz space $`𝒮(\mathrm{I}\mathrm{R}^4)`$. For finite $`\mathrm{\Lambda }_0`$ and in finite volume the theory can be given rigorous meaning starting from the functional integral
$$e^{\frac{1}{\mathrm{}}(L^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\widehat{\phi })+I^{\mathrm{\Lambda },\mathrm{\Lambda }_0})}=𝑑\mu _{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\widehat{\varphi })e^{\frac{1}{\mathrm{}}L^{\mathrm{\Lambda }_0,\mathrm{\Lambda }_0}(\widehat{\varphi }+\widehat{\phi })},$$
(9)
where the factors of $`\mathrm{}`$ have been introduced to allow for a consistent loop expansion in the sequel. In (9) $`d\mu _{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\widehat{\varphi })`$ denotes the (translation invariant) Gaussian measure with covariance $`\mathrm{}\widehat{C}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(x)`$. The normalization factor $`e^{\frac{1}{\mathrm{}}I^{\mathrm{\Lambda },\mathrm{\Lambda }_0}}`$ is due to vacuum contributions. It diverges in infinite volume so that we can take the infinite volume limit only when it has been eliminated . We do not make the finite volume explicit here since it plays no role in the sequel.<sup>6</sup><sup>6</sup>6A rigorous treatment of the thermodynamic limit requires to replace the propagator (5) by a finite volume version, e.g. $`\widehat{C}_V^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(x,y)=\chi _V(x)\widehat{C}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(xy)\chi _V(y),`$ where $`\chi _V`$ is the characteristic function of the volume $`V`$, and to regard the Gaussian measure with covariance $`\widehat{C}_V^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(x,y)`$. In this case the quantity $`I_V^{\mathrm{\Lambda },\mathrm{\Lambda }_0}`$ is obviously well defined, at any order $`l`$ in $`\mathrm{}`$. Then (12) is well-defined. After decomposing $`L_V^{\mathrm{\Lambda },\mathrm{\Lambda }_0}`$ w.r.t. powers of $`\mathrm{}`$ and of the field $`\widehat{\phi }`$, we realize that the coefficient functions $`_{l,n}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}`$ are well-defined in the thermodynamic limit, since they are given as finite sums over UV-regularized connected diagrams. The existence of the thermodynamic limit is of course confirmed by the bounds on the solutions of the FE. It should also be feasible to study the thermodynamic limit itself with the aid of the FE in finite volume, by proving inductively uniform bounds on the (appropriately defined) ”translational invariant part” of the finite volume Green functions and a convergence statement analogous to (18).
The functional $`L^{\mathrm{\Lambda }_0,\mathrm{\Lambda }_0}(\widehat{\phi })`$ is the bare action including counterterms, viewed as a formal power series in $`\mathrm{}`$. Its general form for symmetric<sup>7</sup><sup>7</sup>7The necessary generalizations in the nonsymmetric case will be surveyed in the end of the next section. $`\phi _4^4`$ theory is
$$L^{\mathrm{\Lambda }_0,\mathrm{\Lambda }_0}(\widehat{\phi })=\frac{g}{4!}d^4x\widehat{\phi }^4(x)+$$
$$+d^4x\{\frac{1}{2}a(\mathrm{\Lambda }_0)\widehat{\phi }^2(x)+\frac{1}{2}b(\mathrm{\Lambda }_0)\underset{\mu =0}{\overset{3}{}}(_\mu \widehat{\phi })^2(x)+\frac{1}{4!}c(\mathrm{\Lambda }_0)\widehat{\phi }^4(x)\},$$
(10)
where $`g>0`$ is the renormalized coupling, and the parameters $`a(\mathrm{\Lambda }_0),b(\mathrm{\Lambda }_0),c(\mathrm{\Lambda }_0)`$ fulfill
$$a(\mathrm{\Lambda }_0),b(\mathrm{\Lambda }_0),c(\mathrm{\Lambda }_0)=O(\mathrm{}).$$
(11)
They are directly related to the standard mass, wave function and coupling constant counterterms. Since in the flow equation framework it is not necessary to introduce bare fields in distinction to renormalized ones (our field is the renormalized one in this language), there is a slight difference, which is to be kept in mind only when comparing to other schemes. The Wilson flow equation (FE) is obtained from (9) on differentiating w.r.t. $`\mathrm{\Lambda }`$. It is a differential equation for the functional $`L^{\mathrm{\Lambda },\mathrm{\Lambda }_0}`$ :
$$_\mathrm{\Lambda }(L^{\mathrm{\Lambda },\mathrm{\Lambda }_0}+I^{\mathrm{\Lambda },\mathrm{\Lambda }_0})=\frac{\mathrm{}}{2}\frac{\delta }{\delta \widehat{\phi }},(_\mathrm{\Lambda }\widehat{C}^{\mathrm{\Lambda },\mathrm{\Lambda }_0})\frac{\delta }{\delta \widehat{\phi }}L^{\mathrm{\Lambda },\mathrm{\Lambda }_0}\frac{1}{2}\frac{\delta }{\delta \widehat{\phi }}L^{\mathrm{\Lambda },\mathrm{\Lambda }_0},(_\mathrm{\Lambda }\widehat{C}^{\mathrm{\Lambda },\mathrm{\Lambda }_0})\frac{\delta }{\delta \widehat{\phi }}L^{\mathrm{\Lambda },\mathrm{\Lambda }_0}.$$
(12)
By $`,`$ we denote the standard scalar product in $`L_2(\mathrm{I}\mathrm{R}^4,d^4x)`$. Changing to momentum space and expanding in a formal powers series w.r.t. $`\mathrm{}`$ we write with slight abuse of notation
$$L^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\phi )=\underset{l=0}{\overset{\mathrm{}}{}}\mathrm{}^lL_l^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\phi ).$$
(13)
From $`L_l^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\phi )`$ we then obtain the CAG of loop order $`l`$ in momentum space as <sup>8</sup><sup>8</sup>8The normalization of the $`_{l,n}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}`$ is defined differently from earlier references.
$$(2\pi )^{4(n1)}\delta _{\phi (p_1)}\mathrm{}\delta _{\phi (p_n)}L_l^{\mathrm{\Lambda },\mathrm{\Lambda }_0}|_{\phi 0}=\delta ^{(4)}(p_1+\mathrm{}+p_n)_{l,n}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(p_1,\mathrm{},p_{n1}),$$
(14)
where we have written $`\delta _{\phi (p)}=\delta /\delta \phi (p)`$. Note that our definition of the $`_{l,n}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}`$ is such that $`_{0,2}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}`$ vanishes. The absence of 0-loop two (and one-) point functions is important for the set-up of the inductive scheme, from which we will prove renormalizability below. The FE (12) rewritten in terms of the CAG (14) takes the following form
$$_\mathrm{\Lambda }^w_{l,n}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(p_1,\mathrm{}p_{n1})=\frac{1}{2}_k(_\mathrm{\Lambda }C^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(k))^w_{l1,n+2}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(k,k,p_1,\mathrm{}p_{n1})$$
(15)
$$\underset{\genfrac{}{}{0pt}{}{l_1+l_2=l,w_1+w_2+w_3=w}{n_1+n_2=n}}{}\frac{1}{2}\left[^{w_1}_{l_1,n_1+1}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(p_1,\mathrm{},p_{n_1})(^{w_3}_\mathrm{\Lambda }C^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(p^{}))^{w_2}_{l_2,n_2+1}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(p_{n_1+1},\mathrm{},p_n)\right]_{ssym},$$
$$\text{where }p^{}=p_1\mathrm{}p_{n_1}=p_{n_1+1}+\mathrm{}+p_n.$$
Here we have written (15) directly in a form where also momentum derivatives of the CAG (14) are performed, and we used the shorthand notation
$$^w:=\underset{i=1}{\overset{n1}{}}\underset{\mu =0}{\overset{3}{}}(\frac{}{p_{i,\mu }})^{w_{i,\mu }}\text{ with }w=(w_{1,0},\mathrm{},w_{n1,3}),|w|=w_{i,\mu },w_{i,\mu }\mathrm{I}\mathrm{N}_0.$$
(16)
The symbol $`ssym`$ <sup>9</sup><sup>9</sup>9It is defined differently from the symbol $`sym`$ in , the present conventions being slightly more elegant. means summation over those permutations of the momenta $`p_1,\mathrm{},p_n`$, which do not leave invariant the subsets $`\{p_1,\mathrm{},p_{n_1}\}`$ and $`\{p_{n_1+1},\mathrm{},p_n\}`$. Note that the CAG are symmetric in their momentum arguments by definition. A simple inductive proof of the renormalizability of $`\phi _4^4`$ theory has been exposed several times in the literature , and we will not repeat it in detail. The line of reasoning can be resumed as follows.
The induction hypotheses to be proven are :
A) Boundedness
$$|^w_{l,n}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\stackrel{}{p})|(\mathrm{\Lambda }+m)^{4n|w|}𝒫_1(log\frac{\mathrm{\Lambda }+m}{m})𝒫_2(\frac{|\stackrel{}{p}|}{\mathrm{\Lambda }+m}).$$
(17)
B) Convergence
$$|_{\mathrm{\Lambda }_0}^w_{l,n}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\stackrel{}{p})|\frac{1}{\mathrm{\Lambda }_0^2}(\mathrm{\Lambda }+m)^{5n|w|}𝒫_3(log\frac{\mathrm{\Lambda }_0}{m})𝒫_4(\frac{|\stackrel{}{p}|}{\mathrm{\Lambda }+m}).$$
(18)
Here and in the following the $`𝒫`$ denote (each time they appear possibly new) polynomials with nonnegative coefficients. The coefficients depend on $`l,n,|w|,m`$, but not on $`\stackrel{}{p},\mathrm{\Lambda },\mathrm{\Lambda }_0`$. We used the shorthand $`\stackrel{}{p}=(p_1,\mathrm{},p_{n1})`$ and $`|\stackrel{}{p}|=sup\{|p_1|,\mathrm{},|p_n|\}`$. The statement (18) implies renormalizability, since it proves that the limits $`lim_{\mathrm{\Lambda }_0\mathrm{},\mathrm{\Lambda }0}_{l,n}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\stackrel{}{p})`$ exist to all loop orders $`l`$. But the statement (17) has to be obtained first to prove (18). The inductive scheme to prove the statements proceeds upwards in $`l`$, for given $`l`$ upwards in $`n`$, and for given $`(l,n)`$ downwards in $`|w|`$, starting from some arbitrary $`|w_{max}|3`$. The important point to note is that the terms on the r.h.s. of the FE always are prior to the one on the l.h.s. in the inductive order. So the bound (17) may be used as an induction hypothesis on the r.h.s. Then we may integrate the FE, where terms with $`n+|w|5`$ are integrated down from $`\mathrm{\Lambda }_0`$ to $`\mathrm{\Lambda }`$, since for those terms we have the boundary conditions following from (10)
$$^w_{l,n}^{\mathrm{\Lambda }_0,\mathrm{\Lambda }_0}(p_1,\mathrm{}p_{n1})=0,\text{ for}n+|w|>4,$$
(19)
whereas the terms with $`n+|w|4`$ at the renormalization point - which we choose at zero momentum for simplicity - are integrated upwards from $`0`$ to $`\mathrm{\Lambda }`$, since they are fixed by ($`\mathrm{\Lambda }_0`$-independent) renormalization conditions, fixing the relevant parameters of the theory<sup>11</sup><sup>11</sup>11The simplest choice would be to set $`a_l^R=0,b_l^R=0,c_l^R=0`$, in which case the renormalized coupling is identical to the connected four point function at zero momentum. A shift away from zero momentum would result in nonvanishing terms $`c_l^R`$, just to mention one example of more general choices. :
$$_{l,2}^{0,\mathrm{\Lambda }_0}(p)=a_l^R+b_l^Rp^2+O((p^2)^2),_{0,4}^{0,\mathrm{\Lambda }_0}(0)=g,_{l,4}^{0,\mathrm{\Lambda }_0}(0)=c_l^R,l1.$$
(20)
Symmetry considerations tell us that there are no other nonvanishing renormalization constants apart from $`a_l^R,b_l^R,c_l^R`$, and the Schlömilch or integrated Taylor formula permits us to move away from the renormalization point, treating first $`_{l,4}^{0,\mathrm{\Lambda }_0}`$ and then the momentum derivatives of $`_{l,2}^{0,\mathrm{\Lambda }_0}`$, in descending order. With these remarks on the boundary conditions, and using the bounds on the propagator and its derivatives
$$|^w_\mathrm{\Lambda }C^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(p)|\mathrm{\Lambda }^{3|w|}𝒫(|p|/\mathrm{\Lambda })e^{\frac{p^2+m^2}{\mathrm{\Lambda }^2}},$$
(21)
statement A) above is straightforwardly verified by inductive integration of the FE. Once this has been achieved statement B) follows on applying the same inductive scheme to bound the solutions of the FE, integrated over $`\mathrm{\Lambda }`$ and then derived w.r.t. $`\mathrm{\Lambda }_0`$.
## 3 Temperature independent renormalization of <br>finite temperature $`\phi _4^4`$ theory
We fix the notations recalling at the same time some basic facts about euclidean finite temperature field theory. The scalar field $`\widehat{\phi }(x)`$ becomes periodic in $`x_0`$ at finite temperature with period $`\beta =1/T`$. Correspondingly position space integrals over the zero component of the coordinates are now restricted to the compact interval $`[0,\beta ]`$. Symbols denoting finite temperature quantities will generally be underlined, thus we write
$$\underset{¯}{p}:=(\underset{¯}{p}_0,\stackrel{}{p}):=(2\pi nT,\stackrel{}{p}),n\mathrm{𝖹𝖹},_{\underset{¯}{p}}:=T\underset{n\mathrm{𝖹𝖹}}{}_{\mathrm{I}\mathrm{R}^3}\frac{d^3p}{(2\pi )^3}.$$
(22)
We also introduce the convention
$$\underset{¯}{\overset{^}{\phi }}(x):=_{\underset{¯}{p}}\underset{¯}{\phi }(\underset{¯}{p})e^{i\underset{¯}{p}x},\underset{¯}{\phi }(\underset{¯}{p})=_0^\beta 𝑑x_0_{\mathrm{I}\mathrm{R}^3}d^3x\underset{¯}{\overset{^}{\phi }}(x)e^{i\underset{¯}{p}x},$$
(23)
$$\frac{\delta }{\delta \underset{¯}{\overset{^}{\phi }}(x)}=\frac{(2\pi )^3}{T}_{\underset{¯}{p}}\frac{\delta }{\delta \underset{¯}{\phi }(\underset{¯}{p})}e^{i\underset{¯}{p}x},\frac{\delta }{\delta \underset{¯}{\phi }(\underset{¯}{p})}=\frac{T}{(2\pi )^3}_0^\beta 𝑑x_0_{\mathrm{I}\mathrm{R}^3}d^3x\frac{\delta }{\delta \underset{¯}{\overset{^}{\phi }}(x)}e^{i\underset{¯}{p}x}.$$
(24)
The regularized propagator now takes the form
$$C^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\underset{¯}{p})=\frac{1}{\underset{¯}{p}^2+m^2}\{e^{\frac{\underset{¯}{p}^2+m^2}{\mathrm{\Lambda }_0^2}}e^{\frac{\underset{¯}{p}^2+m^2}{\mathrm{\Lambda }^2}}\}.$$
(25)
The generating functional of the finite temperature CAG will be called $`L^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\underset{¯}{\phi };T)`$. In analogy with (14) we define the CAG through
$$\delta _{\underset{¯}{\phi }(\underset{¯}{p}_1)}\mathrm{}\delta _{\underset{¯}{\phi }(\underset{¯}{p}_n)}L_l^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\underset{¯}{\phi };T)|_{\underset{¯}{\phi }0}=$$
(26)
$$(\frac{T}{(2\pi )^3})^{n1}\delta _{0,\left(\underset{¯}{p}_{1,0}+\mathrm{}+\underset{¯}{p}_{n,0}\right)}\delta ^{\left(3\right)}\left(\stackrel{}{p}_1+\mathrm{}+\stackrel{}{p}_n\right)_{l,n}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\underset{¯}{p}_1,\mathrm{},\underset{¯}{p}_{n1};T).$$
At this stage we could prove renormalizability of the finite temperature theory in the same way as for the zero temperature theory. A slight difference is that the relations (20) are to be replaced by
$$_{l,2}^{0,\mathrm{\Lambda }_0}(\underset{¯}{p};T)=a_l^R(T)+b_l^{R,0}(T)\underset{¯}{p}_0^2+b_l^{R,1}(T)\stackrel{}{p}^{\mathrm{\hspace{0.17em}2}}+O(\underset{¯}{p}^4),$$
$$_{0,4}^{0,\mathrm{\Lambda }_0}(\underset{¯}{p}=0;T)=g,_{l,4}^{0,\mathrm{\Lambda }_0}(\underset{¯}{p}=0;T)=c_l^R(T),l1,$$
(27)
since the space-time $`O(4)`$-symmetry is broken down to a $`\mathrm{𝖹𝖹}_2\times O(3)`$-symmetry which demands a new renormalization condition. However we want to go beyond and prove temperature independent renormalizability, in the sense that the counterterms can be chosen temperature independent. To do so, it is advantageous to pass directly to the difference between the finite and zero temperature theories, which we will do now. Note in this respect that if we wanted to prove the renormalizability of the finite temperature theory, keeping the counterterms fixed at their zero temperature values, would not work within our scheme and procedure : The CAG would become arbitrarily divergent in $`\mathrm{\Lambda }_0`$ with increasing loop order, since integrating relevant terms from $`\mathrm{\Lambda }_0`$ to 0 (instead of integrating them from a renormalization condition fixed at $`\mathrm{\Lambda }=0`$ up to $`\mathrm{\Lambda }_0`$) gives divergent integrals. Thus we rather study the difference functions
$$𝒟_{l,n}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\{\underset{¯}{p}\}):=_{l,n}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\{\underset{¯}{p}\};T)_{l,n}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\{\underset{¯}{p}\}).$$
(28)
We only define and need the difference CAG $`𝒟_{l,n}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}`$ at the external momenta $`(\{\underset{¯}{p}\}):=(\underset{¯}{p}_1,\mathrm{},\underset{¯}{p}_{n1})`$. From the FE (15) and the analogous equation for the $`_{l,n}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\{\underset{¯}{p}\};T)`$ we can derive a FE for the $`𝒟_{l,n}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\{\underset{¯}{p}\})`$ in the following form :
$$_\mathrm{\Lambda }𝒟_{l,n}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\{\underset{¯}{p}\})=\frac{1}{2}_{\underset{¯}{k}}(_\mathrm{\Lambda }C^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\underset{¯}{k}))𝒟_{l1,n+2}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\underset{¯}{k},\underset{¯}{k},\{\underset{¯}{p}\})$$
(29)
$$+\frac{1}{2}\left\{_{\underset{¯}{k}}(_\mathrm{\Lambda }C^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\underset{¯}{k}))_{l1,n+2}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\underset{¯}{k},\underset{¯}{k},\{\underset{¯}{p}\})_k(_\mathrm{\Lambda }C^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(k))_{l1,n+2}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(k,k,\{\underset{¯}{p}\})\right\}$$
$$\underset{\genfrac{}{}{0pt}{}{l_1+l_2=l,}{n_1+n_2=n}}{}\frac{1}{2}\{\left[_{l_1,n_1+1}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\underset{¯}{p}_1,\mathrm{},\underset{¯}{p}_{n_1};T)(_\mathrm{\Lambda }C^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\underset{¯}{p}^{}))𝒟_{l_2,n_2+1}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\underset{¯}{p}_{n_1+1},\mathrm{},\underset{¯}{p}_n)\right]_{ssym}$$
$$+\left[𝒟_{l_1,n_1+1}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\underset{¯}{p}_1,\mathrm{},\underset{¯}{p}_{n_1})(_\mathrm{\Lambda }C^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\underset{¯}{p}^{}))_{l_2,n_2+1}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\underset{¯}{p}_{n_1+1},\mathrm{},\underset{¯}{p}_n)\right]_{ssym}\},$$
where again
$$\underset{¯}{p}^{}=\underset{¯}{p}_1\mathrm{}\underset{¯}{p}_{n_1}=\underset{¯}{p}_{n_1+1}+\mathrm{}+\underset{¯}{p}_n.$$
The boundary conditions we want to impose on the system $`𝒟_{l,n}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}`$ are (from the previous remarks) obviously the following ones :
$$𝒟_{l,n}^{\mathrm{\Lambda }_0,\mathrm{\Lambda }_0}(\underset{¯}{p}_1,\mathrm{},\underset{¯}{p}_{n1})=\mathrm{\hspace{0.17em}0},l,n\mathrm{I}\mathrm{N}.$$
(30)
To start the induction we also note
$$𝒟_{0,n}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\underset{¯}{p}_1,\mathrm{},\underset{¯}{p}_{n1})=\mathrm{\hspace{0.17em}0},n\mathrm{I}\mathrm{N},$$
(31)
at the tree level all difference terms $`𝒟_{0,n}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}`$ vanish. This follows from the fact that restricted to the momenta $`(\underset{¯}{p}_1,\mathrm{}\underset{¯}{p}_{n1})`$ the tree level functions $`_{0,n}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\underset{¯}{p}_1,\mathrm{}\underset{¯}{p}_{n1};T)`$ and $`_{0,n}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\underset{¯}{p}_1,\mathrm{}\underset{¯}{p}_{n1})`$ agree. Now we would like to use the same inductive scheme proceeding upwards in $`l`$, and for given $`l`$ upwards in $`n`$, to prove the finiteness of $`lim_{\mathrm{\Lambda }_0\mathrm{},\mathrm{\Lambda }0}𝒟_{0,n}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}`$. Due to the form of (30) we always integrate the FE for $`𝒟_{l,n}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}`$ from $`\mathrm{\Lambda }_0`$ down to $`\mathrm{\Lambda }`$, since the boundary terms at $`\mathrm{\Lambda }_0`$ always vanish. We want to prove the following
Theorem :
$$|𝒟_{l,n}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\underset{¯}{p}_1,\mathrm{},\underset{¯}{p}_{n1})|(\mathrm{\Lambda }+m)^{sn}𝒫_1(\mathrm{log}\frac{\mathrm{\Lambda }+m}{m})𝒫_2(\frac{|\{\underset{¯}{p}\}|}{\mathrm{\Lambda }+m}),$$
(32)
$$|_{\mathrm{\Lambda }_0}𝒟_{l,n}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\underset{¯}{p}_1,\mathrm{},\underset{¯}{p}_{n1})|\frac{1}{\mathrm{\Lambda }_0^2}(\mathrm{\Lambda }+m)^{sn}𝒫_3(log\frac{\mathrm{\Lambda }_0}{m})𝒫_4(\frac{|\{\underset{¯}{p}\}|}{\mathrm{\Lambda }+m}).$$
(33)
The nonnegative coefficients in the polynomials $`𝒫`$ depend on $`l,n,s,m`$ and (smoothly) on $`T`$, but not on $`\{\underset{¯}{p}\},\mathrm{\Lambda },\mathrm{\Lambda }_0`$. The positive integer $`s\mathrm{I}\mathrm{N}`$ may be chosen arbitrarily.
The finite temperature CAG $`_{0,n}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\underset{¯}{p}_1,\mathrm{},\underset{¯}{p}_{n1};T)`$, when renormalized with the same counterterms as the zero temperature ones, satisfy the same bounds as in (17,18) restricted to $`w=0`$. The coefficients in the polynomials $`𝒫`$ may now depend on $`l,n,m`$ and (smoothly) on $`T`$.
Remark : It is possible to prove the bounds (17,18) also for derivatives of the finite temperature CAG $`_{0,n}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\underset{¯}{p}_1,\mathrm{},\underset{¯}{p}_{n1};T)`$. In the $`p_{i,0}`$-components differentiations then have to replaced by finite differences. However these bounds are not needed in the inductive proof, so we skip them here.
Proof : We first prove (32) and and the statement corresponding to (17) for $`w=0`$, using the inductive scheme indicated previously. Regarding the FE (29) we state that it is compatible with the inductive scheme and that the only term in which (32) cannot be used as an induction hypothesis is the following one :
$$_{\underset{¯}{k}}(_\mathrm{\Lambda }C^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\underset{¯}{k}))_{l1,n+2}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\underset{¯}{k},\underset{¯}{k},\{\underset{¯}{p}\})_k(_\mathrm{\Lambda }C^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(k))_{l1,n+2}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(k,k,\{\underset{¯}{p}\}).$$
(34)
So our sharp $`\mathrm{\Lambda }`$-bound on $`𝒟_{l,n}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}`$ can only be verified if it holds for this difference term. Here we use (17,18) and the Euler-MacLaurin-formula, see e.g. . We can rewrite (34) as
$$\frac{2}{\mathrm{\Lambda }^3}\frac{d^3\stackrel{}{k}}{(2\pi )^4}e^{\frac{\stackrel{}{k}^2+m^2}{\mathrm{\Lambda }^2}}\left[2\pi T\underset{n\mathrm{𝖹𝖹}}{}g(2\pi nT)_{\mathrm{}}^{\mathrm{}}𝑑k_0g(k_0)\right],$$
(35)
where we introduced the function
$$g(k_0)=e^{\frac{k_0^2}{\mathrm{\Lambda }^2}}_{l1,n+2}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(k,k,\{\underset{¯}{p}\})\text{ for }\stackrel{}{k},\{\underset{¯}{p}\}\text{ fixed}.$$
(36)
The Euler-MacLaurin formula for our case can be stated in the form
$$2\pi T\underset{n\mathrm{𝖹𝖹}}{}g(2\pi nT)_{\mathrm{}}^{\mathrm{}}𝑑k_0g(k_0)=\pi T[g(\mathrm{})g(\mathrm{})]$$
(37)
$$+\underset{k=1}{\overset{r+1}{}}\frac{b_{2k}(2\pi T)^{2k}}{(2k)!}[g^{(2k1)}(\mathrm{})g^{(2k1)}(\mathrm{})]+R_{r+1}.$$
Here $`b_{2k}`$ are the Bernoulli numbers. We observe that passing to the limit of an infinite integration interval is justified, since the function $`g(k_0)`$ and its derivatives vanish rapidly at infinity. The remainder $`R_{r+1}`$ obeys the following bound
$$|R_{r+1}|\mathrm{\hspace{0.17em}4}e^{2\pi }T^{2r+3}_{\mathrm{}}^{\mathrm{}}𝑑k_0|g^{(2r+3)}(k_0)|,$$
(38)
therefore we obtain, using again (17,18)
$$|R_{r+1}|T^{2r+3}\frac{(\mathrm{\Lambda }+m)^{2n}}{\mathrm{\Lambda }^{2r+2}}𝒫_1(log\frac{\mathrm{\Lambda }+m}{m})𝒫_2(\frac{|\{k,\underset{¯}{p}\}|}{\mathrm{\Lambda }+m}).$$
(39)
Note that $`r\mathrm{I}\mathrm{N}`$ can be chosen arbitrarily here, and the bound for (34) is thus
$$T^{2r+3}e^{\frac{m^2}{\mathrm{\Lambda }^2}}\frac{(\mathrm{\Lambda }+m)^{2n}}{\mathrm{\Lambda }^{2r+2}}𝒫_1(log\frac{\mathrm{\Lambda }+m}{m})𝒫_2(\frac{|\{k,\underset{¯}{p}\}|}{\mathrm{\Lambda }+m})$$
(40)
$$T^{2r+3}(\mathrm{\Lambda }+m)^{2n2r2}𝒫_3(log\frac{\mathrm{\Lambda }+m}{m})𝒫_4(\frac{|\{k,\underset{¯}{p}\}|}{\mathrm{\Lambda }+m}).$$
After this preparation we consider the induction process : At each loop order we first prove (32), and then (17) for finite $`T`$ and corresponding momenta. This second step is trivial from (17,18) at $`T=0`$, from the definition (28) and from (32) <sup>12</sup><sup>12</sup>12We may choose the bounds for $`s=0`$ from (32,33) when bounding the finite temperature CAG, so that polynomials appearing in the bounds may be chosen s-independent.. We know already the theorem to be true at 0 loop order. This and the form of the FE (29) implies that we do not need a bound on any of the $`_{l,n}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\{\underset{¯}{p}\};T)`$ in the inductive bound on $`𝒟_{l,n}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}`$ at the given loop order $`l`$.
It is instructive to regard how the induction starts at loop order $`l=1`$. Treating first the case $`n=2`$ we find that the only nonvanishing contribution on the r.h.s. of the FE stems from (34), and it is momentum independent, so that integrating over $`\mathrm{\Lambda }`$ we get
$$|𝒟_{1,2}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\underset{¯}{p})|c(\mathrm{\Lambda }+m)^{2r1}$$
with a suitable constant $`c`$, depending on $`r`$. For $`n=4`$ also the last two terms on the r.h.s. of the FE contribute. Using the result for $`𝒟_{1,2}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\underset{¯}{p})`$, integration over $`\mathrm{\Lambda }`$ gives
$$|𝒟_{1,4}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\{\underset{¯}{p}\})|(\mathrm{\Lambda }+m)^{22r1}𝒫(\frac{|\{\underset{¯}{p}\}|}{\mathrm{\Lambda }+m}).$$
From this one inductively obtains the bound for $`n6`$
$$|𝒟_{1,n}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\{\underset{¯}{p}\})|(\mathrm{\Lambda }+m)^{(n2)2r1}𝒫(\frac{|\{\underset{¯}{p}\}|}{\mathrm{\Lambda }+m}).$$
Having bounded the difference functions $`𝒟_{1,n}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}`$ we can bound the CAG $`_{1,n}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(T)=_{1,n}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(T=0)+𝒟_{1,n}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}`$, see (28). Then we may proceed inductively to higher loop orders and verify the inductive bound
$$|𝒟_{l,n}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\{\underset{¯}{p}\})|(\mathrm{\Lambda }+m)^{(n2)2r1}𝒫_1(\mathrm{log}\frac{\mathrm{\Lambda }+m}{m})𝒫_2(\frac{|\{\underset{¯}{p}\}|}{\mathrm{\Lambda }+m}).$$
This proves the first part of the theorem on writing $`s=2r1`$ for $`s`$ odd, and majorizing to obtain even $`s`$. It follows that the $`_{l,n}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(T)`$ may be bounded in agreement with (17,18).
Now we turn to the proof of the statement (33) which implies convergence of the $`𝒟_{l,n}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}`$ for $`\mathrm{\Lambda }_0\mathrm{}`$. The proof is based on the same inductive scheme and starts from the FE (29) integrated over $`\mathrm{\Lambda }`$ from $`\mathrm{\Lambda }_0`$ to $`\mathrm{\Lambda }`$, and then derived w.r.t. $`\mathrm{\Lambda }_0`$. The result is of the form
$$_{\mathrm{\Lambda }_0}𝒟_{l,n}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\{\underset{¯}{p}\})=\text{[RHS of (}\text{29}\text{)]}|_{\mathrm{\Lambda }=\mathrm{\Lambda }_0}+_\mathrm{\Lambda }^{\mathrm{\Lambda }_0}𝑑\lambda _{\mathrm{\Lambda }_0}\text{[RHS of (}\text{29}\text{)]}(\lambda ),$$
(41)
and we denote the RHS of this equation shortly as
$$I_{l,n}^{\mathrm{\Lambda }_0}(\{\underset{¯}{p}\})+I_{l,n}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\{\underset{¯}{p}\}).$$
Since we have imposed $`_{l,n}^{\mathrm{\Lambda }_0,\mathrm{\Lambda }_0}(T)_{l,n}^{\mathrm{\Lambda }_0,\mathrm{\Lambda }_0}`$, and since moreover these terms vanish for $`n6`$, we find
$$I_{l,n}^{\mathrm{\Lambda }_0}(\{\underset{¯}{p}\})=\delta _{n,2}\left[_{\underset{¯}{k}}\frac{e^{\frac{\underset{¯}{k}^2+m^2}{\mathrm{\Lambda }_0^2}}}{\mathrm{\Lambda }_0^3}_k\frac{e^{\frac{k^2+m^2}{\mathrm{\Lambda }_0^2}}}{\mathrm{\Lambda }_0^3}\right]_{l1,4}^{\mathrm{\Lambda }_0,\mathrm{\Lambda }_0}.$$
(42)
Since $`_{l1,4}^{\mathrm{\Lambda }_0,\mathrm{\Lambda }_0}c_{l1}(\mathrm{\Lambda }_0),l>1`$ and $`_{0,4}^{\mathrm{\Lambda }_0,\mathrm{\Lambda }_0}g`$, see (10), we realize that (42) is momentum independent. The difference can be calculated explicitly or bounded again using the Euler-MacLaurin formula, and we obtain
$$|I_{l,n}^{\mathrm{\Lambda }_0}|\delta _{n,2}\mathrm{\Lambda }_0^{22r}𝒫(\mathrm{log}\frac{\mathrm{\Lambda }_0}{m})$$
(43)
for $`r\mathrm{I}\mathrm{N}`$ and a suitable $`𝒫`$ depending on $`r`$. To get a bound on $`I_{l,n}^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\{\underset{¯}{p}\})`$ we apply the derivative in (41) to all entries using the product rule (noting that when applied to $`_\mathrm{\Lambda }C^{\mathrm{\Lambda },\mathrm{\Lambda }_0}`$ it gives zero). In any case the derivative brings down the required factor of $`\mathrm{\Lambda }_0^2`$, either by (18), or by (33) together with the induction hypothesis. Apart from this the bound (33) is obtained similarly as (32), using in particular the Euler-MacLaurin formula for the difference term (34) derived w.r.t. $`\mathrm{\Lambda }_0`$. This proves also (33).
We end this section with two remarks on possible generalizations. First the preceding analysis can be extended to nonsymmetric $`\phi _4^4`$-theory. The action (10) then has to be replaced by
$$\stackrel{~}{L}^{\mathrm{\Lambda }_0,\mathrm{\Lambda }_0}(\widehat{\phi })=L^{\mathrm{\Lambda }_0,\mathrm{\Lambda }_0}(\widehat{\phi })+\frac{h}{3!}d^4x\widehat{\phi }^3(x)+d^4x\{\frac{1}{3!}d(\mathrm{\Lambda }_0)\widehat{\phi }^3(x)+v(\mathrm{\Lambda }_0)\widehat{\phi }(x)\}$$
(44)
with the tree level three-point coupling $`h`$ and $`\mathrm{\Lambda }_0`$-dependent parameters
$$d(\mathrm{\Lambda }_0),v(\mathrm{\Lambda }_0)=O(\mathrm{})$$
(45)
implementing the counterterms necessary to renormalize the one- and three-point functions. Correspondingly we pose additional renormalization conditions
$$_{l,1}^{0,\mathrm{\Lambda }_0}=v_l^R,_{l,3}^{0,\mathrm{\Lambda }_0}(0)=d_l^R\text{ for }l\mathrm{I}\mathrm{N},$$
(46)
to be joined to (20). Then the bounds (17,18) hold again, but are no more trivially fulfilled for $`n`$ odd.<sup>13</sup><sup>13</sup>13These bounds can be improved by replacing $`n`$ by $`\widehat{n}`$, defined to be the smallest even integer greater or equal to $`n`$. Once the theory at $`T=0`$ is bounded, the differences (28) again yield the theory at $`T>0`$. Bounds corresponding to (32,33) are proven proceeding as before, in the symmetric case.
As a second remark, we point out that for the existence of the large cutoff limit $`\mathrm{\Lambda }_0\mathrm{}`$, it is not necessary that the relevant coupling constants are subject to normalization conditions at zero temperature. Equally well we can impose normalization conditions at some temperature $`T_0>0`$. We pointed out that at finite temperature the space-time $`O(4)`$-symmetry is broken down to $`\mathrm{𝖹𝖹}_2\times O(3)`$. Then the 3 independent renormalization constants $`a^R`$, $`b^R`$ and $`c^R`$ at $`T=0`$, (20), become replaced by four parameters $`a^R(T_0)`$, $`b^{R,0}(T_0)`$, $`b^{R,1}(T_0)`$ and $`c^R(T_0)`$ at $`T_0`$, cf. (27), corresponding to four relevant couplings. However, starting from an $`O(4)`$-symmetric zero temperature theory we have proved that
$$L^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\underset{¯}{\phi };T_0)L^{\mathrm{\Lambda },\mathrm{\Lambda }_0}(\underset{¯}{\phi })$$
(47)
has the properties of an irrelevant operator. This implies that for given $`b^{R,0}(T_0)`$ there is a unique choice for $`b^{R,1}(T_0)`$, or vice versa, such that the finite temperature theory stems from an $`O(4)`$-symmetric zero temperature theory. Any different choice would be associated to a zero temperature theory, where $`O(4)`$-symmetry is broken by hand through the renormalization conditions. Note that the $`O(4)`$-symmetric choice is generally not the one where $`b^{R,0}(T_0)=b^{R,1}(T_0)`$: Integration over $`\mathrm{\Lambda }`$, starting from the same counterterms (the $`O(4)`$-symmetric ones) will lead to a finite difference at $`\mathrm{\Lambda }=0`$, since $`O(4)`$-invariance is broken in the propagator. Otherwise stated, the fact that the finite temperature theory stems from an $`O(4)`$-symmetric zero temperature theory, can be simply recognized on inspection of the counterterms, but not on the renormalization conditions.
## 4 Summary
We have presented a proof for the perturbative renormalizability of massive finite temperature $`\phi _4^4`$-theory. The starting point are the bounds (17,18) which prove the renormalizability of the zero temperature theory. In the flow equation framework they serve at the same time as induction hypotheses for the inductive proof. Bounds of this type have by now been rigorously established for nearly all theories of physical interest, including gauge theories, where the restoration of the Ward identities in the final theory pose an additional problem, to be solved by a suitable restriction on the renormalization conditions. Taking due care of the exceptional momentum problem, corresponding bounds can also be established for theories with massless particles.
To extend the bounds to the corresponding finite temperature theories presents no really new problems for the practitioner. The main problem to be solved rather is that the existence of the correlation functions in the large cutoff limit should be proven without changing the counterterms. In our setup this corresponds to posing the boundary conditions (30) for the difference Green functions $`𝒟`$ between the $`T>0`$ and the $`T=0`$ theories. The anounced result is contained in the bounds (32,33). The main new technical tool used to get there is the Euler-MacLaurin formula, generalized to an infinite integration interval for a rapidly decaying integrand. It is applied to the difference terms appearing in the flow equations for the functions $`𝒟`$ that are not bounded by the induction hypothesis alone, (see (34)- (40)). Here it comes to our help that the bounds (17,18) are sufficiently powerful so as to transform momentum derivatives into negative powers of $`\mathrm{\Lambda }`$. Via the Euler-MacLaurin formula it is then possible to gain an arbitrary power in $`\mathrm{\Lambda }`$ paying the corresponding power in $`T`$ (see 39). This achieves (far more than) showing that all difference functions $`𝒟`$ are irrelevant. For the latter a gain of a power of $`\mathrm{\Lambda }^{2+\epsilon }`$ would have sufficed. We emphasize again that our result agrees with the experience and intuition gained from explicit perturbative calculations.
Renormalization is a central issue that is strongly related to the fundamental principles of local quantum field theory. Renormalizability of a field theory gives it a meaning beyond some low energy effective model. The techniques we have presented here for proving renormalizability of a field theory at finite temperature mainly rely on two properties. The first property is renormalizability at zero temperature. The second one is that the difference between the theory at finite and zero temperature act like an irrelevant operator that does not spoil renormalizability. Renormalization group flow equations provide an appropriate tool to put this statement on a strong basis and prove renormalizability for finite temperature. We expect that these methods generalize appropriately to apply to more realistic and complex field theories such as QCD, where both the UV and the IR scale problem are to be attacked.
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# Recent Progress in Neutron Star Theory
## 1 Introduction
Neutron stars are among the most fascinating bodies in our universe. They contain over a solar mass of matter within a radius of $``$ 10 km at densities of order $`10^{15}`$ g/cc. They probe the properties of cold matter at extremely high densities, and have proven to be fantastic test bodies for theories of general relativity. In a broader perspective, neutron stars and heavy ion collisions provide access to the phase diagram of matter at extreme densities and temperatures, that is basic for understanding the very early Universe and several other astrophysical phenomena.
The discovery of the neutron by Chadwick in 1932 prompted Landau to predict the existence of neutron stars. The first theoretical calculations of neutron stars were performed by Oppenheimer and Volkoff in 1939 assuming that they are gravitationally bound states of neutron Fermi gas. The calculated stars had a maximum mass of $``$ 0.7 $`M_{}`$, central densities up to $`6\times 10^{15}`$ g/cm<sup>3</sup> and radii $`10`$ km. For comparison the density of nuclear matter inside a large nucleus like <sup>208</sup>Pb is $``$ 0.16 nucleons/fm<sup>3</sup>, i.e. $``$ 2.7 $`\times 10^{14}`$ g/cm<sup>3</sup> . Their predicted maximum mass was less than the Chandrasekhar mass limit of $``$ 1.4 $`M_{}`$ for white dwarfs made up of iron group nuclei, and having densities up to $`10^9`$ g/cm<sup>3</sup> . The pressure to balance the gravitational attraction in white dwarfs and Oppenheimer-Volkoff neutron stars is supplied by degenerate electron and neutron Fermi gases respectively.
In 1934 Baade and Zwicky suggested that neutron stars may be formed in supernovae in which the iron core of a massive star exceeds the Chandrasekhar limit and collapses. The large amount of energy released in the collapse blows away the rest of the star and the collapsed core may form a neutron star. For efficient production of neutron stars with this mechanism, the maximum mass of neutron stars should exceed 1.4 $`M_{}`$. In the 60’s, using schematic models of nuclear forces, Tsuruta and Cameron showed that they could increase the neutron star masses beyond 1.4 $`M_{}`$.
Bell and Hewish discovered radio pulsars in 1967, and they were soon identified as rotating neutron stars by Gold . The subsequent detection of the Crab pulsar in the remnant of the Crab supernova, observed in China in 1054 A.D., confirmed the link to supernovae, and initiated the present efforts to better understand neutron stars.
### 1.1 A Brief Overview of Observations
Almost 1200 pulsars have been discovered by the turn of this millennium. In these stars the magnetic and rotational axes are misaligned, thus they emit dipole radiation in the form of radio waves that appear to pulse on and off like a lighthouse beacon as the pulsar beam sweeps across the Earth. The rotational energy loss due to dipole radiation is
$$\dot{E}=I\mathrm{\Omega }\dot{\mathrm{\Omega }}=\frac{B^2R^6\mathrm{\Omega }^4\mathrm{sin}^2\theta }{6c^3},$$
(1)
where the moment of inertia for a typical neutron star is $`I10^{45}`$ g cm<sup>2</sup>. Pulsars have magnetic fields $`B`$ of $`10^{12}`$ G, deduced from the observed $`\dot{\mathrm{\Omega }}`$, and independently confirmed by cyclotron absorption lines found in X-ray spectra. Their periods, $`P=2\pi /\mathrm{\Omega }`$, ranging from 1.5 ms to 8.5 s, are increasing with derivatives $`\dot{P}10^{12}10^{21}`$. The pulsar age is approximately given by $`P/2\dot{P}`$ ; most pulsars are old and slowly rotating with relatively small period derivatives, except for a few young pulsars, e.g., those found in the Crab and Vela nebulae.
In 1969 the Crab and the Vela pulsars were observed to ”glitch”, i.e. to suddenly speedup with period changes $`\mathrm{\Delta }P/P`$ of the order of $`10^8`$ and $`10^6`$ respectively . In post-glitch relaxation most of the period increase $`\mathrm{\Delta }P`$ decays. These pulsars have glitched several times since then. The glitches suggest that the neutron stars have a solid crust containing superfluid neutrons. The interesting structure of their crust has been recently reviewed , and we discuss it rather briefly in this report.
The first binary of two pulsars was found by Hulse and Taylor in 1973 and they could determine many of its parameters including both masses, orbital period and period derivative, orbital distance and inclination. General relativity could be tested to an unprecedented accuracy by measuring the inward spiralling of the neutron stars in the Hulse-Taylor binary PSR 1913+16 . The periastron advance in PSR 1913+16 is 4.2 per year as compared to 43” per century for Mercury, which originally was used by Einstein to test his theory of general relativity. Six double neutron star binaries are known so far, and neutron stars in all of them have masses in the range $`1.36\pm 0.08M_{}`$ . They confirm that nuclear forces have a large effect on the structure of neutron stars and increase their maximum mass beyond 1.4 $`M_{}`$. Neutron stars are estimated to have a binding energy of $`10\%`$ of their mass. Thus $`1.5`$ $`M_{}`$ of nuclei are needed to obtain a $`1.35`$ $`M_{}`$ star.
A distinct subclass of radio pulsars are millisecond pulsars with periods $`\stackrel{<}{}100`$ ms. The fastest pulsar known has a period of 1.56 ms . The period derivatives of millisecond pulsars are very small corresponding to low magnetic fields $`10^810^{10}`$G. They are believed to be recycled pulsars, i.e. old pulsars that have been spun up by mass accretion whereby the magnetic fields have decayed. About 80% of the millisecond pulsars are in binaries whereas less than 1% of normal radio pulsars are in binaries. About 20 - almost half of the millisecond pulsars - are found in binaries where the companion is either a white dwarf or a neutron star.
With X-ray detectors on board satellites since the early 1970’s about two hundred X-ray pulsars and bursters have been found of which the rotational period has been determined for about sixty. The X-ray pulsars and bursters are believed to be neutron stars accreting matter from high ($`M\stackrel{>}{}10M_{}`$) and low mass ($`M\stackrel{<}{}1.2M_{}`$) companions respectively. The X-ray pulses are attributed to strong accretion on the magnetic poles emitting X-rays (as northern lights). The observed radiation is pulsed with the rotational frequency of the accreting star. X-ray bursts are thermonuclear explosions of accreted matter on the surface of neutron stars. After accumulating hydrogen on the surface for hours, pressure and temperature become sufficient to trigger a runaway thermonuclear explosion seen as an X-ray burst that lasts a few seconds . Masses of these stars are less accurately measured than for binary pulsars. We mention recent mass determinations for the X-ray pulsar Vela X-1: $`M=1.87_{0.17}^{+0.23}M_{}`$ , and the burster Cygnus X-2: $`M=1.8\pm 0.4)M_{}`$ . They are larger than the typical $`1.36\pm 0.08M_{}`$ masses found in pulsars binaries, presumably due to accreted matter.
A subclass of half a dozen anomalous X-ray pulsars has been discovered. They are slowly rotating, $`P10`$ sec, but rapidly slowing down. This requires huge magnetic fields of $`B10^{14}`$ G and they have appropriately been named “magnetars” . Four gamma ray repeaters discovered so far are also believed to be slowly rotating neutron stars. The magnetars and likely also the gamma ray repeaters reside inside supernova remnants.
Recently, quasi-periodic oscillations (QPO) have been found in 12 binaries of neutron stars with low mass companions. If the QPO originate from the innermost stable orbit of the accreting matter, their observed values imply that the accreting neutron stars have masses up to $`2.3M_{}`$. In this case the QPO’s also constrain the radii of the accreting star.
Non-rotating and non-accreting neutron stars are virtually undetectable but the Hubble space telescope has observed one thermally radiating neutron star . Its surface temperature is $`T6\times 10^5`$ K$`50`$ eV and its distance is less than 120 pc from Earth. Circumstantial evidence indicate a distance of $`80`$ pc which leads to a radius of 12-13 km for this star. In recent years much effort has been devoted to measuring pulsar temperatures, especially with the Einstein Observatory and ROSAT. Surface temperatures of a few pulsars have been measured, and upper limits have been set for many .
¿From the human point of view supernova explosions are rare in our and neighboring galaxies. The predicted rate is 1-3 per century in our galaxy and the most recent one was 1987A in LMC. No neutron star associated with this explosion has been detected; however, 19 neutrinos were detected on earth from 1987A , indicating the formation of a “proto-neutron star”. It has been suggested by Bethe and Brown that an upper limit to the mass of neutron stars can be obtained assuming that the remnant of SN 1987A collapsed into a black hole.
Astrophysicists expect a large abundance of $`10^8`$ neutron stars in our galaxy. At least as many supernova explosions, responsible for all heavier elements present in our Universe today, have occurred. The scarcity of neutron stars in the solar neighborhood may be due to production of black holes or other remnants in supernovae, or due to a high initial velocity (asymmetric “kick”) received during their birth in supernovae. Recently, many neutron stars have been found far away from their supernova remnants; and of the $`1200`$ discovered radio pulsars only about $`10`$ can be associated with the 220 known supernova remnants. Neutron stars thrown out of the galactic plane may be detected by gravitational microlensing experiments designed to search for dark massive objects in the galactic halo.
The recent discovery of afterglow in Gamma Ray Bursters (GRB) allows determination of their very high redshifts ($`z1`$). They imply that GRB occur at enormous distances. Evidence for beaming has been observed , and the estimated energy output is $`10^{53}`$ ergs. Such enormous energies can be produced in neutron star mergers eventually forming black holes. From abundance of binary pulsars one can estimate the rate of neutron star mergers; it is compatible with the rate of GRB of approximately one per day. Another possible mechanism, is a special class of type Ic supernova (hypernovae) where cores collapse to black holes .
The future of neutron star observations looks bright as new windows are about to open. A new fleet of X- and Gamma-ray satellites have and will be launched. With upgraded ground based observatories and detectors for neutrinos and gravitational waves our knowledge of neutron star properties will be greatly improved.
### 1.2 Theory of Neutron Star Matter
Neutron stars are made up of relatively cold, charge neutral matter with densities up to $`7`$ times the equilibrium density $`\rho _0`$ = 0.16 nucleons/fm<sup>3</sup> of charged nuclear matter in nuclei. The matter density is $`>\rho _0`$ over most of the star, apart from the relatively thin crust . The Fermi energy of neutron star matter is in excess of tens of MeV, and hence, at typical temperatures of $`\stackrel{<}{}`$ KeV, thermal effects are a minor perturbation on the gross structure of the star.
Matter at such densities has not yet been produced in the laboratory, its properties must be theoretically deduced from the available terrestrial data with guidance from observed neutron star properties. The quantities of interest are the phase and composition of cold catalyzed neutral dense matter, its energy density $`ϵ(\rho )`$ and pressure $`P(\rho )`$, where $`\rho `$ denotes the baryon number density. The baryon number is conserved in all known interactions, therefore it is convenient to find the composition by minimizing the total energy $`E_T(\rho )`$ per baryon, including rest mass contributions. This gives:
$$ϵ(\rho )=\rho E_T(\rho ),P(\rho )=\rho ^2\frac{E_T(\rho )}{\rho }.$$
(2)
The equation of state (EOS) $`P(ϵ)`$ is found by eliminating $`\rho `$ from the above two.
The gravitational equilibrium of a nonrotating star is described by the Tolman-Oppenheimer-Volkoff (TOV) Eq:
$$\frac{dP(r)}{dr}=\frac{G(ϵ(r)+P(r)/c^2)(m(r)+4\pi r^3P(r)/c^2)}{r^2(12Gm(r)/rc^2)},$$
(3)
where $`G`$ is the gravitational constant, $`P(r)`$ and $`ϵ(r)`$ are the pressure and mass density at radius $`r`$ in the star, and
$$m(r)=_0^r4\pi r^2ϵ(r^{})𝑑r^{},$$
(4)
is the mass inside $`r`$. If we neglect the general relativistic corrections of order $`1/c^2`$ the TOV Eq. reduces to the Newtonian hydrodynamic equation. The TOV Eq. can be easily integrated starting from the central density $`ϵ_c`$ ar $`r=0`$ to find the density profile $`ϵ(r)`$. At the radius $`R`$ of the star $`P(R)=0`$, and $`m(R)=M`$ is the mass of the star as seen from outside. The stability of the star can be deduced from the $`M(ϵ_c)`$ as discussed in , and the equations for rotating stars are given by . The effect of rotation on the structure of most observed neutron stars seems to be rather small, however, it could be significant at periods less than a millisecond .
At densities $`<2\times 10^3\rho _0`$ matter is believed to have the form of a lattice of nuclei in a relativistic degenerate electron gas , qualitatively similar to that of metals. The main focus of the theory reviewed here has been on determining the properties and EOS of matter in the density range $`2\times 10^3\rho _0<\rho <10\rho _0`$ from terrestrial data. In the lower part of this range we expect to find nucleon matter (NM) composed of nucleons and electrons. In contrast to matter in nuclei, it has mostly neutrons with a small fraction of protons and equal number of electrons to maintain charge neutrality. The large Fermi energy, $`\mu _e`$ 100 MeV, of the electron gas limits the fraction of protons in NM.
At higher densities there are several possibilities including condensation of negatively charged pions and kaons, occurrence of hyperons, and the transition from hadronic to quark matter. All these possibilities exploit the large electron Fermi energy of NM, therefore only one of these, if any, may occur and lower the $`\mu _e`$. In addition, neutron star matter can have interesting mixed phase regions in which the mixing phases are charged but the matter is overall neutral .
We begin with a review of NM, and later consider the more exotic possibilities. In the last sections the range of neutron star structures predicted by theory is presented along with a comparison with the observational data.
## 2 Energy-Density Functionals of Nucleon Matter
The simplest description of nuclei is obtained within the mean field approximation. It assumes, following the nuclear shell model, that nucleons occupy single particle orbitals in an average potential well produced by nuclear forces. The energy of the nucleus is assumed to be a functional of the orbitals occupied by the nucleons, and the orbitals are determined variationally as in the Hartree-Fock approximation. In reality the mean field approximation is not exactly valid for nuclei. The observed proton knockout reaction rates indicate that the shell model orbitals are occupied with a probability of $`70\%`$ in the simplest closed shell nuclei like <sup>208</sup>Pb. The differences between the real and the mean field wave-functions, due to correlations induced by nuclear forces, are subsumed in the energy functional as suggested by Kohn in the context of atomic and molecular physics.
The energy density of hypothetical, uniform NM at zero temperature is the main term in the energy functionals. The nucleon orbitals in uniform matter are simple plane waves, and the ground state in mean field approximation is obtained by filling the proton and neutron states up to their Fermi momenta $`k_{F,N}=(3\pi ^2\rho _N)^{1/3}`$, where $`N=n,p`$ for neutrons and protons. The energy density, denoted by $`(\rho _n,\rho _p)`$, includes kinetic and strong interaction contributions, but excludes rest masses and the Coulomb interaction, which destabilizes uniform charged matter. The total density is denoted by $`\rho =\rho _n+\rho _p`$, the asymmetry of the matter is defined as $`\beta =(\rho _n\rho _p)/\rho `$, and the energy per nucleon, $`E(\rho ,\beta )`$, is given by $`(\rho ,\beta )/\rho `$.
Analysis of nuclear properties with the liquid drop model reveals that, in the absence of electromagnetic forces, the ground state of NM is symmetric $`(i.e.\beta =0)`$, has total equilibrium density $`\rho _0=0.16\pm 0.01`$ fm<sup>-3</sup>, and binding energy $`E_0=16\pm 0.5`$ MeV per nucleon. The symmetry energy $`E_{sym}(\rho _0)=34\pm 6`$ MeV, is defined as $`\frac{1}{2}^2E/\beta ^2`$ at equilibrium. The NM energy $`E(\rho ,\beta )`$ can be expanded about its minimum value at $`\beta =0`$ in powers of $`\beta ^2`$, assuming charge symmetry of nuclear forces. In variational as well as Brueckner theories the coefficients of terms with $`\beta ^{n4}`$ are estimated to be small, and $`E(\rho ,\beta )(1\beta ^2)E(\rho ,0)+\beta ^2E(\rho ,1)`$. In this approximation the symmetry energy is the difference between the energy of pure neutron matter and symmetric nuclear matter. The incompressibility $`K_0=240\pm 30`$ MeV of symmetric nuclear matter is defined as $`K_0=k_F^2^2E/k_F^2`$ at equilibrium. The energies of the collective breathing mode vibrations of nuclei are sensitive to $`K_0`$; however, in all stable nuclei the surface effects are significant. It is difficult to extract the density and $`\beta `$-dependence of the incompressibility, and the density-dependence of the symmetry energy from available nuclear data.
Analysis of elastic scattering of nucleons by nuclei shows that the nuclear mean field has a dependence on the energy of the moving nucleon . Over a wide energy range this dependence is approximately linear, suggesting that nucleons in equilibrium nuclear matter have an effective mass $`m^{}0.7m`$, where $`m`$ is the free nucleon mass. This effective mass should not be identified with the Landau effective mass which describes the density of single particle states in a narrow energy interval about the Fermi energy . The Landau $`m^{}`$ in uniform matter is difficult to extract from nuclear data, since nucleons at the Fermi energy are strongly coupled to nuclear surface dynamics. Some of the phenomenological energy functionals are chosen to fit the observed nuclear level densities at the Fermi energy, while others fit the value of $`m^{}(\rho _0,0)`$ obtained from the energy dependence of the optical model potential .
The nonrelativistic functionals based on Skyrme effective interactions generally contain the following terms:
$$(\rho _n,\rho _p)=\tau (1+x_5\rho )+x_1\rho ^2(1+x_2\rho ^\alpha )+\underset{N=n,p}{}\left[x_6\tau _N\rho _N+x_3\rho _N^2(1+x_4\rho ^\alpha )\right].$$
(5)
Here $`\tau _N=0.6k_{F,N}^2\rho _N/m`$ are the kinetic energy densities, and $`\tau =\tau _n+\tau _p`$. The parameters $`x_1`$ to $`x_4`$ and $`\alpha `$ describe the $`\rho `$ and $`\beta `$ dependence of the volume integral of the static part of the effective interaction between nucleons in matter, while the $`x_5`$ and $`x_6`$ describe effective masses produced by the momentum dependence of the effective interaction. In principle the values of the seven parameters in a typical Skyrme functional are constrained by the empirically known values of $`\rho _0,E_0,E_{sym}(\rho _0)`$ and $`K`$ and the choice made for $`m^{}(\rho _0,0)`$. However, since the constraints are insufficient, there are many Skyrme models of the energy functional.
The simple form of the functional (Eq.5) chosen by most Skyrme models is convenient, but the real functional can be much more complex. The analytic form of the energy density predicted by realistic models of nuclear forces, as discussed in the next chapter, has been studied by Ravenhall . A much more elaborate function of the type:
$`(\rho _n,\rho _p)=\rho ^2\left[p_1e^{p_6\rho }+p_2(1e^{p_6\rho })+\left({\displaystyle \frac{p_{10}}{\rho }}+p_{11}\right)e^{(p_9\rho )^2}\right]`$
$`{\displaystyle \frac{1}{4}}(\rho _n\rho _p)^2\left[p_7e^{p_6\rho }+p_8(1e^{p_6\rho })+\left({\displaystyle \frac{p_{12}}{\rho }}+p_{13}\right)e^{(p_9\rho )^2}\right]`$
$`+{\displaystyle \underset{N=n,p}{}}\tau _N\left[1+(p_3\rho +p_5\rho _N)e^{p_4\rho }\right],`$ (6)
is required to reproduce the predicted $`(\rho _n,\rho _p)`$ up to $`\beta =1`$ and $`\rho 1`$ fm<sup>-3</sup>. This functional also explains the nuclear binding energies and the empirically known values for symmetric nuclear matter , however, it is unlikely that the values of all of its thirteen parameters can be obtained by fitting nuclear data.
The energy of NM can be easily calculated from a covariant effective Lagrangian in the mean field approximation, as shown by Walecka , and in the past decade many properties of medium and heavy nuclei have been studied with this approach . The effective Lagrangian used in the recent work has the form:
$`=\overline{\psi }\left[\gamma ^\mu \left(i_\mu g_\omega \omega _\mu g_\rho \stackrel{}{\tau }\stackrel{}{\rho }_\mu \right)mg_\sigma \sigma \right]\psi {\displaystyle \frac{1}{2}}m_\sigma ^2\sigma ^2+{\displaystyle \frac{1}{3}}g_2\sigma ^3+{\displaystyle \frac{1}{4}}g_3\sigma ^4`$
$`+{\displaystyle \frac{1}{2}}m_\omega ^2\omega ^\mu \omega _\mu +{\displaystyle \frac{1}{2}}m_\rho ^2\stackrel{}{\rho }^\mu \stackrel{}{\rho }_\mu {\displaystyle \frac{1}{4}}\mathrm{\Omega }^{\mu \nu }\mathrm{\Omega }_{\mu \nu }{\displaystyle \frac{1}{4}}\stackrel{}{R}^{\mu \nu }\stackrel{}{R}_{\mu \nu },`$ (7)
Hear $`\psi ,\omega _\mu `$ and $`\stackrel{}{\rho }_\mu `$ are respectively the nucleon and $`\omega `$ and $`\rho `$ vector-meson fields. Overhead arrows are used to denote isospin vectors, and $`\mathrm{\Omega }^{\mu \nu }=^\mu \omega ^\nu ^\nu \omega ^\mu `$ etc. The effective scalar field $`\sigma `$ is responsible for nuclear binding, and the $`\sigma ^3`$ and $`\sigma ^4`$ terms are necessary to obtain the empirical incompressibility of nuclear matter . The isovector $`\stackrel{}{\rho }`$ field is required to obtain the empirical symmetry energy. The observed values of the masses $`m,m_\omega `$ and $`m_\rho `$ are used, and the coupling constants $`g_\omega ,g_\rho ,g_\sigma ,g_2`$ and $`g_3`$, as well as the mass $`m_\sigma `$ of the effective scalar field are adjusted to fit the nuclear data. The above Lagrangian, without the $`\sigma ^3`$ and $`\sigma ^4`$ terms but including pion fields and their coupling to the nucleon, is also used to model the two-nucleon interaction discussed in the next section.
The relativistic mean field theory of nuclei is very elegant and often used to study properties of neutron star matter . It has provided important insights into relativistic effects in nuclei and NM. However, the effective mean-field Lagrangian (Eq.7) is unlikely to have a simple physical meaning. The inverse masses of the vector and scalar fields correspond to lengths of $``$ 0.25 and 0.4 fm, which are much smaller than the unit radius $`r_0=(4\pi \rho /3)^{1/3}1.2`$ fm for equilibrium nuclear matter. The naive condition for the validity of the mean field approximation, that $`r_0`$ be much less than the inverse masses of the fields is totally violated in nuclei as well as in neutron stars. Pions are omitted from the effective Lagrangian because they do not contribute to the energy of matter in the mean field approximation. Their higher order contributions are subsumed in the effective scalar field. Therefore, the effective mean-field Lagrangian must be interpreted as a relativistic generalization of Kohn’s energy functional. Eq. (7) assumes the simplest form necessary to fit nuclear data. A more general form, necessary to explain the properties of NM over the wide density-asymmetry range in neutron stars, can have additional fields, isovector scalar for example, density dependent coupling constants to take into account the changes in correlations with density, and field energies containing high powers of the fields, etc.
The energy of low density neutron matter is well determined by realistic models of two-nucleon interaction obtained by fitting the nucleon-nucleon (NN) scattering data. Different models and methods of calculation give very similar results up to $`\rho \rho _0`$, beyond which three-nucleon interactions and relativistic effects, as well as computational difficulties may become appreciable. These energies thus provide a test of the ability of the Skyrme and relativistic mean field theories to find neutron matter properties by extrapolating data on nuclear binding energies, sizes, vibrations, etc. As shown in Fig.1, the neutron matter energies predicted by the various functionals are widely different, and not in agreement with the results of many-body calculations at $`\rho <\rho _0`$. It thus appears likely that the simple forms of effective interactions or Lagrangians used in the present mean field theories are inadequate to predict the properties of neutron star matter by extrapolating the observed nuclear properties. Nevertheless, effective mean-field Lagrangians have been widely used in neutron star studies due to their simplicity .
## 3 Many-Body Theory of Nucleon Matter
Many properties of nuclei and nuclear matter can be understood from the Hamiltonian:
$$H=\underset{i}{}\frac{1}{2m}_i^2+\underset{i<j}{}v_{ij}+\underset{i<j<k}{}V_{ijk}+\mathrm{},$$
(8)
which includes the kinetic energy, two-nucleon interactions denoted by lower case $`v`$ and three-nucleon interactions by capital $`V`$. The ellipsis denote neglected four and higher body interactions. In this section we review the present status of this approach and its limitations.
### 3.1 Models of Two Nucleon Interaction
Our understanding of QCD has not yet progressed enough to predict the two-nucleon interaction $`v_{ij}`$ ab initio. The long range part of $`v_{ij}`$ is known to be mediated by pions, the lightest of all the mesons, and it is denoted by the one pion exchange potential (OPEP) given by:
$`v_{ij}^\pi `$ $`=`$ $`{\displaystyle \frac{f_{\pi NN}^2}{4\pi }}{\displaystyle \frac{m_\pi }{3}}X_{ij}𝝉_i𝝉_j,`$ (9)
$`X_{ij}`$ $`=`$ $`Y_\pi (r_{ij})𝝈_i𝝈_j+T_\pi (r_{ij})S_{ij}.`$ (10)
The pion-nucleon coupling constant $`f_{\pi NN}^2/4\pi =0.075\pm 0.002`$ , and the radial functions associated with the spin-spin and the tensor parts are:
$$Y_\pi (r)=\frac{e^x}{x}\xi _Y(r),T_\pi (r)=\left(1+\frac{3}{x}+\frac{3}{x^2}\right)Y_\pi (r)\xi _T(r),$$
(11)
where $`x=m_\pi r`$. The tensor operator, $`S_{ij}=3𝝈_i\widehat{𝐫}_{ij}𝝈_j\widehat{𝐫}_{ij}𝝈_i𝝈_j`$, and since $`T_\pi (r)Y_\pi (r)`$ in the important $`x\stackrel{<}{}1`$ region, the OPEP is dominated by its tensor part. The complete $`v_{ij}`$ is expressed as $`v_{ij}^\pi +v_{ij}^R`$, where $`v_{ij}^R`$ contains all the other, heavy meson, multiple meson and quark exchange parts. It has to be obtained along with the short range cutoffs $`\xi _Y(r)`$ and $`\xi _T(r)`$ in the OPEP by fitting NN scattering data. In boson-exchange models the $`v^R`$ is approximated by a sum of attractive scalar meson, and repulsive vector meson exchange potentials, while other models use attractive two-pion exchange and repulsive core potentials. The OPEP contains a $`\delta `$function term omitted from Eq.(10). Due to the finite size of nucleons this term acquires a finite range, and is difficult to separate from $`v^R`$.
In the early 1990’s the Nijmegen group carefully examined all the data on elastic NN scattering, at energies below the pion production threshold of $``$ 350 MeV, published between 1955 and 1992. They extracted 1787 proton-proton and 2514 proton-neutron “reliable” data, and showed that these could determine all NN scattering phase shifts and mixing parameters quite accurately, thus claiming that the experimental information on elastic NN scattering is now complete. Additional measurements are being carried out at several laboratories including the Indiana University Cyclotron Facility, and the CELCIUS facility in Uppsala, Sweden, to test the accuracy of the claim and improve on the quality of the data base. Nevertheless the Nijmegen analysis has been a major step.
NN interaction models which fit the Nijmegen data base with a $`\chi ^2/N_{data}1`$ are called “modern”. These include the Nijmegen models called Nijmegen I, II and Reid 93, the Argonne $`v_{18}`$ (A18) and CD-Bonn . In order to fit both the proton-proton and neutron-proton scattering data simultaneously and accurately, these models include a detailed description of the electromagnetic interactions and terms that violate the isospin symmetry of the strong interaction via the differences in the masses of the charged and neutral pions, etc.
These five models use different parameterizations of the $`v_{ij}^R`$, and the Nijmegen-I and CD-Bonn also include nonlocalities suggested by boson-exchange representations. Thus, like the older models, they make different predictions for the many-body systems. However, the differences in their predictions are much smaller than those between older models, and can be partly understood.
The interaction in the spin-isospin $`T,S=0,1`$ has the largest model dependence, which gets carried over to the estimates of the energy of symmetric nuclear matter (sec. 3.4). Fortunately the deuteron structure provides significant information on this interaction. Fig.2 shows the deuteron wave functions obtained with the five modern potentials . The three potentials, Reid 93, Nijmegen II and A18, which are local in each NN partial wave, give essentially the same deuteron wave function. We expect that they will give rather similar matter properties. The Nijmegen I and CD-Bonn potentials have momentum dependent terms associated with heavy meson exchange. These two potentials give larger $`{}_{}{}^{3}S_{1}^{}`$ wave functions at $`r<0.8`$ fm, because they have softer repulsive cores than the local models; however, this effect is not very large as can be seen from Fig.2. At $`r>0.8`$ fm only the CD-Bonn predictions differ from the rest. This is because CD-Bonn has a strongly nonlocal OPEP as suggested by pseudoscalar pion-nucleon coupling, which suppresses the $`{}_{}{}^{3}D_{1}^{}`$ wave function.
These results indicate that the main difference between the preset models of $`v_{ij}`$ is from the assumed nonlocality in OPEP. However, relativistic field theories permit use of OPEP with different nonlocalities related with the Dyson transformation. The three-nucleon interaction, $`V_{ijk}`$, depends upon the choice of OPEP , and the final results obtained after including it should be independent of the choice. Therefore, if relativistic field theories can be used to describe pion exchange forces, one can use either the local OPEP in A18, or the nonlocal one in CD-Bonn. In this case it may be better to use the local representation because it is simpler, and more accurate many-body calculations can be carried out with it.
Realistic models of nuclear forces, the modern as well as the older, predict the existence of a dense toroidal inner core in the deuteron . The density distribution, $`\rho (𝐫^{})`$, of the deuteron in spin projection $`M=0`$ state, in the center of mass frame is shown in Fig.3. Here $`𝐫^{}=\pm 𝐫/2`$ are the nucleon coordinates in the center of mass frame, and $`𝐫`$ is the internucleon distance. This density distribution is symmetric under rotations about the $`Z^{}`$-axis, and the top part shows its cross section in the $`X^{}Z^{}`$ plane as predicted by the A18 model. The bottom part shows the toroidal shape of the equi-density surface for half maximum density. The density peaks on a ring of diameter of $``$ 1.0 fm inside this torus having thickness of $``$ 0.8 fm. The shape is produced by constructive (destructive) interference between the S- and the D-wave functions shown in Fig.2 along the $`X^{}`$ ($`Z^{}`$) axis of the deuteron in the $`M=0`$ state.
The dominant, static part of the $`v_{ij}`$ in the deuteron, obtained by omitting the terms dependent on the angular momentum, is anisotropic due to the tensor part of the OPEP. It strongly depends upon the angle between the unit vector $`𝐫_{ij}`$ and the spin directions. The expectation value of the interaction in the $`S=1,M_S=0`$ state, $`\frac{1}{\sqrt{2}}|+`$, is shown in Fig.4 as a function of $`r_{ij}`$ for $`\theta =0`$ and $`\theta =\pi /2`$. The interaction is attractive for $`\theta =\pi /2`$ and repulsive for $`\theta =0`$, like that between two magnetic dipoles. The OPEP is also shown in Fig.4 by dashed lines. The NN interaction in all states except those with $`T,S=1,0`$ is dominated by the OPEP at $`r_{ij}>1`$ fm. The OPEP is weakest in the $`T,S=1,0`$ states. At small $`r_{ij}`$ the repulsive core of $`v_{ij}^R`$ dominates in all states.
Most of the $`M=0`$ deuteron wave function has also $`M_S=0`$. Thus it is possible to understand the density distribution of the deuteron from the potential shown in Fig.4. The two peaks in the density shown in Fig.3 correspond to the two-nucleons in the deuteron being $``$ 1 fm apart at $`\theta \pi /2`$ where the potential has its minimum value of $`200`$ MeV in the A18 model. Other models also have a deep minimum at this position. The smallness of the density along the $`Z`$-axis is due to the repulsive potential at $`\theta =0`$. Even though these features existed in the potential models of the sixties their experimental confirmation came via a series of measurements of the electromagnetic form factors of the deuteron up to momentum transfers of $``$ 8 fm<sup>-1</sup> conducted since the mid eighties at SLAC and Bates. These and the more recent, high precision measurements carried out at the Jefferson Lab are in good agreement with the predictions of the A18 model, and verify the predicted deuteron structure beyond $`r0.7`$, or equivalently $`r^{}0.35`$ fm.
The $`T,S=1,0`$ two-proton distribution functions, believed to be similar to the two-neutron distribution functions due to isospin symmetry, are predicted to have a dip at $`r0`$ and a peak at $`r1`$ fm, due to the repulsive core and the minima of $`v_{NN}`$ respectively. The experimental information on these spherically symmetric distribution functions is less direct. It comes from sums of longitudinal response functions of light nuclei . The sums observed in <sup>3</sup>He and <sup>4</sup>He are in fair agreement with theory, and show evidence of the predicted structure, however, the relativistic and other corrections to the observed sums are significant.
The observed deuteron form factors and, to a lesser extent, the sums of longitudinal response of light nuclei indicate that the modern two-nucleon potentials and the wave functions they predict have validity at internucleon distances larger than $``$ 0.7 fm. This may appear surprising because the rms charge radius of the proton is known to be 0.8 fm. However, the nucleons seem to have a small and dense core. The charge form factor of the proton, as well as the magnetic form factors of the proton and neutron are well approximated by the dipole $`(1+q^2/q_0^2)^2`$ with $`q_0=840`$ MeV/c. Inverting this form factor gives the proton charge density as $`\rho _{ch}^p(r)=3.3e^{r/0.23fm}`$ fm<sup>-3</sup>. The charge densities of two protons, one fm apart are shown in Fig.5. It should be noted that the charge densities shown can have corrections due to the neglect of relativistic effects in inverting the form factor at $`r\stackrel{<}{}1/m0.21`$ fm, and that recent measurements of the proton form factor show deviations from the dipole form at momenta $`\stackrel{>}{}`$ 4 fm<sup>-1</sup>. They suggest that proton charge density flattens out at $`r<0.3`$ fm. Nevertheless Fig. 5 indicates that nucleons a fm apart can retain their identities. The NN interaction includes the change in the energy due to their overlap, and it has minima near $`r1`$ fm. In absence of the quantum kinetic energy term $`(^2/2m)`$ in the Hamiltonian (Eq.8) the deuteron will shrink to a ring of radius $``$ 0.5 fm, and the equilibrium density of nuclear matter will be $``$ 1 fm<sup>-3</sup>. The density of matter in most neutron stars is less than that.
### 3.2 Models of Three Nucleon Interaction
All realistic models of $`v_{ij}`$, the modern and the older, underbind the triton and other light nuclei and predict too high equilibrium density for symmetric nuclear matter. In both cases the deviation from experiment is not too large, particularly when compared with the expectation values of $`v_{ij}`$. For example, the expectation value $`v_{i<j}`$ in <sup>3</sup>H is about $`50`$ MeV, while the underbinding is by $`<`$ 1 MeV. It is likely that these differences are due to three-nucleon interactions expected and predicted since the fifties .
High precision modeling of the two-nucleon interaction is possible because scattering cross sections can be easily and exactly calculated from the assumed $`v_{ij}`$, and a complete set of $``$ 4000 cross sections has been measured. Such an approach is not practical at present for the three-nucleon interaction $`V_{ijk}`$. In principle deuteron-nucleon scattering and reactions can be used to study $`V_{ijk}`$ in the isospin $`T=1/2`$ state. However, this scattering is dominated by the two-nucleon interaction , and very high precision data is necessary to extract the effects of $`V_{ijk}`$, as is being attempted in new experiments at the Indiana University cyclotron, focused on spin observables.
The large $`v_{i<j}`$ is cancelled to a large extent by the kinetic energy in nuclear binding energies. Thus $`V_{i<j<k}`$ in nuclei is expected to be of the order of 10 $`\%`$ of their binding energy. Hence we can construct realistic models of $`V_{ijk}`$ by fitting binding energies of light nuclei, which can now be calculated with an accuracy of the order of 1 $`\%`$ using Greens Function Monte Carlo (GFMC) methods and the estimated equilibrium properties of nuclear matter. As will be discussed in sec. 3.7, it is not yet possible to calculate the EOS of symmetric nuclear matter with comparable accuracy. The models depend upon the $`v_{ij}`$ used in the Hamiltonian. This is inevitable, because unitary transformations make correlated changes in $`v_{ij}`$ and $`V_{ijk}`$ . Only the combinations of $`v_{ij}`$ and $`V_{ijk}`$ in the Hamiltonian (8) are meaningful.
The information contained in nuclear binding energies and equilibrium properties of nuclear matter is limited. Therefore the realistic models of $`V_{ijk}`$ rely on theory to a much larger extent than the models of $`v_{ij}`$, and contain very few parameters. The Urbana models of $`V_{ijk}`$ contain two isoscalar terms:
$$V_{ijk}=V_{ijk}^{2\pi }+V_{ijk}^R.$$
(12)
The first term represents the Fujita-Miyazawa two-pion exchange interaction:
$$V_{ijk}^{2\pi }=\underset{cyc}{}A_{2\pi }\left(\{𝝉_i𝝉_j,𝝉_i𝝉_k\}\{X_{ij},X_{ik}\}+\frac{1}{4}[𝝉_i𝝉_j,𝝉_i𝝉_k][X_{ij},X_{ik}]\right),$$
(13)
with strength denoted by $`A_{2\pi }`$. The functions $`T_\pi (r_{ij})`$ and $`Y_\pi (r_{ij})`$ in $`X_{ij}`$, Eq.(10), are taken from the A18 model of $`v_{ij}`$. This interaction is due to the pion exchanged by nucleons $`j`$ and $`k`$ being scattered by the nucleon $`i`$ via the $`\mathrm{\Delta }`$ resonance in $`\pi `$-N scattering. In classical terms it is due to the polarization of the quark spins in nucleon $`i`$, due to the pion field of $`j`$ ($`k`$), interacting with $`k`$ ($`j`$); and it is similar to the three-body earth-moon-satellite gravitational interaction due to the polarization of the ocean water on earth by the moon’s gravity.
The $`V_{ijk}^R`$ is purely phenomenological, and has the form:
$$V_{ijk}^R=U_0\underset{cyc}{}T_\pi ^2(r_{ij})T_\pi ^2(r_{ik}).$$
(14)
It was meant to represent the modification of the two-pion exchange part of $`v_{ij}`$ by other particles in matter, however, $``$ 40 $`\%`$ of it is due to relativistic effects discussed in the next subsection. The parameters of the present model U-IX, $`A_{2\pi }=0.0293`$ MeV and $`U_0=0.0048`$ MeV, have been determined from exact GFMC calculations of <sup>3</sup>H and approximate variational calculations of the equilibrium density of nuclear matter with the A18 NN interaction .
The results of essentially exact GFMC calculations with the A18 and U-IX interactions are shown in Table-1. We note that the better known pion exchange parts of these interactions give the largest contributions, but the contributions of the phenomenological parts, $`v_{ij}^R`$ and $`V_{ijk}^R`$, are significant. The column $`\mathrm{\Delta }E_{expt}`$ gives the difference between the experimental energies and the calculated, while $`\mathrm{\Delta }E_{VMC}`$ is that between the variational Monte Carlo (VMC) upper bounds and the exact GFMC energies.
Table-1 shows that the U-IX interaction underbinds $`A=8`$ nuclei, and since <sup>8</sup>He is more underbound than <sup>8</sup>Be, it misrepresents the isospin dependence of $`V_{ijk}`$. The new Illinois models of $`V_{ijk}`$ resolve this problem by including the leading three-pion exchange term, $`V_{ijk}^{3\pi }`$, that is attractive in triplets having isospin $`T=3/2`$, but has little effect on the $`T=1/2`$ triplets in <sup>3</sup>H and <sup>4</sup>He. A much improved fit, with errors $`<2\%`$, to the observed energies is obtained as shown in Fig.6. The three parameters in Pieper’s model IL-2R are the strengths of the $`V_{ijk}^{2\pi }`$, $`V_{ijk}^R`$ and $`V_{ijk}^{3\pi }`$. Calculations of NM properties with this more accurate $`V_{ijk}`$ are in progress; the preliminary results are similar to those with U-IX since the $`V^{3\pi }`$ is much weaker than the $`V^{2\pi }`$ and $`V^R`$. For example, the expectation values of $`V_{ijk}^{2\pi }`$, $`V_{ijk}^R`$ and $`V_{ijk}^{3\pi }`$ of IL-2R model, in <sup>8</sup>Be (<sup>8</sup>He) are respectively $``$ $``$38 ($``$27), +19 (+14) and $``$2 ($``$5) MeV. In the following sections we will review the properties of neutron star matter calculated with A18 and U-IX interactions.
In principle there can be four-nucleon interactions (FNI) neglected in the Hamiltonian (Eq.8). It seems that they are very weak in nuclei. All the models of $`V_{ijk}`$ studied so far reproduce the energy of <sup>4</sup>He with an error $`<0.5\%`$, after fitting the observed energy of <sup>3</sup>H. Since this error is close to the accuracy of the <sup>4</sup>He calculation, there is no indication of FNI in that nucleus. The IL-2R model also gives the experimental energies of $`A=8`$ nuclei within $`1`$ MeV. In <sup>8</sup>Be, for example, the expectation values of $`v_{ij}`$ and $`V_{ijk}`$ are respectively $``$308 and $``$21 MeV respectively. By comparison with experiment we estimate that the possible contribution of FNI in this nucleus is $`<`$ 1 MeV. The $`V_{ijk}`$ presumably has additional smaller terms neglected in IL-2R, but it is difficult to determine their strengths from the nuclear spectra that can be calculated accurately from bare forces at present.
### 3.3 Relativistic Boost Interaction
In all models, the NN scattering data is reduced to the center of mass frame and fitted using phase shifts calculated from the NN interaction, $`v_{ij}`$, in that frame. The $`v_{ij}`$ obtained by this procedure describes the interaction between nucleons having total momentum $`𝐏_{ij}=𝐩_i+𝐩_j=0`$. In general, the interaction between particles depends on their momenta $`𝐩_i`$ and $`𝐩_j`$. For example the electromagnetic interaction between two particles of mass m and charge Q, contains a term with the factor $`𝐩_i𝐩_j/2m^2=p_{ij}^2/2m^2P_{ij}^2/8m^2`$, where $`𝐩_{ij}=(𝐩_i𝐩_j)/2`$ is the relative momentum. The terms containing $`𝐩_{ij}`$ are included in the momentum-dependent parts of $`v_{ij}`$, while the boost interaction $`\delta v(𝐏_{ij})`$ contains parts dependent on the total $`𝐏_{ij}`$. Even though contributions of the boost interaction to the binding energy of SNM and <sup>3</sup>H were estimated by Coester and coworkers in the seventies and eighties , they were first included in studies of dense matter rather recently . Walecka’s relativistic mean field theory naturally contains the boost interactions .
Following the work of Krajcik and Foldy, Friar obtained the following equation relating the boost interaction of order $`P^2`$ to the interaction in the center of mass frame:
$$\delta v(𝐏)=\frac{P^2}{8m^2}v+\frac{1}{8m^2}[𝐏𝐫𝐏,v]+\frac{1}{8m^2}[(\sigma _i\sigma _j)\times 𝐏,v].$$
(15)
The general validity of this equation in relativistic mechanics and field theory was recently demonstrated . Including boost interaction, the nonrelativistic Hamiltonian assumes the form:
$$H^{}=\frac{p_i^2}{2m}+(v_{ij}+\delta v(𝐏_{ij}))+V_{ijk}^{}+\mathrm{},$$
(16)
where the ellipsis denotes the three-body boost, and four and higher body interactions. This $`H^{}`$ contains all terms quadratic in the particle velocities, and is therefore suitable for complete studies in the nonrelativistic limit.
Studies of light nuclei using the VMC method find that the contribution of the two-body boost interaction to the energy is repulsive, with a magnitude which is $``$ 37% of that of $`V_{ijk}^R`$ in the UIX model. The boost interaction thus accounts for a significant part of the $`V_{ijk}^R`$ in Hamiltonians which fit nuclear energies neglecting $`\delta v`$. The $`V_{ijk}^{}`$ in Eq.(16) has a $`V^R`$ of appropriately smaller strength than that in the $`V_{ijk}`$ in Hamiltonian $`H`$ given by Eq.(8).
We should expect additional relativistic corrections to the Hamiltonian (16). However, when nonrelativistic potentials are fit to the experimental data, relativistic effects present in the data are automatically buried in these potentials. In order to study the magnitude of a chosen relativistic correction, such as that due to the approximation of the kinetic energy or the boost interaction, or the nonlocalities of OPEP, it is necessary to refit the same data set, fitted to obtain the nonrelativistic Hamiltonian, and then study the differences. Such comparisons indicate that the relativistic corrections associated with kinetic energies and nonlocalities of OPEP are small, whereas the boost corrections are significant. This is not surprising since the boost interaction was totally omitted from the conventional nonrelativistic nuclear Hamiltonian (Eq.8).
### 3.4 Brueckner Calculations of Nucleon Matter
Calculating the properties of matter from the interaction $`v_{ij}`$ between pairs of its constituents is a well known problem in many-body theory. It is particularly challenging for NM due to the strong spin-isospin dependence of $`v_{ij}`$. In the method developed by Brueckner, Bethe and Goldstone the perturbation expansion of the energy of NM is cast into a series ordered according to the number of independent hole lines (HL). This method has been used extensively since the sixties to study symmetric nuclear matter (SNM) and nuclei. It was also used then to predict properties of pure neutron matter (PNM) soon after the discovery of pulsars. The convergence of the expansion depends upon the choice of the single-particle energies in the assumed unperturbed Hamiltonian. For the hole states with momenta less than $`k_F`$, they are chosen self consistently as suggested by Brueckner and Gammel via the Brueckner-Hartree-Fock (BHF) procedure. The older calculations by Day used kinetic energies for particle states. Since this leads to a discontinuity in the unperturbed single-particle energies at $`k_F`$, Day’s choice is called “discontinuous”. In 1976 the Liège group advocated the “continuous” choice by extending the definition of BHF single-particle energies to particle states. If all the higher order terms of the HL expansion are computed the final results should be independent of the choice.
Detailed calculations of SNM have been carried out with the older Argonne $`v_{14}`$ interaction by the Catania group using both choices. The results of the lowest order 2-HL calculation depend significantly on the choice, however, those including 2+3-HL terms are almost the same for the two choices. They find that the 2-HL results with the continuous (discontinuous) choice are $``$ 15 % below (30 % above) those of the 2+3. Day’s calculations , with the same interaction, include additional 4-HL terms, and give energies below the 2+3 Catania results, closer to the 2-HL continuous. For example, the energies of SNM at $`\rho =0.28`$ fm<sup>-3</sup> are $``$11.3, $``$16.1, $``$18.3 and $`17.8\pm 1.3`$ for the Catania 2-HL discontinuous, 2+3-HL, 2-HL continuous, and Day’s calculations respectively, while the variational upper bound at this density is $`16.2\pm 0.4`$ with the methods described in the next subsection. The size of the error due to truncation of the expansion is estimated in Day’s and variational calculations.
A significant advantage of Brueckner’s method is that it can be easily applied to local as well as nonlocal modern interactions. The present calculations, called lowest-order BHF (LOBHF), include only the 2-HL terms and use the continuous choice . Their results for PNM and SNM are shown in Fig.7; those for matter with intermediate values of proton fraction can be estimated from these by interpolating with $`\beta ^2`$ as metioned in sec. 2. These results provide an estimate of the uncertainty in the predicted matter energy due to that in the NN interaction.
The LOBHF energies for neutron matter are essentially model independent up to $`0.3`$ fm<sup>-3</sup> (Fig.7); at higher densities they deviate partly for the following reason. All the models fit the NN scattering data up to 350 MeV lab energy, i.e. up to maximum relative momenta $`k=2.05`$ fm<sup>-1</sup>. The maximum value of the relative momentum of two hole states in matter is $`k_F`$, and the density of PNM (SNM) at $`k_F=2.05`$ is 0.29 (0.58) fm<sup>-3</sup>. In PNM at $`\rho >0.29`$ the interactions are being used at relative momenta larger than in the fitted data. Baldo et al have shown that the $`{}_{}{}^{3}P_{2}^{}`$ phase shifts predicted by the five modern potentials vary from 8 (A18) to 19 (Nijmegen II) degrees at $`k=3`$ fm<sup>-1</sup>. The average value, $`k_{rms}=\sqrt{3/10}k_F`$, is smaller than $`k_F`$, and exceeds 2.05 fm<sup>-1</sup> at much larger densities of 1.77 (3.54) fm<sup>-3</sup> of PNM (SNM). Presumably, this helps to keep down the model dependence.
The LOBHF energies of SNM have a larger model dependence starting at lower densities (Fig.7). Here the main cause seems to be the assumed nonlocality of $`v_{ij}`$ discussed in sec. 3.1. The local interactions, Nijmegen II, Ried 93 and A18 give similar results, while the most nonlocal CD-Bonn gives the lowest energies. The predicted values of equilibrium $`\rho _0`$ and $`E_0`$ of SNM are respectively 0.31, 0.27, 0.28, 0.27 and 0.37 fm<sup>-3</sup> and $`20.3,17.6,18.7,18.1`$ and $`22.9`$ MeV with Nijmegen I, II, Reid93, A18 and CD Bonn interactions; while the empirical values are 0.16 fm<sup>-3</sup> and $`16`$ MeV. Obviously, the empirical properties of SNM can not be obtained by approximating the nuclear interaction energy by $`v_{i<j}`$. TNI, added to the Hamiltonian to obtain the observed properties, naturally depend upon the choice of $`v_{ij}`$. For example, those to be added to CD Bonn have to have stronger repulsive parts, which dominate at large $`\rho `$. The combinations of $`v_{ij}+V_{ijk}`$, constrained with experimental data, will have smaller model dependence than seen in Fig.7.
Baldo, Bombaci and Burgio have carried out LOBHF calculations with the older Paris and Argonne interactions including Urbana TNI. Like the modern models, these older models give too large $`\rho _0`$ ($`0.26`$ fm<sup>-3</sup>) and $`E_0`$ $`(18`$ MeV) in LOBHF without TNI. By averaging over the position of the third nucleon the TNI is expressed as a density dependent NN interaction to be added to the $`v_{ij}`$. The parameters $`A_{2\pi }`$ and $`U_0`$ were chosen to get closer to the empirical values of $`\rho _0`$ and $`E_0`$; their values of $``$0.0329 and 0.00361 MeV are not too far from those of U-IX ($``$0.0293 and 0.0048). With the Paris + Urbana model they obtain $`\rho _0=0.176`$ fm<sup>-3</sup>, $`E_0=16.0`$ MeV, $`K=281`$ MeV, and $`E_{sym}=33`$ MeV. These values, as well as those obtained using the Argonne interaction instead of Paris, compare rather well with the empirical values given in sec. 2.
Relativistic effects are included in the LOBHF calculations via Dirac-Brueckner approximation suggested by Celenza and Shakin . The calculations include contributions of the boost interactions as well as TNI , and many nucleon interactions generated via the $`Z`$-graphs containing anti-nucleon lines. However, they do not contain contributions of the Fujita-Miyazawa and other TNI due to internal structure of nucleons. Results have been reported by Brockmann and Machleidt for the older Bonn meson exchange NN interaction models; those with the Bonn-A model come close to reproducing the empirical properties of SNM at $`\rho _0`$.
### 3.5 Variational Calculations of Nucleon Matter
Variational calculations of NM with realistic interactions have been carried out since 1970 . The present calculations use variational wave-functions, $`\mathrm{\Psi }_v`$, consisting of a symmetrized product of pair correlation operators, $`F_{ij}`$, operating on the Fermi gas wave-function. In PNM, the $`F_{ij}`$ include four terms generating spatial, $`𝝈_i𝝈_j`$, tensor and spin orbit correlations. The SNM $`F_{ij}`$ have eight terms; the additional four have $`𝝉_i𝝉_j`$ factors.
This wave-function is clearly too simple to accurately describe the ground state of nuclear matter. Monte Carlo studies of few-body nuclei use additional three-body correlations induced by both $`v_{ij}`$ and $`V_{ijk}`$, in the variational wave-function; they reduce the energy of <sup>16</sup>O by $``$ 1 MeV/nucleon. The results shown in Table.1 indicate that the VMC energies of $`A=8`$ nuclei, obtained after including three-body correlation operators, are above the exact GFMC values by $``$ 1 MeV/nucleon. From these we estimate that the present $`\mathrm{\Psi }_v`$ may underbind SNM by a few MeV. In contrast, the three-body correlations have a smaller effect on the energy of pure neutron drops . The variational energy of a drop with eight neutrons, calculated with the simple $`\mathrm{\Psi }_v`$, is greater than the exact value by $``$ 0.5 MeV/nucleon. Thus the variational energies are relatively more accurate for PNM than for SNM. This is as expected, since SNM has much stronger tensor correlations. Despite the aforementioned shortcomings, the simple $`\mathrm{\Psi }_v`$ having only pair correlation operators describes the gross features of the nuclear wave-function rather well. For example, the spin-isospin dependent two-nucleon distribution functions calculated in this approximation are close to the exact distribution functions .
The correlation operators $`F_{ij}`$ are determined from Euler-Lagrange equations that minimize the two-body cluster contribution of an interaction $`\overline{v}_{ij}\lambda _{ij}`$. The interaction $`\overline{v}_{ij}`$ is related to the $`v_{ij}`$ via a parameter $`\alpha `$ meant to simulate the quenching of the spin-isospin interaction between particles i and j, via their interaction with other particles in matter. The operator $`\lambda _{ij}`$ simulates screening effects in matter; it is determined from the ranges $`d_t`$ of tensor correlations, and $`d_c`$ of all the other correlations. The $`\mathrm{\Psi }_v`$ thus depend on three variational parameters: $`\alpha `$, d<sub>c</sub> and d<sub>t</sub>, determined by minimizing the energy. Two additional parameters are used in Ref. to further lower the variational upper bound by small amounts.
The energy expectation value is evaluated using cluster expansion. The one-body term is just the Fermi gas kinetic energy, and the large two-body (2B) term, analogous to the interaction energy in LOBHF, is calculated exactly. The most important of the many-body (MB) cluster contributions are summed using Fermi-hypernetted chain (FHNC) and single operator chain equations , and constrains are imposed to satisfy the fundamental identities of pair distribution functions. The kinetic energy can be calculated using different expressions related by integration by parts. If all MB contributions are calculated, these expressions yield the same result. However, they yield different results when only selected parts of the MB clusters are summed using chain equations. Studies of atomic helium liquids with FHNC summation methods find the exact result to be between the energies obtained using the Jackson-Feenberg (JF) and Pandharipande-Bethe (PB) expressions. The average of these two is used as the result with half the difference as an estimate of the error.
More general pair correlations can be calculated by separately minimizing the two-body cluster contribution to each partial wave, specified by $`l,S,J`$ and the relative momentum $`k`$ . These correlations $`f(l,S,J,k)`$ depend on all the quantum numbers, and yield a lower 2B energy than the $`F_{ij}`$ operator with the same $`\alpha `$, $`d_c`$ and $`d_t`$. The MB contributions cannot be easily calculated with the general $`f(l,S,J,k)`$, however. The differences between optimum $`f(l,S,J,k)`$ and $`F_{ij}`$ can be included via the second order two-particle, two-hole contribution, $`\mathrm{\Delta }E_2`$, in correlated basis perturbation theory . In recent calculations the $`\mathrm{\Delta }E_2`$ is approximated by the difference $`\delta E_{2B}`$ between the 2B cluster energies calculated using $`f(l,S,J,k)`$ and $`F_{ij}`$. However, the values of $`\alpha ,d_c,d_t`$ are determined by minimizing the energy calculated from the $`F_{ij}`$.
The variational and LOBHF energies obtained with the A18 interaction are compared in Fig.8. At densities below 0.6 fm<sup>-3</sup> there is fairly close agreement between them, however, we expect the true results to be a few MeV below the variational upper bound. At higher densities the SNM LOBHF energy is significantly below the variational bound. The convergence of the HL expansion is expected to deteriorate at higher densities, and it may be the cause of the large difference.
### 3.6 Neutral Pion Condensation
Variational calculations with the Hamiltonian (16) with A18+$`\delta v`$+UIX interactions, indicate the occurrence of a phase transition in both PNM and SNM; the energies of the two phases are shown in Fig.9. The tensor correlations have a longer range in the higher density phase (HDP) than in the low density phase (LDP). Detailed studies of the pion fields in the two phases indicate that the HDP has a large enhancement of virtual neutral pions with momenta $`1.5`$ fm<sup>-1</sup>, and therefore this transition is believed to be due to neutral pion condensation in matter. Although the effect of this type of transition on the EOS is relatively small, it can have important consequences for the cooling and evolution of neutron stars .
Since the pioneering work of Migdal and of Sawyer and Scallapino , many investigators have used effective interactions to study the possibility of pion condensation in SNM and PNM. These efforts were recently reviewed by Kunihiro et al . Neutral pion condensation occurs when PNM (SNM) becomes unstable towards the development of a spin (spin-isospin) density wave as discussed in sec. 4.2. The variational wave functions used in the Urbana calculations are not adequate to describe the long range order expected with $`\pi ^0`$-condensation; better calculations may be possible with the quantum Monte Carlo method described in the next subsection. NM with spin or spin-isospin density wave naturally has a pion field of the same wavelength; it causes $`N\mathrm{\Delta }`$ transitions, which help lower the energy of this phase. The effects of the $`\mathrm{\Delta }`$-resonance are absorbed into the interactions in the Hamiltonian (16). In particular, the Fujita-Miyazawa $`V_{ijk}^{2\pi }`$ has a large effect on $`\pi ^0`$ condensation. Without it there is no condensation predicted in SNM, while that in PNM occurs at a higher density of $``$ 0.5 fm<sup>-3</sup>.
This transition was found by Wiringa, Fiks and Fabrocini to occur with the older Argonne $`v_{14}`$ for PNM but not for SNM , while it does not occur in either PNM or SNM with the Urbana $`v_{14}`$ interaction of 1981. In Migdal’s approach , the transition of SNM to the pion condensed phase is inhibited by a positive, short-range $`𝝈_i𝝈_j𝝉_i𝝉_j`$ NN interaction whose strength is represented by the Landau parameter $`g^{}`$. In the case of PNM the sum of the $`𝝈_i𝝈_j`$ and $`𝝈_i𝝈_j𝝉_i𝝉_j`$ interactions occurs since $`𝝉_i𝝉_j=1`$. The contact part of OPEP gives a negative $`𝝈_i𝝈_j𝝉_i𝝉_j`$ potential at small $`r`$, and thus favors pion condensation. The $`𝝈_i𝝈_j𝝉_i𝝉_j`$ part of the modern A18 NN interaction does become negative at small $`r`$, however, it is positive in the older Urbana-Argonne models which do not fit the NN scattering data as well as the A18.
The energies of $`\beta `$-stable LDP and HDP phases have been calculated by interpolation between PNM and SNM . The LDP phase with $`\rho =0.204`$ fm<sup>-3</sup> and proton fraction $`x_p=0.073`$ is found to be in equilibrium with the HDP at $`\rho =0.237`$ and $`x_p=0.057`$. In between there will be mixed phase regions as discussed in sec. 5.
### 3.7 Quantum Monte Carlo Calculations
The variational energy of the ground state of SNM, calculated with the A18+$`\delta v`$+UIX model is $`12`$ MeV, against the empirical value of $`16`$ MeV. As mentioned earlier, this difference is believed to be mostly due to the inadequacy of the present variational wave functions. Accurate calculations of the energies of many-body systems in which the interactions do not depend upon the spins of the particles have been carried out using quantum Monte Carlo (QMC) methods . In these systems one can work with a wave function $`\mathrm{\Phi }(𝐑)`$ that depends only upon the positions of all the particles represented by the configuration vector $`𝐑=𝐫_1,𝐫_2,\mathrm{}𝐫_A`$. A VMC calculation is used to obtain a good approximation $`\mathrm{\Phi }_v(𝐑)`$, and in the GFMC method one operates on $`\mathrm{\Phi }_v(𝐑)`$ with the imaginary time evolution operator $`exp([HE_0]\tau )`$ to project out the exact ground state $`\mathrm{\Phi }_0(𝐑)`$.
The main difficulty in applying QMC methods to nuclear problems is that nuclear forces change spins and isospins of the interacting nucleons, and thus nuclear wave functions contain superpositions of all possible spin-isospin states of $`A`$-nucleons. Their number, $`2^AA!/[Z!(AZ)!]`$, increases very rapidly with $`A`$. For this reason exact QMC calculations have been carried out only for nuclei having up to 8 nucleons , and attempts to calculate $`A=9`$ nuclei are in progress. Carlson has also calculated the ground states of 14 neutrons in a periodic box. Calculations of pure neutron systems are simpler because there is only one isospin state, and all $`𝝉_i𝝉_j`$ can be replaced by the unit operator.
In addition, GFMC calculations of Fermion systems suffer from the “Fermion sign problem”. The real wave functions of simple Fermi systems have nodal surfaces because the $`\mathrm{\Phi }(𝐑)`$ must equal $`\mathrm{\Phi }(𝐑^{})`$ when the configurations $`𝐑`$ and $`𝐑^{}`$ are related by the exchange of a pair of particles. GFMC configurations which diffuse across nodal surfaces, as the system evolves in imaginary time, increase the variance of the calculated quantity, making unconstrained propagation impractical for large systems. In the fixed node method for simple systems this growth of variance is eliminated by restricting the configurations to domains enclosed by the nodal surfaces of the variational wave function. Such calculations generally have an error due to imperfections in that structure, but it is much smaller than that in variational.
A similar problem comes in nuclear GFMC with the additional complexity due to nuclear wave functions having many spin-isospin components, each with a different nodal structure. The growth of the variance is tolerable for $`A7`$, but when $`A8`$ it is necessary to use constrained path methods to control the variance. The constraint can be removed at large $`\tau `$ to test if it influenced the calculated energy significantly. It appears that calculations with $``$ 2% accuracy in the binding energy are possible in this way for systems having up to 14 neutrons.
Auxiliary field diffusion Monte Carlo (AFDMC) seems to be the long-sought breakthrough needed to eliminate the exponential $`(2^A)`$ growth of spin states in GFMC calculations of neutron matter. This method combines two major themes in QMC. Auxiliary-field methods are used in the shell model Monte Carlo calculations , and several condensed matter systems in which the continuous spatial degrees of freedom have been eliminated, while diffusion Monte Carlo is another name for GFMC.
In the approach developed by Schmidt and Fantoni , the spatial parts, (i.e. kinetic energy and spin-independent interactions), of the Hamiltonian are propagated as in GFMC and the spin-dependent interactions between neutrons are replaced by interactions of neutrons with auxiliary fields. Integrating over the auxiliary fields reproduces the original spin-dependent interaction. In addition, a constraint analogous to the fixed-node approximation in GFMC is introduced, by requiring that the real part of the overlap with a trial function remains positive.
More recently, Schmidt and Fantoni have carried out calculations with a realistic Hamiltonian consisting of Argonne $`v_8^{}`$ NN interaction used in the GFMC calculations , and the UIX TNI. The $`v_8^{}`$ contains the main parts of A18, and the difference between the two is treated perturbatively. Results of calculations with 38 neutrons in a periodic box with finite size corrections have been obtained. They are $`5\%`$ below the variational energies obtained with the methods described above. For example, the AFDMC and variational energies for Argonne $`v_8^{}`$ and UIX interactions are 21.8 $`\pm 0.1`$ (65.5 $`\pm 0.1`$) and 23.2 (68.6) MeV per neutron at $`\rho =0.2(0.4)`$ fm<sup>-3</sup>.
The trial functions used to constrain the present AFDMC calculations are rather simple without any spin correlations. In contrast it is possible to use more accurate variational wave functions with spin correlations to constrain the GFMC calculations. Carlson has compared AFDMC and GFMC results for 14 neutrons in a periodic box. At $`\rho =0.15`$ fm<sup>-3</sup> the GFMC energy (220 $`\pm `$ 1) is about 7 % below the AFDMC result of 236.4 $`\pm `$ 1.5 MeV. From these we conclude that the variational PNM energies given in the last section may be $``$ 12 % above the exact values for the A18 and UIX interactions. The error in SNM $`E(\rho )`$ is probably twice as large. As mentioned in the conclusions (sect.7), an overestimation of the $`E(\rho )`$ of neutron matter by 12 % has a rather small effect on the predictions of neutron star properties.
The AFDMC is more accurate than the present variational method, and it is also more versatile. For example, it can be used to study matter with long range spin-isospin order induced by $`\pi ^0`$-condensation discussed in sect 4.2.
## 4 Hadronic and Quark Matter
It is likely that at high densities more general form of matter containing hadrons besides the nucleons, called hadronic matter (HM), has lower energy. The possibilities are that it contains negatively charged mesons like pions or kaons, or other hyperons such as $`\mathrm{\Sigma }^{}`$ or $`\mathrm{\Lambda }`$. Finally, it is expected that at a high enough density there will be a transition to quark matter (QM) in which the quarks are not clustered into nucleons or hadrons.
The interactions between hyperons and nucleons, and between kaons and nucleons are not as well known as those between nucleons, and the energy of quark matter is difficult to calculate realistically. Therefore the transition densities from NM to HM or QM are rather difficult to calculate reliably. We review the recent estimates.
### 4.1 Kaon Condensation
Kaon condensation in dense matter was suggested by Kaplan and Nelson , and has been discussed in many recent publications . Due to the attraction between $`K^{}`$ and nucleons the kaon energy decreases with increasing density, and eventually if it drops below the electron chemical potential in NM, a Bose condensate of $`K^{}`$ will appear. The key quantities of interest are the electron and kaon chemical potentials in NM. The former is obtained from the $`\beta `$-equilibrium condition, $`\mu _e=\mu _n\mu _p`$ relating electron, neutron and proton chemical potentials. The $`\mu _e`$ obtained from the A18+$`\delta v`$+UIX interactions with variational calculations is shown in Fig.10.
In neutron matter at very low densities, the interparticle spacing is much larger than the range of the $`K^{}n`$ interaction, and the kaon interacts many times with the same nucleon before it encounters and interacts with another nucleon. Therefore one can use the scattering length, $`a_{K^{}n}`$, as the “effective” kaon-nucleon interaction. In this low density limit the kaon energy deviates from its rest mass by the Lenz potential, and is given by :
$$\omega _{Lenz}=m_K+\frac{2\pi }{m_R}a_{K^{}n}\rho ,$$
(17)
where $`m_R=m_Km_n/(m_K+m_n)`$ is the kaon-neutron reduced mass. The scattering length extracted from data is $`(0.37i0.57)`$ fm; its imaginary part is due to the open $`\mathrm{\Lambda }\pi ^{}`$ channel in vacuum. In the density region of interest to kaon condensation the kaon energy is too small for the $`K^{}+n\mathrm{\Lambda }+\pi ^{}`$ reaction to occur. Using effective lagrangians based on chiral perturbation theory Brown et. al. estimate $`a_{K^{}n}`$ to be $`0.41`$ fm in absence of reaction channels.
There are two corrections to the Lenz energy at small densities. Including these, the kaon energy $`\omega `$ obtains the form :
$$\omega =m_K+\left(\frac{2m_K}{m_K+\omega }\right)\left(\frac{1}{1a_{K^{}n}\xi \rho }\right)\frac{2\pi }{m_R}a_{K^{}n}\rho ,$$
(18)
where $`\xi \rho `$ is the inverse correlation length. The first correction factor is a relativistic effect obtained from the Klein-Gordon equation, while the second factor is from the theory developed by Ericson and Ericson for propagation of mesons in nuclear matter. The relativistic correction decreases the kaon energy, while the correlation correction increases it.
As the density increases further, and the interparticle spacing becomes of the order of the range of the interaction, the kaon will simultaneously interact with two or more nucleons and the Lenz approximation will break down. In the high density limit the kaon energy deviates from its rest mass by the Hartree potential:
$$\omega _{Hartree}=m_K+\rho v_{K^{}n}(r)d^3r,$$
(19)
where $`v_{K^{}n}`$ is the $`K^{}n`$ interaction potential. As shown in Ref. , the Hartree potential is considerably less attractive than the Lenz potential and thus $`\omega _{Hartree}>\omega _{Lenz}`$. The transition from the Lenz to Hartree limits has been recently studied with a variety of methods including exact calculations for simple cubic crystal model of neutron matter. For reasonable interaction range the transition begins at very low densities ($`<0.1\rho _0`$), and the Hartree limit is essentially reached by $`3\rho _0`$. There are no relativistic corrections to the Hartree energy of kaons condensed in the state with zero momentum provided the $`K^{}n`$ interaction is dominated by the Weinberg-Tomozawa vector potential.
The typical resent results for kaon energy in neutron star matter are shown in Fig.10. The top solid line is obtained with a Wigner-Seitz calculation for pure neutron matter; it is exact in both the low and high density limits, and gives essentially the Hartree energy at $`\rho >3\rho _0`$. The next curve shows the estimated Hartree results for NM containing $``$ 15 % protons, below that is the $`\omega _{Lenz}`$ for neutron matter, while the lowest curves uses $`K^{}N`$ scattering amplitude calculated in matter including relativistic, correlation and proton fraction corrections . The recent estimate for the central density of a 2.0 $`M_{}`$ star is $`5\rho _0`$, and the results shown in Fig.10 indicate that kaon condensation is unlikely at densities lower than that. The heaviest stars ($`M=2.2M_{}`$) made up of NM are predicted to have central densities of $`7\rho _0`$, and the possibility of kaon condensation in their cores can not be ruled out. However, presence of $`\pi ^{}`$, or $`\mathrm{\Sigma }^{}`$ or quark drops will decrease $`\mu _e`$ as discussed in the following subsections and may make kaon condensation unlikely in even the most massive stars.
### 4.2 Charged Pion Condensation
Negatively charged $`\pi ^{}`$-mesons will condense in matter when their chemical potential becomes lower than $`\mu _e`$, as suggested by Migdal and by Sawyer and Scallapino . In the seventies and eighties this possibility was studied by many researchers. Their work has been reviewed in Ref. , and we discuss it rather briefly.
Fig. 10 suggests that, in absence of interactions $`\pi ^{}`$ with zero momentum will condense in matter at a rather low density of $`1.5\rho _0`$ when $`\mu _e`$ exceeds their rest mass of 139 MeV. However, the $`\pi ^{}n`$ S-wave interaction is repulsive, and raises the energy of zero momentum pions sufficiently above the estimated $`\mu _e`$. The recently found , deeply bound $`\pi ^{}`$-nucleus atomic states, are influenced by this repulsion. The $`K^{}n`$ S-wave interaction, on the other hand is attractive, leading to the possibility of kaon condensation discussed in the last section.
The energy of $`\pi ^{}`$, having momenta of the order of 2 fm<sup>-1</sup>, $``$ 400 MeV without interactions, is reduced in matter by the $`\pi ^{}n`$ P-wave interaction due to the coupling of the $`\pi ^{}`$ to $`p`$-$`n`$ and $`N`$-$`\mathrm{\Delta }`$ particle-hole states. Their energies and couplings to $`\pi ^{}`$ are calculated with effective forces described with Landau parameters whose density dependence is not well established. With plausible values for the Landau parameters a second order transition with $`\pi ^{}`$ condensation is predicted at $`2\rho _0`$; however it is not expected to have a large effect on the EOS .
The Japan group also predicts a first order transition to an interesting phase with both $`\pi ^{}`$ and $`\pi ^0`$ condensation at a density of $`4.5\rho _0`$. It has a significant effect on the EOS. Matter in this phase has spin aligned layers as discussed earlier in sec. 3.6 and illustrated in Fig. 11. The condensed $`\pi ^0`$-mesons have momenta in the Z-direction, perpendicular to the layers, while that of $`\pi ^{}`$ is in the X-direction in the plane of the layers. This way, when a proton absorbs a $`\pi ^{}`$ form the condensate and becomes a neutron its spin direction also gets flipped maintaining the attractive interaction with the $`\pi ^0`$ field. It is a challenge to calculate the energy of matter in this interesting phase from bare nuclear forces.
The Illinois group, working with bare nuclear forces, predicts the first order $`\pi ^0`$-condensation at a much lower density of $`1.5\rho _0`$ (sec. 3.6). However, the total decrease in energy of matter at $`5\rho _0`$ due to pion condensation, estimated by the two groups: $``$ 80 by the Japan and $``$ 60 MeV/nucleon by Illinois, are not too different. The recent EOS of the Illinois group contains this energy gain.
### 4.3 Hyperonic Matter
The possibility of hyperons contributing to the ground state of dense matter has been considered since 1959 . The negatively charged $`\mathrm{\Sigma }^{}`$ and $`\mathrm{\Delta }^{}`$ will occur when their chemical potential becomes less than $`\mu _n+\mu _e`$ in matter. With A18+$`\delta v`$+UIX (A18+$`\delta v`$) interactions, the $`\mu _n+\mu _e`$ exceeds the rest mass (1197 MeV) of $`\mathrm{\Sigma }^{}`$ hyperon at density of $``$ 2.2 (2.6) $`\rho _0`$. If the $`\mathrm{\Sigma }^{}`$-nucleon interactions are negligible, they will occur in matter via the weak interaction $`e^{}+n\mathrm{\Sigma }^{}+\nu _e`$ at these densities mainly on account of the large $`\mu _e170`$ MeV. Similarly the neutral hyperons such as the $`\mathrm{\Lambda }`$ will occur when their chemical potentials become less than $`\mu _n`$. For the above interaction models the $`\mu _n`$ exceeds the $`\mathrm{\Lambda }`$ rest mass (1116 MeV) at densities of 2.9 and 3.7 $`\rho _0`$, while the predicted central densities of 1.4 $`M_{}`$ stars are 3.4 and 5.1 $`\rho _0`$.
Much less data exist on hyperon-nucleon (YN) interactions than on the NN, and therefore their models are less constrained. LOBHF calculations using the older Nijmegen soft core YN and either Paris or Argonne $`v_{14}`$ NN interactions show that the thresholds for $`\mathrm{\Sigma }^{}`$ and $`\mathrm{\Lambda }`$ to appear are not much moved by YN interactions. For example, with the Paris NN interaction the threshold densities with (without) YN interactions are 3.0 (2.7) and 3.6 (3.7) $`\rho _0`$ for $`\mathrm{\Sigma }^{}`$ and $`\mathrm{\Lambda }`$ respectively. The Nijmegen group has recently constructed boson exchange interaction models based on YN and NN data base using SU(3) symmetry. LOBHF calculations using these models for NN as well as YN interactions give 2.2 $`\rho _0`$ for the threshold density of $`\mathrm{\Sigma }^{}`$, however for $`\mathrm{\Lambda }`$ it is pushed beyond 7 $`\rho _0`$.
The above LOBHF calculations use only two-body interactions without boost corrections and TNI, and fail to explain the saturation density of SNM as mentioned in sec. 3.4. Both the BHF groups find that the threshold densities are lowered after including TNI effects. For example, the Catania group predicts them to be $``$ 2.1 and 2.4 $`\rho _0`$ with models including TNI and also with Dirac-Brueckner calculations. However, $`\mathrm{\Sigma }^{}NN`$ and $`\mathrm{\Lambda }NN`$ three body forces should also be included along with the TNI for consistency. The binding energies of $`\mathrm{\Lambda }`$-hypernuclei suggest that the $`\mathrm{\Lambda }NN`$ interaction is as strong as the TNI, while there is no data on the $`\mathrm{\Sigma }^{}NN`$.
Ignoring three-body forces and boost interactions in both nucleon and hyperon matter, the LOBHF calculations with the unified NN, YN and YY interaction model indicate that admixtures of $`\mathrm{\Sigma }^{}`$ lower the NM energy by only $`25`$ MeV per nucleon at 5 $`\rho _0`$. At this density the A18+$`\delta v`$+UIX model gives an energy of $`200`$ for $`\beta `$-stable NM. If the energy gain due to $`\mathrm{\Sigma }^{}`$, is not much changed by three-body interactions, we should expect $``$ 10 % effects on the EOS of neutron star matter at $`\rho >2\rho _0`$ due to hyperons; however, if present they would lower the $`\mu _e`$ and increase the proton fraction $`x_p`$ significantly.
### 4.4 Quark Matter
When matter is compressed to densities so high that the quark cores of nucleons overlap substantially, one expects the nucleons to merge and undergo a phase transition to quark matter (QM). The EOS of both HM and QM are necessary to calculate effects of this transition in neutron stars. The interface properties are also needed to study the important mixed phase regions.
At present, lattice QCD can only treat the case of zero baryon chemical potential and is therefore not useful for neutron stars. Lacking a full theory, the simple Bag model is used to estimate the QM EOS. In this model the QM energy has a volume term denoted by the bag constant $`B`$; it represents the difference in the energies of the vacua occupied by hadron and QM, and is responsible for the confinement of quarks within nucleons in nuclei. The term dominating at high densities is the energy of noninteracting u, d and s-quarks; it is calculated neglecting the mass of u, and d-quarks, and using a typical value of $``$ 150 MeV for that of the s-quark. Since the quark Fermi energies are much larger than their masses, the QM properties are not too sensitive to the chosen mass for the s-quark. There is no one-gluon exchange interaction energy between quarks of different flavor, while that between quarks of of the same flavor i is given by $`(2\alpha /3\pi )E_i`$ per quark of flavor i . Here $`E_i`$ is the average kinetic energy per quark, and $`\alpha `$ is the strong interaction coupling constant, assumed to have a value of $``$ 0.5. All higher order gluon-exchange interactions are neglected; their contribution is presumably subsumed in the bag constant whose value is poorly known. Two representative values are $`B=122`$ and 200 MeV fm<sup>-3</sup> .
The equilibrium conditions for uniform QM containing u, d, s-quarks, electrons and muons are:
$$\mu _u+\mu _e=\mu _d=\mu _s,\mu _\mu =\mu _e.$$
(20)
The energy densities of charge neutral QM and NM are compared in Fig.12. In the interesting region of $`\rho 1`$ fm<sup>-3</sup> the total energy density of quark matter is about 1200 MeV fm<sup>-3</sup>, of which only 122 or 200 MeV fm<sup>-3</sup> comes from the bag. Assuming the A18+$`\delta v`$+UIX EOS for NM the first order phase transition to QM is indicated at $``$ 1 baryon/fm<sup>3</sup>.
It has been recently suggested that at high densities QM may have color superconductivity resulting from non-perturbative attraction between quarks. In QM with only u and d-quarks this invariably leads to the possibility of a diquark condensate which breaks global color invariance . The associated color gap is estimated to be of the order of 100 MeV. When the s-quarks are included there are many possible phases .
## 5 Mixtures of Phases in Dense Matter
The phases of matter considered in the past sections are uniform and locally charge neutral, whereas bulk matter needs only to be charge neutral on average. For example, iron metal has positively charged regions occupied by iron nuclei, and the space in between is negatively charged by electrons. Generally the ground state of matter can have a mixture of regions occupied by different phases with the constrain of overall charge neutrality.
Matter in the outer crust of neutron stars, at $`\rho 0.002\rho _0`$, is believed to be like terrestrial matter made up of neutron rich nuclei in electron gas. The inner crust, on the other hand, is believed to have a mixture of regions with positively charged NM composed of neutrons and protons and PNM with only neutrons , both immersed in a nearly uniform electron gas. These mixtures occur in the density range of $``$ 0.002 to 0.6 $`\rho _0`$, beyond which uniform NM is believed to be the ground state. More recently mixed phases of NM and QM , and condensates of pions, kaons , or hyperons in NM have also been considered.
The Coulomb energy of a single phase uniform matter, due to fluctuations in the electron and hadron or quark densities, is negligible; however that of matter with mixed phases is not. For example, matter at $`\rho =0.3\rho _0`$ has drops of NM in PNM. Their size is determined by a competition between Coulomb and surface energies; large drops have too much Coulomb energy and small drops to much surface energy per nucleon. It is necessary to know the energy of the interface to predict the nature of the mixed phase region quantitatively. Those for the interface between NM and PNM have been calculated from energy-density functionals of NM ; while those for interfaces of NM and QM are not well estimated. In this section we discuss QM and NM mixed phases, to avoid duplicating the review by Pethick and Ravenhall , and because they can influence the mass limit of neutron stars.
### 5.1 Equilibrium Conditions for Coexistence of QM and NM
Neglecting surface and Coulomb effects, the equilibrium conditions for the coexistence of of QM and NM at zero temperature are that they have equal pressures, and it costs no energy to convert a neutron or a proton in NM into quarks in QM. The last condition amounts to
$$\mu _n=2\mu _d+\mu _u\mathrm{and}\mu _p=\mu _d+2\mu _u.$$
(21)
The electron density, and hence the $`\mu _e`$, is assumed to be the same in QM and NM, therefore the $`\beta `$-equilibrium conditions are:
$$\mu _n=\mu _p+\mu _e\mathrm{and}\mu _d=\mu _s=\mu _u+\mu _e.$$
(22)
And finally, the charge neutrality condition is given by:
$`f\rho _{QM}^{ch}+(1f)\rho _{NM}^{ch}=e(\rho _e+\rho _\mu ).`$ (23)
Here $`\rho _{QM,NM}^{ch}`$ are the charge densities of QM and NM, $`\rho _{e,\mu }`$ are the electron and muon densities, and $`f`$ is the fraction of space filled by QM. A graphical representation of these equilibrium conditions is given in ref. .
As was first pointed out by Glendenning , the QM and NM mixed phases can span a wide density region. For example, with the A18+$`\delta v`$+UIX NM EOS, the mixed phases begin at 4.6 (3.4) and end at 11.3 (9.1) $`\rho _0`$ for $`B=200`$ (122) MeV fm<sup>-3</sup> (Fig.12). At lower densities we have uniform NM and QM at higher. The QM fractional volume increases almost linearly from 0 to 1 in this transition range over which $`\mu _e0`$. The uniform QM has almost equal number of u, d, and s-quarks, and is essentially charge neutral without electrons. An important consequence of $`\beta `$-equilibrium is that the QM is negatively charged at the beginning of the transition, where $`\mu _e`$ is large. By immersing negatively charged drops of QM in the positively charged NM we can remove some of the high energy electrons and increase the proton fraction in NM. Both of these lower the energy density of matter.
### 5.2 Structure of Mixed Phase Matter
Matter with mixed phases has additional structure due to interfaces dividing the regions occupied by the two phases. The surface and Coulomb energies associated with these interfaces, neglected in above section, raise the energy of the mixed phase matter as well as determine the topology and structural length scales. The surface energy can presumably be estimated from the surface tension, $`\sigma `$. For NM $`\sigma 1`$MeV/fm<sup>2</sup> whereas for QM it poorly known; typical values are in the range 10-100 MeV/fm<sup>2</sup> ).
Denoting the dimensionality of the structures by $`D`$ ($`D=3`$ for droplets and bubbles, $`D=2`$ for rods and $`D=1`$ for plates) the surface energies are generally
$$_S=D\sigma \frac{4\pi }{3}R^2,$$
(24)
and, for $`D=3`$ the Coulomb energy, $`_C=(3/5)Z^2e^2/R`$, where $`R`$ the size of the structure and
$$Ze=(\rho _{QM}^{ch}\rho _{NM}^{ch})\frac{4\pi }{3}R^3,$$
(25)
is the excess charge of the droplet compared with the surrounding medium. General equations for $`_C(R,D)`$ are given in ref. . Minimizing the energy density with respect to $`R`$ we obtain the usual result that $`_S=2_C`$ at equilibrium. Minimizing with respect to the continuous dimensionality as well thus determines both $`R`$ and $`D`$. For droplets ($`D=3`$) the equilibrium radius is found to be:
$$R4.0\mathrm{fm}\left(\frac{\sigma }{1\mathrm{M}\mathrm{e}\mathrm{V}/\mathrm{fm}^2}\right)^{1/3}\left(\frac{\rho _{QM}^{ch}\rho _{NM}^{ch}}{e\rho _0/2}\right)^{2/3}.$$
(26)
A droplet of symmetric NM in vacuum has a surface tension $`\sigma =1`$ MeV$``$fm<sup>-2</sup> for which (26) gives $`R`$ 4 fm, which agrees with the fact that nuclei like <sup>56</sup>Fe are the most stable form of matter at low density. For QM droplets both the surface tension and charge densities are larger but the estimates of $`R`$ are similar.
At the beginning of the mixed phase region we expect that spherical droplets of QM, with $`R5`$ fm, will form a BCC lattice in uniform NM. They would have baryon number of the order of few hundred, and a negative charge of similar magnitude. As density increases and $`f`$ approaches 0.5, the drops would merge and form rods, which merge further on to form sheats. When $`f>0.5`$ the NM sheats break up into rods, and then into drops and eventually disappear when $`f=1`$. This scenario is similar to that in the inner crust; at low densities there are drops of NM occupying a small fraction of space. By $`\rho 0.6\rho _0`$ NM occupies all space via a similar set of mixed phases.
An other effect of the Coulomb and surface energies is that they decrease the density range covered by the mixed phase region. In particular, the lower density edge of this interesting region may be pushed up by almost $`\rho _0`$ if $`\sigma `$ is in the 10 to 50 MeV fm<sup>-2</sup> range . The energy density of the mixed phase matter is also raised by a few MeV fm<sup>-3</sup> in this case. Finally, if $`\sigma `$ were to be large $`(\stackrel{>}{}70`$ MeV fm$`{}_{}{}^{2})`$ the mixed phases may not be energetically favorable, and there will be a simple first order phase transition from NM to QM with a density discontinuity. One should bear in mind that even if the droplet phase were favored energetically, it may not be realized in practice if the time required to nucleate QM drops is too long compared to pulsar ages.
## 6 Neutron Star Observations and Predictions
The gross structure of neutron stars has been predicted using very many EOS, phenomenological as well as based on realistic models of nuclear forces . Of these we consider only those based on realistic models primarily because one can always find phenomenological energy density functionals or Lagrangians which reproduce their EOS.
Typical results for nonrotating stars with maximum mass, and with $`M=1.4M_{}`$, obtained by recent calculations, are listed in Table 2. The results for A18 without boost correction $`\delta v`$ are listed primarily for reference. This correction is unambiguous , and must be added to obtain reliable results. Those for A18+$`\delta v`$ are also to be taken less seriously, because it gives too large value for $`\rho _0`$. The TNI used with the Paris NN interaction is of the Urbana form with parameters determined by reproducing the empirical SNM properties (see sec. 3.4). In A18+UIX and Paris+TNI models the $`\delta v`$ is not considered explicitly; it is approximately subsumed in the TNI fitted to data. Of the three Bonn models, Bonn A comes closest to reproduce the empirical properties of SNM with Dirac-Brueckner (DBHF) method. These calculations include the $`\delta v`$ as well as many-body forces generated via $`Z`$-diagrams.
The A18+$`\delta v`$+UIX, A18+UIX, Paris+TNI and Bonn-A (DBHF) models come close to reproducing the empirical $`\rho _0`$; the later two fit the SNM binding energy; while the former models fit binding energies of light nuclei via exact calculations, since the energy of SNM can not yet be calculated reliably. Nevertheless these four “realistic” models of NM give rather similar results which are not too different from those of the 1988 calculations of Wiringa, Fiks and Fabrocini with the older Urbana-Argonne interactions now replaced with A18 and UIX.
The effect of the possible appearance of QM drops in high density matter has been studied with the A18+$`\delta v`$+UIX model. The $`M_m`$ is reduced to 2.02 and 1.91 $`M_{}`$ for bag-constant values $`B=`$ 200 and 122 MeV Fm<sup>-3</sup> respectively, while the predictions for 1.4 $`M_{}`$ stars remain unchanged. Presence of either kaons or hyperons in dense matter is unlikely to have much of an effect on the 1.4 $`M_{}`$ stars due to their low central density, while that on the mass limit is difficult to estimate quantitatively. For example, if kaons were to condense in matter at $`\rho =5\rho _0`$ and limit $`\rho _c`$ to $`<`$ 5 $`\rho _0`$, the $`M_m`$ of the four realistic models will drop to $``$ 2.0, 2.3, 1.7 and 2.0 $`M_{}`$ respectively; while if hyperons were to lower the energy of matter at $`\rho =5\rho _0`$ by 25 MeV per baryon, the $`M_m`$ would be reduced by $``$ 0.2 $`M_{}`$.
The mass radius relation obtained with models based on the A18 interaction are shown in Fig.13. Results of A18 and A18+$`\delta v`$ models are given primerily for comparison. As expected the harder EOS give larger $`M_m`$ and predict larger radii. The differences between the radii predicted by the realistic models is only $``$ 10 %.
### 6.1 The Mass Limit
The observed mass of Hulse-Taylor pulsar B1913+16 of 1.4411 $`\pm `$ 0.00035 shows that $`M_m>1.44M_{}`$. All the radio pulsars in known neutron star, and neutron star-white dwarf binaries have masses with lower limits less than $`1.44M_{}`$. The X-ray pulsar Vela X-1, which orbits a supergiant, however is consistently estimated to have a larger mass of $`1.9`$. The motion of this star is perturbed from being pure Keplerian, presumably by tidal forces exerted by the neutron star, and its present mass estimate, $`1.87_{0.17}^{+0.23}`$ , indicates that $`M_m>1.7M_{}`$ at 95 % confidance level. Finally, if the QPO’s indeed originate from the innermost stable orbit , then $`M_m>2M_{}`$. These mass limits are compatible with predictions of realistic NM models.
On the other hand there is no evidence that SN 1987A produced a neutron star. Its observed luminosity is now well below the 10<sup>38</sup> ergs/s Eddington limit, suggesting that no neutron star was produced in this supernovae . If we assume that the total mass, $`M_{Tot}`$, of the collapsed core plus the matter that fell back on to the core after the explosion, went into a black hole, then neutron star $`M_m`$ must be less than $`0.9M_{Tot}`$. The factor 0.9 takes into account the $``$ 10 % gravitational binding energy of the neutron star. Bethe and Brown estimate $`M_{Tot}1.73M_{}`$ using supernovae calculations by Wilson and Mayle, and conclude that $`M_m<1.56M_{}`$. Uncertainties in these arguments have been discussed by Zampieri et. al. . If the conclusion is found to be valid, then there must be other explanations for the Vela X-1 observations and the QPO, and the NM prediction for the $`M_m`$ is too large.
### 6.2 Temperatures, Cooling and Radii
Neutron stars are born with interior temperatures of the order $`10^{12}`$ K, but cool rapidly via neutrino emission to temperatures of the order $`10^{10}`$ K within minutes and $`\stackrel{<}{}10^6`$ K in $`10^5`$ yr. Spectra observed in X-ray or UV bands for nearby pulsars have in some cases black-body components from which surface temperatures of order $`T10^6`$ K are extracted for pulsars of age $`10^310^6`$ years. It is, however, unclear how much of the observed radiation is due to pulsar phenomena, to a synchrotron-emitting nebula or to the neutron star itself. In other cases upper limits have been set from the absence of X-rays. The surface temperatures are compatible with predictions from standard modified URCA cooling processes
$$n+nn+p+e^{}+\overline{\nu }_e,n+p+e^{}n+n+\nu _e.$$
(27)
Faster cooling processes as direct URCA or due to quark matter, kaon or hyperon condensates generally lead to considerably lower temperatures . To be consistent with observed surface temperatures the exotic coolant can only exist in a minor portion of the neutron star or it is superfluid whereby cooling is suppressed by factors of $`\mathrm{exp}(\mathrm{\Delta }/T)`$, where $`\mathrm{\Delta }`$ is the pairing gap.
The Hubble Space Telescope (HST) has observed one thermally radiating neutron star RX J185635-3754 with surface temperature $`T6\times 10^5`$ K$`50`$ eV . Its distance is less than 120 pc from Earth and should soon be determined more accurately by HST parallax measurements. Circumstancial evidence indicate a distance of $``$80 pc which leads to a black-body radius of $`1213`$ km from its luminosity and temperature. Such radii would agree well with predictions of realistic NM EOS (Fig.13) for $`M12M_{}`$.
### 6.3 Glitches and Superfluidity
Sudden spin jumps, called glitches, superimposed upon otherwise gradual spin down have been observed in most of the younger isolated pulsars . Since their discovery the Crab and Vela pulsars have each produced about a dozen glitches with period changes $`\mathrm{\Delta }P/P`$ of the order of $`10^8`$ and $`10^6`$ respectively. In post-glitch relaxation most of the period increase $`\mathrm{\Delta }P`$ decays.
Many mechanisms have been proposed to explain the glitches . The most plausible of these attributes glitches to the angular momentum stored in the rotating superfluid neutrons in the inner crust . The magnetic torque slows down the crust and most of the star except for these superfluid neutrons. Their angular momentum is stored in vortices pinned to nuclei in the inner core, until an instability occurs that leads to vortex depinning and sudden angular momentum transfer to the crust, leading to the glitch. At subnuclear densities in the crust, $`{}_{}{}^{1}S_{0}^{}`$ pairing between neutrons leads to gaps of order $``$1 MeV . In NM at $`\rho >\rho _0`$ this pairing gap vanishes, but $`{}_{}{}^{3}P_{2}^{}`$ pairing of neutrons and $`{}_{}{}^{1}S_{0}^{}`$ pairing of protons may occur .
The size of the glitches sets a lower limit on the moment of inertia of the superfluid in the inner crust which in turn sets a lower limit on the neutron star radius for a given mass . Assuming that the mass of Vela pulsar is $`1.4M_{}`$, a conservative limit on its radius is $`R9`$ km; it is compatible with predictions of most EOS.
## 7 Conclusions
Since the discovery of pulsars a significant effort has been devoted to accurately calculate properties of dense NM from realistic models of nuclear forces. Exact calculations of NM are still out or reach, however the new AFDMC methods (sect. 3.7) may eventually succeed. The present variational upper bounds seem to be above the true energies by $``$ 12 %. Such an error does not have serious consequences on the predicted properties of neutron stars. For example, an EOS obtained by reducing the variational energies, without rest mass terms, by 12 % reduces the maximum mass of A18+$`\delta v`$+UIX model by 2.3 % to 2.14 $`M_{}`$, and the radius of 1.4 $`M_{}`$ star by 2.9 % to 11.2 km. Larger uncertainties stem from the fact that the double $`\pi ^0`$ and $`\pi ^{}`$ condensation scenario illustrated in Fig.11 has not yet been calculated with realistic interactions, though it appears unlikely that it will influence the NM EOS by much more than 10 %.
Local models of two-nucleon interaction seem to be now converging. The predictions based on the 1988 calculations with Argonne 14 interaction are not too different from those of the 1998 calculations with the more accurate A18. It also seems likely that the local models give a fairly accurate description of two-nucleon interaction. A concern is that the present models of TNI are based on fits to a rather limited set of data, and are not as precise as the NN-interaction models. However, addition of the UIX TNI to the A18+$`\delta v`$ increases the maximum mass by $``$ 20 % and $`R(1.4)`$ by 13 % (Table 2). These changes may be important but they are not very large.
The present models of kaon-nucleon and hyperon-nucleon interactions are based on very limited data, and we have none on $`K^{}NN`$ and $`\mathrm{\Sigma }^{}NN`$ three-body forces. These could have significant effect on the threshold densities for kaons and hyperons to appear in dense matter. Hopefully advances in QCD and quark-models will provide a more rigorous framework to describe these interactions, and calculate properties of quark matter. The bag model estimates of QM EOS may have significant corrections at densities of interest in neutron stars.
¿From present observations there seem to be three possible scenarios for the limiting mass of neutron stars. If QPO’s are indeed due to accretion from the innermost stable orbit, then the NM predictions of $`M_m2.2M_{}`$ are reasonable, and strange baryons and quark drops do not soften the EOS of matter at $`\rho <7\rho _0`$ significantly. If the Vela X-1 mass measurement is correct, but QPO’s have some other origin, then $`M_m`$ could be $``$ 1.8 $`M_{}`$, indicating some softening of the NM EOS. However, if the present interpretation of QPO’s and Vela X-1 mass measurements are both faulty, and $`M_m`$ is as small as 1.56 $`M_m`$ as estimated from the absence of a neutron star in SN 1987A, then a significant softening of the NM EOS by phase transitions is indicated. Further observations will hopefully clear this situation.
Phase transitions such as NM to QM, can soften the EOS significantly. Fortunately these can have a measurable effect on the spin down of a rapidly rotating star in favorable cases, as has been recently pointed out . Consider the case of a rapidly rotating star whose central density is close to a first order phase transition. As the star slows and the central pressure increases due to decrease of the centrifugal force, the core matter will change its phase and become more dense at a critical angular velocity $`\mathrm{\Omega }_c`$. This decreases the moment of inertia, which assumes the characteristic form:
$$I=I_0\left(1+c_1\mathrm{\Omega }^2c_2(\mathrm{\Omega }_c^2\mathrm{\Omega }^2)^{3/2}+\mathrm{}\right).$$
(28)
for $`\mathrm{\Omega }<\mathrm{\Omega }_c`$. Here, $`c_1`$ and $`c_2`$ are small parameters proportional to the density difference between the two phases, and $`c_2=0`$ for $`\mathrm{\Omega }>\mathrm{\Omega }_c`$.
In order to make contact with observation, the temporal behavior of angular velocities must be considered. The pulsars slow down at a rate given by the loss of rotational energy, believed to be given by: $`d(\frac{1}{2}I\mathrm{\Omega }^2)/dt\mathrm{\Omega }^{n+1}`$, where $`n=3`$ for dipole radiation, Eq. (1) and $`n=5`$ for gravitational radiation. With the moment of inertia given by Eq. (28) the angular velocity can be calculated. The corresponding braking index, $`n(\mathrm{\Omega })=\ddot{\mathrm{\Omega }}\mathrm{\Omega }/\dot{\mathrm{\Omega }}^2`$, depends on the second derivative of the moment of inertia, $`I^{\prime \prime }=dI/d^2\mathrm{\Omega }`$. Using Eq. (28) we obtain:
$`n(\mathrm{\Omega })nc_1\mathrm{\Omega }^2+c_2{\displaystyle \frac{\mathrm{\Omega }^4}{\sqrt{\mathrm{\Omega }_c^2\mathrm{\Omega }^2}}}.`$ (29)
which exhibits a characteristic $`(\mathrm{\Omega }_c\mathrm{\Omega })^{1/2}`$ singularity as $`\mathrm{\Omega }`$ approaches $`\mathrm{\Omega }_c`$ from below. Observations of the braking index of a rapidly rotating, new born pulsar would be very interesting.
All realistic NM EOS predict that the radius of neutron stars with a mass of 1 to 1.5 $`M_{}`$ is $``$ 11 to 12 km. Future high resolution Chandra and XMM space observatories will hopefully be able to measure black-body spectra and detect gravitationally redshifted spectral lines from several stars. Such observations will help determine masses, radii and temperatures uniquely if the distance of the star is known. It is important to know the radius of a 1.4 $`M_{}`$ star, because that would test the EOS in the $`\rho \stackrel{<}{}3\rho _0`$ region in which large modifications of NM EOS are not expected on the basis of our present, naive estimates of kaon-nucleon and $`\mathrm{\Sigma }^{}`$-nucleon interactions.
## 8 Acknowledgements
The authors would like to thank J. Carlson, L. Engvik, S. Fantoni, M. Hjorth-Jensen, F. Lamb, S. Pieper, S. Shapiro and R. Wiringa for discussions and communications. This work has been partly supported by US National Science Foundation under Grant PHY 98-00978.
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# 1 Introduction
## 1 Introduction
The Hera collider enables the study of deep inelastic neutral current scattering $`epeX`$ at very high squared momentum transfers $`Q^2`$, thus probing the structure of $`eq`$ interactions at very short distances. At large scales new phenomena not directly detectable may become observable as deviations from the Standard Model predictions. A convenient tool to assess the experimental sensitivity beyond the maximal available center of mass energy and to parameterise indirect signatures of new physics is the concept of four-fermion contact interactions. Possible sources of such contact terms are either a substructure of the fermions involved or the exchange of a new heavy particle. In the first case a compositeness scale can be related to the size of the composite object, while in the second case the scale parameter is related to the mass and coupling of the exchanged boson.
The principle idea of this contact interaction analysis at Hera is to fix the Standard Model and its parameters, in particular the parton distributions, using experimental data at low $`Q^2`$, where the theory is well established, and to extrapolate the prediction towards high momentum transfers, where deviations due to new physics are expected to be most prominent. In the present paper the differential cross section $`\mathrm{d}\sigma /\mathrm{d}Q^2`$ is analysed over a $`Q^2`$ range of $`20030,000\mathrm{GeV}^2`$ and possible deviations from the Standard Model prediction are searched for in the framework of $`(\overline{e}e)(\overline{q}q)`$ contact interactions. The data are interpreted within conventional scenarios such as model independent compositeness scales of various chiral structures, a classical quark form factor approach and the exchange of heavy leptoquarks. Another investigation concerns the search for low scale quantum gravity effects, which may be observable at Hera via the exchange of gravitons coupling to Standard Model particles and propagating into extra spatial dimensions.
## 2 Data Analysis
The contact interaction analysis is based on the recent $`e^+p`$ neutral current cross section measurements by the H1 experiment discussed in detail in ref. . The data have been collected at a center of mass energy of $`\sqrt{s}=300\mathrm{GeV}`$ and correspond to an integrated luminosity of $`=35.6\mathrm{pb}^1`$, representing a tenfold increase over a previous analysis . The cross section $`\mathrm{d}\sigma /\mathrm{d}Q^2`$ is determined from a purely inclusive measurement of the final state positron with energy $`E_e^{}`$ and polar angle $`\theta _e`$ (defined with respect to the proton direction). The squared momentum transfer is calculated via $`Q^2=4E_eE_e^{}\mathrm{cos}^2(\theta _e/2)`$, where $`E_e`$ is the lepton beam energy. The data are corrected for detector effects and QED radiation and represent the cross section within the kinematic phase space of momentum transfer $`Q^2200\mathrm{GeV}^2`$ and inelasticity $`y=1E_e^{}/E_e\mathrm{sin}^2(\theta _e/2)<0.9`$. The dominant experimental systematics are the uncertainties of the reconstructed positron energy scale, varying between $`0.7\%`$ and $`3\%`$ depending on the detector region, and of the scattering angle, known to $`13`$ mrad. An overall normalisation uncertainty of $`1.5\%`$ is due to the luminosity determination. The differential cross section is shown in figure 1.
The double differential cross section is given in the Standard Model by
$`{\displaystyle \frac{\mathrm{d}^2\sigma (e^+pe^+X)}{\mathrm{d}x\mathrm{d}Q^2}}`$ $`=`$ $`{\displaystyle \frac{2\pi \alpha ^2}{xQ^4}}\left\{Y_+F_2(x,Q^2)Y_{}xF_3(x,Q^2)y^2F_L(x,Q^2)\right\},`$ (1)
where $`x=Q^2/ys`$ is the Bjorken scaling variable and $`Y_\pm =1\pm (1y)^2`$. The generalised proton structure functions $`F_2(x,Q^2)`$, $`F_3(x,Q^2)`$ and $`F_L(x,Q^2)`$ are related to the parton densities and the quark-$`\gamma `$ and quark-$`Z`$ couplings. The cross section calculations are done in the DIS scheme in next-to-leading-order using as standard the CTEQ5D parton parameterisation . Integrating eq. (1) over $`x`$ gives the $`Q^2`$ spectrum which describes the data very well over six orders of magnitude, see figure 1.
In order to derive quantitative tests of the Standard Model and to search for new physics hypotheses, a $`\chi ^2`$ analysis of the data is performed taking the dominant error sources and uncertainties into account. The $`\chi ^2`$ function is defined as
$`\chi ^2`$ $`=`$ $`{\displaystyle \underset{i}{}}\left({\displaystyle \frac{\widehat{\sigma }_i^{exp}f_n\widehat{\sigma }_i^{th}(1_k\mathrm{\Delta }_{ik}(\epsilon _k))}{\mathrm{\Delta }\widehat{\sigma }_i^{exp}f_n}}\right)^2+\left({\displaystyle \frac{f_n1}{\mathrm{\Delta }f_n}}\right)^2+{\displaystyle \underset{k}{}}\epsilon _k^2.`$ (2)
Here $`\widehat{\sigma }_i`$ denotes the experimental or theoretical cross section in the $`Q^2`$ bin $`i`$ and $`f_n`$ is the overall normalisation parameter with an uncertainty $`\mathrm{\Delta }f_n=0.015`$. The experimental error $`\mathrm{\Delta }\widehat{\sigma }_i^{exp}`$ includes statistical and uncorrelated systematic errors added in quadrature. The functions $`\mathrm{\Delta }_{ik}(\epsilon _k)`$ describe for the $`i^{th}`$ bin effects due to correlated systematic errors associated to different sources $`k`$. They depend quadratically on the fit parameters $`\epsilon _k`$, which may be interpreted as pulls, i.e. shifts caused by systematics normalised to their error estimates. There are three sources of correlated systematic errors taken into account: the experimental uncertainties of the positron energy scale and the scattering angle and the uncertainty of the strong coupling entering in the Standard Model prediction (see below).
Concerning cross section calculations the major uncertainty comes from the parton distributions, which are generally provided without error estimates. Different parametrisations in the DIS scheme, MRST 99 and GRV 94 in addition to CTEQ5D, are used to estimate the uncertainties due to various models and assumptions. They do not differ in the shape of the $`Q^2`$ spectrum significantly, but rather in the absolute cross section prediction by up to 2.8%, e.g. comparing CTEQ5D with MRST 99. Several other MRST sets are used for cross checks, like those with different admixtures of quarks and gluons at high $`x`$, or different treatments of strange and charm quarks. All these MRST variants essentially change the overall normalisation of the cross section prediction by less than $`1\%`$, being well below the measurement errors, and introduce no relevant additional $`Q^2`$ dependence. The largest uncertainty comes from the strong coupling constant. Using parton distributions evaluated for couplings differing from the central value of $`\alpha _s(M_Z)=0.118`$ by $`\pm 0.005`$ cause variations of the cross section by $`\pm 1\%`$ at low $`Q^2`$ and up to $`4\%`$ at high $`Q^2`$. These shifts are parameterised and taken into account as correlated systematic error in the $`\chi ^2`$ fit of eq. (2). It should be noted that the applied parton density functions have not been constrained by high $`Q^2`$ data from the Hera experiments. A comparison with a recent QCD analysis in the $`\overline{\text{MS}}`$ scheme , which attempts to provide parton distributions including errors, confirms that the above choice of various parton density functions is a reasonable representation of systematic uncertainties.
A fit of the cross section $`\mathrm{d}\sigma /\mathrm{d}Q^2`$ to the Standard Model expectation using CTEQ5D parton densities yields $`\chi ^2/\mathrm{dof}=12.3/16`$ with a normalisation parameter $`f_n=1.004`$. Limits of a model parameter are derived by varying the parameter until the $`\chi ^2`$ value changes by a certain amount with respect to the Standard Model fit, e.g. $`\chi ^2\chi _{SM}^2=3.84`$ for 95% confidence level (CL). Systematics due to different parton distributions are taken into account by always quoting the most conservative result of the various fits, i.e. the smallest value in case of a lower limit.
## 3 Contact Interaction Phenomenology
New currents or heavy bosons may produce indirect effects through the exchange of a virtual particle interfering with the $`\gamma `$ and $`Z`$ fields of the Standard Model. For particle masses and scales well above the available energy, $`\mathrm{\Lambda }\sqrt{s}`$, such indirect signatures may be investigated by searching for a four-fermion pointlike $`(\overline{e}e)(\overline{q}q)`$ contact interaction. The most general chiral invariant Lagrangian for neutral current vector-like contact interactions can be written in the form
$`_V`$ $`=`$ $`{\displaystyle \underset{q=u,d}{}}\{\eta _{LL}^q(\overline{e}_L\gamma _\mu e_L)(\overline{q}_L\gamma ^\mu q_L)+\eta _{LR}^q(\overline{e}_L\gamma _\mu e_L)(\overline{q}_R\gamma ^\mu q_R)`$ (3)
$`+\eta _{RL}^q(\overline{e}_R\gamma _\mu e_R)(\overline{q}_L\gamma ^\mu q_L)+\eta _{RR}^q(\overline{e}_R\gamma _\mu e_R)(\overline{q}_R\gamma ^\mu q_R)\},`$
where the indices $`L`$ and $`R`$ denote the left-handed and right-handed fermion helicities and the sum extends over up-type and down-type quarks and antiquarks $`q`$. In deep inelastic scattering at high $`Q^2`$ the contributions from the first generation $`u`$ and $`d`$ quarks completely dominate and contact terms arising from sea quarks $`s`$, $`c`$ and $`b`$ are strongly suppressed. Thus, there are eight independent effective coupling coefficients, four for each quark flavour
$`\eta _{ab}^q`$ $``$ $`ϵ{\displaystyle \frac{g^2}{\mathrm{\Lambda }_{ab}^{q2}}},`$ (4)
where $`a`$ and $`b`$ indicate the $`L,R`$ helicities, $`g`$ is the overall coupling strength, $`\mathrm{\Lambda }_{ab}^q`$ is a scale parameter and $`ϵ`$ is a prefactor, often set to $`ϵ=\pm 1`$, which determines the interference sign with the Standard Model currents. The ansatz eq. (3) can be easily applied to any new phenomenon, e.g. $`(eq)`$ compositeness, leptoquarks or new gauge bosons, by an appropriate choice of the coefficients $`\eta _{ab}`$. Scalar and tensor interactions of dimension 6 operators involving helicity flip couplings are strongly suppressed at Hera and therefore not considered.
It has been recently suggested that gravitational effects may become strong at subatomic distances and thus measurable in collider experiments . In such a scenario, which may be realised in string theory, gravity is characterised by a scale $`M_S𝒪(\mathrm{TeV})`$ in $`4+n`$ dimensions. The extra spatial dimensions $`n`$ are restricted to a volume associated with the size $`R`$ and the scales in $`4+n`$ and the ordinary $`4`$ dimensions are related by
$`M_P^2`$ $``$ $`R^nM_S^{2+n},`$ (5)
where $`M_P10^{19}\mathrm{GeV}`$ is the Planck mass. An exciting consequence would be a modification of Newton’s law at distances $`r<R`$, where the gravitational force would rise rapidly as $`F1/r^{2+n}`$ and become strong at the scale $`M_S`$. Experimentally, gravity is essentially not tested in the sub-millimeter range and scenarios with $`n>2`$ extra dimensions at large distances $`R100\mu `$m are conceivable.
In the phenomenology of the Standard Model particles are confined to 4 dimensions while only the graviton propagates as massless spin 2 particle into the $`n`$ extra dimensions. When projected onto the normal 4 dimensional space the graviton appears as a spectrum of Kaluza-Klein excitations with masses $`m^{(j)}=j/R`$, including the zero-mass state. The graviton fields $`G_{\mu \nu }^{(j)}`$ couple to the Standard Model particles via the energy-momentum tensor $`T_{\mu \nu }`$
$`_G`$ $`=`$ $`{\displaystyle \frac{\sqrt{8\pi }}{M_P}}G_{\mu \nu }^{(j)}T^{\mu \nu }.`$ (6)
Summation over the whole tower of Kaluza-Klein states $`j`$ with masses up to the scale $`M_S`$ compensates the huge $`1/M_P`$ suppression and results in an effective contact interaction coupling
$`\eta _G`$ $`=`$ $`{\displaystyle \frac{\lambda }{M_S^4}},`$ (7)
where $`\lambda `$ is the coupling strength of order unity. Note that the scale dependence of gravitational effects is very different from that of conventional contact interactions, eq. (4). In deep inelastic scattering graviton exchange may contribute to the electron-quark subprocess, but the new interaction also induces electron-gluon scattering which is not present in the Standard Model.
It is worth recalling that contact interactions as an effective theory can only be formulated in lowest order. They contribute to the structure functions $`F_2(x,Q^2)`$ and $`xF_3(x,Q^2)`$, but are absent in $`F_L(x,Q^2)`$. On the other hand a cross section calculation in next-to-leading-order QCD appears to be more reliable. This conceptual limitation is less important in the DIS renormalisation scheme, where the expression for the dominant structure function $`F_2`$ is identical and $`xF_3`$ receives only minor corrections in second order.
Contact interaction phenomena are best observed as a modification of the expected $`Q^2`$ dependence and all information is essentially contained in the differential cross section $`\mathrm{d}\sigma /\mathrm{d}Q^2`$, analysed in the present paper. Calculations using the Standard Model prediction, eq. (1), show that for the scenarios under study only those models involving both $`u`$ and $`d`$ quarks with pure $`LL`$ or $`RR`$ couplings and negative interference are slightly more sensitive to an analysis in two variables $`Q^2`$ and $`x`$. With the present luminosity the gain in setting limits on the respective scales would be $`20\%`$ for the negatively interfering $`LL`$ and $`RR`$ composite models and $`10\%`$ for the leptoquark $`S_1^L`$. For all other scenarios the differences are marginal.
## 4 Compositeness Scales
In the Standard Model the fundamental particles – leptons, quarks and gauge bosons – are assumed to be pointlike. A possible fermion compositeness or substructure can be expressed through the $`\eta `$ coefficients of eq. (4) which depend only on the ratio of the coupling constant over the scale. In the present analysis the interference sign is set to $`ϵ=\pm 1`$ for the chiral structures under study, the coupling strength $`g`$ is by convention chosen as $`g^2=4\pi `$ and the compositeness scale $`\mathrm{\Lambda }`$ is assumed to be the same for all up-type and down-type quarks. Among the many possible combinations the following models are investigated: (i) the pure chiral couplings $`LL`$, $`LR`$, $`RL`$ and $`RR`$, (ii) the vectorial couplings $`VV`$, $`AA`$ and $`VA`$, (iii) the mixtures $`LL+RR`$ and $`LR+RL`$.
It is appropriate to analyse the differential cross section in terms of the coupling coefficients $`\eta =ϵ\mathrm{\hspace{0.17em}4}\pi /\mathrm{\Lambda }^2`$. Figure 2 shows the values of $`\chi ^2`$ as a function of $`ϵ/\mathrm{\Lambda }^2`$ from fits to the models under study. In general one observes that the distributions become narrower, i.e. the sensitivity increases, the more chiral structures are involved. The pure chiral couplings prefer negative values of $`\eta `$. This is a consequence of the trend of the data $`(\mathrm{d}\sigma /\mathrm{d}Q^2)/(\mathrm{d}\sigma ^{SM}/\mathrm{d}Q^2)`$ to be slightly low around $`Q^24,00012,000\mathrm{GeV}^2`$ and being followed by an upward fluctuation at higher $`Q^2`$ (see figure 1), which favour a negative interference term. Note that the $`LL`$ and $`RR`$ models and the $`LR`$ and $`RL`$ models are almost indistinguishable in deep inelastic unpolarised $`e^+p`$ scattering. Within each couple the exchanged quantum numbers are the same and therefore the combinations $`LL+RR`$ and $`LR+RL`$ are investigated as well. The data are more sensitive to the $`VV`$, $`AA`$ and $`VA`$ models, where all chiral structures contribute. The most restrictive range of $`ϵ/\mathrm{\Lambda }^2`$ is obtained for the $`VV`$ model, where all contact terms enter with the same sign. Figure 2 also shows that different parton distributions have little influence on the results.
The results of the $`\chi ^2`$ fits are shown in figure 3 and compiled in table 1. Within two standard deviations the couplings $`ϵ/\mathrm{\Lambda }^2`$ are compatible with the Standard Model for all parton density functions used.
Limits on the compositeness scale parameters $`\mathrm{\Lambda }^+`$ and $`\mathrm{\Lambda }^{}`$, corresponding to positive and negative interference, are quoted in table 1 and also presented in figure 3. They vary between $`1.3\mathrm{TeV}`$ and $`5.5\mathrm{TeV}`$ depending on the chiral structure of the model and are in most cases asymmetric with respect to the interference sign. In general the values of $`\mathrm{\Lambda }^+`$ are more restrictive due to the downward trend of the data at intermediate $`Q^2`$ prefering models with negative interference. As an illustration of the sensitivity of the data to compositeness scales figure 4 shows the 95% CL contributions of $`\mathrm{\Lambda }^\pm `$ for the $`VV`$ model using CTEQ5D parton densities.
The results of direct searches for $`(eq)`$ compositeness are compatible with those of other experiments at Hera, Lep and Tevatron. To date the most stringent limits come from $`e^+e^{}`$ experiments with typical cut-off values of $`310\mathrm{TeV}`$ assuming, however, in general the same scale parameter $`\mathrm{\Lambda }`$ for all five active quarks. The Zeus collaboration investigates only models in which at least 2 couplings contribute and derives limits between $`1.7\mathrm{TeV}`$ and $`5\mathrm{TeV}`$ relying solely on the shape of measured distributions. The $`p\overline{p}`$ experiments measure Drell-Yan production and quote limits on $`\mathrm{\Lambda }`$ between $`2.5\mathrm{TeV}`$ and $`6\mathrm{TeV}`$, where the normalisation is based on the number of observed $`Z`$ bosons. Model dependent indirect limits of order $`10\mathrm{TeV}`$ for the pure chiral couplings involving $`u`$ and $`d`$ quarks can be set by atomic parity violation experiments .
## 5 Form Factors
An alternative method to study possible fermion substructures is to assign a finite size of radius $`R`$ to the electroweak charges of leptons and/or quarks while treating the gauge bosons $`\gamma `$ and $`Z`$ still as pointlike particles . A convenient parametrisation is to introduce ‘classical’ form factors $`f(Q^2)`$ at the gauge boson–fermion vertices, which are expected to diminish the Standard Model cross section at high momentum transfer
$`f(Q^2)`$ $`=`$ $`1{\displaystyle \frac{1}{6}}r^2Q^2,`$ (8)
$`{\displaystyle \frac{d\sigma }{dQ^2}}`$ $`=`$ $`{\displaystyle \frac{d\sigma ^{SM}}{dQ^2}}f_e^2(Q^2)f_q^2(Q^2).`$ (9)
The root of the mean-square radius of the electroweak charge distribution, $`R=\sqrt{r^2}`$, is taken as a measure of the particle size. The data are analysed in terms of a single form factor $`f_q`$, i.e. only the quarks are allowed to be extended objects while the positron has no structure by setting $`f_e1`$. This assumption is justified, since the pointlike nature of the electron/positron is already established down to extremely low distances in $`e^+e^{}`$ and $`(g2)_e`$ experiments . The analysis yields an upper limit at 95% CL of the light quark radius of
$`R_q`$ $`<`$ $`1.710^{16}\mathrm{cm}.`$
The result is compatible with those from other measurements. In an analysis of Drell-Yan production of $`e^+e^{}`$ and $`\mu ^+\mu ^{}`$ pairs in $`p\overline{p}`$ scattering the CDF collaboration finds a quark size of $`R_q<110^{16}\mathrm{cm}`$ assuming pointlike leptons. A complementary analysis of the contributions of anomalous magnetic dipole moments to the $`Zq\overline{q}`$ vertex using hadronic $`Z`$ decays gives $`R_q<1.210^{16}\mathrm{cm}`$ for the light $`u`$ and $`d`$ quarks .
## 6 Leptoquarks
Leptoquarks are colour triplet bosons of spin 0 or 1, carrying lepton ($`L`$) and baryon ($`B`$) number and fractional electric charge. They couple to lepton–quark pairs and appear in extensions of the Standard Model which try to establish a connection between leptons and quarks. Leptons and quarks may be arranged in common multiplets, like in Grand Unified Theories or superstring motivated $`E_6`$ models, or they may have a common substructure as in composite models. A fermion number $`F=L+3B`$ is defined, which takes the values $`F=2`$ for leptoquarks coupling to $`e^{}q`$ and $`F=0`$ for leptoquarks coupling to $`e^{}\overline{q}`$. For positrons the fermion number $`F`$ changes by two. The leptoquark mass $`M_{LQ}`$ and its coupling $`\lambda `$ are related to the contact interaction coefficients of eq. (4) via $`g/\mathrm{\Lambda }=\lambda /M_{LQ}`$. The notation and the coupling coefficients $`\eta _{ab}`$ for leptoquarks <sup>1</sup><sup>1</sup>1The coupling coefficients are taken from ref. with the signs corrected (i.e. reversed) for $`F=2`$ scalar and $`F=0`$ vector leptoquarks according to ref. . are given in table 2. The only unknown parameter is the ratio $`M_{LQ}/\lambda `$. Note that the vector leptoquarks have coupling coefficients twice as large in magnitude compared to scalar leptoquarks.
The differential cross section analysis gives no evidence for a virtual leptoquark signal. The resulting lower limits on $`M_{LQ}/\lambda `$ are summarised in table 2 including the full error propagation and a variation of parton densities. In general leptoquarks with positive interference provide stronger limits compared to those with negative interference. This observation is consistent with the results found for compositeness models. The vector leptoquarks which couple to $`u`$ quarks provide the most restrictive limits of $`M_{LQ}/\lambda 𝒪(1\mathrm{TeV})`$. It should be emphasised that upper bounds on the coupling strength $`\lambda `$ can only be set for leptoquark masses exceeding the accessible center of mass energy of Hera. Masses far above 300 GeV are excluded for almost all types of leptoquarks with a coupling of $`\lambda 1`$.
These measurements are a considerable improvement over the previous analysis . But it should be noted that changes in the signs of couplings $`\eta _{ab}^q`$ reduce the sensitivity for $`F=2`$ vector leptoquarks and lead e.g. to weaker limits for $`\stackrel{~}{V}_0^R`$ and $`V_1^L`$ despite the increased luminosity. As an example of the sensitivity of the data to virtual leptoquark exchange figure 5 shows the contributions given by the lower limits on $`M_{LQ}/\lambda `$ for the scalar leptoquark $`S_{1/2}^R`$ and the vector leptoquark $`V_{1/2}^R`$. Both leptoquarks have $`RL`$ couplings to up and down quarks, which differ in magnitude and interference sign.
The present contact interaction results complement the direct leptoquark searches of the H1 collaboration , which have recently been extended beyond the kinematic reach of Hera up to masses of $`M_{LQ}400\mathrm{GeV}`$. The coupling limits derived in both analyses are compatible with each other in the mass region where they overlap. Virtual leptoquark exchange has also been studied in $`e^+e^{}`$ annihilation experiments at Lep . Typical limits on $`M_{LQ}/\lambda `$ are in the range $`0.31.8\mathrm{TeV}`$, but the sensitivity to particular leptoquark types is different from deep inelastic $`e^+p`$ scattering. In most cases the Lep results provide more stringent bounds; the limits for $`S_0^R`$, $`\stackrel{~}{S}_{1/2}^L`$ and $`V_0^R`$ are comparable and those of $`V_{1/2}^R`$ and $`\stackrel{~}{V}_{1/2}^L`$ are superior at Hera.
## 7 Large Extra Dimensions
The contributions of virtual graviton exchange to deep inelastic scattering have been derived from the cross sections given in ref. for $`e^+e^{}`$ collisions by applying crossing relations. The basic processes of elastic $`e^+qe^+q`$ and $`e^+ge^+g`$ scattering can be written as <sup>2</sup><sup>2</sup>2 The following formulae of ref. are used: eq. (79) for the $`eq`$ contributions, eq. (77) for the $`eg`$ contribution replacing the photons by gluons, and eqs. (A.5), (A.7) – (A.9) to expand the functions $`G_i(s/t)`$. The present $`eq`$ results are in agreement with , but the $`eg`$ part differs by a factor of $`1/4`$. The cross section formulae of refs. and cannot be confirmed and the results of both publications are inconsistent with each other.
$`{\displaystyle \frac{\mathrm{d}\sigma (e^+qe^+q)}{\mathrm{d}t}}`$ $`=`$ $`{\displaystyle \frac{\mathrm{d}\sigma ^{SM}}{\mathrm{d}t}}+{\displaystyle \frac{\mathrm{d}\sigma ^G}{\mathrm{d}t}}+{\displaystyle \frac{\mathrm{d}\sigma ^{\gamma G}}{\mathrm{d}t}}+{\displaystyle \frac{\mathrm{d}\sigma ^{ZG}}{\mathrm{d}t}},`$ (10)
$`{\displaystyle \frac{\mathrm{d}\sigma ^G}{\mathrm{d}t}}`$ $`=`$ $`{\displaystyle \frac{\pi \lambda ^2}{32M_S^8}}{\displaystyle \frac{1}{s^2}}\left\{32u^4+64u^3t+42u^2t^2+10ut^3+t^4\right\},`$ (11)
$`{\displaystyle \frac{\mathrm{d}\sigma ^{\gamma G}}{\mathrm{d}t}}`$ $`=`$ $`{\displaystyle \frac{\pi \lambda }{2M_S^4}}{\displaystyle \frac{\alpha e_q}{s^2}}{\displaystyle \frac{(2u+t)^3}{t}},`$ (12)
$`{\displaystyle \frac{\mathrm{d}\sigma ^{ZG}}{\mathrm{d}t}}`$ $`=`$ $`{\displaystyle \frac{\pi \lambda }{2M_S^4}}{\displaystyle \frac{\alpha }{s^2\mathrm{sin}^22\theta _W}}\left\{v_ev_q{\displaystyle \frac{(2u+t)^3}{tm_Z^2}}a_ea_q{\displaystyle \frac{t(6u^2+6ut+t^2)}{tm_Z^2}}\right\},`$ (13)
$`{\displaystyle \frac{\mathrm{d}\sigma (e^+ge^+g)}{\mathrm{d}t}}`$ $`=`$ $`{\displaystyle \frac{\pi \lambda ^2}{8M_S^8}}{\displaystyle \frac{u}{s^2}}\left\{2u^3+4u^2t+3ut^2+t^3\right\}`$ (14)
in an obvious notation of Standard Model (SM), pure graviton (G), $`\gamma G`$ and $`ZG`$ interference contributions. Here $`s`$, $`t=Q^2`$ and $`u`$ are the Mandelstam variables, $`e_q`$ is the quark charge and $`v_f`$ and $`a_f`$ are the vector and axial vector couplings of the fermions to the $`Z`$. The corresponding cross sections for $`e^+\overline{q}`$ scattering are obtained by replacing $`e_qe_q`$ and $`v_qv_q`$ in the expressions above. In order to get the inclusive $`e^+p`$ cross section the subprocesses have to be integrated over the $`x`$ dependent parton distributions, $`q(x)`$, $`\overline{q}(x)`$ and $`g(x)`$, and to be summed up
$`{\displaystyle \frac{\mathrm{d}\sigma (e^+pe^+X)}{\mathrm{d}Q^2}}`$ $`=`$ $`{\displaystyle dx\left\{q(x)\frac{\mathrm{d}\sigma (e^+q)}{\mathrm{d}t}+\overline{q}(x)\frac{\mathrm{d}\sigma (e^+\overline{q})}{\mathrm{d}t}+g(x)\frac{\mathrm{d}\sigma (e^+g)}{\mathrm{d}t}\right\}}.`$ (15)
Note, that expected gravitational effects arising from the gluon contribution are for the highest $`Q^2`$ values of the order percent compared to those coming from the quarks and antiquarks.
The strength of virtual graviton exchange to the cross section contributions is characterised by the ratio $`\lambda /M_S^4`$. The coupling $`\lambda `$ depends on the full theory and is expected to be of order unity. Also the sign of interference with the Standard Model particles is a priori not known. Therefore the coupling is set to $`\lambda =\pm 1`$, following the convention of . The data analysis is similar to the procedures described above. Gravitational effects are searched for by fitting the differential cross section to the above formulae with $`\lambda /M_S^4`$ treated as free parameter. The result of $`\lambda /M_S^4=3.3{}_{3.3}{}^{+4.2}{}_{1.3}{}^{+0.4}\mathrm{TeV}_{}^{4}`$, where the second error reflects the parton density variation, is compatible with the Standard Model expectation. Lower limits at 95% CL on $`M_S`$ for positive and negative coupling are then derived from the change in $`\chi ^2`$ with respect to the Standard Model fit, yielding
$`M_S>0.48\mathrm{TeV}`$ $`\mathrm{for}\lambda =+1,`$
$`M_S>0.72\mathrm{TeV}`$ $`\mathrm{for}\lambda =1.`$
Possible effects of graviton exchange to the data, as given by the exclusion limits, are illustrated in figure 6.
Similar investigations of virtual graviton effects in $`e^+e^{}`$ annihilation into fermion and boson pairs provide comparable limits . From the corresponding reaction of quark pair production scales of $`M_S`$ lower than $`0.50.65\mathrm{TeV}`$ can be excluded. Combining all reactions that lead to two-fermion final states limits approaching $`1\mathrm{TeV}`$ can be set.
## 8 Conclusions
Neutral current deep inelastic cross section measurements are analysed to search for new phenomena mediated through $`(\overline{e}e)(\overline{q}q)`$ contact interactions. No significant signal for compositeness, a quark form factor and virtual leptoquark or graviton exchange is found and the data are used to set limits which supersede and substantially improve former H1 results .
Limits on $`(eq)`$ compositeness are derived within a model independent analysis for scenarios involving one or more chiral couplings. The lower bounds on the scale parameters $`\mathrm{\Lambda }^\pm `$ range between $`1.3\mathrm{TeV}`$ and $`5.5\mathrm{TeV}`$ for a coupling strength $`g=\sqrt{4\pi }`$, depending on the chiral structure of the model.
A different approach to substructures is the concept of form factors. Such an analysis yields an upper limit of the size of the light $`up`$ and $`down`$ quarks of $`R_q<1.710^{16}\mathrm{cm}`$ assuming a pointlike lepton.
A study of virtual leptoquark exchange yields lower limits on the ratio $`M_{LQ}/\lambda `$ which for all types (except one) exceed the collider center of mass energy and approach $`1\mathrm{TeV}`$ for vector leptoquarks with couplings to up quarks. These measurements complement and extend the direct leptoquark searches at Hera to high masses $`M_{LQ}>\sqrt{s}`$.
In a search for possible effects of low scale quantum gravity with gravitons coupling to Standard Model particles and propagating into extra spatial dimensions, lower limits on the effective Planck scale $`M_S`$ of $`0.48\mathrm{TeV}`$ and $`0.72\mathrm{TeV}`$ for positive and negative coupling, respectively, are found.
Acknowledgements. We are grateful to the Hera machine group whose outstanding efforts have made and continue to make this experiment possible. We thank the engineers and technicians for their work in constructing and now maintaining the H1 detector, our funding agencies for financial support, the Desy technical staff for continual assistance, and the Desy directorate for the hospitality extended to the non–Desy members of the collaboration. We gratefully acknowledge valuable discussions with G.F. Giudice.
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# Flexible Linear Polyelectrolytes in Multivalent Salt Solutions: Solubility Conditions
## I Introduction
The precipitation and dissolution of linear polyelectrolytes with the addition of multivalent salt or particles have been extensively studied . In DNA the precipitation provides a promising mechanism to ”pack” long DNA, a critical problem in gene therapy. Moreover, it is correlated with highly accelerated rates of DNA renaturation and cyclisation. Therefore, it is important to determine the concentration of multivalent salt or particles at which the precipitated DNA dissolves in the solution. Furthermore, it is crucial to determine the effective charge of the DNA in the precipitated and dissolved states to determine their interaction with cells.
It has been shown experimentally that the transitions are nearly independent of the details of the linear polyelectrolyte (such as charge density, degree of flexibility, molecular weight, structure function, etc.). That is, flexible (single stranded DNA and Polystyrene Sulphonate) and semiflexible chains (double stranded DNA) can be described by the same thermodynamic model even though flexible chains collapse into amorphous dense spheres and long semiflexible chains into toroid conformations .
We recently develop a thermodynamic model that explains the universal nature of the multivalent counterion induced precipitation transition in low concentration of monovalent salts. In this paper we analyze the redissolution transition by extending our two state model to include large concentrations of multivalent salt.
The experimental diagrams of the logarithm of concentration of monomers $`\varphi `$ versus the logarithm of the concentration of multivalent salt $`m`$ have three regions (A, B and C) separated by 2 transitions lines (L1-L2 and L3) , shown schematically in Fig. 1. Region A corresponds to polyelectrolytes dissolved in water with expanded-stretched conformation at low $`m`$ values. Region B at intermediate multivalent salt concentrations $`m`$ is a solution of collapsed polyelectrolytes (sometimes considered as a coexistence of two phases: one rich and one poor in polyelectrolytes). Region C at large $`m`$ values contains polymers dissolved in water in expanded-coiled conformations. The transition line between the A and B regions, L2, at low monovalent salt concentrations is linear with a slope comparable to $`1/z`$, where $`z`$ is the valence of the counterions of the added salt. In the presence of monovalent salt there is a nearly horizontal transition line, L1, at very low concentrations of monomers. The transition line between regions B and C is horizontal (i.e., independent of the concentration of monomers) in the full regime. In our previous paper we discussed the transition between regions A and B and explained the transition in the regime L2 as the creation of collapsed conformations of polymers whose charge is compensated by condensed multivalent salt. Our model also predicts the transition L1 in the presence of large amounts of monovalent salt. In this paper we study the transition between the regions B and C in solutions with very low (negligible) monovalent salt concentrations.
It is well documented that in polyelectrolyte solutions a fraction of the counterions condensed along the polyelectrolytes to decrease their electrostatic energy . This ion condensation is crucial to understand the precipitation and redissolution transitions. In this paper we compute the fraction of condensed ions as the multivalent salt concentration increases in the two possible polyelectrolyte conformations: collapsed and expanded. Our two state model is based on the fact that in salt free and/or low ionic strength solutions, ion condensation leads to two possible conformation in linear flexible polyelelctrolytes: expanded-stretched with a reduced effective charge and collapsed nearly neutral with the cohesive energy of an ionic glass. Indeed, multivalent counterions lead to large cohesive energies for nearly neutral collapsed structures in region B explaining the transition L2. Since the magnitude of the electrostatic energy of a collapsed chain in region B is much larger than any entropic energy increase (due to the decrease of degrees of freedom resulting from the compaction), our model also describes the precipitation of semiflexible and rigid-rod polyelectrolytes into toroids and bundles, respectively, explaining the universality of the precipitation transition discussed in our previous work.
At the precipitation transition we predict that the expanded chains are slightly charged, and with further addition of multivalent salt, the collapsed chains (in region B) become practically neutral, in excellent agreement with recent electrophoresis experiments . As a result, the aggregates interactions in the dilute regime are negligible justifying a monomolecular collapse model to describe also multi-molecular aggregation.
The effective charge and size of the expanded conformation (in region C) changes with further addition of multivalent salt. The transition L3 occurs at high salt concentrations, and in such conditions the chains are expected to obey random walk or self-avoiding walks statistics due to the screening of the electrostatic interactions . Screening is responsible of the redissolution, as suggested in various models . We show here that the redissolution is actually very sensitive to the relation between the chemical potential of the condensing multivalent particles in the solution $`\mu `$ and the inverse screening length $`\kappa `$, while nearly insensitive to the degree of flexibility and molecular weight of the polyelectrolytes. Screening, for example, may be reduced due to multivalent-monovalent ion associating in the solution, strongly affecting the effective charge of the expanded chains. We determine the effective charge of the chains as a function of $`\mu `$ and $`\kappa `$ including the finite size of the ions and the discrete nature of the charge distribution along the polyelectrolyte, which is essential to compute the electrostatic energies of both the collapsed and expanded states. In a $`\mu \kappa `$ diagram the collapsed region lies between two branches of expanded states: one expanded with a reduced effective charge and the other expanded with an inverted charge. In most common situations the chains precipitate and redissolve within the same branch; i.e., with a reduced effective charge of the same sign of the bare charge. Our model predicts a redissolution transition independent of polyelectrolyte concentration and describes the re-dissolution of flexible, semiflexible and rigid polyelectrolytes. Indeed, the collapsed and expanded states discussed here are akin to the multi-molecular precipitated and dissolved states, respectively, observed in many polyelectrolytes , explaining the rather universal form of the diagram, Fig. 1.
In the Section 2 we summarize the assumptions of our work, and provide a description of each of the phases of the system. In section 3 we construct the free energies of two types of representative states and in Section 4 we give the result of the free energy minimization. In Section 5 we summarize our results, compare them with experiments and emphasize the elements of the theory that have not yet been subject to experimental results. We end with a brief conclusion in section 6.
## II Theoretical model
As in our previous work , we rely in a two state description of the system. Instead of considering all possible states of the system (describing them, for example by a changing scaling exponent), we simply assume that the minimum lies in one of two extremes: collapsed or expanded conformations. The expanded conformation, however, is a function of the concentration of salt in the solution. For example, at low monovalent salt concentrations the chains are expanded rods, and at large concentrations they are expanded coils . The free counterions and solvent can be integrated out provided we take into account the electrostatic interactions between the non-condensed counterions in the solution.
The condensed (collapsed) state is treated as an ionic glass so that its energy can be approximately calculated by techniques from solid state physics ; i.e, we assume that fluctuations play no important role and can be neglected. For a very large chain, at low enough temperatures, the bulk of the collapsed state should acquire an almost crystalline structure. The finite size of the chain, the connectivity of the monomers, and finite temperature effects must induce defects into the structure and the local structure should be glass-like, where the perfect order has not been achieved. An ionic glass structure was indeed observed in recent numerical studies of polyelectrolytes in bad solvents .
In both states of the systems, the collapsed and expanded states, it is necessary to keep track of the finite size of the counterions and monomers. This is clear in the case of the collapsed state, but even in the expanded states there is an important contribution arising from interactions with the nearest neighbors, and this energy can only be calculated by considering the finite size of the particles . For example, a naive coarse-grained model that takes into account only effective charges would assign a zero energy to a neutral cluster, while it can actually have a very strong cohesive energy. The redissolution into an expanded state occurs at high salt concentrations and the effective interactions between monomers are screened. In these conditions the energy associated with interactions between far away segments of the chain can be treated as a small perturbation with respect to that of the neutral polymer state.
The properties of the salt at different concentrations are summarized into two parameters: an effective chemical potential $`\mu `$, and an inverse screening length $`\kappa `$. We will present our results for general values of these parameters and in particular consider the case were the parameters are related by a Debye-Huckel law. We also discuss the effect of multivalent-monovalent ion association, the Bjerrum model , in the diagram.
Let us consider the series of states of a single polymer chain along a line of constant polymer concentration (line P in Fig. 1) with increasing multivalent salt concentration. We start in a state in region A with few counterions condensed. The condensed counterions lie very close to the chain, and cannot be considered as mixed in the solvent. The effective charge induces repulsions between the segments of the chain, and the chain acquires extra stiffness, making it into a rod-like structure.
When a small quantity of multivalent salt is added to the system, it dissociates, and the multivalent counterions are absorbed to the chain because the electrostatic energy to be gained from the condensation is larger (per unit charge) than that of the monovalent counterions. Furthermore, the entropic loss due to condensation is lower for the multivalent counterions. When the charge of the added salt roughly equals the bare charge of the polymers the effective charge of the chains quickly goes to zero. At this point most charges of the chain are compensated by the condensed multivalent counterions and a few monovalent ones. Since a group of charges with total charge zero has a minimal energy when they are set in a compact configuration, a collapsed conformation is acquired. This transition is of course also possible for the case of pure monovalent salt, but it requires either larger concentrations or lower temperatures to offset the entropic contributions. In other words, the determining factor for the transition is the strength of the electrostatic interaction, and the introduction of multivalent salt acts as a jump in parameter space and is not equivalent to a smooth increase in monovalent salt .
In the collapsed state the absence of uncompensated charges in the polymer reduces its solubility in water and polar solvents. It is not clear from experiments, however, whether the collapsed flexible polymers aggregate into large structures or if they form a dense solution of individually collapsed chains. Though we assume monomolecular collapse, the predicted transition points should not exhibit strong dependence on the number of chains aggregated since if the collapsed state is neutral, the residual interactions between chains (promoting or opposing multimolecular aggregates) are energetically small.
Since the collapsed structure is almost neutral, the increase in salt concentrations has little effect in the absolute value of the energy of the chain. The change in the environment affects only the interaction between the positive or negative small excess (or defect) of charge in the polymer, which would be located at the surface of the collapsed structure. The energy of the charges packed inside the structure is dominated by the interaction with their nearest neighbors and is not subject to screening (there are no floating ions between them). On the other hand, the energy of an expanded chain with non-zero effective charge will be greatly reduced by the increased screening.
For the purposes of comparison with the collapsed structures, we can consider a very rigid chain that cannot be collapsed into a compact structure. We can ignore for the moment the fact that this chain would likely form a bundle with other chains, and we can assume that we have isolated it by perhaps mechanical means. At higher multivalent salt concentrations, the rigid polymer continues absorbing ions, and thus it can become neutral and even overcharged (effective charge with opposite sign to the bare charge, as it occur in many systems ). With further increase in the amount of salt the trend can be reversed, and the chain can again become, first neutral, and then simply charged (same charge sign as the bare charge). The increase of salt concentration reduces the entropic penalty for the condensation of counterions. The counterions continue condensing into the polymer as long as there is a reduction in electrostatic energy due to the condensation. The increase in condensation stops when the typical screening length of the free counterions is of the order of the monomer separation, the free counterions can reduce their energy by associating with each other, and a fraction of them will do so instead of associating with the polymer.
At large salt concentrations, the chain reverts from collapsed to expanded, giving the transition $`BC`$. In the expanded state the chain needs not to be neutral. The screening of electrostatic interaction under these conditions strongly reduces the interaction between the uncompensated charge of the polymer, and it is clear that the repulsions will not be strong enough to create a rod-like state. In these conditions an expanded coil state is more likely. In terms of the effective charge of the polymer we find that two scenarios are in principle possible. If the chemical potential of the counterions in the solvent is low enough while the screening length is reduced, the expanded state is overcharged since there is an effective decrease of energy by increasing the condensation to the chain, as recently discussed by Nguyen et al. . On the other hand, when the screening length is small, it is possible that the interactions between coions and counterions in the solvent become important, and it is much harder for the chain to acquire the counterions. This is, the chemical potential for the extraction of counterions from the solvent is large and therefore the effective charge at the redissolution point has the same sign as the bare charge. We show in the next section that though there is a possibility for overcharged polyelectrolytes at the redissolution, in the sense that this state can be preferred to a collapsed one, the salt properties required for this situation ($`\mu `$, $`\kappa `$) are unlikely to be realized in practice for water soluble multivalent ions.
## III Calculation of free energies
We assumed that the polymers are monodisperse with monomer number $`N`$. Each monomer is charged positively, and is monovalent. The original counterions of the polymer also have valence $`1`$. The polymer density is $`\varphi `$. We measure all distances with respect to the monomer size $`b`$ (so that $`b=1`$ in our units). It is necessary for more precise calculations to also give more detailed geometrical information such as the aspect ratio of the monomers and the respective size of all the ions. Here we only need to further consider, besides the size of the monomer, the radius of the multivalent counterions $`a`$. The monovalent salt concentration $`s`$ is set equal to zero, but there are monovalent counterions originally dissociated from the polymer with density $`\varphi `$, and the multivalent salt, with valences $`z:1`$ has concentration $`m`$. All energies are measured in units of $`k_BT`$, (the Boltzmann constant times the temperature), and consider only energies per monomer. We form the dimensionless Bjerrum number $`B`$, equivalent to the Manning parameter, as
$$B=\frac{e^2}{4\epsilon _0\pi bk_BT},$$
(1)
where $`\epsilon _0`$ is the permitivity of the solvent (water). The typical value of this number at room temperature in strongly charged polyelectrolytes is larger than $`1`$ and as high as $`4.2`$ for double stranded DNA. We assume here for simplicity that $`B`$ is a constant, independent of the number of locally associated monomers and ions. Given that in both states water is always present, the corrections in the values of $`ϵ`$ are expected to be minor .
In both expanded and collapsed states we describe the system by a particular geometry and a given fraction $`x`$ of condensed charge. The free energy is separated into a contribution from the solvent and free ions, and another part from the interaction of the monomers and charges condensed to them. The fraction condensed is obtained self-consistently from minimization over all possible condensed fractions. The form of the entropic term is the same in both cases, but it is easier to consider the electrostatic contributions separately.
### A Free counterions contributions
It is convenient to separate the region occupied by the polymer from the rest of the solvent. To calculate the free energy of the solvent region, we first calculate the entropic contribution. We assume that the charge condensed to the chain $`xN`$ is composed by $`x_sN`$ monovalent ions and $`x_mN/z`$ multivalent ions. Thus there are $`(s/\varphi )+1x_s`$ free monovalent counterions per monomer, and thus the energy per monomer in $`k_BT`$ units is
$$F_{s1}=(\frac{m}{\varphi }\frac{x_m}{z})\mathrm{ln}(mx_m\frac{\varphi }{z})+(\frac{s}{\varphi }+1x_s)\mathrm{ln}(s+(1x_s)\varphi ).$$
(2)
In this expression we omitted the contribution of the ions with the same charge of the monomers, since it is assumed that they do not condense, and thus have a constant contribution. We have also omitted the contribution of the solvent, since the concentration of the salts is small compared to the solvent. It is possible to add in this expression an interaction term between the solvent and the ions, but for simplicity we neglect these corrections here; its effect would be only to shift the position of the L3 line in the phase diagram. It is important to consider, however, the effective free energy of interaction between all free ions. A simple way to do this is to use the modified Debye-Huckel form:
$$F_{s2}=\frac{1}{4\pi \varphi }\left(ln(1+\kappa a)\kappa a+\frac{1}{2}(\kappa a)^2\right).$$
(3)
In this expression $`a`$ is the hard-core radius of the salt ions and $`\kappa `$ is the inverse screening length produced by the salt:
$$\kappa ^2=4\pi B(z^2(mx_m\varphi /z)+zm+(2s+(1x_s)\varphi )).$$
(4)
In the very dilute polymer limit, we can simplify these expressions, and use instead
$$\kappa ^2=4\pi z(z+1)Bm,$$
(5)
and expand the total free energy of the solvent $`F_s=F_{s1}+F_{s2}`$ in a series in the overcharge fraction $`y=x1`$. The constant term is irrelevant and we retain only the linear term:
$$F_s=(_yF_s)y=\mu y,$$
(6)
where $`\mu `$ is now the effective chemical potential for the extraction of ions from the solvent into the condensed region. Using the same high salt limit as before we obtain:
$`\mu ={\displaystyle \frac{1}{z}}\mathrm{ln}m+{\displaystyle \frac{1}{8\pi (z+1)m}}{\displaystyle \frac{\kappa ^3}{1+\kappa a}}.`$
### B Expanded state energy
In our previous work , we considered only the energy of rod-like conformations. Since our interest is the regime with large amounts of salt, the screening strongly reduces the rigidity of the chain due to excess charge. It is then adequate to consider a Gaussian conformation (self-avoiding walk conformations give very similar results, as it can be seen from the functional form of the free energy computed below). The electrostatic contribution to the free energy is calculated as a function of the condensed charge, and we set the neutral state $`x=1`$ as reference. The free energy is expressed as the sum of three terms. A constant term reflects the energy of the neutral state, a term linear in the excess charge $`y=(x1)`$ appears due to addition or extraction of a charge into or from the neutral state. Finally, the interaction between segments of the chain with excess charge gives rise to a term quadratic in $`y`$. We write this as:
$$F_e^e=g_0^e+g_1^ey+g_2^ey^2.$$
(7)
A coarse-grained model clearly gives rise to the last term, and the first two can be understood, in that context, as regularizations that take into account the finite size of the particles involved. In the limit of large amounts of salt-added, most of the condensed charge comes from multivalent ions, and thus, we will take $`x=x_m`$, and $`x_s=0`$. The terms are graphically represented in Fig.2.
For the calculation of the energy of the neutral state we basically apply the Wigner principle: in a dense ionic system, the free energy can be approximated by the contribution of nearest neighbors that effectively cancel the charge of the particle. In the case of the collapsed state, below, it is more suitable to consider an approximation based on the calculation related to a full infinite lattice.
The basic neutral cluster of the linear polyelectrolyte consists of one multivalent counterion, and $`z`$ monomers. Since all particles in the cluster (with charges $`q_I`$, both monomers at distances $`r_{ij}`$ from each other, all roughly equal to the monomer size $`r=b=1`$, the electrostatic energy of the cluster is
$$E=\underset{ij}{}B\frac{q_iq_j}{r_{ij}}=B(z^2+z(z1)/2)$$
(8)
Dividing by $`z`$ monomers, we obtain:
$$g_0^e=B(z+1)/2.$$
(9)
When an extra multivalent ion is added to this cluster, the counterion can have roughly the same interaction with the chain monomers, but can locate itself away of the first counterion, say, by putting a monomer in between them so that $`r2`$. Thus we obtain a simple estimate for the energy of the new charged cluster as:
$$E=B(2z^2+z(z1)/2+z^2/2),$$
(10)
and
$$g_1^e=Bz/2.$$
(11)
Finally, the two-body term can be calculated assuming a uniform distribution of charge along the Gaussian chain. A Gaussian chain has a monomer distribution given, approximately, by $`\rho (r)=1/(2\pi r)`$, inside the volume $`V`$ limited by the radius of gyration $`R`$; this is, for $`r<R=N^{1/2}`$. A charge at the center of this distributions interacts with other charges via the screened potential $`Be^{\kappa r}/r`$. The energy for one such charge is given by the average over the volume $`V`$ of this potential:
$$E=V=_V𝑑𝐫\rho (r)B\frac{e^{\kappa r}}{r}=\frac{B}{\kappa }(1e^{\kappa R}),$$
(12)
and by our previous definition,
$$g_2^e=\frac{1}{2}V.$$
(13)
Notice that while this calculation depends on the assumed Gaussian conformation, a simple estimate can also be done by using a general expansion coefficient for the radius of gyration $`RN^\nu `$, ($`\nu =1`$ for rod-like, $`\nu =2/3`$ for self-avoiding walks, and $`\nu =1/2`$ for random walks). In such case, we obtain a scaling for the quadratic precoefficient of the form $`g_2^e\kappa ^{(11/\nu )}`$. While these different models give rise to different estimates, the transition we are interested in occurs when the screening is of the order of the size of the monomers, and so, in our units, $`\kappa 1`$, and the estimates at that point are not so different.
### C Collapsed state energy
We repeat a similar approach for the calculation of the free energy of the collapsed state which is expressed as:
$$F_e^c=g_0^c+g_1^cy+g_2^cy^2.$$
(14)
Inside the collapsed polymer, we can imagine the environment of a charge resembling that of a ionic crystal. If the structure of the crystal is known, the energy per atom in that crystal can be given in terms of the Madelung constant ,
$$E=Me^2/r=MB.$$
(15)
The Madelung constant for an ionic monovalent salt is $`M=1.747`$. There are two important differences in our case. We would like to consider multivalent ions, which reduce the contribution from same-charge interactions, and on the other hand the connectivity of the monomers forces them to be next-neighbors (not the case in ionic crystals), and thus increases the same-charge interactions. A way to obtain a rough estimate for $`M`$ is to consider a simple geometry in which both monomers and counterions lie in columns which in turn form a square lattice when a cross section is taken. The result does not change dramatically when other possible geometries are considered. As a rule of thumb, we can use an effective Madelung constant that is larger than unity, but not much bigger. In the geometry proposed, the number of next neighbors of counterion is $`4z`$. Summing the interaction with $`4z`$ monomers and 4 equally charged next-nearest counterions ($`r=2b`$) we obtain the following estimate of the energy per cluster of $`1`$ counterion and $`z`$ monomers:
$$E=B(4z^2+z(z1)+\frac{1}{2}z(4zz)+\frac{1}{2}\frac{4z^2}{2})$$
(16)
where the interactions with charges outside the cluster are weighted by $`1/2`$ to avoid double counting. Dividing again by $`z`$, we obtain the coefficient of the energy per monomer:
$$g_0^c=B(1+\frac{z}{2}).$$
(17)
In this approximation we have underestimated the Madelung constant for $`z=1`$, for which we obtain $`1.5`$ instead of the exact result $`1.747`$, but it is a reasonable approximation and shows that the energy of this neutral collapsed state is smaller than the neutral expanded state.
In both the collapsed and expanded states, the addition and extraction of charges are not completely symmetric processes. It is energetically unfavorable to add or extract charges from the collapse bulk state, as shown below in the calculation of $`g_2^c`$. Both processes are however better thought of as occurring at the surface of the collapsed region. We have also argued previously , that while the mobility of charges in this type of conformations is not so large, the analogy with a conductor is still useful: all charge excess resides in the surface. While the local conformation at the surface of the collapsed state is not equal to the expanded state, it is simpler to borrow the result for the coefficient related to the addition of one charge from that case. In both situations the new charge is put in close contact with a neutral cluster, but sits slightly closer to the opposite charges. Thus,
$$g_1^c=Bz/2.$$
(18)
Since we assume that the extra charge is sitting at the surface, the interaction energy can be calculated assuming a uniform distribution in a shell surrounding the collapsed polymer. We assume this shell to be spherical with radius $`R`$, and obtain the energy per added charge as
$$g_2^c=E=\frac{1}{2}B\frac{1}{R}=\frac{1}{2}\frac{B}{(3N(1+1/z)/4\pi )^{1/3}}$$
(19)
Here we have considered that the $`N+N/z`$ charged particles are closely packed and that they occupy each a volume of $`b^3`$.
## IV Comparison between states and Phase diagram
After the minimization of the free energies of each state with respect to their respective condensed charge variables, we can compare each other to find the state of the chain for given concentration conditions. This minimization can be done directly from the expressions already given but it is easier to present simpler formulas if we use a set of approximations described below. Numerical results presented later reflect these approximations, but we have checked that numerical minimization of the full expressions do not lead to important qualitative and quantitative changes in the results.
The free energy of each of the states, for given concentrations of monomers and salt correspond to the minima of:
$`F^e`$ $`=`$ $`F_s+F_e^e,`$ (20)
$`F^c`$ $`=`$ $`F_s+F_e^c`$ (21)
for the expanded and collapsed states respectively. We minimize now with respect to the overcharge fraction $`y`$.Taking the first derivatives with respect to this quantity, we obtain the (independent) conditions
$`_yF^e`$ $`=`$ $`\mu +g_1^e+2g_2^ey^e=0,`$ (22)
$`_yF^c`$ $`=`$ $`\mu +g_1^c+2g_2^cy^c=0.`$ (23)
The solutions for the fractions condensed are:
$`y^e`$ $`=`$ $`(\mu +g_1^e)/2g_2^e,`$ (24)
$`y^c`$ $`=`$ $`(\mu +g_1^c)/2g_2^c.`$ (25)
Replacement of this values in the general expressions for the free energy gives the minimized free energies for each state:
$`F^e`$ $`=`$ $`F^e(y^e)=\mu +g_0^e{\displaystyle \frac{(g_1^e+\mu )^2}{4g_2^e}},`$ (26)
$`F^c`$ $`=`$ $`F^c(y^c)=\mu +g_0^c{\displaystyle \frac{(g_1^c+\mu )^2}{4g_2^c}}.`$ (27)
For the moment it is convenient to maintain the expressions of the free energies as functions of the chemical potential instead of the salt concentration. Of the parameters used to describe the different states, only $`g_2^e`$ has an explicit dependence on the concentration (through the inverse screening length $`\kappa `$)l. Also, it is important to note that the coefficient $`g_2^c`$ is extremely large, and for the values of the chemical potential that we need to explore, it turns out that the denominator in the expression for the overcharge is much larger than the numerator, and the overcharging fraction is very close to $`0`$:
$$y^c0.$$
(28)
This result is clear since we have found that there is a very strong penalty for any overcharging (or undercharging) in the collapsed state. The effective charge in a general structure is affected by the environment (namely, by $`\mu `$) but the collapsed structure is only very slightly susceptible to it. Furthermore, the interaction between the charges in the bulk of the collapsed structure is not subject to screening by the free particles outside and its contribution remains constant with changing chemical potential. Thus, the system splits neatly into collapsed neutral structures and a homogeneous mixture of solvent and non-condensed counterions. It is then not surprising that this state adequately represents what in experiment is close to be two separate phases.
Consider now the minimum for the expanded state. The coefficient for charge-charge interactions $`g_2^e`$ varies very strongly with the salt concentration. It remains large when the screening length is large, since in those conditions all uncompensated charges interact strongly. On the other hand at large concentrations of salt the screening drastically reduces the interaction between segments of the chain.
To obtain an expression for the transition point we subtract the free energies, and define the energy difference
$$\mathrm{\Delta }F=F^eF^c.$$
(29)
The system is in the region $`B`$, when $`\mathrm{\Delta }F>0`$, in region $`C`$ when $`\mathrm{\Delta }F<0`$, and the transition curve is defined by $`\mathrm{\Delta }F=0`$. Using the fact, discussed above, that the electrostatic energy of the collapsed state remains essentially constant, we can write the equation $`\mathrm{\Delta }F=0`$ as
$$g_0^e\frac{(g_1^e+\mu )^2}{4g_2^e}=+g_0^c.$$
(30)
One useful way to use and interpret this equation is to consider it as defining a boundary in the $`\mu \kappa ^2`$ plane that corresponds to the properties of the counterions. Within this diagram, a particular form of the relation between concentration and screening for a counterion defines a path in the diagram parametrized by the salt concentration $`m`$. Examples of this are presented in Fig.3. for certain polymer properties values, and different types of $`\mu \kappa ^2`$ relations.
At this point is very easy to explain the observed near independence of the transition on the molecular weight of the chains. Only the coefficients $`g_2`$ carry information on the molecular weight, but as we have seen $`g_2^c`$ drops out of the final equations, and in $`g_2^e`$ the molecular weight provides only an exponentially small correction to the dominant term given by a function of the screening length.
As shown in Fig.3, and as it can be seen from the fact that the equations defining the transition are quadratic in $`\mu `$, the expanded state splits into two regions, each with non-zero effective charge, but with opposite signs on each. This diagram is obtained using the approximations presented above, and a more precise determination of the lines of transition will involve consideration of the monovalent counterions, the finite amount of monomers, and the rod-like state obtained at very low screening. When these considerations are taken into account, the trajectory described by the added salt in the $`\mu \kappa ^2`$ plane will start within the expanded region, cross into the collapsed region (giving rise to the L1-L2 line), and continue there until crossing again, redissolving, into one of the branches with expanded states. The starting point of the trajectory, with zero multivalent salt added $`m=0`$, is located in the branch with natural charge, at a finite $`\kappa `$ value (due to the monovalent counterions), and at $`\mu =\mathrm{}`$. Figure 3 illustrates only the redissolution.
In Fig.3 we show three different curves defining the properties of the counterions: a pure Debye-Huckel case where the size of the counterions are neglected, a modified Debye-Huckel curve with $`a=0.2`$ (that is $`a=0.2b`$), and a third curve that neglects the interactions between the counterions when they are free considers only the entropic term. The reversal of direction in the $`\mu `$ axis of the parametric curve for the Debye-Huckel cases occurs because of the onset of important favorable interactions between the free charges. Thus it is clear that the trajectory can end in the naturally charged branch of the expanded region since there are more free counterions there, and at the same time the repulsive intra-chain interactions are small. The overcharged branch is only accessible if the chemical potential is further reduced at the same time that the screening increases. For most cases with sensible sizes for the counterions the reversal in direction of the parametrized curve comes always before the redissolution, and thus the redissolution is into the naturally charged state.
It is interesting to note that at high enough salt concentrations and high valences, the salt might re-associate on its own. In this case the trajectory in the $`\mu \kappa ^2`$ plane has a different shape than that provided by the Debye-Huckel equation. Indeed, while the results from the Debye-Huckel model are know to provide good approximations for the free energy of the salt in solution they do not properly describe the individual dissociated species. The upturn of the chemical potential essentially indicates a phase transition for the salt/solvent system. In reality there is no phase transition but there is a finite a fraction of ion pairs formed in the solvent (even if they do not reassociate chemically). We have not presented here the curve in the $`\mu \kappa ^2`$ diagram when the association of multivalent ions in the solution is present, as is modeled, for example, in the classical Bjerrum model , or more recently by Fisher, et al. This association gives large corrections to the Debye-Huckel law (for the free energy of the individual species of ions) and in our case provides a $`\mu \kappa `$ curve that monotonically approaches a horizontal line (a limiting $`\mu `$ value), with increasing inverse screening length. These effects are particularly important for large valence salts as in our case and we will discuss them in more detail in a future publication .
## V Comparison with experimental information
With the previously obtained solutions we now contrast our results with the currently known experimental situation. We consider the shape of the phase diagram, the precise location of the collapse transition and the predicted structure of the polymers in both of the phases.
It should be clear that we have already recovered the shape of the $`B`$-$`C`$ transition. Indeed, the (approximate) equation that define the transition does not have any dependence on the monomer concentration, and therefore produces a flat line in the $`\varphi m`$, or $`\varphi \mu `$ diagrams. The transition is defined by the choice of single chains to be in one of two states that must be in equilibrium with the surrounding environment. This environment, represented here by means of an effective chemical potential and screening length, depends strongly on the amount of multivalent salt and only very slightly on the polymer amount. The more precise expressions for the chemical potential recovers a small dependence in the polymer concentration.
There are two important ways in which important changes on the qualitative form of the diagram can arise.The addition of large amounts of monovalent salt introduces new features in the phase diagram, not show in the scheme of Fig. 1. In the introduction we only mentioned the effects on the $`AB`$ transition, but further addition of monovalent salt, creates a different environment of the chains, and changes the transition points and the shape of the transition curve. Secondly, it is clear that as we continue increasing the amount of polymer in the solution, the interactions between the chains start to become important and the transition (if any) is a much more complicated phenomena. From the point of view of our theory it is clear that the increase in concentration of monovalent salt and monomers bring about a breakdown of the assumption that the chemical potential for the overcharging is dependent only on the amount of multivalent salt. The effective environment of the chain becomes more complicated and the chemical potential will now contain important terms coming form the concentrations of all species and from many-body interactions.
Let us consider two concrete numerical examples. Experiments performed on polystyrene sulfonate at room temperature with $`N=410^3`$, Bjerrum number of $`B=2.87`$ ($`b=0.5`$ nm) and an effective ion radius $`a/b=0.2`$, show a transition line at a concentration of $`m=0.2M`$ of $`LaCl_3`$. With these parameters as input, our equation predicts redissolution at a salt concentration of $`m=0.1M`$. Secondly, experiments on double stranded DNA with very large number of base pairs (we take $`N=10^4)`$ and $`B=4.2`$ ($`b=.17`$ nm)in spermine ($`z=4`$), $`a=1`$. obtained redissolution at $`m=.1M`$, while our equations predict a transition at $`m=0.04M`$. This compares favorably with the experiments, especially if we take into account the rough approximations done in the evaluation of the energies of the system.
According to the theory presented, in phase B, the polymer chains are almost neutral with a collapsed conformation. In phase C, the chains are expanded and charged. This suggests, besides other techniques, to confirm these predictions by means of scattering experiments that test the structure of the chains, and of osmotic pressure measurements that can determine the amount of free counterions in each state (as it has been done for semi-dilute solutions).
## VI Conclusions
The redissolution transition observed in multivalent induced precipitated polyelectrolytes with further addition of multivalent particles or salt was predicted extending a previously developed two state model to deal with large salt concentrations. The electrostatic energies of both, collapsed and expanded-coiled, states were computed considering the finite size of the ions and monomers condensed along the chains, and using a mean field approach for the non-condensed ions. We neglected the non-ionic short range interactions of the ions with the solvent, and assume zero monovalent salt.
We found that the redissolution is determined by the properties of the ionic solution. We calculated the effective charge of the chains for the different thermodynamic states of the chain.
### Acknowledgments
We thank E. Raspaud and F. Livolant for useful discussions and for performing electrophoresis experiments , M. Campos for computing electrostatic energies of dense finite size systems of charges, and P. A. Bernikowicz for computing the free energy of aggregates containing $`p=1,2,\mathrm{}`$ chains. This work was sponsored by the National Science Foundation, grants DMR9807601 and DMR9632472.
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# Finite Temperature Excitations of an Inhomogeneous Trapped Bose Gas With Feshbach Resonances*footnote **footnote *Mailing address: Institute of Theoretical Physics, Academia Sinica, Peking 100080, China
## Abstract
We derive and discuss the temperature dependance of the condensate and noncondensate density profile of a Bose-Einstein condensate gas with Feshbach resonance in a parabolic trap. These quantities are calculated self-consistently using the generalized Hartree-Fock-Bogoliubov (HFB) equations within the Bogoliubov approximation. At zero temperature, the HFB equation can be solved by means of a variation method that give the low excitation spectrum. Moreover, within the two-body collision theory, we estimate the relationship between the atom number and the external magnetic field $`B`$, it is in good agreement with the data in recent experiments.
Bose-Einstein condensates\[1-4\] of atomic gases offer new opportunities for studying quantum-degenerate fluids. Almost all the essential properties of Bose-Einstein condensate systems such as the formation and shape of the condensate and noncondensate, the nature of its collective excitations and statistical fluctuations, and the formation and dynamics of vortices are determined by the strength of atom-atom interactions. In contrast with the situation in traditional superfluids, the strength of inter-particle interactions in the atomic condensate can vary over a wide range of values\[5-12\]. In particular, the scattering length that characterizes the atom-atom interactions can be negative, corresponding to an effective inter-atom attraction.
Most recently, in trapped atomic Bose-Einstein condensation, Ketterle’s group reported evidences for modifying the scattering length by magnetic-field-induced Feshbach resonance. Feshbach resonance have been studied twenty one years ago at much higher energies, but the Feshbach resonance energy observed in ultracold atoms can be tuned to near zero. The theoretical studies of the ultracold atoms with Feshbach resonance show that the two-body interactions responsible for the Feshbach resonance produce an additional condensate of molecules, which differs qualitatively from the properties of a single condensate. In this letter, we discuss the temperature dependance of the hybrid atomic/molecular condensate and noncondensate density profile as well as the excitation spectrum. Especially, for an inhomogeneous system with negative scattering length, the excitation spectrum shows a upper constraint on the atom number, which, together with the conventional estimation of the condensate atom number, gives the relationship between the atom number and the external magnetic field $`B`$. The theoretical results are in good agreement with the experiment.
The binary atom Feshbach resonances studied by varying a strong external magnetic field in an alkali-atom trap are hyperfine-induced spin-flip processes that bring the colliding atoms to a bound molecular state of different electron spin. This process can be described by the Hamiltonian
$$H_{FR}=\alpha d^3r\mathrm{\Psi }_m^+(r)\mathrm{\Psi }_a(r)\mathrm{\Psi }_a(r)+h.c.,$$
(1)
where $`\mathrm{\Psi }_m(r),\mathrm{\Psi }_m^+(r)`$ ($`\mathrm{\Psi }_a(r),\mathrm{\Psi }_a^+(r)`$) are the annihilation and creation field operators of the molecules(atoms), $`\alpha `$ stands for the transition matrix element proportional to the overlap of the molecular continuum and bound state wave functions. Usually, the transition matrix element depends on the magnetic field as $`\alpha \sqrt{\lambda _a\mathrm{\Delta }^2/2|BB_0|}`$, where $`B_0`$ is the resonant magnetic field and $`\mathrm{\Delta }`$ characterizes the width of the resonance as a function of $`B`$. The Hamiltonian $`H_{FR}`$ together with the atomic Hamiltonian
$$H_a=d^3r\mathrm{\Psi }_a^+(r)[\frac{^2}{2M}+V_a(r)\mu _a]\mathrm{\Psi }_a(r)+\frac{\lambda _a}{2}d^3r\mathrm{\Psi }_a^+(r)\mathrm{\Psi }_a^+(r)\mathrm{\Psi }_a(r)\mathrm{\Psi }_a(r),$$
(2)
the molecular Hamiltonian
$$H_m=d^3r\mathrm{\Psi }_m^+(r)[\frac{^2}{4M}+V_m(r)\mu _m+ϵ]\mathrm{\Psi }_m(r)+\frac{\lambda _m}{2}d^3r\mathrm{\Psi }_m^+(r)\mathrm{\Psi }_m^+(r)\mathrm{\Psi }_m(r)\mathrm{\Psi }_m(r),$$
(3)
and the atom-molecule interaction Hamiltonian
$$H_{am}=\lambda d^3r\mathrm{\Psi }_a^+(r)\mathrm{\Psi }_m^+(r)\mathrm{\Psi }_m(r)\mathrm{\Psi }_a(r)$$
(4)
forms a total Hamiltonian $`H=H_a+H_m+H_{am}+H_{FR}`$, which governs the dynamics of the system under investigation. Here, $`V_{a(m)}(r)`$ represents the trapped potential for atom (molecule), $`\mu _{a(m)}`$ is the chemical potential of the atoms(molecules), $`\lambda _{a(m)}=\frac{4\pi a_{a(m)}}{(2)M}`$, with $`M`$ being the atomic mass and $`a_{a(m)}`$ the s-wave scattering length of the atom-atom interaction, $`\lambda `$ denotes the coupling constant of the atom-molecule interaction, and $`ϵ`$ is the energy of the intermediate molecular state relative to the continuum of the incident atoms.
Separating out the condensate part in the usual fashion (Bogoliubov approximation) i.e.
$$\mathrm{\Psi }_{a(m)}(r)=\varphi _{a(m)}(r)+\widehat{\phi }_{a(m)}(r),$$
(5)
where $`\varphi _{a(m)}(r)=\mathrm{\Psi }_{a(m)}(r)`$ plays the role of a spatially varying macroscopic Bose field of the atoms(molecules). The possibility that the resonantly formed quasibound atom pairs form a molecular condensate was previously suggested by Timmermans et al.. Using a Raman photonassociation process, the quasibound pairs start to be formed from the atomic condensate and form a molecular condensate. Here, we assume that there are a large number of atoms and molecules in the condensate. It is easy to show that the operators $`\widehat{\phi }_{a(m)}(r)`$ and $`\widehat{\phi }_{a(m)}^+(r)`$ obey the Bose commutation relations
$$[\widehat{\phi }_{a(m)}(r),\widehat{\phi }_{a(m)}^+(r^{})]=\delta (rr^{}).$$
(6)
In terms of $`\widehat{\phi }_{a(m)}`$ and $`\varphi _{a(m)}`$, the Hamiltonian can be expanded as
$`H`$ $`=`$ $`H_0+H^{^{}},`$ (7)
$`H_0`$ $`=`$ $`{\displaystyle }d^3r\{\varphi _a^{}({\displaystyle \frac{^2}{2M}}\mu _a+V_a(r))\varphi _a+{\displaystyle \frac{\lambda _a}{2}}\varphi _a^{}\varphi _a^{}\varphi _a\varphi _a`$ (8)
$`+`$ $`\varphi _m^{}({\displaystyle \frac{^2}{4M}}\mu _m+V_m(r)+ϵ)\varphi _m+{\displaystyle \frac{\lambda _m}{2}}\varphi _m^{}\varphi _m^{}\varphi _m\varphi _m`$ (9)
$`+`$ $`\lambda \varphi _a^{}\varphi _a\varphi _m^{}\varphi _m+\alpha \varphi _m^{}\varphi _a\varphi _a+\alpha \varphi _m\varphi _a^{}\varphi _a^{}\}`$ (10)
$`H^{^{}}`$ $`=`$ $`{\displaystyle }d^3r\{\widehat{\phi }_a^+({\displaystyle \frac{^2}{2M}}\mu _a+V_a(r))\widehat{\phi }_a+2\lambda _a\widehat{\phi }_a^+\widehat{\phi }_a\varphi _a^{}\varphi _a+{\displaystyle \frac{\lambda _a}{2}}(\widehat{\phi }_a^+\widehat{\phi }_a^+\varphi _a\varphi _a+\widehat{\phi }_a\widehat{\phi }_a\varphi _a^{}\varphi _a^{})`$ (11)
$`+`$ $`\widehat{\phi }_m^+({\displaystyle \frac{^2}{4M}}+ϵ\mu _m+V_m(r))\widehat{\phi }_m+2\lambda _m\widehat{\phi }_m^+\widehat{\phi }_m\varphi _m^{}\varphi _m+{\displaystyle \frac{\lambda _m}{2}}(\widehat{\phi }_m^+\widehat{\phi }_m^+\varphi _m\varphi _m+\widehat{\phi }_m\widehat{\phi }_m\varphi _m^{}\varphi _m^{})`$ (12)
$`+`$ $`\lambda \widehat{\phi }_a^+\widehat{\phi }_a\varphi _m\varphi _m^{}+\lambda \widehat{\phi }_m^+\widehat{\phi }_m\varphi _a^{}\varphi _a+\alpha (\widehat{\phi }_a\widehat{\phi }_a\varphi _m^{}+\widehat{\phi }_a^+\widehat{\phi }_a^+\varphi _m)\}.`$ (13)
In derivation of eq.(7,8), the following coupling equations are used,
$`\{{\displaystyle \frac{^2}{2M}}+\lambda _a|\varphi _a|^2+V_a(r)+\lambda |\varphi _m|^2\}\varphi _a+2\alpha \varphi _m\varphi _a^{}=\mu _a\varphi _a,`$ (15)
$`\{{\displaystyle \frac{^2}{4M}}+\lambda _m|\varphi _m|^2+V_m(r)+ϵ+\lambda |\varphi _a|^2\}\varphi _m+\alpha \varphi _a\varphi _a=\mu _m\varphi _m.`$ (16)
This coupling equations may be yielded by the expectation value of the Heisenberg equations
$$i\mathrm{}\dot{\mathrm{\Psi }}_a=[\mathrm{\Psi }_a,H],i\mathrm{}\dot{\mathrm{\Psi }}_m=[\mathrm{\Psi }_m,H],$$
(17)
and replacing the time derivatives by the chemical potentials
$$i\mathrm{}\dot{\varphi }_a\mu _a\varphi _a,i\mathrm{}\dot{\varphi }_m\mu _m\varphi _m.$$
The chemical potential of the molecules is twice the chemical potential of the atoms, in accordance with the condition for chemical equilibrium. The $`\alpha `$terms that couple the equations describe tunneling of pairs of atoms between $`\varphi _m`$ and $`\varphi _a`$ fields, it leads to the form of a second condensate—molecular condensate in an atomic Bose-Einstein condensate\[15-17\]. Using the coupling equations(9), Timmermans et al. investigate the behaviors of the hybrid atomic/molecular condensates near- and off- resonance. The Hamiltonian (8) can be diagonalized by using the Bogoliubov transformation
$`\widehat{\phi }_a(r)={\displaystyle \underset{j}{}}[u_j(r)\alpha _jv_j^{}(r)\alpha _j^+],`$ (18)
$`\widehat{\phi }_a^+(r)={\displaystyle \underset{j}{}}[u_j^{}(r)\alpha _j^+v_j(r)\alpha _j],`$ (19)
$`\widehat{\phi }_m(r)={\displaystyle \underset{j}{}}[x_j(r)\beta _jy_j^{}(r)\beta _j^+],`$ (20)
$`\widehat{\phi }_m^+(r)={\displaystyle \underset{j}{}}[x_j^{}(r)\beta _j^+y_j(r)\beta _j],`$ (21)
where the qusiparticle operators $`\alpha _j`$,$`\alpha _j^+`$,$`\beta _j`$, $`\beta _j^+`$ obey boson commutation relations
$$[\alpha _i,\alpha _j^+]=\delta _{ij},[\alpha _i,\alpha _j]=[\alpha _i^+,\alpha _j^+]=0,$$
$$[\beta _i,\beta _j^+]=\delta _{ij},[\beta _i,\beta _j]=[\beta _i^+,\beta _j^+]=0,$$
$$[\alpha _i^+,\beta _j^+]=[\alpha _i,\beta _j]=[\alpha _i,\beta _j^+]=0,$$
and $`u_j(r),v_j(r),x_j(r),y_j(r)`$ are $`c`$-number functions. Substituting eq.(11) into eq.(8), one yields
$$H^{^{}}=\underset{j}{}E_j\alpha _j^+\alpha _j+\underset{i}{}e_i\beta _i^+\beta _i\underset{j}{}E_jd^3r|v_j(r)|^2\underset{i}{}e_id^3r|y_i(r)|^2$$
(22)
with
$`({\displaystyle \frac{^2}{2M}}+\lambda |\varphi _m|^2+2\lambda _a|\varphi _a|^2\mu _a+V_a(r))u_j(\lambda _a\varphi _a^{}\varphi _a^{}+2\alpha \varphi _m^{})v_j`$ $`=`$ $`E_ju_j,`$ (23)
$`({\displaystyle \frac{^2}{2M}}+\lambda |\varphi _m|^2+2\lambda _a|\varphi _a|^2\mu _a+V_a(r))v_j(\lambda _a\varphi _a^{}\varphi _a^{}+2\alpha \varphi _m^{})u_j`$ $`=`$ $`E_jv_j,`$ (24)
$`({\displaystyle \frac{^2}{4M}}+ϵ+\lambda |\varphi _a|^2+2\lambda _m|\varphi _m|^2\mu _m+V_m(r))x_j\lambda _m\varphi _m^{}\varphi _m^{}y_j=e_jx_j,`$ (25)
$`({\displaystyle \frac{^2}{4M}}+ϵ+\lambda |\varphi _a|^2+2\lambda _m|\varphi _m|^2\mu _m+V_m(r))y_j\lambda _m\varphi _m^{}\varphi _m^{}x_j=e_jy_j.`$ (26)
In order to study the temperature dependence of the excitation spectrum as well as the spatial distribution of the hybrid atom/molecular condensate and noncondensate, we need to solve the coupled mean-field Bogoliubov equations(13), and the condensate equation(9) self-consistently. The calculations procedure can be summarized for an arbitrary confining potential as follows: First of all, we solve eq.(9) self-consistently, once $`\varphi _a`$ and $`\varphi _m`$ are known, the solution of $`u_j,v_j,x_j`$ and $`y_j`$ can be generated. To illustrate this procedure, we present its first step of calculations analytically. The trapped potential considered here is taken to be an isotropic harmonic potential $`V_{a(m)}(r)=\frac{1}{2}M\omega _{a(m)}^2r^2`$, for which $`\varphi _a`$ and $`\varphi _m`$ are spherically symmetric functions,
$$\varphi _{a(m)}(r)=R_{00}(r)Y_{00}(\theta ,\psi ),$$
(27)
with
$$R_{00}(r)=\alpha ^{3/2}\sqrt{\frac{4}{\pi }}exp[\frac{1}{2}\alpha ^2r^2],Y_{00}(\theta ,\psi )=\frac{1}{\sqrt{4\pi }},\alpha =((2)M\omega )^{1/2}.$$
Rather than solving the coupled equations (13) directly, we introduce a new method based on the auxiliary functions
$`u_j`$ $`=`$ $`A_jr|j_a,v_j=B_jr|j_a,`$ (28)
$`x_j`$ $`=`$ $`C_jr|j_m,y_j=D_jr|j_m,`$ (29)
where $`|j_{a(m)}`$ is defined by
$$[\frac{^2}{2(4)M}+V_{a(m)}(r)]|j_{a(m)}=\mathrm{}\omega _{a(m)}(j+\frac{1}{2})|j_{a(m)}.$$
The reason for such selection is that the level shifts caused by atom-atom interactions weakly depend on the shape of the wave function. A combination of Eqs.(13-15) gives
$`(\mathrm{}\omega _a(j+{\displaystyle \frac{1}{2}})+\lambda |\varphi _m|^2+2\lambda _a|\varphi _a|^2\mu _a)A_j(\lambda _a\varphi _a^{}\varphi _a^{}+2\alpha \varphi _m^{})B_j`$ $`=`$ $`E_jA_j,`$ (31)
$`(\mathrm{}\omega _a(j+{\displaystyle \frac{1}{2}})+\lambda |\varphi _m|^2+2\lambda _a|\varphi _a|^2\mu _a)B_j(\lambda _a\varphi _a^{}\varphi _a^{}+2\alpha \varphi _m^{})A_j`$ $`=`$ $`E_jB_j,`$ (32)
$`(\mathrm{}\omega _m(j+{\displaystyle \frac{1}{2}})+ϵ+\lambda |\varphi _a|^2+2\lambda _m|\varphi _m|^2\mu _m)C_j\lambda _m\varphi _m^{}\varphi _m^{}D_j=e_jC_j,`$ (33)
$`(\mathrm{}\omega _m(j+{\displaystyle \frac{1}{2}})+ϵ+\lambda |\varphi _a|^2+2\lambda _m|\varphi _m|^2\mu _m)D_j\lambda _m\varphi _m^{}\varphi _m^{}C_j=e_jD_j,`$ (34)
the eigenfunctions and the corresponding eigenvalues are given by
$`B_j^\pm `$ $`=`$ $`[{\displaystyle \frac{1}{f^\pm (r,j)1}}]^{\frac{1}{2}},A_j^\pm (r)=f^\pm (r,j)B_j^\pm ,`$ (35)
$`D_j^\pm `$ $`=`$ $`[{\displaystyle \frac{1}{g^\pm (r,j)1}}]^{\frac{1}{2}},B_j^\pm (r)=g^\pm (r,j)D_j^\pm ,`$ (36)
and
$`E_j^\pm (r)`$ $`=`$ $`\pm \{(\lambda _a\varphi _a^{}\varphi _a^{}+2\alpha \varphi _m^{})[\mathrm{}\omega _a(j+{\displaystyle \frac{1}{2}})\mu _a+\lambda |\varphi _m|^2+2\lambda _a|\varphi _a|^2]\},`$ (37)
$`e_j^\pm (r)`$ $`=`$ $`\pm \{\lambda _m\varphi _m^{}\varphi _m^{}[\mathrm{}\omega _m(j+{\displaystyle \frac{1}{2}})\mu _m+ϵ+2\lambda _m|\varphi _m|^2]\}.`$ (38)
Here,
$$f^\pm (r,j)=\frac{\lambda _a\varphi _a^{}\varphi _a^{}+2\alpha \varphi _m^{}}{\mathrm{}\omega _a(j+\frac{1}{2})\mu _a+\lambda |\varphi _m|^2+2\lambda _a|\varphi _a|^2(E_j^\pm )^2},$$
and
$$g^\pm (r,j)=\frac{\lambda _m\varphi _m^{}\varphi _m^{}}{\mathrm{}\omega _m(j+\frac{1}{2})\mu _m+ϵ+2\lambda _m|\varphi _m|^2(e_j^\pm )^2}.$$
These explicit solutions enable us to construct the one-body density matrix
$`\rho (r,r^{^{}})`$ $`=`$ $`\rho _a(r,r^{^{}})+2\rho _m(r,r^{^{}}),`$ (39)
$`\rho _a(r,r^{^{}})`$ $`=`$ $`\varphi _a^{}(r)\varphi _a(r^{^{}})`$ (40)
$`+`$ $`{\displaystyle \underset{p=\pm ,i=1}{\overset{\mathrm{}}{}}}[u_i^p(r)u_i^p(r^{^{}})F_i^p+v_i^p(r)v_i^p(r^{^{}})(1+F_i^p)],`$ (41)
$`\rho _m(r,r^{^{}})`$ $`=`$ $`\varphi _m^{}(r)\varphi _m(r^{^{}})`$ (42)
$`+`$ $`{\displaystyle \underset{p=\pm ,i=1}{\overset{\mathrm{}}{}}}[x_i^p(r)x_i^p(r^{^{}})f_i^p+y_i^p(r)y_i^p(r^{^{}})(1+f_i^p)],`$ (43)
where $`F_i^p=\frac{1}{exp(\beta E_i^p)1}`$ and $`f_i^p=\frac{1}{exp(\beta e_i^p)1}`$ are the Bose distribution for the quasiparticle excitations with energies $`E_i^p`$ and $`e_i^p`$, respectively. Setting $`r=r^{^{}}`$, eq.(19) follows the resulting particle density.
We need to point out that eqs. (17) and (18) are results of the first step of the numerical calculations. To complete numerical calculations, we should repeat the above procedures until the eigenvalues $`E_j`$ and $`e_j`$ do not depend on position $`r`$. In what follows, we present a variation method to study the excitations at zero temperature. This method was first introduced in Ref. to study the BEC ground state in the harmonic trap of boson system, and it was generalized in Ref. to investigate the excited states in BEC. Considering eq.(13) as well as
$$[u_j(r)u_j^{}(r)v_j(r)v_j^{}(r)]𝑑r=1,$$
and
$$[x_j(r)x_j^{}(r)y_j(r)y_j^{}(r)]𝑑r=1,$$
which were derived from the Bose commutation relation (6) we arrive at
$`E_j`$ $`=`$ $`{\displaystyle u_j^{}(r)(\frac{^2}{2M}+\lambda |\varphi _m|^2+2\lambda _a|\varphi _a|^2+V_a(r))u_j𝑑r}`$ (44)
$`+`$ $`{\displaystyle v_j^{}(r)(\frac{^2}{2M}+\lambda |\varphi _m|^2+2\lambda _a|\varphi _a|^2+V_a(r))v_j𝑑r}`$ (45)
$``$ $`{\displaystyle u_j^{}(r)(\lambda _a\varphi _a^{}\varphi _a^{}+2\alpha \varphi _m^{})v_j(r)𝑑r}{\displaystyle v_j^{}(r)(\lambda _a\varphi _a^{}\varphi _a^{}+2\alpha \varphi _m^{})u_j(r)𝑑r},`$ (46)
$`e_j`$ $`=`$ $`{\displaystyle x_j^{}(r)(\frac{^2}{4M}+\epsilon +\lambda |\varphi _a|^2+2\lambda _m|\varphi _a|^2+V_m(r))x_j𝑑r}`$ (48)
$`+`$ $`{\displaystyle y_j^{}(r)(\frac{^2}{4M}+\lambda |\varphi _a|^2+\epsilon +2\lambda _m|\varphi _m|^2+V_m(r))y_j𝑑r}`$ (49)
$``$ $`{\displaystyle x_j^{}(r)\lambda _m\varphi _m^{}\varphi _m^{}y_j(r)𝑑r}{\displaystyle y_j^{}(r)\lambda _m\varphi _m^{}\varphi _m^{}x_j(r)𝑑r}.`$ (50)
For simplicity, we study only the case of the spherical harmonic trap. In this case, we may choose the trial wave functions of the excitation components $`u_j(r),v_j(r)`$,$`x_j(r)`$ and $`y_j(r)`$ in the form of the spherical harmonic oscillator wave function $`\xi _{n_r,l,m}`$ with quantum numbers $`(n_r,l,m)`$:
$`\left(\begin{array}{c}u_j(r)\\ v_j^{}(r)\end{array}\right)`$ $`=`$ $`\left(\begin{array}{c}u\\ v^{}\end{array}\right)\xi _{n_r,l,m}(\omega _{n_rlm},r),`$ (56)
$`\left(\begin{array}{c}x_j(r)\\ y_j^{}(r)\end{array}\right)`$ $`=`$ $`\left(\begin{array}{c}x\\ y^{}\end{array}\right)\xi _{n_r,l,m}(\omega _{n_rlm},r),`$ (61)
where $`\omega _{n_rlm}`$ is an adjustable scaling factor of variation. Eqs (20) and (21) show that $`E_j`$ and $`e_j`$ take a similar form, hence we here discuss branches $`E_j`$ of the excitation spectra in detail. For $`(n_r,l,m)=(0,1,0)`$, we have
$$\xi _{0,1,0}=\alpha _{010}^{3/2}[\frac{8}{3\sqrt{\pi }}]^{1/2}\alpha _{010}re^{\alpha _{010}^2r^2/2}Y_{1,0}(\theta ,\psi ).$$
The excitation spectrum in this case is reduced to
$`E`$ $`=`$ $`E[v,\omega _{010}]=(1+2v^2)[{\displaystyle \frac{5}{4}}\mathrm{}\omega _{010}+{\displaystyle \frac{5}{4}}\mathrm{}{\displaystyle \frac{\omega _a^2}{\omega _{010}}}]`$ (62)
$`+`$ $`\lambda (1+2v^2)N_m\omega _m^{3/2}[{\displaystyle \frac{2M}{\pi \mathrm{}}}]^{3/2}[{\displaystyle \frac{\omega _{010}}{\omega _{010}+2\omega _m}}]^{5/2}`$ (63)
$`+`$ $`[2\lambda _a(1+2v^2)2\lambda _av\sqrt{1+v^2}]N_a\omega _a^{3/2}[{\displaystyle \frac{M}{\pi \mathrm{}}}]^{3/2}[{\displaystyle \frac{\omega _{010}}{\omega _{010}+\omega _a}}]^{5/2}`$ (64)
$``$ $`4\alpha v\sqrt{1+v^2}[{\displaystyle \frac{\omega _{010}}{\omega _{010}+\omega _m}}]^{5/2}N_m^{1/2}\omega _m^{3/4}[{\displaystyle \frac{2M}{\pi \mathrm{}}}]^{3/4},`$ (65)
where $`\alpha _{010}^2=\frac{M\omega _{010}}{\mathrm{}}.`$Similarly, for $`(n_r,l,m)=(1,0,0),`$ we have
$$\xi _{1,0,0}=\alpha _{100}^{3/2}[\frac{8}{3\sqrt{\pi }}]^{1/2}(\frac{3}{2}\alpha _{100}^2r^2)re^{\alpha _{100}^2r^2/2}Y_{0,0}(\theta ,\psi ).$$
And
$`E`$ $`=`$ $`E[v,\omega _{100}]=(1+2v^2)[{\displaystyle \frac{7}{4}}\mathrm{}\omega _{100}+{\displaystyle \frac{7}{4}}\mathrm{}{\displaystyle \frac{\omega _a^2}{\omega _{100}}}]`$ (66)
$`+`$ $`\lambda (1+2v^2)N_m\omega _m^{3/2}[{\displaystyle \frac{2M}{\pi \mathrm{}}}]^{3/2}f(\omega _{100},2\omega _m)`$ (67)
$`+`$ $`[2\lambda _a(1+2v^2)2\lambda _av\sqrt{1+v^2}]N_a\omega _a^{3/2}[{\displaystyle \frac{M}{\pi \mathrm{}}}]^{3/2}f(\omega _{100},\omega _a)`$ (68)
$``$ $`4\alpha v\sqrt{1+v^2}f(\omega _{100},\omega _m)N_m^{1/2}\omega _m^{3/4}[{\displaystyle \frac{2M}{\pi \mathrm{}}}]^{3/4},`$ (69)
where
$$f(x,y)=\frac{3}{2}(\frac{x}{x+y})^{3/2}3(\frac{x}{x+y})^{5/2}+\frac{5}{2}(\frac{x}{x+y})^{7/2}.$$
Minimizing the energies of eqs (23) and (24) with respect to the variation parameters $`v`$, $`\omega _{010}`$ and $`\omega _{100}`$, we can determine the excitation spectrum for the mode $`(0,1,0)`$ and $`(1,0,0)`$. The numerical results are illustrated in fig.1 and fig.2. The dashed lines in figures show the excitation spectrum in atom BEC, i.e., $`\alpha =\lambda =0`$. In contrast, the solid line are those of hybrid atomic/molecular condensates near Feshbach resonance. From fig.1 and fig.2 we see that while the excitation frequency for modes $`(0,1,0)`$ increases due to the Feshbach resonance effect, the excitation frequency for the mode $`(1,0,0)`$ decreases. We would like to point out that the numerical results presented here depend on the coupling constant as well as the parameter $`\alpha `$. In Fig.1 and Fig.2, we let $`\alpha =5\lambda _a`$, and $`\lambda _a=0.1`$ (arbitrary units). The other parameters in Fig.1 and Fig.2 are $`\omega _m=1.4\omega _a=7500Hz,`$ $`N_a=N_m=10^6`$.
For clarity, we illustrate the above somewhat formal discussion by considering the binary atom system for a uniform system ($`V_{a(m)}(r)0`$), in this case $`H_{FR}`$ gives a resonant contribution to the atom-atom interaction strength $`a_a`$: $`a_{eff}=a_0(1+\frac{\mathrm{\Delta }}{B_0B})`$, where $`a_0`$ is the off-resonant scattering length, and $`\mathrm{\Delta }`$ characterizes the width of the resonance. For small $`p`$, these excitations are phonons, and their energy tends to zero with $`p`$. Hence,
$$\mu _a=\lambda |\varphi _m|^2+\lambda _a|\varphi _a|^22\alpha \varphi _m^{},$$
(70)
which leads to
$$E_j^2=E^2(p)=(\frac{\mathrm{}^2}{2m})^2p^2(p^2+16\pi na_{eff}).$$
(71)
For a uniform dilute Bose gas with negative scattering length $`a_{eff}`$, eq.(26) implys an instability of those modes with $`p^216\pi n|a_{eff}|`$. For a gas in a trap, however, the wavenumber cannot be arbitrarily small, and the minimum value is of order $`p_{min}\pi /R_0`$( $`R_0`$ is the mean size of the ground state). Hence the system can remain stable if $`\frac{\pi ^2}{R_0^2}16\pi n|a_{eff}|.`$ Since the density is of order $`nN/R_0^2`$, this means that the critical number of the system is
$$N_0\frac{\pi }{16}\frac{R_0}{|a_{eff}|}.$$
(72)
For a positive scatering length $`a_{eff}`$, however, there are not any constraints in $`N`$. The Bogoliubov quasiparticle theory shows that the condensate atoms $`N_0`$ depends on the scattering length and satisfies (for $`(a_{eff}\frac{N}{V})^{\frac{1}{3}}<<1`$).
$$N_0=N(1\frac{8}{3}\sqrt{\frac{Na_{eff}^3}{\pi V}}).$$
(73)
The numerical results of eqs.(27) and (28) are illustrated in Fig.3, which shows the atom number $`N_0`$ vs. external magnetic field $`B`$. The parameters in Fig.3 are $`N/V=N/R^2=10^{15}/cm^3`$, $`\mathrm{\Delta }=0.01mT`$.
To sum up, we have derived a set of four coupled equations of the atomic and molecular excitations within standard Hartree-Fock-Bogoliubov approximation. As shown in eqs. (9) and (13), the $`\alpha `$ terms describing the process that converts atoms into molecules play an important role in atomic/molecular Bose-Einstein condensation. In particular, two low excitation spectrum have been given at zero temperature, which show that the interaction between the hybrid atomic/molecular BEC increase one excitation mode, while they decrease another excitation mode. The mode (0,1,0) comes from the density fluctuation of the condensate like vibrating oscillation, in this sense that the mode (0,1,0) increases near the Feshbach resonance indicates the presence of the Feshbach resonance enhance the density fluctuation like vibrating oscillation in atomic/molecular condensation system, whereas the breath mode(like breathing oscillation) (1,0,0) decrease near the Feshbach resonance. Within the two-body collision regime, we show the atom number remained in BEC vs. the external magnetic field $`B`$, the result is in good agreement with the recent experiment. This work removes from consideration of the case at resonance, since at resonance the Bogoliubov approximation is not available( at resonance, there are few atoms in condensate). The contributions of the noncondensate atoms(molecules) to the excitation spectrum is also ignored (see eq.(9)). These need further investigations.
Figure Captions
Fig.1: Excitation spectrum of Mode $`(0,1,0)`$ vs. number of atoms. Dotted and dashed line indicate those with and without Feshbach resonance, respectively.
Fig.2:Same as fig.1. But for mode $`(1,0,0)`$.
Fig.3:The number of atoms in condensate vs. magnetic field $`B`$.
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# On the Degrees of Irreducible Representations of Hopf Algebras
## 1. Introduction
Let $`H`$ denote a semisimple Hopf algebra over an algebraically closed field $`k`$ of characteristic 0. According to a famous conjecture of Kaplansky \[K\], the degrees of the irreducible representations of $`H`$ are expected to be divisors of $`dim_kH`$. This has been confirmed under the additional hypothesis that $`H`$ is quasi-triangular \[EG\]; see \[S<sub>1</sub>\] for an elegant new proof.
In this note, we show by a very simple argument that the desired divisibility relation $`\frac{dimH}{dimV}`$ also holds for any irreducible $`H`$-module $`V`$ whose character $`\chi _V`$ belongs to the center $`Z(H^{})`$ of $`H^{}`$. Here, as usual, the *character* of $`V`$ is the linear form $`\chi _VH^{}`$ that is defined by $`\chi _V(h)=\mathrm{trace}_{V/k}(h_V)`$, where $`h_V\mathrm{End}_k(V)`$ is given by the action of $`hH`$. Thus:
###### Proposition.
Let $`H`$ denote a semisimple Hopf algebra over an algebraically closed field $`k`$ of characteristic 0. Then $`dim_kV`$ divides $`dim_kH`$ for any irreducible $`H`$-module $`V`$ such that $`\chi _VZ(H^{})`$.
The work presented here was inspired by Y. Sommerhäuser’s lecture “On central character rings”, delivered during the MSRI-workshop on Hopf algebras in October 1999.
The above notations will remain in effect throughout. In addition, $`\mathrm{Irr}(H)`$ will denote a fixed representative set of the isomorphism classes of irreducible (left) $`H`$-modules and $`C(H)`$ is the *character algebra*, that is, the $`k`$-span in $`H^{}`$ of all characters $`\chi _V`$ of $`H`$-modules $`V`$. In general, our notation for Hopf algebras follows \[Mo\].
## 2. Some Preliminaries
Our hypotheses entail that the dual Hopf algebra $`H^{}`$ is semisimple as well; see \[Mo\]. Recall that $`H`$ acts on $`H^{}`$ via $`h\phi ,h^{}=\phi ,h^{}h`$ for $`h,h^{}H`$ and $`\phi H^{}`$. Dually, $`H^{}`$ acts on $`H`$. Some particulars of these actions are explained in the following standard lemma originally due to Masuoka \[M\]; see also \[S<sub>2</sub>, p. 44\].
We let $`\lambda H^{}`$ denote the integral with $`\lambda ,1=1`$ and $`\mathrm{\Lambda }H`$ the integral with $`\lambda ,\mathrm{\Lambda }=1`$; cf. \[S<sub>2</sub>, p. 26\]. Moreover, $`S^{}`$ denotes the antipode of $`H^{}`$.
###### Lemma.
Let $`V\mathrm{Irr}(H)`$ and let $`e_VZ(H)`$ denote the corresponding centrally primitive idempotent of $`H`$, acting as $`\mathrm{Id}_V`$ on $`V`$ and as $`0_W`$ on all $`W\mathrm{Irr}(H)\{V\}`$. Then
$$\frac{dimH}{dimV}e_V=(S^{}\chi _V)\mathrm{\Lambda }\text{and}\frac{dimH}{dimV}e_V\lambda =\chi _V.$$
###### Proof.
The first formula is Corollary 4.6 in \[S<sub>2</sub>\] and the second one is established in the proof of \[M, Lemma (b)\]; cf. also \[S<sub>2</sub>, proof of Proposition 4.5\]. ∎
###### Corollary.
For every idempotent $`\delta =\delta ^2Z(H^{})`$, the element $`\delta \mathrm{\Lambda }C(H^{})H`$ is a character of some $`H^{}`$-module.
###### Proof.
We will apply the Lemma with the roles of $`H`$ and $`H^{}`$ interchanged. Since $`\epsilon ,\mathrm{\Lambda }=dimH`$ holds for the counit $`\epsilon =1_H^{}`$ (e.g., \[S<sub>2</sub>, Prop. 3.4\]), this requires replacing $`\mathrm{\Lambda }`$ by $`\mathrm{\Lambda }^{}=\frac{1}{dimH}\mathrm{\Lambda }`$.
First assume that $`\delta =\delta _M`$ is the centrally primitive idempotent of $`H^{}`$ corresponding to some irreducible $`H^{}`$-module $`M`$. Then, by the Lemma, $`\frac{dimH^{}}{dimM}\delta _M\mathrm{\Lambda }^{}=\chi _MC(H^{})`$. Thus,
$$\delta _M\mathrm{\Lambda }=(dimM)\chi _M=\chi _{M^{dimM}}.$$
In general, $`\delta =_M\delta _M`$ for certain irreducible $`H^{}`$-modules $`M`$, and so $`\delta \mathrm{\Lambda }=\chi _N`$ for the $`H^{}`$-module $`N=_MM^{dimM}`$. ∎
## 3. The Proof
We are now ready to give the proof of the Proposition. So let $`V\mathrm{Irr}(H)`$ be given with $`\chi _VZ(H^{})`$. Our goal is to show that $`\frac{dimH}{dimV}𝒪`$, where $`𝒪`$ denotes the ring of algebraic integers in $`k`$. Since $`\frac{dimH}{dimV}e_V=(S^{}\chi _V)\mathrm{\Lambda }`$, by the Lemma, it suffices to show that $`(S^{}\chi _V)\mathrm{\Lambda }`$ is integral over $``$.
To this end, let $`\delta _1,\mathrm{},\delta _r`$ be the distinct primitive idempotents of $`Z(H^{})`$ and let $`f_i:Z(H^{})k`$ denote the corresponding characters; so $`\zeta =_if_i(\zeta )\delta _i`$ holds for all $`\zeta Z(H^{})`$. In particular, since $`S^{}\chi _VZ(H^{})`$, we have
$$(S^{}\chi _V)\mathrm{\Lambda }=\underset{i}{}f_i(S^{}\chi _V)\delta _i\mathrm{\Lambda }.$$
By the Corollary, each $`\delta _i\mathrm{\Lambda }`$ is the character of some (irreducible) $`H^{}`$-module; so $`\delta _i\mathrm{\Lambda }`$ belongs to the *Grothendieck ring* $`G_0(H^{})=_{M\mathrm{Irr}(H^{})}\chi _MC(H^{})`$. Every element of $`G_0(H^{})`$ satisfies a monic polynomial over $``$. Analogously, the character $`S^{}\chi _V=\chi _V^{}G_0(H)`$ satisfies a monic polynomial over $``$, and so $`f_i(S^{}\chi _V)𝒪`$ holds for all $`i`$. Consequently,
$$(S^{}\chi _V)\mathrm{\Lambda }G_0(H^{})𝒪C(H^{}),$$
which entails that $`(S^{}\chi _V)\mathrm{\Lambda }`$ is integral over $``$, as desired.
## 4. Concluding Remarks
Let $`(.)^{\mathrm{cl}}`$ denote the set of elements in the ring in question that are integral over $``$; e.g.,
$$Z(H)^{\mathrm{cl}}=\underset{V\mathrm{Irr}(H)}{}𝒪e_V.$$
Consider the isomorphism $`f:H^{}\stackrel{}{}H`$, $`f(\phi )=\phi \mathrm{\Lambda }`$. By the Lemma, $`f`$ restricts to isomorphisms $`C(H)\stackrel{}{}Z(H)`$ and $`Z(H^{})\stackrel{}{}C(H^{})`$. The essence of the above proof is that, in fact,
$$Z(H^{})^{\mathrm{cl}}\stackrel{f}{}G_0(H^{})𝒪C(H^{})^{\mathrm{cl}},$$
while Kaplansky’s conjecture is equivalent with
$$G_0(H)\stackrel{f}{}Z(H)^{\mathrm{cl}}.$$
It is tempting to try and consolidate the Etingof-Gelaki result for quasi-triangular Hopf algebras and our Proposition by at least showing that $`f`$ maps the *center* of $`G_0(H)`$ to $`Z(H)^{\mathrm{cl}}`$.
###### Acknowledgment.
The work on this note was started while I was attending the Noncommutative Algebra program at MSRI in the Fall of 1999. I would like to thank the organizers of this program, in particular S. Montgomery, and MSRI staff for their hospitality and support. I would also like to thank Temple University for granting me a research leave during the Fall Semester 1999.
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# Detection and Mapping of Decoupled Stellar and Ionized Gas Structures in the Ultraluminous Infrared Galaxy IRAS 12112+03051
## 1 INTRODUCTION
Ultraluminous infrared galaxies (ULIRGs), with bolometric luminosities $`L_{\mathrm{bol}}`$ $`L_{\mathrm{IR}}10^{12}`$ $`L_{}`$, are the most luminous galaxies in the local Universe. ULIRGs show signs of strong interactions and mergers (Melnick & Mirabel 1990, Clements et al. 1996, Borne et al. 2000a), or even multiple collisions (Borne et al. 2000b), and they have large amounts of gas and dust that significantly obscure the nuclear ionizing sources (Sanders & Mirabel 1996). ULIRGs appear to be the low-redshift analogs of the high-redshift submillimeter galaxies (SCUBA sources), which are responsible for the bulk of the far-infrared background radiation (see Sanders 1999 for a review).
Mid-infrared spectroscopy has shown that ULIRGs with optical H II- and LINER-like spectra are starbursts dominated (Lutz et al. 1998; Genzel et al. 1998; Lutz, Veilleux, & Genzel 1999). However, the increased fraction of optically classified Seyfert galaxies among the more luminous ULIRGs (log$`L_{\mathrm{IR}}12.3`$ $`L_{}`$) is taken as evidence for the presence of dust-enshrouded quasars powering these galaxies (Veilleux, Kim, & Sanders 1999).
The ionization structure of ULIRGs is extended and rather complex, and therefore narrow-slit ($``$ 2$`^{^{\prime \prime }}`$) optical spectroscopy along a given position angle cannot adequately indicate the location, nature, and interplay of the different ionization mechanisms present in these galaxies. The spectral classification of some ULIRGs are already known to depend on the size of the aperture (Kim, Veilleux, & Sanders 1998). Two-dimensional optical spectroscopy using fiber-fed spectrographs is an ideal technique to study the complex stellar and ionized gas structure of the nuclear regions of galaxies in general (Arribas & Mediavilla 2000), and of ULIRGs in particular since it allows a simultaneous and complete mapping of the stellar populations and the ionized gas (e.g., see Colina, Arribas, & Borne 1999).
IRAS 12112+0305 is a low redshift (z= 0.0723) ULIRG (log$`L_{\mathrm{IR}}`$= 12.34 $`L_{}`$), classified in the optical as a LINER (Kim et al. 1998), and with no evidence for a hidden broad-line region (Veilleux, Sanders, & Kim 1999). Mid-infrared diagnostics classify this galaxy as starburst-dominated (Genzel et al. 1998). Near-infrared ground-based images revealed the presence of two nuclei (Carico et al. 1990), while HST optical and near-infrared images show a well defined compact point-like nucleus, as well as several much fainter condensations distributed in an arc-like structure (Borne et al. 2000a; Scoville et al. 2000).
In this letter we highlight new results for IRAS 12112+0305 obtained using integral field optical spectroscopy and HST imaging. Throughout the paper a Hubble constant of 70 km s<sup>-1</sup> Mpc<sup>-1</sup> is assumed.
## 2 OBSERVATIONS AND DATA REDUCTION
Integral field spectroscopy of IRAS 12112+0305 was obtained with the INTEGRAL system (Arribas et al. 1998), a fiber-fed spectrograph mounted at the Naysmith No. 1 platform of the 4.2m William Herschel Telescope. The bundle of fibers consisted of 219 fibers, each 0.9$`^{^{\prime \prime }}`$ in diameter and covering a 16.5$`^{^{\prime \prime }}`$ $`\times `$ 12.3$`^{^{\prime \prime }}`$ field-of-view. The spectra were taken using a 600 line/mm grating, covering the 5000$``$7900Å range, with an effective resolution of 4.8Å. The total integration time was 7500 sec, split into five separate integrations of 1500 sec each, with seeing $`1.0^{^{\prime \prime }}`$.
The reduction has been done following the standard procedures applied to spectra obtained with two-dimensional fiber spectrographs (Arribas et al. 1997 and references therein). The results are presented in Figure 1, together with the HST $`I`$-band image (Borne et al. 2000a) for comparison.
The derived astrophysical properties for the main continuum and line-emitting regions are presented in Table 1. Columns 2 and 3 give the relative positions. Column 4 presents the internal reddening values derived from the H$`\alpha `$/H$`\beta `$ ratio, assuming case B recombination. Column 5 gives the apparent magnitude within an aperture of 0.5$`^{^{\prime \prime }}`$ radius using the HST $`I`$-band image, and column 6 gives the corresponding absolute magnitude after internal reddening correction assuming $`A_\mathrm{I}`$= 1.494 $`\times `$ E(B$``$V). Column 7 shows the observed H$`\alpha `$ flux obtained using an aperture of 2.1$`^{^{\prime \prime }}`$ diameter, while column 8 gives the reddening-corrected H$`\alpha `$ luminosity. The last four columns present the logarithm of the reddening-corrected emission line ratios and the corresponding activity classification.
## 3 RESULTS AND DISCUSSION
### 3.1 Evidence for Decoupled Stellar and Ionized Gas Components
The stellar main body of IRAS 12112+0305 is concentrated in three dominant regions, each separated from each other by about 2$`{}_{}{}^{^{\prime \prime }}3^{^{\prime \prime }}`$ (see the H$`\alpha `$ and H$`\beta `$ continuum images in Figure 1). Two of these regions (called N<sub>s</sub> and N<sub>n</sub> hereinafter; see Table 1 for relative positions) located along position angle PA45, were already detected in the near-infrared, and associated with the nuclei of two galaxies involved in the final stages of a merger (Carico et al. 1990). The INTEGRAL-generated continuum images extend the wavelength coverage towards the blue, clearly showing the presence of a large differential extinction towards the southern nucleus (N<sub>s</sub>). The third region (called R2 hereinafter) is located 3$`^{^{\prime \prime }}`$ north of N<sub>s</sub> at position angle PA20.
The HST $`I`$-band image (Figure 1) shows N<sub>s</sub> as a high-surface brightness compact region which coincides with the K-band point-like nucleus (Scoville et al. 2000). Regions N<sub>n</sub> and R2 are made of several fainter condensations distributed along an arc-like structure of about 5.4$`^{^{\prime \prime }}`$ (i.e., 8 kpc) extent and located 3$`^{^{\prime \prime }}`$ to 5$`^{^{\prime \prime }}`$ (i.e., 4.5 to 7.5 kpc) north$``$northeast of N<sub>s</sub>. The brightest of these condensations, located 3$`^{^{\prime \prime }}`$ ($``$ 4.5 kpc) northeast from N<sub>s</sub>, is most likely associated with the northern nucleus detected in the INTEGRAL-generated continuum images (N<sub>n</sub>) and in the near-infrared (Carico et al. 1990; Scoville et al. 2000).
Although the overall structure of the ionized gas resembles that of the stellar light distribution, the dominant line-emitting regions do not coincide with the nuclei identified above but, on the contrary, are decoupled from them (see \[O III\] and H$`\alpha `$ maps in Figure 1). The brightest line-emitting region (called R1 hereinafter) is located 5$`^{^{\prime \prime }}`$ ($``$ 7.5 kpc) east of N<sub>s</sub> along position angle PA80. This region is associated with a faint $`I`$-band continuum source ($`m_I`$= 20.4) outside the main body of the galaxy (see HST $`I`$-band image in Figure 1). The second-brightest H$`\alpha `$ line-emitting region is associated with region R2, and it is therefore composed of a filament of faint continuum condensations that are also detected in the H$`\beta `$ and \[O III\] maps (Figure 1). The optically dominant nucleus of the galaxy (N<sub>s</sub>) is marginally detected in \[O III\] and appears as the faintest H$`\alpha `$ emission source. The arc-like structure seen in the HST $`I`$-band image corresponds to three well resolved H$`\alpha `$ line-emitting peaks, the faintest of which is associated with the northern nucleus (N<sub>n</sub>) detected in the optical and near-infrared (see discussion above).
### 3.2 Massive Dust-Enshrouded Starbursts as Nuclear Ionizing Sources
The positions of the two apparently faint ionizing nuclei N<sub>n</sub> and N<sub>s</sub> coincide (within the HST absolute astrometry uncertainties of $`\pm `$0.7$`{}_{}{}^{^{\prime \prime }}`$1$`^{^{\prime \prime }}`$) with the two bright compact starburst nuclei detected at radio frequencies (Condon et al. 1991), thus favoring dust-enshrouded nuclear starbursts as the energy sources in these nuclei.
The ionization sources in N<sub>s</sub> and N<sub>n</sub> are highly obscured, with visual extinctions ($`A_\mathrm{V}`$) of 8 and 3.5 magnitudes, respectively (see spectra in Figure 2). The two nuclei appear almost equally bright with $`I`$-band absolute magnitudes of $``$18.9 (N<sub>s</sub>) and $``$18.5 (N<sub>n</sub>). However, when reddening is taken into account, N<sub>s</sub> is found to be ten times more luminous than N<sub>n</sub> (i.e. $`M_I=`$22.8). These magnitudes are within the range of those measured for the nuclei of other luminous and ultraluminous infrared galaxies (Surace et al. 1998).
The line ratios of the ionization sources associated with the nuclei N<sub>s</sub> and N<sub>n</sub> correspond to a mixture of weak-\[O I\] LINER and H II region spectra (see Table 1). The presence of weak-\[O I\] LINERs in the nucleus of galaxies has been taken as evidence for ionization by hot stars in a high-metallicity environment (Filippenko & Terlevich 1992; Shields 1992) or for ionization by a mixture of a low-luminosity AGN and hot stars (Ho et al. 1993). The integrated mid-infrared spectrum of IRAS 12112+0305 does not show any evidence for an AGN (Genzel et al 1998), further supporting the idea that the ionizing sources in the nuclei N<sub>s</sub> and N<sub>n</sub> are associated with dust-enshrouded starbursts. Although N<sub>s</sub> and N<sub>n</sub> are minor contributors to the observed H$`\alpha `$ emission, the reddening-corrected flux emanating from these nuclei dominates the overall H$`\alpha `$ luminosity with a value of 9.1 $`\times `$ 10<sup>42</sup> erg s<sup>-1</sup> (Table 1). If the H$`\alpha `$ flux emitted by the two nuclei were entirely due to stars, then the corresponding star formation rate would amount to about 80 $`M_{}`$ yr<sup>-1</sup> for a Salpeter initial mass function (IMF) with mass limits of 0.1 and 100 $`M_{}`$ (Leitherer et al. 1999).
### 3.3 Brightest Line-Emitting Region: A Tidally-Induced Giant H II Region?
The HST image (Figure 1) shows that the apparently more luminous line-emitting peak (R1) is associated with a faint ($`m_I`$= 20.35) region characterized by a physical size of about 600 parsecs ($``$ 0.4$`^{^{\prime \prime }}`$), an $`I`$-band absolute magnitude of $``$17.7, an internal optical extinction of about one magnitude, and an H$`\alpha `$ luminosity of 8.7 $`\times `$ 10<sup>40</sup> erg s<sup>-1</sup> (see Table 1).
The size, H$`\alpha `$ luminosity and emission line ratios are typical of circumnuclear star-forming regions in nearby spirals (Gonzalez-Delgado & Pérez 1997; Planesas, Colina, & Pérez-Olea 1997), and of giant extragalactic H II regions like NGC 5471 in M101 (Shields 1990). The H$`\alpha `$ luminosity and the equivalent widths of the H$`\beta `$ (74 $`\pm `$ 3Å) and H$`\alpha `$ (450 $`\pm `$ 10 Å) emission lines correspond to that of a 5-Myr old ionizing cluster of 2 $`\times `$ 10<sup>7</sup> $`M_{}`$, assuming a Salpeter IMF with mass limits of 0.1 and 100 $`M_{}`$ (Leitherer et al. 1999). The mass of the ionizing cluster represents only a small fraction ($``$ 3%) of the dynamical mass of this region (upper limit of 7.5 $`\times `$ 10<sup>8</sup> $`M_{}`$), calculated assuming virialization and an instrument-corrected emission line width equivalent to a velocity dispersion of 60 km s<sup>-1</sup> (derived from the line-width measurements of the H$`\beta `$ and \[O III\]5007Å lines).
In summary, the derived properties of this region are characteristic of young massive H II regions and could represent a case for a tidally-induced giant extranuclear self-gravitating star-forming region, or even a dwarf galaxy, decoupled from the much older stellar body of the parent galaxies (Duc & Mirabel 1994, 1998).
## 4 SUMMARY
Integral field optical spectroscopy and HST imaging have revealed the complex stellar and ionized gas structures in the ultraluminous infrared galaxy IRAS 12112+0305. The two nuclei detected in the optical coincide with previously known radio and near-infrared nuclei, but appear as the apparently two faintest line-emitting regions and are consistent with dust-enshrouded (A<sub>V</sub>= 3.5 and 8 mag) weak-\[O I\] LINERs. However, their reddening-corrected H$`\alpha `$ emission dominates the overall H$`\alpha `$ luminosity of the system and, if associated with nuclear starbursts, would correspond to a star formation of 80 $`M_{}`$ yr<sup>-1</sup>.
The observed structure of the ionized gas is decoupled from, and hence does not trace, the stellar light distribution. The brightest line-emitting peak is associated with a faint ($`m_I`$= 20.4) region located 7.5 kpc from the dominant optical nucleus (N<sub>s</sub>), well outside the main stellar body of the system. This region appears to be a recent (tidally induced?) star-forming region containing a young (5 Myr), massive (2 $`\times `$ 10<sup>7</sup> $`M_{}`$) cluster of stars.
The results presented here for IRAS 12112+0305 stress the need for integral-field spectrographs in the study of the complex ionization and stellar light structure of ultraluminous infrared galaxies. Similar studies for high-redshift analogs of ULIRGs (i.e., the SCUBA sources), and for other morphologically complex high-redshift galaxies, will become feasible when integral-field spectrographs on giant telescopes become operational.
L. Colina thanks the Instituto de Astrofísica de Canarias for its hospitality and financial support. K. Borne thanks Raytheon for providing financial support during his Sabbatical Leave. Support for this work was provided by CICYT (Comisión Interministerial de Ciencia y Tecnología) through grants numbers PB98-0340-C02-01 and PB98-0340-C02-02, and by NASA through grant number GO-06346.01-95A from the Space Telescope Science Institute, which is operated by AURA, Inc., under NASA contract NAS5-26555.
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# Circumnuclear Supernova Remnants and H II Regions in NGC 253
## 1 Introduction
The edge-on spiral galaxy NGC 253 is one of the two prototypical, nearby starburst galaxies (along with M82). Its distance of 2.5 Mpc (Turner & Ho, 1985) enables its starburst disk to be studied at very high linear resolution, since 1″ corresponds to only 12 pc. The galaxy contains an inner starburst disk roughly 15″–20″ (180–240 pc) in extent that has been studied in detail, especially by virtue of its emission in the infrared (Forbes, Ward, & DePoy, 1991; Piña et al., 1992; Keto et al., 1993; Forbes et al., 1993; Kalas & Wynn-Williams, 1994; Sams et al., 1994; Böker, Krabbe, & Storey, 1998; Keto et al., 1999) and in molecular lines (Canzian, Mundy, & Scoville, 1988; Israel, White, & Baas, 1995; Jackson et al., 1995; Paglione, Tosaki, & Jackson, 1995; Peng et al., 1996; Frayer, Seaquist, & Frail, 1998). Subarcsecond-resolution imaging at centimeter wavelengths reveals at least 64 compact radio sources in the inner disk (Turner & Ho, 1985; Antonucci & Ulvestad, 1988; Ulvestad & Antonucci, 1991, 1994, 1997); these are thought to be roughly equally divided between H II regions dominated by thermal emission, and supernova remnants dominated by nonthermal emission (Ulvestad & Antonucci, 1997). At the centimeter wavelengths where radio telescopes are most sensitive, the poorer resolution and the increasing strength of the more diffuse galaxy emission lead to confusion that prevents complete source identification in the inner 200-pc starburst.
The 200-pc region clearly dominates the current star formation in NGC 253, as shown by the infrared, millimeter, and centimeter observations. Centimeter imaging of this inner starburst indicates a supernova rate of $``$0.3 yr<sup>-1</sup> (Ulvestad & Antonucci, 1997), which is consistent with results of 0.1–0.3 yr<sup>-1</sup> inferred from models of the infrared emission of the entire galaxy (Rieke et al., 1980; Rieke, Lebofsky, & Walker, 1988). Near- and mid-infrared imaging at arcsecond resolution shows a number of emission peaks that do not generally line up with the radio sources (Keto et al., 1993; Sams et al., 1994; Kalas & Wynn-Williams, 1994), while Hubble Space Telescope images show at least four compact star clusters, which also do not coincide with individual compact radio sources (Watson et al., 1996). Thus, it is apparent that extinction and confusion play major roles in the inner disk at different wavebands, and it may be that only the centimeter and millimeter images reveal the regions of most recent star formation.
Outside the 200-pc starburst disk, NGC 253 may have a significant amount of star formation, as shown by the existence of larger-scale radio emission (Hummel, Smith, & van der Hulst, 1984; Carilli et al., 1992; Beck et al., 1994; Carilli, 1996). Since this part of the galaxy is much less affected by confusion and extinction, there is greater potential for identifying the strongest regions of star formation and studying them in multiple bands. However, recent high-resolution studies of the galaxy outside the inner few hundred parsecs are uncommon, since most of the imaging has been of small fields centered on the main starburst disk. One recent study on the larger scale is that of Vogler & Pietsch (1999), who used ROSAT to identify 73 X-ray sources in the bulge, disk, and halo of NGC 253, attributing most of them to X-ray binary stars. To study the population of compact radio sources on this large scale, we have reprocessed the high-resolution 6-cm and 20-cm Very Large Array (VLA) data obtained in 1987 (Antonucci & Ulvestad, 1988; Ulvestad & Antonucci, 1997) in order to image the entire primary beam of the individual VLA antennas at high resolution; limited computing resources prevented this from being done at the time of the observations. In this paper, we report the locations of the compact radio sources outside the inner 200-pc starburst, compare these locations to images made in other wavebands, and estimate the supernova rate outside the inner disk.
## 2 Observations and Data Analysis
The observations of NGC 253 that we use here were made in 1987, using the VLA (Thompson et al., 1980) in its A configuration. Results from imaging the inner starburst at 6 cm (4860 MHz) were reported by Antonucci & Ulvestad (1988), while the 20-cm (1490 MHz) image of the same region was discussed briefly by Ulvestad & Antonucci (1997). Details of the observations were reported in those papers. Other 6-cm observations had much poorer (u,v) coverage than those made in 1987, so only the 1987 data are considered here.
For self-consistency, the data were re-calibrated and self-calibrated using the same procedures reported previously. We then made large-scale ($`4096\times 4096`$ pixels) images of the radio emission of NGC 253 to reveal the compact radio sources seen throughout the primary beam of the individual 25-m telescopes of the VLA. These images covered $`26^{}\times 26^{}`$ at 20 cm and $`8^{}\times 8^{}`$ at 6 cm. Outside the starburst disk, confusion was not a significant issue, but sensitivity was critical. We therefore used natural weighting of the data in the (u,v) plane (Sramek & Schwab, 1989; Briggs, Schwab, & Sramek, 1999); this scheme produces the maximum sensitivity, at the price of a modest loss in resolution. A zero-spacing (total) flux density was specified in the imaging, in order to partially compensate for the lack of short interferometer spacings. This reduced, but did not completely remove, systematic effects caused by diffuse emission that was not sampled well in our observations.
Table 1 gives the resolution and the r.m.s. noise achieved for the final images. Several effects contribute to reductions in the sensitivity far from the field center; each is discussed briefly here. First, the attenuation of the primary beam of the individual VLA antennas reduces the sensitivity by more than 10% at 6 cm and 20 cm for respective distances from the pointing direction of greater than $`1.^{}7`$ and $`5.^{}4`$. Corrections for this attenuation were made using the AIPS (van Moorsel, Kemball, & Greisen, 1996) software developed by NRAO, and were less than 10% except for several sources (see Section 4.1) far from the main area of radio emission. A second effect, chromatic aberration (“bandwidth smearing”) due to the non-zero observing bandwidth (Bridle & Schwab, 1999), is more deleterious. It reduces the peak flux density by more than a factor of two at distances greater than 2′ from the phase center, and is progressively worse farther out in the field. The actual phase center of the observations was $`10^{\prime \prime }`$ north of the apparent galaxy nucleus, at a B1950 position of $`(\alpha ,\delta )=(00^h45^m05.^s80,25^{}33^{}29.^{\prime \prime }0)`$; since most of the large-scale sources are north of the nucleus, this actually was slightly beneficial in the final data analysis. Compared to the chromatic aberration, the reduction in flux density due to delay smearing (Bridle & Schwab, 1999) over the observations’ 10-second averaging period is negligible. Finally, the lack of short interferometer spacings causes correlated positive or negative emission as high as 0.08–0.1 mJy over areas of $`50\times 50`$ pixels away from the center of the field. The final source detection threshold that was used was four times the quadrature sum of this systematic value and the r.m.s. noise, or approximately 0.4 mJy at each wavelength. However, this threshold rises substantially at distances more than a few arcminutes from the phase center.
## 3 Images and Source Identification
Figure 1 is a portion of the 20-cm image showing the circumnuclear region of NGC 253 at high resolution. The locations of compact X-ray sources detected by Vogler & Pietsch (1999), using the ROSAT High Resolution Imager, are shown as crosses whose sizes indicate the errors in the X-ray positions. We have identified all radio sources that have apparent peak flux densities of at least 0.4 mJy beam<sup>-1</sup> at either 6 cm or 20 cm, outside the central 200-pc starburst. Figures 2a and 2b are respective enlargements of the 20-cm and 6-cm images showing the radio emission from the 200-pc starburst disk and its surroundings, including a number of the circumnuclear compact sources. Note that the 6-cm image shows little detail, due to the larger number of beam areas needed to cover the given region, but it does serve to illustrate the relative locations of the detected sources.
Table 2 lists the source identifications in the circumnuclear region at the two wavelengths. (Each source listed in this table can be seen in the region shown in Figure 1.) Included are the source positions, flux densities at 6 cm and 20 cm, measured sizes, and a few comments. When compact sources are identified at both wavelengths, the higher resolution 6-cm position is quoted; typical position accuracies are $`0.^{\prime \prime }1`$$`0.^{\prime \prime }2`$. At 6 cm, flux densities as low as 0.2 mJy are given for sources that meet the 0.4-mJy threshold at 20 cm, but these flux densities have very large fractional errors. Source flux densities and sizes have been derived by making Gaussian fits to the individual sources in the image plane, and confirmed by integrating over the region of significant emission. Most sources are unresolved in the full-resolution images at both 6 cm and 20 cm, as indicated by a “U” in the table. The sizes of resolved sources are derived from the 6-cm full-resolution image. The total 6-cm flux densities and approximate sizes found at full resolution generally are consistent with values found from a 6-cm image tapered to the 20-cm resolution.
We estimated flux density errors due to several causes. First, the areas of apparent systematic negative or positive flux near the main starburst disk, caused by undersampling in the aperture plane, average as much as 0.08 mJy beam<sup>-1</sup> at 6 cm and 0.10 mJy beam<sup>-1</sup> at 20 cm; this offset is denoted by $`\sigma _o`$. Second, there is a $`5`$% error in the absolute flux density scales of each of the two maps, so $`\sigma _{\mathrm{sc}}=0.05S`$, where $`S`$ is the source flux density. Third, there is a fitting error, $`\sigma _f`$, caused by confusion with underlying diffuse emission and uncertainties in the source fits. Since the newly identified sources reported here are outside the main starburst disk and fairly well isolated, we take the confusion error to be negligible, so $`\sigma _f`$ is simply the uncertainty reported by the least-squares fitting program. This value includes the r.m.s. noise given in Table 1, so the noise is not included separately. We combine these errors in quadrature to get the total flux-density error, $`\sigma `$, using
$$\sigma ^2=\sigma _{o}^{}{}_{}{}^{2}+(0.05S)^2+\sigma _{f}^{}{}_{}{}^{2}.$$
(1)
Thus, the total flux-density error for a 0.4-mJy source is typically 0.09 mJy at 6 cm and 0.11 mJy at 20 cm.
## 4 Discussion
### 4.1 Are All Detected Sources Associated with NGC 253?
There are 22 sources above the 0.4-mJy limit at 20 cm, inside a box 4′ (2.8 kpc) on a side and centered on the nucleus of NGC 253. The density of background sources above this flux-density threshold is expected to be $`10^6`$ sr<sup>-1</sup> (Windhorst et al., 1985). Therefore, we expect about one background source stronger than 0.4 mJy within a field of 16 square arcminutes, implying that almost all of the detected sources are associated with NGC 253. The apparent lack of weak sources at greater distances from the nucleus may be caused primarily by chromatic aberration (see Section 2); sources far from the phase center of the observation have their peak flux densities reduced substantially, and therefore may fall below the detection threshold. Observations using narrow spectral-line channels rather than a broadband continuum would be needed to determine whether other weak compact sources exist farther from the phase center.
Five additional sources were detected at 20 cm within an area 12′ on a side, centered on NGC 253, and are listed in Table 3. These are within the size of the optical galaxy, which has a measured diameter of 25′ to 25th magnitude (de Vaucouleurs et al., 1991). Images of these sources are not shown, since they are substantially degraded by chromatic aberration, which prevents measurement of the source sizes. Therefore, only the total flux densities are given in Table 3. The source locations are best seen in previous lower-resolution radio images of NGC 253, particularly that published by Anantharamaiah & Goss (1990), though none of those images resolves source B2 into the three components distinguished here. Only sources much stronger than 0.4 mJy can be detected in the large region due to the chromatic aberration. The expected density of sources stronger than 15 mJy at 20 cm is $`10^4`$ sr<sup>-1</sup> (Windhorst et al., 1985), so we expect to detect less than one background source above this threshold. Source B3 appears to be associated with a “spur” noted by Carilli et al. (1992), and the three-component source B2 lies along the large-scale galaxy disk, making it probable that at least B2 and B3 are associated with NGC 253.
### 4.2 Circumnuclear Star Formation
The detected circumnuclear sources (Table 2) outside the central starburst have typical flux densities of 0.4–3 mJy at 6 cm and 20 cm, corresponding to radio powers of $`3\times 10^{17}`$ to $`2\times 10^{18}`$ W Hz<sup>-1</sup>. Most of them seem to have relatively steep spectra; for a source with a flux density of 0.4 mJy at 6 cm and a spectral index of $`0.7`$ (with $`S_\nu \nu ^{+\alpha }`$), the total luminosity between 10 MHz (or 100 MHz) and 100 GHz is $`1.1\times 10^{35}`$ erg s<sup>-1</sup>. This luminosity is a factor of 100–1000 lower than the luminosities of the point X-ray sources detected by Vogler & Pietsch (1999), which they generally attribute to X-ray binaries. Since none of the radio sources listed in Table 2 coincides with a compact X-ray source, the radio emission is probably not associated with evolved binaries.
In the main 200-pc starburst disk of NGC 253, Antonucci & Ulvestad (1988) deduced that $`100`$ compact radio sources exist. A supernova rate of $`0.3`$ yr<sup>-1</sup> was estimated from the radio emission (Ulvestad & Antonucci, 1997), while supernova rates of 0.1–0.3 yr<sup>-1</sup> have been derived on other grounds (Rieke et al., 1980; Rieke, Lebofsky, & Walker, 1988). Here, we report an additional 22 circumnuclear compact radio sources that lie outside the central 200-pc starburst, but within 2 kpc (170″) of the galaxy center. Since they do not coincide with X-ray binaries, and most are less than 5–10 pc ($`0.^{\prime \prime }4`$$`0.^{\prime \prime }8`$) in diameter, they most likely are due to supernova remnants and H II regions, as are the sources in the nuclear region. The sources can be compared to the galactic supernova remnant Cas A. For an assumed distance of 2.8 kpc (van den Bergh, 1971) and a flux density of $`500`$ Jy (Baars et al., 1977), Cas A has a 6-cm radio power of $`5\times 10^{17}`$ W Hz<sup>-1</sup>, comparable to the powers of the weaker circumnuclear sources in NGC 253.
There are 10–15 circumnuclear sources with steep radio spectra, likely to be supernova remnants. The number of steep-spectrum sources is uncertain because some are too weak for good spectral-index determination. In addition, since the (u,v) coverage at 6 cm is not matched to that at 20 cm, more flux density may be “resolved out” at 6 cm, which could lead to estimates of spectra that are overly steep. (This resolution effect is the primary reason that spectral indices are not quoted in Table 2.) By comparison, there are at least 32 steep-spectrum sources in the inner 200-pc starburst (Ulvestad & Antonucci, 1997); correcting for confusion would probably increase this number to 50 or more. If the circumnuclear radio sources have the same general character as those in the main starburst disk, a simple comparison with the analysis of Ulvestad & Antonucci (1997) indicates that the estimated circumnuclear supernova rate outside the central starburst is $`0.1`$ yr<sup>-1</sup>. (Of course, some sources outside the strong starburst could have been missed, due primarily to chromatic aberration.) In any case, based solely on the ratio of steep-spectrum sources inside and outside the 200-pc starburst, we infer that at least 20%–30% of the global star-formation and supernova remnants are outside that central starburst. No supernovae have been detected in the circumnuclear region during limited searches made at optical wavelengths in 1988–1991 (Richmond, Filippenko, & Galisky, 1998), and at near-infrared wavelengths in late 1993 (Grossan et al., 1999). However, such searches did not cover long enough time periods to expect supernova detections in NGC 253, and also could have been affected by obscuration and confusion.
Figure 3 shows a complex of sources to the west of the main starburst disk, roughly 800 pc from the galaxy nucleus. Figure 3a is the 20-cm image, showing four individual radio sources. These sources are numbered 4, 6, 7, and 8 in Table 2 and similarly labeled in the figure. Two 6-cm images also are shown. Figure 3b is an image for which the visibility data were tapered and then restored with a point-spread function identical to that at 20 cm, while Figure 3c is a full-resolution image. Inspection of Table 2 reveals that the two strongest 6-cm sources shown in Figure 3b have flat or inverted spectra. Summing the four individual compact sources gives a total 6-cm flux density of $`6.6\pm 0.4`$ mJy and a total 20-cm flux density of $`8.3\pm 0.4`$ mJy, while integration of the images over the region of Figure 3 yields a total 6-cm flux density of 7.7 mJy and a total 20-cm flux density of 9.0 mJy, slightly higher than the fits to the compact sources.
We take the somewhat simplistic step of making a two-component decomposition of the flux density in the radio sources shown in Figure 3. To do so, we assume the presence of a thermal, flat-spectrum component having $`\alpha =0.1`$, and a nonthermal, steep-spectrum component due to supernova remnants. The steep-spectrum sources in the inner starburst of NGC 253 have typical spectral indices of $`\alpha =0.7`$ (Ulvestad & Antonucci, 1997), while galactic supernova remnants have a median spectral index of $`\alpha =0.5`$ (Green, 1996), so there is some range in the estimates of the flat-spectrum component. The simple spectral decomposition gives thermal radio emission containing 5.3 to 5.9 mJy out of a total of 6.6 mJy at 6 cm, for a steep-spectrum component having a spectral index ranging from $`0.5`$ to $`0.7`$. Thus, 80% to 90% of the 6-cm flux density in the region shown in Figure 3 appears to be thermal in origin.
The natural conclusion is that the flat-spectrum radio component is thermal radiation from H II regions energized by young stars. Applying the analysis of the strongest flat-spectrum source in the starburst disk, which was described by Ulvestad & Antonucci (1997), we find that the equivalent of $`70`$ O5 stars is necessary to energize the radio emission within a region of about $`12^{\prime \prime }\times 10^{\prime \prime }`$ ($`144\times 120`$ pc). This is comparable to the number of young stars needed to account for the strongest thermal source in the inner starburst, but those stars are contained in a volume $`6\times 10^5`$ times higher. Therefore, the average thermal gas density in the 130-pc region shown in Figure 3 is $`20`$ cm<sup>-3</sup>, much lower than the values of $`10^4`$ cm<sup>-3</sup> found in the dense part of the inner starburst of NGC 253 (Ulvestad & Antonucci, 1997) and in NGC 5253 (Turner, Ho, & Beck, 1998), or of $`10^3`$ cm<sup>-3</sup> for the typical thermal radio sources in NGC 4038/4039 (Neff & Ulvestad, 2000).
Few high-resolution observations at other wavelengths are available for the region shown in Figure 3. However, a moderate-resolution ($`9^{\prime \prime }\times 17^{\prime \prime }`$) VLA image (Carilli, 1996) shows a source at the same location, with an apparent 3.6-cm flux density between 5 and 6 mJy. This is entirely consistent with the value expected for for the flat-spectrum component deduced above. Near-infrared images shown in Figures 1 and 10 of Engelbracht et al. (1998) indicate a slight enhancement in K (2.2 $`\mu `$m) at the location of the complex of compact radio sources, apparently along an inner spiral arm. This infrared source may represent a highly reddened set of H II regions containing the numerous young stars that energize the local radio complex.
A similar argument can be made for source 15, which lies nearly 2 kpc from the galaxy nucleus, and also appears to have a thermal spectrum. Ionization of this radio source requires the equivalent of about 35 O5 stars in a region about 15–30 pc in diameter, and an average ionized density near 200 cm<sup>-3</sup>. Chromatic aberration makes these estimates somewhat uncertain, but they should be correct to 50% or better.
The relatively low average densities derived for the two strongest thermal-emitting complexes in the circumnuclear region imply that they contain more “normal” star formation, rather than the intense starbursts characteristic of the central 200 parsecs. The number and the density of massive young stars are significantly higher than in Orion and other nearby Galactic O-B associations (Blaauw, 1964), and may be comparable to the richest star-forming regions in our Galaxy, such as W49 (Welch, 1993; de Pree, Mehringer, & Goss, 1997). However, the intensity of massive star formation is considerably lower than in 30 Doradus (Hunter et al., 1995) or in the super star clusters seen in a number of starburst galaxies (O’Connell, Gallagher, & Hunter, 1994; Schweizer & Seitzer, 1998; Turner, Ho, & Beck, 1998; Whitmore et al., 1999; Kobulnicky & Johnson, 1999; Neff & Ulvestad, 2000).
## 5 Summary
We have used archival VLA data to image the circumnuclear region of NGC 253 at arcsecond resolution. Twenty-two compact radio sources have been found in the inner 2 kpc of the galaxy, but outside the well-known 200-pc disk, and most of these are probably associated with regions of recent star formation. The supernova rate inferred outside the central starburst is $``$0.1 yr<sup>-1</sup>; this may be a slight underestimate due to the decreasing sensitivity of the radio observations at distances more than $`2`$′ from the nucleus. Therefore, the region outside the well-studied inner starburst seems to account for a significant fraction of the recent star formation in NGC 253. A collection of sources located 800 pc to the west of the nucleus appears to be a complex of H II regions energized by the equivalent of 70 O5 stars, but with an average ionized gas density $`10^3`$ times lower than that found in the inner starburst of the galaxy.
I thank K. Anantharamaiah, R. Antonucci, N. Mohan, and W. Pietsch for useful discussions and for providing data in advance of publication. I especially thank the referee, Jean Turner, for some perceptive comments and for pointing out errors in Table 2. This research has made use of the NASA/IPAC Extragalactic Database (NED) which is operated by the Jet Propulsion Laboratory, California Institute of Technology, under contract with the National Aeronautics and Space Administration. In addition, this research has made use of NASA’s Astrophysics Data System.
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# Suppression of Magnetic State Decoherence Using Ultrafast Optical Pulses
## A Randomly-spaced pulses
The randomly spaced, radiative pulses act on this two-level system in a manner analogous to the way collisions modify atomic electronic-state coherence. In other word, the pulses do not affect the state populations, but do modify the coherence between the levels. The pulses can be treated in an impact approximation, such that during a collision, the time rate of change of density matrix elements resulting from the pulses is $`\dot{\rho }_{00}=\dot{\rho }_{11}=0`$ and
$$\dot{\rho }_{10}/\rho _{10}=\dot{\rho }_{01}/\rho _{01}=\mathrm{\Gamma }1e^{i\mathrm{\Delta }_s(t_i)\tau _p}=\mathrm{\Gamma },$$
(2)
where $`\mathrm{\Gamma }=T^1`$ is the average pulse rate and we have used the fact that the pulse area is a random number between 0 and 2$`\pi `$. Taking into account the collisional coupling $`V_c(t,b)`$ between the levels, one obtains evolution equations for components of the Bloch vector $`w=\rho _{11}\rho _{00}=2\rho _{11}1`$, $`v=i(\rho _{10}\rho _{01})`$ as
$$dw/dx=U(y)v;\text{ }dv/dx=U(y)wn(y)v,$$
(3)
where $`x=t/\tau _c(b)`$ is a dimensionless time, $`y=b/b_0`$ is a relative impact parameter, and $`U(y)=y^5`$ and $`n(y)=n_0y`$ are dimensionless frequencies. These equations are solved subject to the initial condition $`w(0)=1;`$ $`v(0)=0`$, to obtain the value $`\rho _{11}(x=1,y,n_0)=[w(x=1,y)+1]/2.`$ The relative transition rate $`S`$ is given by
$$S(n_0)=2\pi Nub_0^2_0^{\mathrm{}}y𝑑y\rho _{11}(x=1,y,n_0)/2,$$
(4)
where $`N`$ is the perturber density. A coefficient, $`R(n_0),`$ which measures the suppression of decoherence, can be defined as
$$R(n_0)=_0^{\mathrm{}}y𝑑y\rho _{11}(x=1,y,n_0)/_0^{\mathrm{}}y𝑑y\rho _{11}(x=1,y,0)$$
(5)
Solving Eqs. (3), one finds
$`\rho _{11}(x`$ $`=`$ $`1,y,n_0)=[1{\displaystyle \frac{r_1}{r_2r_1}}(e^{r_1}{\displaystyle \frac{r_1}{r_2}}e^{r_2})]/2;`$ (7)
$`r_{1,2}`$ $`=`$ $`\left(n_0y\pm \sqrt{(n_0y)^24y^{10}}\right)/2.`$ (8)
It is now an easy matter to numerically integrate Eqs. (5) to obtain $`R(n_0)`$. Before presenting the numerical results, we can look at some limiting cases which provide insight into the physical origin of the suppression of decoherence.
A plot of $`\rho _{11}(x=1,y,n_0)`$ as a function of $`y=b/b_0`$ is shown in Fig. 2 for several values of $`n_0.`$ With decreasing $`y`$, $`\rho _{11}`$ increases monotonically to some maximum value $`\rho _{11}(y_m)`$ and then begins to oscillate about $`\rho _{11}=1/2`$ with increasing amplitude. One concludes from such plots that two effects contribute to the suppression of coherence. The first effect, important for large $`n_0`$, is a reduction in the value of $`y_m`$. The second effect, important for $`n_0`$ of order unity, is a decrease in the value of $`\rho _{11}(y_m)`$. Let us examine these two effects separately.
For very large $`n_0,`$ $`n_0^{5/66}1,`$ one can approximate $`\rho _{11}`$ over the range of $`y`$ contributing significantly to the integral (4) as $`\rho _{11}(x=1,y,n_0)\left(1e^{y^{11}/n_0}\right)/2.`$ By evaluating the integrals in (5), one finds a suppression of decoherence ratio given by
$$R(n_0)=0.95/n_0^{2/11}.$$
(9)
The $`n_0^{2/11}`$ dependence is a general result for a collisional interaction that varies as the interatomic separation to the minus 6th power. It can be understood rather easily. The pulses break up the collision into $`n_0y`$ segments, each having a (dimensionless) time duration $`x_b=1/(n_0y)`$. Each segment provides a perturbative contribution to $`\rho _{11}`$ of order $`y^{10}(n_0y)^2`$, provided $`y<y_w`$, where $`y_w`$ is to be determined below. The total population from the entire collision interval varies as $`\rho _{11}y^{10}(n_0y)^2n_0y=y^{11}/n_0`$. Of course, $`\rho _{11}`$ cannot exceed unity. One can define an effective relative Weisskopf radius, $`y_w,`$ as one for which $`\rho _{11}=1`$, namely $`y_w=b_w/b_0=n_0^{1/11}.`$ The total transition rate varies as $`y_w^2n_0^{2/11}`$, in agreement with (9). As $`n_0\mathrm{}`$, the atom is frozen in its initial state.
For values of $`n_0`$ of order unity, the dominant cause of the suppression of decoherence is a decrease in the value of $`\rho _{11}(y_m)`$, rather than the relatively small decrease in $`y_m`$ from its value when $`n_0=0`$. For values $`n_03,`$ approximately 45% of the contribution to the transition rate $`S(n_0)`$ originates from $`y>y_m,`$ and, for these values of $`n_0`$, $`y_m\pi ^{1/5}`$ and $`\rho _{11}(y_m)(1+e^{n_0/2\pi ^{1/5}})/2`$. This allows us to estimate the suppression of decoherence ratio as $`R(n_0)=[0.55+.45(1+e^{n_0/2\pi ^{1/5}})/2],`$ such that $`R(1)=0.93`$, $`R(2)=0.88`$, $`R(3)=0.84.`$ These values are approximately 70% of the corresponding numerical results, indicating that the decrease in $`\rho _{11}(y_m)`$ accounts for approximately 70% of the suppression at low $`n_0`$, with the remaining 30% coming from a decrease in $`y_m`$. The first few collisions are relatively efficient in suppressing decoherence. With increasing $`n_0`$, the suppression process slows, varying as $`n_0^{2/11}`$. In Fig. 3, the suppression of decoherence ratio $`R(n_0)`$, obtained by a numerical solution of Eq. (5), is plotted as a function of $`n_0`$.
## B Uniformly Spaced Pulses
We consider now the case of equally spaced pulses, having effective pulse areas that are randomly chosen, modulo 2$`\pi `$. The time between pulses is $`T`$, and $`n_0=\tau _c(b_0)/T`$. For a relative impact parameter $`y=b/b_0`$, with $`m`$ $`n(y)=n_0y`$ $`m+1`$, where $`m`$ is a positive integer or zero, exactly $`m`$ or $`m+1`$ pulses occur. The effect of the pulses is calculated easily using the Bloch vector. At $`x=0`$, $`w=1`$ and $`v=0`$. The Bloch vector then undergoes free evolution at frequency $`U(y)=`$ $`y^5`$ up until the (dimensionless) time of the first pulse, $`x_s=t_s/\tau _c(b)`$. The pulse randomizes the phase of the Bloch vector, so that the average Bloch vector following the pulse is projected onto the $`w`$ axis. From $`x=x_s`$ to $`x_s`$ $`+T/\tau _c(b)=x_s`$ $`+1/n(y)`$, the Bloch vector again precesses freely and acquires a phase $`UT=y^5/n(y)=y^6/n_0`$, at which time the next pulse projects it back onto the $`w`$ axis. Taking into account the periods of free precession and projection, and averaging over the time $`x_s`$ at which the first pulse occurs, one finds
$`w(y)`$ $`=`$ $`[1n(y)]\mathrm{cos}[y^5]`$ (11)
$`+n(y){\displaystyle _0^1}𝑑x_s\mathrm{cos}[y^5x_s]\mathrm{cos}[y^5(1x_s)];`$
$`0`$ $``$ $`y1/n_0,`$ (12)
$`w(y)`$ $`=`$ $`[m+1n(y)][(m+1)/n(y)1]^1`$ (18)
$`\times {\displaystyle _{1m/n(y)}^{1/n(y)}}dx_s\mathrm{cos}[y^5x_s]\mathrm{cos}^{m1}[y^6/n_0]`$
$`\times \mathrm{cos}[y^5\{1x_s(m1)/n(y)\}]`$
$`+[n(y)m]\left[1m/n(y)\right]^1`$
$`\times {\displaystyle _0^{1m/n(y)}}dx_s\mathrm{cos}[y^5x_s]\mathrm{cos}^m[y^6/n_0]`$
$`\times \mathrm{cos}[y^5\{1x_sm/n(y)\}];`$
$`\text{ }m/n_0`$ $``$ $`y(m+1)/n_0\text{ }\text{ for }m1.`$ (19)
In the limit that $`n_01`$, for all impact parameters that contribute significantly to the transition rate, approximately $`n(y)`$ pulses occur at relative impact parameter $`y`$, implying that $`w(y)\mathrm{cos}^{n(y)}[y^5/n(y)]`$ and
$`R(n_0)`$ $`=`$ $`{\displaystyle \frac{1\mathrm{cos}^{n_0y}[y^6/n_0]}{1\mathrm{cos}[y^5]}}`$ (20)
$``$ $`{\displaystyle \frac{1[1y^{12}/2n_0^2]^{n_0y}}{1\mathrm{cos}[y^5]}}`$ (21)
$``$ $`{\displaystyle \frac{1e^{y^{11}/2n_0}}{1\mathrm{cos}[y^5]}}={\displaystyle \frac{0.84}{n_0^{2/11}}},`$ (22)
which is the same functional dependence found for the randomly spaced pulses. Note that the form {$`1\mathrm{cos}^{n(y)}[y^{15}/n(y)]\}`$ is identical to that found in theories of the Zeno effect .
The suppression of decoherence ratio $`R(n_0)`$, obtained from Eqs. (5) and (19) \[using $`\rho _{11}=(1+w)/2]`$, is plotted in Fig. 3. The fact that it lies below that for randomly spaced pulses is connected with the difference in the average collisional phase shift acquired between radiation pulses for the two models. The oscillations in $`R(n_0)`$ appear to be an artifact of our square pulse collision model. In the absence of the pulses, the first maximum in the transition cross section occurs for $`y_{\mathrm{max}}=(\pi )^{1/5}`$, corresponding to a $`\pi `$ collision pulse. With increasing $`n_0`$, the pulses divide the collision duration into approximately $`n(y)`$ equal intervals. If these pulse intervals are odd or even multiples of $`\pi ,`$ one can enhance or suppress the contribution to the transition rate at specific impact parameters. Numerical calculations carried out for a smooth interatomic potential do not exhibit these oscillations.
## C Discussion
Although the collisional interaction has been modeled as a square pulse, the qualitative nature of the results is unchanged for a more realistic collisional interaction, including level shifts. In fact, for a smooth interatomic potential that allows for an increased number of radiation pulses over the duration of the collisional interaction, the suppression is slightly enhanced from the square pulse values. Although the pulses are assumed to drive only the $`J=1,m=0`$ $`J=0,`$ excited state transition, it is necessary only that the incident pulses produce different phase shifts on the $`J=1,m=0`$ and $`J=1,m=1`$ state amplitudes.
To observe the suppression of decoherence, one could use Yb as the active atom and Xe perturbers. The Weisskopf radius for magnetic decoherence is about 1.0 nm , yielding a decoherence rate of $`10^{10}`$ s<sup>-1</sup> at 500 Torr of Xe pressure at 300$`{}_{}{}^{}C`$, and a collision duration $`\tau _c(b_0)`$ $`2.5`$ ps. Thus, by choosing a pulse train having pulses of duration $`\tau _p=`$100 fs, separated by 0.5 ps, it is possible to have 5 pulses per collision. If an experiment is carried out with an overall time of 100 ps (time from initial excitation to probing of the final state), one needs a train of about 200 pulses. To achieve a phase shift $`\mathrm{\Delta }_s\tau _p`$ of order $`2\pi `$ and maintain adiabaticity, one can take the detuning $`\delta =3\times `$10<sup>13</sup> s<sup>-1</sup> and the Rabi frequency $`\mathrm{\Omega }1\times `$10<sup>14</sup> s<sup>-1</sup> on the $`J=1,m=0`$ $`J=0,`$ excited state transition . The corresponding, power density is $`1.5\times 10^{11}`$ W/cm<sup>2</sup>, and the power per pulse is $`150`$ $`\mu `$J (assuming a $`1`$ mm<sup>2</sup> focal spot size). This is a rather modest power requirement. With 5 pulses/collision duration, one can expect a relative suppression of magnetic state decoherence of order 40%.
Finally, we should like to comment on whether or not the effect described in this work constitutes a Quantum Zeno effect. Normally, the Quantum Zeno effect is presented as a projection of a quantum system onto a given state as a result of a measurement on the system. In the experiment of Itano et al., this ”measurement” is reflected by the presence or absence of spontaneously emitted radiation during each uv ”measurement” pulse. The measurement pulse must be sufficiently long to produce a high likelihood of spontaneous emission whenever the atom is ”projected” into the initial state by the pulse. Following each measurement pulse, the off-diagonal density matrix element for the two states of the rf transition goes to zero. In our experiment involving off-resonant pulses, the number of Rayleigh photons scattered from the $`J=0`$ level during each applied pulse is much less than unity. As such, there is no Quantum Zeno effect, even if suppression of magnetic state decoherence occurs. On average, each pulse having random area destroys the coherence between the $`J=1,m=0`$ and $`J=1,m=\pm 1`$ state amplitudes, but does not kill this coherence for a single atom. With an increasing number of radiation pulses, $`n_0`$, however, both the average value and the standard deviation of the transition probability tends to zero as $`n_0^1`$ for each atom in the ensemble.
The experiment of Itano et al. could be modified to allow for a comparison with the theory presented herein, and to observe the transition into the Quantum Zeno regime. If the pulses that drive the strong transition are replaced by a sequence of off-resonant pulses, each pulse having a duration $`\tau _p`$ much less than the time, $`T_\pi `$, required for the pi pulse to drive the weak transition, and each pulse having an effective area, $`\mathrm{\Delta }_s\tau _p=(\mathrm{\Omega }^2/4\delta )\tau _p`$, that is random in the domain \[0,2$`\pi ],`$ then the pulses will suppress the excitation of the weak transition (it is assumed that $`\mathrm{\Omega }/\delta 1)`$. If the upper state decay rate is $`\gamma `$, then the average number of Rayleigh photons scattered during each pulse is $`n=\left(\mathrm{\Omega }/4\delta \right)^2\gamma \tau _p.`$ For $`n<1`$, there is suppression of the transition rate as in our case, while, for $`n1`$, there is suppression and a Quantum Zeno effect. There is no average over impact parameter, since exactly $`[T_\pi /T`$\] or (\[$`T_\pi /T`$\]+1) pulses in each interval between the pulses, where \[$`x`$\] indicates the integer part of $`x`$.
## D Acknowledgments
PRB is pleased to acknowledge helpful discussions with R. Merlin, A. Rojo and J. Thomas. This research is supported by the National Science Foundation under grant PHY-9800981 and by the U. S. Army Research Office under grants DAAG55-97-0113 and DAAH04-96-0160.
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# Analysis of the experimental results from the EAS installation GAMMA at Mt. Aragats
## 1 <sup>1</sup><sup>1</sup>1Corresponding author. E-mail: martir@lx2.yerphi.amIntroduction
The particle physics and astrophysics aspects of air shower studies are closely interlaced. The problems of sources and nature of the primary radiation and features of the high energy hadron interactions cannot be solved separately. The history of attempts to clarify the origin of the appreciable change in the index of the shower size spectrum (the $`knee)`$, which was found more than 40 years ago and confirmed by other experiments, is the best manifestation of that. Many new experiments were born devoted to the observation of various components of EAS, including Cherenkov light in the atmosphere and experiments under water. However, up to now there is no unequivocal explanation of the reason of this phenomena. The most popular interpretation of the $`knee`$ is the existence of the $`knee`$ in the primary cosmic ray energy spectrum at energy about (3-5)<sup>.</sup>10$`^\text{6}`$ GeV and, as a possible consequence, the change of the mass composition. There are different experimental analysis leading to opposite conclusions about primary composition after $`knee`$ (lighter , normal or heavier ).
At the same time there is an alternative explanation of the $`knee`$, also based on some inconsistencies in the EAS data, connected with a modification in the hadronic interaction properties to the change of the number and energy spectrum of the secondaries. The complex installations, which measure electromagnetic component as well as muon and hadron components of EAS have the best opportunities to explore the nature of the $`knee.`$ HADRON-2 on Tien-Shan and KASCADE are installations of such a kind. The hadron component data obtained from these installations show the fast absorption of the high energy hadrons, which is impossible to be explained by the energy increase of the interaction cross-section only. At the same time, the electromagnetic component data are in an agreement with the QGSJET model and correspond to the mixed composition of primary with progressive changing of the mass composition. In the work on the bases of the Tien-Shan data analysis it is supposed that there is no $`knee`$ in the primary spectrum and $`knee`$ in the EAS size spectrum is explained by the change of the hadronic interaction in the atmosphere. This explanation is based on the results of the EAS electromagnetic component study , according to which there is no $`knee`$ neither in the $`young`$ showers spectrum (showers with $`S<0.75,`$ generated by primary protons interacted in the depth of atmosphere), nor in the $`old`$ showers ($`S>1.05`$) one. And the existence of the $`knee`$ in the same spectra of other experiments could be explained by the missing of the young showers.
On the other hand it is not so clear situation with the so-called $`reverse`$ $`knee`$ in the EAS size spectrum at $`N_e(23)10^7`$, which was discovered in . Subsequently the same irregularity was obtained once again in Tien-Shan experiment .
At present, practically all experiments on research of the cosmic rays in the energy range 10<sup>14</sup>-10<sup>17</sup> eV are pointed to investigate the nature of the $`knee`$. We have to point out that the KASCADE , CASA-MIA and EAS-TOP experiments present the basic advantage to have their data analyzed using multiparameter procedures. Only such an approach to study the EAS characteristics allows some advance in understanding of the reasons of the existence of the $`knee`$ in the spectrum on N<sub>e</sub>. In this paper, we present the experimental results of the GAMMA installation of the experimental complex ANI also working in this field.
The GAMMA installation is located on hillside of the Mt. Aragats in Armenia (3200 m a.s.l.). In comparison with the other large scale experiments with the same goals, the high-altitude observation level of the GAMMA installation provides some advantages, in particular for minimizing the intrinsic fluctuations of the observables due to the stochastic character of the EAS development in atmosphere. This will add some new observations, which would clearly complement the data from KASCADE (near sea-level), CASA-MIA (1200 meters), and EAS-TOP (2000 meters). In the present work, showers are simulated using the CORSIKA code 5.20 in which the hadronic interactions are described by the VENUS model,. On the other hand, to compare the simulation data with experimental ones, the normal mixed composition was used, i.e. proton-40 $`\%`$, $`\alpha `$-21$`\%`$, light-nuclei ($`<`$A$`>`$=14)-14$`\%`$, medium-nuclei ($`<`$A$`>`$=26)-13$`\%`$ and heavy-nuclei ($`<`$A$`>`$=56)-12$`\%`$
## 2 Present status of the GAMMA EAS installation
The GAMMA installation was realized as a part of the project ANI in an attempt to continue experimental study of the hadronic interactions and the energy spectrum and mass composition of the primary cosmic radiation in the energy range of $`10^{14}10^{17}`$ eV.
After some years spent to enlarge the effective area of the muon underground detector, which was necessary to elaborate methodical studies of the detector parameters, to investigate carefully the array response and to determine the precision of the shower parameter estimation, the GAMMA experiment is, now, effectively running.
GAMMA is a central type array and consists of two main parts $`(figure\mathrm{\hspace{0.17em}1})`$:
(i) the surface part, for the registration of the EAS electromagnetic component;
(ii) the large muon underground detector, to register the EAS muon component.
#### 2.0.1 The surface scintillation array
The surface scintillation array constists of 25 groups of 3 scintillation detectors placed on concentric circles with radii of 17, 28, 50 and 70 $`m`$. Each of the detectors has 1 $`m^2`$ of effective area. They are distributed on the full area of $``$1.5$`10^4`$ $`m^2`$. An additional station with 20 detectors of the same type is placed at the 135 $`m`$ from the installation centre. The energy threshold for the registration of the electrons flux, $`E_e^{thres}9.5`$ MeV, is estimated taking into account the thickness and the specificities of the detectors and registration stations. The calibration of detectors has been carried out by periodical (each hour) registration of the cosmic ray background particles . It was taken into account that the most probable value of the registered particles differs from $`0^0`$ and depends on the observation level .
Each of the 25 registration stations is equipped with the timing channel which allows to determine the angular coordinates of the shower axis. The trigger condition for one EAS registration is that in each of four groups of detectors placed at 17 meters from the array centre, two detectors have to register a flux density with $`\rho 3part/m^2`$.
The main EAS parameters: coordinates and angles of the shower axis $`(X,Y,\theta ,\phi )`$, number of electrons $`(N_e)`$ and age parameter $`S_{NKG}`$ are obtained using the SPACE code created on the basis of the statistical methods for the solution of the inverse problem . Taking into account the detector response, GAMMA array geometry and using Monte-Carlo procedure the data bank of pseudo experimental showers was created and treated by the same SPACE code.
As the result of this study, the following accuracies on the EAS parameter estimations are obtained: $`\mathrm{\Delta }X,\mathrm{\Delta }Y3m`$ $`(`$for $`R40m)`$, $`\sigma _\theta 1.5^0`$, $`\sigma _\phi 8^0`$, $`\sigma _{N_e}/<N_e>20\%`$.
Furthermore, the analysis of the shower registration efficiency shows that the showers with particle number $`N_e310^5`$ within 20 $`m`$ and $`N_e10^6`$ within 40 $`m`$ from the array centre are registered with 100% efficiency.
#### 2.0.2 Underground muon detector
The layout of the muon underground detector is shown in $`figure\mathrm{\hspace{0.17em}2}`$. As for the GAMMA surface part identical scintillation detectors with $`S=1m^2`$ were used for the registration of the EAS muon component. It should be noted that the muon scintillation detectors are divided into 2 groups placed under different absorber thickness. This allows to study EAS muons with 2 different energy thresholds: $`E_\mu 5`$GeV (Hall) and $`E_\mu 2.5`$GeV (Tunnel). These thresholds were estimated experimentally using scintillation telescope with a small solid angle. The arrangement of the muon detectors gives the possibility to determine the muon lateral distribution up to $`60m`$ at $`E_\mu `$ 5 GeV and up to $`90m`$ at $`E_\mu `$2.5 GeV.
Some peculiarities of the GAMMA array geometry make necessary to pay attention on the azimuth angle symmetry of data. The surface detectors are placed on the hillside and sometimes the difference between $`Z`$ coordinates of the individual detectors is reaching $`18m`$. On the other hand, the position of muon detector is very asymmetric in regard to the geometric centre of our installation. In this case, the influence of azimuth angle $`\phi `$ could be strong. To check this effect we divided the observed showers into four groups with $`\phi [090]^0,[90180]^0,[180270]^0,[270360]^0`$. The comparison of the average muon lateral distributions for each of these groups by $`\phi `$ is shown in $`figure\mathrm{\hspace{0.17em}3}(a,b)`$ for $`<N_e>=1.2810^6`$. It can be visible, that the difference of the muon lateral distributions for the different ranges of $`\phi `$ is quite negligible. Furthermore, it is necessary to point out that for the small distances (less than $`15m`$ for $`E_\mu `$2.5 GeV and $`8m`$ for $`E_\mu `$5 GeV) there is a significant influence of the punch through of particles (high energy electrons and hadrons) and these distances are excluded from the further analysis.
## 3 Results
We have used the experimental data obtained during 3300 hours operation time. The number of EAS with $`N_e10^5`$ and with zenith angle $`\theta 30^0`$ is $`260.000`$. The effective area for the selection of the EAS axis is $`5.000m^2`$.
### 3.1 Electromagnetic component characteristics
The main parameters of the EAS electromagnetic component at the given observation level are: the total number of the charge particles or shower size $`N_e`$ and the lateral shower age parameter $`S`$. In order to obtain the EAS size spectrum and to investigate its behaviour in the $`knee`$ region the correct estimation of these parameters is needed.
Usually in the energy region of $`10^510^7`$GeVthe scintillation detectors are used to register the EAS charged particles. However it is well known that such type of a detector registers not only the EAS charged particles, but also electrons generated in the absorber above as well as in the scintillator by EAS photons. This contribution to the measured charged particle density depends both on the scintillator thickness and the distance of detector from the shower axis.
The comparison of the charge particle densities seen by scintillation detectors and Geiger counters has been made for different distances from the shower core . In the same way, comparisons with other kind of detectors have been performed, (spark chamber , Geiger counter , neon flash tube , as well as with thin and thick scintillators ). But, up to now, there are not exact estimation of the photon contribution. In the energy range $`10^510^7`$GeV all these investigations show both a smooth rise and a noticeable decrease of the photon contribution to the measured shower size versus the distance to the shower core up to $`100m`$.
Recently the methodical experiment on the MAKET-ANI installation on the Mt. Aragats has been carried out . Comparing the densities measured by scintillators with different thicknesses ($`1.0,\mathrm{\hspace{0.17em}1.5}`$ and $`5.0cm`$), up to $`100m`$, the ratio $`K_\gamma (r)`$ of the measured density $`\rho _{sc}(r)`$to the charged one $`\rho _{ch}(r)`$is described by $`K_\gamma (r)=\rho _{sc}(r)/\rho _{ch}(r)=(r/R_M)^{0.18}`$, where $`r`$ is the distance from the shower core and $`R_M=120m`$ is the Moliere radius for the observation level of $`700g.cm^2`$. According to this result the photon contribution to the charged particle density is $`17\%`$ at $`50m`$ from the shower axis and practically disappears at $`100m`$.
At the same time, simulations show a significant increase of the ratio of the photon number to electrons one versus $`r`$. Calculations for the real conditions of the GAMMA installation have been performed taking into account the absorber and scintillator thicknesses . According to this work the photon contribution to the measured number of the charged particles is practically constant from $`5`$ to $`100m`$ and is $`K_\gamma 1.25`$. The processing of the GAMMA data considering these two results, $`K_\gamma (r)=(r/R_M)^{0.18}`$ and $`K_\gamma =1.25`$, , has been made to study the influence of this factor on the measured EAS features.
In $`figures`$ $`4`$ the experimental electron lateral distributions are shown for three EAS average size intervals $`<N_e>`$=$`2.310^5,`$ $`7.310^5`$ and $`22.810^5`$ as estimated for the two cases of the photon contribution. It can be seen that in both cases the electron lateral distribution can be very well approximated using the NKG function, but because of the different shape of $`K_\gamma `$ versus $`r`$, the average age parameter value $`<S>`$ is different by about 0.16.
The comparison between the measured and simulated electron lateral distributions is shown in $`figure\mathrm{\hspace{0.17em}5}`$ for the same values of $`<N_e>`$. It can be seen a good agreement for the correction factor $`K_\gamma =1.25=const`$. This agreement improves with the increasing of $`<N_e>`$.
It has to be noticed that the age parameter only describes the shape of the charged particles lateral distribution and its values $`<S>`$ $`<1`$ do not contradict to the fact that the shower is over the longitudinal development maximum. $`Figure`$ $`6`$ shows the dependence of the age parameter $`<S>`$ versus the shower size $`<N_e>`$ from our experimental data processed with $`K_\gamma (r)=(r/R_M)^{0.18}`$ and $`K_\gamma =1.25`$ and the results from our simulation and the Tien-Shan experimental data . For $`<N_e>`$ $``$ $`(35)10^5`$, experimental data with $`K_\gamma =1.25`$ are in a good agreement with the simulation results which is not at all the same in the case of $`K_\gamma (r)=(r/R_M)^{0.18}`$ for which the disagreement is large.
For the both variants of processing, the differential size spectra are shown in $`figure`$ $`7a`$ for showers in the zenith angle interval $`\theta 30^0`$, (mean atmospheric depth : $`738g.cm^2`$). In both cases the $`knee`$ can be easily observed. Spectra are approximated in the whole $`N_e`$ interval by formula : $`I(N_e)=AN_e^{\gamma _1}(1+(N_e/N_e^{knee})^k)^{(\gamma _1\gamma _2)/k}`$. The spectra differ mainly at the beginning and the end of the size interval, (by no more than 20%). Spectral indexes difference for GAMMA for both cases is $`\mathrm{\Delta }\gamma =0.36\pm 0.02`$. The $`knee`$ regions approximately coincide at $`N_e^{knee}(1.82.0)10^6`$. The integral intensity for the range of $`N_eN_e^{knee}`$ is $`(1.3\pm 0.4)\mathrm{\hspace{0.17em}\hspace{0.17em}10}^7m^2s^1sr^1`$. This result is in agreement with the data of other experiments .
In order to demonstrate the $`knee`$ region properties more visually spectra are given in $`figure`$ $`7b`$ multiplied by $`N_e^{2.5}`$. For the case of $`K_\gamma =1.25`$ spectral index is less by $`(0.07\pm 0.02)`$ in the whole $`N_e`$ region.
In the $`figure`$ $`8,`$ the size spectra for the two zenith angle intervals are shown at $`<\theta >=14^0`$ and $`<\theta >=31^0`$ to make the qualitative assessment of the behaviour of the spectra at different angles. They are practically parallel and coincide when data with $`<\theta >=31^0`$ are shifted. Such behaviour corresponds to the shower absorption lenght of $`\mathrm{\Lambda }=(230\pm 25)gcm^2`$ below the knee and $`(230\pm 40)gcm^2`$ above up to $`Ne=510^6`$. The permanency of the attenuation length is one evidence of the constant charged particle composition of the showers below and above the $`knee`$ only and says nothing about the cause of the $`knee`$.
In $`figure`$ $`9`$ the spectra are shown with various cuts of the age parameter $`S`$. Relatively $`old`$ shower spectrum $`(S>0.85)`$ practically has no $`knee`$ and is steeper than all shower spectrum. On the other hand, the $`young`$ shower spectrum $`(S<0.75)`$ has the obvious $`knee`$ and is flatter, than the all shower spectrum below the $`knee`$ and is almost parallel to it above. These spectra may be explained with both hypothesises of the $`knee`$ origin noted in the introduction, because of the region of $`N_e(12)10^6`$ is transitional. The $`young`$ showers fraction does not arise above the $`knee`$ because of change of the mass composition to heavier one in the model of primary $`knee`$. In the case of change of the hadron interaction the fraction of the heavy-like events will rise among the showers generated by primary protons.
### 3.2 Muon component characteristics
As it was mentioned above, the underground muon detectors are placed in the hall and tunnel of the GAMMA underground part. Due to the different shieldings of concrete, iron and ground above detectors, muons with two energy thresholds $`E_\mu 5`$GeV (hall) and $`E_\mu 2.5`$GeV (tunnel) are detected. The muon lateral distributions for given thresholds and three $`<N_e>`$ intervals are shown in $`figure\mathrm{\hspace{0.17em}10}a`$. These distributions are limited by the number of detectors and their disposition in the hall and tunnel. Moreover, close to the shower core the signal in the detectors are strongly affected by the contribution of other type of particles due to the punch through effect. So, the correct determination of the muon lateral distribution is available at distances $`[8`$ $``$ $`52]m`$ for $`E_\mu 5`$GeV and $`[2090]m`$ for $`E_\mu 2.5`$GeV. In the interval of $`[050]m`$, the shapes of lateral distribution are very similar for the both muon energy thresholds. This gives the possibility to find a transition coefficient from the density of muons with $`E_\mu 2.5`$GeV to the muon density with $`E_\mu 5`$GeV, which is $`\rho (2.5`$GeV$`)/\rho (5`$GeV$`)1.35`$, for all $`<N_e`$ $`>`$ intervals. In $`figure`$ $`10b`$, the muon lateral distribution with the energy threshold $`E_\mu 5`$GeV for the hall detectors and reconstructed one from the distribution of tunnel detectors are presented. These distributions are well approximated by the Hillas function, , $`\rho _\mu (N_e,r)=0.9(N_e/10^5)r^{(0.75+0.06\mathrm{ln}(N_e/10^5)}\mathrm{exp}(r/80).`$ This formula describes well the experimental data for $`N_e<2.10^6`$. Farther, the dependence of $`\rho _\mu `$ versus $`N_e`$becomes stronger but the shape of the distribution remains the same.
The approximations of the Tien-Shan data , for the same energy threshold $`E_\mu 5`$GeV and the same observation level $`700gcm^2`$ are presented by $`\rho _\mu (N_e,r)=0.95(N_e/10^5)^{0.8}r^{0.75}\mathrm{exp}(r/80)`$. This formula describes the GAMMA experimental data noticeably worse, especially for $`N_e>10^6`$.
As it was mentioned above, the lateral distribution of muons in our experiment,$`(E_\mu 5`$GeV), is studied for distances $`r52m`$ from the shower core where about $`33\%`$ of all muons are contained. Obtaining the total muon number, $`N_\mu ,`$by integration the muon lateral distribution function, the possible extrapolation errors for $`r>52m`$ will give an error in the $`N_\mu `$ estimation. On the other hand, closer to the shower axis, $`r8m,`$ other kinds of particles give a considerable contribution to muon lateral distribution. For this reason, we use the truncated muon number, $`N_\mu ^{trun}`$, which is the number of muons in the ring $`8m<r<52m`$. Experimentally, the truncated muon number is determined as $`N_\mu ^{trun}(r)=(\rho _i^{ex}/w_\mu ^{trun}(r))/K`$ , where:
$`K`$ is number of detectors in the given interval of $`r`$;
$`\rho _i^{ex}`$ is the experimental muon density in the $`i^{th}`$ detector and
$`w_\mu ^{trun}(r)`$ is the truncated muon probability distribution, determined by our muon lateral distribution approximation.
It should be noticed, that according to simulations, $`N_\mu ^{trun}`$($`>`$5 GeV) in our experiment depends on the primary particle mass. $`Figure`$ $`11`$ shows $`N_\mu ^{trun}`$ dependence on $`N_e`$. Up to $`N_e210^6`$ it can be fitted with the expression $`N_\mu ^{trun}N_e^{0.79}`$. For larger $`N_e`$ dependence becomes steeper. It should be noted that because of the shower size threshold $`N_e>10^5`$ the registration of the shower with fixed muon size and 100% efficiency is possible at $`N_\mu ^{trun}10^3`$.
The differential $`N_\mu ^{trun}`$ spectrum is plotted on $`figure`$ $`12`$. It can be seen that there is no obvious $`knee`$, but the data do not contradict to change of spectral index by (0.1-0.2) at $`N_\mu ^{trun}>510^3`$. Contrary to the case of fixed $`N_e`$ the showers at fixed $`N_\mu `$ are enriched by the showers generated by heavy primaries. If the proton spectrum has a $`knee`$ only then the $`N_\mu `$ spectrum will have feebly marked $`knee`$. At energies above the possible $`knee`$ position of the primary iron spectrum, $`N_\mu `$ spectrum will obtain the final index different by $`\mathrm{\Delta }\gamma 0.5`$ from present one. It is necessary to expand our measurements to this energy region.
Situation in the case of the change of hadron interactions for sufficientely high energies depends on details of the secondary spectra but the difference of indexes will be not bigger.
In order to interpret our experimental results the instrumental response has to be taken into account. For this purpose the detector simulation program ARES (ARagats Events Simulation) based on the GEANT package. has been developed for the GAMMA installation. The data from CORSIKA simulated EAS are used as input for the ARES code.
To take into account the detector response to EAS muon component all secondary particles of EAS at the ground level are passed through the absorber and muon detectors. For each shower the deposited energy and the muon number in the individual detector are obtained (the methodical procedure of muon number estimation is described in ). The muon lateral distribution are derived using the CORSIKA simulation data in the primary energy range of 3\*10<sup>5</sup>-10<sup>7</sup>GeV for the different primary groups. $`Figure`$ $`13(a,b)`$ shows the muon lateral distributions in the shower size range 1.78$``$10<sup>5</sup>$`<`$N<sub>e</sub>$`<`$3.16$``$10<sup>5</sup> in the case of the normal mixed composition for the tunnel and hall detectors. The experimental distribution and the CORSIKA/ARES simulation results (with the corresponding muon thresholds) are in a good agreement.
## 4 Conclusions
The study of the EAS electrons and muons $`(E_\mu >5`$GeV and $`E_\mu >2.5`$GeV$`)`$ by the GAMMA array gives reliable information about EAS characteristics in the shower size range of 10<sup>5</sup>-10<sup>7</sup>. There are not noticeable contradictions of our data with other experiments. Using the contribution of the EAS gamma-quanta equal $`K_\gamma =1.25=const`$ an agreement between CORSIKA simulation and experimental data is obtained.
Permanency of the age parameter $`S`$ at $`N_e>210^6`$ and steepening of $`N_\mu ^{trun}`$ depending on $`N_e`$ at the same range by $`N_e`$ conform to the more rapid shower development. The data in the shower size interval $`N_e^{knee}÷10N_e^{knee}`$ are not decisive for the problem of the $`knee`$ origin and enlargement of the study range is necessary for the complex installations to make sure conclusion about the primary composition and spectra. With the object to extend a possibilities of the GAMMA installation it is expected to decrease the registration threshold by $`N_e`$ and to enlarge an effective area of the shower selection with $`N_e>10^6.`$
The present paper is based on the ANI collaboration data bank and express the point of view of the given group of co-authors.
### 4.1 Acknowledgments
The present investigations are embedded in a collaboration between the Moscow Lebedev Institute (Russia), the Yerevan Physics Institute (Armenia) and the Universite de Bordeaux 1 (France).
We give thanks to the Staff of the Aragats Research Station for the assistance during the longterm maintenance of the GAMMA installation.
We are grateful to all colleagues of the Moscow Lebedev Institute and Yerevan Physics Institute Cosmic Ray Divisions who were taking part in the development and creation of the GAMMA installation. We would like to express our special gratitude to prof. S.I. Nikolsky for very useful comments and N.M. Nikolskaya for the creation of the software.
We thank also prof. A.A. Chilingarian, prof. S. Ter-Antonyan and Dr. E.Mnatsakanyan for discussions and we do not forget P. Aguer from the CENBG, (CNRS-In2p3-France) whose help has been useful.
This work was supported by the grant 96-752 of the Armenian Ministry of Industry, by the Russian RFBR 96-02-18098 grant and by the Russian Energy Ministry.
## 5 Figures
Figure 1. Layout of the GAMMA installation
Figure 2. Schematic view of the muon underground detector
Figure 3. Muon lateral distributions at the different azimuth angles $`\phi `$:
a) hall, b) tunnel
Figure 4. Electron lateral distributions for two cases of the photon contribution
a) $`K_\gamma (r)=(r/R_M)^{0.18}`$, b) $`K_\gamma (r)=1.25=const`$
Figure 5. Electron lateral distribution in comparison with the CORSIKA simulation for two cases of the photon contribution
Figure 6. The average age parameter $`<S>`$ versus shower size $`<N_e>`$ at two cases of the photon contribution in comparison with CORSIKA simulation and Tien-Shan data
Figure 7. Differential size spectra at two cases of the photon contribution
Figure 8. Differential size spectra at different zenith angles
Figure 9. Differential size spectra at different interval by age parameter
Figure 10. Muon lateral distribution at different shower sizes in
a) hall, b) tunnel
Figure 11. Average truncated muon number $`<N_\mu ^{tun}>`$ versus number of electrons $`<N_e>`$
Figure 12. Differential spectrum of truncated muon number
Figure 13. Muon lateral distributions in comparison with CORSIKA and ARES simulation results:
a) hall, b) tunnel
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# Borel–Padé vs Borel–Weniger method: a QED and a QCD example
## I Introduction
In this letter we want to present some numerical results which allow us to compare the efficiency of the Borel–Padé method with that of the Borel–Weniger method for resummation of truncated perturbation series (TPS) in some physically significant scenarios. The scenarios we are referring to are those when the function, which we want to find through a resummation, is known to have certain singularity structure in the Borel plane. If there are singularities on the positive axis of the Borel plane, then we implicitly assume that in such cases we either know the correct prescription for integration in the Laplace–Borel integral, or we simply adhere to a certain adopted prescription.
(1) We will first illustrate the efficiency of the two methods on the QED example of the Euler–Heisenberg Lagrangian density, i.e., the one–loop fermion–induced effective action density in a strong uniform electromagnetic field . In this case, the solution is known, and its real part can be written in the following form:
$`\mathrm{Re}\delta \stackrel{~}{}(\stackrel{~}{a};p)`$ $`=`$ $`\mathrm{Re}{\displaystyle _0^{\mathrm{}}}𝑑w\mathrm{exp}\left({\displaystyle \frac{w}{\stackrel{~}{a}}}\right){\displaystyle \frac{(1)}{w}}\left[{\displaystyle \frac{p\mathrm{cos}(w)}{\mathrm{sin}(w+\mathrm{i}ϵ^{})}}\mathrm{coth}(pw)+{\displaystyle \frac{1}{3}}(1p^2){\displaystyle \frac{1}{w^2}}\right],`$ (1)
where we use notations
$$\stackrel{~}{a}\frac{ga}{m^2},\stackrel{~}{b}\frac{gb}{m^2},p\frac{b}{a}\frac{\stackrel{~}{b}}{\stackrel{~}{a}},\delta \stackrel{~}{}\delta /\left(\frac{m^4\stackrel{~}{a}^2}{8\pi ^2}\right).$$
(2)
Here, $`\delta `$ is the actual Lagrangian density induced by the one–loop fluctuations of the fermions in the field; $`g`$ is the field–to–fermion coupling parameter (in QED it is the positron charge $`e_0`$); $`m`$ is the mass of the fermion (electron); $`a`$ and $`b`$ are Lorentz–invariant expressions characterizing the electric and the magnetic fields $`\stackrel{}{}`$ and $`\stackrel{}{}`$, respectively
$$\left(\genfrac{}{}{0pt}{}{a}{b}\right)=\left[\pm \stackrel{}{}^2\stackrel{}{}^2+\sqrt{\left(\stackrel{}{}^2\stackrel{}{}^2\right)^2+4\left(\stackrel{}{}\stackrel{}{}\right)^2}\right]^{1/2}/\sqrt{2}.$$
(3)
Expression (1) can be obtained, for example, directly by integrating out the fermionic degreees of freedom in the path integral expression of the full effective action, then employing the proper–time integral representation for the difference of logarithms, evaluating the traces in the integrand, and subsequently performing Wick rotation by $`\pi /4`$ in the plane of the proper–time $`s`$: $`ags\mathrm{i}w+ϵ^{}`$. We refer to for more details on the latter point. The perturbative expansion of the full solution (1), in powers of $`\stackrel{~}{a}`$, is
$$\delta \stackrel{~}{}^{\mathrm{pert}.}(\stackrel{~}{a};p)=\left[c_1(p)1!\stackrel{~}{a}^2+c_3(p)3!\stackrel{~}{a}^4+c_5(p)5!\stackrel{~}{a}^6+\mathrm{}\right],$$
(4)
with coefficients
$`c_1(p)`$ $`=`$ $`{\displaystyle \frac{1}{45}}\left[(1p^2)^2+7p^2\right],c_3(p)={\displaystyle \frac{1}{945}}\left[2(1p^2)^3+13p^2(1p^2)\right],\mathrm{etc}.`$ (5)
In the case of the pure magnetic field (p.m.f.), the corresponding expressions are simpler
$`\delta \stackrel{~}{}(\stackrel{~}{b})_{\stackrel{~}{a}=0}{\displaystyle \frac{8\pi ^2\delta _{a=0}}{m^4\stackrel{~}{b}^2}}={\displaystyle _0^{\mathrm{}}}𝑑w\mathrm{exp}\left({\displaystyle \frac{w}{\stackrel{~}{b}}}\right){\displaystyle \frac{(1)}{w}}\left[{\displaystyle \frac{\mathrm{coth}(w)}{w}}{\displaystyle \frac{1}{3}}{\displaystyle \frac{1}{w^2}}\right],`$ (6)
$$\delta \stackrel{~}{}^{\mathrm{pert}.}(\stackrel{~}{b})_{\stackrel{~}{a}=0}=\left[\stackrel{~}{c}_11!\stackrel{~}{b}^2+\stackrel{~}{c}_33!\stackrel{~}{b}^4+\mathrm{}\right],\stackrel{~}{c}_1=\frac{1}{45},\stackrel{~}{c}_3=\frac{2}{945},\mathrm{}$$
(7)
We can now use (4)–(5), and (7), as a laboratory for resummation methods, since the full (resummed) solutions (1) and (6) are known. Since (1) and (6) are Laplace–Borel integrals, it is natural to use these examples for testing combined resummation techniques which involve Borel transformation. Borel transform $`B_L`$ of series (4) is
$$B_L(w;p)=c_1(p)w+c_3(p)w^3+c_5(p)w^5+\mathrm{},$$
(8)
and analogously for (7). In Ref. , we used Borel–Padé technique for resummation, i.e., we applied various Padé approximants $`[N/M]_B(w;p)`$ to (8)<sup>3</sup><sup>3</sup>3 $`[N/M]_B(w;p)`$, being ratio of polynomials in $`w`$ of powers $`N`$ and $`M`$, respectively , is based solely on the truncated perturbation series (TPS) of (8) involving only terms with $`c_n`$: $`nN+M`$. and then employed the Laplace–Borel integral to obtain the resummed value
$$BP^{[\mathrm{N}/\mathrm{M}]}\left[\delta \stackrel{~}{}^{\mathrm{pert}.}\right](\stackrel{~}{a};p)=_0^{\mathrm{}}𝑑w\mathrm{exp}\left(\frac{w}{\stackrel{~}{a}}\right)[N/M]_\mathrm{B}(w;p).$$
(9)
The integration over poles in (9) was carried out according to the Cauchy principal value prescription, since the full solution (1) requires it.<sup>4</sup><sup>4</sup>4 Various QCD and QED applications of the Borel–Padé approach with the principal value prescription can be found in . The novel method of Ref. is, in addition, well suited for obtaining the imaginary part of $`\delta `$.
Recently, the authors of proposed the use of Weniger (delta sequence) transformations as an alternative to the use of Padé approximants, for direct resummation of truncated perturbation series. For a truncated perturbation series (TPS) of the form
$$F_{[n+1]}(z)=\underset{0}{\overset{n+1}{}}\gamma _jz^j$$
(10)
it is defined as
$$\delta _n^{(0)}(\zeta ;\gamma _0,\mathrm{},\gamma _{n+1})=\frac{{\displaystyle \underset{j}{\overset{n}{}}}(1)^j\left({\displaystyle \genfrac{}{}{0pt}{}{n}{j}}\right){\displaystyle \frac{(\zeta +j)_{n1}}{(\zeta +n)_{n1}}}{\displaystyle \frac{z^{nj}F_{[j]}(z)}{\gamma _{j+1}}}}{{\displaystyle \underset{j}{\overset{n}{}}}(1)^j\left({\displaystyle \genfrac{}{}{0pt}{}{n}{j}}\right){\displaystyle \frac{(\zeta +j)_{n1}}{(\zeta +n)_{n1}}}{\displaystyle \frac{z^{nj}}{\gamma _{j+1}}}},$$
(11)
where $`(\zeta +j)_{n1}\mathrm{\Gamma }(\zeta +j+n1)/\mathrm{\Gamma }(\zeta +j)`$ are the Pochhammer symbols and $`\zeta =1`$ is usually taken. The approximant (11) is a ratio of two polynomials in $`z`$ of power $`n`$ each, and when expanded back in powers of $`z`$ it reproduces all the terms of $`F_{[n+1]}`$.
The authors applied (11) directly to the TPS’s of $`\delta \stackrel{~}{}^{\mathrm{pert}.}(\stackrel{~}{a};p)/\stackrel{~}{a}^2`$ of (4), and when re-expanding the approximant in powers of $`\stackrel{~}{a}`$ they were able to predict the next coefficient in the series with a better precision than the one provided by the corresponding diagonal (or almost diagonal) Padé approximant. Further, in the case of the pure magnetic field they showed that the method (11), when applied directly to the TPS’s in $`\stackrel{~}{b}`$ of the induced Lagrangian density,<sup>5</sup><sup>5</sup>5 The approximants (11) applied to the TPS’s of the series (7) divided by $`\stackrel{~}{b}^2`$. gave better results of resummation than the corresponding Padé approximants.
We now combine the method (11) with the Borel transformation (4) $``$ (8), and compare the results of resummation obtained in this way with the results of the corresponding Borel–Padé approximants of Ref. . Formula (11) is applied to the Borel transform (8) divided by $`w`$. We identify $`zw^2`$ (we thank the authors of for pointing out that this clarification was missing in the original version of the preprint). In the ensuing Borel–Weniger approximant, we integrate in the Laplace–Borel integral over the poles of the integrand with the Cauchy principal value prescription, just as in Borel–Padé approximant (9), in accordance with the full known solution (1).
The results of these calculations are presented in Figs. 1(a)–(d), as functions of the electric field strength parameter $`\stackrel{~}{a}`$, for various values of $`p\stackrel{~}{b}/\stackrel{~}{a}=0.,0.5,1.5,5.0`$. In Fig. 2 we present the analogous results for the case of the pure magnetic field (p.m.f.), as function of the magnetic field parameter $`\stackrel{~}{b}`$. N3 and $`[3/4]`$ denote the Borel–Weniger and the Borel–Padé resummations based on the truncated Borel transform (8) with the first four nonzero terms (i.e., three terms beyond the leading order); N5 and $`[5/6]`$ are based on the first six terms in (8). Comparison with the exact solutions, also present in the Figures, shows that Borel–Padé is better than the corresponding Borel–Weniger, except in the case of $`p=5.0`$ (electric field combined with a much stronger magnetic field). Fig. 2 suggests that Borel–Padé is better than Borel–Weniger for resummation of functions whose Borel transforms have singularities only outside the positive axis. Further, comparison of Fig. 2 with the results of Table I of Ref. suggests that Borel–Padé and Borel–Weniger methods are much more efficient than Weniger method in resumming series with singularities in the Borel plane. Weniger method in the p.m.f. case is better than Padé method .
We can also do analogous calculations for the induced energy densities $`\delta U`$
$`\delta 𝒰`$ $`=`$ $`a{\displaystyle \frac{\mathrm{Re}\delta }{a}}|_b\mathrm{Re}\delta ,`$ (12)
$`\delta \stackrel{~}{𝒰}(\stackrel{~}{a};p)`$ $`=`$ $`\mathrm{Re}{\displaystyle _0^{\mathrm{}}}𝑑w\mathrm{exp}\left({\displaystyle \frac{w}{\stackrel{~}{a}}}\right){\displaystyle \frac{(1)}{w}}\left[{\displaystyle \frac{pw}{\mathrm{sin}^2(w+\mathrm{i}ϵ^{})}}\mathrm{coth}(pw)+{\displaystyle \frac{1}{3}}(1+p^2)+{\displaystyle \frac{1}{w^2}}\right],`$ (13)
$`\delta \stackrel{~}{𝒰}^{\mathrm{pert}.}`$ $`=`$ $`\left[d_1(p)1!\stackrel{~}{a}^2+d_3(p)3!\stackrel{~}{a}^4+\mathrm{}\right],`$ (14)
$`d_1(p)`$ $`=`$ $`{\displaystyle \frac{1}{45}}\left[3+5p^2p^4\right],d_3(p)={\displaystyle \frac{1}{945}}\left[10+21p^27p^4+2p^6\right],\mathrm{etc}.`$ (15)
where $`\delta \stackrel{~}{𝒰}8\pi ^2\delta 𝒰/(m^4\stackrel{~}{a}^2)`$. In that case, the simple Borel transform has a double–pole structure on the positive real axis, and the Padé and Weniger approximants have trouble simulating such multiple poles adequately. Therefore, we employ a slightly modified Borel transform in the case of the induced energy densities
$$MB_U(w;p)=d_1(p)\frac{w^2}{2}+d_3(p)\frac{w^4}{4}+d_5(p)\frac{w^6}{6}+\mathrm{},$$
(16)
which has no multiple–pole structure – all the poles are simple. The (modified) Laplace–Borel integral in this case is
$$\delta \stackrel{~}{𝒰}(\stackrel{~}{a};p)=\frac{1}{\stackrel{~}{a}}_0^{\mathrm{}}𝑑w\mathrm{exp}\left(\frac{w}{\stackrel{~}{a}}\right)MB_U(w;p),$$
(17)
where again the Cauchy principal value has to be taken, once $`MB_U(w;p)`$ is replaced in (17) by its Padé or Weniger approximants. For details, we refer to Ref. where Borel–Padé was employed also for the induced energy densities. Weniger formula (11) is now applied to the modified Borel transform (16) divided by $`w^2`$. The results are presented in Figs. 3(a)–(d), as functions of $`\stackrel{~}{a}`$ at fixed $`p=0.,0.5,1.5,5.0`$, respectively.<sup>6</sup><sup>6</sup>6 In the case of the pure magnetic field, the energy density is the same as the Lagrangian density, except for the sign change. We present the solutions of Borel–Weniger and Borel–Padé based on the first four (N3, \[4/4\]) and six (N5, \[6/6\]) nonzero terms of the modified Borel transform of the energy density. We see that for the induced energy density the situation is less clear. In the cases $`p=0`$, $`0.5`$ and $`5.0`$ the Borel–Padé and Borel–Weniger resummations are apparently of comparable quality, while at $`p=1.5`$ the Borel–Padé appears to work better.
We can see these trends also if we compare the perturbation coefficients predicted by these two methods with the exact ones. These results are written in Table I for the case of the Lagrangian density (predicted $`c_9`$ and $`c_{13}`$) and in Table II for the case of the energy density (predicted $`d_9`$ and $`d_{13}`$) . Predictions of Borel–Padé and Borel–Weniger are of comparable quality in the cases of $`p=0.,0.5,5.0`$ for energy density and in the case of $`p=5.0`$ for Lagrangian density. In other cases, predictions of Borel–Padé are better.
In fact, in the case $`p=5.0`$ of the energy density, the modified Borel–Weniger is slightly, but discernibly, better than the modified Borel–Padé. Comparing predictions of Table I (for $`p=0.0`$) with those of Tables II and III of Ref. suggests strongly that the discussed Borel–Padé and Borel–Weniger methods are better than Weniger method in predicting the coefficients $`c_n`$. Weniger method is better than Padé method in predicting $`c_n`$’s .
(2) The second example to compare the efficiency of the Borel–Padé and Borel–Weniger methods will be taken from QCD, and it will have to do with the “fixing” of a pole of a Borel transform rather than with a resummation. We look at the Bjorken polarized sum rule (BjPSR), which involves the isotriplet combination of the first moments over $`x_{\mathrm{Bj}}`$ of proton and neutron polarized structure functions
$$_0^1𝑑x_{\mathrm{Bj}}\left[g_1^{(p)}(x_{\mathrm{Bj}};Q_{\mathrm{ph}}^2)g_1^{(n)}(x_{\mathrm{Bj}};Q_{\mathrm{ph}}^2)\right]=\frac{1}{6}|g_A|\left[1S(Q_{\mathrm{ph}}^2)\right].$$
(18)
Here, $`p^2=Q_{\mathrm{ph}}^2`$$`<0`$ is $`\gamma ^{}`$ momentum transfer. At $`Q_{\mathrm{ph}}^2=3\mathrm{G}\mathrm{e}\mathrm{V}^2`$ where three quarks are assumed active ($`n_f=3`$), and if taking $`\overline{\mathrm{MS}}`$ scheme and renormalization scale (RScl) $`Q_0^2=Q_{\mathrm{ph}}^2`$, we have the following TPS of the BjPSR observable $`S(Q_{\mathrm{ph}}^2)`$ available :
$`S_{[2]}(Q_{\mathrm{ph}}^2;Q_0^2=Q_{\mathrm{ph}}^2;c_2^{\mathrm{MS}},c_3^{\mathrm{MS}})`$ $`=`$ $`a_0(1+3.583a_0+20.215a_0^2),`$ (19)
$`\mathrm{with}:a_0=a(\mathrm{ln}Q_0^2;c_2^{\mathrm{MS}},c_3^{\mathrm{MS}},\mathrm{}),n_f`$ $`=`$ $`3,c_2^{\mathrm{MS}}=4.471,c_3^{\mathrm{MS}}=20.99.`$ (20)
Here we denoted by $`a`$ the strong coupling parameter $`a\alpha _s/\pi `$.
It is known from that the Borel transform $`B_S(z)`$ of $`S`$ has the lowest positive pole at $`z_{\mathrm{pole}}=1/\beta _0=4/9`$ (leading infrared renormalon) and that this pole has a much stronger residuum than the highest negative pole at $`z_{\mathrm{pole}}=1/\beta _0`$ (leading ultraviolet renormalon). The question we raise here is: How well can Padé and Weniger approximants to the Borel transform $`B_S(z)`$ determine the next coefficient $`r_3`$ of the term $`r_3a_0^4`$ in the TPS (19), via the requirement that $`z_{\mathrm{pole}}=4/9`$? For that, we have to know well the actual $`r_3`$. That term can be determined reasonably well on the basis of two approximants discussed in – the effective charge approximant (ECH) $`𝒜_S^{(\mathrm{ECH})}(c_3)`$ with $`c_320.`$, and another, also RScl– and scheme–independent approximant $`𝒜_{S^2}^{1/2}(c_3)`$ with $`c_315.5`$. These two approximants give the correct location of the leading infrared renormalon pole, and when we expand them back in powers of $`a_0`$ we obtain $`r_3129.4`$ and $`r_3130.8`$, respectively. Therefore, we can estimate with high confidence the actual $`r_3`$: $`r_3=130.\pm 1`$.
It is important to consider the RScl– and scheme–invariant Borel transform when we want to apply Padé or Weniger approximants to it, so that the predicted values of $`r_3`$ will be independent of the RScl– and scheme in which we work at the intermediate stage. Such a Borel transform has been used in , and we use its variant $`\stackrel{~}{B}_S(z)`$ as specified in \[cf. Eqs. (18)–(20) there\]. Such a Borel transform reduces (up to a $`z`$–dependent nonsingular factor) to the usual Borel transform in the approximation of the one–loop evolution. The resulting power expansion of $`\stackrel{~}{B}_S(z)`$ up to $``$$`z^3`$ will depend on the coefficient $`r_3`$
$$\stackrel{~}{B}_S(z)=1+\frac{32}{81}(\gamma 1)y+0.02078\mathrm{}y^2+\frac{8}{729}(21.88\mathrm{}+\frac{1}{6}r_3)y^3+𝒪(y^4),$$
(21)
where $`\gamma =0.577\mathrm{}`$ is Euler constant, and $`y2\beta _0z`$. If we apply $`[2/1]`$ and $`[1/2]`$ Padé approximants to the TPS (21) and demand $`z_{\mathrm{pole}}=1/\beta _0`$ ($`y_{\mathrm{pole}}=2`$), we obtain predictions $`r_3=137.0`$ and $`r_3=128.0`$, respectively. The prediction of $`[1/2]`$ is significantly better, and this could possibly be explained with the more involved denominator structure of $`[1/2]`$ in comparison to $`[2/1]`$. When applying to (21) Weniger formula (11) ($`\delta _2^{(0)}`$ with $`\zeta =1`$), we obtain $`r_3=135.3`$. This is further away from the actual value of $`130.\pm 1.`$ than the prediction of $`[1/2]`$. In both $`[1/2]`$ and $`\delta _2^{(0)}`$, the denominators are polynomials of quadratic degree in $`z`$.
To summarize this QCD example: We applied Padé and Weniger approximants to a (TPS of a) Borel transform of the Bjorken polarized sum rule and demanded that the leading infrared renormalon pole be reproduced correctly. Weniger approximant $`\delta _2^{(0)}`$ then apparently gives a somewhat worse prediction for the next coefficient than the corresponding Padé approximant $`[1/2]`$.
The work of G.C. was supported by the Korean Science and Engineering Foundation (KOSEF). The work of J.-Y.Y. was supported by the German Federal Ministry of Science (BMBF).
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# Stability properties of trapped Bose-Fermi gases mixture
## Abstract
The stability of Bose-Fermi gases trapped in an isotropic potentials at ultracold temperature is strongly influenced by the interaction between the fermions and the bosons. At zero temperature, the stability criterion is given in this paper using variation method, the results show that whether a fermion-boson mixture is stable depends mainly on the interaction between the fermions and the bosons. For finite temperature, however, the stability is not only related to the coupling constants, but also to the temperature. The stability conditions for finite temperature are also derived and discuss in details in this paper.
PACS number(s):03.75.Fi, 05.30.Fk,05.30.Jp
Since the realization of dilute alkali atomic vapor condensates(Bose-Einstein condensation or BEC) in 1995, large efforts have been make to study many-body effects and macroscopic properties of the gases, which may be more transparently demonstrated in BEC than in other many-body systems. For fermionic atomic vapor, however, it is difficult to achieve a degenerate gas. Since the evaporative cooling of a pure fermionic gas is ineffective at temperature sufficiently low due to the suppression of $`s`$ wave scattering between identical fermions. As theory and experiment advanced, a new rich phenomenology has appeared in which new conditions arise, which are not accessible in other BEC systems. One of the most stunning of these is the recent experimental demonstration of a condensate mixture composed of two spin states of $`{}_{}{}^{87}Rb`$. The realization of two condensates mixture is related to the sympathetic cooling mechanism, i.e., the exchange of energy due to elastic collisions between atoms of cooled and thermal samples. Most recently, B.DeMarco and D.S.Jin report their observation of degenerate Fermi gas using an evaporative cooling strategy. Although the strategy uses a two-component Fermi gas, the mixture of Bose and Fermi gas attracts a lot of attention from the viewpoint of both experiment and theoretical study.
The mixed system of Bose and Fermi particles is itself an interesting subject for investigation. The hydrogen deuterium system has been studied at the early stage of these investigations, and there is now a lot of literature devoted to the properties of pure degenerate trapped atomic Fermi gases\[5-9\].
In a recent paper, M$`\varphi `$lmer has used a simple mean field models to study the spatial distribution of a Bose-Fermi gas mixture at $`T=0K`$ within Thomas-Fermi approximation. The results show that the distributions depend strongly on the relative sign and magnitude of the boson-boson and boson-fermion scattering lengths. Here, we shall study the Bose-Fermi gas mixture using a variation method at zero temperature, this method was first introduced in to study the BEC ground state in a harmonic trap of a Bose system, and later generalized by H.Shi and W.M.Zheng to study BEC with attractive interactions. In addition, we study the stability of the Bose-Fermi gas mixture at finite temperature. The results show that there is a region of temperature in which the phase separation of the mixture happens. And the span of the region depends on the coupling constants.
To begin, we consider a second-quantized grand canonical Hamiltonian of interacting Bose and Fermi gases
$`H`$ $`=`$ $`H_b+H_f+V_{bf},`$ (1)
$`H_b`$ $`=`$ $`{\displaystyle 𝑑r\varphi ^+(r)(\frac{p^2}{2m_b}\mu _b+\frac{1}{2}m_b\omega _br^2)\varphi (r)}+{\displaystyle \frac{g_{bb}}{2}}{\displaystyle 𝑑r𝑑r^{^{}}\varphi ^+(r)\varphi ^+(r^{^{}})\varphi (r^{^{}})\varphi (r)},`$ (2)
$`H_f`$ $`=`$ $`{\displaystyle 𝑑r\psi ^+(r)(\frac{p^2}{2m_f}\mu _f+\frac{1}{2}m_f\omega _fr^2)\psi (r)},`$ (3)
$`V_{bf}`$ $`=`$ $`g_{bf}{\displaystyle 𝑑r𝑑r^{^{}}\varphi ^+(r)\psi ^+(r^{^{}})\delta (rr^{^{}})\psi (r^{^{}})\varphi (r)},`$ (4)
where $`\varphi (r)`$ and $`\psi (r)`$ denote boson and fermion field operators with masses $`m_b`$ and $`m_f`$, respectively. For weakly interacting dilute gases, the interactions between the bosonic atoms are modeled by $`\delta `$ potentials and the interactions among the fermionic atoms are neglected, since the interactions between atoms at very low temperature is suppressed for polarized systems. $`g_{bb}`$ and $`g_{ba}`$ stand for boson-boson and boson-fermion coupling constant, respectively.
$$g_{bb}=\frac{4\pi \mathrm{}^2}{m_b}a_{bb},g_{bf}=\frac{2\pi \mathrm{}^2}{m_{bf}}a_{bf},$$
$`a_{bb}`$($`a_{bf}`$) are $`s`$-wave scattering length between boson and boson (boson and fermion), and $`m_{bf}`$ is a reduced mass of the boson and the fermion. The chemical potentials $`\mu _b`$ and $`\mu _f`$ are determined through the conditions
$$N_b=𝑑r\varphi ^+(r)\varphi (r),N_f=𝑑r\psi ^+(r)\psi (r).$$
(5)
At $`T=0`$, self-consistent mean field theory, assuming that all $`N`$ bosonic particles in a gas populated the same state denoted by single particle wave function $`\mathrm{\Phi }(r)`$, lead to a nonlinear Schrödinger equation (or the Gross-Pitaevskii equation) for $`\mathrm{\Phi }(r)=\varphi (r)`$
$$[\frac{\mathrm{}^2}{2m_b}^2+\frac{1}{2}m_b\omega _b^2r^2+g_{bb}n_b(r)]\mathrm{\Phi }(r)=E_b\mathrm{\Phi }(r),$$
(6)
we here omit quantities $`g_{bf}n_f(r)`$, which is smaller than $`g_{bb}n_b(r)`$ in the case of $`N_b>>N_f`$. In order to get a degenerate fermionic gas, the boson particles appear in the system only as a coolant, so the number of bosons is always much larger than the number of fermions. In the same approximation, the fermionic wave function is given by a Slater determinant
$$\mathrm{\Psi }(r_1,r_2,\mathrm{},r_{N_f})=\frac{1}{\sqrt{N_f!}}\left[\begin{array}{cccc}\mathrm{\Psi }_1(r_1)& \mathrm{\Psi }_1(r_2)& \mathrm{}& \mathrm{\Psi }_1(r_{N_f})\\ \mathrm{\Psi }_2(r_1)& \mathrm{\Psi }_2(r_2)& \mathrm{}& \mathrm{\Psi }_2(r_{N_f})\\ & & \mathrm{}& \\ & & \mathrm{}& \\ & & \mathrm{}& \\ \mathrm{\Psi }_{N_f}(r_1)& \mathrm{\Psi }_{N_f}(r_2)& \mathrm{}& \mathrm{\Psi }_{N_f}(r_{N_f})\end{array}\right],$$
(7)
where $`\mathrm{\Psi }_i(r)`$ is the single particle states determined by Hartree-Fock self-consistent equation
$$[\frac{\mathrm{}^2}{2m_f}^2+\frac{1}{2}m_f\omega _f^2r^2+g_{bf}n_b(r)]\mathrm{\Psi }_i(r)=E_i\mathrm{\Psi }_i(r).$$
(8)
The density of the fermions is given by
$$n_f(r)=|\mathrm{\Psi }(r)|^2.$$
(9)
In the semiclassical (Thomas-Fermi) approximation, the particle are assigned classical position and momenta, but the effects of quantum statistics are taken into account. Under this approximation, the Eqs.(3) and (5) for the boson and fermion wave function are equivalent to
$`{\displaystyle \frac{1}{2}}m_b\omega _b^2r^2+g_{bb}n_b(r)`$ $`=`$ $`\mu _b,`$ (10)
$`{\displaystyle \frac{\mathrm{}^2}{2m_f}}[6\pi ^2n_f(r)]^{\frac{2}{3}}+{\displaystyle \frac{1}{2}}m_f\omega _f^2r^2+g_{bf}n_b(r)`$ $`=`$ $`e_F.`$ (11)
The main conclusion of this equations is discussed in Ref.. We obtain $`n_b(r)=\frac{1}{g_{bb}}(\mu _b\frac{1}{2}m_b\omega _b^2r^2)`$ from the first line of Eqs(7). Substituting $`n_b(r)`$ into the second line of Eqs(7), we yield
$$\frac{\mathrm{}^2}{2m_f}[6\pi ^2n_f(r)]^{\frac{2}{3}}+\frac{1}{2}m_f\omega _f^2r^2+\frac{g_{bf}}{g_{bb}}(\mu _b\frac{1}{2}m_b\omega _b^2r^2)=e_F,$$
(12)
this equation shows that the fermions experience a potential minimum in the center of the trap if $`g_{bf}/g_{bb}<m_f\omega _f^2/m_b\omega _b^2,`$ in this case the entire distribution behaves like a fermionic core within the Bose condensate. The fermion density is a constant throughout the Bose condensate if $`g_{bf}/g_{bb}=m_f\omega _f^2/m_b\omega _b^2.`$ Whereas the fermions are repelled from the center of the trap and localized near the edge of the Bose condensate if $`g_{bf}/g_{bb}>m_f\omega _f^2/m_b\omega _b^2,`$ i.e. a phase separation occurs in this system. We would like to note that the distribution of BEC remains unchanged in the above discussions, since we assume $`N_b>>N_f`$. To drive Eq.(8), we assume that the Thomas-Fermi approximation(TFA) is valid. The coupling constant $`g_{bb}`$ and $`g_{bf}`$ may take any value as long as the TFA is available, and the phase separation depend mainly on ratio $`g_{bf}/g_{bb}`$. In what follows we discuss the separation of the bosonic and fermionic parts from the other aspect for zero temperature by using variation method, the results are indeed different from those under TFA. We note the solution of Eq.(5) requires prior knowledge of the boson density profile $`n_b=|\mathrm{\Phi }(r)|^2`$. To obtain the density profile, we have to solve the Gross-Pitaevskii equation (3). There are a large number of literatures devoted to solve the Gross-Pitaevskii equation, we here use a variation method to solve the problem. For a isotropic trapping potential, we may assume the trial wave function for $`\mathrm{\Phi }(r)`$ in Eq.(3) to be
$$\mathrm{\Phi }(r)=\sqrt{N_b}\omega ^{\frac{3}{4}}(\frac{m_b}{\pi \mathrm{}})^{\frac{3}{4}}e^{m_b\omega r^2/2\mathrm{}},$$
(13)
where $`\omega `$ is the effective frequency and is taken as a variational parameter. Substituting Eq.(9) into Eq.(3), we obtained the ground-state energy
$$E_b[\mathrm{\Phi }]=E_b(\omega )=\frac{3}{4}N_b\mathrm{}\omega +\frac{3}{4}N_b\mathrm{}\frac{\omega _b^2}{\omega }+g_{bb}N_b^2(\frac{\omega m_b}{2\pi \mathrm{}})^{\frac{3}{2}}.$$
(14)
If $`E_b(\omega )`$ is plotted as a function of $`\omega `$ , one sees that a stable local minimum exists only up to a certain maximum number of atoms for $`g_{bb}<0`$. The critical point occurs where
$$\frac{E_b(\omega )}{\omega }|_{(\omega =\omega _c,N_b=N_{bc})}=0,\text{ and }\frac{^2E_b(\omega )}{\omega ^2}|_{\omega =\omega _c,N_b=N_{bc}}>0.$$
(15)
Here, $`\omega _c`$ stands for the variational parameter that minimizes the ground state energy. Using equation (10), for $`g_{bb}<0`$ the critical number of bosons is given by
$$N_b^c<2\mathrm{}\omega _b^2(\omega _c)^{\frac{5}{2}}\frac{1}{|g_{bb}|}(\frac{2\pi \mathrm{}}{m_b})^{\frac{3}{2}}.$$
(16)
where $`\omega _c`$ satisfies
$$\mathrm{}\omega _c^2\mathrm{}\omega _b^2+2g_{bb}N_b(\frac{m_b}{2\pi \mathrm{}})^{\frac{3}{2}}\omega _c^{\frac{5}{2}}=0.$$
(17)
Parameter $`\omega _b166Hz`$ relevant to the experiment gives $`N_b^c1400`$, which is in good agreement with the experiment. The solution $`\omega _c`$ of Eq.(13) against $`g_{bb}`$ is plotted in Fig.1, which shows that as $`|g_{bb}|`$ increases, the variation parameter $`\omega _c`$ decreases, and it has a maximum equal to $`\omega _b`$ at $`g_{bb}=0`$. We will use this solution to study the stability of the mixture at zero temperature below.
We may determine the ground-sate energy functional of the fermions provided $`n_b(r)`$ is known. In terms of the fermion distribution $`n_f(r)`$, the energy functional $`E_f`$ of the fermions is given by
$$E_f=E_f[n_f(r)]=\frac{d^3r}{6\pi ^2}\frac{\mathrm{}^2}{2m_f}[6\pi ^2n_f(r)]^{5/3}+\frac{1}{2}d^3rm_f\omega _f^2r^2n_f(r)+d^3rg_{bf}n_f(r)n_b(r),$$
(18)
since the interaction between the bosons and the fermions is rather week, we may consider a Gaussian function as a trial fermions’ distribution
$$n_f(r)=N_f\mathrm{\Omega }^{3/2}(\frac{m_f}{\pi \mathrm{}})^{3/2}e^{\frac{m_f\mathrm{\Omega }(rr_f)^2}{\mathrm{}}}.$$
(19)
Here $`r_f`$ and $`\mathrm{\Omega }`$ is treated as variation parameters. Substituting this wave functions into Eq.(14), one obtains
$$E_f=E_f(\mathrm{\Omega },r_f,N_f,\omega _c)=P+\frac{1}{2}m_f\omega _f^2r_f^2N_f+hN_bN_f(\pi G)^{\frac{3}{2}}e^{Gr_f^2}$$
(20)
with
$$P=P(\mathrm{\Omega },N_f)=(\frac{3}{5})^{\frac{3}{2}}\frac{1}{\pi }(6\pi ^2)^{\frac{2}{3}}\mathrm{}\mathrm{\Omega }N_f^{\frac{5}{3}}+\frac{3\mathrm{}\omega _f^2}{4\mathrm{\Omega }}N_f$$
$$G=G(\mathrm{\Omega },\omega _c)=\frac{m_fm_b\omega _c\mathrm{\Omega }}{\mathrm{}(m_f\mathrm{\Omega }+m_b\omega _c)}.$$
As known, a physical state corresponds to a stable or metastable point of the energy functional. If a separation of the fermion and boson component occurs, then $`r_{fc}`$ that minimizes the energy $`E_f`$ takes a positive nonzero value. In other words, there are no separations between the two components when energy $`E_f`$ exhibits a minimal value at $`r_f=0`$. When distribution function is restricted to the form of the trial function (15) we may write the conditions of a minimal energy in terms of derivatives of the energy with respect to the adjustable variation parameters of the trial function. We show them as follows
$`{\displaystyle \frac{E_f}{\mathrm{\Omega }}}|_{\mathrm{\Omega }=\mathrm{\Omega }_c}=0,{\displaystyle \frac{E_f}{r_f}}|_{r_f=r_{fc}}=0,`$ (21)
$`{\displaystyle \frac{^2E_f}{\mathrm{\Omega }^2}}{\displaystyle \frac{^2E_f}{r_f^2}}({\displaystyle \frac{^2E_f}{\mathrm{\Omega }r_f}})^2>0,`$ (22)
the stationary conditions (there are no separation between the boson and fermion ) are
$`Y`$ $`=`$ $`Y(N_b,\omega _c,g_{bf})=[\mathrm{}{\displaystyle \frac{\omega _f^2}{\mathrm{\Omega }_c^3}}+g_{bf}N_b\pi ^{\frac{3}{2}}\sqrt{G(\mathrm{\Omega }_c,\omega _c)}{\displaystyle \frac{^2G(\mathrm{\Omega }_c,\omega _c)}{\mathrm{\Omega }^2}}+{\displaystyle \frac{g_{bf}N_b\pi ^{3/2}}{2G(\mathrm{\Omega }_c,\omega _c)}}({\displaystyle \frac{G(\mathrm{\Omega }_c,\omega _c)}{\mathrm{\Omega }}})^2]`$ (23)
$`\times `$ $`[m_f\omega _f^22g_{bf}N_b\pi ^{3/2}G^{5/2}(\mathrm{\Omega }_c,\omega _c)]>0,`$ (24)
where $`\mathrm{\Omega }_c`$ determined by $`E_f/\mathrm{\Omega }|_{\mathrm{\Omega }_c}=0`$ satisfies the following equation
$$\frac{2}{3\pi }(6\pi ^2)^{\frac{2}{3}}(\frac{3}{5})^{\frac{3}{2}}\mathrm{}N_f^{\frac{2}{3}}\frac{1}{2}\mathrm{}(\frac{\omega _f}{\mathrm{\Omega }_c})^2+g_{bf}N_b\pi ^{\frac{3}{2}}\sqrt{G(\mathrm{\Omega }_c,\omega _c)}\frac{G(\mathrm{\Omega }_c,\omega _c)}{\mathrm{\Omega }}=0,$$
(25)
the solution of Eq.(19) as a function of $`g_{bf}`$ is shown in Fig.2, a magnification part of the curve near $`g_{bf}=0`$ (but $`g_{bf}>0`$) is give in the inset. This curve indicates that the fermions prefer to occupy the trap centre for $`g_{bf}<0`$, and the larger the coupling constant $`|g_{bf}|(g_{bf}<0)`$, the sharper the distribution of the fermions. However, as we show below, the fermions and the bosons can not always coexist even if $`g_{bf}<0`$. $`\frac{E_f}{r_f}|_{r_f=r_{fc}}=0`$ has two solutions, one solution is $`r_{fc}=0`$ and the another is
$$r_{fc}=\sqrt{\frac{1}{G}ln(2g_{bf}N_b\pi ^{\frac{3}{2}}G^{\frac{5}{2}}(\mathrm{\Omega }_c,\omega _c))ln(m_f\omega _f^2)}.$$
(26)
For $`g_{bf}<0`$ or
$$0<g_{bf}<\frac{m_f\omega _f^2}{2N_b\pi ^{\frac{3}{2}}G^{\frac{5}{2}}(\mathrm{\Omega }_c,\omega _c)}$$
i.e. the interaction between fermion and the boson is attractive or weekly repulsive, the solution $`r_{fc}=0`$ holds, which indicate that there is not separation between the fermions and the bosons. For $`g_{bf}>\frac{m_f\omega _f^2}{2N_b\pi ^{\frac{3}{2}}G^{\frac{5}{2}}(\mathrm{\Omega }_c,\omega _c)}`$ the fermions experience a effective potential minimum at $`r_{fc}=\sqrt{\frac{1}{G}ln(2g_{bf}N_b\pi ^{\frac{3}{2}}G^{\frac{5}{2}}(\mathrm{\Omega }_c,\omega _c))ln(m_f\omega _f^2)}>0,`$ the Bose condensate is surrounded by a shell of fermions in this case. $`Y(N_b,g_{bf},\omega _c)`$ as functions of the coupling constant $`g_{bf}`$ are shown in figure 3. We see that whether the boson-fermion mixture is stable depends not only on the coupling constant $`g_{bf}`$ and $`g_{bb}`$(through $`\omega _c)`$, but also on $`N_b`$ and $`N_f`$(through $`\mathrm{\Omega }_c`$), i.e., the stability of the mixture system depends on the number of both boson and fermion system. For example, in Fig.3-a we show $`Y`$ given by Eq.(8) as a function of the coupling constant $`g_{bf}`$ for fixed $`N_f=100,N_b=10000`$, while Fig.3-b is for the same parameters as in Fig.3-a except for $`N_b=1000`$, it is obvious that the region of $`g_{bf}`$ in which the system has no phase separation has been broadened with $`N_b`$ decreases (for fixed $`N_f`$). The inset present the dependence of $`Y`$ on $`g_{bf}`$ at a larger scale of $`g_{bf}.`$ It is interesting to compare the above mentioned results with those obtained by treating the fermions in the Thomas-Fermi approximation, this is done in Ref., and we note that the semiclassical description gives a qualitatively correct description and it reliably predicts the phase separation.
Now we tune our attention to discuss the above problem at finite temperature. First of all, we consider the homogenous case, for the boson and fermion system, thermodynamical properties are trivial if there are not interaction between them. But in this case the sympathetic cooling scheme does not take any effect and the degenerate fermions in a trapped potential have not been achieved. The thermodynamical properties may be changed when the interaction between the fermions and bosions is turn on, then a new phenomenon, the phase separation, may occur in this system. for a homogeneous fermion and boson mixture system, the Helmholts free energy can be written as
$$\beta F=\frac{V}{\lambda _2^3}f_{\frac{5}{2}}(z_f)+\frac{1}{2}g_{ff}\rho _fN_f\lambda _f^2+ln(1z_b)\frac{V}{\lambda _b^3}g_{\frac{5}{2}}(z_b)+2g_{bb}\rho _bN_b\lambda _b^2+g_{bf}(\lambda _b^2+\lambda _f^2)N_fN_b/V,$$
(27)
where index $`f`$ refers to the fermionic component, whereas index $`b`$ stands for the bosonic one, $`N_i`$ is the number of particles in component $`i`$, $`\lambda _i`$ denotes the thermal wave length of component $`i`$, $`f_n(z)`$ and $`g_n(z)`$ represent the Fermi and Bose integral, respectively.The equation (21) is based on the pseudopotential form of the atom-atom interaction, and may be assumed accurate when the system is dilute. i.e. $`\rho _ig_{ii}^3<<1`$ and $`g_{ii}/\lambda _i<<1`$, where $`\rho _i`$ is the density of the component $`i`$. This condition is well satisfies for the samples of alkali atoms in experiments to date.
From eq.(21) we obtain the chemical potential for each component straightforwardly,
$`\beta \mu _b`$ $`=`$ $`\beta \mu _b^0+4g_{bb}\rho _b\lambda _b^2+g_{bf}(\lambda _b^2+\lambda _f^2)N_f/V,`$ (28)
$`\beta \mu _f`$ $`=`$ $`\beta \mu _f^0+g_{ff}\rho _f\lambda _f^2+g_{bf}(\lambda _b^2+\lambda _f^2)N_b/V,`$ (29)
where $`\mu _i^0`$ are the chemical potentials of ideal gas. There are three terms in each chemical potential, the second term comes from the interaction within the component and the third term is from the interaction between the fermion and boson component. As known, an homogenous binary mixture is stable only when the symmetric matrix $`\widehat{\mu }`$ given by
$$\widehat{\mu }=\left[\begin{array}{cc}\frac{\mu _b}{\rho _b}\hfill & \frac{\mu _b}{\rho _f}\hfill \\ \frac{\mu _f}{\rho _b}\hfill & \frac{\mu _f}{\rho _f}\hfill \end{array}\right]$$
(30)
is non-negatively definite, in other words, all eigenvalues of matrix $`\widehat{\mu }`$ given in Eq.(23) are non-negative. Mathematically, for homogeneous fermion and boson mixture the stability conditions are
$$\frac{\mu _b}{\rho _b}0,\frac{\mu _f}{\rho _f}0,$$
(31)
and
$$det\left[\begin{array}{cc}\frac{\mu _b}{\rho _b}\hfill & \frac{\mu _b}{\rho _f}\hfill \\ \frac{\mu _f}{\rho _b}\hfill & \frac{\mu _f}{\rho _f}\hfill \end{array}\right]0.$$
(32)
For ideal gas, we have $`\rho _b=\frac{1}{\lambda _b^3}g_{\frac{3}{2}}(z_b),`$ $`\rho _f=\frac{1}{\lambda _f^3}f_{\frac{3}{2}}(z_f),`$ this leads to
$$\beta \frac{\mu _f^0}{\rho _f}=\frac{\lambda _f^3}{f_{\frac{1}{2}}(z_f)},\beta \frac{\mu _b^0}{\rho _b}=\frac{\lambda _b^3}{g_{\frac{1}{2}}(z_b)},$$
(33)
It follows from eqs (24) and (25) that
$`4g_{bb}\lambda _b^2+{\displaystyle \frac{\lambda _b^3}{g_{\frac{1}{2}}(z_b)}}0,`$ (34)
$`g_{ff}\lambda _f^2+{\displaystyle \frac{\lambda _f^3}{f_{\frac{1}{2}}(z_f)}}0,`$ (35)
and (36)
$`Z(T,g_{bf},g_{ff},g_{bb})=Z=(4g_{bb}\lambda _b^2+{\displaystyle \frac{\lambda _b^3}{g_{\frac{1}{2}}(z_b)}})(g_{ff}\lambda _f^2+{\displaystyle \frac{\lambda _f^3}{f_{\frac{1}{2}(z_f)}}})g_{bf}^2(\lambda _b^2+\lambda _f^2)^20.`$ (37)
It is well known that a homogeneous imperfect gas with attractive interaction is not stable. The fermions in this kind of gas could form BCS state, which consists two fermionic particles interacting with each other but not with the other fermions from the Fermi gas, whereas bosons with attractive interaction could collapse into liquid. Hence, we here discuss the system with repulsive interactions. It is obvious that the stability condition (27) and (28) hold always for $`g_{bb}>0`$, $`g_{ff}>0`$. We would like to point out that the stability conditions (27-29) do not involve the densities of the both components. At first sight, this seems to be confusion, in fact, there is no contradiction. One can demonstrate that at low density the Helmholtz free energy of the bogoliubov gas reduce to a quadratic form in $`N_b`$ and $`N_f`$. To have a minimum, this form should be positive definite, i.e., $`\text{det}\frac{^2F}{N_bN_f}0.`$ Therefore, the corresponding stability criterion involves only density-independent constants in the order of approximation used. This criterion is similar to the stability conditions for two-component Bose-Einstein condensate in a trapped untracold gas\[19-25\]. When $`T\mathrm{},\lambda _i0`$, hence $`Z\frac{1}{\rho _b\rho _f}`$. Thus at high temperature, the homogeneous binary gas mixture is always stable and no phase separation occur. In the case considered here, Fermi temperature $`T_F=\frac{h^2}{2mk_B}(\frac{3N_f}{8\pi V})^{\frac{2}{3}}`$ is much lower than BEC temperature $`T_c=\frac{h^2}{2\pi mK_B}(\frac{N_b}{2.612V})^{\frac{2}{3}},`$ i.e.,as temperature decreases, it first passes the BEC transition point $`T_c`$. When $`TT_c`$, $`g_{\frac{1}{2}}(1)\mathrm{}`$, so
$$Z(T,g_{bf},g_{bb},g_{ff})4g_{bb}\lambda _b^2(g_{ff}\lambda _f^2+\frac{\lambda _f^3}{f_{\frac{1}{2}}(z_f)})g_{bf}^2(\lambda _b^2+\lambda _f^2)^2.$$
In particular, when $`T<<T_F`$,i.e., the temperature is much smaller than the Fermi temperature of the fermion system, the stability condition becomes (setting $`m_f=m_b`$)
$$g_{bb}g_{ff}g_{bf}^20,$$
(38)
which does not depend on temperature and coincides with the stability conditions of two-component BEC. Although it is difficult to reach this region of very low temperature, yet it attracts much more attention. Because both superfluidity and shell effects are expected to occur at temperature much smaller than the Fermi temperature. $`Z`$ given by Eq.(29) as a function of the temperature is shown in Fig.4, we see that the system is always stable when $`T0`$ and $`T\mathrm{}`$, and the system is unstable for $`T_{c1}<T<T_{c2}`$, where $`T_{c1}`$ and $`T_{c2}`$ are roots of $`Z(T,g_{bb},g_{bf},g_{ff})=0`$. In particular, $`T_{c1}`$ and $`T_{c2}`$ depend on $`g_{bf}`$ $`g_{ff}`$ and $`g_{bb}`$. As $`g_{bf}`$ decreases (for fixed $`g_{bb}`$ and $`g_{ff}`$), $`T_{c1}`$ tends to $`T_{c2}`$ (in Fig.4 going from dotted line to solid line). The critical temperature $`T_{c1}`$ and $`T_{c2}`$ characterize the onset of the phase separation, which is quite different from the Bose-Einstein condensation and the degenerate fermions. The critical temperature of BEC and of the onset of degenerate fermionic gas depend mainly on the density of the system $`N_i/V(i=b,f)`$. Especially, the BEC and the degenerate fermionic gas may happen even if $`g_{bf}=0`$. For the phase separation, however, nothing will happen if $`g_{bf}=0`$. For a fixed temperature and the coupling constant $`g_{bf}`$, $`Z`$ vs. $`g_{bb}`$ and $`g_{ff}`$ is shown in Fig.5, which represents the dependence of the stability on the interaction strength inside each component.
Until now, we considered only a homogeneous Fermi-Bose gas mixture at finite temperature. In reality, however, experiments with ultracold atoms are performed by trapping and cooling in an external potential that can be generally modeled by an isotropic harmonic oscillator $`V(r)=\frac{m}{2}\omega _t^2r^2`$, where $`\omega _t`$ is the trapping frequency. An exact criterion for the stability of an inhomogeneous Bose-fermi mixture should involve calculating the Helmholts free energy as a function at all eigenstates of the trapping potential. Fortunately, in the system considered here it is a good approximation to take use of the local-density approximation, which treats the system as being locally homogeneous. This requires that the level spacing $`\mathrm{}\omega _t`$ of the trapping potential is much smaller than the Fermi energy. Of course, the local density approximation always breaks down at the edge of the gas cloud where the density vanishes and the effective Fermi energy becomes zero. In this approximation, the stability conditions can still be calculated by means of the equations derived above, with the understanding that now the effective chemical potentials are spatially dependent through
$$\mu _b=\mu _b^0\frac{1}{2}m_b\omega _b^2r^2,\mu _f=\mu _f^0\frac{1}{2}m_f\omega _f^2r^2.$$
Thus a local stability condition is the same as given in Eq.(29) but replacing $`z_i(i=1,2)`$ by
$$\stackrel{~}{z}_1=z_1e^{\frac{\beta }{2}m_b\omega _b^2r^2},\text{ and }\stackrel{~}{z}_2=z_2e^{\frac{\beta }{2}m_f\omega _f^2r^2}.$$
As shown in inset of Fig.4, the region of temperature in which the system is unstable decrease for the case of $`r0`$. As compared with the case without trapped potential, the total energy of the system increases for it in a trap. Alternatively, within the TFA, the chemical potentials decrease in this process. So this effect is equal to be that the particle number of the system has a loss. In this sense, the system is more stable than before.
In summary, we considered a dilute Bose-Fermi gas mixture in an isotropic trap. The atom can interact via s-wave scattering except within the fermions. These interactions strongly affect the stability of the system at zero and finite temperature. In addition, the stability conditions depend on the ratio rate $`N_b/N_f`$, the larger the ratio rate, the smaller the region of stability. For finite temperature, however, the stability conditions depends not only on the interactions,but also on the temperature. The region $`T_{c1}TT_{c2}`$ in which the system is unstable depend on the strength of the interaction between and within the bosons and the fermions. For an anisotropic trap, the stability conditions remain unchanged, whereas somewhat would be changed for zero temperature compared with the case of isotropic trap. To study the effects, we should introduced the another variation parameter in Eqs (9) and (5) to characterize the BEC and the fermions in this trap. Consequently, the stability condition (18) for zero temperature changes and the phase separation could different for different orientation. These need further investigations.
ACKNOWLEDGEMENT:
We thank Dr. Li You for his stimulating and helpful discussions.
Figure captions:
Fig. 1:The parameter $`\omega _c`$ which minimizes the energy functional versus the coupling constant $`g_{bb}`$. The trapped frequency $`\omega _b=166Hz`$ and the number of the bosonic atom $`N_b=1000`$.
Fig. 2:The parameter $`\mathrm{\Omega }_c`$ which minimizes the energy functional as a function of the coupling constant $`g_{bf}`$. The parameters chosen are $`g_{bb}=0.05`$ in units of $`\mathrm{}\omega _fa^3`$ ($`a=\sqrt{\mathrm{}/\omega _bm_b}`$) and all the coupling constants are chosen in this units hence forth, $`N_f=100`$, $`\omega _f=166Hz`$. Scatter and solid line correspond to different number of bosonic atom, as specified in the figure. The inset presents the enlarged part of the curve near $`g_{bf}=0(>0)`$.
Fig.3:Plot of $`Y`$ given by Eq.(18) as a function of coupling constant $`g_{bf}`$. The parameter chosen are a:$`N_f=100,N_b=1000.`$ b:$`N_f=100,N_b=10000.`$ The curve for a larger scale of $`g_{bf}`$ is presented as an inset in the figure.
Fig.4:Plot of $`Z`$ given by Eq.(29) as a function of temperature $`T`$. The parameters chosen are $`N_b=1000,N_f=10000,g_{bb}=0.05,g_{ff}=0.01`$. Dashed-dotted line:$`g_{bf}=0.3`$, dotted line $`g_{bf}=0.02`$, solid line $`g_{bf}=0.01`$. The dotted line in the inset is the same as the dotted line in the figure, while the solid line in the inset is for the gases in a trap with trapped frequency $`166Hz`$.
Fig.5:Plot of $`Z`$ as a function of $`g_{ff}`$ and $`g_{bb}`$. The parameters chosen are temperature $`T=0.1T_F`$, $`g_{bf}=0.2`$.
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# Explicit regulator maps on polylogarithmic motivic complexes
## 1 Introduction
1. Regulator maps on the level of complexes. Let $`X`$ be an algebraic variety. Beilinson \[B1\] defined the rational motivic cohomology of $`X`$ via the algebraic K-theory of $`X`$ by the formula
$$H_{Mot}^i(X,(n)):=gr_n^\gamma K_{2ni}(X)$$
Beilinson \[B2\] and Lichtenbaum \[L\] conjectured that the weight $`n`$ motivic cohomology of $`X`$ should appear as cohomology groups of some complexes, called the weight $`n`$ motivic complexes of $`X`$.
If $`X`$ is defined over $``$ Beilinson \[B1\] constructed the regulator map to the Deligne cohomology of $`X`$:
$$H_{Mot}^i(X,(n))H_D^i(X,(n))$$
The regulator map plays key role in the (hypothetical) formulas for special values of $`L`$-functions of motives over number fields \[B1\].
Our point is that the regulator map should be explicitly defined on the level of complexes. So for any algebraic variety over $``$ one should have homomorphisms of complexes
$$\text{weight }n\text{ motivic complex of }X\text{weight }n\text{ Deligne complex of }X$$
(1)
The cone of this map, shifted by $`1`$, defines the Arakelov motivic complex:
$$\mathrm{R}\mathrm{\Gamma }_X^𝒜(n):=\mathrm{Cone}\left(\text{the map (}\text{1}\text{)}\right)[1]$$
and so its cohomology should be called the Arakelov motivic cohomology of $`X`$. The group $`H^{2n}(\mathrm{R}\mathrm{\Gamma }_X^𝒜(n))`$ is canonically isomorphic to the group of codimension $`n`$ Arakelov cycles on $`X()`$ (see s. 3.2). This isomorphism is transparent for a version of the Deligne complex recalled in s. 3.1. Examples are given in the section 3. (When $`X`$ is defined over $``$ we should take the cone of the regulator map to the so called real Deligne cohomology of $`X_{}`$, see s. 3.2).
Motivic complexes are objects of the derived category. Several candidates for motivic complexes are known, each with its own charm. They should be quasiisomorphic. Explicit regulator maps should be defined for each of them.
First of all there are Bloch’s higher Chow complexes (\[Bl\]). They satisfy many of the expected properties of motivic complexes. Explicit regulator maps from these complexes to the Deligne complexes were constructed in \[G4\] using the Chow polylogarithm construction given in \[G3\].
The goal of the present paper is to define regulator maps for another version of motivic complexes, the polylogarithmic complexes (\[G1-2\]). Our regulator maps are defined very explicitly via the classical polylogarithm functions with some funny combinations of Bernoulli numbers serving as the coefficients.
Combining this with Beilinson’s conjecture on regulators we get, as a bonus, a precise conjecture on special values of $`L`$-functions of varieties over number fields. If the variety in question is spectrum of a number field it reduces precisely to Zagier’s conjecture. So our conjecture is in the same relationship to Beilinson’s conjecture as Zagier’s conjecture \[Z\] to the Borel theorem \[Bo\].
The regulator maps for the polylogarithmic complexes of weights $`n3`$ were constructed in \[G1-2\] and played an important role in the proofs of Zagier’s and Deninger’s conjectures on $`\zeta `$-functions at $`s=3`$ (loc. cit., \[G5\]).
2. Polylogarithmic complexes. Let $`F`$ be an arbitrary field. In \[G1-2\] we defined a complex $`\mathrm{\Gamma }(F;n)`$ of the following shape
$$_n\stackrel{\delta }{}_{n1}F^{}\stackrel{\delta }{}_{n2}\mathrm{\Lambda }^2F^{}\stackrel{\delta }{}\mathrm{}\stackrel{\delta }{}_2\mathrm{\Lambda }^{n2}F^{}\stackrel{\delta }{}\mathrm{\Lambda }^nF^{}$$
(2)
called the polylogarithmic complex. Here the group $`_n`$ sits in degree $`1`$, and the differential is of degree $`1`$. We conjectured that this complex is quasiisomorphic to the weight $`n`$ motivic complex of $`\mathrm{Spec}(F)`$, so one should have
$$H^i(\mathrm{\Gamma }(F;n))\stackrel{\mathrm{?}}{=}gr_n^\gamma K_{2ni}(F)$$
(3)
There is a good evidence this is so for small weights, see s. 1.3 below.
Now let $`F=(X)`$ be the field of rational functions on a complex algebraic variety $`X`$. Our main results are explicit formulas for the regulator maps from the polylogarithmic complexes to the Deligne complexes of $`\mathrm{Spec}((X))`$, see theorems 2.5 and 2.6. The group $`_n`$ is directly related to the properties of the classical polylogarithms. In particular there is a homomorphism
$$_n()(n1)$$
(4)
given in terms of the classical polylogarithms. Surprisingly its generalization to a homomorphism of complexes is quite complicated.
Our regulator maps enjoy compatibility with the residues property (condition (d) in theorem 2.2), which would guarantee that they can be extended from the generic point of $`X`$ to $`X`$ itself. However there is a serious difficulty in the definition of polylogarithmic complex $`\mathrm{\Gamma }(X;n)`$ for a general variety $`X`$ and $`n>3`$, (see p. 240 in \[G1\]). It would be resolved if homotopy invariance of complexes (2) will be known (see conjecture 1.39 in \[G1\]). As a result we have unconditional definition of the polylogarithmic complexes $`\mathrm{\Gamma }(X;n)`$ only in the following cases:
a) $`X=Spec(F)`$, $`F`$ is an arbitrary field.
b) $`X`$ is an regular curve over any field, and $`n`$ is arbitrary.
c) $`X`$ is an arbitrary regular scheme, but $`n3`$.
Compatibility with the residues provides the regulator map on the level of complexes in all these cases, and it would provide it in general if the mentioned above difficulty will be resolved.
3. Comparison with Beilinson’s regulator map. A homomorphism from the motivic cohomology to $`H^i(\mathrm{\Gamma }(X;n))`$ has been constructed in the following cases:
(1) $`F`$ is an arbitrary field, $`n3`$ (\[G1-2,5\]) and $`n=4`$, $`i>1`$ (to appear).
(2) $`X`$ is a curve over a number field, $`n3`$ (\[G5\]) and $`n=4`$, $`i>1`$ (to appear).
In all these cases we proved that this homomorphism followed by the regulator map on polylogarithmic complexes (when $`F=(X)`$ in (1)) coincides with Beilinson’s regulator. This proves the traditionally difficult “surjectivity” property: the image of the regulator map on these polylogarithmic complexes contains the image of Beilinson’s regulator map in the Deligne cohomology.
The results (2) combined with the results of R. de Jeu \[RdJ1-2\] prove that the image of Beilinson’s regulator map in the Deligne cohomology coincides with the image of the regulator map on polylogarithmic complexes in the case (2).
4. Remarks. We give a detailed proof of the main result. It is a rather involved but direct calculation. However the reader may skip it because there exists a completely different, more conceptual, approach to this result. We outline it in s. 2.8 to show the main theorem in a general framework. The details will be published elsewhere. This approach needs quite elaborate and rather sofisticated machinery, so the elementary proof presented in the last section might be the quickest way to check the main theorem.
Acknowledgment. The author gratefully acknowledges the support of the NSF grant DMS-9800998. I am grateful to Spencer Bloch for discussions on Arakelov motivic complexes.
I am delighted to contribute this paper to the volume telated to 20-th anniversary of Spencer Bloch’s Irvine lectures, which always have been a great source of inspiration for me.
## 2 The main result
1. Classical polylogarithms. Recall their definition:
$$Li_1(z):=\mathrm{log}(1z);Li_n(z):=_0^zLi_{n1}(t)\frac{dt}{t},n2.$$
They are multivalued analytic functions, but admit the following single-valued cousins:
$$\widehat{}_n(z):=\pi _n(\underset{k=0}{\overset{n1}{}}\beta _kLi_{nk}(z)\mathrm{log}^{nk}|z|)$$
where
$$\pi _n(a+ib)=\{\begin{array}{cc}a\hfill & n\text{ odd}\hfill \\ ib\hfill & n\text{ even}.\hfill \end{array}$$
and
$$\beta _k:=\frac{2^kB_k}{k!},\underset{k0}{}\beta _kt^k=\frac{2t}{e^{2t}1}$$
so $`\beta _{2m+1}=0`$ for $`m1`$, and
$$\beta _0=1,\beta _1=1,\beta _2=\frac{1}{3},\beta _4=\frac{1}{45},\beta _6=\frac{2}{945},\mathrm{}$$
These functions were written by Zagier \[Z\]. Their Hodge-theoretic interpretation was given by Beilinson and Deligne \[BD\]. For example for the dilogarithm it is the Bloch-Wigner function
$$\widehat{}_2(z):=i_2(z):=\pi _2\left(Li_2(z)\right)+i\mathrm{arg}(1z)\mathrm{log}|z|$$
2. The groups $`_n(F)`$ and polylogarithmic complexes(see s.1.4 in \[G2\]). For a set $`X`$ denote by $`[X]`$ the free abelian group generated by symbols $`\{x\}`$ where $`x`$ run through all elements of the set $`X`$. Let $`F`$ be an arbitrary field. We define inductively subgroups $`_n(F)`$ of $`[P_F^1]`$, $`n1`$ and set
$$_n(F):=[P_F^1]/_n(F)$$
By definition $`_1(F):=(\{x\}+\{y\}\{xy\},(x,yF^{});\{0\};\{\mathrm{}\})`$. Then $`_1(F)=F^{}`$. Let $`\{x\}_n`$ be the image of $`\{x\}`$ in $`_n(F)`$. Consider homomorphisms
$`[P_F^1]\stackrel{\delta _n}{}\{\begin{array}{ccc}_{n1}(F)F^{}\hfill & :\hfill & n3\hfill \\ \mathrm{\Lambda }^2F^{}\hfill & :\hfill & n=2\hfill \end{array}`$ (7)
$`\delta _n:\{x\}\{\begin{array}{ccc}\{x\}_{n1}x\hfill & :\hfill & n3\hfill \\ (1x)x\hfill & :\hfill & n=2\hfill \end{array}\delta _n:\{\mathrm{}\},\{0\},\{1\}0`$ (10)
Set $`𝒜_n(F):=\mathrm{Ker}\delta _n.`$ Any element $`\alpha (t)=\mathrm{\Sigma }n_i\{f_i(t)\}[P_{F(t)}^1]`$ has a specialization $`\alpha (t_0):=\mathrm{\Sigma }n_i\{f_i(t_0)\}[P_F^1]`$ at each point $`t_0P_F^1`$.
###### Definition 2.1.
$`_n(F)`$ is generated by elements $`\{\mathrm{}\},\{0\}`$ and $`\alpha (0)\alpha (1)`$ where $`\alpha (t)`$ runs through all elements of $`𝒜_n(F(t))`$.
Then $`\delta _n(_n(F))=0`$ (\[G1\], 1.16). So we get homomorphisms
$$\delta _n:_n(F)_{n1}(F)F^{},n3;\delta _2:_2(F)\mathrm{\Lambda }^2F^{}$$
and finally the mentioned in the introduction complex $`\mathrm{\Gamma }(F,n)`$:
$$_n\stackrel{\delta }{}_{n1}F^{}\stackrel{\delta }{}_{n2}\mathrm{\Lambda }^2F^{}\stackrel{\delta }{}\mathrm{}\stackrel{\delta }{}_2\mathrm{\Lambda }^{n2}F^{}\stackrel{\delta }{}\mathrm{\Lambda }^nF^{}$$
where $`\delta :\{x\}_p_{i=1}^{np}y_i\delta _p(\{x\}_p)_{i=1}^{np}y_i`$.
Let $`F=`$. Set $`\widehat{}_n(m_i\{z_i\}_n):=m_i\widehat{}_n(z_i)`$. One can prove that $`\widehat{}_n(_n()))=0`$ (\[G2\], theorem 1.13). So we are getting homomorphism (4).
3. The residue homomorphism for complexes $`\mathrm{\Gamma }(F,n)`$ (s. 1.14 in \[G1\]). Let $`F=K`$ be a field with a discrete valuation $`v`$, the residue field $`k_v`$ and the group of units $`U`$. Let $`u\overline{u}`$ be the projection $`Uk_v^{}`$. Choose a uniformizer $`\pi `$. There is a homomorphism $`\theta :\mathrm{\Lambda }^nK^{}\mathrm{\Lambda }^{n1}k_v^{}`$ uniquely defined by the following properties $`(u_iU)`$:
$$\theta (\pi u_1\mathrm{}u_{n1})=\overline{u}_1\mathrm{}\overline{u}_{n1};\theta (u_1\mathrm{}u_n)=0$$
It is clearly independent of $`\pi `$. Define a homomorphism $`s_v:[P_K^1][P_{k_v}^1]`$ by setting $`s_v\{x\}=\{\overline{x}\}\text{ if }x\text{ is a unit}`$ and $`0`$ otherwise. It induces a homomorphism $`s_v:_m(K)_m(k_v)`$. Put
$$_v:=s_v\theta :_m(K)\mathrm{\Lambda }^{nm}K^{}_m(k_v)\mathrm{\Lambda }^{nm1}k_v^{}$$
It defines a morphism of complexes $`_v:\mathrm{\Gamma }(K,n)\mathrm{\Gamma }(k_v,n1)[1]`$.
4. Main result: a preliminary form. Let $`𝒜^i(\eta _X)`$ be the space of real smooth $`i`$-forms at the generic point $`\eta _X:=\mathrm{Spec}(X)`$ of a complex variety $`X`$. Denote by $`X^{(1)}`$ the set of the codimension one closed irreducible subvarieties in $`X`$. Let $`𝒟`$ be the de Rham differential on distributions, and $`d`$ the de Rham differential on $`𝒜^i(\eta _X)`$. A typical example:
$$d(di\mathrm{arg}z)=0;𝒟(di\mathrm{arg}z)=2\pi i\delta (z)$$
(11)
The difference $`𝒟d`$ is the de Rham residue homomorphism. It is defined on distributions smooth at the generic point of $`X()`$. Its value on such a distribution is concentrated on a union of codimension one subvarieties.
###### Theorem 2.2.
There exist a homomorphism of complexes
$$\begin{array}{ccccccc}_n((X))& \stackrel{\delta }{}& _{n1}((X))(X)^{}& \stackrel{\delta }{}& \mathrm{}& \stackrel{\delta }{}& ^n(X)^{}\\ & & & & & & \\ r_n(1)& & r_n(2)& & & & r_n(n)\\ & & & & & & \\ 𝒜^0(\eta _X)(n1)& \stackrel{d}{}& 𝒜^1(\eta _X)(n1)& \stackrel{d}{}& \mathrm{}& \stackrel{d}{}& 𝒜^{n1}(\eta _X)(n1)\end{array}$$
such that
(a) $`r_n(1)(\{f\}_n)=\widehat{}_n(f)`$.
(b) $`dr_n(n)(f_1\mathrm{}f_n)+\pi _n(d\mathrm{log}f_1\mathrm{}d\mathrm{log}f_n)=0`$.
(c) The differential form $`r_n(m)()`$ defines a distribution on $`X()`$.
(d) The homomorphism $`r_n()`$ is compatible with residues:
$$𝒟r_n(m)r_n(m+1)\delta =2\pi i\underset{YX^{(1)}}{}r_{n1}(m1)_{v_Y},m<n$$
$$𝒟r_n(n)\pi _n(d\mathrm{log}f_1\mathrm{}d\mathrm{log}f_n)=2\pi i\underset{YX^{(1)}}{}r_{n1}(n1)_{v_Y}$$
where $`v_Y`$ is the valuation on the field $`(X)`$ defined by a divisor $`Y`$.
Remark. This result has been formulated in \[G4\], see theorem 4.3.
The part (d) means that $`r_n()`$ sends the residue homomorphism $`_{v_Y}`$ to the de Rham residue homomorphism.
Here are two examples. Set
$$\alpha (f,g):=\mathrm{log}|f|d\mathrm{log}|g|+\mathrm{log}|g|d\mathrm{log}|f|$$
Example 1: n=3. Set
$`r_3(2):\{f\}_2g\widehat{}_2(f)di\mathrm{arg}g{\displaystyle \frac{1}{3}}\mathrm{log}|g|\alpha (1f,f)`$
Example 2: n=4.
$`r_4(2):\{f\}_3g\widehat{}_3(f)di\mathrm{arg}g{\displaystyle \frac{1}{3}}\widehat{}_2(f)\mathrm{log}|g|d\mathrm{log}|f|`$
$`r_4(3):\{f\}_2g_1g_2+\widehat{}_2(f)di\mathrm{arg}g_1di\mathrm{arg}g_2{\displaystyle \frac{1}{3}}\alpha (1f,f)`$
$`\left(\mathrm{log}|g_1|di\mathrm{arg}|g_2|\mathrm{log}|g_2|di\mathrm{arg}|g_1|\right)+{\displaystyle \frac{1}{3}}\widehat{}_2(f)d\mathrm{log}|g_1|d\mathrm{log}|g_2|`$
5. The numbers $`\beta _{k,p}`$. Define for any integers $`p1`$ and $`k0`$ the numbers
$$\beta _{k,p}:=(1)^p(p1)!\underset{0i[\frac{p1}{2}]}{}\frac{1}{(2i+1)!}\beta _{k+p2i}$$
For instance
$$\beta _{k,1}=\beta _{k+1};\beta _{k,2}=\beta _{k+2};\beta _{k,3}=2(\frac{1}{1!}\beta _{k+3}+\frac{1}{3!}\beta _{k+1});$$
$$\beta _{k,4}=3!(\frac{1}{1!}\beta _{k+4}+\frac{1}{3!}\beta _{k+2});\beta _{k,5}=4!(\frac{1}{1!}\beta _{k+5}+\frac{1}{3!}\beta _{k+3}+\frac{1}{5!}\beta _{k+1})$$
One has recursions
$$2p\beta _{k+1,2p}=\beta _{k,2p+1}\frac{1}{2p+1}\beta _{k+1};(2p1)\beta _{k+1,2p1}=\beta _{k,2p}$$
(12)
These recursions together with $`\beta _{k,1}=\beta _{k+1}`$ determine the numbers $`\beta _{k,p}`$.
###### Lemma 2.3.
Let $`m1`$. Then
$$\beta _{0,2m}=\beta _{0,2m+1}=\frac{1}{2m+1}$$
(13)
$$\beta _{1,2m1}=\frac{1}{(2m1)(2m+1)},\beta _{1,2m}=0$$
(14)
Proof. Let us prove formula
$$\beta _{0,2m}:=(2m1)!\left(\frac{1}{1!}\beta _{2m}+\frac{1}{3!}\beta _{2m2}+\mathrm{}+\frac{1}{(2m1)!}\beta _2\right)=\frac{1}{2m+1}$$
(15)
Since $`\beta _0=1`$, $`\beta _1=1`$ and $`\beta _{2m+1}=0`$ for $`m>0`$ this formula, as one easily checks, is equivalent to the following one:
$$\underset{p1}{}\frac{1}{p!}\beta _{2m+1p}=0$$
The left hand side is the coefficient in $`t^{2m+1}`$ of the power series
$$(e^t1)\frac{2t}{e^{2t}1}=\frac{2t}{e^t+1}$$
The right hand side of the last equation is an almost even function in $`t`$: if we denote it by $`F(t)`$ then $`F(t)=F(t)+2t`$.
Formula (15) immediately implies that
$$\beta _{1,2m1}:=(2m2)!\left(\frac{1}{1!}\beta _{2m}+\frac{1}{3!}\beta _{2m2}+\mathrm{}+\frac{1}{(2m1)!}\beta _2\right)=\frac{1}{(2m1)(2m+1)}$$
Using $`\beta _{2m+1}=0`$ for $`m>0`$ and $`\beta _1=1`$ we have
$$\beta _{0,2m+1}=(2m)!\left(\frac{1}{1!}\beta _{2m+1}+\frac{1}{3!}\beta _{2m1}+\mathrm{}+\frac{1}{(2m1)!}\beta _1+\frac{1}{(2m+1)!}\beta _1\right)=\frac{1}{2m+1}$$
and
$$\beta _{1,2m}=(2m1)!\left(\frac{1}{1!}\beta _{2m+1}+\frac{1}{3!}\beta _{2m1}+\mathrm{}+\frac{1}{(2m1)!}\beta _3\right)=0$$
The lemma is proved.
###### Proposition 2.4.
For any $`p1`$ one has
$$\beta _{n2,p+1}(n1)\beta _{n1,p}\underset{k=1}{\overset{n3}{}}\beta _{k,p}\beta _{nk1}=0$$
(16)
Proof. We will do it by induction on $`p`$ using the recursion relations (12). If $`p=1`$ the formula we need to prove boils down to the identity
$$\underset{k=2}{\overset{n2}{}}\beta _i\beta _{ni}+n\beta _n=0$$
(17)
which is easy to check using the generating functions. Let us denote by $`()_{n,p}`$ the left hand side of (16). Then $`()_{n,p}+()_{n+1,p1}=0`$. Indeed, one has
$$()_{n,p}+()_{n+1,p1}=\beta _{n2,p+1}+(p1)\beta _{n1,p}$$
(18)
$$(n1)\beta _{n1,p}n(p1)\beta _{n,p1}\underset{k=1}{\overset{n3}{}}\beta _{k,p}\beta _{nk1}(p1)\underset{k=0}{\overset{n3}{}}\beta _{k+1,p1}\beta _{nk1}$$
Let us assume first that $`p`$ is even. Using recursions (12) we write (18) as
$$\frac{1}{p+1}\beta _{n1}\beta _{n1,p}(p1)\beta _{n,p1}(p1)\beta _{1,p1}\beta _{n1}$$
To prove that it is zero we use again recursions as well as the first formula in (14). Now let $`p`$ be an odd number. Then the recursion relations show that (18) equals to
$$\beta _{n1,p}(p1)\beta _{n,p1}+\frac{n1}{p}\beta _n+\frac{1}{p}\underset{k=1}{\overset{n3}{}}\beta _{k+1}\beta _{nk1}(p1)\beta _{1,p1}\beta _{n1}$$
Using the identity (18) together with the recursions and the second formula in (14) we see that this expression is also equal to zero. The proposition is proved.
6. Construction of the homomorphism $`r_n()`$. Let us define differential $`1`$-forms $`\widehat{}_{p,q}`$ on $`P^1\backslash \{0,1,\mathrm{}\}`$ for $`q1`$ as follows:
$$\widehat{}_{p,q}(z):=\widehat{}_p(z)\mathrm{log}^{q1}|z|d\mathrm{log}|z|,p2$$
(19)
$$\widehat{}_{1,q}(z):=\alpha (1z,z)\mathrm{log}^{q1}|z|$$
It provides a distribution on $`P^1`$. Moreover for any rational function $`f`$ on a complex variety $`X`$ the $`1`$-form $`\widehat{}_{p,q}(f)`$ provides a distribution on $`X()`$.
A useful notation. Set
$$𝒜_m\left\{\underset{i=1}{\overset{2p}{}}d\mathrm{log}|g_i|\underset{i=2p+1}{\overset{m}{}}di\mathrm{arg}g_j\right\}:=$$
$$\mathrm{Alt}_m\left\{\frac{1}{(2p)!(m2p)!}\underset{i=1}{\overset{2p}{}}d\mathrm{log}|g_i|\underset{i=2p+1}{\overset{m}{}}di\mathrm{arg}g_j\right\}$$
and
$$𝒜_m\left\{\mathrm{log}|g_1|\underset{i=2}{\overset{p}{}}d\mathrm{log}|g_i|\underset{i=p+1}{\overset{m}{}}di\mathrm{arg}g_j\right\}:=$$
$$\mathrm{Alt}_m\left\{\frac{1}{(p1)!(mp)!}\mathrm{log}|g_1|\underset{i=2}{\overset{p}{}}d\mathrm{log}|g_i|\underset{i=p+1}{\overset{m}{}}di\mathrm{arg}g_j\right\}$$
So $`𝒜_m(F(g_1,..,g_m))`$ is a weighted alternation (we divide by the order of the stabilizer of the term we alternate).
Now we are ready for the precise formulation of our main result.
###### Theorem - Construction 2.5.
Let $`f,g_1,\mathrm{},g_m`$ be rational functions on a complex variety $`X`$. Then the following formula provides maps satisfying all the properties of theorem 2.2.
$$r_{n+m}(m+1):\{f\}_ng_1\mathrm{}g_m$$
$$\widehat{}_n(f)𝒜_m\left\{\underset{p0}{}\frac{1}{2p+1}\underset{i=1}{\overset{2p}{}}d\mathrm{log}|g_i|\underset{j=2p+1}{\overset{m}{}}di\mathrm{arg}g_j\right\}+$$
(20)
$$\underset{k1}{}\underset{1pm}{}\beta _{k,p}\widehat{}_{nk,k}(f)𝒜_m\left\{\mathrm{log}|g_1|\underset{i=2}{\overset{p}{}}d\mathrm{log}|g_i|\underset{j=p+1}{\overset{m}{}}di\mathrm{arg}g_j\right\}$$
(21)
Here are several examples.
Example 1. $`m=1`$, $`n`$ is arbitrary.
$$\{f\}_ng\widehat{}_n(f)di\mathrm{arg}g\underset{k=1}{\overset{n1}{}}\beta _{k+1}\widehat{}_{nk,k}(f)\mathrm{log}|g|$$
Example 2. $`m=2`$, $`n`$ is arbitrary.
$$\{f\}_ng_1g_2\widehat{}_n(f)\left\{di\mathrm{arg}g_1di\mathrm{arg}g_2+\frac{1}{3}d\mathrm{log}|g_1|d\mathrm{log}|g_2|\right\}$$
$$\underset{k=1}{\overset{n1}{}}\beta _{k+1}\widehat{}_{nk,k}(f)(\mathrm{log}|g_1|di\mathrm{arg}g_2\mathrm{log}|g_2|di\mathrm{arg}g_1)$$
$$+\underset{k1}{}\beta _{k+2}\widehat{}_{nk,k}(f)(\mathrm{log}|g_1|d\mathrm{log}|g_2|\mathrm{log}|g_2|d\mathrm{log}|g_1|)$$
Example 3. The homomorphism $`r_5()`$.
$$\{f\}_4g\widehat{}_4(f)di\mathrm{arg}g\frac{1}{3}\widehat{}_3(f)d\mathrm{log}|f|\mathrm{log}|g|+\frac{1}{45}\alpha (1f,f)\mathrm{log}^2|f|\mathrm{log}|g|$$
$$\{f\}_3g_1g_2\widehat{}_3(f)\{di\mathrm{arg}g_1di\mathrm{arg}g_2+\frac{1}{3}d\mathrm{log}|g_1|d\mathrm{log}|g_2|\}$$
$$\frac{1}{3}\widehat{}_2(f)d\mathrm{log}|f|\left(\mathrm{log}|g_1|di\mathrm{arg}g_2\mathrm{log}|g_2|di\mathrm{arg}g_1\right)$$
$$\frac{1}{45}\alpha (1f,f)\mathrm{log}|f|\left(\mathrm{log}|g_1|d\mathrm{log}|g_2|\mathrm{log}|g_2|d\mathrm{log}|g_1|\right)$$
$$\{f\}_2g_1g_2g_3$$
$$\widehat{}_2(f)𝒜_3\left(di\mathrm{arg}g_1di\mathrm{arg}g_2di\mathrm{arg}g_3+\frac{1}{3}d\mathrm{log}|g_1|d\mathrm{log}|g_2|di\mathrm{arg}g_3\right)$$
$$\alpha (1f,f)𝒜_3\left(\frac{1}{3}\mathrm{log}|g_1|di\mathrm{arg}g_2di\mathrm{arg}g_3+\frac{1}{15}\mathrm{log}|g_1|d\mathrm{log}|g_2|d\mathrm{log}|g_3|\right)$$
Remark. The morphism of complexes $`r_n()`$ is not defined uniquely by its properties if $`n>3`$. For example if $`n=4`$ we can have a map homotopic to our regulator map by using the homotopy which is given by the homomorphism
$$_2((X))\mathrm{\Lambda }^2(X)^{}𝒜^0(\eta _X),\{f\}_3g\widehat{}_2(f)\mathrm{log}|f|\mathrm{log}|g|$$
(the other components of the homotopy are zero). Indeed, the residue map for it is zero, and it takes values in $`(2)`$.
7. Theorem 2.5 from the point of view of the Deligne complex. Recall that a $`p`$-distribution on a manifold $`X`$ is a linear continuous functional on the space of $`(dim_{}Xp)`$-forms with compact support. Denote by $`𝒟_X^p`$ the space of all real $`p`$-distributions on $`X`$.
Let $`X`$ be a regular variety over $``$. The $`n`$-th Beilinson-Deligne complex $`\underset{¯}{}_𝒟(n)_X`$ can be defined as a total complex associated with the following bicomplex of sheaves in classical topology on $`X()`$:
$$\begin{array}{ccccccccccc}(𝒟_X^0& \stackrel{d}{}& 𝒟_X^1& \stackrel{d}{}& \mathrm{}& \stackrel{d}{}& 𝒟_X^n& \stackrel{d}{}& 𝒟_X^{n+1}& \stackrel{d}{}& \mathrm{})(n1)\\ & & & & & & & & & & \\ & & & & & & \pi _n& & \pi _n& & \\ & & & & & & & & & & \\ & & & & & & \mathrm{\Omega }_{X,\mathrm{log}}^n& \stackrel{}{}& \mathrm{\Omega }_{X,\mathrm{log}}^{n+1}& \stackrel{}{}& \end{array}$$
Here $`𝒟_X^0`$ placed in degree 1 and $`(\mathrm{\Omega }_{X,\mathrm{log}}^{},)`$ is the de Rham complex of holomorphic forms with logarithmic singularities at infinity. We will denote by $`\underset{¯}{}_𝒟(n)(U)`$ the complex of the global sections.
Theorem 2.2 can be reformulated as follows
###### Theorem 2.6.
Let $`X`$ be a complex algebraic variety. Set $`\stackrel{~}{r}_n(i):=r_n(i)`$ for $`i<n`$ and let
$$\stackrel{~}{r}_n(n):\mathrm{\Lambda }^n(X)^{}𝒜^{n1}(\eta _X)(n1)\mathrm{\Omega }_{\mathrm{log}}^n(\eta _X)$$
$$f_1\mathrm{}f_nr_n(n)(f_1\mathrm{}f_n)+d\mathrm{log}f_1\mathrm{}d\mathrm{log}f_n$$
(22)
Then we get a homomorphism of complexes
$$\stackrel{~}{r}_n():\mathrm{\Gamma }((X);n)\underset{¯}{}_𝒟(n)(\eta _X)$$
(23)
compatible with the residues as explained in the part (d) of theorem 2.2.
Indeed, condition b) of theorem 2.2 just means that the right hand side of (22) is a cycle in the Deligne complex $`\underset{¯}{}(n)_𝒟(\eta _X)`$.
There is a natural Dolbeault resolution of the complex $`\mathrm{\Omega }_{\mathrm{log}}^n`$. Using it in the complex of sheaves $`\underset{¯}{}_𝒟(n)_X`$ to replace the subcomplex $`\mathrm{\Omega }_{\mathrm{log}}^n`$ we get a complex $`\underset{¯}{}_𝒟^{}(n)_X`$ of fine sheafs on $`X()`$. The property (d) would allow us to extend the homomorphism $`\stackrel{~}{r}_n()`$ to a morphism of complexes
$$\stackrel{~}{r}_n():\mathrm{\Gamma }(X;n)\underset{¯}{}_𝒟^{}(n)(X())$$
For small weights this was explained in detail in \[G2\].
If $`X`$ is a variety over $``$, then
$$H_𝒟^i(X_{};(n))=H_𝒟^i(X(),(n))^{\overline{F}_{\mathrm{}}}$$
where $`\overline{F}_{\mathrm{}}`$ is the de Rham involution, i.e. the composition of the involution $`F_{\mathrm{}}`$ on $`X()`$ induced by the complex conjugation with the complex conjugation of coefficients. If $`X`$ is a variety over $``$ then the regulator map is defined as
$$H^i\mathrm{\Gamma }(X;n)H_𝒟^i(X_{}/_{};(n))$$
(24)
###### Conjecture 2.7.
The image of the regulator map (24) in the Deligne cohomology coincides with the image of Beilinson’s regulator map.
If $`X=\mathrm{Spec}(F)`$ where $`F`$ is a number field then the only nontrivial case is $`i=1`$, and we get a version of Zagier’s conjecture. The next case is when $`X`$ is a curve over number field, and $`i=2`$. Then we come to conjecture 1.5 in \[G5\], see also theorem 4.4 in \[G5\] and s. 7 in \[G7\].
There is a version of conjecture 2.7 expressing the special value of the corresponding $`L`$-function at $`s=n`$ via the regulator map on the polylogarithmic complex. Its specialization for elliptic curves is conjecture 1.10 in \[G7\], see also theorem 4.7 in \[G5\]. For the $`n=2`$ it is a theorem of Bloch for elliptic CM curves and of Beilinson for elliptic curves over $``$, for $`n=3`$ it was conjectured by Deninger and proved by the author in \[G5\] for elliptic curves over $``$, and it is also proved now by the author for $`n=4`$ for elliptic curves over $``$ (to appear).
8. Generalizations. According to the Tannakian formalism the category of mixed $``$-Hodge-Tate structures is equivalent to the category of graded comodules over a certain Lie coalgebra $`_{}()`$ positively graded by the weights. There exists a natural homomorphism of groups
$$_n()_n()$$
provided by the Hodge realization of the polylogarithmic motive. One has a homomorphism of $``$-vector spaces
$$p_n:_n()(n1)$$
called the Lie-period map, see \[D\] for its construction. The composition
$$_n()_n()(n1)$$
coincides with the homomorphism $`\widehat{}_n`$, see \[BD\].
We generalize this picture considering variations of mixed Hodge-Tate structures over complex varieties and working on the level of complexes.
First of all, the category of variations of mixed $``$-Hodge-Tate structures over $`\eta _X=\mathrm{Spec}(X)`$, where $`X`$ is a complex algebraic variety, is equivalent to the category of graded comodules over a Lie coalgebra $`_{}(\eta _X)`$.
The Hodge realization functor provides canonical homomorphism
$$l_n:_n((X))_n(\eta _X)$$
(25)
Applying it componentwise we get a morphism of complexes (\[G1-2\])
$$l_{nk}^kl_1:_{nk}((X))\mathrm{\Lambda }^k(X)^{}_{nk}(\eta _X)\mathrm{\Lambda }^k_1(\eta _X)$$
The vector space on the right is a subspace in $`\left(\mathrm{\Lambda }^{k+1}_{}(\eta _X)\right)_n`$ where the subscript $`n`$ means the weight $`n`$ part. The standard cochain complex of a Lie coalgebra $`_{}`$ looks as follows
$$_{}(\eta _X)\mathrm{\Lambda }^2_{}(\eta _X)\mathrm{\Lambda }^3_{}(\eta _X)\mathrm{}$$
(26)
It is a complex of graded $``$-vector spaces. We get a morphism of complexes
$$\begin{array}{ccccccc}_n((X))& \stackrel{\delta }{}& _{n1}((X))(X)^{}& \stackrel{\delta }{}& \mathrm{}& \stackrel{\delta }{}& ^n(X)^{}\\ & & & & & & \\ & & & & & & \\ & & & & & & \\ _n(\eta _X)& \stackrel{}{}& \left(\mathrm{\Lambda }^2_{}(\eta _X)\right)_n& \stackrel{}{}& \mathrm{}& \stackrel{}{}& \mathrm{\Lambda }^n_1(\eta _X)\end{array}$$
The bottom line of this diagram is the degree $`n`$ part of the complex (26).
The Lie-period provides a natural homomorphism of $``$-vector spaces
$$p_n:_n(\eta _X)𝒜^0(\eta _X)(n1)$$
(27)
The following result tells us that it is a beginning of a morphism of complexes.
###### Theorem 2.8.
There exists a morphism of complexes
$$\text{the weight }n\text{ part of (}\text{26}\text{)}\text{the weight }n\text{ Deligne complex }\underset{¯}{}_𝒟(n)(\eta _X)$$
(28)
whose degree $`1`$ component $`_{}(\eta _X)𝒜^0(\eta _X)`$ coincides with the map (27).
A detailed account on this result and related issues would double the size of this paper. We will do it in a different place.
Combining homomorphism (28) with the previous map of complexes
$$\mathrm{\Gamma }(\eta _X;n)\text{the weight }n\text{ part of (}\text{26}\text{)}$$
we get a homomorphism of complexes from theorem 2.6.
## 3 Arakelov motivic complexes: examples
1. Another version of Deligne’s complexes. We need a complex introduced twenty years ago by Deligne and quasiisomorphic to the complex $`\underset{¯}{}_𝒟^{}(n)(X())`$.
Let $`𝒟_X^{p,q}=𝒟^{p,q}`$ be the abelian group of complex valued distributions of type $`(p,q)`$ on $`X()`$. Consider the following cohomological bicomplex, where the group $`𝒟^{0,0}`$ is placed in degree $`1`$:
$$\begin{array}{ccccccccc}& & & & & & & & 𝒟_{cl}^{n,n}\\ & & & & & & & & \\ & & & & & & & 2\overline{}& \\ & & & & & & & & \\ 𝒟^{0,n1}& \stackrel{}{}& 𝒟^{1,n1}& \stackrel{}{}& \mathrm{}& \stackrel{}{}& 𝒟^{n1,n1}& & \\ & & & & & & & & \\ \overline{}& & \overline{}& & & & \overline{}& & \\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& & \\ \overline{}& & \overline{}& & & & \overline{}& & \\ & & & & & & & & \\ 𝒟^{0,1}& \stackrel{}{}& 𝒟^{1,1}& \stackrel{}{}& \mathrm{}& \stackrel{}{}& 𝒟^{n1,1}& & \\ & & & & & & & & \\ \overline{}& & \overline{}& & & & \overline{}& & \\ & & & & & & & & \\ 𝒟^{0,0}& \stackrel{}{}& 𝒟^{1,0}& \stackrel{}{}& \mathrm{}& \stackrel{}{}& 𝒟^{n1,0}& & \end{array}$$
Let $`Tot^{}`$ be the total complex of this bicomplex. It is concentrated in degrees $`[1,2n]`$. The complex $`C_𝒟^{}(X;(n))`$ is subcomplex of $`Tot^{}`$ defined as follows. Intersect the part of the complex $`Tot^{}`$ coming from the $`n\times n`$ square in the diagram with the complex of $`(n1)`$-valued distributions. Consider the subgroup $`𝒟_{cl,}^{n,n}(n)𝒟_{cl}^{n,n}`$ of the $`(n)`$-valued distributions of type $`(n,n)`$. They form a subcomplex in $`Tot^{}`$ because $`\overline{}`$ sends $`(n1)`$-valued distributions to $`(n)`$-valued distributions. This is the complex $`C_𝒟^{}(X;(n))`$ with the differential denoted by $`D`$. It is a truncation of the complex considered by Deligne.
###### Proposition 3.1.
The complex $`C_𝒟^{}(X;n)`$ is quasiisomorphic to the truncated Deligne complex $`\tau _{2n}\underset{¯}{}_𝒟^{}(n)(X())`$.
Proof. See \[C\] or \[G3\].
Now if $`X`$ is a variety over $``$, then
$$C_𝒟^{}(X_{};(n)):=C_𝒟^{}(X,n)^{\overline{F}_{\mathrm{}}}H_𝒟^i(X_{};(n))=H^i\left(C_𝒟^{}(X_{};(n))\right)$$
2. The Arakelov motivic complexes. By definition the weight $`n`$ Arakelov motivic complex $`\mathrm{\Gamma }^𝒜(X;n)`$ is defined as follows
$$\mathrm{R}\mathrm{\Gamma }_X^𝒜(n):=\mathrm{Cone}\left(\mathrm{R}\mathrm{\Gamma }_X(n)\stackrel{\mathrm{Reg}}{}𝒞_D(X_{},;(n))\right)[1]$$
where $`\mathrm{R}\mathrm{\Gamma }_X(n)`$ is the weight $`n`$ motivic complex and $`\mathrm{Reg}`$ is the regulator map. It is an object of the derived category. When $`X`$ is regular scheme over $``$ then
$$\mathrm{R}\mathrm{\Gamma }_X_{}^𝒜(n):=\mathrm{Cone}\left(\mathrm{R}\mathrm{\Gamma }_X(n)\stackrel{\mathrm{Reg}}{}𝒞_D(X_{};(n))\right)[1]$$
Remark. For the application to special values of $`L`$-functions we should take a model of $`X`$ over $``$ and have a similar construction for finite primes $`p`$ where it has semistable reduction.
Below we discuss complexes $`\mathrm{\Gamma }^𝒜(X;n)`$ with the polylogarithmic complex $`\mathrm{\Gamma }(X;n)`$ serving as the motivic complex $`\mathrm{R}\mathrm{\Gamma }_X(n)`$. We will abuse notations by writing $`(𝒟^{p,q}𝒟^{q,p})_{}(k)`$ for the subgroup of $`(k)`$-valued currents in $`𝒟^{p,q}𝒟^{q,p}`$. The last groups of the complex $`\mathrm{\Gamma }^𝒜(X;n)`$ look as follows (see also s. 3.6 below):
$$\begin{array}{cccccc}\mathrm{}& _{Y_2X^{(n2)}}K_2(Y_2)& & _{Y_1X^{(n1)}}(Y_1)^{}& & _{YX^{(n)}}\\ & & & & & \\ \mathrm{}& r_n(n2)& & r_n(n1)& & r_n(n)\\ & & & & & \\ \mathrm{}& (𝒟_{}^{n2,n1}𝒟^{n1,n2})_{}(n1)& \stackrel{+\overline{}}{}& 𝒟_{}^{n1,n1}(n1)& \stackrel{2\overline{}}{}& 𝒟_{cl,}^{n,n}(n)\end{array}$$
The $`(2n)`$-th cohomology group of this complex looks as follows:
$$\frac{\{\text{codim. }n\text{ cycle}Y,g𝒟_{cl,}^{n1,n1}(n1)\}\text{such that}2\overline{}g+\delta (Y)=0}{\{\mathrm{div}(f),\mathrm{log}|f|\delta _{Y_1}\},f(Y_1)^{};\{0;u+\overline{}v\}}$$
(29)
where
$$u𝒟^{n2,n1};v𝒟^{n1,n2}$$
This is the group of codimension $`n`$ Arakelov cycles on $`X()`$. To compare with \[S\] notice that
$$dd^{}=\frac{1}{2\pi i}\overline{},2\overline{}\mathrm{log}|f|=\delta _{\mathrm{div}(f)}$$
If we modify this definition by replacing the last group $`𝒟_{cl,}^{n,n}(n)`$ by its quotient by smooth forms $`𝒟_{cl,}^{n,n}(n)/𝒜_{cl,}^{n,n}(n)`$ we get the group of codimension $`n`$ Arakelov cycles defined by Gillet and Soulé \[S\]:
$$\frac{\{YX^{(n)},g𝒟_{cl,}^{n1,n1}(n1)\}\text{such that}2\overline{}g+\delta (Y)\text{is smooth}}{\{\mathrm{div}(f),\mathrm{log}|f|\delta _{Y_1}\},f(Y_1)^{};\{0;u+\overline{}v\}}$$
(30)
Such a modification does not seem natural in our context, but it worked well in the Gillet-Soulé theory. The reason to have several versions of Arakelov groups is the desire to have a theory of Chern classes with values in Arakelov motivic cohomology (and finally the higher Riemann-Roch theorem) for various versions of the category of vector bundles with some kind of hermitian metric discussed below.
a) Consider holomorphic vector bundles with flat hermitian metrics. They have Chern classes with values in the groups (29). In particular holomorphic line bundles with flat hermitian metrics form an abelian group under the tensor product, denoted $`\stackrel{~}{\mathrm{Pic}}_0(X)`$, which sits in an exact sequence
$$0^{\pi _0(X())}\stackrel{~}{\mathrm{Pic}}_0(X)\mathrm{Pic}_0(X)0$$
Indeed, a flat hermitian hermitian metric in a line bundle over a connected manifold is determined up to a constant. There is an isomorphism
$$\stackrel{~}{\mathrm{Pic}}_0(X)\stackrel{=}{}H^2\left(\mathrm{\Gamma }^𝒜(X;1)\right)$$
given by the first Chern class as follows. If $`s`$ is a section of a holomorphic line bundle over $`X`$ with a hermitian metric $`||||`$ then the pair $`(\mathrm{log}s,\mathrm{div}(s))`$ provides the corresponding class in $`H^2\left(\mathrm{\Gamma }^𝒜(X;1)\right)`$. Indeed, according to the Poincaré-Lelong formula for a meromorphic section $`s`$ of any holomorphic line bundle one has
$$2\overline{}\mathrm{log}||s||\delta _{\mathrm{div}(s)}=c_1(L,||||)$$
where on the right is the Chern form related to the hermitian structure on the line bundle, which is zero if and only if the metric is flat.
b) To have Chern classes for holomorphic vector bundles with arbitrary hermitian metrics we are more or less forced to the Gillet-Soulé definition. However in this case the groups of Arakelov cycles are infinite dimensional.
c) Choose a Kahler metric on $`X()`$ and consider only such metrics (harmonic metrics) on holomorphic vector bundles over $`X()`$ whose Chern forms are harmonic with respect to the choosen metric. Then we have Chern classes yet to another modification of the Arakelov group where $`𝒟_{cl,}^{n,n}(n)`$ replaced by its quotient by the image of the following map provided by the Hodge theory:
$$H^{2n}(X(),(n))𝒜_{cl,}^{n,n}(n)$$
3. The weight one. The regulator map on the weight one motivic complex looks as follows:
$$\begin{array}{ccc}(X)^{}& & _{YX^{(1)}}\\ & & \\ r_1(1)& & r_1(2)\\ & & \\ 𝒟^{0,0}& \stackrel{2\overline{}}{}& 𝒟_{cl,}^{1,1}(1)\end{array}$$
$$r_1(2):Y_12\pi i\delta _{Y_1},r_1(1):f\mathrm{log}|f|$$
Here the top line is the weight $`1`$ motivic complex, sitting in degrees $`[1,2]`$. Shifting by $`1`$ the total complex associated with this bicomplex we get the weight $`1`$ Arakelov motivic complex $`\mathrm{\Gamma }^𝒜(X;1)`$.
Here are examples of the regulator maps for the weights $`n3`$.
4. The weight two. The regulator map on the weight two motivic complexes looks as follows.
$$\begin{array}{ccccccc}_2((X))& \stackrel{\delta }{}& \mathrm{\Lambda }^2(X)^{}& \stackrel{}{}& _{YX^{(1)}}(Y)^{}& \stackrel{}{}& _{YX^{(2)}}\\ & & & & & & \\ r_2(1)& & r_2(2)& & r_2(3)& & r_2(4)\\ & & & & & & \\ 𝒟_{}^{0,0}(1)& \stackrel{D}{}& (𝒟^{0,1}𝒟^{1,0})_{}(1)& \stackrel{D}{}& 𝒟_{}^{1,1}(1)& \stackrel{2\overline{}}{}& 𝒟^{2,2}(2)_{cl,}\end{array}$$
where we set
$$r_2(4):Y_2(2\pi i)^2\delta _{Y_2}$$
$$r_2(3):(Y_1,f)2\pi i\mathrm{log}|f|\delta _{Y_1}$$
$$r_2(2):fg\mathrm{log}|f|di\mathrm{arg}g+\mathrm{log}|g|di\mathrm{arg}f$$
$$r_2(1):\{f\}_2\widehat{}_2(f)$$
To prove that we get a morphism of complexes we use theorem 2.2. The following argument is needed to check the commutativity of the second square. The de Rham differential of the distribution $`r_2(2)(fg)`$ is
$$𝒟\left(\mathrm{log}|f|di\mathrm{arg}g+\mathrm{log}|g|di\mathrm{arg}f\right)=$$
$$\pi _2(d\mathrm{log}f\mathrm{log}g)+2\pi i(\mathrm{log}|g|\delta (f)\mathrm{log}|f|\delta (g))$$
This does not coincide with $`r_2(3)(fg)`$, but the difference is
$$(𝒟r_2(2)r_2(3))(fg)=\pi _2(d\mathrm{log}f\mathrm{log}g)(𝒟^{0,2}𝒟^{2,0})_{}(1)$$
Defining the differential $`D`$ on the second group of the complex $`𝒞_𝒟(X_{},(2))`$ we take the de Rham differential and throw away from it precisely these components. Therefore the middle square is commutative.
5. The weight three. The weight three motivic complex $`\mathrm{\Gamma }(X;3)`$ is the total complex of the following bicomplex: (the first group is in degree $`1`$)
$$\begin{array}{ccccc}_3((X))& & _2((X))(X)^{}& & \mathrm{\Lambda }^3(X)^{}\\ & & & & \\ & & _{Y_1X^{(1)}}_2((Y_1))& & _{Y_1X^{(1)}}\mathrm{\Lambda }^2(Y_1)^{}\\ & & & & \\ & & & & _{Y_2X^{(2)}}(Y_2)^{}\\ & & & & \\ & & & & _{Y_3X^{(3)}}(Y_3)^{}\end{array}$$
The Deligne complex $`𝒞_𝒟(X_{},(3))`$ looks as follows:
$$\begin{array}{ccccccc}& & & & & & 𝒟_{cl,}^{3,3}(3)\\ & & & & & 2\overline{}& \\ 𝒟^{0,2}& \stackrel{}{}& 𝒟^{1,2}& \stackrel{}{}& 𝒟^{2,2}& & \\ \overline{}& & \overline{}& & \overline{}& & \\ 𝒟^{0,1}& \stackrel{}{}& 𝒟^{1,1}& \stackrel{}{}& 𝒟^{1,2}& & \\ \overline{}& & \overline{}& & \overline{}& & \\ 𝒟^{0,0}& \stackrel{}{}& 𝒟^{1,0}& \stackrel{}{}& 𝒟^{2,0}& & \end{array}$$
(Recall that the $`3\times 3`$ square in this diagram consists of $`(2)`$-valued distributions). We construct the regulator map from the motivic complex $`\mathrm{\Gamma }(X;3)`$ to the Deligne complex $`𝒞_𝒟(X_{},(3))`$ by setting
$$r_3(6):Y_3(2\pi i)^3\delta _{Y_3}$$
$$r_3(5):(Y_2,f)(2\pi i)^2\mathrm{log}|f|\delta _{Y_2}$$
$$r_3(4):(Y_1,fg)2\pi i(\mathrm{log}|f|di\mathrm{arg}g+\mathrm{log}|g|di\mathrm{arg}f)\delta _{Y_1}$$
$$r_3(3):(Y_1,\{f\}_2)2\pi i\widehat{}_2(f)\delta _{Y_1}$$
$$r_3(3):f_1f_2f_3\mathrm{Alt}_3\left(\frac{1}{6}\mathrm{log}|f_1|d\mathrm{log}|f_2|d\mathrm{log}|f_3|+\frac{1}{2}\mathrm{log}|f_1|di\mathrm{arg}f_2di\mathrm{arg}f_3\right)$$
$$r_3(2):\{f\}_2g\widehat{}_2(f)di\mathrm{arg}g\frac{1}{3}\mathrm{log}|g|\left(\mathrm{log}|1f|d\mathrm{log}|f|+\mathrm{log}|f|d\mathrm{log}|1f|\right)$$
$$r_3(1):\{f\}_3\widehat{}_3(f)$$
6. The general case. Let $`d:=\mathrm{dim}X`$. Then the complex $`\mathrm{\Gamma }(X;n)`$ is the total complex of the following bicomplex:
$$\mathrm{\Gamma }((X);n)_{Y_1X^{(1)}}\mathrm{\Gamma }((Y_1);n1)[1]$$
$$_{Y_2X^{(2)}}\mathrm{\Gamma }((Y_2);n2)[1]\mathrm{}_{Y_dX^{(d)}}\mathrm{\Gamma }((Y_d);nd)[d]$$
where the arrows are provided by the residue maps, see \[G1\], p 239-240.
To define the regulator map from this gadget to the complex $`𝒞_𝒟(X_{};(n))`$ we specify it for each of $`\mathrm{\Gamma }((Y_k);nk)[k]`$ where $`k=0,\mathrm{},d`$. Namely, we take the constructed in theorem (2.5) homomorphism $`r_{nk}()`$ for $`\eta _{Y_k}`$ and multiply it by $`(2\pi i)^{nk}\delta _{Y_{nk}}`$. Notice that the distribution $`\delta _Y`$ depends only on the generic point of a subvariety $`Y`$. Its image in the bicomplex described in the section 3.1 is in
$$𝒟^{p,q};pk,qk,p+qn+k$$
which is the triangle symmetric with respect to the diagonal in the $`n\times n`$ square presented on the picture below.
Compatibility with the differentials is a corollary of theorem 2.5 and the following remark. If we would consider the $`n\times n`$ square in the Deligne complex as a part of the Dolbeault complex then we do not get a morphism of complexes. For $`k<n`$ the descrepency lies in the subgroup
$$𝒟^{k,n+1}𝒟^{n+1,k}$$
(31)
and come from the elements $`(Y_k,f_1\mathrm{}f_{nk})`$. However since the restriction of $`D`$ to the subgroup $`𝒟^{k,n}𝒟^{n,k}`$ equals to the de Rham differential $`𝒟`$ modulo the components (31) this does not create a problem.
## 4 Proofs
1. The differential equation for the function $`_n(z)`$. The following proposition was stated without a proof as formula 1.14 in \[G2\].
###### Proposition 4.1.
The differential equation for $`\widehat{}_n(f)`$ for $`n3`$ is:
$$d\widehat{}_n(f)=\widehat{}_{n1}(f)di\mathrm{arg}z\underset{k=2}{\overset{n1}{}}\beta _k\widehat{}_{nk,k}(f)$$
(32)
Proof. Consider the generating series
$$Li(z;t):=\underset{n1}{}Li_n(z)t^{n1};\widehat{}(z;t):=\underset{n1}{}\widehat{}_n(z)t^{n1}$$
###### Lemma 4.2.
Differential equations (32) together with the formulas
$$d\widehat{}_1(z)=d\mathrm{log}|1z|,d\widehat{}_2(z)=\mathrm{log}|1z|di\mathrm{arg}z+\mathrm{log}|z|di\mathrm{arg}(1z)$$
(33)
are equivalent to the formula
$$d\widehat{}(z;t)=d\mathrm{log}|1z|+\mathrm{log}|z|di\mathrm{arg}(1z)t+\widehat{}(z;t)di\mathrm{arg}zt$$
(34)
$$\left\{\frac{2|\mathrm{log}z|t}{e^{2\mathrm{log}|z|t}1}+\mathrm{log}|z|t1\right\}\left(\widehat{}(z;t)\frac{d\mathrm{log}|z|}{\mathrm{log}|z|}+d\mathrm{log}|1z|\right)$$
(35)
Proof. We check directly that the differential equation (33) appears as the coefficients of (35) in $`t^0`$ and $`t^1`$. From now on we will work modulo the monomials $`t^0`$ and $`t^1`$. The second term on the right of (32) (written for $`n3`$) provides the total contribution of
$$\underset{n3}{}\left(\underset{k2}{\overset{n2}{}}\beta _k\widehat{}_{nk}(z)\mathrm{log}^{k1}|z|d\mathrm{log}|z|+\beta _{n1}\alpha (1z,z)\mathrm{log}^{n2}|z|\right)t^{n1}=$$
$$\underset{2k<n}{}\beta _k\widehat{}_{nk}(z)\mathrm{log}^{k1}|z|t^{n1}d\mathrm{log}|z|$$
$$\underset{n3}{}\beta _{n1}\mathrm{log}^{n1}|z|t^{n1}d\mathrm{log}|1z|=$$
$$\left\{\frac{2\mathrm{log}zt}{e^{2\mathrm{log}|z|t}1}+\mathrm{log}|z|t1\right\}\left(\widehat{}(z;t)\frac{d\mathrm{log}|z|}{\mathrm{log}|z|}+d\mathrm{log}|1z|\right)$$
which is precisely (35). The first term on the right of (32) gives the last term in (34). The lemma is proved.
Let us prove formula (34) - (35). One has for $`n2`$:
$$d\pi _n(Li_m(z))=\frac{1}{2}d\left(Li_m(z)+(1)^{n1}\overline{Li_m(z)}\right)=$$
$$\pi _n(Li_{m1}(z))d\mathrm{log}|z|+\pi _{n1}(Li_{m1}(z))di\mathrm{arg}z$$
(36)
Define an operator $`P`$ acting on the generating series for a sequence of holomorphic functions $`f_n(z)`$ by
$$P\left(\underset{n1}{}f_n(z)t^{n1}\right):=\underset{n1}{}\pi _n(f_n)t^{n1}$$
Then
$$\widehat{}(z;t)=P\left\{\underset{0k<n}{}\beta _kLi_{nk}(z)\mathrm{log}^k|z|t^{n1}\right\}=$$
$$P\left\{Li(z;t)\frac{2\mathrm{log}|z|t}{e^{2\mathrm{log}|z|t}1}\right\}$$
(37)
One has
$$d\left(\frac{2x}{e^{2x}1}\right)=\left(\frac{2x}{e^{2x}1}\frac{4x^2e^{2x}}{(e^{2x}1)^2}\right)\frac{dx}{x}=\frac{2x}{e^{2x}1}\left(12x\frac{2x}{e^{2x}1}\right)\frac{dx}{x}$$
Thus since the second factor in (37) is a real function we get
$$d\widehat{}(z;t)=d_1P\left(Li(z;t)\frac{2\mathrm{log}|z|t}{e^{2\mathrm{log}|z|t}1}\right)$$
$$+P\left\{Li(z;t)\frac{2\mathrm{log}|z|t}{e^{2\mathrm{log}|z|t}1}\left(12\mathrm{log}|z|t\frac{2\mathrm{log}|z|t}{e^{2\mathrm{log}|z|t}1}\right)\right\}\frac{d\mathrm{log}|z|}{\mathrm{log}|z|}$$
where $`d_1`$ is the differential applied to $`Li(z;t)`$ only. Since $`\frac{2x}{e^{2x}1}+x1`$ is an even function in $`x`$ we rewrite this as
$$d_1P\left(Li(z;t)\frac{2\mathrm{log}|z|t}{e^{2\mathrm{log}|z|t}1}\right)$$
(38)
$$\widehat{}(z;t)\left(\frac{2\mathrm{log}|z|t}{e^{2\mathrm{log}|z|t}1}+\mathrm{log}|z|t1\right)\frac{d\mathrm{log}|z|}{\mathrm{log}|z|}$$
(39)
$$P\left\{Li(z;t)\frac{2\mathrm{log}|z|t}{e^{2\mathrm{log}|z|t}1}t\right\}d\mathrm{log}|z|$$
(40)
We handle this as follows.
i) (39) matches half of (34).
ii) The $`Li_1`$-part of (38) is
$$d_1P\left(Li_1(z)\frac{2\mathrm{log}|z|t}{e^{2\mathrm{log}|z|t}1}\right)=\frac{2\mathrm{log}|z|t}{e^{2\mathrm{log}|z|t}1}d\mathrm{log}|1z|$$
It matches the second half of (35) modulo coefficients in $`t^0`$ and $`t^1`$.
iii) The rest of (38) equals to
$$P\left\{Li(z;t)\frac{2\mathrm{log}|z|t}{e^{2\mathrm{log}|z|t}1}d\mathrm{log}zt\right\}=$$
(41)
$$P\left\{\widehat{}(z;t)(z;t)\right\}di\mathrm{arg}zt+P\left\{Li(z;t)\frac{2\mathrm{log}|z|t}{e^{2\mathrm{log}|z|t}1}t\right\}d\mathrm{log}|z|$$
(42)
Thus the first term in (42) matches the right term in (34), and the second is canceled with (40). The proposition is proved.
2. Proof of theorem 2.5. i) Let us show that the map
$$\{f\}\widehat{}_{nk,k}(f)$$
(43)
provides a group homomorphism $`_n((X))𝒜_\eta ^1(X())`$. The arguments are similar to the proof of theorem 1.15 in \[G2\]. Suppose first that $`nk>1`$. By theorem 1.15 in \[G2\] the map $`\{f\}\widehat{}_n(f)`$ provides a group homomorphism $`_n((X))𝒜_\eta ^0(X())`$. Consider the following map
$$_n((X))_{nk}((X))S^k(X)^{},\{f\}_n\{f\}_{nk}f^k$$
(44)
It can be defined as a composition
$$\{f\}_n\{f\}_{n1}f\{f\}_{n2}ff\mathrm{}\{f\}_{nk}f\mathrm{}f$$
Each of these maps is of type $`\delta id`$ and thus is a homomorphism of abelian groups. The composition followed by the homomorphism
$$_{nk}((X))S^k(X)^{}𝒜_\eta ^1(X())$$
$$\{f\}_{nk}g_1\mathrm{}g_k\widehat{}_{nk}(f)\frac{1}{k}d(\mathrm{log}|g_1|\mathrm{}\mathrm{log}|g_k|)$$
leads to the map (43). In the case $`nk=1`$ we present map (43) as a composition
$$_n((X))\stackrel{(\text{44})}{}_2((X))S^{n2}(X)^{}\stackrel{a}{}𝒜_\eta ^1(X())$$
$$a:\{f\}_2g_1\mathrm{}g_{n2}\alpha (1f,f)\mathrm{log}|g_1|\mathrm{}\mathrm{log}|g_{n2}|$$
and use the fact that $`\{f\}_2\alpha (1f,f)`$ is a group homomorphism.
ii) The properties (a) and (b) in theorem 2.2 are true by the very definitions.
iii) It is very easy to see that the property (d) is also true by the very definitions: it basically reflects the fact that the numerical coefficients in the formula for $`r_{n+m}(m+1)`$ do not depend on $`m`$ and uses formula (11).
iv) It remains to prove that the diagram in theorem 2.2 is commutative. Let us show first that its right square
$$\begin{array}{ccc}_2((X))\mathrm{\Lambda }^{n2}(X)^{}& \stackrel{\delta }{}& \mathrm{\Lambda }^n(X)^{}\\ & & \\ r_n(n1)& & r_n(n)\\ & & \\ 𝒜_\eta ^{n2}(X())& \stackrel{d}{}& 𝒜_\eta ^{n1}(X())\end{array}$$
is commutative. We change the numeration by putting $`n:=m+2`$.
###### Proposition 4.3.
$$r_{m+2}(m+1):\{f\}_2g_1\mathrm{}g_m$$
$$\widehat{}_2(f)𝒜_m\left\{\underset{p0}{}\frac{1}{2p+1}\underset{i=1}{\overset{2p}{}}d\mathrm{log}|g_i|\underset{j=2p+1}{\overset{m}{}}di\mathrm{arg}g_j\right\}$$
$$\alpha (1f,f)𝒜_m\left(\underset{p0}{}\frac{1}{(2p+1)(2p+3)}\mathrm{log}|g_1|\underset{i=2}{\overset{2p+1}{}}d\mathrm{log}|g_i|\underset{j=2p+2}{\overset{m}{}}di\mathrm{arg}g_j\right)$$
Proof. It follows immediately from the definition and lemma 2.3:
Let us compute
$$d\left(r_{m+2}(m+1)(\{f\}_2g_1\mathrm{}g_m)\right)$$
Using proposition 4.3 we get
$$\left(\mathrm{log}|1f|di\mathrm{arg}f+\mathrm{log}|f|di\mathrm{arg}(1f)\right)$$
(45)
$$𝒜_m\left\{\underset{p0}{}\frac{1}{2p+1}\underset{i=1}{\overset{2p}{}}d\mathrm{log}|g_i|\underset{j=2p+1}{\overset{m}{}}di\mathrm{arg}g_j\right\}$$
$$+2d\mathrm{log}|1f|d\mathrm{log}|f|$$
$$𝒜_m\left(\underset{p0}{}\frac{1}{(2p+1)(2p+3)}\mathrm{log}|g_1|\underset{i=2}{\overset{2p+1}{}}d\mathrm{log}|g_i|\underset{j=2p+2}{\overset{m}{}}di\mathrm{arg}g_j\right)$$
(46)
$$+\alpha (1f,f)𝒜_m\left(\underset{p0}{}\frac{1}{2p+3}\underset{i=1}{\overset{2p+1}{}}d\mathrm{log}|g_i|\underset{j=2p+2}{\overset{m}{}}di\mathrm{arg}g_j\right)$$
(47)
We can write the differential form $`r_m(g_1\mathrm{}g_m)`$ as follows:
$$𝒜_m\left\{\underset{p0}{}\frac{1}{2p+1}\mathrm{log}|g_1|\underset{i=2}{\overset{2p+1}{}}d\mathrm{log}|g_i|\underset{j=2p+2}{\overset{m}{}}di\mathrm{arg}g_j\right\}$$
(48)
In particular this makes transparent the property
$$dr_m(g_1\mathrm{}g_m)+\pi _m(d\mathrm{log}g_1\mathrm{}\mathrm{log}g_m)=0$$
Now computing
$$dr_{m+2}((1f)fg_1\mathrm{}g_m)$$
(49)
using formula (48) and comparing the result with formulas (45) - (47) we see the following:
1) Formula (45) matches the part of (49) where the contribution of $`(1f)f`$ is given by $`\mathrm{log}|1f|di\mathrm{arg}|f|+\mathrm{log}|f|di\mathrm{arg}|1f|`$.
2) Formula (47) matches the part of (49) where the contribution of $`(1f)f`$ is given by $`\mathrm{log}|1f|d\mathrm{log}|f|+\mathrm{log}|f|d\mathrm{log}|1f|`$.
Before we continue any further let us note that
$$d\mathrm{log}|1f|di\mathrm{arg}(f)=d\mathrm{log}|f|di\mathrm{arg}(1f)$$
(50)
$$d\mathrm{log}|1f|d\mathrm{log}|f|=di\mathrm{arg}(1f)di\mathrm{arg}f$$
(51)
Indeed,
$$0=d\mathrm{log}(1f)d\mathrm{log}f=\left(d\mathrm{log}|1f|+di\mathrm{arg}(1f)\right)\left(d\mathrm{log}|f|+d\mathrm{arg}if\right)$$
Therefore using (50) we see the following:
3) the sum of the terms of (49) where the contribution of $`(1f)f`$ is given either by $`d\mathrm{log}|1f|di\mathrm{arg}|f|`$ or by $`di\mathrm{arg}(1f)fd\mathrm{log}|f|`$ is zero.
4) It remains to show that (46) matches the part of (49) where the contribution of $`(1f)f`$ is given by $`d\mathrm{log}|1f|d\mathrm{log}|f|`$ or $`di\mathrm{arg}(1f)fdi\mathrm{arg}f`$.
Let us substitute
$$\frac{2}{(2p+1)(2p+3)}=\frac{1}{2p+1}\frac{1}{2p+3}$$
into (46) and split the formula we get into two parts. The first part, denoted $`(\text{46})_I`$, is the one where the terms appear with the coefficient $`\frac{1}{2p+1}`$. The second part, denoted $`(\text{46})_{II}`$, is the rest. It corresponds to $`\frac{1}{2p+3}`$.
Using (51) to calculate the part of (49) where $`(1f)f`$ contributes via $`di\mathrm{arg}(1f)fdi\mathrm{arg}f`$ we see that it matches with $`(\text{46})_I`$. The other part, $`(\text{46})_{II}`$, matches the part of (49) where $`(1f)f`$ contributes $`d\mathrm{log}|1f|d\mathrm{log}|f|`$. So we achieved our goal 4). The commutativity of the last square is proved.
Now let us prove the commutativity of the squares different from the last one, i.e. ($`n>2`$)
$$\begin{array}{ccc}_n((X))\mathrm{\Lambda }^m(X)^{}& \stackrel{\delta }{}& _{n1}((X))\mathrm{\Lambda }^{m+1}(X)^{}\\ & & \\ r_{n+m}(m+1)& & r_{n+m}(m+2)\\ & & \\ 𝒜_\eta ^m(X())& \stackrel{d}{}& 𝒜_\eta ^{m+1}(X())\end{array}$$
One has
$$dr_{n+m}(m+1)(\{f\}_ng_1\mathrm{}g_m)=$$
(52)
$$d\widehat{}_n(f)𝒜_m\left\{\underset{p0}{}\frac{1}{2p+1}\underset{i=1}{\overset{2p}{}}d\mathrm{log}|g_i|\underset{j=2p+1}{\overset{m}{}}di\mathrm{arg}g_j\right\}$$
(53)
$$+\underset{k1}{}\underset{1pm}{}\beta _{k,p}d\widehat{}_{nk,k}(f)𝒜_m\left\{\mathrm{log}|g_1|\underset{i=2}{\overset{p}{}}d\mathrm{log}|g_i|\underset{j=p+1}{\overset{m}{}}di\mathrm{arg}g_j\right\}$$
(54)
$$\underset{k1}{}\underset{1pm}{}p\beta _{k,p}\widehat{}_{nk,k}(f)𝒜_m\left\{\underset{i=1}{\overset{p}{}}d\mathrm{log}|g_i|\underset{j=p+1}{\overset{m}{}}di\mathrm{arg}g_j\right\}$$
(55)
Let us compare this expression with
$$r_{n+m}(m+2)(\{f\}_{n1}fg_1\mathrm{}g_m)=(\text{20})_{m+2}+(\text{21})_{m+2}$$
(56)
where $`(\text{20})_{m+2}`$ (resp. $`(\text{21})_{m+2}`$) stays for the term similar to $`(\text{20})`$ (resp. $`(\text{21})`$) in the definition of $`r_{n+m}(m+1)`$ (see theorem-construction 2.5).
Observe that a priori $`f`$ can contribute to $`(\text{20})_{m+2}`$ via $`d\mathrm{log}|f|`$ or $`di\mathrm{arg}f`$, and to $`(\text{21})_{m+2}`$via $`\mathrm{log}|f|`$, $`d\mathrm{log}|f|`$ or $`di\mathrm{arg}f`$. A careful reader should keep track of these five different cases during the proof.
We split the job in two parts:
A) Check that the diagram is commutative if we pay attention only to the terms where $`d\mathrm{log}|1f|`$ or $`di\mathrm{arg}(1f)`$ appears. We call such terms $`A`$-terms.
B) Then we check that the diagram is commutative modulo the terms with $`d\mathrm{log}|1f|`$ or $`di\mathrm{arg}(1f)`$.
Part A). Here is the contribution from (53) - (55). Recall that $`n3`$.
From (53) using (32) we get:
$$\beta _{n1}\mathrm{log}^{n1}|f|d\mathrm{log}|1f|𝒜_m\left\{\underset{p0}{}\frac{1}{2p+1}\underset{i=1}{\overset{2p}{}}d\mathrm{log}|g_i|\underset{j=2p+1}{\overset{m}{}}di\mathrm{arg}g_j\right\}$$
(57)
From (54):
$$\underset{k=1}{\overset{n3}{}}\underset{1pm}{}\beta _{k,p}\beta _{nk1}\mathrm{log}^{n2}|f|d\mathrm{log}|1f|d\mathrm{log}|f|$$
(58)
$$𝒜_m\left\{\mathrm{log}|g_1|\underset{i=2}{\overset{p}{}}d\mathrm{log}|g_i|\underset{j=p+1}{\overset{m}{}}di\mathrm{arg}g_j\right\}$$
(59)
$$+\underset{1pm}{}\beta _{n2,p}di\mathrm{arg}(1f)d\mathrm{log}|f|\mathrm{log}^{n2}|f|(\text{59})$$
(60)
$$\underset{1pm}{}(n1)\beta _{n1,p}d\mathrm{log}|1f|d\mathrm{log}|f|\mathrm{log}^{n2}|f|(\text{59})$$
(61)
From (55):
$$\underset{1pm}{}p\beta _{n1,p}\mathrm{log}^{n1}|f|d\mathrm{log}|1f|$$
(62)
$$𝒜_m\left\{\underset{i=1}{\overset{p}{}}d\mathrm{log}|g_i|\underset{j=p+1}{\overset{m}{}}di\mathrm{arg}g_j\right\}$$
(63)
The contribution of $`(\text{20})_{m+2}`$ is zero. The factor $`f`$ from $`\{f\}_{n1}fg_1\mathrm{}g_m`$ can contribute into $`(\text{21})_{m+2}`$ via $`\mathrm{log}|f|`$, $`di\mathrm{arg}f`$ and $`d\mathrm{log}|f|`$.
The terms where $`f`$ contributes via $`\mathrm{log}|f|`$ are
$$\underset{k1}{}\underset{1pm}{}\beta _{k,p+1}\widehat{}_{n1k,k}(f)\mathrm{log}|f|(\text{63})$$
(64)
Since we are concerned only with the $`A`$-terms, this boils down to
$$\underset{1pm}{}\beta _{n2,p+1}\widehat{}_{1,n2}(f)\mathrm{log}|f|(\text{63})=$$
(65)
$$\underset{1pm}{}\beta _{n2,p+1}\mathrm{log}^{n1}|f|d\mathrm{log}|1f|(\text{63})$$
Using formulas (12) we see that it matches (57) + (62).
The terms where $`f`$ contributes into $`(\text{21})_{m+2}`$ via $`di\mathrm{arg}f`$ are
$$\underset{k1}{}\underset{1pm}{}\beta _{k,p}\widehat{}_{n1k,k}(f)di\mathrm{arg}f(\text{59})$$
(66)
Its contribution modulo $`B`$-terms is
$$\underset{1pm}{}\beta _{n2,p}\mathrm{log}^{n2}|f|d\mathrm{log}|1f|di\mathrm{arg}f(\text{59})$$
(67)
Using (50) we see that it matches (60).
Similar considerations for terms where $`f`$ contributes into $`(\text{21})_{m+2}`$ via $`d\mathrm{log}|f|`$ minus ((61) + (58)) leads to the left hand side of (16), which is zero according to proposition 2.4. The part A) is proved.
Part B). We will write $`X\stackrel{A}{=}Y`$ if $`XY=0`$ modulo the $`A`$-terms.
a) Notice that the contribution of $`f`$ to $`(\text{21})_{m+2}`$ via $`d\mathrm{log}|f|`$ is zero since $`\widehat{}_{nk,k}(f)d\mathrm{log}|f|=0`$ modulo $`A`$-terms.
b) The contribution to (53) of the term $`\widehat{}_{n1}(f)di\mathrm{arg}f`$ in formula (32) for $`d\widehat{}_n(f)`$ matches the part of term $`(\text{20})_{m+2}`$ where $`f`$ contributes via $`di\mathrm{arg}f`$.
c) One has
$$d\widehat{}_{nk,k}(f)\stackrel{A}{=}\widehat{}_{nk1,k}(f)di\mathrm{arg}f$$
(68)
Thus (54) matches the part of $`(\text{21})_{m+2}`$ where $`f`$ contributes via $`di\mathrm{arg}f`$.
d) The contribution to (53) of the other term in $`d\widehat{}_n(f)`$ modulo $`A`$ is
$$\underset{p0}{}\underset{k=2}{\overset{n2}{}}\frac{\beta _k}{2p+1}\widehat{}_{nk,k1}(f)$$
(69)
$$𝒜_m\left\{\mathrm{log}|f|\underset{i=1}{\overset{2p}{}}d\mathrm{log}|g_i|\underset{j=2p+1}{\overset{m}{}}di\mathrm{arg}g_j\right\}$$
(70)
The only other terms in (52) with factor (70) are the terms of (55) where $`k>1`$ and the number of $`d\mathrm{log}|g_i|`$’s is even, i.e.
$$\underset{k1}{}2p\beta _{k,2p}\widehat{}_{nk,k1}(f)(\text{70})$$
(71)
On the other hand, in (56) the terms with factor (70) are
$$\underset{k1}{}\beta _{k,2p+1}\widehat{}_{n1k,k}(f)(\text{70})$$
(72)
which are precisely the terms of $`(\text{21})_{m+2}`$ where $`f`$ contributes via $`\mathrm{log}|f|`$, and $`p`$ in $`\beta _{k,p}`$ is odd.
Thus out of (69), (71) and (72) we see that $`\widehat{}_{nk,k1}(f)(\text{70})`$ appears in our check of the commutativity of the diagram with factor
$$\beta _{k1,2p+1}+\frac{\beta _k}{2p+1}+2p\beta _{k,2p}\stackrel{(\text{12})}{=}0$$
e) The term of $`(\text{20})_{m+2}`$ where $`f`$ contributes via $`d\mathrm{log}|f|`$, i.e.
$$\underset{p0}{}\frac{1}{2p+1}\widehat{}_{n1}(f)𝒜_{m+1}\left\{d\mathrm{log}|f|\underset{i=1}{\overset{2p1}{}}d\mathrm{log}|g_i|\underset{j=2p}{\overset{m}{}}di\mathrm{arg}g_j\right\}$$
(73)
matches the term (55) with $`k=1`$, $`p`$: odd, i.e.
$$\underset{p0}{}(2p1)\beta _{1,2p1}\widehat{}_{n1,1}(f)𝒜_m\left\{\underset{i=1}{\overset{2p1}{}}d\mathrm{log}|g_i|\underset{j=2p}{\overset{m}{}}di\mathrm{arg}g_j\right\}$$
(74)
Indeed, according to (14) $`(2p1)\beta _{1,2p1}=\frac{1}{2p+1}`$.
f) All the terms of (55) with $`k=1`$ and even $`p`$ are zero since $`\beta _{1,2p}=0`$ by (14).
g) The terms of (55) with odd $`p`$ and $`k>1`$ look as follows:
$$\underset{p0}{}\underset{k2}{}(2p+1)\beta _{k,2p+1}\widehat{}_{nk,k1}(f)$$
(75)
$$𝒜_{m+1}\left\{\mathrm{log}|f|\underset{i=1}{\overset{2p+1}{}}d\mathrm{log}|g_i|\underset{j=2p+2}{\overset{m}{}}di\mathrm{arg}g_j\right\}$$
(76)
They match with the terms of $`(\text{21})_{m+2}`$ where $`f`$ contributes via $`\mathrm{log}|f|`$ and $`p`$ is even, see (64).
The the part B) of the commutativity proof is now complete. To double check that we list all the possibilities for contribution of $`f`$ in (56), and indicate for each of them the part of the proof where it was handled.
Contribution of $`f`$ to $`(\text{20})_{m+1}`$:
via $`d\mathrm{log}|f|`$: see e); via $`di\mathrm{arg}f`$: see b).
Contribution of $`f`$ to $`(\text{21})_{m+2}`$:
via $`\mathrm{log}|f|`$, $`p`$ odd: see d); via $`\mathrm{log}|f|`$, $`p`$ even: see g);
via $`d\mathrm{log}|f|`$: see a); via $`di\mathrm{arg}f`$: see c).
Contribution of $`f`$ to (53): see b) and d).
Contribution of $`f`$ to (54): see c).
Contribution of $`f`$ to (55):
$`k=1`$, $`p`$ odd: see e); $`k>1`$, $`p`$ odd: see g);
$`k=1`$, $`p`$ even: see f); $`k>1`$, $`p`$ even: see d).
The main theorem is proved.
Department of mathematics, Brown University, Providence, RI 02912. e-mail sasha@math.brown.edu
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# New results of 116Cd double 𝛽 decay study with 116CdWO4 scintillators
## I Introduction
Neutrinoless (0$`\nu `$) double $`\beta `$ decay is forbidden in the Standard Model (SM) since it violates lepton number ($`L`$) conservation. However many extensions of the SM incorporate $`L`$ violating interactions and thus could lead to the 0$`\nu `$2$`\beta `$ decay . Currently, besides conventional neutrino ($`\nu `$) exchange mechanism, there are many other possibilities to trigger this process . In that sense neutrinoless 2$`\beta `$ decay has a great conceptual importance due to the strong Schechter-Valle theorem obtained in a gauge theory of the weak interaction that a non-vanishing 0$`\nu `$2$`\beta `$ decay rate requires neutrino to be massive Majorana particle, independent of which mechanism induces it. Therefore, at present 0$`\nu `$2$`\beta `$ decay is considered as a very powerful test of new physical effects beyond the SM, and even the absence of this process – at the present level of sensitivity – would help to restrict or narrow this wide choice of theoretical models. At the same time 0$`\nu `$2$`\beta `$ decay is very important in the light of the measured deficit of the atmospheric neutrinos flux and the result of the LSND accelerator experiment , which could be explained by means of the neutrino oscillations requiring in turn nonzero neutrino masses ($`m_\nu `$). However oscillation experiments are sensitive to neutrino mass difference, while measured 0$`\nu `$2$`\beta `$ decay rate can give the absolute value of the $`m_\nu `$ scale, and hence provide a crucial test of $`m_\nu `$ models.
Despite of numerous attempts to observe 0$`\nu `$2$`\beta `$ decay from 1948 up to present days this process still remains unobserved. The highest $`T_{1/2}`$(0$`\nu `$) limits were set in direct experiments with several nuclides: $`T_{1/2}10^{22}`$ yr for <sup>82</sup>Se , <sup>100</sup>Mo , <sup>116</sup>Cd ; $`T_{1/2}10^{23}`$ yr for <sup>130</sup>Te and <sup>136</sup>Xe ; and $`T_{1/2}10^{25}`$ yr for <sup>76</sup>Ge .
With the aim to enlarge the number of 2$`\beta `$ decay candidate nuclides studied at a sensitivity comparable with that for <sup>76</sup>Ge and <sup>136</sup>Xe (neutrino mass limit of 0.5–2 eV), cadmium tungstate crystal scintillators, enriched in <sup>116</sup>Cd to 83%, were developed and exploited in <sup>116</sup>Cd research . The measurements were carried out in the Solotvina Underground Laboratory in a salt mine 430 m underground ($``$1000 m w. e.) . In the first phase of the experiment only one <sup>116</sup>CdWO<sub>4</sub> crystal (15.2 cm<sup>3</sup>) was used. It was viewed by a photomultiplier tube (PMT) through a light-guide 51 cm long and placed inside a plastic scintillator ($``$38$`\times `$115 cm) which served as veto counter. A passive shield of high purity (HP) copper (5 cm), lead (23 cm) and polyethylene (16 cm) surrounded the plastic counter. The background rate in the energy range 2.7–2.9 MeV ($`Q_{2\beta }`$=2805 keV ) was equal $``$0.6 counts/yr$``$kg$``$keV. With 19175 h statistics the half-life limit for 0$`\nu `$2$`\beta `$ decay of <sup>116</sup>Cd was set as $`T_{1/2}`$(0$`\nu `$) $``$ 3.2$`\times `$1$`0^{22}`$ yr (90% C.L.) , which corresponds to the restriction on the neutrino mass $`m_\nu 3.9`$ eV . Limits on 0$`\nu `$2$`\beta `$ decay with emission of one (M1) or two (M2) Majorons were obtained too: $`T_{1/2}`$(0$`\nu `$M1) $``$ 1.2$`\times `$1$`0^{21}`$ yr and $`T_{1/2}`$(0$`\nu `$M2) $``$ 2.6$`\times `$1$`0^{20}`$ yr (90% C.L.) .
In the present paper new and advanced results of <sup>116</sup>Cd research obtained with the help of an upgraded apparatus are described.
## II New set-up with four <sup>116</sup>CdWO<sub>4</sub> detectors
### A Set-up and measurements
In order to enhance the sensitivity of the <sup>116</sup>Cd experiment, the following improvements were scheduled: increase of the number of <sup>116</sup>Cd nuclei, reduction of the background and improvement of the data taking and processing . With this aim the new set-up with four enriched <sup>116</sup>CdWO<sub>4</sub> crystals (total mass 339 g) has been mounted in the Solotvina Laboratory in August 1998. All materials used in the installation were previously tested and selected for low radioactive impurities in order to reduce their contributions to background.
In the new apparatus, a scheme of which is shown in fig. 1, four enriched crystals are viewed by the PMT (EMI9390) through one light-guide 10 cm in diameter and 55 cm long, which is composed of two glued parts: quartz 25 cm long and plastic scintillator (Bicron BC-412) 30 cm long. The <sup>116</sup>CdWO<sub>4</sub> crystals are surrounded by an active shield made of 15 natural CdWO<sub>4</sub> scintillators of large volume ($``$200 cm<sup>3</sup> each) with total mass of 20.6 kg. Due to the high purity of the CdWO<sub>4</sub> crystals and their large density ($``$8 g/cm<sup>3</sup>) this active shield reduces background effectively. The veto crystals are viewed – by a low background PMT ($``$17 cm) – through an active plastic light-guide ($``$17$`\times `$49 cm). In turn the whole array of CdWO<sub>4</sub> counters is placed inside an additional active shield made of polystyrene-based plastic scintillator with dimensions 40$`\times `$40$`\times `$95 cm. Thus, together with both active light-guides (connected with enriched and natural crystals on opposite sides), a complete 4$`\pi `$ active shield of the main <sup>116</sup>CdWO<sub>4</sub> detectors is provided.
The outer passive shield consists of HP copper (thickness $`3`$$`6`$ cm), lead ($`22.5`$$`30`$ cm) and polyethylene (16 cm). Two plastic scintillators (120$`\times `$130$`\times `$3 cm) are installed above the passive shield to provide a cosmic muons veto. Because air in the Solotvina Laboratory can be contaminated by radon (at the level $``$ 30 Bq/m<sup>3</sup>) the set-up is isolated carefully against air penetration. All cavities inside the shield are filled up by pieces of plexiglass, and HP Cu shield is sealed with the help of silicon glue and enclosed inside a tight mylar envelope.
The new event-by-event data acquisition is based on two IBM personal computers (PC) and a CAMAC crate with electronic units. For each event the following information is stored on the hard disc of the first computer: the amplitude (energy) of a signal, its arrival time and the following additional tags: the coincidence between different detectors; the signal of radio-noise detection system; triggers for light emitting diode (LED) and pulse shape digitizer. The second computer records the pulse shape of the <sup>116</sup>CdWO<sub>4</sub> scintillators in the energy range 0.25–5 MeV. This complementary system is developed on the basis of a fast 12 bit ADC (Analog Devices AD9022) and is connected with computer by parallel digital I/O board (PC-DIO-24 from National Instruments) . Two PC-DIO-24 boards are used to link both computers and establish – with the help of proper software – a one-to-one correspondence between the pulse shape data recorded by the second computer and the information stored in the first PC.
The energy scale and resolution of the main detector – four enriched crystals taken as a whole – were determined in the measurements with different sources (<sup>22</sup>Na, <sup>40</sup>K, <sup>60</sup>Co, <sup>137</sup>Cs, <sup>207</sup>Bi, <sup>226</sup>Ra, <sup>232</sup>Th and <sup>241</sup>Am). The energy dependence of the resolution can be expressed (for the energy above 50 keV) as $`FWHM(`$keV) =$`\sqrt{226+\text{1}6.6E+6.42\times 10^3E^2}`$, where energy $`E`$ is in keV. For instance, the resolution ($`FWHM`$) was equal to 14.5% at 1064 keV and 11% at 2615 keV. The full energy peaks are well fitted in the energy region 0.06–2.6 MeV by a Gaussian function with typical value $`\chi ^2`$ = 0.8–1.9. Moreover, the calibration spectra of the <sup>232</sup>Th source were simulated with the help of GEANT3.21 package and event generator DECAY4 (the last defines initial kinematics of the events). The simulated <sup>232</sup>Th spectra are in good agreement with the measured ones confirming the assumption of a Gaussian peak shape. In particular for the 2615 keV peak of <sup>208</sup>Tl – which is close to the 2$`\beta `$ decay energy of <sup>116</sup>Cd – the value of $`\chi ^2`$ =1.3.
Also, the relative light yield for for $`\alpha `$ particles as compared with that for electrons ($`\alpha /\beta `$ ratio) and energy resolution were measured with $`\alpha `$ sources and corrected by using time-amplitude analysis (see below) as following: $`\alpha /\beta `$ = 0.15(1) $`+`$7$`\times `$1$`0^6E_\alpha `$ and $`FWHM_\alpha (`$keV) = $`0.053E_\alpha `$ ($`E_\alpha `$ is in keV). The routine calibration is carried out weekly with a <sup>207</sup>Bi source (570, 1064 and 1770 keV) and once per two weeks with <sup>232</sup>Th (2615 keV). The dead time of the spectrometer and data acquisition is monitored permanently with the help of an LED optically connected with the main PMT. The actual dead time value is $``$4.2$`\%`$ ($``$3% is owing to random coincidence between the main and shield detectors; $``$1.2% is caused by miscounts of the data acquisition).
The background spectrum measured during 4629 h in the new installation with four <sup>116</sup>CdWO<sub>4</sub> crystals is given in fig. 2, where the old data obtained with one <sup>116</sup>CdWO<sub>4</sub> crystal of 121 g are also shown for comparison. As it is visible from this figure, the background is decreased in the whole energy range, except for the $`\beta `$ spectrum of <sup>113</sup>Cd ($`Q_\beta `$ = $`316`$ keV), whose abundance in <sup>116</sup>CdWO<sub>4</sub> crystals is $``$2% . In the energy region $`2.5`$$`3.2`$ MeV – where the peak of 0$`\nu `$2$`\beta `$ decay of <sup>116</sup>Cd is expected – the background rate is reduced to a value of 0.03 counts/yr$``$kg$``$keV (only 4 events in the energy window $`2.5`$$`3.2`$ MeV were detected during 4629 h), twenty times lower than in the previous set-up. It is achieved, first, due to improvement of passive and active shield, and secondly, as a result of data processing advance (time-amplitude and pulse-shape analysis), which are described below.
### B Time-amplitude analysis of the data
The energy and arrival time of each event can be used for analysis and selection of some decay chains in <sup>232</sup>Th, <sup>235</sup>U and <sup>238</sup>U families (see f. e. ref. ). As an example (important in the following for the background rejection in the energy range of 0$`\nu 2\beta `$ decay), we consider here in detail the time-amplitude analysis of the following sequence of $`\alpha `$ decays from <sup>232</sup>Th family: <sup>220</sup>Rn ($`Q_\alpha `$ = $`6.40`$ MeV, $`T_{1/2}`$ = $`55.6`$ s) $``$ <sup>216</sup>Po ($`Q_\alpha `$ = $`6.91`$ MeV, $`T_{1/2}`$ = $`0.145`$ s) $``$ <sup>212</sup>Pb. Because the energy of $`\alpha `$ particles from <sup>220</sup>Rn decay corresponds to $``$1.2 MeV in $`\beta /\gamma `$ scale of <sup>116</sup>CdWO<sub>4</sub> detector, the events in the energy region $`0.7`$$`1.8`$ MeV were used as triggers. Then all events (within $`0.9`$$`1.9`$ MeV) following the triggers in the time interval $`10`$$`1000`$ ms (containing 94.5% of <sup>216</sup>Po decays) were selected. The spectra of the <sup>220</sup>Rn and <sup>216</sup>Po $`\alpha `$ decays obtained in this way from data – as well as the distribution of the time intervals between the first and second events – are presented in fig. 3. It is evident from this figure that the selected spectra and time distribution are in an excellent agreement with those expected from $`\alpha `$ particles of <sup>220</sup>Rn and <sup>216</sup>Po. Using these results and taking into account the efficiency of the time-amplitude analysis and the number of accidental coincidences (3 pairs from 218 selected), the determined activity of <sup>228</sup>Th (<sup>232</sup>Th family) inside the <sup>116</sup>CdWO<sub>4</sub> crystals is as low as 38(3) $`\mu `$Bq/kg.
The same technique was applied to the sequence of $`\alpha `$ decays from the <sup>235</sup>U family: <sup>223</sup>Ra ($`Q_\alpha `$ = $`5.98`$ MeV, $`T_{1/2}`$ = $`11.44`$ d) $``$ <sup>219</sup>Rn ($`Q_\alpha `$ = $`6.95`$ MeV, $`T_{1/2}`$ = $`3.96`$ s) $``$ <sup>215</sup>Po ($`Q_\alpha `$ = $`7.53`$ MeV, $`T_{1/2}`$ = $`1.78`$ ms) $``$ <sup>211</sup>Pb. For the fast couple (<sup>219</sup>Rn $``$ <sup>215</sup>Po) all events within $`0.8`$$`1.8`$ MeV were used as triggers, while a time interval $`1`$$`10`$ ms (65.7% of <sup>215</sup>Po decays) and an energy window 0.9–2.0 MeV were set for the second events. The obtained $`\alpha `$ peaks correspond to an activity of 5.5(14) $`\mu `$Bq/kg for the <sup>227</sup>Ac impurity in the crystals.
As regard the <sup>226</sup>Ra chain (<sup>238</sup>U family) the following sequence of $`\beta `$ and $`\alpha `$ decays was analyzed: <sup>214</sup>Bi ($`Q_\beta `$ = $`3.27`$ MeV, $`T_{1/2}`$ = $`19.9`$ m) $``$ <sup>214</sup>Po ($`Q_\alpha `$ = $`7.83`$ MeV, $`T_{1/2}=164.3`$ $`\mu `$s) $``$ <sup>210</sup>Pb. For the first and second events the energy threshold was equal 0.1 MeV, and a time interval of $`100`$$`1000`$ $`\mu `$s (64.1% of <sup>214</sup>Po decays) was used. While the obtained spectrum of the first pulses agrees with the model of the $`\beta `$ decay of <sup>214</sup>Bi, and the distribution of the time intervals between the first and second events can be fitted by an exponent with $`T_{1/2}`$ = $`140_{20}^{+30}`$ $`\mu `$s (in reasonable agreement with the <sup>214</sup>Po half-life value), the spectrum of the second events is continuous, contrary to the anticipated $`\alpha `$ peak of <sup>214</sup>Po. Probably, part of this continuous distribution can be explained by <sup>226</sup>Ra contamination of the materials neighboring the <sup>116</sup>CdWO<sub>4</sub> crystals (optical grease, teflon, Mylar, radon in air), while another part is caused by <sup>226</sup>Ra decays in the crystals. Under such an assumption activity limits for <sup>226</sup>Ra contaminations are derived as $``$0.13(3) Bq/kg for optical grease, $``$8 mBq/kg for teflon, $``$1.8 $`\mu `$Bq/dm<sup>2</sup> for Mylar, and $``$ 5 $`\mu `$Bq/kg for <sup>116</sup>CdWO$`_4,`$ whose values do not contradict bounds obtained earlier . To prove these assumptions, the events belonging to the <sup>214</sup>Bi $``$ <sup>214</sup>Po $``$ <sup>210</sup>Pb chain were independently searched for in the time window of $`5`$$`88`$ $`\mu `$s (28.9% of <sup>214</sup>Po decays) with the help of pulse shape analysis (see below). For both events the energy threshold was $``$0.3 MeV. The result obtained (<sup>226</sup>Ra activity in the <sup>116</sup>CdWO<sub>4</sub> crystals $``$ 14 $`\mu `$Bq/kg) is similar to that of the time-amplitude analysis.
Finally, all couples of events found for <sup>232</sup>Th, <sup>235</sup>U and <sup>238</sup>U families as described above were eliminated from the measured data.
### C Pulse-shape discrimination
The pulse shape of the <sup>116</sup>CdWO<sub>4</sub> scintillators in the energy region of $`0.25`$$`5`$ MeV is digitized by a 12 bit ADC and stored in 2048 channels with 50 ns channel’s width. Due to different shapes of scintillation signal for various kinds of sourcesIt is known, that scintillation efficiency and pulse shape of inorganic crystals depend on the local density of the energy released, hence allowing to identify the incoming radiation (see e. g. ref. ). ($`\alpha `$ particles, protons, $`\gamma `$ quanta and cosmic muons were investigated), the pulse-shape (PS) discrimination method based on the optimal digital filter was developed and clear discrimination between $`\gamma `$ rays (electrons) and $`\alpha `$ particles was achieved .
The pulse shapes of enriched crystals were measured for $`\alpha `$ particles with an <sup>241</sup>Am source and for $`\gamma `$ rays with <sup>60</sup>Co, <sup>137</sup>Cs, <sup>207</sup>Bi and <sup>232</sup>Th sources in the special calibration runsBecause $`\gamma `$ rays interact with matter by mean of the energy transfer to electrons, it was assumed that pulse shapes for electrons and $`\gamma `$-s are the same. This statement was proved in the measurement with conversion electrons of <sup>207</sup>Bi by using the signal of the thin plastic scintillator (placed between source and detector) as signature of the electron hitting the <sup>116</sup>CdWO<sub>4</sub> crystal.. To provide an analytic description of the $`\alpha `$ or $`\gamma `$ signals $`f_\alpha (t)`$ and $`f_\gamma (t)`$ the pulse shape resulting from the average of a large number of individual events has been fitted with the sum of three (for $`\alpha `$ particles) or two (for $`\gamma `$-s ) exponents, giving the reference pulse shapes $`\overline{f}_\alpha (t)`$ and $`\overline{f}_\gamma (t)`$ (see for more details ref. ). In the data processing the digital filter is applied to each experimental signal $`f(t)`$ with aim to obtain the numerical characteristic of its shape (shape indicator, $`SI`$) defined as: $`SI=_kf(t_k)P(t_kt_o)`$, where the sum is over all time bins (from $`k`$ = 1 to $`k`$ = 2048), $`f(t_k)`$ is the digitized amplitude of a given signal (normalized to its area) at the time $`t_k`$. The weight function $`P(t_kt_o)`$ is determined as $`P(t)=\{\overline{f}_\alpha (t)\overline{f}_\gamma (t)\}/\{\overline{f}_\alpha (t)+\overline{f}_\gamma (t)\},`$ and $`t_o`$ is the time origin of the signal. The measured with sources $`SI`$ distributions are well described by a Gaussian functions, whose mean values and standard deviations $`\sigma _\alpha `$ and $`\sigma _\gamma `$ have a slight energy dependence<sup>§</sup><sup>§</sup>§For the $`\gamma `$-s (300–3200 keV) $`SI_\gamma `$= 18.09 – (4.5$`\times `$10$`{}_{}{}^{5}E_{\gamma }^{}`$), $`\sigma _\gamma `$= 2.61 – (4.7$`\times `$10$`{}_{}{}^{4}E_{\gamma }^{})+707/E_\gamma `$, while for the $`\alpha `$ particles (4000–6000 keV) $`SI_\alpha `$= 29.0; $`\sigma _\alpha `$= 5.11 – (5.52$`\times `$10$`{}_{}{}^{4}E_{\alpha }^{})+5520/E_\alpha `$. Here all variables are dimensionless ($`E_\gamma `$ and $`E_\alpha `$ are expressed in keV).. For 0.9 MeV $`\gamma `$ quanta $`SI_\gamma `$ = $`18\pm 3`$, while for 4.8 MeV $`\alpha `$ particles $`SI_\alpha `$ = $`29.0\pm 3.6`$. It allows us to determine the efficiency of the PS event selection for the different chosen intervals of $`SI`$ values ($`\pm \sigma `$, $`\pm 2\sigma `$, etc.).
The PS selection technique ensures the very important possibility to discriminate ”illegal” events: double pulses, $`\alpha `$ events, etc., and thus to suppress background. An example of a double pulse is shown in fig. 4a. Value of the shape indicator for the full signal is $`SI=`$$`47`$; for the first pulse $`SI_1`$ = 18.4 (hence it corresponds to $`\gamma `$ or $`\beta `$ particle), for the second pulse $`SI_2`$ = 37.4 ($`\alpha `$ particle). The energy release is 1.97 MeV, and without PS analysis it would be a candidate event for $`2\nu 2\beta `$ decay of <sup>116</sup>Cd.
Since the shape indicator characterizes the full signal, it is also useful to examine the pulse front edge. For example, it was found that at least 99% of ”pure” $`\gamma `$ events (measured with calibration <sup>232</sup>Th source) satisfy the following restriction on pulse rise time : $`\mathrm{\Delta }t`$($`\mu `$s) $``$ 1.24 – 0.5$`E_\gamma `$ \+ 0.078$`E_\gamma ^2`$, where $`E_\gamma `$ is dimensionless variable expressed in MeV. Hence, this filter was applied to the background data, and all events, which do not pass the test, were excluded from the residual $`\beta /\gamma `$ spectrum.
The results of PS analysis of the data are presented in fig. 5. The initial (without PS selection) spectrum of the <sup>116</sup>CdWO<sub>4</sub> scintillators in the energy region $`1.2`$$`4`$ MeV – collected during 4629 h in anticoincidence with active shield – is depicted in fig. 5a, while the spectrum after PS selection of the $`\beta /\gamma `$ events, whose $`SI`$ lie in the interval $`SI_\gamma 3.0\sigma _\gamma `$ $`SISI_\gamma +2.4\sigma _\gamma `$ and $`\mathrm{\Delta }t`$($`\mu `$s) $``$ 1.24 – 0.5$`E_\gamma `$ \+ 0.078$`E_\gamma ^2`$ (98% of $`\beta /\gamma `$ events), is shown in fig. 5b. From these figures the background reduction due to pulse-shape analysis is evident. Further, fig. 5c represents the difference between spectra in fig. 5a and 5b. These events, at least for the energy above 2 MeV, can be produced by <sup>228</sup>Th activity from the intrinsic contamination of the <sup>116</sup>CdWO<sub>4</sub> crystals (measured by the time-amplitude analysis as described above). Indeed, two decays in the fast chain <sup>212</sup>Bi ($`Q_\beta `$ = $`2.25`$ MeV) $``$ <sup>212</sup>Po ($`Q_\alpha `$ = $`8.95`$ MeV, $`T_{1/2}`$ = $`0.3`$ $`\mu `$s) $``$ <sup>208</sup>Pb can not be time resolved in the CdWO<sub>4</sub> scintillator (with an exponential decay time $``$15 $`\mu `$s ) and will result in one event. The example of such an event – recorded by the PS acquisition system – is depicted in fig. 4b. To determine the residual activity of <sup>228</sup>Th in the crystals, the response function of <sup>116</sup>CdWO<sub>4</sub> detectors for the <sup>212</sup>Bi $``$ <sup>212</sup>Po $``$ <sup>208</sup>Pb chain was simulated with the help of GEANT3.21 code and event generator DECAY4. The simulated function is shown in fig. 5c, from which one can see that the high energy part of the experimental spectrum is well reproduced ($`\chi ^2`$ = $`1.3`$) by the expected response for <sup>212</sup>Bi $`^{212}`$Po $`^{208}`$Pb decaysThe rest of spectrum below 1.9 MeV (fig. 5c) can be explained as high energy tail of the PS selected $`\alpha `$ particles (see fig. 6).. Corresponding activity of <sup>228</sup>Th inside the <sup>116</sup>CdWO<sub>4</sub> crystals, deduced from the fit in the 1.9–3.7 MeV energy region, is 37(4) $`\mu `$Bq/kg, that is in a good agreement with the value determined by the time-amplitude analysis of the chain <sup>220</sup>Rn $``$ <sup>216</sup>Po $``$ <sup>212</sup>Pb. Besides, the front edge analysis of 80 events with the energy $`2.0`$$`4.2`$ MeV ($`SI`$ $`SI_\gamma `$ \+ 2.54$`\sigma _\gamma `$; $`\mathrm{\Delta }t`$ $``$ 0.2 $`\mu `$s) was fulfilled and the half-life derived from the average time delay between the first and second part of the signal (see fig. 4b) is $`T_{1/2}`$ = $`0.31(6)`$ $`\mu `$s, in agreement with the <sup>212</sup>Po table value $`T_{1/2}`$ = $`0.299(2)`$ $`\mu `$s .
Fig. 6 represents the spectrum after PS selection of the background events, whose $`SI`$ lie in the interval $`SI_\gamma +2.4\sigma _\gamma <`$ $`SI`$ $`<`$ $`SI_\alpha +2.4\sigma _\alpha `$ ($``$90% of $`\alpha `$ events). The obtained distribution with maximum at 0.95 MeV is well reproduced by the model, which includes all $`\alpha `$ particles from chains in <sup>232</sup>Th and <sup>238</sup>U families. The total $`\alpha `$ activity of the <sup>116</sup>CdWO<sub>4</sub> crystals deduced from fig. 6 is 1.4(3) mBq/kg. This value can be adjusted with the activities determined by the time-amplitude analysis under usual (for crystals) assumption that secular radioactive equilibriums in some chains of <sup>232</sup>Th and <sup>238</sup>U families (like, f. e. <sup>230</sup>Th $``$ <sup>226</sup>Ra chain) are broken.
## III Results and discussion
### A Two-neutrino double beta decay of <sup>116</sup>Cd
To determine the half-life of two-neutrino 2$`\beta `$ decay of <sup>116</sup>Cd, the background in the energy interval $`900`$$`2900`$ keV was simulated with the help of GEANT3.21 package and event generator DECAY4. In addition to <sup>116</sup>Cd two neutrino 2$`\beta `$ decay distribution, only three components shown in fig. 2 were used to build up the background model: <sup>40</sup>K contamination of the enriched and natural CdWO<sub>4</sub> scintillators, whose activity limits of less than 4 mBq/kg were established earlier , and external $`\gamma `$ background caused by <sup>232</sup>Th and <sup>238</sup>U contamination of the PMTs (one PMT for <sup>116</sup>CdWO<sub>4</sub> crystals; one for CdWO<sub>4</sub>; two for plastic active shield)The radioactive impurities of all PMTs used in the installation were previously measured by R&D low background set-up as (0.4–2.2) Bq/PMT and (0.1– 0.2) Bq/PMT for <sup>226</sup>Ra and <sup>228</sup>Th activity, respectively .. This simple background model describes experimental data in the chosen energy interval $`900`$$`2900`$ keV reasonably well ($`\chi ^2`$ =$`1.3`$) and gives the following results: the activities of <sup>40</sup>K inside the enriched and natural CdWO<sub>4</sub> crystals are equal 0.8(2) and 2.1(3) mBq/kg, respectively; the half-life of two-neutrino 2$`\beta `$ decay of <sup>116</sup>Cd is $`T_{1/2}`$(2$`\nu `$) = 2.6(1)$`\times 10^{19}`$ yr (only statistical uncertainties are given, while systematical errors are pointed below).
Taking advantage of the high statistics in our experiment (approximately 3600 events of <sup>116</sup>Cd two neutrino 2$`\beta `$ decay are contained within the interval $`900`$$`2900`$ keV), we can prove our model with the help of experimental $`2\nu 2\beta `$ decay Kurie plot: $`K(\epsilon )=[S(\epsilon )/\left\{(\epsilon ^4+10\epsilon ^3+40\epsilon ^2+60\epsilon +30)\epsilon \right\}]^{1/5}`$, where $`S`$ is the number of events with the energy $`\epsilon `$ (in electron mass units) in the experimental spectrum after background subtraction. For the real $`2\nu 2\beta `$ decay events such a Kurie plot should be the straight line $`K(\epsilon )(Q_{2\beta }\epsilon ),`$ where $`Q_{2\beta }`$ is the 2$`\beta `$ energy release. From fig. 7, where our experimental Kurie plot is depicted, one can see that in the region $`1.1`$$`2.4`$ MeV it is well fitted by the straight line with $`Q_{2\beta }=2790(87)`$ keV (table value is $`Q_{2\beta }=2805(4)`$ keV). To take into account the energy resolution of the detector the experimental spectrum was also fitted by the convolution of the theoretical $`2\nu 2\beta `$ distribution $`\rho (\epsilon )=A\epsilon (\epsilon ^4+10\epsilon ^3+40\epsilon ^2+60\epsilon +30)(Q_{2\beta }\epsilon )^5`$ with the detector resolution function (Gaussian with $`FWHM`$ determined as described above) with $`A`$ and $`Q_{2\beta }`$ values as free parameters. This fit in the energy region $`1.2`$$`2.8`$ MeV yields a very similar value $`Q_{2\beta }=2779(52)`$ keV and an $`A`$ corresponding to $`T_{1/2}`$(2$`\nu )=`$ 2.5(3)$`\times `$1$`0^{19}`$ yr, thus justifying our assumption that experimental data in the region above 1.2 MeV are related mainly with <sup>116</sup>Cd two neutrino $`2\beta `$ decay. In fact, a signal to background ratio deduced from our data is 4:1 for the energy interval $`1.2`$$`2.9`$ MeV, and 15:1 for the energy range $`1.6`$$`2.9`$ MeV, which are higher than those reached up-to-date in other $`2\beta `$ decay experiments .
To estimate systematical uncertainties of the measured half-life, different origins of errors were taken into account, whose contributions are listed in Table I. The final value is equal to:
$`T_{1/2}(2\nu )=2.6\pm 0.1`$(stat)$`{}_{}{}^{+0.7}{}_{0.4}{}^{}`$(syst)$`\times `$1$`0^{19}`$ yr.
Our result is in agreement with those measured earlier ($`T_{1/2}(2\nu )=2.6_{0.5}^{+0.9}\times 10^{19}`$ yr and $`T_{1/2}(2\nu )=2.7_{0.4}^{+0.5}`$(stat)$`{}_{}{}^{+0.9}{}_{0.6}{}^{}`$(syst)$`\times 10^{19}`$ yr ) and disagrees to some extent with the value $`T_{1/2}(2\nu )=3.75\pm 0.35`$(stat)$`\pm 0.21`$(syst)$`\times 10^{19}`$ yr from ref. <sup>\**</sup><sup>\**</sup>\**Note, that in the quite small detection efficiency ($`1.73\%`$) was calculated by the Monte Carlo method without experimental test, thus perhaps systematical error could be higher than quoted value..
### B New limits for 0$`\nu `$2$`\beta `$ decay of <sup>116</sup>Cd to ground state of <sup>116</sup>Sn
To estimate the half-life limit for the neutrinoless decay mode, a simple background model was used. In fact, in the $`0\nu 2\beta `$ decay energy region only three background contributions are important: (i) external $`\gamma `$ background from U/Th contamination of the PMTs; (ii) tail of the $`2\nu 2\beta `$ decay spectrum; and (iii) internal background distribution expected from the <sup>212</sup>Bi $``$ <sup>212</sup>Po $``$ <sup>208</sup>Pb decay (<sup>228</sup>Th chain). As it was shown above, two decays in the fast chain <sup>212</sup>Bi $``$ <sup>212</sup>Po $``$ <sup>208</sup>Pb really create the background in the region of $`0\nu 2\beta `$ decay (see fig. 5c). For the activity of <sup>228</sup>Th inside the <sup>116</sup>CdWO<sub>4</sub> crystals two values were obtained: 38(3) $`\mu `$Bq/kg (time-amplitude method) and 37(4) $`\mu `$Bq/kg (pulse-shape analysis). Hence, in the limit of statistical errors, we do not find an indication of a failure for the rejection of $`\alpha `$ pulses by our PS analysis.
The high energy part of the experimental spectrum of the <sup>116</sup>CdWO<sub>4</sub> crystals measured in anticoincidence with the shielding detectors and after the time-amplitude and pulse-shape selection is shown in fig. 8. The peak of $`0\nu 2\beta `$ decay is absent, thus from the data we obtain a lower limit of the half-life: $`lim`$ $`T_{1/2}=\mathrm{ln}2N\eta t/limS,`$ where $`N=4`$.66$`\times `$1$`0^{23}`$ is number of <sup>116</sup>Cd nuclei, $`t`$ is the measuring time ($`t=4629`$ h), $`\eta `$ is the total detection efficiency for $`0\nu 2\beta `$ decay, and $`limS`$ is the number of events in the peak which can be excluded with a given confidence level. The value of the detection efficiency $`\eta _{MC}=0.83`$ was calculated by the DECAY4 and GEANT3.21 codes, while the efficiency of the PS analysis $`\eta _{PS}=0.98`$ was determined as described above, thus the total efficiency $`\eta =\eta _{MC}\eta _{PS}=0.81`$. To obtain the value of $`limS,`$ the part of the spectrum in the $`1.9`$$`3.8`$ MeV region was fitted by the sum of the simulated $`0\nu 2\beta `$ peak and three background functions (external $`\gamma `$ rays from PMT-s contamination; contribution from <sup>212</sup>Bi $``$ <sup>212</sup>Po $``$ <sup>208</sup>Pb intrinsic chain; and $`2\nu 2\beta `$ decay tail). This fit gives the value of $`S=1.1\pm 1.2`$ counts, which corresponds – in accordance with Feldman-Cousins procedure for results close to the edge of physically accepted area recommended by the Particle Data Group (PDG) – to a $`limS=0.9(0.3)`$ counts with 90%(68%) C.L., and subsequently to $`T_{1/2}`$(0$`\nu `$2$`\beta )1.5(4.6`$)$`\times `$1$`0^{23}`$ yr at 90%(68%) C.L. However, because of the low statistics in the energy range where the effect is expected, the obtained values can be cross-checked in a more simple and explicit way. Indeed, in the energy interval $`2.6`$$`3.1`$ MeV (containing 91% of 0$`\nu `$2$`\beta `$ peak) there is only one measured event, while the background expected on the basis of the GEANT simulation, in the same energy region is 3.2$`{}_{}{}^{+2.1}{}_{1.1}{}^{}`$ counts (1.9$`\pm `$0.7 events from PMT contamination; 0.4$`\pm `$0.1 events from 2$`\nu `$2$`\beta `$ distribution; 0.9$`{}_{}{}^{+2}{}_{0.9}{}^{}`$ counts from mentioned <sup>212</sup>Bi $``$ <sup>212</sup>Po $``$ <sup>208</sup>Pb chain). Following the PDG recommendation we can derive from these numbers the excluded limit as $`limS=1.8(0.5)`$ with 90%(68%) C.L., which leads to $`T_{1/2}`$(0$`\nu `$2$`\beta )0.7(2.5`$)$`\times `$1$`0^{23}`$ yr at 90%(68%) C.L. confirming the previous estimate. Finally, the following values were accepted as conservative half-life limits for neutrinoless $`2\beta `$ decay of <sup>116</sup>Cd:
$`T_{1/2}`$(0$`\nu `$2$`\beta )0.7(2.5`$)$`\times `$1$`0^{23}`$ yr, 90%(68%) C.L.
Using calculations , one can obtain restrictions on the neutrino mass and right-handed admixtures in the weak interaction: $`m_\nu 3.0`$ eV, $`\eta 3.`$9$`\times `$10<sup>-8</sup>, $`\lambda 3.`$4$`\times `$1$`0^6`$ at 90% C.L., and neglecting right-handed contribution $`m_\nu 2.6(1.4)`$ eV at 90% (68%) C.L. On the basis of calculations we get a similar result: $`m_\nu 2.4(1.3)`$ eV at 90%(68%) C.L. In accordance with ref. the value of the R-parity violating parameter of minimal SUSY standard model is restricted by our $`T_{1/2}`$ limit to $`\epsilon 8.8(6.4`$)$`\times `$1$`0^4`$ at 90%(68%) C.L. (calculations give more stringent restrictions: $`\epsilon 3.4(2.4`$)$`\times `$1$`0^4`$).
### C The bounds on 0$`\nu `$2$`\beta `$ decay of <sup>116</sup>Cd to excited levels of <sup>116</sup>Sn
Not only ground state (g.s.) but also excited levels of <sup>116</sup>Sn with $`E_{lev}Q_{2\beta }`$ can be populated in 0$`\nu `$2$`\beta `$ decay of <sup>116</sup>Cd. In this case one or several $`\gamma `$ quanta, conversion electrons and/or e<sup>+</sup>e<sup>-</sup> pairs will be emitted in a deexcitation process, in addition to two electrons emitted in 2$`\beta `$ decay. The response functions of <sup>116</sup>CdWO<sub>4</sub> detectors for $`0\nu 2\beta `$ decay to the first and second excited levels of <sup>116</sup>Sn (2$`{}_{}{}^{+}{}_{1}{}^{}`$ with $`E_{lev}=1294`$ keV and 0$`{}_{}{}^{+}{}_{1}{}^{}`$ with $`E_{lev}=1757`$ keV) were simulated with the help of GEANT3.21 and DECAY4 codes. The full absorption of all emitted particles should result in the peak with $`E=Q_{2\beta }`$ (practically the same peak as it is expected for the 0$`\nu 2\beta `$ decay of <sup>116</sup>Cd to the g.s. of <sup>116</sup>Sn). Calculated full peak efficiencies are: $`\eta (2_1^+)=0.14`$ and $`\eta (0_1^+)=0.07`$. These numbers and the value of $`limS=1.8(0.5)`$ with 90%(68%) C.L. (determined for the g.s.$``$g.s. transition) give the following restrictions on half-lives of <sup>116</sup>Cd neutrinoless $`2\beta `$ decay to excited levels of <sup>116</sup>Sn:
$`T_{1/2}`$(g.s.$``$ $`2_1^+)1.3(4.8`$)$`\times `$1$`0^{22}`$ yr, 90%(68%) C.L.,
$`T_{1/2}`$(g.s.$``$ $`0_1^+)0.7(2.4`$)$`\times `$1$`0^{22}`$ yr, 90%(68%) C.L.
### D Neutrinoless 2$`\beta `$ decay with Majoron(s) emission
The procedure to obtain half-life limits for 0$`\nu `$2$`\beta `$ decay with one (two) Majoron(s) emission was carried out in two steps as follows. First, because in the measured spectrum contributions of <sup>40</sup>K are negligible above the energy 1.6 MeV, the data were fitted in the energy region 1.6–2.8 MeV for 0$`\nu `$M1 mode (1.6–2.6 MeV for 0$`\nu `$M2) by using only three theoretical distributions: $`\gamma `$ background from measured PMT-s contamination (<sup>226</sup>Ra and <sup>232</sup>Th chains) and two neutrino 2$`\beta `$ decay of <sup>116</sup>Cd, as background, and 0$`\nu `$2$`\beta `$ decay with one (two) Majoron(s) emission, as effect. With this simple model the $`\chi ^2`$ value was equal to 1.1 both for 0$`\nu `$M1 and 0$`\nu `$M2 fits. As a result, the number of events under a theoretical 0$`\nu `$M1 curve was determined as $`9\pm 21`$, giving no statistical evidence for the effect. It leads to an upper limit of 41(26) events at 90%(68%) C.L., that together with an efficiency value $`\eta =0.905`$ corresponds to the half-life limit $`T_{1/2}`$(0$`\nu `$M1) $`3.7(5.9`$)$`\times `$1$`0^{21}`$ yr. A similar procedure for 0$`\nu `$2$`\beta `$ decay with two Majorons emission gives $`T_{1/2}`$(0$`\nu `$M2) $`5.9(9.4`$)$`\times `$1$`0^{20}`$ yr at 90%(68%) C.L. The part of the experimental spectrum and theoretical 0$`\nu `$M1 and 0$`\nu `$M2 distributions with half-lives equal to these limits are shown in fig. 8. On one hand, the obtained results can be treated as conservative because the accepted background model consists of only two origins, the external background from U/Th contamination of the PMT and 2$`\nu `$2$`\beta `$ decay distribution, while in the energy region of interest some other sources of background, such as the mentioned <sup>212</sup>Bi $``$ <sup>212</sup>Po $``$ <sup>208</sup>Pb chain contribution from the intrinsic impurities of the crystals, could enlarge the 0$`\nu `$M1 and 0$`\nu `$M2 limits. On the other hand, the uncertainty of the 2$`\nu `$2$`\beta `$ decay half-life of <sup>116</sup>Cd could lead to the overestimated value of the 0$`\nu `$M1 bound. To avoid the last possibility we have estimated the $`T_{1/2}`$(0$`\nu `$M1) limit in a more explicit way. With this aim the energy interval 2.5–2.8 MeV – where the tail of 2$`\nu `$2$`\beta `$ decay distribution is practically zero (see fig. 8) – was used as the most sensitive region for the 0$`\nu `$M1 double $`\beta `$ decay search. In this energy interval the number of measured counts is 3, while the expected contribution (in the same energy range) from the PMT contamination is 3.2$`\pm `$1.0 events and from 2$`\nu `$2$`\beta `$ decay is 1.5$`\pm `$0.5 events, thus the total expected background is 4.7$`\pm `$1.1 counts. Following the PDG recommendation we can derive from these numbers the excluded limit for the effect being sought as 3.0 events with 90% C.L. Taking into account that the interval 2.5–2.8 MeV contains 8.9% of full 0$`\nu `$M1 curve, it yields a limit of $`T_{1/2}`$(0$`\nu `$M1) $``$4.5$`\times `$1$`0^{21}`$ yr (90% C.L.), confirming the preceding estimate:
$`T_{1/2}`$(0$`\nu `$M1) $`3.7(5.9`$)$`\times `$1$`0^{21}`$ yr, 90%(68%) C.L.,
$`T_{1/2}`$(0$`\nu `$M2) $`5.9(9.4`$)$`\times `$1$`0^{20}`$ yr, 90%(68%) C.L.
Both present half-life limits are more stringent than those established in our previous measurement during 19986 h and in the NEMO experiment .
The probability of neutrinoless 2$`\beta `$ decay with Majoron emission can be expressed as: $`\left\{T_{1/2}(0\nu M1)\right\}^1`$= $`<g_M>^2`$NME$`^2G,`$ where $`<g_M>`$ is the effective Majoron-neutrino coupling constant, NME is the nuclear matrix element and $`G`$ is the kinematical factor. Using our result $`T_{1/2}`$(0$`\nu `$M1) $`3.7(5.9`$)$`\times `$1$`0^{21}`$ yr and values of $`G`$ and NME calculated in the QRPA model with proton-neutron pairing we obtain $`g_M12(9.5`$)$`\times `$1$`0^5`$ ($`g_M6`$.5(5.4)$`\times `$1$`0^5`$ on the basis on calculation ) with 90%(68%) C.L., which is one of the best restriction up-to-date obtained in the direct 2$`\beta `$ decay experiments .
## IV Conclusion
The search for <sup>116</sup>Cd double $`\beta `$ decay with enriched <sup>116</sup>CdWO<sub>4</sub> scintillators has entered in a new phase. The set-up with four <sup>116</sup>CdWO<sub>4</sub> crystals (339 g) is running since October 1998. Improved passive shield, new active shield made of fifteen CdWO<sub>4</sub> crystals (total mass 20.6 kg), as well as time-amplitude and pulse-shape analysis of the data result in the reduction of the background rate in the $`2.5`$$`3.2`$ MeV region to 0.03 counts/yr$``$kg$``$keV. This reduction, together with an about threefold increase in the number of <sup>116</sup>Cd nuclei, leads to the substantial sensitivity enhancement of the <sup>116</sup>Cd experiment by more than one order of magnitude. Due to that the neutrino mass limit of $`m_\nu 2.6(1.4)`$ eV at 90%(68%) C.L. was set after the first 4629 h run.
In August 1999 one of our <sup>116</sup>CdWO<sub>4</sub> crystals was annealed at high temperature, and its light output was increased by $``$13%. The PMT of the main <sup>116</sup>CdWO<sub>4</sub> detectors was changed by the special low background EMI tube (5 inches in diameter) with the RbCs photocathode, whose spectral response better fits the CdWO<sub>4</sub> scintillation light. As a result, the spectrometric parameters of four crystals taken as a whole were improved. In particular, the energy resolution of the main detector is now 11.4% at 1064 keV and 8.6% at 2615 keV (comparing with those before this upgrading: 14.5% and 11%). Besides, the PS discrimination ability of the detector was improved too, as it is visible from fig. 9, where the $`SI`$ distribution of the measured background events – before and after the last upgrading – is depicted. It is expected that after approximately 5 years of measurements the half-life limit $`T_{1/2}`$(0$`\nu `$2$`\beta )`$ 4$`\times `$1$`0^{23}`$ yr will be reached which corresponds to $`m_\nu 1.2`$ eV. The bounds on neutrinoless 2$`\beta `$ decay with Majorons emission and 2$`\beta `$ transitions to the excited levels of <sup>116</sup>Sn would be improved too.
###### Acknowledgements.
The authors express their gratitude to M. Bini and O. Vihliy for their efforts to develop and test new data acquisition system for the experiment.
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# KAON PHOTOPRODUCTION OF THE DEUTERON AND THE 𝑃-MATRIX APPROACH TO THE 𝑌𝑁 FINAL STATE INTERACTION
Strangeness photoproduction of the deuteron is investigated theoretically making use of the covariant reaction formalism and the $`P`$–matrix approach the final state hyperon-nucleon interaction. Remarkably simple analytical expression for the amplitude is obtained. Pronounced effects due to final state interaction are predicted.
PACS numbers: 25.20 Lj, 25.30 Rw, 13.60. Le, 13.60 Rj
Up to now most of the information on the hyperon–nucleon $`(YN)`$ interaction has been obtained either from hypernuclei or from $`K^{}d`$ and $`\pi ^+d`$ reactions. After several decades of studies our knowledge on the $`YN`$ system is still far from being complete.
Recently the interest to the $`\mathrm{\Lambda }N`$$`\mathrm{\Sigma }N`$ system has flared again in connection to the expected CEBAF experimental results on the kaon photoproduction on the deuterium . The final state $`YN(Y=\mathrm{\Lambda },\mathrm{\Sigma })`$ interaction (FSI) plays an important role in the $`\gamma dK^+YN`$ reaction. Therefore high resolution photoproduction experiments can substantially deepen our understanding of the $`YN`$ dynamics.
The problem of FSI in $`\gamma dK^+Yn`$ reaction has been addressed by several authors starting from the pioneering paper by Renard and Renard . Two novel features differ the present work from the previous studies. First, covariant formalism based on direct evaluation of Feynman diagrams is used which allows to analyse the data beyond the region of low energy and low momentum transfer. Second, the $`YN`$ interaction is described within the $`P`$-matrix approach which takes into account the subnuclear degrees of freedom and disentangle the dynamical singularities from kinematical threshold effects . The $`P`$–matrix analysis of the $`YN`$ interaction was presented in (see also ) and we shall use the set of parameters from .
In our previous publication we have presented some preliminary results without proving the central assertion, namely that the FSI effects allow remarkably simple evaluation within the $`P`$-matrix approach. The proof of this statement along with the presentation of the covariant reaction formalism contributes the core of the present publication. Tooled with the methods presented below one can easily treat FSI within any other approach making use of the known relation between $`P`$–matrix and potential approaches .
The double differential cross section of the reaction $`\gamma dK^+Yn`$ reads
$$\frac{d^2\sigma }{d|𝐩_K|d\mathrm{\Omega }_K}=\frac{1}{2^{11}\pi ^5}\frac{𝐩_K^2}{kM_dE_K}\frac{\lambda ^{1/2}(s_2,m_Y^2,m_n^2)}{s_2}𝑑\mathrm{\Omega }_{Yn}^{}|T|^2.$$
(1)
Here $`k`$, $`𝐩_K^2`$, $`E_K`$ and $`\mathrm{\Omega }_K`$ correspond to the deuteron rest system with $`z`$-axis defined by the incident photon beam direction $`𝐤`$. The solid angle $`\mathrm{\Omega }_{Yn}^{}`$ is defined in the $`Yn`$ center-of-momentum system, $`s_2=(p_Y+p_N)^2,`$ $`\lambda (x,y,z)=x^22(y+z)x+(yz)^2`$.
The fully covariant analogue of (1) valid in any reference frame has the form
$$\frac{d^2\sigma }{ds_2dt_1}=\frac{1}{2^{10}\pi ^4\lambda (s,0,M_d^2)}𝑑t_2𝑑s_1\frac{|T(t_1,s_1,s_2,t_2)|^2}{[\mathrm{\Delta }_4(t_1,s_1,s_2,t_2)]^{1/2}},$$
(2)
where $`s=(k+p_d)^2,s_1=(p_K+p_Y)^2,s_2=(p_Y+p_N)^2,t_1=(kp_K)^2,t_2=(p_dp_n)^2,`$ and $`\mathrm{\Delta }_4`$ is a $`4\times 4`$ symmetric Gram determinat . The region of integration in (2) has to satisfy $`\mathrm{\Delta }_40`$. The number of essential final state Lorentz scalar variables is 4¡ namely $`t_1,s_1,s_2,t_2`$.
The amplitude $`T`$ will be approximated by the two dominant diagrams, namely the tree (pole, or plane waves) graph and the loop (triangle) diagram with $`YN`$ FSI.Consider first the tree diagram. Two blocks entering into it have to be specified: (i) the elementary photoproduction amplitude $`M^{\gamma K}`$ on the proton, and (ii) the deuteron vertex $`\mathrm{\Gamma }_d`$. The elementary amplitude used in the present calculation was derived from the tree level effective Lagrangian with the account of several resonances in $`s,t`$ and chanels . This amplitude has the following decomposition over invariant terms
$$M^{\gamma K}=\overline{u}_Y\underset{j=1}{\overset{6}{}}𝒜_j_j(s^{},t^{},u^{})u_p,$$
(3)
where $`s^{}=(k+p_p)^2`$, $`t^{}=(kp_K)^2`$, $`u^{}=(kp_Y)^2`$.
The decomposition of the deuteron vertex function $`\mathrm{\Gamma }_d`$ in independent Lorentz structures has the form
$`\mathrm{\Gamma }_d`$ $`=`$ $`\sqrt{m_N}\left[(p_p+p_n)^2M_d^2\right]\left[\phi _1(t_2){\displaystyle \frac{(p_pp_n)_\mu }{2m_N^2}}+\phi _2(t_2){\displaystyle \frac{1}{m_N}}\gamma _\mu \right]^\mu (p_d,\lambda ).`$ (4)
Here $`^\mu (p_d,\mu )`$ is the polarization 4-vector of the deuteron with momentum $`p_d`$ and polarization $`\lambda `$. The vertex (4) implies that the deuteron as well as the spectator neuteron are on mass shell while the proton is off its mass shell. Now we can write the following expression for the tree diagram
$$T^{(t)}=\overline{u}_Y\left\{\left(\underset{j=1}{\overset{6}{}}𝒜_j_j(s^{},t^{},u^{})\right)S(p_p)\mathrm{\Gamma }_d\right\}u_n^c,$$
(5)
where $`S(p_p)`$ is the proton propagator and $`u_n^c`$ is a charge conjugated neutron spinor. The deuteron vertex in (5) may be substituted by the relativistic deuteron wave function according to $`\psi _d=[2(2\pi )^3M_d]^{1/2}S(p_p)\mathrm{\Gamma }_d`$ . Then (5) can be rewritten as
$$T^{(t)}=[(2\pi )^32M_d]^{1/2}M^{\gamma K}\psi _d,$$
(6)
where $`\psi _d`$ is the relativistic deuteron wave functions discussed at length in , and where summation over magnetic quantum numbers is tacitly assumed.
The loop diagram with $`YN`$ rescattering is given by the expression
$$T^{(l)}=\frac{d^4p_n}{(2\pi )^4}\overline{u}_Y(p_Y^{})\left\{\left(\underset{j=1}{\overset{6}{}}𝒜_j_j\right)S(p_p)\mathrm{\Gamma }_dCS(p_n)T_{YN}S(p_Y)\right\}\overline{u}(p_N^{}).$$
(7)
Here $`C`$ is the charge-conjugation matrix, $`T_{YN}`$ is a hyperon-nucleon vertex, this vertex being “dressed” by corresponding spinors constitutes the hyperon-nucleon amplitude $`F_{YN}`$.
The comprehensive treatment of the loop diagram will be presented in the forthcoming publication while here we resort to a simple approximation with the aim to exposure the effects of the FSI. Only positive frequency components are kept in the propagators $`S(p_n)`$ and $`S(p_Y)`$ in (7), while the propagator $`S(p_p)`$ together with $`\mathrm{\Gamma }_d`$ is lumped into the relativistic deuteron wave function $`\psi _d`$ . Then integration over the time component $`dp_n^0`$ is performed. Thus we arrive at the following expression for $`T^{(l)}`$
$$T^{(l)}=[(2\pi )^32M_d]^{1/2}\frac{d𝐩_n}{(2\pi )^3}\frac{M^{\gamma K}\psi _d(p_n)F_{YN}(E_{YN};q,q^{})}{𝐪^2𝐪^2i0},$$
(8)
where $`𝐪`$ and $`𝐪^{}`$ are the $`YN`$ c.m. momenta before and after the $`YN`$ FS, $`𝐪=𝐩_n\frac{1}{2}(𝐤𝐩_k)`$. The quality $`F_{YN}(E_{YN};qq^{})`$ is the half-off-shell YN amplitude at $`E_{YN}=𝐪_{}^{}{}_{}{}^{2}/2\mu _{YN}𝐪^2/2\mu _{YN}`$. The use of the nonrelativistic propagator in (8) is legitimate since FSI is essential at low relative YN momenta. The arguments of the elementary amplitude $`M^{\gamma K}`$ are specified in (3) but in the kinematical region where $`YN`$ FSI is essential $`M^{\gamma K}`$ can be considered as point-like and hence $`M^{\gamma k}`$ can be factored out of the integral (8) at the values of the arguments fixed by the energies and momenta of the initial and final particles (i.e. at the values of the arguments corresponding to the plane-waves diagram).
Next we consider the $`YN`$ rescattering amplitude $`F_{YN}`$ and recall the interpretation of the $`P`$-matrix in terms of the underlying coupled–channel quark–hadron potential . Namely, the $`P`$-matrix description is equivalent to the coupling between hadron and quark channels via the nonlocal energy dependent potential of the form
$$V_{hqh}=\underset{n}{}\frac{f_n(r)f_n(r^{})}{EE_n},$$
(9)
where $`E_n`$ are the energies of the six-quark ”primitives” ($`P`$-matrix poles), and the form-factors are given by $`f_n(r)=\lambda _n\delta (rb)`$, where $`b`$ is the bag radius and the coupling constants $`\lambda _n`$ are related to the residues of the $`P`$-matrix via
$$P=P_0+\underset{n}{}\frac{\lambda _n^2}{E_nE}.$$
(10)
As it was shown in a single primitive at $`E_n=2.34GeV`$ is sufficient for a high quality description of the existing $`YN`$ experimental data. In order to avoid lengthy equations we consider in the present note the region close to $`\mathrm{\Lambda }N`$ threshold while the preliminary results which included $`\mathrm{\Lambda }NN`$ coupling were announced in and will be treated within the present approach in a forthcoming detailed publication.
The momentum–space half-of-shell amplitude $`F_{YN}(𝐪,𝐩;E_{YN})`$ corresponding to the potential (9) reads
$$F(E_{YN};q,q^{})=\lambda _n^2b^2\frac{\mathrm{sin}qb}{qb}d^1(E_{YN})\frac{sinq^{}b}{q^{}b},$$
(11)
$$d(E_{YN})=E_{YN}E_n+\frac{\lambda _n^2}{q}e^{iq^{}b}\mathrm{sin}q^{}b.$$
(12)
This form of the amplitude allows to perform momentum integration in (8) analytically with an accuracy sufficient to display the effects due to FSI. In the complex $`q`$ plane the integral (8) picks up contributions from the poles of the propagator and the sigularities of the deuteron wave function. In a separate detailed publication we show both numerically and using models for $`\psi _d`$ that the contribution of the deuteron wave function singularities does not exceed 20% (this conclusion can be immediately verified considering the simplest singularity of $`\psi _d`$ at $`q^2=\alpha ^2)`$. Thus substituting (11-12) into (8) and performing integration with the above remark in mind we get
$$T^{(l)}[(2\pi )^32M_d]^{1/2}M^{\gamma K}\psi _d(p_n^{})\left[\frac{1}{f(q^{})}1\right],$$
(13)
where $`p_n^{}`$ is related to $`q^{}`$ in the same way as $`p_n`$ to $`q`$ (see above), and where
$$f(q^{})=1\frac{\lambda _n^2b}{\mathrm{\Delta }q^2/2\mu _{YN}}e^{iq^{}b}\frac{\mathrm{sin}q^{}b}{q^{}b},$$
(14)
and $`\mathrm{\Delta }=E_nm_Ym_N,\mu _{YN}`$ is the $`YN`$ reduced mass. One can easily verify that $`f(q^{})`$ is the Jost function corresponding to the potential (9) (some authors use the notation $`f(q^{})`$). From (14) and (6) a remarkably simple expression for the sum of the tree and loop diagrams follows
$$T^{(t)}+T^{(l)}=T^{(t)}/f(q^{}),$$
(15)
where the final state $`YN`$ momentum $`q^{}`$ is expressed in terms of the $`YN`$ invariant mass $`s_2^{1/2}`$ as $`q^{}=\frac{\lambda ^{1/2}(s_2,m_Y^2,m_n^2)}{2s_2^{1/2}}.`$
Expression (14) is physically absolutely transparent: one immediately recognizes the standard enhancement factor given by the inverse Jost function corresponding to the potential (9). Inclusion of the $`\mathrm{\Lambda }NN`$ coupling is straightforward . The above equations may be used beyond the $`P`$-matrix approach to the $`YN`$ interaction since there is a well trotted path connecting $`P`$-matrix and $`S`$-matrix approaches.
In conclusion we present the results of the calculation obtained using equations covariant equation (2), Eqs. (6), and (14). Use has been made of the elementary photoproduction amplitude from , the deuteron vertex function taken from the relativistic Gross model and the plane-waves diagram with this input was calculated in . In Fig.1 the double differential cross section (1) is shown as a function of the photon energy in the $`\mathrm{\Lambda }n`$ invariant mass region close to the threshold ($`2.05GeV\sqrt{s_{\mathrm{\Lambda }n}}2.10GeV)`$. The enchanement of the cross section close to threshold due to FSI is quite pronounced. We remind that according to the $`{}_{}{}^{3}S_{1}^{}`$ $`\mathrm{\Lambda }n`$ scattering length is 1.84 fm in line with some versions of the Nijmegen potential .
The author would like to thank V.A.Karmanov for fruitful discussions and suggestions. Valuable remarks by T.Mizutani and A.E.Kudryavtsev are gratefylly acknowledged. This work was possible due to stimulating contacts with C.Fayard, G.H.Lamot, F.Rouvier and B.Saghai. Hospitality and financial support from the University Claude Bernard and DAPNIA (saclay) as well as support from RFFI grant 970216406 are gratefully acknowledged.
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Fig. 1. The double diffential cross section as a function of the photon energy for
$`p_K=0.861GeV`$, $`\theta _{\gamma K}=0^0`$. The dased line is the plane waves contribution,
the solid line incorporates FSI according to Eq. (14).
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# 1 Introduction
## 1 Introduction
It has been appreciated for a long time that the construction of consistent superstring theories requires appropriate projections on the underlying conformal field theory, most prominently the GSO projection. A careful implementation of these projections is e.g. required when one computes the massless spectrum of such theories. A clear conceptual understanding becomes even more mandatory when it comes to specifying boundary conditions for open strings. The analysis of the interplay between projections in superstring theories and boundary conditions for open strings is our main concern in this paper. We will concentrate on compactifications of type II superstring theories in which the internal part is an $`N=\mathrm{\hspace{0.17em}2}`$ rational conformal field theory, among them in particular the Gepner models. Our approach enables us e.g. to derive formulas for Gepner model boundary states entirely from well-established principles. Where comparable with the literature, these results differ from the formulas obtained elsewhere except for some particularly simple models.
Let us be more explicit. In the construction of a superstring compactification one starts with specifying the vacuum configuration. This amounts to choosing a conformal field theory $`𝒞_{\mathrm{int}}`$ for the ‘inner’ or ‘internal’ sector. $`𝒞_{\mathrm{int}}`$ must satisfy a number of consistency constraints, such as possessing the correct Virasoro central charge, enough supersymmetry on the world sheet, and modular invariance. Afterwards, additional projections need to be imposed on $`𝒞_{\mathrm{int}}`$. This includes in particular the GSO projection, which ensures space-time supersymmetry. But there is another generic projection, too, which in the case of flat backgrounds looks quite innocent and which nevertheless will play an important role below. Namely, the total conformal field theory in question is a tensor product of the inner sector with the flat space-time part, and the constraint that necessitates a projection is that the spin structures for all fermionic fields on the world sheet must be aligned. In other words, the fermionic fields have to be either in the Ramond or in the NeveuSchwarz sector simultaneously in each factor of the tensor product. As we will see, the interplay between these two projections has quite non-trivial consequences in non-flat backgrounds, in particular for the description of boundary states.
Now it is well-known that just projecting out states from a conformal field theory typically destroys its consistency, like e.g. modular invariance of the torus partition function. The projection therefore must be compensated by some additional manipulations, such as including new, twisted, sectors. For instance, in the case of the Gepner construction, the inner sector conformal field theory $`𝒞_{\mathrm{int}}`$ one starts with can be written as a tensor product $`𝒞_{k_1k_2\mathrm{}k_r}=𝒞_{k_1}\mathrm{}𝒞_{k_r}`$ of $`N=\mathrm{\hspace{0.17em}2}`$ minimal models. On this theory $`𝒞_{k_1k_2\mathrm{}k_r}`$ one imposes fermion alignment and the GSO projection, but at the same time includes additional states that do not appear in the spectrum of the original tensor product theory . Put differently, the torus partition function of the full Gepner conformal field theory $`𝒞^{\left(\mathrm{Gep}\right)}`$ – i.e. the theory that is obtained by these manipulations of projecting out old and of adding new states – contains non-diagonal (‘twisted’) contributions when viewed in terms of primary fields of the $`N=\mathrm{\hspace{0.17em}2}`$ tensor product. In particular the vacuum field of the original theory $`𝒞_{k_1k_2\mathrm{}k_r}`$ not only gets combined with itself, but also with other fields – to be called simple currents – of the original theory. (The corresponding states play in fact a crucial role in the space-time physics. The associated vertex operators provide e.g. the gravitini.) At the chiral level, this means that the chiral symmetry algebra of the Gepner model is extended beyond that of the $`N=\mathrm{\hspace{0.17em}2}`$ tensor product, namely precisely by including the relevant simple currents into the algebra . Note that from the point of view of the extended chiral algebra, the partition function is diagonal. <sup>1</sup><sup>1</sup>1 For simplicity, here we restrict ourselves to modular invariants of A-type for the $`N=\mathrm{\hspace{0.17em}2}`$ minimal models. Bulk spectra for other modular invariants have been computed in . For a recent discussion of boundary states in bulk theories with modular invariants of D-type or E-type (based on the results of ), see .
In short, the Gepner construction amounts to extending the underlying tensor product of $`𝒞_{\mathrm{int}}`$ with the theory $`𝒞_{\mathrm{s}\mathrm{t}}`$ that describes the surviving $`D`$ flat non-compact space-time dimensions by certain simple currents. For concreteness we refer to this new theory $`𝒞^{\left(\mathrm{Gep}\right)}`$ as the Gepner extension. To separate generic aspects which are related to the space-time part from aspects that depend on the chosen inner sector it turns out to be simpler, and conceptually clearer, to break up the extension into two separate steps. Thus we first perform a suitable extension on the inner sector $`𝒞_{\mathrm{int}}`$ alone. This way we arrive at a theory to which we refer as the CalabiYau extension $`𝒞^{\left(\mathrm{CY}\right)}`$. In models that possess a geometric interpretation as a sigma model on a CalabiYau manifold, it is this theory $`𝒞^{\left(\mathrm{CY}\right)}`$, rather than the original theory $`𝒞_{\mathrm{int}}`$ (e.g. the mere tensor product $`𝒞_{k_1k_2\mathrm{}k_r}`$ of minimal models) that should be compared with the geometrical data. The proper combination of $`𝒞^{\left(\mathrm{CY}\right)}`$ with the space-time theory $`𝒞_{\mathrm{s}\mathrm{t}}`$ then still requires further projections. Thus in a second step, we tensor $`𝒞^{\left(\mathrm{CY}\right)}`$ with $`𝒞_{\mathrm{s}\mathrm{t}}`$, and thereafter perform yet another extension that involves both $`𝒞^{\left(\mathrm{CY}\right)}`$ and $`𝒞_{\mathrm{s}\mathrm{t}}`$. As we will see, this latter extension is completely straightforward. This allows us to concentrate our attention on $`𝒞^{\left(\mathrm{CY}\right)}`$.
Simple current fields possess a variety of nice properties which allow for a very general and powerful treatment of arbitrary projections in which the chiral algebra gets enlarged . Such simple current extensions have often been compared to orbifold constructions. For our purposes it is, however, indispensable not to mix up the two operations of simple current extension and orbifolding. While the respective closed string partition functions indeed display a certain similarity – both correspond to projecting out some states and adding new ‘twisted’ states – there is a significant difference at the level of the chiral symmetry algebras and, as a consequence, at the level of chiral conformal field theory. Briefly, in a simple current extension the chiral algebra $`𝔄`$ gets larger – the new algebra $`𝔄_{\mathrm{ext}}`$ consists of the old one plus the simple current fields – while in the orbifold construction it gets smaller – the new algebra $`𝔄^G`$ is the fixed point subalgebra of the old one with respect to the orbifold group $`G`$. Accordingly, a simple current extension of a given theory has, generically, fewer primary fields (inequivalent representations of the chiral algebra) than the original theory; the ‘twisted states’ that appear in the partition function correspond to left-right asymmetric combinations of ordinary $`𝔄`$-representations. On the other hand, an orbifold has in general more primary fields than its mother theory, and the additional states correspond to new fields which appear already at the chiral level and carry ‘twisted representations’ of the original chiral algebra $`𝔄`$. The differences between the chiral aspects of the two constructions become particularly relevant when it comes to the study of boundary effects. Still, these two types of constructing a new conformal field theory from a given one are closely related – they are in fact each other’s inverse. The simple currents form an abelian group $`𝒢`$ under the fusion product, and it can be shown that the operations of extension by a group $`𝒢`$ of simple currents and of taking the abelian orbifold with respect to the character group $`G=𝒢^{}`$ are precisely inverse to each other. <sup>2</sup><sup>2</sup>2 It follows e.g. that the sizes of the stabilizer groups of the simple current and of the orbifold action are complementary, i.e. full simple current orbits correspond to orbifold fixed points and vice versa. For more quantitative statements, see section 7 of .
The reason for emphasizing the differences between the various extended theories that arise in a string compactification rests in the following observation. Once one works with the appropriate conformal field theory $`𝒞^{\left(\mathrm{CY}\right)}`$, the standard results for boundary conditions in (unitary) conformal field theories can be employed. In particular, Cardy’s construction of boundary states for boundary conditions that preserve the full chiral algebra can be applied directly.
The main points of this paper are the following. After establishing the necessary information about the Gepner and CalabiYau extensions (section 2), in section 3 we analyze in detail which symmetries of a Gepner model must be preserved and which ones can be broken by a given boundary condition. We also recall the recent increase of understanding of symmetry breaking boundary conditions (see , and also for applications to WZW models) and apply those results to Gepner models. We thereby obtain all boundary conditions that preserve full $`N=\mathrm{\hspace{0.17em}2}`$ world sheet and half of space-time supersymmetry, the so-called A-type boundary conditions. This includes in particular the boundary states recently obtained in . Our analysis reveals that within the boundary conditions of A-type, different ‘automorphism types’ appear, so that the A-type conditions can be naturally partitioned into several subsets. Boundary operators that change the automorphism type of the boundary conditions do not respect the GSO projection and therefore, generically, describe unstable brane configurations. Explicit formulas for all boundary states of Gepner models which preserve the full extended algebra are given.
In section 4 we turn to boundary conditions for which the action of the chiral algebra of the inner sector $`𝒞_{\mathrm{int}}`$ is twisted by some automorphism, in particular the B<sup>C</sup>-type conditions which are based on the mirror automorphism of the $`N=\mathrm{\hspace{0.17em}2}`$ superconformal algebra. As the relevant chiral algebra of the Gepner model is larger than that of $`𝒞_{\mathrm{int}}`$, it is necessary to lift the automorphism to the simple current extension $`𝒞^{\left(\mathrm{CY}\right)}`$. In the analysis of both A- and B-type conditions we encounter the problem that such a lift is typically not unique; we employ arguments from quantum Galois theory to describe this non-uniqueness (more details are provided in the appendix). Finally, in section 5 we comment on the relationship between various “singular” structures encountered in Gepner models and their geometric counterparts and mention some open problems.
## 2 The Gepner extension and the CalabiYau extension
### 2.1 The bosonic string map
The simple current machinery was mostly developed for unitary conformal field theories. But since for maintaining the world sheet supersymmetry we must align the superghosts as well, the conformal field theory of our interest is definitely not unitary. To deal with this problem, we make use of the bosonic string map . This stratagem allows us to map the non-unitary chiral conformal field theory of our primary interest to another chiral conformal field theory that is unitary and that possesses the same topological data, i.e. modular matrices, but also braiding and fusing matrices. In particular, both in the open and closed string sector we can then work with ordinary partition functions rather than with supersymmetric partition functions. It is worthwhile to point out that while the bosonic string map was originally designed to construct heterotic theories, we use it here to simplify the description of type II superstring theories which are supersymmetric both in the left and the right chiral part.
Concretely, the fermions of the flat $`D`$-dimensional space-time theory $`𝒞_{\mathrm{s}\mathrm{t}}`$ together with the superghosts can be described by the lorentzian lattice D<sub>D/2,1</sub>; the first $`D/2`$ components come from the bosonization of the space-time fermions, and the last one (with opposite sign in the kinetic energy) from the bosonization of the superghost system. The bosonic string map $`B`$ then amounts to replacing the non-unitary conformal field theory D<sub>D/2,1</sub> by the unitary conformal field theory D<sub>D/2+3</sub>. Both of these theories have four primary fields, corresponding to the four conjugacy classes of the D-type simple Lie algebras; the map exchanges the characters for the zero ($`o`$) and vector ($`v`$) conjugacy classes and multiplies the characters for the spinor ($`s`$) and conjugate spinor ($`c`$) conjugacy classes by $`1`$. Thus $`B`$ is encoded in the matrix
$$B=\left(\begin{array}{cccc}\hfill 0& \hfill 11& \hfill 0& \hfill 0\\ \hfill 11& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 11& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill 11\end{array}\right),$$
(2.1)
where $`11`$ is a unit matrix in the state space of the additional conformal field theory with which the theory for fermions and superghosts gets tensored, i.e. of the inner sector theory $`𝒞_{\mathrm{int}}`$ and the bosonic part of $`𝒞_{\mathrm{s}\mathrm{t}}`$. Denoting by a tilde quantities before the string map (‘supersymmetric quantities’), and without tilde the ones after the string map (‘ordinary CFT quantities’), we thus have, schematically, $`\stackrel{~}{\chi }=B\chi `$. For the modular transformation matrices this amounts to
$$\stackrel{~}{S}=BSB^1,\stackrel{~}{T}=BTB^1.$$
(2.2)
Given the modular invariant torus partition function $`Z`$ of the ordinary conformal field theory, which satisfies $`[Z,S]=\mathrm{\hspace{0.17em}0}=[Z,T]`$, it follows immediately that
$$\stackrel{~}{Z}:=BZB^1$$
(2.3)
is modular invariant on the supersymmetric side.
The $`N=\mathrm{\hspace{0.17em}2}`$ superconformal algebra contains a $`𝔲(1)`$ current subalgebra. Via spectral flow, space-time supersymmetry is achieved when all $`𝔲(1)`$ charges with respect to this subalgebra are odd integers. This condition can be fulfilled by a suitable projection on the allowed representations; this is precisely what the GSO projection does. Because of the exchange between $`o`$ and $`v`$ and the $`r`$-dependence of the $`𝔲(1)`$ charge of $`D_r`$-spinors, the bosonic string map (2.1) changes all charges with respect to the $`𝔲(1)`$ subalgebra of the $`N=\mathrm{\hspace{0.17em}2}`$ algebra by $`1mod2`$ . This means in particular that while in the supersymmetric theory the GSO projection is to odd integral $`𝔲(1)`$ charges, in the bosonic theory it is to even integral $`𝔲(1)`$ charges.
### 2.2 Simple current extensions
Starting from the tensor product theory $`𝒞_{\mathrm{s}\mathrm{t}}𝒞_{\mathrm{int}}`$, the Gepner construction proceeds by projecting out certain states and adding new ones . As already mentioned, technically this can be realized by the procedure of simple current extension. Basically, the simple current extension of a conformal field theory with chiral algebra $`𝔄`$ by some group $`𝒢`$ of simple currents of integral conformal weight has the following effects .
When fused with any other primary field $`\lambda `$ of the theory, a simple current $`\mathrm{J}`$ yields just a single field $`\mathrm{J}\lambda `$. Thus a simple current is invertible in the fusion ring. The group $`𝒢`$ then acts on the fusion ring of the $`𝔄`$-theory, and the simple current extension amounts to dividing out this action of $`𝒢`$.
The projection amounts to keep only those fields $`\lambda `$ which obey $`Q_\mathrm{J}(\lambda )=\mathrm{\hspace{0.17em}0}`$ for all $`\mathrm{J}𝒢`$, where
$$Q_\mathrm{J}(\lambda ):=\mathrm{\Delta }_\lambda +\mathrm{\Delta }_\mathrm{J}\mathrm{\Delta }_{\mathrm{J}\lambda }mod=\mathrm{\Delta }_\lambda \mathrm{\Delta }_{\mathrm{J}\lambda }mod$$
(2.4)
is the so-called monodromy charge of the field $`\lambda `$ of the $`𝔄`$-theory with respect to the simple current (with integral conformal weight) $`\mathrm{J}𝒢`$.
To obtain the primary fields of the extended theory we must organize the $`𝔄`$-fields that survive the projection into orbits $`[\lambda ]`$ under the fusion product with the currents in $`𝒢`$.
The diagonal modular invariant of the extended theory reads
$$Z_{\mathrm{ext}}=\underset{\genfrac{}{}{0pt}{}{[\lambda ]}{Q_\mathrm{J}(\lambda )=0\mathrm{J}𝒢}}{}|𝒮_\lambda ||\underset{\mathrm{J}𝒢/𝒮_\lambda }{}\chi _{\mathrm{J}\lambda }(\tau )|^2,$$
(2.5)
where $`𝒮_\lambda 𝒢`$ is the so-called stabilizer of $`\lambda `$, i.e. the subgroup $`𝒮_\lambda `$ consisting of those elements of $`𝒢`$ which leave $`\lambda `$ fixed under the fusion product of the $`𝔄`$-theory. Note that (2.5) is non-diagonal when viewed in terms of the primaries of the original theory. The terms of the form $`\chi _\lambda \chi _{\mathrm{J}\lambda }^{}`$ indicate the inclusion of twisted states which are needed to ensure modular invariance. (For an analysis in the WZW case, see .) Also, both the stabilizer subgroup $`𝒮_\lambda `$ and the monodromy charge are well defined for orbits $`[\lambda ]`$, not only for individual fields $`\lambda `$.
When an orbit $`[\lambda ]`$ has a non-trivial stabilizer, the factor of $`|𝒮_\lambda |`$ in the partition function (2.5) seems to indicate that the corresponding states occur several times. An additional ‘quantum number’ distinguishing those states is provided by a character of $`𝒮_\lambda `$, i.e. by $`\psi _\lambda 𝒮_\lambda ^{}`$. Accordingly, the primary fields of the extended theory are completely labeled as
$$\lambda _{\mathrm{ext}}=[\lambda ,\psi _\lambda ].$$
(2.6)
It is worth stressing that, while the prescription for projecting out states and adding new ones is already in itself sufficient for obtaining the spectrum of the model, to determine the complete modular properties of the model (and a fortiori for obtaining boundary conditions) it is indispensable to take proper care of such additional quantum numbers. Naively, in the case of a non-trivial stabilizer the projection rules appear to require the inclusion of the same state several times into the partition function. This would spoil unitarity of the modular S-matrix of the theory. The puzzle is resolved by realizing that those seemingly identical states are indeed distinguished by a further quantum number. Simple currents constitute a convenient conceptual framework for summarizing the required additional information.
The modular S-matrix $`S_{\mathrm{ext}}`$ of the extended theory can be expressed in terms of the modular S-matrix $`S`$ of the $`𝔄`$-theory and of similar matrices $`S^\mathrm{J}`$ with $`\mathrm{J}𝒢`$. The latter describe the modular S-transformation of one-point chiral blocks (of the $`𝔄`$-theory) on the torus with insertion $`\mathrm{J}`$ . Explicitly , the matrix elements of $`S_{\mathrm{ext}}`$ (labeled, according to the above, by $`𝒢`$-orbits of monodromy charge zero $`𝔄`$-primaries $`\mu `$, supplemented by a character $`\psi _\mu `$ of the stabilizer $`𝒮_\mu `$) read
$$(S_{\mathrm{ext}})_{[\lambda ,\psi _\lambda ],[\mu ,\psi _\mu ]}=\frac{\left|𝒢\right|}{\left|𝒮_\lambda \right|\left|𝒮_\mu \right|}\underset{\mathrm{J}𝒮_\lambda 𝒮_\mu }{}\psi _\lambda (\mathrm{J})\psi _\mu (\mathrm{J})^{}S_{\lambda ,\mu }^\mathrm{J},$$
(2.7)
where $`S^\mathrm{\Omega }S`$ is the ordinary S-matrix. When there are no fixed points (i.e., orbits with non-trivial stabilizer), then this expression collapses to
$$(S_{\mathrm{ext}})_{[\lambda ],[\mu ]}=|𝒢|S_{\lambda ,\mu },$$
(2.8)
so in this particular case the original S-matrix already contains all information about $`S_{\mathrm{ext}}`$.
Actually the formula (2.7) does not cover the most general situation. In full generality we rather have to account for the fact that the implementation of symmetries in quantum systems is typically only projective. This can also happen for the symmetries studied here. Quantitatively, the effect is described by a two-cocycle on the stabilizer group $`𝒮_\mu `$. What is remarkable is that this two-cocycle can be computed entirely in terms of the matrices $`S^\mathrm{J}`$. One can show that the projectivity is properly taken into account by replacing $`𝒮_\mu `$ by the subgroup
$$𝒰_\mu 𝒮_\mu $$
(2.9)
of $`𝒮_\mu `$ on which the two-cocycle vanishes; $`𝒰_\mu `$ is called the untwisted stabilizer of $`\mu `$. <sup>3</sup><sup>3</sup>3 The formula for the extended S-matrix then reads
$$(S_{\mathrm{ext}})_{[\lambda ,\widehat{\psi }_\lambda ],[\mu ,\widehat{\psi }_\mu ]}=|𝒢|[|𝒮_\lambda ||𝒰_\lambda ||𝒮_\mu ||𝒰_\mu |]^{1/2}\underset{\mathrm{J}𝒰_\lambda 𝒰_\mu }{}\widehat{\psi }_\lambda (\mathrm{J})\widehat{\psi }_\mu (\mathrm{J})^{}S_{\lambda ,\mu }^\mathrm{J},$$
where $`\widehat{\psi }_\mu `$ is a character of $`𝒰_\mu 𝒮_\mu 𝒢`$. Note that in this general situation even the labeling of primary fields is different from the case where $`𝒰_\mu `$ coincides with $`𝒮_\mu `$ for all $`\mu `$; in place of the label (2.6) we now have $`\lambda _{\mathrm{ext}}=[\lambda ,\widehat{\psi }_\lambda ]`$. Our notation for the simple current orbits is actually adapted to the general situation, as $`𝒢_\mu /𝒰_\mu `$ acts non-trivially on the characters $`\widehat{\psi }_\mu `$ when $`𝒰_\mu `$ is a proper subgroup of $`𝒮_\mu `$.
Also, while the arguments in were not sufficient to prove the formula rigorously, a proof is possible by combining them with the results of on the uniqueness of the modularisation of a premodular category. Independently, various aspects of the formula can be tested directly . For instance, with the help of the computer program kac (see http://norma.nikhef.nl/$`\stackrel{~}{}`$t58/kac.html) it was checked in a huge number of cases that it produces non-negative integers when inserted into the Verlinde formula. Moreover, manifestly the formula requires only information about the chiral conformal field theory, and even only information about topological aspects of the chiral theory. In particular it does not involve any knowledge about boundary conditions.
While in Gepner models with diagonal or charge conjugation invariant one always has $`𝒰_\mu =𝒮_\mu `$, for models where the extension of the $`N=\mathrm{\hspace{0.17em}2}`$ tensor product is by a larger group – e.g. corresponding to taking non-diagonal modular invariants of the $`N=\mathrm{\hspace{0.17em}2}`$ minimal models – cases where $`𝒰_\mu `$ is a proper subgroup of $`𝒮_\mu `$ can and do arise. This must e.g. be taken into account when analyzing the boundary states introduced in . In Gepner models with diagonal (or charge conjugation) torus partition function – the situation of our main interest below – one always has $`𝒰_\mu =𝒮_\mu `$, so that henceforth we will ignore this modification.
Via the Verlinde formula, the fusion rules of the extended theory can then be expressed through the fusion rules of the $`𝔄`$-theory and the fixed point quantities $`𝒮_\lambda `$ and $`S^\mathrm{J}`$. For instance, for $`𝒢=_2=\{\mathrm{\Omega },\mathrm{J}\}`$ one finds
$$\begin{array}{cc}(\mathrm{N}_{\mathrm{ext}})_{[\lambda ,\psi ],[\lambda ^{},\psi ^{}]}^{[\lambda ^{\prime \prime },\psi ^{\prime \prime }]}\hfill & =\frac{1}{\left|𝒮_\lambda \right|\left|𝒮_\lambda ^{}\right|\left|𝒮_{\lambda ^{\prime \prime }}\right|}\text{[}\mathrm{N}_{\lambda ,\lambda ^{}}^{\lambda ^{\prime \prime }}+\mathrm{N}_{\lambda ,\mathrm{J}\lambda ^{}}^{\lambda ^{\prime \prime }}+2\underset{\genfrac{}{}{0pt}{}{\mu }{\mathrm{J}\mu =\mu }}{}\text{(}\psi ^{}\psi ^{\prime \prime }\frac{S_{\lambda ,\mu }S_{\lambda ^{},\mu }^\mathrm{J}S_{\lambda ^{\prime \prime },\mu }^\mathrm{J}}{S_{\mathrm{\Omega },\mu }}\hfill \\ & \\ & \text{ }+\psi \psi ^{\prime \prime }\frac{S_{\lambda ,\mu }^\mathrm{J}S_{\lambda ^{},\mu }S_{\lambda ^{\prime \prime },\mu }^\mathrm{J}}{S_{\mathrm{\Omega },\mu }}+\psi \psi ^{}\frac{S_{\lambda ,\mu }^\mathrm{J}S_{\lambda ^{},\mu }^\mathrm{J}S_{\lambda ^{\prime \prime },\mu }^{}}{S_{\mathrm{\Omega },\mu }}\text{)}\text{]}.\hfill \end{array}$$
(2.10)
As is clear from (2.7), the rows and columns of the $`S^\mathrm{J}`$-matrices are labeled by only those primaries of the $`𝔄`$-theory that are fixed under $`\mathrm{J}`$. Thus unless at least two of the fields $`\lambda `$, $`\lambda ^{}`$ and $`\lambda ^{\prime \prime }`$ are fixed under $`\mathrm{J}`$, the corresponding terms in (2.10) vanish.
It is worth emphasizing that a simple current extension amounts to nothing else than to a change of the underlying chiral conformal field theory. This fact, which is somewhat hidden in other treatments of projections (compare e.g. ), is of central importance for gaining a better understanding of boundary states in Gepner models. Namely, it implies in particular that all the features that are revealed in the analysis of these projections can be understood in a manner that is completely independent from our (yet uncomplete) understanding of world sheet boundary effects. Since the relevant results have undergone extensive physical and mathematical consistency checks which do not use any information about boundary conditions, we can safely exploit these structures as an input in the construction of boundary states for Gepner models.
A crucial property of superconformal field theories is that the supersymmetry automatically leads to the existence of certain simple currents. First of all, independently on the number of world sheet supersymmetries, every superconformal field theory has a distinguished simple current $`v`$: the generator of world sheet supersymmetry, which has order two and conformal dimension $`\mathrm{\Delta }=\mathrm{\hspace{0.17em}3}/2`$. The monodromy charge with respect to $`v`$ is 0 for primary fields in the NeveuSchwarz sector and 1/2 for primaries in the Ramond sector. The ‘superpartner’ of a primary field $`\lambda `$ is given by fusion product of $`\lambda `$ with this simple current $`v`$.
In case the superconformal field theory has extended ($`N=\mathrm{\hspace{0.17em}2}`$) supersymmetry, there is yet another simple current $`s_{\mathrm{int}}`$: the Ramond ground state $`R_0`$ with highest $`𝔲(1)`$ charge. This can be seen after expressing the $`𝔲(1)`$ current $`J`$ of the $`N=\mathrm{\hspace{0.17em}2}`$ algebra in terms of a canonical free boson, $`J(z)=\mathrm{i}\sqrt{c/3}\varphi (z)`$. Then the Ramond ground state is given by $`\mathrm{exp}(\mathrm{i}\sqrt{c/12}\varphi )`$. Its conformal dimension is $`\mathrm{\Delta }=c/24`$, as befits a Ramond ground state, and it has the correct $`𝔲(1)`$ charge $`c/6`$. In this formulation, it is easy to see that the monodromy charge with respect to this simple current equals half of the superconformal $`𝔲(1)`$ charge of a field.
### 2.3 The Gepner extension
In this subsection we display the simple current extension that leads from an internal $`N=\mathrm{\hspace{0.17em}2}`$ superconformal theory $`𝒞_{\mathrm{int}}`$ to a consistent string background $`𝒞^{\left(\mathrm{Gep}\right)}`$. Let us point out that this extension is not only applicable for the original Gepner models, but likewise for any other $`N=\mathrm{\hspace{0.17em}2}`$ compactification in which the internal theory is a rational conformal field theory, for instance for KazamaSuzuki models.
The (flat) space-time bosons and Virasoro ghosts will play no role in what follows, and accordingly we suppress their contribution to $`𝒞_{\mathrm{s}\mathrm{t}}`$. What then remains of the space-time theory are the fermions and superghosts. After the bosonic string map, these are described by the level one WZW theory based on the Lie algebra D<sub>D/2+3</sub>; this (unitary) conformal field theory $`\mathrm{D}_{D/2+3}`$ has four primary fields, which we label as
$$\varpi \{o,s,v,c\}.$$
(2.11)
The Gepner extension is the extension of the tensor product theory $`\mathrm{D}_{D/2+3}𝒞_{\mathrm{int}}`$ by a certain simple current group $`𝒦^{\left(\mathrm{Gep}\right)}`$, which accounts for fermion alignment and GSO projection. The group $`𝒦^{\left(\mathrm{Gep}\right)}`$ is generated by some number $`r`$ of order-two currents
$$\mathrm{v}_i:=(v;v_{i,\mathrm{int}})$$
(2.12)
together with
$$\mathrm{s}_{\mathrm{tot}}:=(s;s_{\mathrm{int}});$$
(2.13)
the fields $`\mathrm{v}_i`$ will be referred to as alignment currents (rather than as vector currents, as is done e.g. in ), and $`\mathrm{s}_{\mathrm{tot}}`$ as the total spinor current. When the inner sector $`𝒞_{\mathrm{int}}`$ can itself be written as a tensor product (e.g. of $`N=\mathrm{\hspace{0.17em}2}`$ minimal models as in the original Gepner construction), then each tensor factor provides us with one of the currents $`v_{i,\mathrm{int}}`$, which then is a non-trivial field in the $`i`$th factor, tensored with the identity field of all other factors of $`𝒞_{\mathrm{int}}`$. The subspace with lowest conformal weight of the representation space of the bosonic part of the $`N=\mathrm{\hspace{0.17em}2}`$ algebra corresponding to the field $`v_{i,\mathrm{int}}`$ has a dimension which is a multiple of two; it contains the two supercurrents $`G_i^\pm `$ of the $`i`$th tensor factor.
The order $`N_s`$ of $`\mathrm{s}_{\mathrm{tot}}`$ is model dependent. Also, depending on the model the resulting group $`𝒦^{\left(\mathrm{Gep}\right)}`$ is either the direct product of the $`_{N_s}`$ generated by the total spinor current $`\mathrm{s}_{\mathrm{tot}}`$ and of the $`_2`$ groups generated by the alignment currents $`\mathrm{v}_i`$, or else some quotient of that direct product group. The latter happens when the field $`(\mathrm{s}_{\mathrm{tot}})^{N_s/2}`$ is contained in $`(_2)^r`$, <sup>4</sup><sup>4</sup>4 For instance, in $`N=\mathrm{\hspace{0.17em}2}`$ minimal models with odd level $`k`$, the spinor current $`s`$ (see formula (3.25) below) satisfies $`s^{2k+4}=v`$. This results in the equality $`\mathrm{s}_{\mathrm{tot}}^{N_s/2}=_{i=1}^r\mathrm{v}_i`$ when $`𝒞_{\mathrm{int}}`$ is the tensor product of $`r`$ minimal models with only odd levels. For more details, see subsection 3.5. and in that case the corresponding quotient is, as an abstract abelian group, the direct product of $`_{N_s}`$ with $`r1`$ copies of $`_2`$. Thus the simple current group of the Gepner extension has the structure
$$𝒦^{\left(\mathrm{Gep}\right)}=_{N_s}\times (_2)^{r\eta }\mathrm{with}\eta \{0,1\}.$$
(2.14)
The extension by the group generated by the currents $`\mathrm{v}_i`$ guarantees world sheet supersymmetry. Indeed, the total $`N=\mathrm{\hspace{0.17em}2}`$ superconformal algebra must split into two modules of its bosonic subalgebra, the vacuum module and a module containing the supercurrents $`G^\pm =_{i=1}^rG_i^\pm `$; this is so only after extension by the $`\mathrm{v}_i`$. (Concretely, e.g. the two terms in $`G_1^\pm +G_2^\pm G_1^\pm \mathbf{\hspace{0.17em}1}+\mathbf{\hspace{0.17em}1}G_2^\pm `$ lie in two distinct irreducible modules of $`𝔄_{\mathrm{int}}`$, and these modules get combined into an irreducible module of the extended algebra precisely due to the extension by $`\mathrm{v}_1\mathrm{v}_2=(o;v_{1,\mathrm{int}}v_{2,\mathrm{int}})`$.) The extension by the total spinor current $`\mathrm{s}_{\mathrm{tot}}`$ implements the GSO projection and hence ensures space-time supersymmetry. The monodromy charge (2.4) with respect to the current $`\mathrm{s}_{\mathrm{tot}}`$ can be shown to coincide (modulo $``$) with half of the superconformal $`𝔲(1)`$ charge of a state. Also, a change from the $`o`$ to the $`v`$ conjugacy class results in a change of this monodromy charge by $`1/2mod`$, and in the Ramond sector the same effect results from the $`r`$-dependence of the $`𝔲(1)`$ charge of the spinors of $`D_r`$. Recalling the form (2.1) of the bosonic string map, we thus see again that it changes the effect of the GSO projection from a projection to odd integral $`𝔲(1)`$ charges to a projection to even integral $`𝔲(1)`$ charges.
Let us remark that the abelian orbifold construction that brings us back from the Gepner model $`𝒞^{\left(\mathrm{Gep}\right)}`$ to the tensor product $`\mathrm{D}_{D/2+3}𝒞_{\mathrm{int}}`$ consists of orbifolding by the group $`K^{\left(\mathrm{Gep}\right)}`$ that is generated by the automorphisms
$$\begin{array}{c}\begin{array}{c}\mathrm{v}_i\mathrm{v}_i,\mathrm{v}_j\mathrm{v}_j\text{ for }ji\hfill \\ \mathrm{s}_{\mathrm{tot}}\mathrm{s}_{\mathrm{tot}},\hfill \end{array}\}\text{ for }i=\mathrm{\hspace{0.17em}1},2,\mathrm{},r,\mathrm{and}\hfill \\ \\ \text{ }\mathrm{v}_j\mathrm{v}_j\text{ for all }j,\mathrm{s}_{\mathrm{tot}}\mathrm{exp}(2\pi \mathrm{i}/N_s)\mathrm{s}_{\mathrm{tot}}.\hfill \end{array}$$
(2.15)
Also note that the latter map corresponds to a shift $`\varphi \varphi +4\pi \sqrt{3/c}/N_s`$ of the free boson $`\varphi `$ in terms of which the spinor current can be written as $`\mathrm{s}_{\mathrm{tot}}=\mathrm{exp}(\mathrm{i}\sqrt{c/12}\varphi )`$.
### 2.4 The CalabiYau extension
The original Gepner construction – reformulated in the previous subsection in terms of simple current extensions – involves both the flat space-time part $`\mathrm{D}_{D/2+3}`$ and the inner sector $`𝒞_{\mathrm{int}}`$. This tends to obscure the connection with the geometric formulation in terms of sigma models on CalabiYau manifolds. The object in the Gepner model that corresponds to the compactification manifold in the geometric setting is not simply the conformal field theory $`𝒞_{\mathrm{int}}`$ – e.g. the tensor product $`𝒞_{k_1k_2\mathrm{}k_r}`$ of minimal models – but rather an extension of this tensor product, which will be specified shortly; we shall denote it by $`𝒞^{\left(\mathrm{CY}\right)}`$ and call it the CalabiYau extension. The connection to geometric compactifications is usually derived using a LandauGinzburg description of the minimal models (for a different line of arguments, see ), and indeed this construction involves a non-trivial projection on the tensor product, commonly referred to as forming a LandauGinzburg orbifold . In a second step, one combines this extended inner sector $`𝒞^{\left(\mathrm{CY}\right)}`$ with flat space-time (i.e. tensors with the theory $`\mathrm{D}_{D/2+3}`$) and then performs an additional extension, which has more similarities to the GSO projection in ten flat dimensions. Unlike the step from $`𝒞_{\mathrm{int}}`$ to $`𝒞^{\left(\mathrm{CY}\right)}`$, this further extension is completely straightforward. Let us stress that the procedure that we call the CalabiYau extension can be performed independently of any connection of the internal sector to a classical geometrical compactification.
Inspecting the Gepner extension, one observes that the group $`𝒦^{\left(\mathrm{Gep}\right)}`$ contains many currents that have a trivial space-time part and therefore effectively define an extension of the internal theory $`𝒞_{\mathrm{int}}`$ alone. We denote the group of these currents by $`𝒦^{\left(\mathrm{CY}\right)}`$. Using the fusion rules of the $`\mathrm{D}_{D/2+3}`$ theory (which are of the form $`_2\times _2`$ when we compactify to $`D=d+2=\mathrm{\hspace{0.17em}6}`$ dimensions, and $`_4`$ for compactifications to $`D=\mathrm{\hspace{0.17em}4}`$ or $`8`$), we find that $`𝒦^{\left(\mathrm{CY}\right)}`$ is generated by all products of any two of the currents $`v_{i,\mathrm{int}}`$, and hence by
$$\mathrm{w}_i:=v_{1,\mathrm{int}}v_{i,\mathrm{int}},i\{2,3,\mathrm{},r\}$$
(2.16)
which will again be called alignment currents, together with the current <sup>5</sup><sup>5</sup>5 The presence of the power of $`v_{1,\mathrm{int}}`$ accounts for the fact that $`s^2=v^{d/2}`$ in the $`\mathrm{D}_{D/2+3}`$ theory. More explicitly, for compactifications to $`D=\mathrm{\hspace{0.17em}6}`$, the group $`𝒦^{\left(\mathrm{CY}\right)}`$ contains all products of an even number of currents $`v_{i,\mathrm{int}}`$ and all products of $`(s_{\mathrm{int}})^2`$ with the former, while for compactifications to $`D=\mathrm{\hspace{0.17em}4}`$ or 8 in addition to all products of an even number of currents $`v_{i,\mathrm{int}}`$ one has all products of $`(s_{\mathrm{int}})^2`$ with an odd number of $`v_{i,\mathrm{int}}`$.
$$\mathrm{u}:=s_{\mathrm{int}}^2(v_{1,\mathrm{int}})^{d/2}.$$
(2.17)
This group has the structure
$$𝒦^{\left(\mathrm{CY}\right)}=_{N_s/2}\times (_2)^{r1\eta },$$
(2.18)
where the contribution $`\eta \{0,1\}`$ in the exponent accounts for the possibility that $`\mathrm{u}^{N_s/4}`$ is contained in the product of the $`r1`$ $`_2`$ groups that are generated by the currents $`\mathrm{w}_i`$ (compare the remarks before formula (2.14)).
The theory $`𝒞^{\left(\mathrm{CY}\right)}`$ that is obtained upon extension of $`𝒞_{\mathrm{int}}`$ by $`𝒦^{\left(\mathrm{CY}\right)}`$ inherits a simple current (sub)group $`\{o^{\left(\mathrm{CY}\right)},s^{\left(\mathrm{CY}\right)},v^{\left(\mathrm{CY}\right)},c^{\left(\mathrm{CY}\right)}\}`$ that has the same fusion rules (i.e., $`_2\times _2`$ and $`_4`$, respectively) as the $`\mathrm{D}_{D/2+3}`$ theory. The final projections that bring us from $`\mathrm{D}_{D/2+3}𝒞^{\left(\mathrm{CY}\right)}`$ to the Gepner extension $`𝒞^{\left(\mathrm{Gep}\right)}`$ amount to extending by the simple current $`(v;v^{\left(\mathrm{CY}\right)})`$, which aligns fermions in $`\mathrm{D}_{D/2+3}`$ and $`𝒞^{\left(\mathrm{CY}\right)}`$, and in addition by either $`(s;s^{\left(\mathrm{CY}\right)})`$ or $`(c;s^{\left(\mathrm{CY}\right)})`$. This very last extension is the true analog of the GSO projection in flat space-time. In particular, when one combines the left and right halves of the theory, the choice between type IIA and IIB theories is equivalent to a choice between the extension by $`(s;s^{\left(\mathrm{CY}\right)})`$ both on the left and the right (or equivalently, by $`(c;s^{\left(\mathrm{CY}\right)})`$ on both sides), or else by $`(s;s^{\left(\mathrm{CY}\right)})`$ on one side and by $`(c;s^{\left(\mathrm{CY}\right)})`$ on the other. We remark that all currents in the extension of $`\mathrm{D}_{D/2+3}𝒞^{\left(\mathrm{CY}\right)}`$ to $`𝒞^{\left(\mathrm{Gep}\right)}`$ act freely. Thus the simple formula (2.8) for the modular S-matrix of the extension applies, and hence as announced this extension is straightforward.
The dependence of the precise structure of the $`𝒦^{\left(\mathrm{CY}\right)}`$ extension on the number of compactified dimensions can be traced back to the fact that the internal spectral flow operator, mapping the R to the NS sector, changes the $`𝔲(1)`$ charge by $`c_{\mathrm{int}}/6`$, where $`c_{\mathrm{int}}=\mathrm{\hspace{0.17em}12}3d/2`$ is the central charge of $`𝒞_{\mathrm{int}}`$. Thus, while in the NS sector we always project onto integral internal $`𝔲(1)`$ charge, the internal $`𝔲(1)`$ charges in the R sector are either integers or in $`+1/2`$, depending on the external dimension being $`6`$ or $`4,8`$, respectively. The integrality of the $`𝔲(1)`$ charges in the NS sector is necessary to make a correspondence between chiral primary fields and differential forms in the geometric compactification possible, and is therefore highly welcome. In fact, $`s_{\mathrm{int}}^2`$ can be identified model independently as the square of the Ramond ground state with maximal $`𝔲(1)`$ charge.
As we shall see below, there is yet another intermediate theory, to be denoted by $`𝒞_{\mathrm{wsusy}}`$, that is of interest. This is the theory that one obtains from the inner sector $`𝒞_{\mathrm{int}}`$ by extending with the subgroup $`(_2)^{r1}𝒦^{\left(\mathrm{CY}\right)}`$ generated by the alignment currents $`\mathrm{w}_i`$ only, i.e. by enforcing only fermion alignment, and hence world sheet supersymmetry, on the internal theory. $`𝒞_{\mathrm{wsusy}}`$ will play an important role in the analysis of boundary conditions. We summarize the relation between the various extensions schematically as
$$\mathrm{D}_{D/2+3}𝒞_{\mathrm{int}}\mathrm{D}_{D/2+3}𝒞_{\mathrm{wsusy}}\mathrm{D}_{D/2+3}𝒞^{\left(\mathrm{CY}\right)}𝒞^{\left(\mathrm{Gep}\right)}.$$
(2.19)
Note that the group that furnishes the extension from $`𝒞_{\mathrm{wsusy}}`$ to $`𝒞^{\left(\mathrm{CY}\right)}`$ is the cyclic group generated by the image $`\mathrm{U}`$ of the simple current $`\mathrm{u}`$ (3.24) in the $`𝒞_{\mathrm{wsusy}}`$-theory. The order of $`\mathrm{U}`$ can differ from the order $`N_s/2`$ of $`\mathrm{u}`$ by a factor of two; it is given by $`N_s^{}/2`$ with
$$N_s^{}=2^\eta N_s,$$
(2.20)
where $`\eta `$ is the integer introduced in formula (2.14).
We can – and will – simplify the discussion and restrict our attention in the sequel to the intermediate theory $`𝒞^{\left(\mathrm{CY}\right)}`$ (and later on also to $`𝒞_{\mathrm{wsusy}}`$). As the additional extension to $`𝒞^{\left(\mathrm{Gep}\right)}`$ is so simple, we do not loose any essential features when doing so. In particular the issue of fixed points arises always only in the study of $`𝒞^{\left(\mathrm{CY}\right)}`$. To conclude this section, let us emphasize that the construction described above is model independent and does not rely on specific aspects of the $`N=\mathrm{\hspace{0.17em}2}`$ superconformal field theory used in the inner sector.
## 3 A-type boundary conditions
### 3.1 Intermediate chiral algebras
Already in closed string theory the simple current extension leading to $`𝒞^{\left(\mathrm{CY}\right)}`$ must be taken into account properly. In particular, a careful treatment of fixed points is compulsory to find the correct massless spectrum, compare e.g. . When open strings are present, for the following reason an even deeper understanding of fixed point resolution is required. In the computation of the massless spectrum one just counts states, thus a detailed understanding of the underlying partition functions suffices. In contrast, as was first realized by Cardy , in the construction of boundary states the modular S-matrix enters directly; therefore a complete knowledge of this matrix in simple current extensions is required as well. Moreover, open string partition functions (annulus amplitudes) also implicitly contain the modular matrices, so that not even the open string spectrum can be obtained correctly without proper resolution of the fixed points. In this context, it is an important observation that typically the CalabiYau extension, and hence also the Gepner extension, do possess fixed points. For instance, as will be discussed in subsection 3.5, in the case of a tensor product $`𝒞_{\mathrm{int}}=𝒞_{k_1k_2\mathrm{}k_r}`$ of $`N=\mathrm{\hspace{0.17em}2}`$ minimal models, fixed points arise precisely if at least one level $`k_i`$ is even.
The group of automorphisms of the $`N=\mathrm{\hspace{0.17em}2}`$ superconformal algebra is the Lie group O$`(2)`$. This group has two connected components, and any element of the component not connected to the identity can be obtained by composing an element of the identity component with the mirror automorphism (see formula (4.2)). Those automorphisms of the chiral algebra of a tensor product of internal models that respect the $`N=\mathrm{\hspace{0.17em}2}`$ structure are generically given by O$`(2)`$ as well (but additional permutation symmetries are present when some of the factors of the tensor product are identical). Accordingly, in such string compactifications one conventionally distinguishes between two classes of boundary states: Those which correspond to an automorphism in the identity component of O$`(2)`$, and those corresponding to an automorphism in the other component. In the literature, the former states are often collectively referred to as A-type boundary states, while latter are said to be of B-type. A-type boundary conditions leave the chiral algebra $`𝔄_{\mathrm{int}}`$ of the inner sector $`𝒞_{\mathrm{int}}`$ invariant; insisting that also an $`N=\mathrm{\hspace{0.17em}2}`$ subalgebra of the extension $`𝒞^{\left(\mathrm{CY}\right)}`$ remains unbroken, the A-type automorphism must be the identity map. In this section we study in detail boundary conditions $`|a_A^{\left(\mathrm{CY}\right)}`$ which do preserve $`𝔄_{\mathrm{int}}`$, i.e. which satisfy
$$\text{(}Y_n\mathbf{\hspace{0.17em}1}+(1)^{\mathrm{\Delta }_Y1}\mathbf{\hspace{0.17em}1}Y_n\text{)}|a_A^{\left(\mathrm{CY}\right)}=0$$
(3.1)
for every field $`Y(z)=_nY_nz^{n\mathrm{\Delta }_Y}`$ of conformal weight $`\mathrm{\Delta }_Y`$ in $`𝔄_{\mathrm{int}}`$.
We start our analysis by recalling that the total chiral algebra $`𝔄^{\left(\mathrm{CY}\right)}`$ of the CalabiYau extension $`𝒞^{\left(\mathrm{CY}\right)}`$ is obtained from $`𝔄_{\mathrm{int}}`$ by a simple current extension with the group $`𝒦^{\left(\mathrm{CY}\right)}`$. A generic boundary state $`|a_A^{\left(\mathrm{CY}\right)}`$ will not preserve all of $`𝔄^{\left(\mathrm{CY}\right)}`$, but only some subalgebra $`𝔄_a`$ containing $`𝔄_{\mathrm{int}}`$. This subalgebra cannot be arbitrary, though. First of all, we are interested in conformally invariant boundary conditions only, and hence the Virasoro subalgebra of $`𝔄^{\left(\mathrm{CY}\right)}`$ must be preserved. This is automatically satisfied by the boundary states $`|a_A^{\left(\mathrm{CY}\right)}`$, since the Virasoro algebra is already contained in the inner sector algebra $`𝔄_{\mathrm{int}}`$. But in addition the preserved subalgebra must have enough structure to allow for the construction of conformal blocks and, based on them, of correlation functions. One therefore has to require that $`𝔄_a`$ must again be a vertex operator algebra. To get an overview over all boundary conditions of our present interest, we thus look for all vertex operator algebras $`𝔄`$ that lie between $`𝔄_{\mathrm{int}}`$ and $`𝔄^{\left(\mathrm{CY}\right)}`$:
$$𝔄_{\mathrm{int}}𝔄𝔄^{\left(\mathrm{CY}\right)}.$$
(3.2)
For general vertex operator algebras, this classification of subalgebras would be a hopeless problem. Here, however, we know that the inner sector chiral algebra $`𝔄_{\mathrm{int}}`$ can be characterized as the subalgebra of $`𝔄^{\left(\mathrm{CY}\right)}`$ that is left pointwise fixed by the group $`K^{\left(\mathrm{CY}\right)}`$ of automorphisms of $`𝔄^{\left(\mathrm{CY}\right)}`$, described in formula (2.15), which is dual to $`𝒦^{\left(\mathrm{CY}\right)}`$. This observation allows us to employ a basic result from the Galois theory for vertex operator algebras , which tells us that the possible chiral algebras between $`𝔄_{\mathrm{int}}`$ and $`𝔄^{\left(\mathrm{CY}\right)}`$ are in one-to-one correspondence with subgroups of the group $`K^{\left(\mathrm{CY}\right)}`$. We conclude that to every boundary state $`|a_A^{\left(\mathrm{CY}\right)}`$ of $`𝒞^{\left(\mathrm{CY}\right)}`$ we can associate a subgroup $`K_a^{\left(\mathrm{CY}\right)}`$ of $`K^{\left(\mathrm{CY}\right)}`$, such that the subalgebra of $`𝔄^{\left(\mathrm{CY}\right)}`$ that is preserved by the boundary condition is the fixed point algebra with respect to $`K_a^{\left(\mathrm{CY}\right)}`$.
### 3.2 Boundary states and automorphism types
Boundary conditions that preserve a fixed point algebra of the bulk symmetries with respect to some finite abelian group of automorphisms have been studied in . When applied to the present situation, the pertinent results of may be summarized as follows.
Using notation from the unextended theory $`𝒞_{\mathrm{int}}`$, the boundary states can be labeled in a way much similar to the labeling (2.6) of primary fields of the extended theory $`𝒞^{\left(\mathrm{CY}\right)}`$, namely as
$$a=[\mu ,\psi _\mu ].$$
(3.3)
The difference is that, unlike in (2.6), here $`\mu `$ can be any primary field label of $`𝒞_{\mathrm{int}}`$, i.e. now there is no restriction on the monodromy charge. <sup>6</sup><sup>6</sup>6 At this point, the possibility of having untwisted stabilizers $`𝒰_\mu 𝒮_\mu `$ (see formula (2.9)) must in general be taken into account. Then $`\psi _\mu `$ gets in fact replaced by a character $`\widehat{\psi }_\mu `$ of the untwisted stabilizer $`𝒰_\mu `$, and the simple current orbit is obtained by an action that also changes $`\widehat{\psi }_\mu `$ in a non-trivial way. Also, the prefactor of $`\stackrel{~}{S}`$ gets changed analogously as in the formula in footnote 3.
When $`K_a^{\left(\mathrm{CY}\right)}`$ is non-trivial, then $`|a_A^{\left(\mathrm{CY}\right)}`$ can no longer be written as a linear combination of Ishibashi states of $`𝒞^{\left(\mathrm{CY}\right)}`$. This is simply due to the fact that the preserved chiral symmetry is then not big enough to guarantee that all $`𝔄^{\left(\mathrm{CY}\right)}`$-descendants are reflected at the boundary in the same way, so that different descendants must be treated differently. The boundary state can, however, still be expressed in terms of suitable generalizations $`|\mathrm{}`$ of the Ishibashi states of the unextended theory $`𝒞_{\mathrm{int}}`$. These states are labeled by a pair $`(\lambda ,\psi _\lambda )`$, where now again the restriction of zero monodromy charge is to be imposed on the primary field $`\lambda `$ (but no simple current orbit is taken any longer). <sup>7</sup><sup>7</sup>7 Here $`\psi _\lambda `$ is always a character of the full stabilizer, even when the untwisted stabilizer is strictly smaller than the stabilizer .
Heuristically the situation can be understood as follows. On the side of Ishibashi states, only states with vanishing monodromy charge are present; in the orbifold language of , only states in the untwisted sector appear. This means that the orbifold element in ‘space’ direction on the torus is always trivial, whereas we still project, i.e. we still deal with non-trivial group elements in ‘time’ direction. According to Cardy’s ideas, the boundary states are obtained by a modular S-transformation from the Ishibashi states. After that transformation we only have the trivial group element in ‘time’ direction, which implies that there is no projection, so that orbits appear. However, in ‘space’ direction we now have non-trivial elements, and therefore the twisted sectors show up in the description of the boundary states.
To make these heuristic ideas quantitative, we introduce a matrix $`\stackrel{~}{S}`$ that takes over the role that the usual S-matrix plays in the Cardy case. As shown in , such a matrix indeed exists. The following structures were uncovered.
The expansion of the boundary state $`|a_A^{\left(\mathrm{CY}\right)}`$ with respect to the generalized Ishibashi states reads
$$|[\mu ,\psi _\mu ]_A^{\left(\mathrm{CY}\right)}=\underset{\lambda ,\psi _\lambda }{}\frac{\stackrel{~}{S}_{(\lambda ,\psi _\lambda ),[\mu ,\psi _\mu ]}}{\sqrt{\stackrel{~}{S}_{(\mathrm{\Omega }),[\mu ,\psi _\mu ]}}}|(\lambda ,\psi _\lambda ).$$
(3.4)
The matrix $`\stackrel{~}{S}`$ appearing here can be expressed as
$$\stackrel{~}{S}_{(\lambda ,\psi _\lambda ),[\mu ,\psi _\mu ]}=\frac{\left|𝒢\right|}{\left|𝒮_\lambda \right|\left|𝒮_\mu \right|}\underset{\mathrm{J}𝒮_\lambda 𝒮_\mu }{}\psi _\lambda (\mathrm{J})\psi _\mu (\mathrm{J})^{}S_{\lambda ,\mu }^\mathrm{J},$$
(3.5)
which is similar to (2.7) (but remember that now also twisted sectors are allowed in the second label). Complete information on the boundary state, like brane tensions or RR charges, is encoded in this matrix $`\stackrel{~}{S}`$.
Note the similarity between the result (3.4) and Cardy’s formula for symmetry preserving boundary conditions, in which the modular S-matrix appears in place of $`\stackrel{~}{S}`$. It turns out that the matrix $`\stackrel{~}{S}`$ has still more in common with the modular S-matrix. Indeed, as first realized in , a subset of the sewing constraints for correlation functions of a rational conformal field theory can be isolated which leads to a simple non-linear equation for the bulk-boundary coefficients for excitation of the vacuum field on the boundary. As pointed out in , this equation means that the reflection coefficients constitute one-dimensional representations of a certain finite-dimensional associative algebra, which generalizes the fusion rule algebra and is called the classifying algebra. <sup>8</sup><sup>8</sup>8 Extending Cardy’s work, it was shown in that the structure constants of this algebra are traces on suitable spaces of conformal blocks.
These results allow us to introduce the notion of an elementary boundary condition; this furnishes by definition an irreducible representation of the classifying algebra. Thus the elementary boundary conditions are in one-to-one correspondence with the one-dimensional irreducible representations of the classifying algebra. The matrix $`\stackrel{~}{S}`$ as given by (3.5) diagonalizes the structure constants of the classifying algebra, analogously as the modular S-matrix diagonalizes the fusion rules. In string theory, on the other hand, one must in addition introduce (ChanPaton) multiplicities for boundary conditions. Thus the space of all boundary conditions to be considered in string theory forms a cone over the elementary boundary conditions, and generically one is dealing with higher-dimensional, and hence necessarily reducible, representations of the classifying algebra. It is also quite common that some of the solutions that are present as elementary boundary conditions in the conformal field theory possess a zero ChanPaton multiplicity, i.e. do not appear at the string theory level at all.
Based on the properties of the matrix $`\stackrel{~}{S}`$, the space of boundary conditions for $`𝔄^{\left(\mathrm{CY}\right)}`$ that preserve $`𝔄_{\mathrm{int}}`$ can be analyzed in detail . One important result is that each of the boundary states studied here possesses a definite automorphism type. This means that it can be written as a linear combination of twisted Ishibashi states, where the twist is a fixed automorphism $`\omega `$ of $`𝔄^{\left(\mathrm{CY}\right)}`$. Such a twisted Ishibashi state $`|[\lambda ,\psi _\lambda ]_\omega `$ obeys the twisted Ward identity
$$\text{(}𝒴_n\mathbf{\hspace{0.17em}1}+(1)^{\mathrm{\Delta }_𝒴1}\mathbf{\hspace{0.17em}1}\omega (𝒴_n)\text{)}|[\lambda ,\psi _\lambda ]_\omega =0$$
(3.6)
for every field $`𝒴`$ of conformal weight $`\mathrm{\Delta }_𝒴`$ in $`𝔄^{\left(\mathrm{CY}\right)}`$. (This generalizes the usual definition of Ishibashi states of Dirichlet-type for the free boson. In the terminology of , the relation (3.6) says that $`\omega `$ provides the glueing condition of the boundary condition.)
The subset of A<sub>ω</sub>-type boundary states – that is, of all boundary states which preserve $`𝔄_{\mathrm{int}}`$ and which have some prescribed automorphism type $`\omega `$ – corresponds to a subalgebra of the classifying algebra. The structure constants of this subalgebra can be expressed in terms of traces of the action of $`\omega `$ on the space of the relevant three-point conformal blocks. This again generalizes the Cardy case, in which the structure constants are just the fusion rules, which in turn are nothing else than traces of the identity automorphism on spaces of three-point blocks.
The fact that the boundary states that preserve $`𝔄_{\mathrm{int}}`$ come in several different automorphism types can also be understood as follows. The identity map on $`𝔄_{\mathrm{int}}`$ can be lifted to an automorphism of $`𝔄^{\left(\mathrm{CY}\right)}`$ in several distinct ways. The group of these lifts of automorphisms is just the quantum Galois group for the extension. Each element of this group gives us an automorphism type. Let us stress that the fact that all the boundary conditions that are commonly referred to as A-type do possess an automorphism type is non-trivial indeed. It reflects the equally non-trivial statement of quantum Galois theory that all intermediate algebras are obtained as fixed algebras.
Once the boundary states are given explicitly, i.e. once the matrix $`\stackrel{~}{S}`$ is known, all annulus amplitudes can be computed by sandwiching a string propagator $`q^{L_0+\stackrel{~}{L}_0c/12}`$ between the two appropriate boundary states. But already from the general expression (3.5) above (and from general properties of the matrices $`S^\mathrm{J}`$), it can be established via representation theoretic arguments that in full generality the annulus coefficients are non-negative integers, as befits the coefficients of the open string partition function. The completeness and associativity properties of the annulus coefficients can be shown to be satisfied as well. Further, one can write the annulus amplitude for two boundary conditions of automorphism types $`\omega _1`$ and $`\omega _2`$ as a sum of characters $`\chi _{[\nu ,\psi _\nu ]^{}}^{}`$ of the extension of $`𝒞_{\mathrm{int}}`$ by the subgroup
$$𝒦_{\omega _1\omega _2}:=\{\mathrm{J}𝒦^{\left(\mathrm{CY}\right)}|Q_\mathrm{J}(\omega _1)=0=Q_\mathrm{J}(\omega _2)\}$$
(3.7)
of $`𝒦^{\left(\mathrm{CY}\right)}`$. (Here the isomorphism between $`K^{\left(\mathrm{CY}\right)}`$ and the dual group $`(𝒦^{\left(\mathrm{CY}\right)})^{}`$ is used, i.e. the automorphisms are regarded as characters of $`𝒦^{\left(\mathrm{CY}\right)}`$.) The coefficients then read <sup>9</sup><sup>9</sup>9 Here again we suppress the changes that arise when genuine untwisted stabilizer groups are present; see formula (6.23) of .
$$\mathrm{A}_{[\mu _1,\psi _1],[\mu _2,\psi _2]}^{[\nu ,\psi _\nu ]^{}}=\frac{\left|𝒦_{\omega _1\omega _2}\right|}{\left|𝒦^{\left(\mathrm{CY}\right)}\right|}\underset{[\lambda ,\psi _\lambda ^{}]^{}}{}\frac{\left|𝒮_\lambda \right|}{\left|𝒮_\lambda ^{}\right|}\underset{\psi _\lambda \psi _\lambda ^{}}{}\stackrel{~}{S}_{(\lambda ,\psi _\lambda ),[\mu _1,\psi _1]}^{}\stackrel{~}{S}_{(\lambda ,\psi _\lambda ),[\mu _2,\psi _2]}S_{[\lambda ,\psi _\lambda ^{}]^{},[\nu ,\psi _\nu ]^{}}^{}/S_{[\lambda ,\psi _\lambda ^{}]^{},[\mathrm{\Omega }]^{}}^{},$$
(3.8)
where $`S^{}`$ is the modular S-matrix of the $`𝒦_{\omega _1\omega _2}`$-extension of $`𝒞_{\mathrm{int}}`$, and where the second summation is over all $`𝒮_\lambda `$-characters $`\psi _\lambda `$ that restrict to the given character $`\psi _\lambda ^{}`$ of the subgroup $`𝒮_\lambda ^{}=𝒮_\lambda 𝒦_{\omega _1\omega _2}`$.
Another important conclusion to be drawn is that the monodromy charge constitutes a grading of the annulus coefficients, in the following sense. It follows from the result (3.8) that
$$\mathrm{A}_{[\mu _1,\psi _1],[\mu _2,\psi _2]}^{[\nu ,\psi _\nu ]^{}}=\mathrm{e}^{2\pi \mathrm{i}[Q_\mathrm{J}(\mu _2)Q_\mathrm{J}(\mu _1)+Q_\mathrm{J}(\nu )]}\mathrm{A}_{[\mu _1,\psi _1],[\mu _2,\psi _2]}^{[\nu ,\psi _\nu ]^{}},$$
(3.9)
so that the annulus coefficient $`\mathrm{A}_{[\mu _1,\psi _1],[\mu _2,\psi _2]}^{[\nu ,\psi _\nu ]^{}}`$ vanishes unless $`Q_\mathrm{J}(\nu )=Q_\mathrm{J}(\mu _1)Q_\mathrm{J}(\mu _2)`$ for all $`\mathrm{J}𝒦^{\left(\mathrm{CY}\right)}`$. Thus all open string states that appear in the annulus amplitude for two boundary conditions of automorphism types $`\omega _1`$ and $`\omega _2`$ have a common monodromy charge $`Q_\mathrm{J}`$ with respect to any current $`\mathrm{J}`$ in $`𝒦^{\left(\mathrm{CY}\right)}`$, and this common value is given by
$$\mathrm{exp}(2\pi \mathrm{i}Q_\mathrm{J})=\omega _1^1\omega _2(\mathrm{J}).$$
(3.10)
### 3.3 World sheet supersymmetry and space-time supersymmetry
We now analyze what boundary conditions in string compactifications of Gepner type preserve the (super)symmetries that have to be imposed to obtain a consistent bulk theory. To this end we study in more detail the twisted Ward identities (3.6). Again we start our discussion with boundary states that preserve the chiral symmetry algebra $`𝔄_{\mathrm{int}}`$ of the inner sector $`𝒞_{\mathrm{int}}`$. As already mentioned, such boundary states are often collectively said to be of A-type, see for instance . However, as we have detailed above, the inner sector $`𝒞_{\mathrm{int}}`$ of the string compactification needs to be extended to $`𝒞^{\left(\mathrm{CY}\right)}`$, which does not have $`𝔄_{\mathrm{int}}`$ as chiral symmetry, but rather its simple current extension $`𝔄^{\left(\mathrm{CY}\right)}`$. To identify boundary states that are relevant for string theory, in particular those that can be given a geometric interpretation as D-branes, it is necessary to understand what part of the extended algebra $`𝔄^{\left(\mathrm{CY}\right)}`$ is preserved or broken by a given boundary state.
Recall that the simple current extension from $`𝒞_{\mathrm{int}}`$ to $`𝒞^{\left(\mathrm{CY}\right)}`$ consists of a part that ensures world sheet supersymmetry and another part necessary for space-time supersymmetry. Accordingly, the automorphism type $`\omega _a`$ of a boundary condition that preserves $`𝔄_{\mathrm{int}}`$ carries information concerning both the extension by the alignment currents $`\mathrm{w}_i=v_{1,\mathrm{int}}v_{i,\mathrm{int}}`$ $`(i=\mathrm{\hspace{0.17em}2},3,\mathrm{},r`$) and the extension by $`\mathrm{u}=s_{\mathrm{int}}^2v_{1,\mathrm{int}}^{d/2}`$. These individual pieces of information can be thought of as measuring which supersymmetries are broken or conserved by the boundary state. More concretely, we can attribute to every boundary state $`|a_A^{\left(\mathrm{CY}\right)}`$ an element $`\omega _a`$ of the orbifold group $`K^{\left(\mathrm{CY}\right)}`$ that is dual to $`𝒦^{\left(\mathrm{CY}\right)}`$. As an abstract abelian group, this is again $`_{N_s/2}\times _2^{r1\eta }`$. The automorphism $`\omega _a`$ occurs in the twisted Ward identities (3.6) satisfied by $`|a_A^{\left(\mathrm{CY}\right)}`$; explicitly, we have
$$\begin{array}{cc}\omega _a(\mathrm{w}_i)=\zeta _{a,i}\mathrm{w}_i\hfill & \mathrm{with}\zeta _{a,i}\{\pm 1\},\hfill \\ & \\ \omega _a(\mathrm{u})=\mathrm{e}^{2\mathrm{i}\vartheta _a}\mathrm{u}\hfill & \mathrm{with}2\vartheta _a\mathrm{\hspace{0.17em}2}\pi /(N_s/2),\hfill \end{array}$$
(3.11)
and hence
$$\begin{array}{c}\text{(}(\mathrm{w}_i)_n\mathbf{\hspace{0.17em}1}+\zeta _{a,i}(1)^{\mathrm{\Delta }_{\mathrm{w}_i}1}\mathbf{\hspace{0.17em}1}(\mathrm{w}_i)_n\text{)}|a_A^{\left(\mathrm{CY}\right)}=0,\hfill \\ \\ \text{(}\mathrm{u}_n\mathbf{\hspace{0.17em}1}+\mathrm{e}^{2\mathrm{i}\vartheta _a}(1)^{\mathrm{\Delta }_\mathrm{u}1}\mathbf{\hspace{0.17em}1}\mathrm{u}_n\text{)}|a_A^{\left(\mathrm{CY}\right)}=0.\hfill \end{array}$$
(3.12)
Also recall that for each tensor factor of the internal $`N=\mathrm{\hspace{0.17em}2}`$ theory the field $`v_{i,\mathrm{int}}`$ contains the supersymmetry charges $`G_i^\pm `$. Now $`N=\mathrm{\hspace{0.17em}1}`$ world sheet supersymmetry plays the role of a gauge symmetry of perturbative superstring theory. In order not to destroy this constitutive feature of superstrings, $`N=\mathrm{\hspace{0.17em}1}`$ world sheet supersymmetry must not be broken by any boundary state that is present (i.e., has a non-zero ChanPaton multiplicity) at the string theory level. Concretely, we have the relation $`\{Q_{\mathrm{BRST}},\beta \}=G`$ between the $`N=\mathrm{\hspace{0.17em}1}`$ supercurrent $`G`$, the BRST charge $`Q_{\mathrm{BRST}}`$ and the superghost $`\beta `$. When combined with the identities
$$(Q_{\mathrm{BRST}}\mathbf{\hspace{0.17em}1}\mathrm{𝟏}Q_{\mathrm{BRST}})|a_A^{\left(\mathrm{CY}\right)}=0\mathrm{and}(\beta _r\mathbf{\hspace{0.17em}1}+\mathrm{i}ϵ\mathbf{\hspace{0.17em}1}\beta _r)|a_A^{\left(\mathrm{CY}\right)}=0$$
(3.13)
($`ϵ\{\pm 1\}`$) which encode BRST invariance of the boundary state and the boundary condition for the superghost (which is model independent and independent of the chosen boundary condition for the theory $`𝒞^{\left(\mathrm{CY}\right)}`$), this relation implies that we must also have
$$(G_r\mathbf{\hspace{0.17em}1}+\mathrm{i}ϵ\mathbf{\hspace{0.17em}1}G_r)|a_A^{\left(\mathrm{CY}\right)}=0.$$
(3.14)
Comparison of this Ward identity with the result (3.12) then tells us that we must require the invariance property $`\omega _a(G)=G`$ of the $`N=\mathrm{\hspace{0.17em}1}`$ supercurrent.
The extended theory $`𝒞^{\left(\mathrm{CY}\right)}`$ actually possesses $`N=\mathrm{\hspace{0.17em}2}`$ world sheet supersymmetry, with supercurrents $`G^\pm =_iG_i^\pm `$, into which the $`N=\mathrm{\hspace{0.17em}1}`$ supersymmetry can be embedded in several different ways, namely as $`G=(\mathrm{e}^{\mathrm{i}\gamma }G^++\mathrm{e}^{\mathrm{i}\gamma }G^{})/\sqrt{2}`$ for any $`\gamma `$. By inspecting the operator products of the supercurrents, it follows that there exists some $`N=\mathrm{\hspace{0.17em}1}`$ subalgebra that is preserved if and only if in the relations (3.11) we have $`\zeta _{a,i}=\mathrm{\hspace{0.17em}1}`$ for all $`i`$, i.e. if and only if the part of the automorphism type that concerns the alignment currents is trivial. (This appears to have been ignored in part of the literature.) As a matter of fact, in that case the boundary condition preserves all $`N=\mathrm{\hspace{0.17em}1}`$ subalgebras, and hence even the whole $`N=\mathrm{\hspace{0.17em}2}`$ algebra.
Thus from now on we only admit those automorphism types which obey
$$\omega _a(\mathrm{w}_i)=\mathrm{w}_i.$$
(3.15)
(Put differently, in the string theory any conformal field theory boundary condition which violates this relation is assigned ChanPaton multiplicity zero.) According to the Ward identities (3.12), the automorphism type of any of the remaining boundary conditions is then completely characterized by a single number. This number is essentially given by $`\vartheta \mathrm{\hspace{0.17em}2}\pi /N_s`$; but we also have to take into account that the order $`N_s^{}/2`$ of the image $`\mathrm{U}`$ of $`\mathrm{u}`$ in the theory $`𝒞_{\mathrm{wsusy}}`$ can be different from the order $`N_s/2`$ of $`\mathrm{u}`$ itself. Thus the boundary conditions must rather be labeled by
$$\theta =2\pi n/N_s^{}\mathrm{with}n\{0,1,\mathrm{},N_s^{}1\}.$$
(3.16)
We shall denote the corresponding automorphism type of branes by A<sub>θ</sub>. Note that here we use $`\theta `$ instead of $`2\theta `$ as label, even though from the point of view of the CalabiYau extension, the automorphism types $`\theta `$ and $`\theta +\pi `$ cannot be distinguished. We do so because in full string theory (i.e., for $`𝒞^{\left(\mathrm{Gep}\right)}`$ rather than $`𝒞^{\left(\mathrm{CY}\right)}`$), the two automorphism types differ on the final GSO projection.
This restriction to boundary conditions that preserve not only $`𝔄_{\mathrm{int}}`$, but even its extension by the alignment currents i.e. the algebra $`𝔄_{\mathrm{wsusy}}`$, also ensures that the annulus partition function between any two boundary conditions has an open string spectrum with world sheet supersymmetry.
In contrast to world sheet supersymmetry, space-time supersymmetry is not an indispensable ingredient of a superstring theory. Also, it is not determined by chiral considerations alone, in the sense that the space-time supersymmetry generators $`Q`$ are given by closed string operators in which left- and right-movers are combined. Indeed, they arise as zero modes
$$Q=\mathrm{d}^2zs(z,\overline{z})$$
(3.17)
of fields in the full conformal field theory obtained by putting together both chiral halves. In an $`N=\mathrm{\hspace{0.17em}2}`$ string compactification, the field $`s(z,\overline{z})`$ consists of a spectral flow operator on one chiral half combined with the vacuum field of the other chiral half (for more explicit expressions, see e.g. ). Thus there are both left- and right-moving supersymmetry charges, $`Q_L`$ and $`Q_R`$. In the context of open strings, a preserved space-time supersymmetry is a certain linear combination of a left- and right-moving charge that annihilates the boundary state :
$$(Q_L+Q_R^{})|a_A^{\left(\mathrm{CY}\right)}=0.$$
(3.18)
If $`s`$ is any left-moving spectral flow operator, $`\omega _a(s)`$ will again have the properties of a spectral flow, and therefore the sum of the corresponding left and right-moving zero modes will still constitute a conserved supersymmetry. That different boundary states in Gepner models conserve different space-time supersymmetry charges has been observed in . What is new about the analysis above is the observation that this corresponds to dealing with boundary conditions of distinct automorphism types. Indeed, the fact that within the classification into A- and B-type there exist subclasses of boundary states with different automorphism type has so far not been appreciated in the literature.
Thus all boundary conditions in Gepner models that have been considered so far in the literature, and for which we have identified an automorphism type labeled by an element $`\theta \mathrm{\hspace{0.17em}2}\pi /N_s^{}`$, in fact do preserve half of space-time supersymmetry. In short, in all situations studied in this paper, the presence of space-time supersymmetry in the closed string sector (ensured by the GSO projection) together with the preservation of $`N=\mathrm{\hspace{0.17em}1}`$ world sheet supersymmetry by a boundary state already guarantees that the boundary state is BPS. On the other hand, the open string spectrum, encoded in the annulus partition function, will be space-time supersymmetric only if the automorphism types of the boundary conditions on the two boundary components of the annulus are equal. (Without reference to automorphism types, a condition of this type for space-time supersymmetry of the partition function was also derived in .) In that case all annuli can be expressed in terms of characters of the CalabiYau extension $`𝒞^{\left(\mathrm{CY}\right)}`$, and as a consequence the GSO projection – which is a chiral issue – guarantees in particular the absence of tachyons in the spectrum of open string states. Geometrically, the difference $`\theta _1\theta _2`$ between the labels of the two automorphism types can be given the intuitive interpretation of an angle between two branes. The open string spectrum is space-time supersymmetric if the angle between the two branes is zero. In contrast, when the angle is non-zero, then absence of tachyons in the open string spectrum is not guaranteed any longer.
Thus in order to guarantee the absence of open string tachyons, among the A<sub>θ</sub>-type boundary conditions it is generically necessary to restrict to those with a single fixed value of $`\theta `$. For instance, when we decide to keep a boundary state with $`\theta =\mathrm{\hspace{0.17em}0}`$, then we typically have to dismiss all boundary states with $`\theta \mathrm{\hspace{0.17em}0}`$. However, in the case of type I theories, an orientifold projection may stabilize the brane (for reviews see ); in such circumstances, A<sub>θ</sub> boundary conditions with several distinct values of $`\theta `$ could coexist. Whether this happens or not, and for what choices of $`\theta `$, are model dependent questions. (The answer can in particular depend on bulk moduli.) Furthermore, when we also include B-type conditions, it can happen that requiring absence of open string tachyons restricts the allowed $`\theta `$ value of A<sub>θ</sub> conditions even when only a single $`\theta `$ is kept.
### 3.4 A<sub>0</sub>-type boundary conditions
The special class of boundary conditions of type $`\text{A}\text{0}\mathrm{A}_{\theta =0}`$ are just those of ‘trivial automorphism type’ $`\omega =\mathrm{id}`$, i.e. those for which the identity map of $`𝔄^{\left(\mathrm{CY}\right)}`$ is used as the extension of the identity of $`𝔄_{\mathrm{int}}`$. These boundary conditions preserve the whole algebra $`𝔄^{\left(\mathrm{CY}\right)}`$. Put differently, they are precisely the boundary conditions that were studied long ago by Cardy for an arbitrary rational conformal field theory. In the case of tensor products of minimal models, various aspects of A<sub>0</sub>-type conditions have already been studied in . As far as the space-time aspects of string theory are concerned, the A<sub>0</sub>-type boundary conditions are not distinguished in any specific manner among the larger set of A-type boundary conditions, which as explained above preserve both world sheet and space-time supersymmetry, too. However, from a pure world sheet point of view, A<sub>0</sub>-type boundary conditions are special in that they preserve the full chiral algebra $`𝔄^{\left(\mathrm{CY}\right)}`$ of the CalabiYau extension. This means that they are directly accessibly by Cardy’s method, and we therefore still treat them separately here.
The only data that enter Cardy’s construction of boundary states are the entries of the modular S-matrix of the conformal field theory, i.e. in our case the matrix (2.7) for the CalabiYau extension $`𝒞^{\left(\mathrm{CY}\right)}`$. (In terms of the classifying algebra mentioned above, this result is a manifestation of the fact that the invariant subalgebra that corresponds to the A<sub>0</sub>-type boundary conditions is nothing but the fusion rule algebra of $`𝒞^{\left(\mathrm{CY}\right)}`$.) These data are well under control (some explicit formulas will be presented below). Besides the boundary states, also other aspects of these special boundary conditions are well understood (compare e.g. for a general discussion of correlation functions). Therefore the case of trivial automorphism type – which was also used as a starting point in the constructions in – is absolutely under control.
In particular, the A<sub>0</sub>-type boundary states are in natural one-to-one correspondence with the primary fields of the CalabiYau extension $`𝒞^{\left(\mathrm{CY}\right)}`$, and they can be expanded in the Ishibashi states of $`𝒞^{\left(\mathrm{CY}\right)}`$:
$$|[\mu ,\psi _\mu ]_{A_0}^{\left(\mathrm{CY}\right)}=\underset{[\lambda ,\psi _\lambda ]}{}\frac{(S^{\left(\mathrm{CY}\right)})_{[\lambda ,\psi _\lambda ],[\mu ,\psi _\mu ]}}{\sqrt{(S^{\left(\mathrm{CY}\right)})_{[\mathrm{\Omega }],[\mu ,\psi _\mu ]}}}|[\lambda ,\psi _\lambda ].$$
(3.19)
Similarly, in the formula (3.8) for the annulus coefficients we now have $`𝒦_{\omega _1\omega _2}=𝒦^{\left(\mathrm{CY}\right)}`$ as well as $`\stackrel{~}{S}=S^{}=S^{\left(\mathrm{CY}\right)}`$, so that it reduces to
$$\mathrm{A}_{[\mu _1,\psi _1],[\mu _2,\psi _2]}^{[\nu ,\psi _\nu ]}=\underset{[\lambda ,\psi _\lambda ]}{}(S^{\left(\mathrm{CY}\right)})_{[\lambda ,\psi _\lambda ],[\mu _1,\psi _1]}^{}(S^{\left(\mathrm{CY}\right)})_{[\lambda ,\psi _\lambda ],[\mu _2,\psi _2]}(S^{\left(\mathrm{CY}\right)})_{[\lambda ,\psi _\lambda ],[\nu ,\psi _\nu ]}/(S^{\left(\mathrm{CY}\right)})_{[\lambda ,\psi _\lambda ],[\mathrm{\Omega }]}.$$
(3.20)
By comparison with the Verlinde formula for $`𝒞^{\left(\mathrm{CY}\right)}`$ we then learn that the annulus coefficients indeed coincide with the structure constants of the fusion algebra of $`𝒞^{\left(\mathrm{CY}\right)}`$.
### 3.5 Tensor products of minimal models
Let us now specialize to the Gepner models proper, where the inner sector $`𝒞_{\mathrm{int}}`$ is a tensor product $`𝒞_{k_1k_2\mathrm{}k_r}=𝒞_{k_1}\mathrm{}𝒞_{k_r}`$ of $`N=\mathrm{\hspace{0.17em}2}`$ minimal models $`𝒞_{k_i}`$ at levels $`k_i_{>0}`$. We also restrict our attention to the A<sub>0</sub>-type boundary conditions. In this special case it is particularly easy to make the formula for $`\stackrel{~}{S}`$, and hence the description of boundary states, fully explicit. We first recall the following information about such Gepner models that will be needed in the sequel.
It ist convenient to think of an $`N=\mathrm{\hspace{0.17em}2}`$ minimal model at level $`k`$ as a coset construction $`𝔰𝔲(2)_k\times 𝔲(1)_4/𝔲(1)_{2h}`$, where $`h:=k+2`$. Accordingly, we denote its primary fields by $`\mathrm{\Phi }_q^{l,s}`$. Then the labels $`l,s,q`$ are in the ranges $`l\{0,1,\mathrm{},k\}`$, $`s\{0,1,2,3\}`$, and $`q\{0,1,\mathrm{},2h1\}`$, subject to the parity selection rule $`l+sq\mathrm{\hspace{0.17em}2}`$, and to the ‘field identification’ $`\mathrm{\Phi }_q^{l,s}\mathrm{\Phi }_{q+h}^{kl,s+2}`$. This labeling of primary fields refers to the bosonic subalgebra of the $`N=\mathrm{\hspace{0.17em}2}`$ superconformal algebra. For example, the two world sheet supercurrents are not regarded as descendants of the vacuum but, rather, both correspond to the primary field $`\mathrm{\Phi }_0^{0,2}\mathrm{\Phi }_{\pm h}^{k,0}`$. The identity primary field is $`\mathrm{\Omega }=\mathrm{\Phi }_0^{0,0}`$. (For more details about minimal models see, for example, .)
The primary fields of $`𝒞_{k_1k_2\mathrm{}k_r}`$ are then labeled as
$$(\mathrm{\Phi }_{q_1}^{l_1,s_1},\mathrm{\Phi }_{q_2}^{l_2,s_2},\mathrm{},\mathrm{\Phi }_{q_r}^{l_r,s_r})$$
(3.21)
where $`\mathrm{\Phi }_{q_i}^{l_i,s_i}`$ is a primary field of the $`i`$th minimal model. For brevity, below we will also use the notation
$$(\lambda ,\sigma ,\xi )(l_1,s_1,q_1,\mathrm{},l_r,s_r,q_r)$$
(3.22)
for these collections of labels.
The simple current group $`𝒦^{\left(\mathrm{CY}\right)}`$ of the CalabiYau extension is of the form (2.18), with $`r`$ the number of minimal model factors. It is generated by the $`r1`$ order-two currents
$$\mathrm{w}_i:=(v,\mathrm{\Omega },\mathrm{},\mathrm{\Omega },v,\mathrm{\Omega },\mathrm{},\mathrm{\Omega }),i\{2,3,\mathrm{},r\},$$
(3.23)
where the $`v`$-entries are in the first (say) and $`i`$th minimal model, together with the combination
$$\mathrm{u}:=(s,s,\mathrm{},s)^2(v,\mathrm{\Omega },\mathrm{},\mathrm{\Omega })^{d/2}=(s^2v^{d/2},s^2,\mathrm{},s^2).$$
(3.24)
Here $`v`$ stands for the minimal model primary field $`\mathrm{\Phi }_0^{0,2}`$ that contains the two world sheet supercurrents, and $`s=\mathrm{\Phi }_1^{0,1}`$ is the simple current in the Ramond sector whose action provides the spectral flow. These fields
$$s=\mathrm{\Phi }_1^{0,1}\mathrm{and}v=\mathrm{\Phi }_0^{0,2}$$
(3.25)
have conformal weight $`c/24=k/8(k+2)`$ and $`3/2`$, respectively.
The order of $`\mathrm{u}`$ is
$$\mathrm{ord}(\mathrm{u})=N_s/2=\mathrm{scm}_{i=1,2\mathrm{}r}\{\eta _ih_i\}$$
(3.26)
with $`h_ik_i+2`$ and $`\eta _i=\mathrm{\hspace{0.17em}1}`$ for $`k_i`$ even, $`\eta _i=\mathrm{\hspace{0.17em}2}`$ for $`k_i`$ odd.
Exploiting our knowledge about the minimal model fusion rules, it is straightforward combinatorics to establish the following group and fixed point structure of $`𝒦^{\left(\mathrm{CY}\right)}`$.
When all levels $`k_i`$ are odd, then we have
$$\mathrm{u}^{N_s/4}=\underset{i=1}{\overset{r}{}}\mathrm{w}_i,$$
(3.27)
and hence $`𝒦^{\left(\mathrm{CY}\right)}=_{N_s/2}\times _2^{r2}`$. In this case there are no fixed points.
In contrast, when at least one level is even, then we have $`𝒦^{\left(\mathrm{CY}\right)}=_{N_s/2}\times _2^{r1}`$, and fixed points do occur. In fact there is then a unique simple current $`\mathrm{L}𝒦^{\left(\mathrm{CY}\right)}`$ having fixed points. $`\mathrm{L}`$ is given by
$$\mathrm{L}=\mathrm{u}^{N_s/4}\underset{i=1}{\overset{r}{}}(\mathrm{w}_i)^{ϵ_i},$$
(3.28)
where the value of $`ϵ_i\{0,1\}`$ depends on the power of 2 contained in $`N_s/2`$. Namely, according to (3.26), $`N_s/2`$ is always even; when it is also divisible by 4, then
$$ϵ_i=\{\begin{array}{cc}1\hfill & \text{when the power of 2 in }h_i\text{ is maximal},\hfill \\ 0\hfill & \mathrm{else}.\hfill \end{array}$$
(3.29)
When $`N_s/2`$ is not divisible by 4, then we have instead
$$ϵ_i=\{\begin{array}{cc}1\hfill & \text{when }h_i\text{ is odd},\hfill \\ 0\hfill & \mathrm{else}.\hfill \end{array}$$
(3.30)
It also follows that $`\mathrm{L}`$ has order 2, and that with the help of field identification it can be rewritten as
$$\mathrm{L}=(\mathrm{\Phi }_0^{0,0},\mathrm{},\mathrm{\Phi }_0^{0,0},\mathrm{\Phi }_0^{k_{r^{}+1},0},\mathrm{},\mathrm{\Phi }_0^{k_r,0}).$$
(3.31)
Here without loss of generality we assume that the $`h_i`$ have been ordered in such a way that those containing the maximal power of 2 are the last $`rr^{}`$ ones, i.e. are labeled by $`i\{r^{}+1,\mathrm{},r\}`$. Fixed points of $`\mathrm{L}`$ are all fields (3.21) which obey $`l_i=k_i/2`$ for every $`i=r^{}+1,\mathrm{},r`$, while all other labels are arbitrary (except, of course, for the parity selection rule $`l_i+s_iq_i\mathrm{\hspace{0.17em}2}`$ and the restriction to zero monodromy charge with respect to $`𝒦^{\left(\mathrm{CY}\right)}`$).
Employing the general formula (2.7), we are thus in a position to display the modular S-matrix of the CalabiYau extension. For the tensor product theory $`𝒞_{k_1k_2\mathrm{}k_r}`$ we have the tensor product of the S-matrices of the individual factors:
$$S_{(\lambda ,\sigma ,\xi ),(\lambda ^{},\sigma ^{},\xi ^{})}=2^r\underset{i=1}{\overset{r}{}}S_{l_i,l_i^{}}^{𝔰𝔲(2)_{k_i}}S_{s_i,s_i^{}}^{𝔲(1)_4}\text{(}S_{q_i,q_i^{}}^{𝔲(1)_{2h_i}}\text{)}^{}$$
(3.32)
(the $`r`$ factors of $`2`$ stem from field identification in each of the $`r`$ minimal models). The fixed point matrix $`S^\mathrm{L}`$ reads
$$S_{(\lambda ,\sigma ,\xi )^\mathrm{L},(\lambda ^{},\sigma ^{},\xi ^{})^\mathrm{L}}^\mathrm{L}=2^r\underset{i=1}{\overset{r}{}}S_{s_i,s_i^{}}^{𝔲(1)_4}\left(S_{q_i,q_i^{}}^{𝔲(1)_{2h_i}}\right)^{}\underset{i=1}{\overset{r^{}}{}}S_{l_i,l_i^{}}^{𝔰𝔲(2)_{k_i}}\underset{i=r^{}+1}{\overset{r}{}}S_{k_i/2,k_i/2}^{\mathrm{J}𝔰𝔲(2)_{k_i}}.$$
(3.33)
(This is only defined when both fields involved are fixed points, which is indicated by attaching the superscript of the labels $`(\mathrm{})^\mathrm{L}`$.) Here $`S^{\mathrm{J}𝔰𝔲(2)_k}`$ is the fixed point ‘matrix’ for the simple current $`l=k`$ of the $`𝔰𝔲(2)`$ WZW model at level $`k`$. As $`k`$ acts by fusion as $`kl=kl`$ and hence has just a single fixed point $`k/2`$, the quantity $`S^{\mathrm{J}𝔰𝔲(2)_k}`$ is in fact a single number, and in this simple situation the results of tell us immediately that $`S_{k/2,k/2}^{\mathrm{J}𝔰𝔲(2)_k}=\mathrm{e}^{2\pi \mathrm{i}3k/16}`$. Moreover, $`_{i=r^{}+1}^rk_i/16`$ must be a multiple of $`1/12`$ , so that we have $`_{i=r^{}+1}^rS_{k_i/2,k_i/2}^{\mathrm{J}𝔰𝔲(2)_{k_i}}=\mathrm{e}^{\pi \mathrm{i}p/2}`$ with $`p\{0,1,2,3\}`$ (inspecting the list of Gepner models, one sees that all values of $`p`$ occur). For further simplification, notice that
$$\underset{i=r^{}+1}{\overset{r}{}}(S^{\mathrm{J}𝔰𝔲(2)_{k_i}})^2=(1)^{_{i=r^{}+1}^r3k_i/4}=(1)^{_{i=r^{}+1}^rk_i/4}$$
(3.34)
and that $`_{i=r^{}+1}^rk_i/4`$ is always an integer, as follows from the condition on the central charge.
Putting this information together, we see that the modular S-matrix of $`𝒞^{\left(\mathrm{CY}\right)}`$ reads
$$\begin{array}{c}S_{[(\lambda ,\sigma ,\xi ),\psi ],[(\lambda ^{},\sigma ^{},\xi ^{}),\psi ^{}]}^{\left(\mathrm{CY}\right)}=2^{2r2\eta }N_s\underset{i=1}{\overset{r}{}}S_{s_i,s_i^{}}^{𝔲(1)_4}\left(S_{q_i,q_i^{}}^{𝔲(1)_{2h_i}}\right)^{}\underset{i=1}{\overset{r^{}}{}}S_{l_i,l_i^{}}^{𝔰𝔲(2)_{k_i}}\hfill \\ \\ \text{ }\text{[}\text{(}1\frac{3}{4}\underset{i=r^{}+1}{\overset{r}{}}\delta _{l_i,k_i/2}\delta _{l_i^{},k_i/2}\text{)}\underset{i=r^{}+1}{\overset{r}{}}S_{l_i,l_i^{}}^{𝔰𝔲(2)_{k_i}}+\frac{1}{4}\psi \psi ^{}\underset{i=r^{}+1}{\overset{r}{}}\delta _{l_i,k_i/2}\delta _{l_i^{},k_i/2}S_{k_i/2,k_i/2}^{\mathrm{J}𝔰𝔲(2)_{k_i}}\text{]}\hfill \\ \\ \text{ }=\frac{2^{r2\eta }N_s}{_{i=1}^rh_i}\underset{i=1}{\overset{r}{}}\mathrm{e}^{\pi \mathrm{i}(q_iq_i^{}/h_is_is_i^{}/2)}\underset{i=1}{\overset{r^{}}{}}\mathrm{sin}\text{(}\frac{\left(l_i+1\right)\left(l_i^{}+1\right)\pi }{h_i}\text{)}\hfill \\ \\ \text{ }\text{[}\text{(}1\frac{3}{4}\underset{i=r^{}+1}{\overset{r}{}}\delta _{l_i,k_i/2}\delta _{l_i^{},k_i/2}\text{)}\underset{i=r^{}+1}{\overset{r}{}}\mathrm{sin}\text{(}\frac{\left(l_i+1\right)\left(l_i^{}+1\right)\pi }{h_i}\text{)}\hfill \\ \\ \text{ }+2^{2+(rr^{})/2}(1)^{_{i=r^{}+1}^rk_i/4}\psi \psi ^{}\underset{i=r^{}+1}{\overset{r}{}}\delta _{l_i,k_i/2}\delta _{l_i^{},k_i/2}h_i^{1/2}\text{]}.\hfill \end{array}$$
(3.35)
Recall that $`\eta =\mathrm{\hspace{0.17em}1}`$ when all levels are odd, in which case there are no fixed points, and $`\eta =\mathrm{\hspace{0.17em}0}`$ else. Upon insertion of (3.35) into the Verlinde formula, one obtains the fusion rule coefficients of $`𝒞^{\left(\mathrm{CY}\right)}`$; for $`\eta =\mathrm{\hspace{0.17em}0}`$ we arrive at the following expression (for $`\eta =\mathrm{\hspace{0.17em}1}`$ the formula is similar, but without the complications involving fixed points):
$$\begin{array}{c}(\mathrm{N}^{\left(\mathrm{CY}\right)})_{[(\lambda ,\sigma ,\xi ),\psi ],[(\lambda ^{},\sigma ^{},\xi ^{}),\psi ^{}]}^{[(\lambda ^{\prime \prime },\sigma ^{\prime \prime },\xi ^{\prime \prime }),\psi ^{\prime \prime }]}=\frac{1}{\left|𝒮_\lambda \right|\left|𝒮_\lambda ^{}\right|\left|𝒮_{\lambda ^{\prime \prime }}\right|}\underset{n=0,\mathrm{},N_s/41}{}\underset{\genfrac{}{}{0pt}{}{ϵ_j=0,1j=1,2\mathrm{},r}{_jϵ_j=0mod2}}{}\hfill \\ \\ \text{ }\underset{i=1}{\overset{r^{}}{}}\text{(}\delta _{\mathrm{\Delta }s_i+2n+2ϵ_i}\delta _{\mathrm{\Delta }q_i+2n}\mathrm{N}_{l_i,l_i^{}}^{l_i^{\prime \prime }}+\delta _{\mathrm{\Delta }s_i+2n+2ϵ_i+2}\delta _{\mathrm{\Delta }q_i+2n+h_i}\mathrm{N}_{l_i,k_il_i^{}}^{l_i^{\prime \prime }}\text{)}\hfill \\ \\ \text{ }\text{{}_{i=r^{}+1}^r\text{(}\delta _{\mathrm{\Delta }s_i+2n+2ϵ_i}\delta _{\mathrm{\Delta }q_i+2n}\mathrm{N}_{l_i,l_i^{}}^{l_i^{\prime \prime }}+\delta _{\mathrm{\Delta }s_i+2n+2ϵ_i+2}\delta _{\mathrm{\Delta }q_i+2n+h_i}\mathrm{N}_{l_i,k_il_i^{}}^{l_i^{\prime \prime }}\text{)}\hfill \\ \\ \text{ }+_{i=r^{}+1}^r\text{(}\delta _{\mathrm{\Delta }s_i+2n+2ϵ_i}\delta _{\mathrm{\Delta }q_i+2n}\mathrm{N}_{l_i,k_il_i^{}}^{l_i^{\prime \prime }}+\delta _{\mathrm{\Delta }s_i+2n+2ϵ_i+2}\delta _{\mathrm{\Delta }q_i+2n+h_i}\mathrm{N}_{l_i,l_i^{}}^{l_i^{\prime \prime }}\text{)}\hfill \\ \\ \text{ }+2_{i=r^{}+1}^r\delta _{\mathrm{\Delta }s_i+2n+2ϵ_i}\delta _{\mathrm{\Delta }q_i+2n}\text{[}\psi \psi ^{}_{i=r^{}+1}^r(S^{\mathrm{J}𝔰𝔲(2)_{k_i}})^2\delta _{l_i,k_i/2}\delta _{l_i^{},k_i/2}\frac{S_{l_i^{\prime \prime },k_i/2}^{}}{S_{0,k_i/2}}\hfill \\ \\ \text{ }+\psi ^{}\psi ^{\prime \prime }_{i=r^{}+1}^r\delta _{l_i^{},k_i/2}\delta _{l_i^{\prime \prime },k_i/2}\frac{S_{l_i,k_i/2}}{S_{0,k_i/2}}+\psi \psi ^{\prime \prime }_{i=r^{}+1}^r\delta _{l_i,k_i/2}\delta _{l_i^{\prime \prime },k_i/2}\frac{S_{l_i^{},k_i/2}}{S_{0,k_i/2}}\text{]}\hfill \\ \\ \text{ }+2_{i=r^{}+1}^r\delta _{\mathrm{\Delta }s_i+2n+2ϵ_i+2}\delta _{\mathrm{\Delta }q_i+2n+h_i}\text{[}\psi \psi ^{}_{i=r^{}+1}^r(S^{\mathrm{J}𝔰𝔲(2)_{k_i}})^2\delta _{l_i,k_i/2}\delta _{l_i^{},k_i/2}\frac{S_{k_il_i^{\prime \prime },k_i/2}^{}}{S_{0,k_i/2}}\hfill \\ \\ \text{ }+\psi ^{}\psi ^{\prime \prime }_{i=r^{}+1}^r\delta _{l_i^{},k_i/2}\delta _{l_i^{\prime \prime },k_i/2}\frac{S_{k_il_i,k_i/2}}{S_{0,k_i/2}}+\psi \psi ^{\prime \prime }_{i=r^{}+1}^r\delta _{l_i,k_i/2}\delta _{l_i^{\prime \prime },k_i/2}\frac{S_{k_il_i^{},k_i/2}}{S_{0,k_i/2}}\text{]}\text{}}.\hfill \end{array}$$
(3.36)
Concerning the notation, the following remarks are in order. First, we put
$$\mathrm{\Delta }s_i:=s_i+s_i^{}s_i^{\prime \prime },\mathrm{\Delta }q_i:=q_i+q_i^{}q_i^{\prime \prime }.$$
(3.37)
Second, the factors $`\delta _x\delta _{x,0}`$ represent the fusion coefficients of the various $`𝔲(1)`$ factors; in particular, an appropriate periodicity in their subscripts is understood. Third, we have separated the $`𝒦^{\left(\mathrm{CY}\right)}`$-summation into a part involving the simple currents without fixed points ($`_{n,ϵ_i}`$) and one that implements the order-two simple current $`\mathrm{L}`$ (the expression in curly brackets), compare formula (2.10). Further, the innermost (pairwise) summation takes into account the field identification in the minimal models. Finally, $`\psi \psi (\mathrm{L})\{\pm 1\}`$ corresponds to the two irreducible characters of the group $`_2=\{\mathrm{\Omega },\mathrm{L}\}`$.
The terms in formula (3.36) that involve factors of $`\psi `$ can be simplified by using the identity (3.34) and noting that
$$\frac{S_{l,k/2}}{S_{0,k/2}}=\mathrm{sin}\text{(}\frac{\pi }{2}(l+\mathrm{\hspace{0.17em}1})\text{)}=\{\begin{array}{cc}\hfill (1)^{l/2}& \text{if }l\text{ is even},\hfill \\ \hfill 0& \text{if }l\text{ is odd}.\hfill \end{array}$$
(3.38)
Both of the expressions in square brackets are thus non-zero only if the labels $`l_i,l_i^{},l_i^{\prime \prime }`$ are even for all $`i\{r^{}+1,r^{}+2,\mathrm{},r\}`$, in which case they read
$$\psi \psi ^{}(1)^{_{i=r^{}+1}^rk_i/4+l_i^{\prime \prime }/2}\mathrm{\Pi }\mathrm{\Pi }^{}+\psi ^{}\psi ^{\prime \prime }(1)^{_{i=r^{}+1}^rl_i/2}\mathrm{\Pi }^{}\mathrm{\Pi }^{\prime \prime }+\psi \psi ^{\prime \prime }(1)^{_{i=r^{}+1}^rl_i^{}/2}\mathrm{\Pi }\mathrm{\Pi }^{\prime \prime },$$
(3.39)
where we introduced $`\mathrm{\Pi }:=_{i=r^{}+1}^r\delta _{l_i,k_i/2}`$ and analogously $`\mathrm{\Pi }^{}`$ and $`\mathrm{\Pi }^{\prime \prime }`$. We also mention that the annulus coefficients for annuli with two boundary conditions of A<sub>0</sub>-type are just given by the fusion rules, so that the expression (3.36) directly provides us with the multiplicities of the corresponding open string states.
### 3.6 A<sub>0</sub>-type boundary conditions for minimal model tensor products
Having collected these results about minimal model tensor products, we are now in a position to write down the A<sub>0</sub>-type boundary states for the corresponding CalabiYau extensions $`𝒞^{\left(\mathrm{CY}\right)}`$. According to Cardy’s results, the labeling of the A<sub>0</sub>-type boundary conditions is precisely the same as for the primary fields of $`𝒞^{\left(\mathrm{CY}\right)}`$. Let us make this explicit. We start with the collection $`(\lambda ,\sigma ,\xi )(l_1,s_1,q_1,\mathrm{},l_r,s_r,q_r)`$ of labels ranging over $`l_i\{0,1,\mathrm{},k_i\}`$, $`s_i\{0,1,2,3\}`$, and $`q_i\{0,1,\mathrm{},2h_i\}`$. We then implement the various projections by imposing the following selections and identifications (both on bulk fields and on boundary conditions).
Selections:
We impose $`l_i+s_i+q_i\mathrm{\hspace{0.17em}2}`$ (minimal model selection rule), $`s_1+s_i\mathrm{\hspace{0.17em}2}`$ for all $`i=\mathrm{\hspace{0.17em}2},3,\mathrm{},r`$ (fermion alignment), and $`Q_{i=1}^r(q_i/h_i+s_i/2)s_1d/4`$ (charge projection).
Identifications:
We restrict the $`l_i`$ to the range $`0l_ik_i/2`$ (minimal model field identification). Representatives for the orbits with respect to the alignment currents $`\mathrm{w}_i`$ can be labeled by a single $`s\{0,1\}`$ when $`d/2+r`$ is odd, and $`s\{0,1,2,3\}`$ when $`d/2+r`$ is even; all the $`s_i`$ are equal to $`s`$ $`mod\mathrm{\hspace{0.17em}2}`$.
Implementation of the identification implied by the current $`\mathrm{u}`$ is more difficult; it involves the divisibility properties of the heights $`h_i`$, which in general do not have a simple structure. In special cases, for instance when $`\mathrm{lcm}_i\{h_i\}=h_j`$ for some $`j`$, the corresponding label $`q_j`$ can be set to zero using this identification.
Explicit formulas for boundary states of Gepner models were first presented in , where the boundary states were expressed in terms of the modular S-matrices of the $`𝔰𝔲(2)`$ and $`𝔲(1)`$ building blocks of the minimal models. However, as explained above, the chiral algebra $`𝔄^{\left(\mathrm{CY}\right)}`$ of the Gepner model is much larger than the chiral algebra $`𝔄_{k_1k_2\mathrm{}k_r}`$ of the tensor product $`𝒞_{k_1k_2\mathrm{}k_r}`$ of minimal models. Accordingly, in the bulk the modular transformations are described by the extended S-matrix $`S^{\left(\mathrm{CY}\right)}`$. The non-trivial information contained in $`S^{\left(\mathrm{CY}\right)}`$ that cannot be obtained from the tensor product S-matrix alone stems from the presence of fixed points under the CalabiYau extension. Once $`S^{\left(\mathrm{CY}\right)}`$ has been determined, the usual description of boundary conditions that preserve the full chiral algebra can be applied, though it is now to be formulated with the help of the matrix $`S^{\left(\mathrm{CY}\right)}`$ rather than $`S`$. It is therefore not guaranteed that the boundary states can be written in the form presented in . That the modular S-matrix must be ‘properly resolved’ for the construction of boundary states in Gepner models was emphasized in . It was also noticed in that some of the boundary states given in are not consistent with all projections in Gepner models, which explains certain discrepancies between the results of and . More recently, it has been established in that some of the boundary states of are not elementary.
In A-type boundary states for Gepner models with $`D=\mathrm{\hspace{0.17em}4}`$ and $`r=\mathrm{\hspace{0.17em}5}`$ were constructed by implementing the simple current extension (including in particular fixed point resolution) directly on the boundary states of the underlying $`N=\mathrm{\hspace{0.17em}2}`$ tensor product. On the other hand, the results reported above clearly also allow to obtain these A-type boundary states by performing the extension already in the bulk. When doing so, one obtains the boundary states by merely combining standard results for simple current extensions in the bulk (which lead in particular to the formula (3.35) for the modular S-matrix $`S^{\left(\mathrm{CY}\right)}`$ of the CalabiYau extension) with the general results of Cardy for symmetry preserving boundary conditions in arbitrary rational conformal field theories. Let us point out that this way we get the A-type boundary states for every Gepner model, i.e. for $`D=\mathrm{\hspace{0.17em}4},6,8`$ and for any allowed $`r`$. <sup>10</sup><sup>10</sup>10 To make contact to geometry, $`D/2+r`$ must be odd. This can, however, always be achieved by introducing a trivial $`k=\mathrm{\hspace{0.17em}0}`$ minimal model as an additional factor of the tensor product. Moreover, when proceeding in this manner the boundary states are completely determined, up to over-all normalization, by the algorithm. Thus there is e.g. no need to invoke integrality of the annulus coefficients in order to fix the relative strength of the contributions from twisted and untwisted sectors. Indeed, integrality of the annulus coefficients involving only A<sub>0</sub>-type boundary conditions coefficients is just the Verlinde formula as applied to the extended theory, i.e. to the CalabiYau extension.
More concretely, the A<sub>0</sub>-type boundary states are given by the expression (3.19),
$$|[(\lambda ,\sigma ,\xi ),\psi ]_{A_0}^{\left(\mathrm{CY}\right)}=\underset{[(\lambda ^{},\sigma ^{},\xi ^{}),\psi ^{}]}{}\frac{(S^{\left(\mathrm{CY}\right)})_{[(\lambda ,\sigma ,\xi ),\psi ],[(\lambda ^{},\sigma ^{},\xi ^{}),\psi ^{}]}}{\sqrt{(S^{\left(\mathrm{CY}\right)})_{[0,0,0],[\lambda ^{},\sigma ^{},\xi ^{}),\psi ^{}]}}}|[(\lambda ^{},\sigma ^{},\xi ^{}),\psi ^{}]$$
(3.40)
with $`S^{\left(\mathrm{CY}\right)}`$ as presented in formula (3.35) and with the summation ranging over all primaries of $`𝒞^{\left(\mathrm{CY}\right)}`$, <sup>11</sup><sup>11</sup>11 We may wish to split the boundary state $`|a_A^{\left(\mathrm{CY}\right)}`$ into its contributions involving only the Ishibashi states for full $`𝒦^{\left(\mathrm{CY}\right)}`$-orbits and only those for fixed points, respectively. Then the number $`_{i=r^{}+1}^rh_i^{1/2}`$ that is built in $`S^{\left(\mathrm{CY}\right)}`$ (see the last line of formula (3.35)) appears as a factor between the two parts. In this factor was obtained by imposing integrality of annulus coefficients. and the annulus coefficients for two boundary conditions of A<sub>0</sub>-type coincide with the fusion rules (3.36). In the special case that $`(\lambda ,\sigma ,\xi )`$ is not a fixed point of the current $`\mathrm{L}`$ (3.28), the expression (3.40) for the boundary state reduces to
$$\begin{array}{cc}|[(\lambda ,\sigma ,\xi )]_{A_0}^{\left(\mathrm{CY}\right)}\hfill & =\text{[}\frac{2^{r2\eta }N_s}{_{j=1}^rh_j}\text{]}^{1/2}\underset{[(\lambda ^{},\sigma ^{},\xi ^{}),\psi ^{}]}{}\hfill \\ & \\ & \text{ }\underset{i=1}{\overset{r}{}}\mathrm{e}^{\pi \mathrm{i}(q_iq_i^{}/h_is_is_i^{}/2)}\frac{\mathrm{sin}\text{(}(l_i+1)(l_i^{}+1)\pi /h_i\text{)}}{\text{(}\mathrm{sin}\text{(}(l_i^{}+1)\pi /h_i\text{)}\text{)}^{1/2}}|[(\lambda ^{},\sigma ^{},\xi ^{}),\psi ^{}].\hfill \end{array}$$
(3.41)
For the boundary state in the full Gepner model, one has to combine this expression with the space-time part of the boundary state, include a phase factor $`S_{\varpi ,\varpi ^{}}^{\mathrm{D}_{D/2+3}}/\sqrt{S_{o,\varpi ^{}}^{\mathrm{D}_{D/2+3}}}`$, a sign from undoing the bosonic string map, and a factor of $`2`$ from the remaining projections. The resulting formula is then essentially the one reported in ; but even in this special case our result differs in the power of 2 that appears in the prefactor.
In contrast, when $`(\lambda ,\sigma ,\xi )`$ is a fixed point of $`\mathrm{L}`$, then the additional terms in the formula (3.35) for $`S^{\left(\mathrm{CY}\right)}`$ contribute, and the expression for the boundary state gets a bit lengthier. Boundary states $`|[(\lambda ,\sigma ,\xi ),\psi ]_{A_0}^{\left(\mathrm{CY}\right)}`$ corresponding to resolved fixed points have been studied in . Apart from the proper power of 2, our result differs from the one presented in also by the absence of a factor of $`\sqrt{N_s}`$ in the $`\psi `$-dependent terms.
### 3.7 General A-type boundary conditions
Boundary conditions of automorphism type A<sub>θ</sub> with $`\theta \mathrm{\hspace{0.17em}0}`$ cannot be obtained from Cardy’s results alone. In this subsection, we explain more explicitly how the generalizations developed in allow to structurize and construct all A-type boundary conditions. We are faced with the following situation. Given the theory $`𝒞^{\left(\mathrm{CY}\right)}`$, obtained from $`𝒞_{\mathrm{int}}`$ by a simple current extension with $`𝒦^{\left(\mathrm{CY}\right)}`$, we want to construct all boundary conditions that preserve at least the chiral symmetry algebra of $`𝒞_{\mathrm{wsusy}}`$. Here $`𝒞_{\mathrm{wsusy}}`$ is the theory obtained from $`𝒞_{\mathrm{int}}`$ by extension with the alignment currents $`\mathrm{w}_i`$ alone; this means in particular that all boundary conditions preserve world sheet supersymmetry.
The solution to this problem is as follows. The set of all such boundary conditions is the set of irreducible representations of the classifying algebra $`(𝒞^{\left(\mathrm{CY}\right)}:𝒞_{\mathrm{wsusy}})`$. This set of boundary conditions can be divided into subsets of definite automorphism type, here labeled by the angle $`2\theta `$ with $`\theta /2\pi /N_s^{}`$. (Recall from subsection 3.3 that $`\theta `$ can be interpreted as specifying which space-time supersymmetry is preserved by a boundary condition.) Furthermore, each of these subsets furnishes the set of irreducible representations of an individual classifying algebra $`_{2\theta }`$ which, just like $``$, is a semisimple commutative associative algebra over $``$ , and one has
$$=\underset{\theta }{}_{2\theta }$$
(3.42)
as a direct sum of algebras over $``$ . We remark that the same situation arises in much simpler models in statistical mechanics, too. In the critical three-state Potts model, for instance, there are eight boundary conditions which come in two automorphism types that are distinguished by the reflection condition for the $`W_3`$-current . Accordingly, in that case the classifying algebra is a direct sum $`_+_{}`$ of a six- and a two-dimensional algebra; the irreducible representations of $`_+`$ provide the three fixed and the three mixed boundary conditions of the Potts model, while the ones of $`_{}`$ provide the free and ‘new’ conditions.
In general there is no simple relationship between the various individual classifying algebras $`_{2\theta }`$. In the case of Gepner models, however, it turns out that there exist symmetries generated by simple currents of the theory $`𝒞_{\mathrm{wsusy}}`$ (the so-called ‘phase symmetries’), which act on the set of boundary conditions and thereby relate individual classifying algebras for different $`\theta `$ with each other. Because of those symmetries, the boundary conditions for $`\theta \mathrm{\hspace{0.17em}0}`$ look still very similar to the Cardy case. It follows in particular that the combinatorics of fixed points and their resolution do not depend on the value of $`\theta `$. (This can already be deduced from formula (3.31) for the fixed point current $`\mathrm{L}`$.) Those symmetries between boundary conditions with different $`\theta `$ have been implicitly used in for a nice organization of A-type boundary conditions. Note, however, that this simplification is intrinsically linked to the special symmetries of $`N=\mathrm{\hspace{0.17em}2}`$ minimal models and cannot be expected to be present in general.
For a more detailed description, we recall that the group that furnishes the extension from $`𝒞_{\mathrm{wsusy}}`$ to $`𝒞^{\left(\mathrm{CY}\right)}`$ is generated by the image $`\mathrm{U}`$ of the simple current $`\mathrm{u}`$ (3.24) in the $`𝒞_{\mathrm{wsusy}}`$-theory, and that this current has order $`N_s^{}/2`$ with $`N_s^{}`$ as defined in formula (2.20), i.e. $`N_s^{}=N_s/2^\eta `$ with $`\eta \{0,1\}`$. The primary fields of $`𝒞_{\mathrm{wsusy}}`$ are labeled by $`\mathrm{\Lambda }:=[(\lambda ,\sigma ,\xi )]^{}`$ where the prime indicates that orbits are taken with respect to the group generated by the alignment currents only (which do not have fixed points); thus in particular $`\mathrm{U}=[\mathrm{u}]^{}`$. A distinguished basis $`\{\stackrel{~}{\mathrm{\Phi }}_\mathrm{\Lambda }\}`$ of the classifying algebra $``$ is then labeled by the set of all $`𝒞_{\mathrm{wsusy}}`$-fields $`\mathrm{\Lambda }`$ that have vanishing monodromy charge with respect to $`\mathrm{U}`$, together with a $`_2`$ character accounting for a possibly non-trivial stabilizer. Further, natural bases of the individual classifying algebras $`_{2\theta }`$ are provided by twisted sums $`\stackrel{~}{\mathrm{\Phi }}_\mathrm{\Lambda }^\theta :=_{j=0}^{(N_s^{}/2)1}\mathrm{e}^{2\mathrm{i}j\theta }\stackrel{~}{\mathrm{\Phi }}_{\mathrm{U}_{}^j\mathrm{\Lambda }}`$ of the basis elements $`\stackrel{~}{\mathrm{\Phi }}_\mathrm{\Lambda }`$ of $``$. The set of all boundary conditions, on the other hand, is labeled by the set of orbits $`[\mathrm{\Lambda },\psi ]^{\prime \prime }`$ of fields of $`𝒞_{\mathrm{wsusy}}`$ with respect to the extension by $`\mathrm{U}`$ (again with proper account for stabilizers), but without restriction on the value of $`Q_\mathrm{U}`$. The set of boundary conditions with fixed automorphism type A<sub>θ</sub> then corresponds to orbits $`[\mathrm{\Lambda },\psi ]^{\prime \prime }`$, with the same minimal model and fermion alignment selection and identification rules as for $`\theta =\mathrm{\hspace{0.17em}0}`$, but with different charge projection condition
$$Q_\mathrm{U}(\mathrm{\Lambda })=\theta /\pi mod.$$
(3.43)
Hereby in particular the counting of all A<sub>θ</sub>-type boundary states is reduced to simple, if lengthy, combinatorics, which can (and should) be directly implemented in a computer algorithm.
## 4 Remarks on B-type boundary conditions
So far we have restricted our attention exclusively to A-type boundary conditions. Recall that these need not preserve all of the chiral algebra $`𝔄^{\left(\mathrm{CY}\right)}`$ of the CalabiYau extension, but at least its subalgebra $`𝔄_{\mathrm{wsusy}}`$, which in turn contains $`𝔄_{\mathrm{int}}`$, the chiral algebra of the unextended inner sector $`𝒞_{\mathrm{int}}`$ (see the chain (2.19) of embeddings). Now we rather want to study boundary conditions that possess a non-trivial glueing condition already for the subalgebra $`𝔄_{\mathrm{int}}`$. To have a well-defined total supercurrent for the tensor product, and for compatibility with the fermion alignment, we must use the same automorphism of the $`N=\mathrm{\hspace{0.17em}2}`$ algebra for each factor of the tensor product $`𝒞_{\mathrm{int}}`$. As mentioned in subsection 3.1, generically the automorphism group we have to consider is the non-connected Lie group O$`(2)`$. In the connected component of the identity of O$`(2)`$ we can only use the identity element $`\omega =\mathrm{id}`$ itself; this gives rise to all A-type boundary conditions, of all subtypes discussed above. Here we are interested in automorphisms from the other connected component, which are characterized by
$$J_nJ_n,G_r^+\mathrm{e}^{\mathrm{i}\gamma }G_r^{},G_r^{}\mathrm{e}^{\mathrm{i}\gamma }G_r^+$$
(4.1)
with $`\gamma `$. Unlike in the case of the identity component, where only the identity could be chosen, we can allow for any arbitrary value of $`\gamma `$. Boundary conditions obtained this way are referred to as B-type conditions. As a distinguished subset, they include those where the breaking is induced by the mirror automorphism of the total $`N=\mathrm{\hspace{0.17em}2}`$ algebra, which is obtained for $`\gamma =\mathrm{\hspace{0.17em}0}`$:
$$J_nJ_n,G_r^\pm G_r^{}.$$
(4.2)
We will use the name B<sup>C</sup>-type for these specific boundary conditions; the subscript $`C`$ reminds of charge conjugation.
As compared to A-type conditions, we have to face a new problem, which also arises in other circumstances. Namely, we are only given an automorphism of the subalgebra $`𝔄_{\mathrm{int}}`$, but not an automorphism of the full chiral algebra $`𝔄^{\left(\mathrm{CY}\right)}`$ of our interest. We have already seen above in the case $`\omega =\mathrm{id}`$ that there may exist several different lifts of an automorphism to a larger algebra. But at a more fundamental level, there is even no guarantee that any of the automorphisms of B-type can be lifted at all. Indeed, as explained in appendix A, in a general simple current extension there can exist an obstruction to lifting a given automorphism.
Fortunately, in the specific situation of interest to us here, there is in fact no such obstruction. This follows from the fact that A- and B-type conditions get exchanged by the mirror map. In short, the B-type conditions of a Gepner model are well-defined for the full symmetry algebra $`𝔄^{\left(\mathrm{CY}\right)}`$ because they correspond to the A-type conditions of the mirror model, which by the results above are fully under control. In the bulk, the mirror map just amounts to applying charge conjugation in the inner sector. But at least in many Gepner models it can also be described alternatively by forming the modular invariant that is associated to a suitable group of simple currents of non-integral conformal weight. These currents implement the ‘phase symmetries’ of the minimal models. (In the literature this is again usually referred to as an orbifold construction.) Accordingly, the methods of , which show how to deal with boundary conditions for torus partition functions associated to simple currents of half-integral conformal weight, should be helpful for analyzing the B-type conditions.
But still the concrete details in the description of the B-type conditions are rather involved. When formulated in terms of the chiral algebra $`𝔄_{\mathrm{int}}`$, the complication manifests itself in the fact that the subalgebra that is preserved by the boundary conditions is neither contained in $`𝔄_{\mathrm{int}}`$ nor does it contain $`𝔄_{\mathrm{int}}`$. Roughly, one simultaneously tries to reduce the algebra – by taking the orbifold with respect to the automorphism group $`\mathrm{\Gamma }`$ in question, e.g. with respect to the $`_2`$ generated by the mirror automorphism (4.2) – and to extend it – by the simple current extension with the group $`𝒦^{\left(\mathrm{CY}\right)}`$. More concretely, we want to know a lift $`\widehat{\mathrm{\Gamma }}`$ of the orbifold group $`\mathrm{\Gamma }`$ to $`𝔄^{\left(\mathrm{CY}\right)}`$ and a chiral algebra extension $`\widehat{𝒦}`$ of $`(𝔄_{\mathrm{int}})^\mathrm{\Gamma }`$ to $`(𝔄^{\left(\mathrm{CY}\right)})^{\widehat{\mathrm{\Gamma }}}`$ such that the diagram
$$\begin{array}{c}\text{ }𝔄^{\left(\mathrm{CY}\right)}\\ \\ 𝒦^{\left(\mathrm{CY}\right)}\text{ }\widehat{\mathrm{\Gamma }}\text{ }\\ \\ \text{ }𝔄_{\mathrm{int}}(𝔄^{\left(\mathrm{CY}\right)})^{\widehat{\mathrm{\Gamma }}}\\ \\ \mathrm{\Gamma }\text{ }\widehat{𝒦}\text{ }\\ \\ (𝔄_{\mathrm{int}})^\mathrm{\Gamma }\text{ }\end{array}$$
(4.3)
is commutative. Then in particular the map indicated by the dashed arrow is well defined. On the other hand, for our study of boundary conditions the crucial issue is the relation between the largest ($`𝔄^{\left(\mathrm{CY}\right)}`$) and smallest ($`(𝔄_{\mathrm{int}})^\mathrm{\Gamma }`$) chiral algebra in the diagram. The latter is the orbifold of the former by some automorphism group $`\mathrm{\Gamma }^{\left(\mathrm{CY}\right)}`$, and conversely, $`𝔄^{\left(\mathrm{CY}\right)}`$ is a certain extension $``$ of $`(𝔄_{\mathrm{int}})^\mathrm{\Gamma }`$:
$$\begin{array}{cc}& 𝔄^{\left(\mathrm{CY}\right)}\\ & \\ \hfill 𝒦^{\left(\mathrm{CY}\right)}& \\ & \\ \hfill 𝔄_{\mathrm{int}}\text{ }& \\ & \\ \hfill \mathrm{\Gamma }& \\ & \\ & (𝔄_{\mathrm{int}})^\mathrm{\Gamma }\\ & \\ & \text{ }|\\ & \\ & |\\ & \\ & \mathrm{\Gamma }^{\left(\mathrm{CY}\right)}\text{ }\\ & \\ & \text{ }\end{array}$$
(4.4)
Note that $`\mathrm{\Gamma }^{\left(\mathrm{CY}\right)}`$ is not the direct product $`\mathrm{\Gamma }\times 𝒢`$ of the orbifold group $`\mathrm{\Gamma }`$ (here $`_2`$) and the simple current group $`𝒢`$ (here $`𝒦^{\left(\mathrm{CY}\right)}`$). In fact, the current $`\mathrm{u}`$ (3.24) is not invariant under the mirror automorphism, so that even if one dealt with a group constructed out of $`\mathrm{\Gamma }`$ and $`𝒢`$, it could definitely not be their direct product. Closer inspection shows that the mirror orbifold $`(𝔄_{\mathrm{int}})^\mathrm{\Gamma }`$ of the tensor product theory is an orbifold of the CalabiYau extension $`𝔄^{\left(\mathrm{CY}\right)}`$ by a non-abelian group $`\mathrm{\Gamma }^{\left(\mathrm{CY}\right)}`$. Conversely, in order to obtain the orbifold of $`𝒞^{\left(\mathrm{CY}\right)}`$ from the mirror orbifold of $`𝒞_{\mathrm{int}}`$ one must consider an extension $``$ of the chiral algebra by fields not all of which are simple currents, though still all have integral quantum dimension. One may expect that $`\widehat{\mathrm{\Gamma }}`$ is a subgroup of $`\mathrm{\Gamma }^{\left(\mathrm{CY}\right)}`$, and that the fields in $`\widehat{𝒦}`$ form a subset of those in $``$. Some of the necessary mathematical tools for attacking this problem have recently been established (see e.g. ), but they are not sufficient to obtain a complete description of the associated B-type boundary conditions.
However, we can still draw some general conclusions about B-type conditions by again invoking mirror symmetry. As a matter of fact, the statement that A- and B-type conditions get exchanged by the mirror symmetry has to be refined, because the mirror automorphism (4.2) only involves the total $`N=\mathrm{\hspace{0.17em}2}`$ superconformal algebra of the tensor product, whereas as noticed above for the complete specification of A-type (and likewise of B-type) boundary conditions we have to prescribe the action of the automorphism on the full chiral algebra of the CalabiYau extension. <sup>12</sup><sup>12</sup>12 In addition, one should be aware of the fact that initially the mirror map refers to the situation before applying the bosonic string map. However, the mirror map is compatible with the bosonic string map, so that we can directly apply it here. The point is that, in conformal field theory terminology, the mirror map is nothing but charge conjugation. But as a consequence of formula (2.2) we have $`\stackrel{~}{C}B=BC`$, where $`C`$ is the charge conjugation matrix of the bosonic theory (i.e., after the bosonic string map $`B`$) and $`\stackrel{~}{C}`$ the charge conjugation matrix of the supersymmetric theory (before the bosonic string map). To do so, we use the fact that, as a special case of more general T-duality relations , those boundary conditions of a conformal field theory with torus partition function $`Z`$ which leave invariant precisely the subalgebra of the chiral algebra that is pointwise fixed under the charge conjugation automorphism $`\omega ^C`$ can be put in one-to-one correspondence with the boundary conditions that preserve the full chiral algebra of a different conformal field theory, namely the one with torus partition function $`CZ`$. <sup>13</sup><sup>13</sup>13 The simplest instance of this phenomenon is the exchange between Neumann and Dirichlet conditions in the theory of a single free boson. It follows that the A<sub>0</sub>-type boundary conditions of a Gepner model correspond to the B<sup>C</sup>-type boundary conditions of its mirror model, and vice versa. How this generalizes to the more general A- and B-type conditions will be discussed elsewhere. In particular, this consideration tells us that there is no obstruction to a lift of the mirror automorphism of the total $`N=\mathrm{\hspace{0.17em}2}`$ algebra to the full chiral algebra of the CalabiYau extension.
Finally let us mention that whenever the tensor product $`𝒞_{\mathrm{int}}`$ contains identical factors, there are additional automorphisms, beyond those belonging to O(2) discussed above, that can be modded out without spoiling world sheet supersymmetry, namely the permutation symmetries that interchange the identical factors.
## 5 Fixed points and singularities
We conclude this paper with a few comments on the relationship between fixed point structures in Gepner models and various other ‘singular’ structures that occur in the analysis of string compactifications on CalabiYau spaces and in Gepner models. The following list summarizes various such structures.
It is generally believed that a Gepner model describes the exact solution of string theory compactified on a certain CalabiYau manifold, at a specific point of its moduli space. There is a general prescription for finding the polynomial constraints that provide the embedding of the CalabiYau manifold in a weighted projective space, see e.g. . When carrying this construction through, one encounters the problem that the variety defined by those polynomial constraints is not smooth, but has singularities. It is only after resolving these singularities that one obtains the CalabiYau manifold.
The moduli of the CalabiYau manifold are related to the $`(c,c)`$ and $`(a,c)`$ rings of the $`N=\mathrm{\hspace{0.17em}2}`$ superconformal field theory . For instance, deformations of the complex structure of the CalabiYau space correspond to Gepner model fields that are chiral primaries with respect to both the left- and right-moving chiral algebras, with total $`𝔲(1)`$ charges $`1`$. Fields that have in addition identical left- and right-moving labels in each minimal model can be related to polynomial deformations of the complex structure, i.e. to a change in the defining polynomial. But it has been pointed out long ago that polynomial deformations give neither a complete nor an unambiguous enumeration of complex structure deformations. In the Gepner model, this can be related to the existence of twisted, i.e. left-right asymmetric, modes in the $`(c,c)`$ ring, see e.g. . <sup>14</sup><sup>14</sup>14 These states can be obtained directly by the algorithms incorporating the simple current extension, e.g. by the computer programs that were used to produce the bulk spectra in .
In the study of boundary states in Gepner models and their comparison with geometric D-branes and bundles in the corresponding CalabiYau manifolds , it was noticed that several boundary states constructed in are not elementary, in the sense that the annulus coefficient of the vacuum field is larger than one . In it has been argued that this can be understood if one assumes that the relevant boundary states correspond to branes carrying reducible bundles.
There have been speculations that (some of) these different singular structures are intimately related. Now in order for any definite relationship between geometrical and conformal field theory structures to exist, at least the combinatorial data characterizing them should be similar. In the situation at hand, this is actually not the case.
Let us first recall some of our results concerning A-type boundary conditions. As we have explained, the construction of A-type boundary states in Gepner models is completely under control, and only existing technology is needed. <sup>15</sup><sup>15</sup>15 In particular, the fixed points can already be understood entirely at the chiral level in closed string theory. It is therefore e.g. unnecessary to study open string partition functions in order to fix normalization factors (as advocated in ). Rather, the annulus coefficients are non-negative integers by construction. Also, nothing is gained by separating boundary states or partition functions into a part from an ‘untwisted’ and one from a ‘twisted’ sector; all the relevant structures are already present in the closed string theory. We have also seen that the fixed points in Gepner models are always of $`_2`$ type, i.e. the only non-trivial stabilizer group that occurs is of the form $`\{\mathrm{\Omega },\mathrm{L}\}`$. On the other hand, singularities in the construction of the CalabiYau manifold as a complete intersection in a weighted projective space can occur if the weights have common divisors, and are locally of the type $`^d/_{\mathrm{}}`$ where $`d=\mathrm{\hspace{0.17em}2},3`$ and $`\mathrm{}`$ can be any integer, see, for example, .
It is also unlikely that there is any relationship between twisted modes in the $`(c,c)`$ ring and singularities in the construction of the CalabiYau space. By inspecting lists of twisted $`(c,c)`$ fields for particular Gepner models, it is easily checked that they do not display any $`_{\mathrm{}}`$ structure. And indeed, the existence of non-polynomial deformations is not related to a singularity occuring in the construction of the manifold, but to the presence of an obstruction in the relevant cohomologies of the family of smooth manifolds . On the conformal field theory side, twisted modes may also be related to the presence of an enhanced symmetry. For instance, compactification on K3 leads to $`N=\mathrm{\hspace{0.17em}4}`$ world sheet supersymmetry, and half of the modes in the $`(c,c)`$ and $`(c,a)`$ rings are twisted. Incidentally, in such a situation boundary conditions might break different parts of this extended symmetry, so that a finer classification of boundary conditions is possible.
Furthermore, fixed points in the Gepner extension do not seem to be related to the presence of twisted modes in the $`(c,c)`$ ring. In fact, the Gepner extension has fixed points whenever at least one level is even. But many of those Gepner models do not possess any twisted modes in their $`(c,c)`$ ring. A simple example is provided by the Gepner model $`(2,2)`$ which does have fixed points, but simply corresponds to a torus embedded in one-dimensional projective space, without any singularities.
Finally we would like to point out that in our opinion several interesting problems concerning boundary states in Gepner models and branes on CalabiYau compactifications are still open. For example, the construction of states charged under twisted modes in the $`(c,c)`$ ring, as well as a better understanding of composite boundary states and their connection to reducible bundles, is highly desirable.
## Appendix A Lifting orbifold automorphisms
In the main text we have seen that when analyzing B-type boundary conditions one must study the interplay between simple current extensions and orbifolds. Here we investigate this issue, which is of interest also in other situations, in its own right. Thus consider an arbitrary rational conformal field theory with chiral algebra $`𝔄`$, a group $`𝒢`$ of simple currents (with integral conformal weight) of the theory, and a group $`\mathrm{\Gamma }`$ of automorphisms of $`𝔄`$. Every automorphism $`\omega `$ of $`𝔄`$ induces a permutation of the (labels for) primary fields, which is an automorphism of the fusion rules and which we denote by $`\omega ^{}`$. By including the simple currents into the chiral algebra, one obtains an extended theory with chiral algebra $`𝔄_{\mathrm{ext}}𝔄`$, while by dividing out the automorphisms in $`\mathrm{\Gamma }`$ one obtains an orbifold theory with chiral algebra $`𝔄^\mathrm{\Gamma }𝔄`$. Our goal is then to make sense of the symbol $`𝔄_{\mathrm{ext}}^\mathrm{\Gamma }`$. For simplicity and definiteness we restrict to the case $`\mathrm{\Gamma }=_2`$, i.e. there is only a single non-trivial automorphism $`\omega `$ and it has order two.
A necessary prerequisite for attaining this goal is to lift the automorphism $`\omega `$ to some automorphism $`\omega _{\mathrm{ext}}`$ of the extended chiral algebra $`𝔄_{\mathrm{ext}}`$. As a matter of fact, our first task should be to determine whether such a lift is possible at all. Indeed there can be an obstruction, and this will be studied below. But for the moment let us restrict to those cases where the lifting of $`\omega `$ is not obstructed. In that case we need to investigate the uniqueness of the lift. We already know from the discussion in the main text that typically even the identity automorphism of $`𝔄`$ will possess several distinct lifts to $`𝔄_{\mathrm{ext}}`$. To be more concrete, we use the fact that we can characterize $`𝔄`$ as the fixed algebra of $`𝔄_{\mathrm{ext}}`$ under an action of the group $`G_{\mathrm{ext}}:=𝒢^{}`$ of characters of the simple current group $`𝒢`$. Further, $`𝔄^\omega `$ is the fixed algebra in $`𝔄`$ under $`\omega `$. Thus we can characterize $`𝔄^\omega `$ equally well as the subalgebra of the big algebra $`𝔄_{\mathrm{ext}}`$ under the combined action of $`G_{\mathrm{ext}}`$ and $`\omega _{\mathrm{ext}}`$, i.e. under the action of the group $`H_{\mathrm{ext}}`$ generated by $`G_{\mathrm{ext}}`$ and $`\omega _{\mathrm{ext}}`$. Note that we do not assume here that $`\omega _{\mathrm{ext}}`$ has order two.
We can then again employ the result from the Galois theory for vertex operator algebras that the possible chiral algebras between $`𝔄^\omega `$ and $`𝔄_{\mathrm{ext}}`$ are in one-to-one correspondence with subgroups of the group $`H_{\mathrm{ext}}`$. Now since $`𝔄`$ is the fixed point subalgebra of $`𝔄_{\mathrm{ext}}`$ under $`G_{\mathrm{ext}}`$, for every $`gG_{\mathrm{ext}}`$ the element $`g_\omega :=\omega _{\mathrm{ext}}^1g\omega _{\mathrm{ext}}H_{\mathrm{ext}}`$ acts on $`𝔄`$ as $`\omega ^11\omega =\mathrm{\hspace{0.17em}1}`$, and by the Galois correspondence every such $`g_\omega `$ already lies in $`G_{\mathrm{ext}}`$. This tells us that $`G_{\mathrm{ext}}`$ is a normal abelian subgroup of $`H_{\mathrm{ext}}`$, and the conjugation by $`\omega _{\mathrm{ext}}`$ acts on $`G_{\mathrm{ext}}`$ by an outer automorphism $`\widehat{\omega }`$, i.e. $`\omega _{\mathrm{ext}}^1g\omega _{\mathrm{ext}}=\widehat{\omega }(g)`$. As $`\omega _{\mathrm{ext}}^2G_{\mathrm{ext}}`$ and $`G_{\mathrm{ext}}`$ is abelian, the automorphism $`\widehat{\omega }`$ has order two. So we have the structure <sup>16</sup><sup>16</sup>16 It is instructive to think of $`H_{\mathrm{ext}}`$ like of a compact Lie group with two connected components. The ‘identity component’ is $`G_{\mathrm{ext}}`$, and it has a natural unit element, while the other component, consisting of the automorphisms of $`𝔄_{\mathrm{ext}}`$ whose restriction to $`𝔄`$ acts like $`\omega `$, does not possess a natural base point.
$$0G_{\mathrm{ext}}H_{\mathrm{ext}}_2\mathrm{\hspace{0.17em}0}.$$
(A.1)
As $`\omega _{\mathrm{ext}}`$ does not necessarily have order two, the term ‘orbifold of $`𝔄_{\mathrm{ext}}`$ by $`\omega _{\mathrm{ext}}`$’ is to be interpreted as the orbifold by the cyclic subgroup of $`H_{\mathrm{ext}}`$ that is generated by $`\omega _{\mathrm{ext}}`$. But those elements of this subgroup that are of the form $`(\omega _{\mathrm{ext}})^{2n}`$ form a cyclic subgroup $`G_{\mathrm{ext}}^0`$ of $`G_{\mathrm{ext}}`$, so that we may as well perform the orbifolding stepwise, first by $`G_{\mathrm{ext}}^0`$ and afterwards by $`\omega _{\mathrm{ext}}`$ which then has order two on $`(𝔄_{\mathrm{ext}})^{G_{\mathrm{ext}}^0}`$. It follows that at the price of possibly working with a different simple current extension than the original one, we may restrict our attention to the case when $`\omega _{\mathrm{ext}}`$ has order two.
A simple illustration of the non-uniqueness of the lift is provided by the following $`c=\mathrm{\hspace{0.17em}1}`$ theories. The original theory is the rational free boson $`X`$ at compactification radius $`R^2=mn^2`$ with $`n`$ integral and $`m`$ even integral. Thus the chiral algebra $`𝔄`$ is generated by the current $`j=\mathrm{i}X`$, giving rise to $`𝔲(1)_{mn^2}`$, and by the two fields $`\mathrm{\Phi }_\pm =\mathrm{exp}(\pm \mathrm{i}\sqrt{m}nX)`$; there are $`mn^2`$ primary fields, which we label by the integers from 0 to $`mn^21`$ and each of which is a simple current. Extending this theory by the simple currents $`𝒢=\{\mathrm{}mn|\mathrm{}=\mathrm{\hspace{0.17em}0},1,\mathrm{},n1\}_n`$, one obtains the theory of a free boson with compactification radius $`R^2=m`$ and chiral algebra $`𝔄_{\mathrm{ext}}`$ generated by $`𝔲(1)_m`$ and $`\mathrm{\Phi }_\pm ^{\mathrm{ext}}=\mathrm{exp}(\pm \mathrm{i}\sqrt{m}X)`$. On the other hand, dividing out the charge conjugation automorphism $`\omega `$ from $`𝔄`$ one arrives at the $`_2`$ orbifold of the free boson, with $`mn^2/2+7`$ primary fields. The map $`\omega `$ acts on the fields generating $`𝔄`$ as
$$\omega (j)=j,\omega (\mathrm{\Phi }_\pm )=\mathrm{\Phi }_{}.$$
(A.2)
It can be lifted to $`𝔄_{\mathrm{ext}}`$ as
$$\omega _{\mathrm{ext}}(j)=j,\omega _{\mathrm{ext}}(\mathrm{\Phi }_\pm ^{\mathrm{ext}})=\zeta ^{\pm 1}\mathrm{\Phi }_{}^{\mathrm{ext}},$$
(A.3)
where $`\zeta `$ is an arbitrary $`n`$th root of unity. In terms of the free boson $`X`$, this reads
$$\omega _{\mathrm{ext}}(X)=X+\frac{2\pi \mathrm{}}{n\sqrt{m}}$$
(A.4)
with some $`\mathrm{}\{0,1,\mathrm{},n1\}`$. This example also displays nicely that even the identity map of $`𝔄`$ can typically be lifted in several inequivalent ways. Clearly, for each $`n`$th root $`\zeta `$ of unity, the map
$$\mathrm{id}_{\mathrm{ext}}(j)=j,\mathrm{id}_{\mathrm{ext}}(\mathrm{\Phi }_\pm ^{\mathrm{ext}})=\zeta ^{\pm 1}\mathrm{\Phi }_\pm ^{\mathrm{ext}}$$
(A.5)
(acting on the free boson field as $`XX+\mathrm{\hspace{0.17em}2}\pi \mathrm{}/n\sqrt{m}`$ for $`\zeta =\mathrm{exp}(2\pi \mathrm{i}\mathrm{}/n)`$) is an automorphism of $`𝔄_{\mathrm{ext}}`$ and restricts to the identity map on $`𝔄`$. Note that while the order of the map given by (A.3) is 2 independently of the value of $`\zeta `$, this is no longer true for the map (A.5).
Let us now come back to the issue of existence of $`\omega _{\mathrm{ext}}`$. Under the most general circumstances such a lift may actually not exist at all. First, clearly the compatibility condition $`\omega ^{}(𝒢)=𝒢`$ must be satisfied. In the sequel we assume that this is the case. (In the situation of our interest, this condition is indeed met. Also, the condition is automatically fulfilled whenever an extension is by all integer spin simple currents of a given theory.) But even with this assumption the existence of a lift $`\omega _{\mathrm{ext}}`$ is not guaranteed. Rather, one has to study the relation between the subgroup
$$𝒢_0:=\{\mathrm{J}𝒢|\omega ^{}(\mathrm{J})=\mathrm{J}\}$$
(A.6)
of $`𝒢`$ and the group $`𝒢^\omega `$ of simple currents of the $`𝔄^\omega `$-theory. By general orbifold rules , each $`\mathrm{J}𝒢_0`$ gives rise to two simple currents $`\mathrm{J}_\pm `$ in the untwisted sector of the orbifold theory $`𝔄^\omega `$. The fields $`\mathrm{J}_\pm `$ form a subgroup $`𝒢_0^\omega `$ of $`𝒢^\omega `$. Thus $`𝒢_0^\omega `$ is a $`_2`$-extension of $`𝒢_0`$, i.e. there is an exact sequence
$$0_2𝒢_0^\omega \stackrel{\pi }{}𝒢_00,$$
(A.7)
where the projection $`\pi `$ acts as $`\pi (\mathrm{J}_+)=\mathrm{J}=\pi (\mathrm{J}_{})`$. But $`𝒢_0^\omega `$ is not necessarily a direct product of $`𝒢_0`$ with $`_2`$; the obstruction is expressed by an element $`[ϵ]`$ in $`H^2(𝒢_0,_2)`$.
Inspecting the fusion rules among the fields $`\mathrm{J}_\pm `$ in the orbifold theory $`𝔄^\omega `$, one finds that this cohomology class $`[ϵ]`$ has the following conformal field theory interpretation. The associated commutator cocycle $`\overline{ϵ}(\mathrm{J}_1,\mathrm{J}_2):=ϵ(\mathrm{J}_1,\mathrm{J}_2)/ϵ(\mathrm{J}_2,\mathrm{J}_1)`$ (which only depends on the class $`[ϵ]`$ and not on the choice of a representative $`ϵ`$) can be expressed as
$$\overline{ϵ}(\mathrm{J}_1,\mathrm{J}_2)=\underset{\dot{\mu }}{}S_{\mathrm{J}_1,\dot{\mu }}^{(0)}S_{\mathrm{J}_2,\dot{\mu }}^{(0)}(S_{\mathrm{J}_1\mathrm{J}_2,\dot{\mu }}^{(0)})^{}/S_{\mathrm{\Omega },\dot{\mu }}^{(0)}.$$
(A.8)
Here $`S^{(0)}`$ denotes the matrix that governs the modular behavior of the differences $`\chi _{\mu _+}\chi _\mu _{}`$ of orbifold characters in the untwisted sector coming from the same $`𝔄`$-primary. More specifically, under $`\tau 1/\tau `$ these differences become linear combinations of characters $`\chi _{\dot{\mu }}`$ in the twisted sector, with the coefficients given by $`S^{(0)}`$ (for more details, see ). By the consistency of the orbifold fusion rules, $`\overline{ϵ}`$ can only take the values $`\pm 1`$ and satisfies $`\overline{ϵ}(\mathrm{J}_1,\mathrm{J}_2)=\overline{ϵ}(\mathrm{J}_2,\mathrm{J}_1)`$, and therefore indeed determines a unique class $`[ϵ]`$ in $`H^2(𝒢_0,_2)`$. In the special case where $`𝒢_0=𝒢`$, for the construction of $`𝔄_{\mathrm{ext}}^\omega `$ we must pick, for each $`\mathrm{J}𝒢`$, one of the fields $`\mathrm{J}_\pm 𝒢^\omega `$ of the $`\omega `$-orbifold, in such a manner that the chosen set of representatives closes under fusion. Thus we must find a section $`\sigma :𝒢^0𝒢^\omega `$ for the exact sequence (A.7). Such a section exists only if the extension is trivial, i.e. if $`[ϵ]=\mathrm{\hspace{0.17em}1}`$. In conclusion, there is an obstruction to the lift of $`\omega `$ to an automorphism $`\omega _{\mathrm{ext}}`$ of $`𝔄_{\mathrm{ext}}`$; moreover, it is controlled by the twisted sector of the orbifold, and hence computable.
For various classes of orbifold constructions the presence of an obstruction can be decided without too much effort. For instance, from the results of it follows immediately that there is no obstruction when $`\omega `$ is an automorphism of the chiral algebra of a WZW model that comes from an inner automorphism of the underlying simple Lie algebra. Similarly, no obstruction is present for the charge conjugation orbifold of a single free boson <sup>17</sup><sup>17</sup>17 When $`n`$ as introduced before formula (A.2) is odd, then $`𝒢_0=\{0\}`$, and the statement is trivial. When $`n`$ is even, then $`𝒢_0=\{0,mn^2/2\}_2`$, while the simple current group $`𝒢^\omega `$ of the orbifold is $`_2\times _2`$, so $`𝒢^\omega `$ is a trivial extension of $`𝒢_0`$. and for arbitrary permutation orbifolds. As a matter of fact, we do not know of any orbifold construction where the obstruction is present. It is tempting to expect that the obstruction is indeed absent in all cases that appear in conformal field theory, but so far we do not have any general argument to this effect. In any case, in all the applications in the main text we are able to show that the obstruction is absent. Therefore in the present paper we do not attempt to push this issue further.
The fields in $`𝒢`$ which are not contained in the subgroup $`𝒢_0`$ (A.6) come in pairs $`\mathrm{J}`$ and $`\omega ^{}(\mathrm{J})`$, and each such pair gives rise to a single field in $`𝔄^\omega `$, which has quantum dimension 2, i.e. is not a simple current any longer. As a consequence, these fields do not have a direct influence on the presence of an obstruction. On the other hand, even when there is no obstruction, it turns out to be quite non-trivial to describe such non-simple current fields in sufficiently explicit terms in concrete models. In particular, consistency of the fusion rules of the $`𝔄`$-theory does not seem to be of any help.
Acknowledgement
We are grateful to Jens Fjelstad, Jürg Fröhlich, Lennaert Huiszoon, Peter Kaste, Wolfgang Lerche, Andy Lütken, and Bert Schellekens for helpful discussions and comments.
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# Spin-flip transitions between Zeeman sublevels in semiconductor quantum dots
## I Introduction.
Quantum dots (QD’s) in semiconductor heterostructures provide a unique opportunity to study the properties of the electron quantum states in detail and manipulate the electrons in these ”artificial atoms” in a controllable way (see reviews ). The shape and size of quantum dots can be varied by changing the gate voltage. This also tunes the number of electrons in the dot. Besides, the electronic states can be significantly modified by a magnetic field applied perpendicular to the plane of the heterostructure.
Quantum dots are considered as possible candidates for building a quantum computer. The crucial point of the idea is the necessity to couple dots coherently and keep coherence on sufficiently long time scales. In this respect, there is a great demand in the theoretical estimation of the typical spin dephasing time of the electron in the QD. In our previous work we have shown that the localized character of the electron wave functions in the QD’s suppresses the most effective intrinsic spin-flip mechanisms related to the absence of inversion symmetry in GaAs-like crystals. This leads to an unusually low rate of spin-flip transitions. However, in Ref. we concentrated on the case of inelastic transitions between the neighbouring quantized energy levels in the dot which corresponds to a relatively large energy transfer. On the other hand, the quantum bit was proposed to involve two Zeeman sublevels of the same orbital level. Therefore, in the present work we consider the transitions between such sublevels. Since the transition involves a fairly small energy transfer, the resuls are very different from those of Ref. .
As in Ref. , we concentrate on the spin-flip processes due to the spin-orbit interaction. This is the main source of the spin-flips for the three- and two-dimensional electron states in the GaAs-type crystal without an inversion center. Besides, in such a polar-type crystal one finds a strong coupling of electrons to the bosonic environment via the piezo-electric interaction with acoustic phonons. The combination of these two mechanisms provides an effective spin-lattice relaxation of free carriers in $`A_{III}B_V`$ semiconductors and heterostructures. We show, however, that the spin-lattice relaxation for the electron localized in the QD is much less effective.
We have calculated the rates for the different spin-orbit related mechanisms which cause a spin-flip in the course of the phonon-assistant transition between the Zeeman sublevels. Besides, we have estimated the spin-flip rates due to several other mechanisms, for example, due to the fluctuating magnetic field produced by the fluctuating electron density in the leads or due to modulation of the hyperfine coupling with nuclei by lattice vibrations.
## II Spin-orbit mechanisms.
We consider the case of a strong confinement in the z-direction and the typical lateral dot size is of the order of thousand $`\AA `$ and much larger than the width of the 2D layer in the z-direction. We begin with the following one-electron Hamiltonian that is derived from the Kane model (see Ref. ) and describes the 2D electrons of the conduction band in the presence of magnetic field $`𝐁`$, lateral confining potential $`U(𝐫)`$ and the phonons:
$`\widehat{}={\displaystyle \frac{\widehat{𝐩}^2}{2m}}+U(𝐫)+U_{ph}(𝐫,t)+{\displaystyle \frac{1}{2}}g\mu _B\widehat{𝝈}𝐁+{\displaystyle \underset{i=1}{\overset{3}{}}}\widehat{}_i;`$ (1)
$`\widehat{}_1=\beta (\widehat{\sigma }_x\widehat{p}_x+\widehat{\sigma }_y\widehat{p}_y);\beta ={\displaystyle \frac{2}{3}}p_z^2{\displaystyle \frac{\mathrm{\Delta }}{(2mE_g)^{1/2}m_{cv}E_g}};\widehat{}_2={\displaystyle \frac{1}{2}}V_0\widehat{𝝈}\widehat{𝝋};\widehat{}_3=\stackrel{~}{g}\mu _B{\displaystyle \underset{ik}{}}u_{ik}\widehat{\sigma }_iB_k`$ (2)
Here $`\widehat{𝐩}=i\mathrm{}\mathbf{}+(e/c)𝐀`$ is the 2D electron momentum operator, $`m`$ the effective mass, $`\widehat{𝝈}`$ the Pauli matrices. Axes $`x,y,z`$ coincide with the main crystallographic ones with the z-axis along the normal to the 2D plane (the orientation). The magnetic field has an arbitrary direction. The third term describes the spin-independent interaction with the phonons, including the piezoelectric ones. The fourth term is the Zeeman energy. The other three terms describe all possible spin-orbit effects. $`\widehat{}_1`$ stems from the absence of the inversion symmetry in the bulk. Velocity $`\beta `$ takes the values in the interval $`(1÷3)10^5cm/s`$ for GaAs heterostructures. $`\widehat{}_2`$ describes the spin-orbit splitting of the electron spectrum due to the strain field produced by the acoustic phonons. There $`\widehat{\phi }_x=(1/2)\{u_{xy},\widehat{p}_y\}_+`$, $`\widehat{\phi }_y=(1/2)\{u_{yx},\widehat{p}_x\}_+`$, $`\widehat{\phi }_z=(1/2)\{u_{zx},\widehat{p}_x\}_+(1/2)\{u_{zy},\widehat{p}_y\}_+`$, where $`\{,\}_+`$ denotes the anticommutator, $`u_{ij}`$ is the lattice strain tensor, and $`V_0=8\times 10^7cm/s`$. In GaAs the electron $`g`$-factor ($`g=0.44`$) differs strongly from the free electron value $`g_0=2`$ owing to the strong spin-orbit interaction which mixes the valence band and conduction band states . The admixture depends on the lattice deformation. Coefficient $`\stackrel{~}{g}`$ can be found within the Kane approach, $`\stackrel{~}{g}=(2m_0/\sqrt{3}m)(\mathrm{\Delta }/E_g)(d/E_g)`$, $`d=4.5`$ eV is one of the three deformation constants describing the strain effect on the hole band splitting , $`\stackrel{~}{g}10.4`$.
The three terms in Hamiltonian (2) correspond to the three distinct mechanisms of the spin flip. The first mechanism is due to the spin-orbit admixture of state with an opposite spin. While without the spin-orbit interaction the Zeeman sublevels correspond to the orbital state with the spin up or down, the spin-orbit terms provide a small admixture of the state of the opposite spin to each sublevel. This enables the phonon-assistant transition between the two states. This mechanism corresponds to term $`\widehat{}_1`$.
The second and the third mechanisms are described by the $`\widehat{}_2`$ and $`\widehat{}_3`$ terms and correspond to two distinct kinds of direct spin-phonon coupling. Below we show that the admixture mechanism is actually a dominant one.
### A Admixture mechanism.
Let us show that for this mechanism the matrix element of the $`U_{ph}(𝐫,t)`$ operator for the spin-flip transition between the Zeeman levels is proportional to the product of the Zeeman energy and the phonon strain field. Since we deal with a small energy transfer, we consider only the interaction with piezo-phonons, hence, for mode $`𝐪\alpha `$, $`\alpha =l,t`$ , we have
$$U_{ph}^{𝐪\alpha }(𝐫,t)=\sqrt{\mathrm{}/2\rho \omega _{𝐪\alpha }}\mathrm{exp}(i\mathrm{𝐪𝐫}i\omega _{𝐪\alpha }t)eA_{𝐪\alpha }b_{𝐪\alpha }^++c.c.;A_{𝐪\alpha }=\xi _i\xi _k\beta _{ikj}e_{𝐪\alpha }^j,$$
(3)
where the effective piezo-electric modulus $`A_{𝐪\alpha }`$ of wave $`𝐪\alpha `$ has been introduced, $`\beta _{ikj}`$ is the piezo-tensor, $`𝝃=𝐪/q`$, $`𝐪`$ the phonon wave vector, $`𝐞`$ the phonon unit polarization vector and $`\rho `$ is the crystal mass density. For the crystal of cubic symmetry without an inversion center (class $`T_d`$) tensor $`\beta _{ikj}`$ has only non-zero components (all of them equal to each other) when all three indexes $`i,k,j`$ are different, $`\beta _{xyz}=\beta _{xzy}=\mathrm{}.=h_{14}`$. For GaAs $`eh_{14}=1.2\times 10^7eV/cm`$, see, for example, Ref..
The matrix element for the spin-flip transition between the Zeeman sublevels of orbital level $`n`$ with emission of phonon $`𝐪\alpha `$ is:
$$nU_{ph}^{𝐪\alpha }n=\underset{kn}{}[\frac{(U_{ph}^{𝐪\alpha })_{nk}(\widehat{}_1)_{kn}^{}}{E_nE_kg\mu _BB}+\frac{(\widehat{}_1)_{nk}^{}(U_{ph}^{𝐪\alpha })_{kn}}{E_nE_k+g\mu _BB}],$$
(4)
where states $`n,k`$ and corresponding energies $`E_n,E_k`$ are determined by first two terms in Hamiltonian (2). The spin quantization axis coincides with the magnetic field vector. In the absence of a magnetic field the two terms in Eq.(4) cancel each other since $`(\widehat{}_1)_{nk}^{}=(\widehat{}_1)_{kn}^{}`$ and the matrix elements involving the phonon operator are symmetric with respect to the interchange of indexes $`n`$ and $`k`$. This ”Van Vleck cancellation” is a consequence of Kramers’ theorem and reduces the matrix element by a factor of $`g\mu _BB/\mathrm{}\omega _0`$, $`\mathrm{}\omega _0`$ being the typical distance between the orbital levels in the dot. Note that this cancellation occurs for a spin-orbit Hamiltonian of an arbitrary form. For instance, it could include the third order terms in the lateral momentum operator. This is in strong contrast with the cancellation of the linear in the $`\beta `$ terms in the matrix elements for the spin-flip transition between different orbital levels , which results from the fact that spin-orbit terms $`\widehat{}_1`$ are linear in the lateral momentum operators, $`\widehat{p}_{x,y}`$. Expanding in the above formula with respect to the Zeeman energy, using relation $`(\widehat{p}_i)_{nk}=(im/\mathrm{})(E_nE_k)(x_i)_{nk}`$ and the condition that the phonon wave length is much larger than the dot size (i.e. $`g\mu _BB\sqrt{ms^2\mathrm{}\omega _0}`$, $`s`$ is the sound velocity), we obtain the efective spin-flip Hamiltonian which acts on the subspace of the Zeeman sublevels of orbital level $`n`$:
$$\widehat{H}_{so}^{(n)}=g\mu _BB\frac{m\beta }{\mathrm{}e}[\widehat{\sigma }_x\alpha _{xx}^{(n)}E_x\widehat{\sigma }_y\alpha _{yy}^{(n)}E_y+\frac{(\alpha _{xy}^{(n)}+\alpha _{yx}^{(n)})}{2}(\widehat{\sigma }_xE_y\widehat{\sigma }_yE_x)],$$
(5)
where $`E_{x,y}=_{x,y}U_{ph}(x,y)/e`$ is the phonon-induced electric field in the location of the dot. Here we introduce polarizability tensor $`\widehat{\alpha }`$ that may depend on $`B_z`$. It is given by:
$$\alpha _{ik}^{(n)}(B_z)=2e^2\underset{mn}{}\frac{(x_i)_{nm}(x_k)_{mn}}{E_nE_m}$$
(6)
Effective Hamiltonian (5) is a very general one and can be used to calculate the spin-flip rates for arbitrary states and dots. To specify, we will consider only parabolic elliptic dots with the main axes along the $`x,y`$ symmetry axes, $`\omega _{x,y}`$ being the oscillator frequencies. Then the symmetry of the kinetic coefficients ensures that $`(\alpha _{xy}^{(n)}+\alpha _{yx}^{(n)})=0`$. We have to calculate the spin flip matrix elements $`<+1/2\widehat{\sigma }_{x,y}1/2>`$ over functions $`\mathrm{\Psi }_\mu ,\mu =\pm 1/2`$, which are the eigenfunctions of operator $`\widehat{\sigma }_z^{}`$, where the $`z^{}`$ axis is directed along the magnetic field vector. These functions are expressed through the eigenfunctions $`\chi _m`$ of $`\widehat{\sigma }_z`$ operator: $`\mathrm{\Psi }_\mu =_{m=\pm 1/2}D_{\mu m}^{(1/2)}(\phi ,\vartheta ,0)\chi _m`$, where $`D^{(1/2)}`$ is the finite rotations matrix and $`\phi ,\vartheta `$ are the azimuthal and polar angles presenting $`𝐁`$ in the spherical coordinates. We substitute $`E_{x,y}`$ in terms of the boson creation/annihilation operators. Then for the square modulus of the spin flip matrix element that involves emission of a phonon with wave vector $`𝐪`$ we obtain:
$`\widehat{H}_{so}^{}(𝐪\alpha )^2=\left({\displaystyle \frac{g\mu _BBm\beta }{\mathrm{}e}}\right)^2A_{𝐪\alpha }^2\left({\displaystyle \frac{\mathrm{}}{2\rho \omega _{𝐪\alpha }}}\right)\times `$ (7)
$`\{(\alpha _{xx}^2q_x^2+\alpha _{yy}^2q_y^2){\displaystyle \frac{(1+\mathrm{cos}^2\vartheta )}{2}}{\displaystyle \frac{\mathrm{sin}^2\vartheta }{2}}[(\alpha _{xx}^2q_x^2\alpha _{yy}^2q_y^2)\mathrm{cos}2\phi 2\alpha _{xx}\alpha _{yy}q_xq_y\mathrm{sin}2\phi ]\}`$ (8)
The summing up over all $`𝐪`$ yields the rate due to the first mechanism
$`\mathrm{\Gamma }_1={\displaystyle \frac{2\pi }{\mathrm{}}}{\displaystyle \frac{d^3q}{(2\pi )^3}\underset{\alpha =l,t}{}C_\alpha \widehat{H}_{so}^{}(𝐪\alpha )^2\delta (\mathrm{}s_\alpha qg\mu _BB)}=`$ (9)
$`={\displaystyle \frac{(g\mu _BB)^5}{35\pi \rho \mathrm{}^4}}\left({\displaystyle \frac{h_{14}m\beta }{e\mathrm{}}}\right)^2[(\alpha _{xx}^2+\alpha _{yy}^2)(1+\mathrm{cos}^2\vartheta )(\alpha _{xx}^2\alpha _{yy}^2)\mathrm{sin}^2\vartheta \mathrm{cos}2\phi ]\left({\displaystyle \frac{1}{s_l^5}}+{\displaystyle \frac{4}{3s_t^5}}\right).`$ (10)
Here $`C_l=1,C_t=2`$ and $`s_l,s_t`$ are the longitudinal and transverse sound velocities. The anisotropy factors used are: $`A_{𝐪,l}^2=36h_{14}^2\mathrm{cos}^2\theta \mathrm{sin}^4\theta \mathrm{sin}^2\varphi \mathrm{cos}^2\varphi `$ where $`\varphi ,\theta `$ are the azimuthal and polar angles of vector $`𝐪`$. $`<A_{𝐪,t}^2>=4h_{14}^2<(\xi _x\xi _ye_z+\xi _x\xi _ze_y+\xi _y\xi _ze_x)^2>=2h_{14}^2[\mathrm{cos}^2\theta \mathrm{sin}^2\theta +\mathrm{sin}^4\theta (19\mathrm{cos}^2\theta )\mathrm{sin}^2\varphi \mathrm{cos}^2\varphi ]`$, where $`<\mathrm{}>`$ means averaging over the orientations of the $`𝐞`$ vector in the plane which is perpendicular to $`𝐪`$. The averaging is done by the formula: $`<e_ie_k>=(1/2)(\delta _{ik}\xi _i\xi _k)`$. As usual, in the case of finite temperature Eq.(10) should be multiplied by factors $`N_\omega +1(N_\omega )`$ for the transition with emission (absorbsion) of a phonon, $`N_\omega =1/(e^{\mathrm{}\omega /T}1),\mathrm{}\omega =g\mu _BB`$. Thus, in the case of high temperature $`Tg\mu _BB`$ the spin-flip rate will be proportional to $`(g\mu _BB)^4T`$.
In the particular case of a circular dot $`\omega _x=\omega _y=\omega _0`$ we have $`\alpha _{xx}(B_z)=\alpha _{yy}(B_z)=\alpha _{xx}(0)=e^2/m\omega _0^2`$. Then, for instance, for the transition between the Zeeman sublevels of the ground state of the circular dot with emission of a piezo- phonon we obtain:
$$\mathrm{\Gamma }_1=\frac{(g\mu _BB)^5}{\mathrm{}(\mathrm{}\omega _0)^4}\mathrm{\Lambda }_p(1+\mathrm{cos}^2\vartheta );\mathrm{\Lambda }_p\frac{2}{35\pi }\frac{(eh_{14})^2\beta ^2}{\rho \mathrm{}}\left(\frac{1}{s_l^5}+\frac{4}{3s_t^5}\right).$$
(11)
The dimensionless constant $`\mathrm{\Lambda }_p`$ shows the strength of the effective spin-piezo-phonon coupling in the heterostructure and ranges from $`710^3`$ to $`610^2`$ depending on $`\beta `$. The spin-flip rate exhibits a very strong dependence on the Zeeman energy and lateral confinement energy $`\omega _0`$. To give a number, $`\mathrm{\Gamma }_11.510^3s^1`$ for $`\mathrm{}\omega _0=10K`$ and a relatively large magnetic field $`B=1T`$.
Formula (4) is written with allowance for the wave function corrections of the first order with respect to the spin-orbit Hamiltonian. The corrections of the second order are described by the following spin-orbit Hamiltonian:
$$\widehat{}_{\sigma _z}=\frac{m\beta ^2}{\mathrm{}}\widehat{\sigma }_z(x\widehat{p}_yy\widehat{p}_x)$$
(12)
Then, using this Hamiltonian in formula (4) instead of $`\widehat{}_1`$, we can get a nonzero contribution to the spin flip matrix element even with zero Zeeman splitting in the denominator (but with taking into account the orbital magnetic field). Keeping again only the term which is linear in $`\mathrm{𝐪𝐫}`$ in the expansion of exponent $`\mathrm{exp}(i\mathrm{𝐪𝐫})`$, for the rate finally we obtain:
$$\mathrm{\Gamma }_,^{(n)}=\frac{2}{35\pi }\frac{(g\mu _BB)^3(eh_{14})^2}{\mathrm{}^4\rho }\left(\frac{m\beta ^2}{\mathrm{}}\right)^2\mathrm{sin}^2\vartheta \left((D_x^{(n)})^2+(D_y^{(n)})^2\right)\left(\frac{1}{s_l^5}+\frac{4}{3s_t^5}\right),$$
(13)
where
$$D_x^{(n)}=2Re\underset{mn}{}\frac{x_{nm}(\widehat{L}_z+(eB_zr^2/2c))_{mn}}{E_nE_m},\widehat{L}_z=i\mathrm{}\left(x\frac{}{y}y\frac{}{x}\right).$$
(14)
In the absence of the magnetic field quantities $`D_x,D_y`$ are identically equal to zero. Using the properties of the matrix elements for the linear oscillators we obtain that in the case of elliptic (circular) dots $`D_x=D_y=0`$. Keeping the term which is quadratic in $`\mathrm{𝐪𝐫}`$ in the expansion of the exponent $`\mathrm{exp}(i\mathrm{𝐪𝐫})`$, we obtain a non-zero contribution but the corresponding spin-flip rate is smaller than contribution $`\mathrm{\Gamma }_1`$ by a factor of $`(\beta /s)^2(\omega _c/\omega _0)^2<1`$, here $`\omega _c=eB_z/mc`$. It is also clear that in the case of irregular dots quantities $`D_{x,y}`$ are not equal to zero. The ratio of the corresponding rate and $`\mathrm{\Gamma }_1`$ can be estimated as $`\tau ^2(m_0\beta a/\mathrm{})^2`$, where $`\tau `$ is a dimensionless parameter which describes the deviation from ellipticity and $`a`$ is a dot size. Even when $`\tau 1`$ this ratio is of the order of unity for a typical dot size $`a10^3\AA `$. Therefore, for $`\tau 1`$ we can expect that contribution Eq.(13) is much smaller than $`\mathrm{\Gamma }_1`$.
Note that, besides term $`\widehat{}_1`$ which is linear in the 2D momentum, the initial Hamiltonian also contains the term which is cubic in the momentum: $`(1/2)\widehat{\sigma }_x\{\widehat{p}_x,\widehat{p}_y^2\}_+(1/2)\widehat{\sigma }_y\{\widehat{p}_y,\widehat{p}_x^2\}_+`$. Again, in the presence of the orbital magnetic field we could get some contribution to the spin-flip rate. To this end, we need to calculate quantities $`\stackrel{~}{D}_x,\stackrel{~}{D}_y`$ obtained from $`D_x,D_y`$ by replacing operator $`\widehat{L}_z+(eB_zr^2/2c)`$ by $`(1/2)\{\widehat{p}_x,\widehat{p}_y^2\}_+`$ or $`(1/2)\{\widehat{p}_y,\widehat{p}_x^2\}_+`$. In the case of elliptic (circular) dots we obtain $`\stackrel{~}{D}_x=\stackrel{~}{D}_y=0`$ because of the symmetry.
### B Direct spin-phonon coupling.
Using the standard presentation for the strain tensor in terms of the acoustic phonon modes, we calculate the matrix element of $`\widehat{}_2`$ for the electron spin-flip transition between the Zeeman sublevels of orbital state $`\mathrm{\Phi }`$ with emission of a phonon with momentum $`𝐪`$:
$$M_,=\frac{V_0}{4}\left(\frac{\mathrm{}}{2\rho \omega _q}\right)^{1/2}[q_xe_y+q_ye_x]\mathrm{\Phi }\frac{1}{2}\{(\widehat{p}_x+i\widehat{p}_y),\mathrm{exp}(i\mathrm{𝐪𝐫})\}_+\mathrm{\Phi }.$$
(15)
For simplicity, here we set $`𝐁z`$. Similar expressions were obtained in Ref. for a different problem. The total spin-flip rate is given by the Fermi golden rule:
$$\mathrm{\Gamma }_2=\frac{\pi \mathrm{}V_0^2}{16\rho g\mu _BB}\frac{d^3q}{(2\pi )^3}(q_x^2+q_y^2)\mathrm{\Phi }\frac{1}{2}\{(\widehat{p}_x+i\widehat{p}_y),\mathrm{exp}(i\mathrm{𝐪𝐫})\}_+\mathrm{\Phi }^2\delta (\mathrm{}sqg\mu _BB).$$
(16)
The relevant phonon wave length is much larger than the dot size, which allows for further simplifications. We concentrate on a circular dot with confining frequency $`\omega _0`$. For the orbital states with $`n=0`$ and $`l=0,\pm 1`$ (the ground and the first two excited states):
$$\mathrm{\Gamma }_2=\frac{V_0^2(g\mu _BB)^5}{240\pi \rho s^7\mathrm{}^4}[l+\frac{\omega _c}{2\sqrt{\omega _0^2+(\omega _c^2/4)}}(l+1)!]^2.$$
(17)
The spin-flip rate produced by term $`\widehat{}_3`$ does not depend on the structure of the orbital state and is given by
$$\mathrm{\Gamma }_3\frac{(\mu _B\stackrel{~}{g}B)^2(\mu _BgB)^3}{\rho s_t^5\mathrm{}^4}$$
(18)
Let us now compare the rates $`\mathrm{\Gamma }_{1,2,3}`$ obtained. All of them are proportional to the fifth power of energy splitting $`g\mu _BB`$, so that their ratio hardly depends on the magnetic field. First, we note that the ratio of $`\mathrm{\Gamma }_3`$ and $`\mathrm{\Gamma }_2`$ is of the order of $`(\stackrel{~}{g}s_t/gV_0)^27.810^31`$. So that $`\mathrm{\Gamma }_2`$ is more important. The ratio of $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_2`$ is of the order of $`(eh_{14}/mV_0/\mathrm{})^2(m\beta ^2ms_t^2/(\mathrm{}\omega _0)^4)`$. For $`\mathrm{}\omega _01÷10K`$ the ratio is of the order of $`10^6÷10^2`$ and increases only for larger dots that have smaller $`\omega _0`$. Thus, we conclude that the admixture mechanism dominates.
### C Two-phonon processes.
The calculated rate $`\mathrm{\Gamma }_1`$ is small partly because of the small phonon density of the states at the scale of the Zeeman energy. On the other hand, for the case of the spin-flip transitions between the Zeeman levels of usual impurity the two phonon processes under some conditions may become more important than the single phonon processes. At sufficiently small Zeeman splitting the contribution of the single phonon processes is very small, and with increasing temperature the role of the processes when one phonon is absorbed and the other is emitted is increased. It is also true for the case of a quantum dot and here we give some formulas which describe the contribution of such two phonon processes for GaAs quantum dots in several limiting cases. We also indicate the conditions under which these contributions can be important.
If we treat the interaction with the phonons in the second order, we obtain processes in which a phonon is scattered from state $`𝐩`$ to state $`𝐪`$ while the electron spin flips. The effective matrix element contains transitions to an excited orbital state with the emission or absorption of a phonon and then transitions back to the ground state with the absorption or emission of a phonon. The spin may flip either in the first or second transition. The matrix element is :
$$<V_2>\left(\frac{\mathrm{}}{\rho s\sqrt{pq}}\right)(eh_{14})^2[N_p(N_q+1)]^{1/2}\underset{a}{}\left\{\frac{[H_{p,q}^++H_{p,q}^{}]}{\mathrm{\Delta }_a\mathrm{}sq}+\frac{[H_{q,p}^++H_{q,p}^{}]}{\mathrm{\Delta }_a+\mathrm{}sp}\right\},$$
(19)
$$H_{p,q}^\pm =(\mathrm{\Psi }_0^+,\mathrm{exp}(i\mathrm{𝐩𝐫})\mathrm{\Psi }_a^\pm )(\mathrm{\Psi }_a^\pm ,\mathrm{exp}(i\mathrm{𝐪𝐫})\mathrm{\Psi }_0^{}),$$
where $`N_p`$ is the Bose distribution function and $`\mathrm{\Delta }_a`$ is the energy separation between the ground state whose wave function is $`\mathrm{\Psi }_0`$ and the excited state whose wave function for spin up, say, is $`\mathrm{\Psi }_a^+`$. We can neglect the Zeeman energy in the denominators since no Van Vleck cancellation occurs here. We consider again interaction with piezo-phonons since deformation phonons become important at very high temperature (see below). For simplicity we consider here only the case when the magnetic field is perpendicular to the 2D plane and study the relaxation of $`S_z`$ spin component.
As it was shown in Ref., there is a cancellation of the linear in the $`\beta `$ terms in the matrix elements of type $`(\mathrm{\Psi }_0^+,\mathrm{exp}(i\mathrm{𝐩𝐫})\mathrm{\Psi }_a^{})`$ for the spin-flip transition between different orbital levels. This is a consequence of the fact that spin-orbit terms $`\widehat{}_1`$ are linear in the lateral momentum operators, $`\widehat{p}_{x,y}`$. For that reason, quantities $`H_{p,q}^\pm `$ are proportional to the first power of $`\beta `$ only if one takes into account the Zeeman splitting in the electron spectrum. We consider the temperature in the interval $`g\mu _BBT\mathrm{}\omega _0`$, where $`\mathrm{}\omega _0`$ is the characteristic energy distance between the levels in the dot. Since $`\mathrm{}sp\mathrm{}sq<T`$, then we can neglect the phonon energies in the denominators while calculating the contribution to $`<V_2>`$ proportional to the first power of $`\beta `$. It is apparent that the spin-flip rate has a different temperature dependence for the temperatures smaller and larger than $`T_0\sqrt{ms^2\mathrm{}\omega _0}`$. At this characteristic temperature the phonon wave length is equal to lateral dot size $`\lambda `$. For GaAs at $`\mathrm{}\omega _010K`$ temperature $`T_01K`$. Let us give the estimate for the spin-flip rate in the case $`TT_0`$, when $`p\lambda ,q\lambda 1`$. Here the estimate for $`H_{p,q}^\pm `$ is $`H_{p,q}^\pm (\beta /\lambda \omega _0)(g\mu _BB/\mathrm{}\omega _0)(\lambda q)^3`$. Then the relaxation rate is:
$$\mathrm{\Gamma }_1^{(2)}(T)=(2\pi /\mathrm{})\underset{p,q}{}<V_2>^2\delta [\mathrm{}s(pq)g\mu _BB]\frac{\mathrm{\Lambda }_p^2}{\mathrm{}}\frac{s^2}{\beta ^2}\frac{(g\mu _BB)^2(ms^2)^{5/2}}{(\mathrm{}\omega _0)^{7/2}}\left(\frac{T}{T_0}\right)^9$$
(20)
In performing the integral over $`𝐩`$, we have neglected $`g\mu _BB`$ in comparison to $`\mathrm{}sp`$.
In the case $`TT_0`$ the momentum components parallel to the plane are estimated as $`q_{}\lambda 1`$, otherwise the matrix elements are exponentially small. As to the $`q_z`$ values, the contribution of the region $`q_zT/\mathrm{}sq_{}`$ is much smaller than that where $`q_zq_{}T/\mathrm{}s`$. Such is the case even without regard for the fact that for the orientation of the 2D plane the effective piezo-modulus $`A_{𝐪\alpha }`$ introduced above has a smallness $`q_{}/q_z1`$. Thus, calculating the contribution from $`q_zq_{}T/\mathrm{}s`$ and taking into account that $`N_p=T/\mathrm{}sp1`$, we obtain for the spin-flip rate in the case $`T_0T\mathrm{}\omega _0`$:
$$\mathrm{\Gamma }_2^{(2)}(T)\frac{\mathrm{\Lambda }_p^2}{\mathrm{}}\frac{s^2}{\beta ^2}\frac{(g\mu _BB)^2(ms^2)^{5/2}}{(\mathrm{}\omega _0)^{7/2}}\left(\frac{T}{T_0}\right)^2$$
(21)
The contribution from the deformation phonons is much smaller. In the case $`TT_0`$ the characteristic $`q_zT/\mathrm{}sq_{}1/\lambda `$, i.e. the deformation phonons are emitted almost perpendicular to the 2D plane. Then, for the deformation potential contribution to the spin-flip rate we obtain: $`\mathrm{\Gamma }_d(\mathrm{\Lambda }_d^2/\mathrm{})(\beta ^2/s^2)(g\mu _BB)^2T^3/(\mathrm{}\omega _0)^3ms^2`$, where $`\mathrm{\Lambda }_d(1/2\pi )(\mathrm{\Sigma }^2m^2/\rho \mathrm{}^3s)`$ is the dimensionless constant which shows the strength of the electron interaction with deformation phonons. For GaAs $`\mathrm{\Lambda }_d10^5`$. Even at $`T\mathrm{}\omega _0`$ the value of $`\mathrm{\Gamma }_d/\mathrm{\Gamma }_2^{(2)}(\mathrm{}\omega _0/ms^2)^{5/2}(\mathrm{\Lambda }_d\beta ^2/\mathrm{\Lambda }_ps^2)^2`$ is much smaller than unity for any realistic $`\mathrm{}\omega _0`$. For example, at $`\mathrm{}\omega _0=30K`$ this ratio is $`0.03`$.
Let us compare the two-phonon contribution $`\mathrm{\Gamma }_2^{(2)}`$ with $`\mathrm{\Gamma }_1T/g\mu _BB`$. We see that the two-phonon contribution $`\mathrm{\Gamma }_2^{(2)}`$ prevails at sufficiently small Zeeman splittings: $`g\mu _BB<ms^2(\mathrm{\Lambda }_ps^2/\beta ^2)^{1/2}(T/T_0)^{1/2}`$. Taking the maximal temperature $`T\mathrm{}\omega _0`$, we obtain $`g\mu _BB<[(\mathrm{\Lambda }_ps^2/\beta ^2)\sqrt{ms^2/\mathrm{}\omega _0}]^{1/2}T_0`$. For $`\mathrm{}\omega _010K`$ we see that this contribution is more important for magnetic fields smaller than approximately 0.4 T (where the estimate for the spin-flip time is of the order of ms). On the other hand, at $`TT_0`$ we obtain for the same $`\mathrm{}\omega _0`$ that $`g\mu _BB<0.03K`$ (i.e. the two-phonon contribution is more important for magnetic fields smaller than $`1kG`$). For these fields the characteristic spin-flip time is of the order of 1s, i.e. it is still long.
The general conclusion is that at sufficiently low temperatures (much smaller than $`\mathrm{}\omega _0`$) the characteristic Zeeman splittings below which the two-phonon contribution to the spin-flip rate dominates are small and corresponding spin-flip times are unusually long (see the estimates above).
## III Other mechanisms of the spin flip.
Let us discuss briefly other mechanisms of the spin flip. The spin transitions between the Zeeman sublevels of the impurity state in semiconductors ( mostly Si) were extensively studied quite a long time ago. Except spin-orbit coupling, several other mechanisms were proposed, such as: 1) modulation of the hyperfine coupling with nuclei by lattice vibrations, 2) the spin-spin interaction between the bound electron and the conduction electron in the leads, 3) the spin-current interaction, when the bound electron spin flip is caused by the fluctuating magnetic field of the conduction electrons, 4) an exchange scattering process which flips the spins of both the conduction electron and the bound donor electron. Whereas the spin-orbit interaction strongly depends on the crystal symmetry and is different for Si and GaAs, the other mechanisms are quite general in nature and we can profit from the discussion in .
Mechanism (4) requires an overlap of the wave functions of the electrons in the leads and in the dot. In the context of QD it is considered in Ref. . The corresponding rates are not intrinsic to the dot since they are proportional to the barrier transparencies. They can be tuned to arbitrary low values. Refs.() have demonstrated that the spin-flip rates associated with mechanisms (2,3) are very small. As an example, we give the rate estimation for mechanism (3). The Bio-Savar formula relates the magnetic field and current fluctuations in the leads so that $`<H^2>_\omega (1/c^2a^2)<I^2>_\omega `$, $`a`$ being the characteristic distance between the electron in the dot and the electrons in the leads. Using the Nyquist formula for the correlator of the currents we estimate $`<H^2>_\omega (1/c^2a^2)\mathrm{}\omega \mathrm{coth}(\mathrm{}\omega /2T)(1/R)`$, $`R`$ being the typical resistance of the leads or the dot environment. Thus, the corresponding spin-flip relaxation rate is estimated as
$$\mathrm{\Gamma }_4\mu _B^2<H^2>_\omega /\mathrm{}^2\omega \left(\frac{\lambda _c}{a}\right)^2\frac{\mathrm{}}{e^2R}$$
(22)
where $`\lambda _c=e^2/m_0c^22.810^{13}`$ cm is the classical electron radius, $`\mathrm{}\omega =g\mu _BB`$. Rate $`\mathrm{\Gamma }_4`$ is proportional to the first power of Zeeman splitting so that it may formally compete with $`\mathrm{\Gamma }_1`$ at sufficiently small splittings. However, this occurs at splittings that are so small that the corresponding rates are not observable. To give an example, we choose $`R=1Ohm`$ and $`\mathrm{}\omega _0=1K`$, which corresponds to $`a1.1510^5cm`$. Then rate $`\mathrm{\Gamma }_4`$ would dominate if splitting $`\mathrm{}\omega 2.510^3K`$. This corresponds to the rates lower than $`810^4s^1`$!
As to mechanism (1), i.e. modulation of the hyperfine coupling with nuclei by lattice vibrations, the relative strength of this mechanism and the spin-orbit interaction can be different for different materials. For example, in the case of Si where the spin-orbit interaction is much weaker than in GaAs, the dominant mechanism of the spin-flip for the case of the phonon assisted transitions between the Zeeman levels of usual impurity (the situation studied in Ref. ) was found to be the modulation of the hyperfine coupling with nuclei by lattice vibrations. In the case of GaAs, however, our conclusion is that the dominant mechanism is the admixture mechanism of the spin-orbit interaction. This conclusion was reached for the first time in Ref., where the calculations used essentially followed those in Ref.. Here we give the result obtained in Ref. for the rate due to the modulation of the hyperfine coupling with nuclei by lattice vibrations
$$\mathrm{\Gamma }_h(g\mu _BB)^3\gamma ^2\omega _N^2/\mathrm{}^2s^5\rho ,$$
(23)
where $`\omega _N(v_0A^2/a^2z_0)^{1/2}/\mathrm{}`$ is the electron spin precession frequency in the random field of unpolarized nuclei, $`v_0`$ is the unit cell volume, $`A`$ the hyperfine interaction constant, $`a`$ the dot lateral size and $`z_0`$ the electron wave function extension in the z-direction. Finally, $`\gamma (1/m)(dm/d\mathrm{\Delta })`$ is the change in effective mass $`m`$ with dilation, see also Ref. . For GaAs QD with $`a10^3\AA `$ and $`z_010^2\AA `$, $`\omega _N`$ can reach $`10^8s^1`$. Let us compare spin-flip rate Eq.(23) with $`\mathrm{\Gamma }_1`$ (this comparison was done earlier in Ref.). Even taken for $`\gamma 50`$ (see also ) we can easily see that $`\mathrm{\Gamma }_h`$ will compete with $`\mathrm{\Gamma }_1`$ at the Zeeman splitting $`\gamma \omega _N(\mathrm{}\omega _0)^2/\beta eh_{14}`$ which is so small that the corresponding rate is not observable. For example, for $`\mathrm{}\omega _0=10K`$ the splitting is of the order of $`10^5`$ K. Therefore, the admixture mechanism of the spin-orbit interaction is the dominant one.
It should be noted, that in this work we have not considered the electron spin relaxation mechanism which is through the hyperfine interaction related to the internal nuclear dynamics. The latter is due to the dipole-dipole interaction between the nuclei which does not conserve the total nuclear spin . This mechanism might be important at low magnetic fields. However, this problem is not simple and needs a seperate investigation.
Finally, we mention the experimental studies of spin relaxation in n-type GaAs quantum dots. Such an experiment has been recently carried out . The non-equilibrium tunneling current through excited states in an AlGaAs/GaAs quantum dot was studied using a pulse-excitation technique which measures the energy relaxation time from the excited state to the ground state. Some excited states showed a relaxation time which was much longer than a few $`\mu s`$, while the other showed time much shorter than a few ns. This great difference in relaxation times was ascribed to the fact that some inelastic transitions are accompanied by the spin flip. For these transitions the relaxation time was so long that the method used in the above mentioned paper only allowed to give some estimation (much longer than a few $`\mu s`$). Though the transitions studied by T.Fujisawa et al. could in general involve the spin flip transitions between the states with different orbital structures (this situation was considered in our previous paper ), the experimental data confirm the general statement that the spin-flip processes in n-type quantum dots can be really slow.
In conclusion, we have calculated the rates for the phonon-assisted spin-flip transitions between the Zeeman sublevels in a quantum dot for all possible mechanisms and shown that the admixture mechanism of the spin-orbit interaction is a dominant one. The corresponding spin-flip rate $`\mathrm{\Gamma }_1`$ (see Eqs.(10,11)) exhibits a strong dependence on Zeeman energy and at small magnetic fields takes very low values (up to seconds).
This work is part of the research program of the ”Stichting voor Fundamenteel Onderzoek der Materie (FOM)”. We acknowledge the support of the Netherlands Organization for Scientific Research (NWO) in the framework of Dutch-Russian collaboration and the NEDO project NTDP-98. We are grateful to L.P. Kouwenhoven, T.H. Oosterkamp, G.E.W. Bauer, T. Fujisawa, Y. Tokura, Y. Hirayama and D. Loss for useful discussions.
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# Lepton polarization and CP–violating effects in 𝐵→𝐾^∗𝜏⁺𝜏⁻ decay in standard and two Higgs doublet models
## 1 Introduction
Rare $`B`$ meson decays, induced by flavor–changing neutral current (FCNC) $`bs`$ transitions, is one of the most promising research area in high energy physics. Theoretical interest to the rare $`B`$ decays lies in their role as a potential precision testing ground for the standard model (SM) at loop level and experimentally these decays will provide quantitative information about the Cabibbo–Kobayashi–Maskawa (CKM) matrix elements $`V_{td},V_{ts}`$ and $`V_{tb}`$. Besides, these rare decays have the potential of establishing new physics beyond SM, such as two Higgs doublet model (2HDM), minimal supersymmetric extension of the SM (MSSM), left–right models, etc.
At present, the main interest is focused on the rare $`B`$ meson decays, for which the SM predicts ”large” branching ratios and that can be potentially measurable at working $`B`$ factories and LHC. The rare $`BK^{}\mathrm{}^+\mathrm{}^{}`$ $`(\mathrm{}=e,\mu ,\tau )`$ decays are such ones. At the same time this decay constitutes a suitable tool in looking for new physics beyond the SM. At quark level the process is described by the $`bs\mathrm{}^+\mathrm{}^{}`$ transition. This transition in framework of the SM and its various extensions have been extensively investigated .
One efficient way in establishing new physics is the measurement of lepton polarization. This problem is widely discussed in literature for the $`bs\mathrm{}^+\mathrm{}^{}`$ decay . Note that all previous studies for the lepton polarization, except the work , have been limited to SM and its minimal extensions. In the analysis of the $`\tau `$ lepton polarization for the $`bs\tau ^+\tau ^{}`$ decay was performed in a model independent way. In this work twelve (ten local and two nonlocal) four–Fermi operator interactions were introduced in a model independent way instead of three independent structures which are present in SM.
It is well known that the theoretical analysis of the inclusive decays is rather easy but their experimental detection is quite difficult. However for exclusive decays the situation is contrary to this case, i.e., their experimental study is easy but theoretical investigation is difficult. This is due to the fact that the description of the exclusive decays requires form factors, i.e., the matrix elements of the effective Hamiltonian between initial $`B`$ and final meson states. This problem is related to to the nonperturbative sector of the QCD and it can be solved only by means of a nonperturbative approach.
These matrix elements have been investigated in framework of different approaches such as chiral theory , three point QCD sum rules method , relativistic quark model by the light–front formalism , effective heavy quark theory and light cone QCD sum rules .
The aim of the present paper is to perform a comprehensive study the lepton polarizations and CP–violating asymmetry in the exclusive $`BK^{}\mathrm{}^+\mathrm{}^{}`$ $`(\mathrm{}=\mu ,\tau )`$ decay in the SM and three versions of the 2HDM. Note that this decay in framework of 2HDM (model I and model II) were investigated in (second reference), where $`P_L`$ and $`P_T`$ and $`A_{FB}`$ were studied. Here in this work we extend these considerations by studying in model III, including normal polarization of $`\tau `$ lepton and CP–violating asymmetry, using the current limits for the parameters of 2HDM coming from low energy experiments like $`B^0`$$`\overline{B}^0`$ mixing, $`\rho `$–parameter analysis and $`bs\gamma `$ decay.
The paper is organized as follows. In section 2, starting from a general form of four–Fermi interactions we derive model independent expressions for the longitudinal, transversal and normal polarizations. In section 3, we apply the above–mentioned general results of the lepton polarizations to SM and to three types of 2HDM (so called models I, II and III). A brief summary of our results is presented in this section.
## 2 Lepton polarizations
In this section we present the expressions for the longitudinal, transversal and normal polarizations of the $`\tau `$ lepton in a model independent way. For this aim we follow where the matrix element for $`bs\tau ^+\tau ^{}`$ transition is given in terms of twelve most general independent four–Fermi interactions:
$``$ $`=`$ $`{\displaystyle \frac{G\alpha }{\sqrt{2}\pi }}V_{tb}V_{ts}^{}\{C_{SL}\overline{s}i\sigma _{\mu \nu }{\displaystyle \frac{q^\nu }{q^2}}(m_sL)b\overline{\mathrm{}}\gamma _\mu \mathrm{}+C_{BR}\overline{s}i\sigma _{\mu \nu }{\displaystyle \frac{q^\nu }{q^2}}(m_bR)b\overline{\mathrm{}}\gamma _\mu \mathrm{}`$ (1)
$`+`$ $`C_{LL}\overline{s}_L\gamma _\mu b_L\overline{\mathrm{}}_L\gamma ^\mu \mathrm{}_L+C_{LR}\overline{s}_L\gamma _\mu b_L\overline{\mathrm{}}_R\gamma ^\mu \mathrm{}_R+C_{RL}\overline{s}_R\gamma _\mu b_R\overline{\mathrm{}}_L\gamma ^\mu \mathrm{}_L`$
$`+`$ $`C_{RR}\overline{s}_R\gamma _\mu b_R\overline{\mathrm{}}_R\gamma ^\mu \mathrm{}_R+C_{LRLR}\overline{s}_Lb_R\overline{\mathrm{}}_L\mathrm{}_R+C_{RLLR}\overline{s}_Rb_L\overline{\mathrm{}}_L\mathrm{}_R`$
$`+`$ $`C_{LRRL}\overline{s}_Lb_R\overline{\mathrm{}}_R\mathrm{}_L+C_{RLRL}\overline{s}_Rb_L\overline{\mathrm{}}_R\mathrm{}_L+C_T\overline{s}\sigma _{\mu \nu }b\overline{\mathrm{}}\sigma ^{\mu \nu }\mathrm{}`$
$`+`$ $`iC_{TE}ϵ^{\mu \nu \alpha \beta }\overline{s}\sigma _{\mu \nu }b\overline{\mathrm{}}\sigma _{\alpha \beta }\mathrm{}\},`$
where $`C_{XX}`$ are the coefficient of the four–Fermi interactions. Among them, there are two non–local Fermi interactions denoted by $`C_{SL}`$ and $`C_{BR}`$, which correspond to $`2C_7^{eff}`$ in the SM. Two ($`C_{LL}`$$`C_{LR}`$) of the four vector type interactions ($`C_{LL}`$, $`C_{LR}`$, $`C_{RL}`$ and $`C_{RR}`$) are also present in the SM (in the forms of $`C_{LL}=C_9^{eff}C_{10}`$ and $`C_{LR}=C_9^{eff}+C_{10}`$). $`C_{LRLR}`$, $`C_{RLLR}`$, $`C_{LRRL}`$ and $`C_{RLRL}`$ are the scalar type interactions and the last two terms with coefficients $`C_T`$ and $`C_{TE}`$ are the tensor type interactions.
For simplicity we will take the mass of the strange quark to be zero and neglect tensor type interactions, since it was indicated in that the physical observables are not sensitive to the presence of the tensor type interactions. Having established the matrix element for the $`bs\tau ^+\tau ^{}`$ transition, our next problem is calculation of the matrix elements, $`K^{}\left|\overline{s}\gamma _\mu (1\pm \gamma _5)b\right|B`$, $`K^{}\left|\overline{s}i\sigma _{\mu \nu }q^\nu (1+\gamma _5)b\right|B`$ and $`K^{}\left|\overline{s}(1\pm \gamma _5)b\right|B`$, in order to be able to calculate the physically measurable quantities at hadronic level. These matrix elements can be written in terms of the form factors in the following way:
$`K^{}(p_K^{},\epsilon )\left|\overline{s}\gamma _\mu (1\pm \gamma _5)b\right|B(p_B)=`$
$`ϵ_{\mu \nu \rho \sigma }\epsilon ^\nu p_K^{}^\rho q^\sigma {\displaystyle \frac{2V(q^2)}{m_B+m_K^{}}}\pm i\epsilon _\mu ^{}(m_B+m_K^{})A_1(q^2)i(p_B+p_K^{})_\mu (\epsilon ^{}q){\displaystyle \frac{A_2(q^2)}{m_B+m_K^{}}}`$
$`iq_\mu {\displaystyle \frac{2m_K^{}}{q^2}}(\epsilon ^{}q)\left[A_3(q^2)A_0(q^2)\right],`$
$`K^{}(p_K^{},\epsilon )\left|\overline{s}i\sigma _{\mu \nu }q^\nu (1+\gamma _5)b\right|B(p_B)=`$ (3)
$`4ϵ_{\mu \nu \rho \sigma }\epsilon ^\nu p_K^{}^\rho q^\sigma T_1(q^2)+2i\left[\epsilon _\mu ^{}(m_B^2m_K^{}^2)(p_B+p_K^{})_\mu (\epsilon ^{}q)\right]T_2(q^2)`$
$`+2i(\epsilon ^{}q)\left[q_\mu (p_B+p_K^{})_\mu {\displaystyle \frac{q^2}{m_B^2m_K^{}^2}}\right]T_3(q^2),`$
where $`\epsilon `$ is the polarization vector of $`K^{}`$ meson and $`q=p_Bp_K^{}`$ is the momentum transfer. In order to ensure finiteness of (2) at $`q^2=0`$, we assume that $`A_3(q^2=0)=A_0(q^2=0)`$. To calculate the matrix element $`K^{}\left|\overline{s}(1\pm \gamma _5)b\right|B`$, we multiply both sides of Eq. (2) by $`q_\mu `$ and use the equation of motion. Neglecting the mass of the strange quark, we get
$`K^{}(p_K^{},\epsilon )\left|\overline{s}(1\pm \gamma _5)b\right|B(p_B)={\displaystyle \frac{1}{m_b}}\{i(\epsilon ^{}q)(m_B+m_K^{})A_1(q^2)`$ (4)
$`\pm i(m_Bm_K^{})(\epsilon ^{}q)A_2(q^2)\pm 2im_K^{}(\epsilon ^{}q)[A_3(q^2)A_0(q^2)]\}.`$
Using the equation of motion, the form factor $`A_3`$ can be expressed as a linear combination of the form factors $`A_1`$ and $`A_2`$ (see )
$`A_3(q^2)={\displaystyle \frac{m_B+m_K^{}}{2m_K^{}}}A_1(q^2){\displaystyle \frac{m_Bm_K^{}}{2m_K^{}}}A_2(q^2).`$ (5)
Using this relation we obtain
$`K^{}(p_K^{},\epsilon )\left|\overline{s}(1\pm \gamma _5)b\right|B(p_B)={\displaystyle \frac{1}{m_b}}\left\{2im_K^{}(\epsilon ^{}q)A_0(q^2)\right\}.`$ (6)
From Eqs. (1), (2), (3) and (6) we get the following expression for the matrix element of the $`BK^{}\mathrm{}^+\mathrm{}^{}`$ decay
$`={\displaystyle \frac{G\alpha }{4\sqrt{2}\pi }}V_{tb}V_{ts}^{}`$ (7)
$`\times \{\overline{\mathrm{}}\gamma _\mu (1\gamma _5)\mathrm{}[ϵ_{\mu \nu \rho \sigma }\epsilon ^\nu p_K^{}^\rho q^\sigma (2A_1)i\epsilon _\mu ^{}B_1+i(\epsilon ^{}q)(p_B+p_K^{})_\mu B_2+iq_\mu (\epsilon ^{}q)B_3]`$
$`+\overline{\mathrm{}}\gamma _\mu (1+\gamma _5)\mathrm{}\left[ϵ_{\mu \nu \rho \sigma }\epsilon ^\nu p_K^{}^\rho q^\sigma (2C_1)i\epsilon _\mu ^{}D_1+i(\epsilon ^{}q)(p_B+p_K^{})_\mu D_2+iq_\mu (\epsilon ^{}q)D_3\right]`$
$`+\overline{\mathrm{}}(1\gamma _5)\mathrm{}i(\epsilon ^{}q)B_4+\overline{\mathrm{}}(1+\gamma _5)\mathrm{}i(\epsilon ^{}q)B_5\},`$
where
$`A_1`$ $`=`$ $`(C_{LL}+C_{RL}){\displaystyle \frac{V(q^2)}{m_B+m_K^{}}}2C_{BR}{\displaystyle \frac{m_b}{q^2}}T_1,`$
$`B_1`$ $`=`$ $`(C_{LL}C_{RL})(m_B+m_K^{})A_12C_{BR}{\displaystyle \frac{m_b}{q^2}}(m_B^2m_K^{}^2)T_2,`$
$`B_2`$ $`=`$ $`{\displaystyle \frac{C_{LL}C_{RL}}{m_B+m_K^{}}}A_22C_{BR}{\displaystyle \frac{m_b}{q^2}}\left[T_2+{\displaystyle \frac{q^2}{m_B^2m_K^{}^2}}T_3\right],`$
$`B_3`$ $`=`$ $`(C_{LL}C_{RL}){\displaystyle \frac{2m_K^{}}{q^2}}(A_3A_0)2C_{BR}{\displaystyle \frac{m_b}{q^2}}T_3,`$
$`C_1`$ $`=`$ $`A_1(C_{LL}C_{LR},C_{RL}C_{RR}),`$
$`D_1`$ $`=`$ $`B_1(C_{LL}C_{LR},C_{RL}C_{RR}),`$
$`D_2`$ $`=`$ $`B_2(C_{LL}C_{LR},C_{RL}C_{RR}),`$
$`D_3`$ $`=`$ $`B_3(C_{LL}C_{LR},C_{RL}C_{RR}),`$
$`B_4`$ $`=`$ $`(C_{LRRL}C_{RLRL})\left({\displaystyle \frac{2m_K^{}}{m_b}}A_0\right),`$
$`B_4`$ $`=`$ $`(C_{LRLR}C_{RLLR})\left({\displaystyle \frac{2m_K^{}}{m_b}}A_0\right),`$
At this point we would like to make the following comment. The main difference from the SM case is, we have six different structures (after setting $`m_s=0`$ and neglecting tensor interactions) in the inclusive channel, while under the same conditions we have only two new structures, namely scalar type interactions proportional to $`B_4`$ and $`B_5`$. On the other hand there appears no any new structure in the 2HDM.
Having established this matrix element, let us now consider the final lepton polarization. We define the following three orthogonal unit vectors:
$`\stackrel{}{e}_L`$ $`=`$ $`{\displaystyle \frac{\stackrel{}{p}_1}{|\stackrel{}{p}_1|}},`$
$`\stackrel{}{e}_N`$ $`=`$ $`{\displaystyle \frac{\stackrel{}{p}_K^{}\times \stackrel{}{p}_1}{|\stackrel{}{p}_K^{}\times \stackrel{}{p}_1|}},`$
$`\stackrel{}{e}_T`$ $`=`$ $`\stackrel{}{e}_L\times \stackrel{}{e}_N,`$ (9)
where $`\stackrel{}{p}_1`$ and $`\stackrel{}{p}_K^{}`$ are the three–momenta of the lepton $`\mathrm{}^{}`$ and $`K^{}`$ meson, respectively, in the center of mass of final leptons. The differential decay rate for any spin direction $`\stackrel{}{n}`$ of the $`\mathrm{}^{}`$ lepton, where $`\stackrel{}{n}`$ is a unit vector in the $`\mathrm{}^{}`$ rest frame, can be expressed in the following form
$`{\displaystyle \frac{d\mathrm{\Gamma }(\stackrel{}{n})}{dq^2}}={\displaystyle \frac{1}{2}}\left({\displaystyle \frac{d\mathrm{\Gamma }}{dq^2}}\right)_0\left[1+\left(P_L\stackrel{}{e}_L+P_N\stackrel{}{e}_N+P_T\stackrel{}{e}_T\right)\stackrel{}{n}\right],`$ (10)
where the subscript ”0” corresponds to the unpolarized differential decay rate whose explicit form will be presented below. $`P_L`$, $`P_N`$ and $`P_T`$ are recognized as the longitudinal, normal and transversal polarizations, respectively. It follows from the definition of unit vectors $`\stackrel{}{e}_i`$ that $`P_T`$ obviously lies in the decay plane whose orientation is determined by the vectors $`\stackrel{}{p}_1`$ and $`\stackrel{}{p}_K^{}`$ and $`P_N`$ is perpendicular to this plane.
The expression for the unpolarized differential decay rate in Eq. (10) can be written as:
$`\left({\displaystyle \frac{d\mathrm{\Gamma }}{dq^2}}\right)_0`$ $`=`$ $`{\displaystyle \frac{G^2\alpha ^2}{2^{14}\pi ^5m_B}}\left|V_{tb}V_{ts}^{}\right|^2\lambda ^{1/2}v`$ (11)
$`\times `$ $`\{32\lambda m_B^4[{\displaystyle \frac{1}{3}}(m_B^2sm_{\mathrm{}}^2)(|A_1|^2+|C_1|^2)+2m_{\mathrm{}}^2\text{Re}(A_1C_1^{})]`$
$`+`$ $`96m_{\mathrm{}}^2\text{Re}(B_1D_1^{}){\displaystyle \frac{4}{r}}m_B^2m_{\mathrm{}}\lambda \text{Re}[(B_1D_1)(B_4^{}B_5^{})]`$
$`+`$ $`{\displaystyle \frac{8}{r}}m_B^2m_{\mathrm{}}^2\lambda \{\text{Re}[(B_3^{}+D_2^{}D_3^{})B_1]+\text{Re}[(B_2^{}B_3^{}+D_3^{})D_1]\text{Re}(B_4B_5^{})]\}`$
$`+`$ $`{\displaystyle \frac{4}{r}}m_B^4m_{\mathrm{}}(1r)\lambda \left\{\text{Re}[(B_2D_2)(B_4^{}B_5^{})]\right\}`$
$`+`$ $`{\displaystyle \frac{8}{r}}m_B^4m_{\mathrm{}}^2(1r)\lambda \left\{\text{Re}[(B_2D_2)(B_3^{}D_3^{})]\right\}`$
$``$ $`{\displaystyle \frac{8}{r}}m_B^4m_{\mathrm{}}^2\lambda (2+2rs)\text{Re}(B_2D_2^{}){\displaystyle \frac{4}{r}}m_B^4m_{\mathrm{}}s\lambda \text{Re}[(B_3D_3)(B_4^{}B_5^{})]`$
$``$ $`{\displaystyle \frac{4}{r}}m_B^4m_{\mathrm{}}^2s\lambda \left[\left|B_3\right|^2+\left|D_3\right|^22\text{Re}(B_3D_3^{})\right]+{\displaystyle \frac{2}{r}}m_B^2(m_B^22m_{\mathrm{}}^2)\lambda \left[\left|B_4\right|^2+\left|B_5\right|^2\right]`$
$``$ $`{\displaystyle \frac{8}{3rs}}m_B^2\lambda \left[m_{\mathrm{}}^2(22r+s)+m_B^2s(1rs)\right]\left[\text{Re}(B_1B_2^{})+\text{Re}(D_1D_2^{})\right]`$
$`+`$ $`{\displaystyle \frac{4}{rs}}\left[2m_{\mathrm{}}^2(\lambda 6rs)+m_B^2s(\lambda +12rs)\right]\left[\left|B_1\right|^2+\left|D_1\right|^2\right]`$
$`+`$ $`{\displaystyle \frac{4}{3rs}}m_B^4\lambda \{m_B^2s\lambda +m_{\mathrm{}}^2[2\lambda +3s(2+2rs)]\}\left[\right|B_2|^2+\left|D_2|^2\right]\}.`$
The polarizations $`P_L`$, $`P_N`$ and $`P_T`$ are defined as:
$`P_i(q^2)={\displaystyle \frac{{\displaystyle \frac{d\mathrm{\Gamma }}{dq^2}}(\stackrel{}{n}=\stackrel{}{e}_i){\displaystyle \frac{d\mathrm{\Gamma }}{dq^2}}(\stackrel{}{n}=\stackrel{}{e}_i)}{{\displaystyle \frac{d\mathrm{\Gamma }}{dq^2}}(\stackrel{}{n}=\stackrel{}{e}_i)+{\displaystyle \frac{d\mathrm{\Gamma }}{dq^2}}(\stackrel{}{n}=\stackrel{}{e}_i)}}.`$ (12)
After lengthy calculations we get the following general expressions for the longitudinal, transversal and normal polarizations of the $`\mathrm{}^{}`$ lepton (for $`m_s=0`$ and neglecting the tensor interaction)
$`P_L`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{\Delta }}}v\{{\displaystyle \frac{4}{3r}}\lambda ^2m_B^6\left[\right|B_2|^2\left|D_2|^2\right]+{\displaystyle \frac{4}{r}}\lambda m_B^2m_{\mathrm{}}\text{Re}[(B_1D_1)(B_4^{}+B_5^{})]`$
$``$ $`{\displaystyle \frac{4}{r}}\lambda m_B^4m_{\mathrm{}}(1r)\text{Re}[(B_2D_2)(B_4^{}+B_5^{})]+{\displaystyle \frac{32}{3}}\lambda m_B^6s\left[\left|A_1\right|^2\left|C_1\right|^2\right]`$
$``$ $`{\displaystyle \frac{2}{r}}\lambda m_B^4s\left[\left|B_4\right|^2\left|B_5\right|^2\right]+{\displaystyle \frac{4}{r}}\lambda m_B^4m_{\mathrm{}}s\text{Re}[(B_3D_3)(B_4^{}+B_5^{})]`$
$``$ $`{\displaystyle \frac{8}{3r}}\lambda m_B^4(1rs)[\text{Re}(B_1B_2^{})\text{Re}(D_1D_2^{})]+{\displaystyle \frac{4}{3r}}\lambda m_B^2(\lambda +12rs)\left[\right|B_1|^2\left|D_1|^2\right]\},`$
$`P_T`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{\Delta }}}\sqrt{\lambda }\pi \{8m_B^3m_{\mathrm{}}\sqrt{s}\text{Re}[(A_1+C_1)(B_1^{}+D_1^{})]`$
$`+`$ $`{\displaystyle \frac{1}{2r}}m_B^3m_{\mathrm{}}(1+3r+s)\sqrt{s}\left[2\text{Re}(B_1D_2^{})2\text{Re}(B_2D_1^{})\right]`$
$`+`$ $`{\displaystyle \frac{1}{r\sqrt{s}}}m_Bm_{\mathrm{}}(1rs)\left[\left|B_1\right|^2\left|D_1\right|^2\right]`$
$`+`$ $`{\displaystyle \frac{1}{r\sqrt{s}}}m_Bm_{\mathrm{}}^2(1rs)\left[2\text{Re}(B_1B_5^{})2\text{Re}(D_1B_4^{})\right]`$
$`+`$ $`{\displaystyle \frac{1}{2r}}m_B^3m_{\mathrm{}}(1rs)\sqrt{s}\text{Re}[2(B_1+D_1)(B_3^{}D_3^{})]`$
$`+`$ $`{\displaystyle \frac{1}{r\sqrt{s}}}m_B^3m_{\mathrm{}}^2\lambda \left[2\text{Re}(B_2B_5^{})+2\text{Re}(D_2B_4^{})\right]`$
$`+`$ $`{\displaystyle \frac{1}{r\sqrt{s}}}m_B^5m_{\mathrm{}}(1r)\lambda \left[\left|B_2\right|^2\left|D_2\right|^2\right]+{\displaystyle \frac{1}{2r}}m_B^5m_{\mathrm{}}\lambda \sqrt{s}\text{Re}[2(B_2+D_2)(B_3^{}D_3^{})]`$
$`+`$ $`{\displaystyle \frac{1}{2r\sqrt{s}}}m_B^3m_{\mathrm{}}[(1rs)(1r)+\lambda ]\left[2\text{Re}(B_1B_2^{})+2\text{Re}(D_1D_2^{})\right]`$
$`+`$ $`{\displaystyle \frac{1}{2r\sqrt{s}}}m_B(1rs)(2m_{\mathrm{}}^2+m_B^2s)\left[2\text{Re}(D_1B_5^{})2\text{Re}(B_1B_4^{})\right]`$
$`+`$ $`{\displaystyle \frac{1}{2r\sqrt{s}}}m_B^3\lambda (2m_{\mathrm{}}^2+m_B^2s)[2\text{Re}(D_2B_5^{})+2\text{Re}(B_2B_4^{})]\},`$
$`P_N`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{\Delta }}}\pi vm_B^3\sqrt{\lambda }\sqrt{s}\{8m_{\mathrm{}}\text{Im}(B_1^{}C_1+A_1^{}D_1)`$ (13)
$`+`$ $`{\displaystyle \frac{1}{r}}m_B^2\lambda \left[\text{Im}(m_{\mathrm{}}B_3m_{\mathrm{}}D_3B_4)B_2^{}+\text{Im}(B_5B_3+D_3)D_2^{}\right]`$
$`+`$ $`{\displaystyle \frac{1}{r}}m_{\mathrm{}}(1+3rs)\text{Im}[(D_1+B_1)(B_2^{}D_2^{})]`$
$`+`$ $`{\displaystyle \frac{1}{r}}(1rs)[\text{Im}(B_4m_{\mathrm{}}B_3+m_{\mathrm{}}D_3)B_1^{}+\text{Im}(m_{\mathrm{}}B_3B_5m_{\mathrm{}}D_3)D_1^{}]\},`$
where $`\mathrm{\Delta }`$ is the expression within the curly parenthesis of the unpolarized differential decay rate in Eq. (11). These expressions for the longitudinal, transversal and normal polarizations are general and model independent (if the tensor interaction is neglected). It follows from the expressions of $`P_T`$ and $`P_N`$ that they are proportional to the lepton mass and therefore they are nonvanishing only for the $`\tau `$ lepton. In this work we also analyze the CP–violating asymmetry, which is defined as
$`A_{CP}(q^2)={\displaystyle \frac{\left({\displaystyle \frac{d\mathrm{\Gamma }}{dq^2}}\right)_0\left({\displaystyle \frac{d\overline{\mathrm{\Gamma }}}{dq^2}}\right)_0}{\left({\displaystyle \frac{d\mathrm{\Gamma }}{dq^2}}\right)_0+\left({\displaystyle \frac{d\overline{\mathrm{\Gamma }}}{dq^2}}\right)_0}},`$
where $`(d\mathrm{\Gamma }/dq^2)_0`$ is the unpolarized differential decay rate given by Eq. (11) and $`(d\overline{\mathrm{\Gamma }}/dq^2)_0`$ is the unpolarized differential decay rate for the antiparticle channel. Note that in SM , CP–violating asymmetry is equal to zero (or suppressed very strongly), since all form factors are real (see below), Wilson coefficients $`C_7^{eff}`$ and $`C_{10}`$ are real and only $`C_9^{eff}`$ contains a strong phase. But this strong phase can not lead to CP–violation itself. Using these general expressions we can study the sensitivity of the $`\tau `$–lepton polarizations on the new Wilson coefficients. Furthermore one can investigate how strongly these polarizations deviate from the SM predictions and for which Wilson coefficient this departure is more essential. But in the present work we will apply these general results to concrete models, namely to the SM and three type of 2HDM, i.e., models I, II (about models I and II, see for example and III. Note that in models I and II, the flavor changing neutral currents which appear at tree level are avoided by imposing ad hoc symmetry . The phenomenological consequence of the 2HDM without this discrete symmetry has been investigated in (see also ). One novel feature of model III is existence of new weak phase which appears in Yukawa interaction of fermions with Higgs fields (see below). Existence of this new weak phase can lead to sizeable CP violation in $`BK^{}\mathrm{}^+\mathrm{}^{}`$ decay. Therefore if in future experiments sizeable CP violation in the $`BK^{}\mathrm{}^+\mathrm{}^{}`$ decay is discovered, it is an unambiguous indication of the existence of new physics beyond SM, since in the SM the CP asymmetry suppressed very strongly.
Making the following replacements in the expressions given in Eq. (8), the explicit forms of $`A_i`$, $`B_i`$, $`C_i`$ and $`D_i`$ can be obtained in SM and 2HDM easily.
1. SM
$`C_{LL}`$ $`=`$ $`C_9^{eff}(m_b)C_{10}(m_b),`$
$`C_{RL}`$ $`=`$ $`0,`$
$`C_{BR}`$ $`=`$ $`2C_7^{eff}(m_b),`$
$`C_{LR}`$ $`=`$ $`C_9^{eff}(m_b)+C_{10}(m_b),`$
$`C_{RR}`$ $`=`$ $`0,`$
$`C_{LRRL}`$ $`=`$ $`C_{RLLR}=C_{LRLR}=C_{RLRL}=0.`$ (14)
2. 2HDM
$`C_{LL}`$ $`=`$ $`C_9^{eff\mathrm{\hspace{0.17em}2}HDM}(m_b)C_{10}^{2HDM}(m_b),`$
$`C_{RL}`$ $`=`$ $`0,`$
$`C_{BR}`$ $`=`$ $`2C_7^{eff\mathrm{\hspace{0.17em}2}HDM}(m_b),`$
$`C_{LR}`$ $`=`$ $`C_9^{eff\mathrm{\hspace{0.17em}2}HDM}(m_b)+C_{10}^{2HDM}(m_b),`$
$`C_{RR}`$ $`=`$ $`0,`$
$`C_{LRRL}`$ $`=`$ $`C_{Q_1},`$
$`C_{RLRL}`$ $`=`$ $`C_{Q_2},`$
$`C_{LRLR}`$ $`=`$ $`C_{Q_1},`$
$`C_{RLLR}`$ $`=`$ $`C_{Q_2}.`$ (15)
The coefficients $`C_i^{2HDM}(m_W)`$ ($`i=7,9`$ and $`10`$) to the leading order are given by (see for example )
$`C_7^{2HDM}(m_W)`$ $`=`$ $`x{\displaystyle \frac{(75x8x^2)}{24(x1)^3}}+{\displaystyle \frac{x^2(3x2)}{4(x1)^4}}\mathrm{ln}x`$ (16)
$`+`$ $`\left|\lambda _{tt}\right|^2\left({\displaystyle \frac{y(75y8y^2)}{72(y1)^3}}+{\displaystyle \frac{y^2(3y2)}{12(y1)^4}}\mathrm{ln}y\right)`$
$`+`$ $`\lambda _{tt}\lambda _{bb}\left({\displaystyle \frac{y(35y)}{12(y1)^2}}+{\displaystyle \frac{y(3y2)}{6(y1)^3}}\mathrm{ln}y\right),`$
$`C_9^{2HDM}(m_W)`$ $`=`$ $`{\displaystyle \frac{1}{sin^2\theta _W}}B(m_W)+{\displaystyle \frac{14sin^2\theta _W}{sin^2\theta _W}}C(m_W)`$ (17)
$`+`$ $`{\displaystyle \frac{x^2(2519x)}{36(x1)^3}}+{\displaystyle \frac{3x^4+30x^354x^2+32x8}{18(x1)^4}}\mathrm{ln}x+{\displaystyle \frac{4}{9}}`$
$`+`$ $`|\lambda _{tt}|^2[{\displaystyle \frac{14sin^2\theta _W}{sin^2\theta _W}}{\displaystyle \frac{xy}{8}}({\displaystyle \frac{1}{y1}}{\displaystyle \frac{1}{(y1)^2}}\mathrm{ln}y)`$
$``$ $`y({\displaystyle \frac{47y^279y+38}{108(y1)^3}}{\displaystyle \frac{3y^36y^3+4}{18(y1)^4}}\mathrm{ln}y)],`$
$`C_{10}^{2HDM}(m_W)`$ $`=`$ $`{\displaystyle \frac{1}{sin^2\theta _W}}\left(B(m_W)C(m_W)\right)`$ (18)
$`+`$ $`\left|\lambda _{tt}\right|^2{\displaystyle \frac{1}{sin^2\theta _W}}{\displaystyle \frac{xy}{8}}\left({\displaystyle \frac{1}{y1}}+{\displaystyle \frac{1}{(y1)^2}}\mathrm{ln}y\right),`$
$`C_{Q_1}(m_W)`$ $`=`$ $`{\displaystyle \frac{m_bm_{\mathrm{}}}{m_{h^0}^2}}{\displaystyle \frac{1}{\left|\lambda _{tt}\right|^2}}{\displaystyle \frac{1}{sin^2\theta _W}}{\displaystyle \frac{x}{4}}\{(sin^2\alpha +hcos^2\alpha )f_1(x,y)+`$ (19)
$`+`$ $`\left[{\displaystyle \frac{m_{h^0}^2}{m_W^2}}+\left(sin^2\alpha +hcos^2\alpha \right)(1z)\right]f_2(x,y)+`$
$`+`$ $`{\displaystyle \frac{sin^22\alpha }{2m_{H^\pm }^2}}[m_{h^0}^2{\displaystyle \frac{(m_{h^0}^2+m_{H^0}^2)^2}{2m_{H^0}^2}}]f_3(y)\},`$
$`C_{Q_2}(m_W)`$ $`=`$ $`{\displaystyle \frac{m_bm_{\mathrm{}}}{m_{H^\pm }^2}}{\displaystyle \frac{1}{\left|\lambda _{tt}\right|^2}}\left\{f_1(x,y)+\left[1+{\displaystyle \frac{m_{H^\pm }^2m_{A^0}^2}{m_W^2}}\right]f_2(x,y)\right\},`$ (20)
where
$`x`$ $`=`$ $`{\displaystyle \frac{m_t^2}{m_W^2}},y={\displaystyle \frac{m_t^2}{m_{H^\pm }^2}},z={\displaystyle \frac{x}{y}},h={\displaystyle \frac{m_{h^0}^2}{m_{H^0}^2}},`$
$`B(x)`$ $`=`$ $`{\displaystyle \frac{x}{4(x1)}}+{\displaystyle \frac{x}{4(x1)^2}}\mathrm{ln}x,`$
$`C(x)`$ $`=`$ $`{\displaystyle \frac{x}{4}}\left({\displaystyle \frac{x6}{2(x1)}}+{\displaystyle \frac{3x+2}{2(x1)^2}}\mathrm{ln}x\right),`$
$`f_1(x,y)`$ $`=`$ $`{\displaystyle \frac{x\mathrm{ln}x}{x1}}{\displaystyle \frac{y\mathrm{ln}y}{y1}},`$
$`f_2(x,y)`$ $`=`$ $`{\displaystyle \frac{x\mathrm{ln}y}{(zx)(x1)}}+{\displaystyle \frac{\mathrm{ln}z}{(z1)(x1)}},`$
$`f_3(y)`$ $`=`$ $`{\displaystyle \frac{1y+y\mathrm{ln}y}{(y1)^2}},`$ (21)
$`sin^2\theta _W=0.23`$ is the Weinberg angle, $`h^0,H^0`$ and $`A^0`$ are two scalar and pseudoscalar Higgs fields, respectively. The coefficients $`\lambda _{tt}`$ and $`\lambda _{bb}`$ for model I and model II of the 2HDM are:
$`\lambda _{tt}=\mathrm{cot}\beta ,\lambda _{bb}=\mathrm{cot}\beta ,\text{for model I},`$
$`\lambda _{tt}=\mathrm{cot}\beta ,\lambda _{bb}=+\mathrm{tan}\beta ,\text{for model II},`$ (22)
while in model III $`\lambda _{tt}`$ or $`\lambda _{bb}`$ is complex, i.e.,
$`\lambda _{tt}\lambda _{bb}\left|\lambda _{tt}\lambda _{bb}\right|e^{i\varphi }.`$
From Eqs. (16)–(20) we observe that the SM results for the Wilson coefficients $`C_7^{SM}(m_W)`$, $`C_9^{SM}(m_W)`$ amd $`C_{10}^{SM}(m_W)`$ (and correspondingly at $`\mu =m_b`$ scale) can all be obtained from 2HDM results by making the following replacements
$`C_{Q_1}0`$ , $`C_{Q_2}0,`$
$`C_7^{SM}(m_W)`$ $`=`$ $`C_7^{2HDM}(y0),`$
$`C_9^{SM}(m_W)`$ $`=`$ $`C_9^{2HDM}(y0),`$
$`C_{10}^{SM}(m_W)`$ $`=`$ $`C_{10}^{2HDM}(y0),`$
The evolution of the Wilson coefficients from the higher scale $`\mu =m_W`$ down to the low energy scale $`\mu =m_b`$ is described by the renormalization group equation. The coefficients $`C_7^{eff}(\mu ),C_9^{eff}(\mu ),C_{10}(\mu )`$ at the scale $`𝒪(\mu =m_b)`$ are calculated in and $`C_{Q_1}`$ and $`C_{Q_2}`$ at the same scale to leading order are calculated in . The Wilson coefficient $`C_{10}`$ is not modified as we move from $`\mu =m_W`$ to $`\mu =m_b`$ scale, i.e., $`C_{10}(m_b)C_{10}^{2HDM}(m_W)`$. In order to calculate $`C_9^{2HDM}`$ at $`m_b`$ scale, it is enough to make the replacement $`C_9^{SM}(m_W)C_9^{2HDM}(m_W)`$ and then solve the corresponding renormalization group equation. Hence, including the NLO QCD corrections, $`C_9^{eff}(m_b)`$ can be written as:
$`C_9^{eff}(\mu )=C_9^{2HDM}(\mu )\left[1+{\displaystyle \frac{\alpha _s(\mu )}{\pi }}\omega (\widehat{s})\right]`$ (23)
$`+g(\widehat{m}_c,\widehat{s})\left[3C_1(\mu )+C_2(\mu )+3C_3(\mu )+C_4(\mu )+3C_5(\mu )+C_6(\mu )\right]`$
$`{\displaystyle \frac{1}{2}}g(0,\widehat{s})\left(C_3(\mu )+3C_4(\mu )\right){\displaystyle \frac{1}{2}}g(1,\widehat{s})\left(4C_3+4C_4+3C_5+C_6\right)`$
$`{\displaystyle \frac{1}{2}}g(0,\widehat{s})\left(C_3+3C_4\right)+{\displaystyle \frac{2}{9}}\left(3C_3+C_4+3C_5+C_6\right),`$
where $`\widehat{m}_c=m_c/m_b,\widehat{s}=p^2/m_b^2`$, and
$`\omega \left(\widehat{s}\right)={\displaystyle \frac{2}{9}}\pi ^2{\displaystyle \frac{4}{3}}Li_2\left(\widehat{s}\right){\displaystyle \frac{2}{3}}\mathrm{ln}\left(\widehat{s}\right)\mathrm{ln}\left(1\widehat{s}\right)`$ (24)
$`{\displaystyle \frac{5+4\widehat{s}}{3\left(1+2\widehat{s}\right)}}\mathrm{ln}\left(1\widehat{s}\right){\displaystyle \frac{2\widehat{s}\left(1+\widehat{s}\right)\left(12\widehat{s}\right)}{3\left(1\widehat{s}\right)^2\left(1+2\widehat{s}\right)}}\mathrm{ln}\left(\widehat{s}\right)+{\displaystyle \frac{5+9\widehat{s}6\widehat{s}^2}{3\left(1\widehat{s}\right)\left(1+2\widehat{s}\right)}}`$
represents the $`𝒪\left(\alpha _s\right)`$ correction from the one gluon exchange in the matrix element of $`O_9`$, while the function $`g(\widehat{m}_c,\widehat{s})`$ arises from one loop contributions of the four–quark operators $`O_1`$$`O_6`$, whose form is
$`g(y_i,\widehat{s})`$ $`=`$ $`{\displaystyle \frac{8}{9}}\mathrm{ln}\left(\widehat{m}_i\right)+{\displaystyle \frac{8}{27}}+{\displaystyle \frac{4}{9}}y_i`$ (25)
$`{\displaystyle \frac{2}{9}}(2+y_i)\sqrt{\left|1y_i\right|}\{\mathrm{\Theta }(1y_i)(\mathrm{ln}{\displaystyle \frac{1+\sqrt{\left|1y_i\right|}}{1\sqrt{\left|1y_i\right|}}}i\pi )`$
$`+\mathrm{\Theta }(y_i1)2\mathrm{arctan}{\displaystyle \frac{1}{\sqrt{y_i1}}}\},`$
where $`y_i=4\widehat{m}_{i}^{}{}_{}{}^{2}/\widehat{p}^2`$. The Wilson coefficient $`C_9^{eff}`$ receives also long distance contributions, which have their origin in the real $`c\overline{c}`$ intermediate states, i.e., $`J/\psi `$, $`\psi ^{}`$, $`\mathrm{}`$. The $`J/\psi `$ family is introduced by the Breit–Wigner distribution for the resonances through the replacement
$`g(\widehat{m}_c,\widehat{s})g(\widehat{m}_c,\widehat{s}){\displaystyle \frac{3\pi }{\alpha _{em}^2}}\kappa {\displaystyle \underset{V_i=J/\psi _i,\psi ^{},\mathrm{}}{}}{\displaystyle \frac{m_{V_i}\mathrm{\Gamma }(V_i\mathrm{}^+\mathrm{}^{})}{(p^2m_{V_i}^2)+im_{V_i}\mathrm{\Gamma }_{V_i}}},`$ (26)
where the phenomenological parameter $`\kappa =2.3`$ is chosen in order to reproduce correctly the experimental value of the branching ratio (see for example )
## 3 Numerical analysis
In this section we would like to present our numerical results. The main free parameters $`\lambda _{tt},\lambda _{bb}`$ of the 2HDM are restricted from $`BX_s\gamma `$ decay, $`B^0`$$`\overline{B}^0`$ mixing, $`\rho `$ parameter and neutron electric–dipole moment , that yields $`\left|\lambda _{bb}\right|=50`$, $`\left|\lambda _{tt}\right|0.03`$. Throughout the numerical analysis for the mass of the Higgs bosons we have used $`m_{h^0}=80GeV,m_{H^\pm }=250GeV,m_{A^0}=250GeV`$ and $`m_{H^0}=150GeV`$.
For the values of the form factors, we have used the results of , where the radiative corrections to the leading twist contribution and $`SU(3)`$ breaking effects are also taken into account. The $`q^2`$ dependence of the form factors can be represented in terms of three parameters as
$`F(q^2)={\displaystyle \frac{F(0)}{1a_F\frac{q^2}{m_B^2}+b_F\left(\frac{q^2}{m_B^2}\right)^2}},`$
where, the values of parameters $`F(0)`$, $`a_F`$ and $`b_F`$ for the $`BK^{}\mathrm{}^+\mathrm{}^{}`$ decay are listed in Table 1.
In Fig. (1) we present the dependence of the CP–violating asymmetry on $`q^2`$ and on the weak phase $`\varphi `$ for the $`BK^{}\tau ^+\tau ^{}`$ decay in model III, since we have already noted that in SM and in models I and II the CP–violating asymmetry is practically zero. We observe that CP asymmetry differs from zero in the region $`0<\varphi <2\pi `$, except at $`\varphi =0,\pi `$ and $`2\pi `$, and its value in the region $`0<\varphi <\pi `$ ($`\pi <\varphi <2\pi `$) is negative (positive).
Fig .(2) depicts the dependence of the averaged CP asymmetry $`A_{CP}`$ (here and in all of the following discussions, by the averaged values of the physical quantities we mean integration over $`q^2`$ in the region $`14GeV^2q^2(m_Bm_K^{})^2`$) on the weak phase angle $`\varphi `$ in model III, taking into account short and long distance contributions. It follows from this figure that $`A_{CP}`$ varies in the range (-0.04, 0.04) which is different from zero and it definitely is an indication of the existence of new physics beyond SM, since $`A_{CP}`$ is practically equal to zero in the SM.
In Fig. (3), the dependence of $`P_L`$ on $`q^2`$ and the weak phase angle $`\varphi `$ without long distance effects in model III is presented. From this figure one can see that, for $`q^2>14GeV^2`$, $`P_L`$ varies in the range (-0.65, -0.8) which is larger than the SM prediction. This is due to the fact that the ”new” contribution which comes from the charged Higgs boson gives constructive interference to the SM results.
In Fig. (4) we present the averaged longitudinal polarization $`P_L`$ on the weak phase angle $`\varphi `$, taking into account short and long distance contributions. For completeness the predictions of SM, model I and model II on $`P_L`$ are also presented. It is observed from this figure that $`P_L`$ in model III as modulo, is larger than the ones predicted by SM, model I and model II. Therefore an observation of $`\left|P_L\right|0.65`$ is another conclusive confirmation of the existence of new physics beyond SM.
In Figs. (5) and (6) we present the dependence of the transversal and normal polarizations of the $`\tau `$ lepton on $`q^2`$ and on the weak phase angle $`\varphi `$, respectively, without long distance effects in model III. The dependence of the averaged transversal and normal polarizations on the weak phase angle, taking into account short and long distance contributions, are depicted in Figs. (7) and (8), respectively. For sake of completeness we presented also the predictions of SM, model I and model II of the same physical quantity in both figures. Fig. (7) clearly depicts that, the prediction of model III on $`P_T`$ as modulo, is approximately five times smaller than the ones predicted by SM, model I and model II. However the situation is totally different for $`P_N`$, having a range of values in the region (0.10, 0.15) in model III, it is approximately two or three times larger than the ones predicted by SM, model I and model II.
Here we would like to make the following remark. It follows from Eq. (13) that $`P_N`$ is defined as the imaginary part of the form factors and of the corresponding Wilson coefficients $`C_7^{eff},C_9^{eff}`$, $`C_{10}`$, $`C_{Q_1}`$ and $`C_{Q_2}`$. In SM $`C_7^{eff}`$ and $`C_{10}`$ are real and only $`C_9^{eff}`$ has imaginary part. On the other side all theoretical methods predict these form factors to be real quantities. For this reason, if in future experiments a different value for $`P_N`$ were observed compared to the SM prediction, it is an indication of unambiguous information about the existence of the above–mentioned CP–violating phase in theory.
In conclusion, we have investigated the exclusive $`BK^{}\tau ^+\tau ^{}`$ decay in the SM and in three different versions of the 2HDM. From the results we have obtained we conclude that the combined analysis of the CP–violating asymmetry and $`\tau `$ lepton polarization effects are very useful tools in looking for new physics beyond SM.
## Figure captions
Fig. 1 The dependence of the CP–violating asymmetry $`A_{CP}`$ of the $`\tau `$ lepton on $`q^2`$ and on the weak phase $`\varphi `$ in model III.
Fig. 2 The dependence of the averaged CP asymmetry $`A_{CP}`$ of the $`\tau `$ lepton on the weak phase $`\varphi `$ in model III, taking into account short and long distance contributions.
Fig. 3 The dependence of the longitudinal polarization $`P_L`$ of the $`\tau `$ lepton on $`q^2`$ and on the weak phase $`\varphi `$ in model III, taking into account only the short distance contribution in $`C_9^{eff}`$.
Fig. 4 The dependence of the averaged longitudinal polarization $`P_L`$ of $`\tau `$ lepton on the weak phase $`\varphi `$, taking into account short and long distance contributions.
Fig. 5 The dependence of the transversal polarization $`P_T`$ of the $`\tau `$ lepton on $`q^2`$ and on the weak phase $`\varphi `$ in model III, taking into account only the short distance contribution in $`C_9^{eff}`$.
Fig. 6 The dependence of the normal polarization $`P_N`$ of the $`\tau `$ lepton on $`q^2`$ and on the weak phase $`\varphi `$ in model III, taking into account only the short distance contribution in $`C_9^{eff}`$.
Fig. 7 The dependence of the averaged transversal polarization $`P_T`$ of the $`\tau `$ lepton on the weak phase $`\varphi `$, taking into account short and long distance contributions.
Fig. 8 The dependence of the averaged normal polarization $`P_N`$ of the $`\tau `$ lepton on the weak phase $`\varphi `$, taking into account short and long distance contributions.
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# BRST Invariance and Renormalisability of the SU(2)×U(1) Electroweak Theory with Massive W Z Bosons
## Abstract
Since the SU(n) gauge theory with massive gauge bosons has been proven to be renormalisable we reinvestigate the renormalisability of the SU<sub>L</sub>(2) $`\times `$ U<sub>Y</sub>(1) electroweak theory with massive W Z bosons. We expound that with the constraint conditions caused by the W Z mass term and the additional condition chosen by us we can performed the quantization and construct the ghost action in a way similar to that used for the massive SU(n) theory. We also show that when the $`\delta `$ functions appearing in the path integral of the Green functions and representing the constraint conditions are rewritten as Fourier integrals with Lagrange multipliers $`\lambda _a`$ and $`\lambda _y`$, the BRST invariance is kept in the total effective action consisting of the Lagrange multipliers, ghost fields and the original fields. Furthermore, by comparing with the massless theory and with the massive SU(n) theory we find the general form of the divergent part of the generating functional for the regular vertex functions and prove the renormalisability of the theory. It is also clarified that the renormalisability of the theory with the W Z mass term is ensured by that of the massless theory and the massive SU(n) theory.
PACS numbers: 03.65.Db, 03.80.+r, 11.20.Dj
I. Introduction
Although the negative answer to the problem of renormalisability of a SU(n) theory with massive gauge bosons is widely known, such theories continue to be studied (see for example Refs. \[1-8\]). However, since the negative answer had not been voted down, it was naturally difficult to investigate the possibility of directly adding a mass term to the SU<sub>L</sub>(2) $`\times `$ U<sub>Y</sub>(1) theory. Recently, the renormalisability of the massive SU(n) gauge theory has been proven . Therefore we will reinvestigate the SU<sub>L</sub>(2) $`\times `$ U<sub>Y</sub>(1) theory of S.L.Glashow with the mass term of the W Z fields. The study of the theory including the mass term of the matter fields as well w1ll be reported in Ref. .
In order to make appropriate the mass ratio, the W Z mass term must contain a product of the SU<sub>L</sub>(2) and U<sub>Y</sub>(1) fields and thus cause constraint conditions containing products of such fields. Next, such a mass term is invariant under an infinitesimal gauge transformation with $`\delta \theta _1`$ and $`\delta \theta _2`$ equal to zero and $`\delta \theta _3`$ equal to $`\delta \theta _y`$, where $`\theta _a`$ and $`\theta _1`$ are the parameters of the gauge group. Therefore an additional constraint condition should be properly chosen. We will expound that with the constraint conditions caused by the W Z mass term and the additional condition chosen by us we can performed the quantization and construct the ghost action in a way similar to that used for the massive SU(n) theory . We will also show that when the $`\delta `$ functions appearing in the path integral of the Green functions and representing the constraint conditions are rewritten as Fourier integrals with Lagrange multipliers $`\lambda _a`$ and $`\lambda _y`$, the BRST invariance is kept in the total effective action consisting of the Lagrange multipliers, ghost fields and the original fields.
As the constraint conditions contain the products of the SU<sub>L</sub>(2) and U<sub>Y</sub>(1) fields, the divergent part of the generating functional $`\mathrm{\Gamma }`$ for the regular vertex functions is dependent on the classical fields of the Lagrange multipliers $`\lambda _a`$ and $`\lambda _y`$ when the generating functional for the Green functions contains the sources of these Lagrange multipliers. The problem of whether such a generalized form of the theory is renormalisable becomes complicated. However, we are not interested in using the Green functions involving $`\lambda _a`$ or $`\lambda _y`$. Thus we can avoid introducing the sources of these Lagrange multipliers to the generating functional for the Green functions. An equivalent and convenient procedure is to derive the Slavnov–Taylor identities and the additional identities for $`\mathrm{\Gamma }`$ with the help of the generalized form of the theory and then let vanish the functional derivatives of $`\mathrm{\Gamma }`$ with respect to the classical fields of these Lagrange multipliers. In this way the divergent part of $`\mathrm{\Gamma }`$ will be shown to satisfy the same equations appearing in the massless theory. Furthermore, by comparing with the massless theory and with the massive SU(n) theory we will be able to find the general form of the divergent part of $`\mathrm{\Gamma }`$ and prove the renormalisability of the theory. Meanwhile it will be clarified that the renormalisability of the theory with the W Z mass term is ensured by that of the massless theory and the massive SU(n) theory.
In section $`2`$ we will find the constraint conditions caused by the W Z mass term. The additional constraint condition will also be chosen. The method of quantization will be explained in section $`3`$. Setion $`4`$ is devoted to prove the renormalisability of the theory. Concluding remarks will be given in the final section.
II. Original and Additional Constraint Conditions
For the sake of comvenience we assume in the present work that the matter fields consist only of the electron and electron-neutrino fields and are often denoted by $`\psi (x)`$ and $`\overline{\psi }(x)`$. The former stands for the purely left-handed neutrino field $`\nu _L`$, the left- and right-handed parts of the electron field namely $`e_L`$, $`e_R`$, and the latter stands for $`\overline{\nu }_L`$, $`\overline{e}_L`$ and $`\overline{e}_R`$. Next let $`W_{a\mu }(x)`$, $`W_{y\mu }(x)`$ be the SU<sub>L</sub>(2) and U<sub>Y</sub>(1) gauge fields and $`g`$, $`g_1`$ be the coupling constants. Thus the W Z mass term in the Lagrangian is
$`_{WM}={\displaystyle \frac{1}{2}}M^2W_{a\mu }W_a^\mu +{\displaystyle \frac{1}{2}}M^2\left({\displaystyle \frac{g_1}{g}}\right)^2W_{y\mu }W_y^\mu M^2\left({\displaystyle \frac{g_1}{g}}\right)W_{3\mu }W_y^\mu ,`$ (2.1)
or
$$_{WM}=\frac{1}{2}M^2W_{1\mu }(x)W_1^\mu (x)+\frac{1}{2}M^2W_{2\mu }(x)W_2^\mu (x)+\frac{1}{2}M_z^2Z_\mu (x)Z^\mu (x),$$
where $`M_z^2`$ stands for $`g^2(g^2+g_1^2)M^2`$, and $`Z_\mu (x)`$, $`A_\mu (x)`$ are the field functions of Z boson and photon, namely
$`Z_\mu ={\displaystyle \frac{1}{\sqrt{(g^2+g_1^2)}}}(gW_{3\mu }g_1W_{y\mu }),`$ (2.2)
$`A_\mu ={\displaystyle \frac{1}{\sqrt{(g^2+g_1^2)}}}\epsilon (g_1W_{3\mu }+gW_{y\mu }),`$ (2.3)
where $`\epsilon `$ is $`1`$ or $`1`$.
The original Lagrangian of the SU<sub>L</sub>(2) $`\times `$ U<sub>Y</sub>(1) electroweak theory with the mass term $`_{WM}`$ is
$`=_\psi +_{\psi W}+_{WM}+_{WL}+_{WY},`$ (2.4)
where $`_\psi `$ describe the pure matter fields, $`_{\psi W}`$ is the coupling term between the matter and gauge fields and
$`_{WL}={\displaystyle \frac{1}{4}}F_{a\mu \nu }F_a^{\mu \nu },`$ (2.5)
$`_{WY}={\displaystyle \frac{1}{4}}B_{\mu \nu }B^{\mu \nu },`$ (2.6)
with
$`F_{a\mu \nu }=_\mu W_{a\nu }_\nu W_{a\mu }gC_{abc}W_{b\mu }W_{c\nu },`$ (2.7)
$`B_{\mu \nu }=_\mu W_{y\nu }_\nu W_{y\mu }.`$ (2.8)
$`C_{abc}`$ stands for the structure constants of SU<sub>L</sub>(2) with $`C_{123}`$ equal to $`1`$.
Denote by $`\theta _a(x),\theta _y(x)`$ the parameters of the gauge group. Thus, under an infinitesimal gauge transformation, the fields $`W_a^\mu `$, $`W_y^\mu `$, $`\psi `$ and $`\overline{\psi }`$ transform as
$`\delta W_a^\mu (x)={\displaystyle \frac{1}{g}}^\mu \delta \theta _a(x)C_{abc}W_c^\mu (x)\delta \theta _b(x),`$
$`\delta W_y^\mu (x)={\displaystyle \frac{1}{g_1}}^\mu \delta \theta _y(x),`$
$`\delta \nu _L(x)={\displaystyle \frac{i}{2}}\delta \theta _1(x)e_L(x)+{\displaystyle \frac{1}{2}}\delta \theta _2(x)e_L(x)+{\displaystyle \frac{i}{2}}\delta \theta _3(x)\nu _L(x){\displaystyle \frac{i}{2}}\delta \theta _y(x)\nu _L(x),`$
$`\delta e_L(x)={\displaystyle \frac{i}{2}}\delta \theta _1(x)\nu _L(x){\displaystyle \frac{1}{2}}\delta \theta _2(x)\nu _L(x){\displaystyle \frac{i}{2}}\delta \theta _3(x)e_L(x){\displaystyle \frac{i}{2}}\delta \theta _y(x)e_L(x),`$
$`\delta e_R(x)=i\delta \theta _y(x)e_R(x),`$
$`\delta \overline{\nu }_L(x)={\displaystyle \frac{i}{2}}\delta \theta _1(x)\overline{e}_L(x)+{\displaystyle \frac{1}{2}}\delta \theta _2(x)\overline{e}_L(x){\displaystyle \frac{i}{2}}\delta \theta _3(x)\overline{\nu }_L(x)+{\displaystyle \frac{i}{2}}\delta \theta _y(x)\overline{\nu }_L(x),`$
$`\delta \overline{e}_L(x)={\displaystyle \frac{i}{2}}\delta \theta _1(x)\overline{\nu }_L(x){\displaystyle \frac{1}{2}}\delta \theta _2(x)\overline{\nu }_L(x)+{\displaystyle \frac{i}{2}}\delta \theta _3(x)\overline{e}_L(x)+{\displaystyle \frac{i}{2}}\delta \theta _y(x)\overline{e}_L(x),`$
$`\delta \overline{e}_R(x)=i\delta \theta _y(x)\overline{e}_R(x).`$
Therefore the action transforms as
$`\delta {\displaystyle d^4x(x)}=\delta {\displaystyle d^4x_{WM}(x)}`$
$`={\displaystyle }d^4x\{({\displaystyle \frac{M^2}{g}}_\mu W_1^\mu (x)+{\displaystyle \frac{M^2}{g}}g_1W_{2\mu }(x)W_y^\mu (x))\delta \theta _1`$
$`+\left({\displaystyle \frac{M^2}{g}}_\mu W_2^\mu (x){\displaystyle \frac{M^2}{g}}g_1W_{1\mu }(x)W_y^\mu (x)\right)\delta \theta _2`$
$`+({\displaystyle \frac{M^2}{g}}_\mu W_3^\mu (x){\displaystyle \frac{M^2}{g^2}}g_1_\mu W_y^\mu (x))(\delta \theta _3\delta \theta _y)\}.`$ (2.9)
Since the classical equations of motion make the action invariant under an arbitrary infinitesimal transformation of the field functions, they certainly make the W Z mass term invariant under an arbitrary infinitesimal gauge transformation. This means that when $`M`$ is not equal to zero, the classical equations of motion leads to the following constraint conditions
$`{\displaystyle \frac{M^2}{g}}_\mu W_1^\mu (x)+{\displaystyle \frac{M^2}{g}}g_1W_{2\mu }(x)W_y^\mu (x)=0,`$ (2.10)
$`{\displaystyle \frac{M^2}{g}}_\mu W_2^\mu (x){\displaystyle \frac{M^2}{g}}g_1W_{1\mu }(x)W_y^\mu (x)=0,`$ (2.11)
$`{\displaystyle \frac{M^2}{g}}_\mu W_3^\mu (x){\displaystyle \frac{M^2}{g^2}}g_1_\mu W_y^\mu (x)=0.`$ (2.12)
These are the original constraint conditions. As it can be seen from (2.9) that the W Z mass term is invariant under an infinitesimal gauge transformation with $`\delta \theta _1`$ and $`\delta \theta _2`$ equal to zero and $`\delta \theta _3`$ equal to $`\delta \theta _y`$. For this reason, $`_\mu W_3^\mu `$ and $`_\mu W_y^\mu `$ appear in one constraint. We now choose an additional condition and replace (2.12) with
$`{\displaystyle \frac{M^2}{g}}_\mu W_3^\mu (x)+{\displaystyle \frac{M^2}{g}}g_1W_{3\mu }(x)W_y^\mu (x)=0,`$ (2.13)
$`_\mu W_y^\mu (x)+gW_{3\mu }(x)W_y^\mu (x)=0.`$ (2.14)
III. Quantization and BRST Invariance
Write (2.10), (2.11) and (2.13),(2.14) as
$`\mathrm{\Phi }_a(x)=0,\mathrm{\Phi }_y(x)=0,`$ (3.1)
with
$`\mathrm{\Phi }_1(x)=_\mu W_1^\mu (x)+g_1W_{2\mu }(x)W_y^\mu (x),`$ (3.2)
$`\mathrm{\Phi }_2(x)=_\mu W_2^\mu (x)g_1W_{1\mu }(x)W_y^\mu (x),`$ (3.3)
$`\mathrm{\Phi }_3(x)=_\mu W_3^\mu (x)+g_1W_{3\mu }(x)W_y^\mu (x),`$ (3.4)
$`\mathrm{\Phi }_y(x)=_\mu W_y^\mu (x)+gW_{3\mu }(x)W_y^\mu (x).`$ (3.5)
Taking the constraint conditions (3.1) into account one should write the path integral of the Green functions inolving only the original fields as
$`{\displaystyle \frac{1}{N_0}}{\displaystyle 𝒟[𝒲,\overline{\psi },\psi ]\mathrm{\Delta }[𝒲,\overline{\psi },\psi ]\underset{a^{},x^{}}{}\delta \left(\mathrm{\Phi }_a^{}(x^{})\right)\delta \left(\mathrm{\Phi }_y(x^{})\right)W_{a\mu }(x)W_{b\nu }(y)\mathrm{}\mathrm{exp}\{\mathrm{i}I\}},`$ (3.6)
where
$`I={\displaystyle d^4x(x)},`$
$`N_0={\displaystyle 𝒟[𝒲,\overline{\psi },\psi ]\mathrm{\Delta }[𝒲,\overline{\psi },\psi ]\underset{a^{},x^{}}{}\delta \left(\mathrm{\Phi }_a^{}(x^{})\right)\delta \left(\mathrm{\Phi }_y(x^{})\right)\mathrm{exp}\{\mathrm{i}I\}}.`$
The weight factor $`\mathrm{\Delta }[𝒲,\overline{\psi },\psi ]`$ is to be determined. Since only the field functions which satisfy the constraint conditions can play roles in the integral (3.6), the value of the Lagrangian can be changed for the field functions which do not satisfy these conditions. In view of the fact that the conditions (3.1) make the action invariant with respect to the infinitesimal gauge trasformation, we now imagine to replace the mass term $`_{WM}`$ in (3.6) with a gauge invariant mass term which is equal to $`_{WM}`$ when the conditions (3.1) are satisfied. Thus, analogous to the case in the Fadeev–Popov method \[1,11-16\], $`\mathrm{\Delta }[𝒲,\overline{\psi },\psi ]`$ should be gauge invariant and make the following equation valid for an arbitrary gauge invariant quantity $`𝒪(𝒲,\overline{\psi },\psi )`$
$`{\displaystyle 𝒟[𝒲,\overline{\psi },\psi ]\mathrm{\Delta }[𝒲,\overline{\psi },\psi ]\underset{a^{},x^{}}{}\delta \left(\mathrm{\Phi }_a^{}(x^{})\right)\delta \left(\mathrm{\Phi }_y(x^{})\right)𝒪(𝒲,\overline{\psi },\psi )\mathrm{exp}\{\mathrm{i}\stackrel{~}{I}\}}`$
$`{\displaystyle 𝒟[𝒲,\overline{\psi },\psi ]𝒪(𝒲,\overline{\psi },\psi )\mathrm{exp}\{\mathrm{i}\stackrel{~}{I}\}}.`$
where $`\stackrel{~}{I}`$ is a gauge invariant action constructed by replacing $`_{WM}`$ with the imagined mass term. This means that the weight factor $`\mathrm{\Delta }[𝒲,\overline{\psi },\psi ]`$ can be determined according to the Fadeev–Popov equation of the following form
$`\mathrm{\Delta }[𝒲,\overline{\psi },\psi ]{\displaystyle \underset{z}{}d\mathrm{\Omega }(z)\underset{\sigma ,x}{}\delta \left(\mathrm{\Phi }_\sigma ^\mathrm{\Omega }(x)\right)}=1.`$ (3.7)
where $`\sigma `$ stands for $`1,2,3,y`$, $`\mathrm{\Phi }_\sigma ^\mathrm{\Omega }(x)`$ is the result of acting on $`\mathrm{\Phi }_\sigma (x)`$ with a gauge transformation having the parameters of the element $`\mathrm{\Omega }(x)`$ of the gauge group, $`d\mathrm{\Omega }(z)`$ is the volume element of the group integral. It follows that with the F–P ghost fields $`C_a(x)`$, $`C_y(x)`$, $`\overline{C}_a(x)`$, $`\overline{C}_y(x)`$ as new variables, one can express the ghost Lagrangian as
$`^{(C)}(x)=\overline{C}_a(x)\mathrm{\Delta }\mathrm{\Phi }_a(x)+\overline{C}_y(x)\mathrm{\Delta }\mathrm{\Phi }_y(x),`$ (3.8)
where $`\mathrm{\Delta }\mathrm{\Phi }_a(x)`$, $`\mathrm{\Delta }\mathrm{\Phi }_y(x)`$ are defined by the BRST transformtion of $`\mathrm{\Phi }_a(x)`$ and $`\mathrm{\Phi }_y(x)`$ so that
$`\delta _B\mathrm{\Phi }_a(x)=\delta \zeta \mathrm{\Delta }\mathrm{\Phi }_a(x),\delta _B\mathrm{\Phi }_y(x)=\delta \zeta \mathrm{\Delta }\mathrm{\Phi }_y(x),`$ (3.9)
where $`\delta \zeta `$ is an infinitesimal fermionic parameter independent of $`x`$. The BRST transformation of the gauge fields or matter fields is nothing but the infinitesimal gauge transformation with $`\delta \theta _a`$ and $`\delta \theta _y`$ equal to $`g\delta \zeta C_a`$ and $`g_1\delta \zeta C_y`$ respectively. Namely
$`\delta _BW_a^\mu (x)=\delta \zeta \mathrm{\Delta }W_a^\mu (x)=\delta \zeta D_{ab}^\mu C_b(x),`$ (3.10)
$`\delta _BW_y^\mu (x)=\delta \zeta \mathrm{\Delta }W_y^\mu (x)=\delta \zeta ^\mu C_y(x),`$ (3.11)
$`\delta _B\psi (x)=\delta \zeta \mathrm{\Delta }\psi (x),\delta _B\overline{\psi }(x)=\delta \zeta \mathrm{\Delta }\overline{\psi }(x),`$ (3.12)
where
$`D_{ab}^\mu (x)=\delta _{ab}^\mu +gf_{abc}A_c^\mu (x),`$
$`\mathrm{\Delta }\nu _L(x)={\displaystyle \frac{i}{2}}gC_1(x)e_L(x){\displaystyle \frac{1}{2}}gC_2(x)e_L(x){\displaystyle \frac{i}{2}}gC_3(x)\nu _L(x)+{\displaystyle \frac{i}{2}}g_1C_y(x)\nu _L(x),`$
$`\mathrm{\Delta }e_L(x)={\displaystyle \frac{i}{2}}gC_1(x)\nu _L(x)+{\displaystyle \frac{1}{2}}gC_2(x)\nu _L(x)+{\displaystyle \frac{i}{2}}gC_3(x)e_L(x)+{\displaystyle \frac{i}{2}}g_1C_y(x)e_L(x),`$
$`\mathrm{\Delta }e_R(x)=ig_1C_y(x)e_R(x),`$
$`\mathrm{\Delta }\overline{\nu }_L(x)={\displaystyle \frac{i}{2}}gC_1(x)\overline{e}_L(x){\displaystyle \frac{1}{2}}gC_2(x)\overline{e}_L(x)+{\displaystyle \frac{i}{2}}gC_3(x)\overline{\nu }_L(x){\displaystyle \frac{i}{2}}g_1C_y(x)\overline{\nu }_L(x),`$
$`\mathrm{\Delta }\overline{e}_L(x)={\displaystyle \frac{i}{2}}gC_1(x)\overline{\nu }_L(x)+{\displaystyle \frac{1}{2}}gC_2(x)\overline{\nu }_L(x){\displaystyle \frac{i}{2}}gC_3(x)\overline{e}_L(x){\displaystyle \frac{i}{2}}g_1C_y(x)\overline{e}_L(x),`$
$`\mathrm{\Delta }\overline{e}_R(x)=ig_1C_y(x)\overline{e}_R(x).`$
$`C_a(x)`$ and $`C_y(x)`$ are also transformed as usual
$`\delta _BC_a(x)=\delta \zeta \mathrm{\Delta }C_a(x)=\delta \zeta {\displaystyle \frac{g}{2}}C_{abc}C_b(x)C_c(x),`$
$`\delta _BC_y(x)=0.`$
Now we can write $`\mathrm{\Delta }\mathrm{\Phi }_a(x)`$, $`\mathrm{\Delta }\mathrm{\Phi }_y(x)`$ as
$`\mathrm{\Delta }\mathrm{\Phi }_1=_\mu \mathrm{\Delta }W_1^\mu (x)+g_1\mathrm{\Delta }W_2^\mu (x)W_{y\mu }(x)+g_1W_{2\mu }(x)\mathrm{\Delta }W_y^\mu (x),`$ (3.13)
$`\mathrm{\Delta }\mathrm{\Phi }_2=_\mu \mathrm{\Delta }W_2^\mu (x)g_1\mathrm{\Delta }W_1^\mu (x)W_{y\mu }(x)g_1W_{1\mu }(x)\mathrm{\Delta }W_y^\mu (x),`$ (3.14)
$`\mathrm{\Delta }\mathrm{\Phi }_3=_\mu \mathrm{\Delta }W_3^\mu (x)+g_1\mathrm{\Delta }W_3^\mu (x)W_{y\mu }(x)+g_1W_{3\mu }(x)\mathrm{\Delta }W_y^\mu (x),`$ (3.15)
$`\mathrm{\Delta }\mathrm{\Phi }_y=_\mu \mathrm{\Delta }W_y^\mu (x)+g\mathrm{\Delta }W_3^\mu (x)W_{y\mu }(x)+gW_{3\mu }(x)\mathrm{\Delta }W_y^\mu (x),`$ (3.16)
Since $`\mathrm{\Delta }W_a^\mu `$, $`\mathrm{\Delta }W_y^\mu `$, $`\mathrm{\Delta }\psi (x)`$, $`\mathrm{\Delta }\overline{\psi }(x)`$ and $`\mathrm{\Delta }C_a(x)`$ are BRST invariant, it is easy to see that $`\mathrm{\Delta }\mathrm{\Phi }_a(x)`$ and $`\mathrm{\Delta }\mathrm{\Phi }_y(x)`$ are also BRST invariant.
One can further generalized the theory by regarding as new variables the Lagrange multipliers $`\lambda _a(x)`$ and $`\lambda _y(x)`$ associated with the constraint conditions. Thus the total effective Lagrangian and action consist of these Lagrange multipliers, ghosts and the original variables, namely
$`_{\mathrm{eff}}(x)=(x)+^{(C)}(x)+\lambda _a(x)\mathrm{\Phi }_a(x)+\lambda _y(x)\mathrm{\Phi }_y(x),`$ (3.17)
$`I_{\mathrm{eff}}={\displaystyle d^4x_{\mathrm{eff}}(x)}.`$ (3.18)
Correspondingly, the path integral of the generating functional for the Green functions is
$`𝒵[\overline{\eta },\eta ,\overline{\chi },\chi ,J,j]={\displaystyle \frac{1}{N_\lambda }}{\displaystyle 𝒟[\overline{\psi },\psi ,𝒲,\overline{C},C,\lambda ]\mathrm{exp}\left\{\mathrm{i}\left(I_{\mathrm{eff}}+I_s\right)\right\}},`$ (3.19)
where $`N_\lambda `$ is a constant, $`I_s`$ is the source term in the action. They are defined by
$`N_\lambda ={\displaystyle 𝒟[\overline{\psi },\psi ,𝒲,\overline{C},C,\lambda ]\mathrm{exp}\left\{\mathrm{i}I_{\mathrm{eff}}\right\}},`$
$`I_s={\displaystyle }d^4x\{\overline{\eta }(x)\psi (x)+\overline{\psi }(x)\eta (x)+\overline{\chi }_a(x)C_a(x)+\overline{C}_a(x)\chi _a(x)+\overline{\chi }_y(x)C_y(x)`$
$`+\overline{C}_y(x)\chi _y(x)+J_a^\mu (x)W_{a\mu }(x)+J_y^\mu (x)W_{y\mu }(x)+j_a(x)\lambda _a(x)+J_y(x)\lambda _y(x)\},`$ (3.20)
where $`\overline{\eta }(x),\eta (x)\mathrm{}`$ stand for the sources. In particular, $`j_a(x)`$, $`j_y(x)`$ are the sources of $`\lambda _a(x)`$, $`\lambda _y(x)`$, respectively.
We now check the BRST invariance of the effective action $`I_{eff}`$ defined by (3.17) and (3.18). With $`\overline{C}_a(x)`$, $`\overline{C}_y(x)`$ transforming as
$$\delta _B\overline{C}_a(x)=\delta \zeta \lambda _a(x),\delta _B\overline{C}_y(x)=\delta \zeta \lambda _y(x),$$
and noticing the invariance of $`\mathrm{\Delta }\mathrm{\Phi }_a,\mathrm{\Delta }\mathrm{\Phi }_y`$, one has
$$\delta _Bd^4x^{(C)}(x)=d^4x\left\{\lambda _a(x)\delta _B\mathrm{\Phi }_a(x)\lambda _y(x)\delta _B\mathrm{\Phi }_y(x)\right\}.$$
Therefore
$$\delta _BI_{eff}=\delta _BI_{WM}+d^4x\left\{\left(\delta _B\lambda _a(x)\right)\mathrm{\Phi }_a(x)+\left(\delta _B\lambda _y(x)\right)\mathrm{\Phi }_y(x)\right\}.$$
From this and the expression of $`\delta _BI_{WM}`$, it can be shown that the effective action is invariant, when the transformation of $`\lambda _a(x)`$ and $`\lambda _y(x)`$ are defined as
$`\delta _B\lambda _1(x)=\delta \zeta M^2C_1(x),`$
$`\delta _B\lambda _2(x)=\delta \zeta M^2C_2(x),`$
$`\delta _B\lambda _3(x)=\delta \zeta M^2C_3(x)\delta \zeta {\displaystyle \frac{g_1}{g}}M^2C_y(x),`$
$`\delta _B\lambda _y(x)=\delta \zeta {\displaystyle \frac{g_1^2}{g^2}}M^2C_y(x)\delta \zeta {\displaystyle \frac{g_1}{g}}M^2C_3(x).`$
IV. Renormalisability
Let $`W_{a\mu }(x),W_{y\mu }(x)`$, $`C_a(x),C_y(x),\mathrm{}`$ stand for the renormalized field founctions, $`g,g_1`$ and $`M`$ be renormalized parameters. By introducing the source terms of the composite field functions $`\mathrm{\Delta }W_a^\mu `$, $`\mathrm{\Delta }W_y^\mu `$, $`\mathrm{\Delta }C_a(x)`$, $`\mathrm{\Delta }\psi (x)`$, $`\mathrm{\Delta }\overline{\psi }(x)`$ and the sources $`K_\mu ^a(x)`$, $`K_\mu ^y(x)`$, $`L_a(x)`$, $`n_\alpha (x)`$, $`l_\alpha (x)`$, $`p_\alpha (x)`$, $`n_\alpha ^{}(x)`$, $`l_\alpha ^{}(x)`$ and $`p_\alpha ^{}(x)`$, the effective Lagrangian without counterterm becomes
$`_{eff}^{[0]}(x)`$ $`=`$ $`\lambda _a(x)\mathrm{\Phi }_a(x)+\lambda _y(x)\mathrm{\Phi }_y(x)+_{WL}(x)+_{WY}(x)`$ (4.1)
$`+_{WM}(x)+^{(C)}(x)+_\psi (x)+_{\psi W}(x)`$
$`+K_\mu ^a(x)\mathrm{\Delta }W_a^\mu (x)+K_\mu ^y(x)\mathrm{\Delta }W_y^\mu (x)+L_a(x)\mathrm{\Delta }C_a(x)`$
$`+n_\alpha (x)\mathrm{\Delta }\nu _{L\alpha }(x)+l_\alpha (x)\mathrm{\Delta }e_{L\alpha }(x)+p_\alpha (x)\mathrm{\Delta }e_{R\alpha }(x)`$
$`+n_\alpha ^{}(x)\mathrm{\Delta }\overline{\nu }_{L\alpha }(x)+l_\alpha ^{}(x)\mathrm{\Delta }\overline{e}_{L\alpha }(x)+p_\alpha ^{}(x)\mathrm{\Delta }\overline{e}_{R\alpha }(x).`$
The complete effective Lagrangian is the sum of $`_{eff}^{[0]}`$ and the counterterm $`_{count}`$
$`_{\mathrm{eff}}=_{\mathrm{eff}}^{[0]}+_{count}.`$ (4.2)
With (4.1), the generating functional for Green functions is defined as
$`𝒵^{[0]}[\overline{\eta },\eta ,\overline{\chi },\chi ,J,j,K,L,n,l,p,n^{},l^{},p^{}]={\displaystyle \frac{1}{N}}{\displaystyle 𝒟[\overline{\psi },\psi ,𝒲,\overline{C},C,\lambda ]\mathrm{exp}\left\{\mathrm{i}\left(I_{eff}^{[0]}+I_s\right)\right\}},`$ (4.3)
$`I_{eff}^{[0]}`$ is the effective action $`d^4x_{eff}^{[0]}(x)`$, $`N`$ is a constant to make $`𝒵^{[0]}`$ equal to $`1`$ in the absence of
$`I_s={\displaystyle }d^4x\{\overline{\eta }(x)\psi (x)+\overline{\psi }(x)\eta (x)+\overline{\chi }_a(x)C_a(x)+\overline{C}_a(x)\chi _a(x)+\overline{\chi }_y(x)C_y(x)`$
$`+\overline{C}_y(x)\chi _y(x)+J_a^\mu (x)W_{a\mu }(x)+J_y^\mu (x)W_{y\mu }(x)+j_a(x)\lambda _a(x)+j_y(x)\lambda _y(x)\},`$
where $`\overline{\eta }\psi `$ and $`\overline{\psi }\eta `$ stand for
$`\overline{\eta }\psi =\overline{\eta }_\alpha ^{(\nu )}\nu _{L\alpha }+\overline{\eta }_\alpha ^{(l)}e_{L\alpha }+\overline{\eta }_\alpha ^{(r)}e_{R\alpha },`$
$`\overline{\psi }\eta =\overline{\nu }_{L\alpha }\eta _\alpha ^{(\nu )}+\overline{e}_{L\alpha }\eta _\alpha ^{(l)}+\overline{e}_{R\alpha }\eta _\alpha ^{(r)}.`$
Denoting by $`𝒲^{[0]}`$ and $`\mathrm{\Gamma }^{[0]}`$ the generating functionals for connected Green functions and regular vertex functions respectively, one has
$`𝒵^{[0]}=\mathrm{exp}\left\{\mathrm{i}𝒲^{[0]}[\overline{\eta },\eta ,\overline{\chi },\chi ,J,j,K,L,n,l,p,n^{},l^{},p^{}]\right\},`$ (4.4)
$`\mathrm{\Gamma }^{[0]}[\stackrel{~}{\psi },\stackrel{~}{\overline{\psi }},\stackrel{~}{W},\stackrel{~}{\overline{C}},\stackrel{~}{C},\stackrel{~}{\lambda },K,L,n,l,p,n^{},l^{},p^{}]`$
$`=𝒲^{[0]}{\displaystyle }d^4x[J_a^\mu \stackrel{~}{W}_{a\mu }+J_y^\mu \stackrel{~}{W}_{y\mu }+j_a\stackrel{~}{\lambda }_a+j_y\stackrel{~}{\lambda }_y+\overline{\chi }_a\stackrel{~}{C}_a+\stackrel{~}{\overline{C}}_a\chi _a+\overline{\chi }_y\stackrel{~}{C}_y`$
$`+\stackrel{~}{\overline{C}}_y\chi _y+\overline{\eta }^{(\nu )}\stackrel{~}{\nu }_L+\overline{\eta }^{(l)}\stackrel{~}{e}_L+\overline{\eta }^{(r)}\stackrel{~}{e}_R+\stackrel{~}{\overline{\nu }}_L\eta ^{(\nu )}+\stackrel{~}{\overline{e}}_L\eta ^{(l)}+\stackrel{~}{\overline{e}}_R\eta ^{(r)}],`$ (4.5)
where $`\stackrel{~}{W}_{a\mu }`$, $`\stackrel{~}{\nu }_L`$, $`\mathrm{}`$ are the so-called classical fields defined by
$`\stackrel{~}{W}_{a\mu }(x)={\displaystyle \frac{\delta 𝒲^{[0]}}{\delta J_a^\mu (x)}},\stackrel{~}{\lambda }_a(x)={\displaystyle \frac{\delta 𝒲^{[0]}}{\delta j_a(x)}},\stackrel{~}{C}_a(x)={\displaystyle \frac{\delta 𝒲^{[0]}}{\delta \overline{\chi }_a(x)}},`$
$`\stackrel{~}{\overline{C}}_a(x)={\displaystyle \frac{\delta 𝒲^{[0]}}{\delta \chi _a(x)}},\stackrel{~}{W}_{y\mu }(x)={\displaystyle \frac{\delta 𝒲^{[0]}}{\delta J_y^\mu (x)}},\stackrel{~}{\lambda }_y={\displaystyle \frac{\delta 𝒲^{[0]}}{\delta j_y(x)}},`$
$`\stackrel{~}{C}_y(x)={\displaystyle \frac{\delta 𝒲^{[0]}}{\delta \overline{\chi }_y(x)}},\stackrel{~}{\overline{C}}_y(x)={\displaystyle \frac{\delta 𝒲^{[0]}}{\delta \chi _y(x)}},\stackrel{~}{\nu }_{L\alpha }(x)={\displaystyle \frac{\delta 𝒲^{[0]}}{\delta \overline{\eta }_\alpha ^{(\nu )}(x)}},`$
$`\stackrel{~}{e}_{L\alpha }(x)={\displaystyle \frac{\delta 𝒲^{[0]}}{\delta \overline{\eta }_\alpha ^{(l)}(x)}},\stackrel{~}{e}_{R\alpha }(x)={\displaystyle \frac{\delta 𝒲^{[0]}}{\delta \overline{\eta }_\alpha ^{(r)}(x)}},\stackrel{~}{\overline{\nu }}_{L\alpha }(x)={\displaystyle \frac{\delta 𝒲^{[0]}}{\delta \eta _\alpha ^{(\nu )}(x)}},`$
$`\stackrel{~}{\overline{e}}_{L\alpha }(x)={\displaystyle \frac{\delta 𝒲^{[0]}}{\delta \eta _\alpha ^{(l)}(x)}},\stackrel{~}{\overline{e}}_{R\alpha }(x)={\displaystyle \frac{\delta 𝒲^{[0]}}{\delta \eta _\alpha ^{(r)}(x)}}.`$
Therefore
$`J_a^\mu (x)={\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta \stackrel{~}{W}_{a\mu }(x)}},j_a(x)={\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta \stackrel{~}{\lambda }_a(x)}},\overline{\chi }_a(x)={\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta \stackrel{~}{C}_a(x)}},`$
$`\chi _a(x)={\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta \stackrel{~}{\overline{C}}_a(x)}},J_y^\mu (x)={\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta \stackrel{~}{W}_{y\mu }(x)}},j_y(x)={\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta \stackrel{~}{\lambda }_y(x)}},`$
$`\overline{\chi }_y(x)={\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta \stackrel{~}{C}_y(x)}},\chi _y(x)={\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta \stackrel{~}{\overline{C}}_y(x)}},\eta _\alpha ^{(\nu )}(x)={\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta \stackrel{~}{\overline{\nu }}_{L\alpha }(x)}},`$
$`\eta _\alpha ^{(l)}(x)={\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta \stackrel{~}{\overline{e}}_{L\alpha }(x)}},\eta _\alpha ^{(r)}(x)={\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta \stackrel{~}{\overline{e}}_{R\alpha }(x)}},\overline{\eta }_\alpha ^{(\nu )}(x)={\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta \stackrel{~}{\nu }_{L\alpha }(x)}},`$
$`\overline{\eta }_\alpha ^{(l)}(x)={\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta \stackrel{~}{e}_{L\alpha }(x)}},\overline{\eta }_\alpha ^{(r)}(x)={\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta \stackrel{~}{e}_{R\alpha }(x)}}.`$
Besides, for $`K_\mu ^a,L_a`$ $`\mathrm{}`$, the spectators in the Legendre transtrormation, one has
$`{\displaystyle \frac{\delta 𝒲^{[0]}}{\delta K_\mu ^a(x)}}={\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta K_\mu ^a(x)}},{\displaystyle \frac{\delta 𝒲^{[0]}}{\delta K_\mu ^y(x)}}={\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta K_\mu ^y(x)}},{\displaystyle \frac{\delta 𝒲^{[0]}}{\delta L_a(x)}}={\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta L_a(x)}},`$
$`{\displaystyle \frac{\delta 𝒲^{[0]}}{\delta n_\alpha (x)}}={\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta n_\alpha (x)}},{\displaystyle \frac{\delta 𝒲^{[0]}}{\delta l_\alpha (x)}}={\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta l_\alpha (x)}},{\displaystyle \frac{\delta 𝒲^{[0]}}{\delta p_\alpha (x)}}={\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta p_\alpha (x)}},`$
$`{\displaystyle \frac{\delta 𝒲^{[0]}}{\delta n_\alpha ^{}(x)}}={\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta n_\alpha ^{}(x)}},{\displaystyle \frac{\delta 𝒲^{[0]}}{\delta l_\alpha ^{}(x)}}={\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta l_\alpha ^{}(x)}},{\displaystyle \frac{\delta 𝒲^{[0]}}{\delta p_\alpha ^{}(x)}}={\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta p_\alpha ^{}(x)}}.`$
In order to find the Slavnov–Taylor identity satisfied by the generating functional for the regular vertex functions, we change the variables in the path integral of $`𝒵^{[0]}`$ as follows
$`W_a^\mu (x)W_a^\mu (x)+\delta \zeta \mathrm{\Delta }W_a^\mu (x),W_y^\mu (x)W_y^\mu (x)+\delta \zeta \mathrm{\Delta }W_y^\mu (x),`$
$`C_a(x)C_a(x)+\delta \zeta \mathrm{\Delta }C_a(x),C_y(x)C_y(x),`$
$`\overline{C}_a(x)\overline{C}_a(x)\delta \zeta \lambda _a(x),\overline{C}_y(x)\overline{C}_y(x)\delta \zeta \lambda _y(x),`$
$`\psi (x)\psi (x)+\delta \zeta \mathrm{\Delta }\psi (x),\overline{\psi }(x)\overline{\psi }(x)+\delta \zeta \mathrm{\Delta }\overline{\psi }(x),`$
$`\lambda _a(x)\lambda _a(x),\lambda _y(x)\lambda _y(x).`$
The volume element of the path integral does not change and the changes in $`I_s`$ and $`_{WM}`$ lead to
$`{\displaystyle }d^4x\{{\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta K_\mu ^a(x)}}{\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta \stackrel{~}{W}_a^\mu (x)}}+{\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta K_\mu ^y(x)}}{\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta \stackrel{~}{W}_y^\mu (x)}}+{\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta L_a(x)}}{\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta \stackrel{~}{C}_a(x)}}`$
$`+{\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta \stackrel{~}{\nu }_{L\alpha }(x)}}{\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta n_\alpha (x)}}+{\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta \stackrel{~}{e}_{L\alpha }(x)}}{\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta l_\alpha (x)}}+{\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta \stackrel{~}{e}_{R\alpha }(x)}}{\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta p_\alpha (x)}}`$
$`+{\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta \stackrel{~}{\overline{\nu }}_{L\alpha }(x)}}{\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta n_\alpha ^{}(x)}}+{\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta \stackrel{~}{\overline{e}}_{L\alpha }(x)}}{\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta l_\alpha ^{}(x)}}+{\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta \stackrel{~}{\overline{e}}_{R\alpha }(x)}}{\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta p_\alpha ^{}(x)}}`$
$`\stackrel{~}{\lambda }_a(x){\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta \stackrel{~}{\overline{C}}_a(x)}}\stackrel{~}{\lambda }_y(x)+{\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta \stackrel{~}{\overline{C}}_y(x)}}\mathrm{\Delta }_{WM}(x)^{[0]}\}=0,`$ (4.6)
where
$$\mathrm{\Delta }_{WM}(x)^{[0]}=\frac{1}{N𝒵^{[0]}}𝒟[\overline{\psi },\psi ,𝒲,\overline{C},C]\mathrm{\Delta }_{WM}(x)\mathrm{exp}\left\{\mathrm{i}\left(I_{\mathrm{eff}}^{[0]}+I_s\right)\right\}.$$
With the definition of $`\mathrm{\Delta }_{WM}(x)`$
$$\delta _B_{WM}(x)=\delta \zeta \mathrm{\Delta }_{WM}(x),$$
one can write
$`\mathrm{\Delta }_{WM}(x)^{[0]}`$ $`=`$ $`M^2\stackrel{~}{W}_{a\mu }(x){\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta K_\mu ^a(x)}}+M^2\left({\displaystyle \frac{g_1}{g}}\right)^2\stackrel{~}{W}_{y\mu }(x){\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta K_\mu ^y(x)}}`$
$`M^2{\displaystyle \frac{g_1}{g}}\stackrel{~}{W}_{y\mu }(x){\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta K_\mu ^3(x)}}M^2{\displaystyle \frac{g_1}{g}}\stackrel{~}{W}_{3\mu }(x){\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta K_\mu ^y(x)}}.`$
Next, from the invariance of the path integral of $`𝒵^{[0]}`$ with respect to the translation of the integration variables $`\overline{C}_a(x)`$, $`\overline{C}_y(x)`$, $`\lambda _a(x)`$ and $`\lambda _y(x)`$, one can get a set of auxiliary identities
$`{\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta \stackrel{~}{\overline{C}}_1(x)}}_\mu {\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta K_\mu ^1(x)}}g_1\stackrel{~}{W}_{y\mu }{\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta K_\mu ^2(x)}}g_1\stackrel{~}{W}_{2\mu }{\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta K_\mu ^y(x)}}=0,`$ (4.7)
$`{\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta \stackrel{~}{\overline{C}}_2(x)}}_\mu {\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta K_\mu ^2(x)}}+g_1\stackrel{~}{W}_{y\mu }{\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta K_\mu ^1(x)}}+g_1\stackrel{~}{W}_{1\mu }{\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta K_\mu ^y(x)}}=0,`$ (4.8)
$`{\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta \stackrel{~}{\overline{C}}_3(x)}}_\mu {\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta K_\mu ^3(x)}}g_1\stackrel{~}{W}_{y\mu }{\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta K_\mu ^3(x)}}g_1\stackrel{~}{W}_{3\mu }{\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta K_\mu ^y(x)}}=0,`$ (4.9)
$`{\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta \stackrel{~}{\overline{C}}_y(x)}}_\mu {\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta K_\mu ^y(x)}}g\stackrel{~}{W}_{y\mu }{\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta K_\mu ^3(x)}}g\stackrel{~}{W}_{3\mu }{\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta K_\mu ^y(x)}}=0,`$ (4.10)
and
$`{\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta \stackrel{~}{\lambda }_a(x)}}=\mathrm{\Phi }_a(x)^{[0]},{\displaystyle \frac{\delta \mathrm{\Gamma }^{[0]}}{\delta \stackrel{~}{\lambda }_y(x)}}=\mathrm{\Phi }_y(x)^{[0]}.`$ (4.11)
where
$`\mathrm{\Phi }_a(x)^{[0]}={\displaystyle \frac{1}{N𝒵^{[0]}}}{\displaystyle 𝒟[\overline{\psi },\psi ,𝒲,\overline{C},C,\lambda ]\mathrm{\Phi }_a(x)\mathrm{exp}\left\{\mathrm{i}\left(I_{\mathrm{eff}}^{[0]}+I_s\right)\right\}},`$ (4.12)
$`\mathrm{\Phi }_y(x)^{[0]}={\displaystyle \frac{1}{N𝒵^{[0]}}}{\displaystyle 𝒟[\overline{\psi },\psi ,𝒲,\overline{C},C,\lambda ]\mathrm{\Phi }_y(x)\mathrm{exp}\left\{\mathrm{i}\left(I_{\mathrm{eff}}^{[0]}+I_s\right)\right\}}.`$ (4.13)
Let $`\stackrel{~}{\mathrm{\Phi }}_a(x)`$, $`\stackrel{~}{\mathrm{\Phi }}_y(x)`$, $`\stackrel{~}{}_{WM}`$ be the results obtained from $`\mathrm{\Phi }_a(x)`$, $`\mathrm{\Phi }_y(x)`$, $`_{WM}`$ by replacing the field functions with the classical field functions and define
$`\overline{\mathrm{\Gamma }}^{[0]}=\mathrm{\Gamma }^{[0]}{\displaystyle d^4x\left\{\stackrel{~}{\lambda }_a(x)\stackrel{~}{\mathrm{\Phi }}_a(x)+\stackrel{~}{\lambda }_y(x)\stackrel{~}{\mathrm{\Phi }}_y(x)+\stackrel{~}{}_{WM}\right\}},`$ (4.14)
Thus, from (4.6)–(4.11), one gets
$`{\displaystyle }d^4x\{{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}^{[0]}}{\delta K_\mu ^a(x)}}{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}^{[0]}}{\delta \stackrel{~}{W}_a^\mu (x)}}+{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}^{[0]}}{\delta K_\mu ^y(x)}}{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}^{[0]}}{\delta \stackrel{~}{W}_y^\mu (x)}}+{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}^{[0]}}{\delta L_a(x)}}{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}^{[0]}}{\delta \stackrel{~}{C}_a(x)}}`$
$`+{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}^{[0]}}{\delta \stackrel{~}{\nu }_{L\alpha }(x)}}{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}^{[0]}}{\delta n_\alpha (x)}}+{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}^{[0]}}{\delta \stackrel{~}{e}_{L\alpha }(x)}}{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}^{[0]}}{\delta l_\alpha (x)}}+{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}^{[0]}}{\delta \stackrel{~}{e}_{R\alpha }(x)}}{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}^{[0]}}{\delta p_\alpha (x)}}`$
$`+{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}^{[0]}}{\delta \stackrel{~}{\overline{\nu }}_{L\alpha }(x)}}{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}^{[0]}}{\delta n_\alpha ^{}(x)}}+{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}^{[0]}}{\delta \stackrel{~}{\overline{e}}_{L\alpha }(x)}}{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}^{[0]}}{\delta l_\alpha ^{}(x)}}+{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}^{[0]}}{\delta \stackrel{~}{\overline{e}}_{R\alpha }(x)}}{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}^{[0]}}{\delta p_\alpha ^{}(x)}}\}=0.`$ (4.15)
and
$`{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}^{[0]}}{\delta \stackrel{~}{\lambda }_a(x)}}=\mathrm{\Phi }_a(x)^{[0]}\stackrel{~}{\mathrm{\Phi }}_a(x),{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}^{[0]}}{\delta \stackrel{~}{\lambda }_y(x)}}=\mathrm{\Phi }_y(x)^{[0]}\stackrel{~}{\mathrm{\Phi }}_y(x),`$ (4.16)
$`{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}^{[0]}}{\delta \stackrel{~}{\overline{C}}_1(x)}}_\mu {\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}^{[0]}}{\delta K_\mu ^1(x)}}g_1\stackrel{~}{W}_{y\mu }{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}^{[0]}}{\delta K_\mu ^2(x)}}g_1\stackrel{~}{W}_{2\mu }{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}^{[0]}}{\delta K_\mu ^y(x)}}=0,`$ (4.17)
$`{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}^{[0]}}{\delta \stackrel{~}{\overline{C}}_2(x)}}_\mu {\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}^{[0]}}{\delta K_\mu ^2(x)}}+g_1\stackrel{~}{W}_{y\mu }{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}^{[0]}}{\delta K_\mu ^1(x)}}+g_1\stackrel{~}{W}_{1\mu }{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}^{[0]}}{\delta K_\mu ^y(x)}}=0,`$ (4.18)
$`{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}^{[0]}}{\delta \stackrel{~}{\overline{C}}_3(x)}}_\mu {\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}^{[0]}}{\delta K_\mu ^3(x)}}g_1\stackrel{~}{W}_{y\mu }{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}^{[0]}}{\delta K_\mu ^3(x)}}g_1\stackrel{~}{W}_{3\mu }{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}^{[0]}}{\delta K_\mu ^y(x)}}=0,`$ (4.19)
$`{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}^{[0]}}{\delta \stackrel{~}{\overline{C}}_y(x)}}_\mu {\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}^{[0]}}{\delta K_\mu ^y(x)}}g\stackrel{~}{W}_{y\mu }{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}^{[0]}}{\delta K_\mu ^3(x)}}g\stackrel{~}{W}_{3\mu }{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}^{[0]}}{\delta K_\mu ^y(x)}}=0.`$ (4.20)
As $`\mathrm{\Phi }_a(x),\mathrm{\Phi }_y(x)`$ contain the products of the SU<sub>L</sub>(2) and U<sub>Y</sub>(1) fields, (4.16) is complicated unless the generating functional for the Green functions does not contain the sources of the Lagrange multipliers $`\lambda _a`$ and $`\lambda _y`$. Actuarely we are not interested in using the Green functions involving $`\lambda _a`$ or $`\lambda _y`$. Our intention to use the generalized form of the theory containing the sources of these Lagrange multipliers is to study the Renormalisability of the theory for which such sources are absent from the generating functional for the Green functions and therefore $`\mathrm{\Phi }_a(x)^{[0]}`$ and $`\mathrm{\Phi }_y(x)^{[0]}`$ are equal to zero. We now, according to $`(4.11)`$, let vanish $`\frac{\delta \mathrm{\Gamma }^{[0]}}{\delta \stackrel{~}{\lambda }_a(x)}`$ and $`\frac{\delta \mathrm{\Gamma }^{[0]}}{\delta \stackrel{~}{\lambda }_a(x)}`$ to make $`\mathrm{\Phi }_a(x)^{[0]}`$ and $`\mathrm{\Phi }_y(x)^{[0]}`$ equal to zero. This means
$`\stackrel{~}{\mathrm{\Phi }}_a(x)=0,\stackrel{~}{\mathrm{\Phi }}_y(x)=0.`$ (4.21)
and
$`{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}^{[0]}}{\delta \stackrel{~}{\lambda }_a(x)}}=0,{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}^{[0]}}{\delta \stackrel{~}{\lambda }_y(x)}}=0,`$ (4.22)
In the following we will denote by $`\overline{\mathrm{\Gamma }}^{[0]}[\psi ,\overline{\psi },W,\overline{C},C,\lambda ,K,L,n,l,p,n^{},l^{},p^{}]`$ the functional that is obtained from $`\overline{\mathrm{\Gamma }}^{[0]}[\stackrel{~}{\psi },\stackrel{~}{\overline{\psi }},\stackrel{~}{W},\stackrel{~}{\overline{C}},\stackrel{~}{C},\stackrel{~}{\lambda },K,\mathrm{}]`$ by replacing the classical field functions with the usual field functions. Assume that the dimensional regularization method is used and the Slavnov–Taylor identity and the auxiliary identities are guaranteed. Denote the tree part and one loop part of $`\overline{\mathrm{\Gamma }}^{[0]}`$ by $`\overline{\mathrm{\Gamma }}_0^{[0]}`$ and $`\overline{\mathrm{\Gamma }}_1^{[0]}`$ respectively. $`\overline{\mathrm{\Gamma }}_0^{[0]}`$ is thus the modified action $`\overline{I}_{eff}^{[0]}`$ obtained from $`I_{eff}^{[0]}`$ by excluding the mass term and $`(\lambda _a,\lambda _y)`$ terms. From $`(4.15)`$ and $`(4.17)(4.22)`$ one has
$`\mathrm{\Phi }_a(x)=0,\mathrm{\Phi }_y(x)=0,`$ (4.23)
$`{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}^{[0]}}{\delta \lambda _a(x)}}=0,{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}^{[0]}}{\delta \lambda _y(x)}}=0,`$ (4.24)
$`\mathrm{\Lambda }_{op}\overline{\mathrm{\Gamma }}_0^{[0]}=0,`$
and
$`\overline{\mathrm{\Gamma }}_0^{[0]}\overline{\mathrm{\Gamma }}_1^{[0]}+\overline{\mathrm{\Gamma }}_1^{[0]}\overline{\mathrm{\Gamma }}_0^{[0]}=\mathrm{\Lambda }_{op}\overline{\mathrm{\Gamma }}_1^{[0]}=0,`$ (4.25)
$`\mathrm{\Sigma }_a(x)\overline{\mathrm{\Gamma }}^{[0]}=0,\mathrm{\Sigma }_y(x)\overline{\mathrm{\Gamma }}^{[0]}=0.`$ (4.26)
where $`\mathrm{\Lambda }_{op}`$,$`\mathrm{\Sigma }_a(x)`$ and $`\mathrm{\Sigma }_y(x)`$ are defined by
$`\mathrm{\Lambda }_{op}={\displaystyle }d^4x\{{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}_0^{[0]}}{\delta K_\mu ^a(x)}}{\displaystyle \frac{\delta }{\delta W_a^\mu (x)}}+{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}_0^{[0]}}{\delta W_a^\mu (x)}}{\displaystyle \frac{\delta }{\delta K_\mu ^a(x)}}+{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}_0^{[0]}}{\delta K_\mu ^y(x)}}{\displaystyle \frac{\delta }{\delta W_y^\mu (x)}}+{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}_0^{[0]}}{\delta W_y^\mu (x)}}{\displaystyle \frac{\delta }{\delta K_\mu ^y(x)}}`$
$`+{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}_0^{[0]}}{\delta L_a(x)}}{\displaystyle \frac{\delta }{\delta C_a(x)}}+{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}_0^{[0]}}{\delta C_a(x)}}{\displaystyle \frac{\delta }{\delta L_a(x)}}+{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}_0^{[0]}}{\delta \nu _{L\alpha }(x)}}{\displaystyle \frac{\delta }{\delta n_\alpha (x)}}+{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}_0^{[0]}}{\delta n_\alpha (x)}}{\displaystyle \frac{\delta }{\delta \nu _{L\alpha }(x)}}`$
$`+{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}_0^{[0]}}{\delta e_{L\alpha }(x)}}{\displaystyle \frac{\delta }{\delta l_\alpha (x)}}+{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}_0^{[0]}}{\delta l_\alpha (x)}}{\displaystyle \frac{\delta }{\delta e_{L\alpha }(x)}}+{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}_0^{[0]}}{\delta e_{R\alpha }(x)}}{\displaystyle \frac{\delta }{\delta p_\alpha (x)}}+{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}_0^{[0]}}{\delta p_\alpha (x)}}{\displaystyle \frac{\delta }{\delta e_{R\alpha }(x)}}`$
$`+{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}_0^{[0]}}{\delta \overline{\nu }_{L\alpha }(x)}}{\displaystyle \frac{\delta }{\delta n_\alpha ^{}(x)}}+{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}_0^{[0]}}{\delta n_\alpha ^{}(x)}}{\displaystyle \frac{\delta }{\delta \overline{\nu }_{L\alpha }(x)}}+{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}_0^{[0]}}{\delta \overline{e}_{L\alpha }(x)}}{\displaystyle \frac{\delta }{\delta l_\alpha ^{}(x)}}+{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}_0^{[0]}}{\delta l_\alpha ^{}(x)}}{\displaystyle \frac{\delta }{\delta \overline{e}_{L\alpha }(x)}}`$
$`+{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}_0^{[0]}}{\delta \overline{e}_{R\alpha }(x)}}{\displaystyle \frac{\delta }{\delta p_\alpha ^{}(x)}}+{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}_0^{[0]}}{\delta p_\alpha ^{}(x)}}{\displaystyle \frac{\delta }{\delta \overline{e}_{R\alpha }(x)}}\},`$ (4.27)
$`\mathrm{\Sigma }_1(x)={\displaystyle \frac{\delta }{\delta \overline{C}_1(x)}}_\mu {\displaystyle \frac{\delta }{\delta K_\mu ^1(x)}}g_1W_{y\mu }{\displaystyle \frac{\delta }{\delta K_\mu ^2(x)}}g_1W_{2\mu }{\displaystyle \frac{\delta }{\delta K_\mu ^y(x)}},`$ (4.28)
$`\mathrm{\Sigma }_2(x)={\displaystyle \frac{\delta }{\delta \overline{C}_2(x)}}_\mu {\displaystyle \frac{\delta }{\delta K_\mu ^2(x)}}+g_1W_{y\mu }{\displaystyle \frac{\delta }{\delta K_\mu ^1(x)}}+g_1W_{1\mu }{\displaystyle \frac{\delta }{\delta K_\mu ^y(x)}},`$ (4.29)
$`\mathrm{\Sigma }_3(x)={\displaystyle \frac{\delta }{\delta \overline{C}_3(x)}}_\mu {\displaystyle \frac{\delta }{\delta K_\mu ^3(x)}}g_1W_{y\mu }{\displaystyle \frac{\delta }{\delta K_\mu ^3(x)}}g_1W_{3\mu }{\displaystyle \frac{\delta }{\delta K_\mu ^y(x)}},`$ (4.30)
$`\mathrm{\Sigma }_y(x)={\displaystyle \frac{\delta }{\delta \overline{C}_y(x)}}_\mu {\displaystyle \frac{\delta }{\delta K_\mu ^y(x)}}gW_{y\mu }{\displaystyle \frac{\delta }{\delta K_\mu ^3(x)}}gW_{3\mu }{\displaystyle \frac{\delta }{\delta K_\mu ^y(x)}}.`$ (4.31)
The meaning of the notation $`AB`$ is the same as in the common use, namely
$`AB={\displaystyle }d^4x\{{\displaystyle \frac{\delta A}{\delta K_\mu ^a(x)}}{\displaystyle \frac{\delta B}{\delta W_a^\mu (x)}}+{\displaystyle \frac{\delta A}{\delta K_\mu ^y(x)}}{\displaystyle \frac{\delta B}{\delta W_y^\mu (x)}}+{\displaystyle \frac{\delta A}{\delta L_a(x)}}{\displaystyle \frac{\delta B}{\delta C_a(x)}}`$
$`+{\displaystyle \frac{\delta A}{\delta \nu _{L\alpha }(x)}}{\displaystyle \frac{\delta B}{\delta n_\alpha (x)}}+{\displaystyle \frac{\delta A}{\delta e_{L\alpha }(x)}}{\displaystyle \frac{\delta B}{\delta l_\alpha (x)}}+{\displaystyle \frac{\delta A}{\delta e_{R\alpha }(x)}}{\displaystyle \frac{\delta B}{\delta p_\alpha (x)}}`$
$`+{\displaystyle \frac{\delta A}{\delta \overline{\nu }_{L\alpha }(x)}}{\displaystyle \frac{\delta B}{\delta n_\alpha ^{}(x)}}+{\displaystyle \frac{\delta A}{\delta \overline{e}_{L\alpha }(x)}}{\displaystyle \frac{\delta B}{\delta l_\alpha ^{}(x)}}+{\displaystyle \frac{\delta A}{\delta \overline{e}_{R\alpha }(x)}}{\displaystyle \frac{\delta B}{\delta p_\alpha ^{}(x)}}\}.`$ (4.32)
$`(4.24)(4.26)`$ are of course satisfied by the finite part and the pole part of $`\overline{\mathrm{\Gamma }}_1^{[0]}`$. Thus the equations of the pole part $`\overline{\mathrm{\Gamma }}_{1,div}^{[0]}`$ are
$`{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}_{1,div}^{[0]}}{\delta \lambda _a(x)}}=0,{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}_{1,div}^{[0]}}{\delta \lambda _y(x)}}=0,`$ (4.33)
$`\mathrm{\Lambda }_{op}\overline{\mathrm{\Gamma }}_{1,div}^{[0]}=0,`$ (4.34)
$`\mathrm{\Sigma }_a(x)\overline{\mathrm{\Gamma }}_{1,div}^{[0]}=0,\mathrm{\Sigma }_y(x)\overline{\mathrm{\Gamma }}_{1,div}^{[0]}=0.`$ (4.35)
Obviously, the same equations should be found for a SU<sub>L</sub>(2)$`\times `$U<sub>Y</sub>(1) theory without the mass term if the same constraint conditions are chosen.
If $`M=0`$, then it is known from the renormalisability of the theory that $`\overline{\mathrm{\Gamma }}_{1,div}^{[0]}`$ is a combination of the following terms
$`T_{GL}=g{\displaystyle \frac{\overline{\mathrm{\Gamma }}_0^{[0]}}{g}},T_{GY}=g_1{\displaystyle \frac{\overline{\mathrm{\Gamma }}_0^{[0]}}{g_1}},`$
$`T_{WL}={\displaystyle d^4x\left\{W_a^\mu (x)\frac{\delta \overline{\mathrm{\Gamma }}_0^{[0]}}{\delta W_a^\mu (x)}+L_a(x)\frac{\delta \overline{\mathrm{\Gamma }}_0^{[0]}}{\delta L_a(x)}\right\}},`$
$`T_{WY}={\displaystyle d^4xW_y^\mu (x)\frac{\delta \overline{\mathrm{\Gamma }}_0^{[0]}}{\delta W_y^\mu (x)}},`$
$`T_{CK}={\displaystyle d^4x\left\{\overline{C}_a(x)\frac{\delta \overline{\mathrm{\Gamma }}_0^{[0]}}{\delta \overline{C}_a(x)}+C_a(x)\frac{\delta \overline{\mathrm{\Gamma }}_0^{[0]}}{\delta C_a(x)}+K_\mu ^a(x)\frac{\delta \overline{\mathrm{\Gamma }}_0^{[0]}}{\delta K_\mu ^a(x)}\right\}},`$
$`T_{CKY}={\displaystyle d^4x\left\{\overline{C}_y(x)\frac{\delta \overline{\mathrm{\Gamma }}_0^{[0]}}{\delta \overline{C}_(x)}+C_y(x)\frac{\delta \overline{\mathrm{\Gamma }}_0^{[0]}}{\delta C_y(x)}+K_\mu ^y(x)\frac{\delta \overline{\mathrm{\Gamma }}_0^{[0]}}{\delta K_\mu ^y(x)}\right\}},`$
$`T_{\nu L}={\displaystyle d^4x\left\{\nu _{L\alpha }(x)\frac{\delta \overline{\mathrm{\Gamma }}_0^{[0]}}{\delta \nu _{L\alpha }(x)}+\overline{\nu }_{L\alpha }(x)\frac{\delta \overline{\mathrm{\Gamma }}_0^{[0]}}{\delta \overline{\nu }_{L\alpha }(x)}\right\}},`$
$`T_{eL}={\displaystyle d^4x\left\{e_{L\alpha }(x)\frac{\delta \overline{\mathrm{\Gamma }}_0^{[0]}}{\delta e_{L\alpha }(x)}+\overline{e}_{L\alpha }(x)\frac{\delta \overline{\mathrm{\Gamma }}_0^{[0]}}{\delta \overline{e}_{L\alpha }(x)}\right\}},`$
$`T_{eR}={\displaystyle d^4x\left\{e_{R\alpha }(x)\frac{\delta \overline{\mathrm{\Gamma }}_0^{[0]}}{\delta e_{R\alpha }(x)}+\overline{e}_{R\alpha }(x)\frac{\delta \overline{\mathrm{\Gamma }}_0^{[0]}}{\delta \overline{e}_{R\alpha }(x)}\right\}},`$
$`T_{nn^{}}={\displaystyle d^4x\left\{n_\alpha (x)\frac{\delta \overline{\mathrm{\Gamma }}_0^{[0]}}{\delta n_\alpha (x)}+n_\alpha ^{}(x)\frac{\delta \overline{\mathrm{\Gamma }}_0^{[0]}}{\delta n_\alpha ^{}(x)}\right\}},`$
$`T_{ll^{}}={\displaystyle d^4x\left\{l_\alpha (x)\frac{\delta \overline{\mathrm{\Gamma }}_0^{[0]}}{\delta l_\alpha (x)}+l_\alpha ^{}(x)\frac{\delta \overline{\mathrm{\Gamma }}_0^{[0]}}{\delta l_\alpha ^{}(x)}\right\}},`$
$`T_{pp^{}}={\displaystyle d^4x\left\{p_\alpha (x)\frac{\delta \overline{\mathrm{\Gamma }}_0^{[0]}}{\delta p_\alpha (x)}+p_\alpha ^{}(x)\frac{\delta \overline{\mathrm{\Gamma }}_0^{[0]}}{\delta p_\alpha ^{}(x)}\right\}}.`$
With these terms one can form five solutions of equations $`(4.33)(4.35)`$, which can be chosen as
$`T_{(1)}=T_{WL}T_{GL}T_{CK},`$ (4.36)
$`T_{(2)}=T_{WY}T_{GY}T_{CKY},`$ (4.37)
$`T_{(3)}=T_{CK}+T_{CKY}+T_{nn^{}}+T_{ll^{}}+T_{pp^{}},`$ (4.38)
$`T_{(4)}=T_{\nu L}+T_{eL}T_{nn^{}}T_{ll^{}},`$ (4.39)
$`T_{(5)}=T_{eR}T_{pp^{}}.`$ (4.40)
Note that $`T_{(3)}`$ is $`2\left(\overline{\mathrm{\Gamma }}_0^{[0]}I_{WL}I_{WY}I_\psi I_{\psi W}\right)`$. $`T_{(1)}`$ is a combination of $`I_{WL}`$, $`T_{(3)}`$ and $`d^4xC_y(x)\frac{\delta \overline{\mathrm{\Gamma }}_0^{[0]}}{\delta C_y(x)}`$. $`T_{(2)}`$ is a combination of $`I_{WY}`$ and $`d^4xC_y(x)\frac{\delta \overline{\mathrm{\Gamma }}_0^{[0]}}{\delta C_y(x)}`$. The sum of $`T_{(4)}`$ and $`T_{(5)}`$ is $`2\left(I_\psi +I_{\psi W}\right)`$. $`d^4xC_y(x)\frac{\delta \overline{\mathrm{\Gamma }}_0^{[0]}}{\delta C_y(x)}`$ and $`T_{(5)}`$ can be easily checked to satisfy $`(4.34)(4.35)`$. In addition to $`(4.36)`$$`(4.40)`$, a new term appearing in $`\overline{\mathrm{\Gamma }}_{1,\mathrm{div}}^{[0]}`$ when $`M0`$ should includ $`M^2`$ as a factor and also satisfies $`(4.34)(4.35)`$. Only $`I_{WM}`$ can be a candidate. Is such a term can really appear ? Imagine a limiting case that the matter fields and the U$`{}_{Y}{}^{}(1)`$ fields are absent. Thus the constraint conditions become Lorentz conditions and the above five solutions become two, namely, $`\left(T_{WL}T_{GL}\right)`$ and $`T_{CK}`$. This combination of $`T_{WL}`$ and $`T_{GL}`$ are due to the restriction of the constraint condition containing $`^\mu W_{y\mu }`$ and therefore should be decomposed into two independent terms when the U$`{}_{Y}{}^{}(1)`$ fields are absent. In fact, it is known that a SU(n) theory with massive gauge Bosons is renormalisability and that when the matter fields are absent $`\overline{\mathrm{\Gamma }}_{n+1,div}^{[n]}`$ of such a theory is a combination of three independent terms $`T_{WL}`$, $`T_{GL}`$ and $`T_{CK}`$. It follows that for the present theory $`\overline{\mathrm{\Gamma }}_{1,div}^{[0]}`$ does not cantain the mass term $`I_{MW}`$ neither and can be expressed as
$`\overline{\mathrm{\Gamma }}_{1,div}^{[0]}=\alpha _1^{(1)}T_{(1)}+\alpha _2^{(1)}T_{(2)}+\alpha _3^{(1)}T_{(3)}+\alpha _4^{(1)}T_{(4)}+\alpha _5^{(1)}T_{(5)},`$ (4.41)
where, $`\alpha _1^{(1)},\mathrm{},\alpha _5^{(1)}`$ are constants of order $`(\mathrm{})^1`$ and are divergent when the space-time dimension tends to $`4`$.
In order to cancel the one loop divergence the counterterm of order $`\mathrm{}^1`$ in the action should be chosen as
$`\delta I_{count}^{[1]}=\overline{\mathrm{\Gamma }}_{1,div}^{[0]},`$ (4.42)
Since
$`\overline{I}_{eff}^{[0]}=\overline{\mathrm{\Gamma }}_0^{[0]},`$ (4.43)
it is known from $`(4.41)`$ that the sum of $`\overline{I}_{eff}^{[0]}`$ and $`\delta I_{count}^{[1]}`$, to order of $`\mathrm{}^1`$, can be written as
$`\overline{I}_{eff}^{[1]}`$ $`[\psi ,\overline{\psi },W,C,\overline{C},K,L,n,l,p,n^{},l^{},p^{},g,g_1]`$ (4.44)
$`=\overline{I}_{\mathrm{eff}}^{[0]}[\psi ^{[0]},\overline{\psi }^{[0]},W^{[0]},C^{[0]},\overline{C}^{[0]},K^{[0]},L^{[0]},n^{[0]},n^{{}_{}{}^{}[0]},\mathrm{},g^{[0]},g_1^{[0]}],`$
where the bare fields and the bare parameters (to order $`(\mathrm{})^1`$) are defined as
$`W_{a\mu }^{[0]}=(Z_3^{[1]})^{1/2}W_{a\mu }=\left(1\alpha _1^{(1)}\right)W_{a\mu },L_a^{[0]}=(Z_3^{[1]})^{1/2}L_a,`$ (4.45)
$`W_{y\mu }^{[0]}=(Z_3^{{}_{}{}^{}[1]})^{1/2}W_{y\mu }=\left(1\alpha _2^{(1)}\right)W_{y\mu },`$ (4.46)
$`C_a^{[0]}=(\stackrel{~}{Z}_3^{[1]})^{1/2}C_a=\left(1\alpha _3^{(1)}+\alpha _1^{(1)}\right)C_a,`$ (4.47)
$`\overline{C}_a^{[0]}=(\stackrel{~}{Z}_3^{[1]})^{1/2}\overline{C}_a,K_\mu ^{a[0]}=(\stackrel{~}{Z}_3^{[1]})^{1/2}K_\mu ^a,`$ (4.48)
$`C_y^{[0]}=(\stackrel{~}{Z}_3^{{}_{}{}^{}[1]})^{1/2}C_y=\left(1\alpha _3^{(1)}+\alpha _2^{(1)}\right)C_y,`$ (4.49)
$`\overline{C}_y^{[0]}=(\stackrel{~}{Z}_3^{{}_{}{}^{}[1]})^{1/2}\overline{C}_y,K_\mu ^{y[0]}=(\stackrel{~}{Z}_3^{[1]})^{1/2}K_\mu ^y,`$ (4.50)
$`\nu _L^{[0]}=(Z_{\nu L}^{[1]})^{1/2}\nu _L=\left(1\alpha _4^{(1)}\right)\nu _L,\overline{\nu }_L^{[0]}=(Z_{\nu L}^{[1]})^{1/2}\overline{\nu }_L,`$ (4.51)
$`e_L^{[0]}=(Z_{eL}^{[1]})^{1/2}e_L=(Z_{\nu L}^{[1]})^{1/2}e_L,\overline{e}_L^{[0]}=(Z_{eL}^{[1]})^{1/2}\overline{e}_L,`$ (4.52)
$`e_R^{[0]}=(Z_{eR}^{[1]})^{1/2}e_R=\left(1\alpha _5^{(1)}\right)e_R,\overline{e}_R^{[0]}=(Z_{eR}^{[1]})^{1/2}\overline{e}_R,`$ (4.53)
$`n^{[0]}=(Z_{(n)}^{[1]})^{1/2}n=\left(1\alpha _3^{(1)}+\alpha _4^{(1)}\right)n,n^{{}_{}{}^{}[0]}=(Z_{(n)}^{[1]})^{1/2}n^{},`$ (4.54)
$`l^{[0]}=(Z_{(l)}^{[1]})^{1/2}l=(Z_{(n)}^{[1]})^{1/2}l,l^{{}_{}{}^{}[0]}=(Z_{(l)}^{[1]})^{1/2}l^{},`$ (4.55)
$`p^{[0]}=(Z_{(p)}^{[1]})^{1/2}p=\left(1\alpha _3^{(1)}+\alpha _5^{(1)}\right)p,p^{{}_{}{}^{}[0]}=(Z_{(p)}^{[1]})^{1/2}p^{},`$ (4.56)
$`g^{[0]}=Z_g^{[1]}g=(Z_3^{[1]})^{1/2}g,g_1^{[0]}=Z_g^{{}_{}{}^{}[1]}g_1=(Z_3^{{}_{}{}^{}[1]})^{1/2}g_1.`$ (4.57)
Next, defined
$`\mathrm{\Phi }_1^{[0]}=^\mu W_{1\mu }^{[0]}+g_1^{[0]}W_{2\mu }^{[0]}W_y^{\mu [0]},`$
$`\mathrm{\Phi }_2^{[0]}=^\mu W_{2\mu }^{[0]}g_1^{[0]}W_{1\mu }^{[0]}W_y^{\mu [0]},`$
$`\mathrm{\Phi }_3^{[0]}=^\mu W_{3\mu }^{[0]}+g_1^{[0]}W_{3\mu }^{[0]}W_y^{\mu [0]},`$
$`\mathrm{\Phi }_y^{[0]}=^\mu W_{y\mu }^{[0]}+g^{[0]}W_{3\mu }^{[0]}W_y^{\mu [0]}.`$
From $`(4.45),(4.46)`$ and $`(4.57)`$ one has
$$g^{[0]}W_a^{\mu [0]}=gW_a^\mu ,g_1^{[0]}W_y^{\mu [0]}=g_1W_y^\mu ,$$
and
$`\mathrm{\Phi }_a^{[0]}=(Z_3^{[1]})^{1/2}\mathrm{\Phi }_a,\mathrm{\Phi }_y^{[0]}=(Z_3^{{}_{}{}^{}[1]})^{1/2}\mathrm{\Phi }_y.`$ (4.58)
Thus by adding $`I_{WM}`$ and the $`\lambda `$ terms into $`\overline{I}_{eff}^{[1]}`$ and forming
$`I_{eff}^{[1]}=\overline{I}_{eff}^{[1]}+I_{WM}+{\displaystyle d^4x\left\{\lambda _a(x)\mathrm{\Phi }_a(x)+\lambda _y(x)\mathrm{\Phi }_y(x)\right\}},`$ (4.59)
one gets
$`I_{eff}^{[1]}[\psi ,\overline{\psi },W,C,\overline{C},\lambda ,K,L,n,l,p,n^{},l^{},p^{},g,g_1,M]`$
$`=I_{eff}^{[0]}[\psi ^{[0]},\overline{\psi }^{[0]},W^{[0]},C^{[0]},\overline{C}^{[0]},\lambda ^{[0]},K^{[0]},L^{[0]},n^{[0]},n^{{}_{}{}^{}[0]},\mathrm{},g^{[0]},g_1^{[0]},M^{[0]}],`$ (4.60)
where
$`M^{[0]}=(Z_3^{[1]})^{1/2}M,\lambda _a^{[0]}=(Z_3^{[1]})^{1/2}\lambda _a,\lambda _a^{[0]}=(Z_3^{{}_{}{}^{}[1]})^{1/2}\lambda _y.`$ (4.61)
Obviously, if the action $`I_{eff}^{[1]}`$ is used to replace $`I_{eff}^{[0]}`$ in $`(4.3)`$ and define $`𝒵^{[1]}`$, $`\mathrm{\Gamma }^{[1]}`$ as well as
$`\overline{\mathrm{\Gamma }}^{[1]}=\mathrm{\Gamma }^{[1]}I_{WM}{\displaystyle d^4x\left\{\lambda _a(x)\mathrm{\Phi }_a(x)+\lambda _y(x)\mathrm{\Phi }_y(x)+_{WM}\right\}},`$ (4.62)
then one has
$`\overline{\mathrm{\Gamma }}^{[1]}[\psi ,\overline{\psi },W,C,\overline{C},\lambda ,K,L,n,l,p,n^{},l^{},p^{},g,g_1,M]`$
$`=\overline{\mathrm{\Gamma }}^{[0]}[\psi ^{[0]},\overline{\psi }^{[0]},W^{[0]},C^{[0]},\overline{C}^{[0]},\lambda ^{[0]},K^{[0]},L^{[0]},n^{[0]},n^{{}_{}{}^{}[0]},\mathrm{},g^{[0]},g_1^{[0]},M^{[0]}].`$ (4.63)
From this it is easy to check that, to order $`\mathrm{}^1`$, $`\overline{\mathrm{\Gamma }}^{[1]}`$ is finite. Moreover, by changing into bare fields and bare parameters the fields and parameters in $`(4.15)`$ $`(4.22)`$ and then transforming them back into the renormalized fields and renormalized parameters according to $`(4.45)`$$`(4.59)`$, one can see that, under condition $`(4.23)`$, $`\overline{\mathrm{\Gamma }}^{[1]}`$ also satisfies
$`\mathrm{\Lambda }_{op}\overline{\mathrm{\Gamma }}^{[1]}=0,`$ (4.64)
$`{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}^{[1]}}{\delta \lambda _a(x)}}=0,{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}^{[1]}}{\delta \lambda _y(x)}}=0,`$ (4.65)
$`\mathrm{\Sigma }_a(x)\overline{\mathrm{\Gamma }}^{[1]}=0,\mathrm{\Sigma }_y(x)\overline{\mathrm{\Gamma }}^{[1]}=0.`$ (4.66)
It is now clear that the renormalisability of the theory can be verified by the inductive method. The following is an outline of the proof. Assume that up to $`n`$ loop the theory has been proved to be renormalisable by introducing the counterterm
$$I_{\mathrm{count}}^{[n]}=\underset{l=1}{\overset{n}{}}\delta I_{\mathrm{count}}^{[l]},$$
where $`\delta I_{count}^{[l]}`$ is the counterterm of order $`\mathrm{}^l`$ and has the form of (4.41),(4.42). Therefore the modified generating functional $`\overline{\mathrm{\Gamma }}^{[n]}`$ for the regular vertex, defined by the action
$$I_{\mathrm{eff}}^{[n]}=I_{\mathrm{eff}}^{[0]}+I_{\mathrm{count}}^{[n]}$$
satisfied equations $`(4.64)(4.66)`$ (under $`(4.23)`$) and, to order $`\mathrm{}^n`$, is finite. This also means that the fields or parameters in each of the following brackets have the same renormalization factor:
$$(W_{a\mu }^{[0]},L_a),(C_a,\overline{C}_a,K_\mu ^a),(C_y,\overline{C}_y,K_\mu ^y),(\nu _L,\overline{\nu }_L,e_L,\overline{e}_L),(e_R,\overline{e}_R),(n,n^{},l,l^{}),(p,p^{}),(\lambda ,M,g),$$
and that
$`Z_g^{{}_{}{}^{}[n]}(Z_3^{{}_{}{}^{}[n]})^{1/2}=1,Z_g^{[n]}(Z_3^{[n]})^{1/2}=1,`$
$`Z_3^{[n]}\stackrel{~}{Z}_3^{[n]}=\stackrel{~}{Z}_3^{{}_{}{}^{}[n]}\stackrel{~}{Z}_3^{{}_{}{}^{}[n]}=Z_{\nu L}^{[n]}Z_{(n)}^{[n]}=Z_{eR}^{[n]}Z_{(p)}^{[n]}.`$
We have to proved that by using a counterterm of order $`\mathrm{}^{n+1}`$ which also has the form of (4.41),(4.42), $`\overline{\mathrm{\Gamma }}^{[n+1]}`$ can be make satisfy $`(4.64)(4.66)`$ and finite to order $`\mathrm{}^{n+1}`$, where $`\overline{\mathrm{\Gamma }}^{[n+1]}`$ is the modified generating functional for the regular vertex, determined by the action
$$I_{\mathrm{eff}}^{[n+1]}=I_{\mathrm{eff}}^{[n]}+\delta I_{\mathrm{count}}^{[n+1]}.$$
Denote by $`\overline{\mathrm{\Gamma }}_k^{[n]}`$ the part of order $`\mathrm{}^k`$ in $`\overline{\mathrm{\Gamma }}^{[n]}`$. For $`kn`$, $`\overline{\mathrm{\Gamma }}_k^{[n]}`$ is equal to $`\overline{\mathrm{\Gamma }}_k^{[k]}`$, because it can not contain the contribution of a counterterm of order $`\mathrm{}^{k+1}`$ or higher. Thus on expanding $`\overline{\mathrm{\Gamma }}^{[n]}`$ to order $`\mathrm{}^{n+1}`$ one has
$$\overline{\mathrm{\Gamma }}^{[n]}=\underset{k=0}{\overset{n}{}}\overline{\mathrm{\Gamma }}_k^{[k]}+\overline{\mathrm{\Gamma }}_{n+1}^{[n]}+\mathrm{}.$$
Using this and extracting the terms of order $`\mathrm{}^{(n+1)}`$ from the equations satisfied by $`\overline{\mathrm{\Gamma }}^{[n]}`$, namely $`(4.64)(4.66)`$, one finds
$`\mathrm{\Lambda }_{op}\overline{\mathrm{\Gamma }}_{n+1}^{[n]}=0,`$ (4.67)
$`{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}_{n+1}^{[n]}}{\delta \lambda _a(x)}}=0,{\displaystyle \frac{\delta \overline{\mathrm{\Gamma }}_{n+1}^{[n]}}{\delta \lambda _y(x)}}=0,`$ (4.68)
$`\mathrm{\Sigma }_a(x)\overline{\mathrm{\Gamma }}_{n+1}^{[n]}=0,\mathrm{\Sigma }_y(x)\overline{\mathrm{\Gamma }}_{n+1}^{[n]}=0,`$ (4.69)
Let $`\overline{\mathrm{\Gamma }}_{n+1,div}^{[n]}`$ stand for the pole part of $`\overline{\mathrm{\Gamma }}_{n+1}^{[n]}`$. By repeating the steps going from $`(4.33)`$ to $`(4.41)`$, one can arrive at
$`\overline{\mathrm{\Gamma }}_{n+1,div}^{[n]}=\alpha _1^{(n+1)}T_{(1)}+\alpha _2^{(n+1)}T_{(2)}+\alpha _3^{(n+1)}T_{(3)}+\alpha _4^{(n+1)}T_{(4)}+\alpha _5^{(n+1)}T_{(5)},`$ (4.70)
where $`\alpha _1^{(n+1)},\mathrm{},\alpha _5^{(n+1)}`$ are constants of order $`(\mathrm{})^{n+1}`$. Therefore, in order to cancel the $`n+1`$ loop divergence the counterterm of order $`\mathrm{}^{n+1}`$ should be chosen as
$`\delta I_{count}^{[n+1]}=\overline{\mathrm{\Gamma }}_{n+1,div}^{[n]}[\psi ,\overline{\psi },W,C,\overline{C}].`$ (4.71)
Adding this counterterm, the mass term and the $`\lambda `$ terms to $`\overline{I}_{\mathrm{eff}}^{[n]}`$, one can express the effective action of order $`\mathrm{}^{n+1}`$ as
$`I_{\mathrm{eff}}^{[n+1]}[\psi ,\overline{\psi },W,C,\overline{C},\lambda ,K,L,n,l,p,n^{},l^{},p^{},g,g_1,M]`$
$`=I_{\mathrm{eff}}^{[0]}[\psi ^{[0]},\overline{\psi }^{[0]},W^{[0]},C^{[0]},\overline{C}^{[0]},\lambda ^{[0]},K^{[0]},L^{[0]},n^{[0]},n^{{}_{}{}^{}[0]},\mathrm{},g^{[0]},g_1^{[0]},M^{[0]}],`$ (4.72)
where the bare fields and the bare parameters (to order $`(\mathrm{})^{n+1}`$) are defined as
$`W_{a\mu }^{[0]}=(Z_3^{[n+1]})^{1/2}W_{a\mu }=\left((Z_3^{[n]})^{1/2}\alpha _1^{(n+1)}\right)W_{a\mu },L_a^{[0]}=(Z_3^{[n+1]})^{1/2}L_a,`$ (4.73)
$`W_{y\mu }^{[0]}=(Z_3^{{}_{}{}^{}[n+1]})^{1/2}W_{y\mu }=\left((Z_3^{{}_{}{}^{}[n]})^{1/2}\alpha _2^{(n+1)}\right)W_{y\mu },`$ (4.74)
$`C_a^{[0]}=(\stackrel{~}{Z}_3^{[n+1]})^{1/2}C_a=\left((\stackrel{~}{Z}_3^{[n]})^{1/2}+(\alpha _3^{(n+1)}+\alpha _1^{(n+1)})\right)C_a,`$ (4.75)
$`\overline{C}_a^{[0]}=(\stackrel{~}{Z}_3^{[n+1]})^{1/2}\overline{C}_a,K_\mu ^{a[0]}=(\stackrel{~}{Z}_3^{[n+1]})^{1/2}K_\mu ^a,`$ (4.76)
$`C_y^{[0]}=(\stackrel{~}{Z}_3^{{}_{}{}^{}[n+1]})^{1/2}C_y=\left((\stackrel{~}{Z}_3^{{}_{}{}^{}[n]})^{1/2}+(\alpha _3^{(n+1)}+\alpha _2^{(n+1)})\right)C_y,`$ (4.77)
$`\overline{C}_y^{[0]}=(\stackrel{~}{Z}_3^{{}_{}{}^{}[n+1]})^{1/2}\overline{C}_y,K_\mu ^{y[0]}=(\stackrel{~}{Z}_3^{[n+1]})^{1/2}K_\mu ^y,`$ (4.78)
$`\nu _L^{[0]}=(Z_{\nu L}^{[n+1]})^{1/2}\nu _L=\left((Z_{\nu L}^{[n]})^{1/2}\alpha _4^{(n+1)}\right)\nu _L,\overline{\nu }_L^{[0]}=(Z_{\nu L}^{[n+1]})^{1/2}\overline{\nu }_L,`$ (4.79)
$`e_L^{[0]}=(Z_{eL}^{[n+1]})^{1/2}e_L=(Z_{\nu L}^{[n+1]})^{1/2}e_L,\overline{e}_L^{[0]}=(Z_{eL}^{[n+1]})^{1/2}\overline{e}_L,`$ (4.80)
$`e_R^{[0]}=(Z_{eR}^{[n+1]})^{1/2}e_R=\left((Z_{eR}^{[n]})^{1/2}\alpha _5^{(n+1)}\right)e_R,\overline{e}_R^{[0]}=(Z_{eR}^{[n+1]})^{1/2}\overline{e}_R,`$ (4.81)
$`n^{[0]}=(Z_{(n)}^{[n+1]})^{1/2}n=\left((Z_{(n)}^{[n]})^{1/2}+(\alpha _3^{(n+1)}+\alpha _4^{(n+1)})\right)n,n^{{}_{}{}^{}[0]}=(Z_{(n)}^{[n+1]})^{1/2}n^{},`$ (4.82)
$`l^{[0]}=(Z_{(l)}^{[n+1]})^{1/2}l=(Z_{(n)}^{[n+1]})^{1/2}l,l^{{}_{}{}^{}[0]}=(Z_{(l)}^{[n+1]})^{1/2}l^{},`$ (4.83)
$`p^{[0]}=(Z_{(p)}^{[n+1]})^{1/2}p=\left((Z_{(p)}^{[n]})^{1/2}\alpha _3^{(n+1)}+\alpha _5^{(n+1)}\right)p,p^{{}_{}{}^{}[0]}=(Z_{(p)}^{[n+1]})^{1/2}p^{},`$ (4.84)
$`g^{[0]}=Z_g^{[n+1]}g=(Z_3^{[n+1]})^{1/2}g,g_1^{[0]}=Z_g^{{}_{}{}^{}[n+1]}g_1=(Z_3^{{}_{}{}^{}[n+1]})^{1/2}g_1,`$ (4.85)
$`g^{[0]}=Z_g^{[n+1]}g=(Z_3^{[n+1]})^{1/2}g,g_1^{[0]}=Z_g^{{}_{}{}^{}[n+1]}g_1=(Z_3^{{}_{}{}^{}[n+1]})^{1/2}g_1,`$ (4.86)
$`M^{[0]}=Z_M^{[n+1]}M=(Z_3^{[n+1]})^{1/2}M,`$ (4.87)
and $`\lambda _a^{[0]},\lambda _y^{[0]}`$ are
$`\lambda _a^{[0]}=(Z_3^{[n+1]})^{1/2}\lambda _a,\lambda _y^{[0]}=(Z_3^{{}_{}{}^{}[n+1]})^{1/2}\lambda _y.`$ (4.88)
Therefore, in terms of such bare fields and bare parameters, $`\overline{\mathrm{\Gamma }}^{[n+1]}`$ can be expressed as
$`\overline{\mathrm{\Gamma }}^{[n+1]}`$ $`[W,C,\overline{C},\psi ,\overline{\psi },K,L,n,l,p,n^{},l^{},p^{},g,g_1,M]`$ (4.89)
$`=\overline{\mathrm{\Gamma }}^{[0]}[W^{[0]},C^{[0]},\overline{C}^{[0]},\psi ^{[0]},\overline{\psi }^{[0]},K^{[0]},L^{[0]},n^{[0]},n^{{}_{}{}^{}[0]},\mathrm{},g^{[0]},g_1^{[0]},M^{[0]}].`$
From this one can conclude that $`\overline{\mathrm{\Gamma }}^{[n+1]}`$, under $`(4.23)`$, satisfies $`(4.64)`$$`(4.66)`$ and is finite to order $`\mathrm{}^{n+1}`$. Since the theory can be renormalized to one loop the renormalisability has been proven.
V. Concluding Remarks
By taking into account the original constraint conditions and the additional condition we have carried out the quantization of the SU<sub>L</sub>(2) $`\times `$ U<sub>Y</sub>(1) electroweak theory with the W Z mass term and construct the ghost action in a way similar to that used for the massive SU(n) theory . We have also shown that when the $`\delta `$ functions appearing in the path integral of the Green functions and representing the constraint conditions are rewritten as Fourier integrals with Lagrange multipliers $`\lambda _a`$ and $`\lambda _y`$, the total effective action consisting of the Lagrange multipliers, ghost fields and the original fields is BRST invariant. Furthermore, by comparing with the massless theory and with the massive SU(n) theory we have found the general form of the divergent part of the generating functional for the regular vertex functions and proven the renormalisability of the theory. It has also been clarified that the renormalisability of the theory with the W Z mass term is ensured by the renormalisability of the massless theory and the massive SU(n) theory.
If the harmlessness of the W Z mass term had been proven at the begining of 1960s, the SU<sub>L</sub>(2) $`\times `$ U<sub>Y</sub>(1) electroweak theory without the Higgs mechanism would have been deeply studied and tested. Today, the standard model of the electroweak theory has achieved great successes and the whereabouts of the Higgs Bosons is still unknown. It is therefore reasonable to ask if such successes really depends on the Higgs mechanism and to pay attention to the theory without the Higgs mechanism.
ACKNOWLEDGMENTS
We are grateful to Professor Yang Li-ming for helpful discussions. This work was supported by National Natural Science Foundation of China and supported in part by Doctoral Programm Foundation of the Institution of Higher Education of China.
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# Possible Origin of Antimatter Regions in the Baryon Dominated Universe
## I Introduction
The statement that our Universe is baryon asymmetrical as a whole is quite firmly established observational fact of contemporary cosmology. Indeed, if large regions of matter and antimatter coexist, then annihilations would take place at the borders between them. If the typical size of such a domain was small enough, then the energy released by these annihilations would result in a diffuse $`\gamma `$ –ray background, in distortions of the spectrum of the cosmic microwave radiation and light element abundance, neither of which is observed (see for review e.g. ). Recent analysis of this problem for baryon symmetric Universe demonstrates that the size of regions should exceed $`1000`$ Mpc., being comparable with the modern cosmological horizon. It therefore seems that the Universe is fundamentally matter–antimatter asymmetric. However the arguments used in do not exclude the case when the Universe is composed almost entirely of matter with relatively small insertions of primordial antimatter. Thus we may expect the existence of macroscopically large antimatter regions in the Universe, that differs drastically from the case of baryon symmetric Universe. We call the region filled with antimatter in the baryon dominated Universe, as antizillah. Of course the existence of antizillahs is not rigorous requirement of baryosynthesis, but some modification of baryogenesis scenarios will result in formation of domains with different sign of baryon charge (see for example ). The only condition which is necessary to satisfy is the amount of antibaryons within antizillahs must be small comparing to the total baryon number of the Universe.
At the first glance it is not difficult to have some amount of antizillahs if we simply suppose that in the early Universe when the baryon excess is generated the C–and CP–violation have different sign in different space regions . This may be achieved, for example, in models with two different sources of CP–violation, explicit and spontaneous one. However, any spontaneous CP– violation processes are a result of early phase transition of first or second order what implies very small size of primordial antizillahs . For example if the antizillahs are formed in the second order phase transition, their size at the moment of formation is determined by $`l_i1/(\lambda T_c)`$, where $`T_c`$ is so called Ginsburg temperature (the critical temperature at which the phase transition take place) and $`\lambda `$ is the selfinteraction coupling constant of field which breaks CP symmetry . In the result of expansion the modern sizes of domains would reach $`l_0l_i(T_c/T_0)=1/(\lambda T_0)10^{21}pc/\lambda `$, where $`T_0`$ is the present temperature of the background radiation.
On the other hand it has been revealed that the average displacement of the antizillah’s boundary caused by annihilation with surrounding matter is about $`0.5pc`$ at the end of radiation dominated (RD) epoch. Therefore any primordial antizillah having initial size up to $`0.5pc`$ or more at the end of RD stage is survived to the contemporary epoch and in the case of successive homogeneous expansion has the size $`1kpc`$ or more. Any primordial antizillah with scale less then critical survival size $`l_c1kpc`$ at contemporary epoch must be eaten up by annihilation process. Thus it is the serious problem which any model with thermal phase transition encounters to create primordial antizillah with the size exceeding the critical survival size $`l_c`$ to avoid complete annihilation.
There is an additional problem for baryosynthesis with surviving antizillah’s sizes. The point is that any phase transition is accompanied by formation of topological defects. If we blow–up the region with different signs of charge symmetry, we automatically blow–up the scale of respective topological defect structure. If the structure decays sufficiently late in the observable part of Universe, the contribution of energy density of such topological defects could be sufficiently high to contradict with observations. It can be easily estimated that the structure with the scale corresponding to the survival size enters the horizon and starts to decay at $`T0.1MeV`$, i.e. in the period of Big bang nucleosynthesis. To remove these unwanted relics sufficiently early it is necessary to have a mechanism for symmetry restoration. This mechanism implies that the baryogenesis is going on within rather narrow time interval .
In the present paper we have elaborated the issue for inhomogeneous baryosynthesis without the difficulties pointed above. The proposed approach is based on the mechanism of spontaneous baryogenesis . This mechanism implies the existence of complex scalar field carrying baryonic charge with explicitly broken $`U(1)`$ symmetry. The baryon/antibaryon number excess is produced, when the phase of this additional field moves along the valley of its potential .
It is supposed that the vacuum energy responsible for inflation is driven by any scalar inflaton field, and additional complex field coexists with the inflaton. Due to the fact that vacuum energy during inflational period is too large, the tilt of potential is vanished. This implies that the phase of the field behaves as ordinary massless Nambu–Goldstone (NG) boson and the radius of NG potential is firmly established by the scale of spontaneous $`U(1)`$ symmetry breaking. Owing to quantum fluctuations of massless field at the de Sitter background the phase is varied in different regions of the Universe. When the vacuum energy decreases the tilt of potential becomes topical, and pseudo NG (PNG) field starts oscillate. As the field rolls down in one direction during the first oscillation, it preferentially creates baryons over antibaryons, while the opposite is true as it rolls down in the opposite direction. Thus to have globally baryon dominated Universe one must have the phase sited in the point, corresponding to the positive baryon excess generation, just at the beginning of inflation (when the size of the modern Universe crosses the horizon). Then subsequent quantum fluctuations can move the phase to the appropriate position causing the antibaryon excess production. If it takes place not too late after the inflation begins, the size of antizillah may exceed the critical surviving size $`l_c`$.
The main idea of proposed issue is based on the existence of quantum fluctuations along the effectively massless angular direction of baryonic charged scalar field. Thus, more general, the considered issue of generation of antizillahs is applicable practically to all mechanisms of baryogenesis where the number density and sign of baryon asymmetry depend on the angular component of complex scalar field. The advantage of the mechanism of spontaneous baryogenesis considered here is the quite simple unambiguous inflation dynamics of scalar field generated baryon charge. This fact allows to establish quantitatively definite relationship between the effects of inflation and generation of baryon (antibaryon) excess in inhomogeneous baryogenesis. However, this relationship may be too rigid for the realistic model of antimatter domain formation compatible with the whole set of astrophysical constraints. The consistent picture may need more sophisticated scenarios. The principal possibility for such scenario can be considered on the base of Affleck Dine (AD) baryogenesis mechanism that still receives a lot of attention .
AD baryogenesis also involves the cosmological evolution of effective scalar field, which carries baryonic charge, being composed of supersymmetric partners of electrically neutral quark and lepton combinations. The important feature of supersymmetric extensions of standard model is the existence of ”flat directions” in field space, on which the scalar potential vanishes . We will refer for the definiteness to the flat directions of minimal standard supersymmetric model (MSSM) . Thus, if the some component of scalar field lies along a flat direction, this component can be considered as a free massless complex scalar so called AD field . At the level of renormalizable terms, ”flat directions” are generic, but supersymmetry breaking and nonrenormalizable operators lift the ”flat directions” and sets the scale for their potential. During the inflational period the AD field develops non–zero vacuum expectation value and subsequently when the Hubble rate becomes of the order of the curvature of AD potential, the condensate starts to oscillate around its present minimum. Baryon asymmetry can be induced in such condensate only if there exists phase shift between real and imaginary parts of the AD field. Such shift and consequently B and CP violation is provided by A–term in the potential which parameterizes MSSM ”flat direction” . The resulting sign and number density of baryon asymmetry depends on the magnitude of initial phase of AD field and on phase shift created by A-term at the relaxation period . Therefore the de–Sitter fluctuations can generate antizillahs in the baryon asymmetric Universe in the similar way to the spontaneous baryogenesis if the angular direction of AD field is characterized the mass that is much smaller that the Hubble constant $`H`$ during inflation. It takes place if there are no of order $`H`$ corrections to the A–term .
The early dynamics of AD field is quite complicated owing to the non–trivial background energy density driving inflation in MSSM. Moreover AD potential can get corrections from the vacuum energy that removes its minimum from the original one . In general there are two types of inflation in MSSM, D–term or F–term inflation (see for review ), depending on the type of vacuum contributing the energy density during de-Sitter stage. In the case of D–term inflation AD fields and inflaton slow roll coherently (in the absence of order $`H^2`$ corrections to the mass squared term of AD potential). It implies that the radius of effectively massless angular AD direction is determined by the immediate value of inflaton field. For the case of F–term inflation the AD scalar will get an order $`H^2`$ negative mass squared term causing the minimum of AD potential. The AD field is closed to the minimum during the F–inflation stage and this minimum determines the radius of circle valley of effectively massless angular direction.
The conclusion from this explicit example based on the MSSM is following. For any complicated inflation dynamics of baryon charged field it is possible to simulate appropriate massless direction that behaves similar to the circle valley of NG potential. This fact makes the proposed issue for generation of antizillahs viable not only for spontaneous baryogenesis mechanism, but for the all mechanisms dealing with effectively massless angular directions during inflation .
The paper is organized as to following. In section II we discuss the quantum behavior of nondominant $`U(1)`$ symmetric scalar field at the inflation period. We estimate the amplitude and space scale of fluctuations of the phase for this field without PNG tilt. The size distribution of these fluctuations determines the size distribution of antizillahs. The section III contains calculations of baryon/antibaryon net excess production at the relaxation of phase when the tilt of Mexican hat potential becomes topical. We summarize our conclusions and discuss some problems of the considered scenarios in section IV.
## II Phase Distribution for NG Field at The Inflation Period
We start our consideration with the discussion of evolution of $`U(1)`$ symmetric scalar field which coexists with inflaton at the inflation epoch. The quantum fluctuations of such field during the inflation stage cause the perturbations for the phase marking the Nambu–Goldstone vacuum. In our model this phase determines the sign and value of baryon excess, so the size distribution of domains containing the appropriate phase values, caused by that fluctuations, coincide with the size distribution of antizillahs.
Thus to estimate the number density of antimatter regions with sizes exceeding the critical survival size $`l_c`$ in the baryogenesis model under consideration we have to deal with long – wave quantum fluctuations of the NG boson field at the inflation period. Various aspects of this question have been examined in the numerous papers in the connection with cosmology of invisible axion. Also the de–Sitter quantum fluctuations have been analyzed in the framework of AD baryogenesis .
The effective potential of the complex field is taken in the usual form
$$V(\chi )=m_\chi ^2\chi ^{}\chi +\lambda _\chi (\chi ^{}\chi )^2+V_0,$$
(1)
where the field $`\chi `$ can be represented in the form
$$\chi =\frac{f}{\sqrt{2}}\mathrm{exp}\left(\frac{i\alpha }{f}\right)$$
(2)
The $`U(1)`$ symmetry breaking implies that the radial component of the field $`\chi `$ acquires a nonvanishing classical part, $`f=m_\chi /\sqrt{\lambda _\chi }`$ and field $`\alpha `$ in eq. (2) becomes a massless NG scalar field with a vanishing effective potential, $`V(\alpha )=0`$. In this case, $`\chi `$ has the familiar Mexican–hat potential, and the degenerated vacua correspond to the circle of radius $`f`$. Throughout present paper we deal with dimensionless angular field $`\theta =\alpha /f`$.
We concern here the possibility to store appropriate phase value in the domain with the size exceeding the critical survival size. Such value of phase plays the role of starting point for clockwise movement, which is going to generate antibaryon excess when the tilt of potential breaking $`U(1)`$ explicitly, will turn to be topical.
We assume that the Hubble constant varies slowly during inflation. Also we use well established behavior of quantum fluctuations on the de Sitter background . It implies that vacuum fluctuations of every scalar field grow exponentially in the inflating Universe. When the wavelength of a particular fluctuation becomes greater than $`H^1`$ the average amplitude of this fluctuation freezes out at some nonzero value because of the large friction term in the equation of motion of the scalar field, whereas its wavelength grows exponentially. In the other words such a frozen fluctuation is equivalent to the appearance of classical field that does not vanish after averaging over macroscopic space intervals. Because the vacuum must contain fluctuations of every wavelengths, inflation leads to the creation of more and more new regions containing the classical field of different amplitudes with scale greater than $`H^1`$. The averaged amplitude of such NG field fluctuations generated during each time interval $`H^1`$ is given by
$$\delta \alpha =\frac{H}{2\pi }$$
(3)
During such time interval the universe expands by a factor of $`e`$. Since the NG field is massless during inflation period (the PNG tilt is vanish yet), one can see that the amplitude of each frozen fluctuation is not changed in time at all and the phases of each wave are random. Thus the quantum evolution of NG field looks like one–dimensional Brownian motion along the circle valley corresponding to the bottom of NG potential. This statement implies that the values of the phase $`\theta `$ in different regions become different, and the corresponding variance grows as
$$(\delta \theta )^2=\frac{H^3t}{4\pi ^2f^2}$$
(4)
that means that dispersion grows as $`\sqrt{(\delta \theta )^2}=\frac{H}{2\pi f}\sqrt{N}`$, where N is the number of e–folds. In the other words the phase $`\theta `$ makes quantum step with the scale $`\frac{H}{2\pi f}`$ at each e–fold, and the total number of steps during some time interval $`\mathrm{\Delta }t`$ is given by $`N=H\mathrm{\Delta }t`$.
Let us consider the scale $`k^1=H_0^1=3000h^1Mpc`$ which is the biggest cosmological scale of interest. We suppose that Universe is baryon asymmetric in this scale which leaves the horizon at definite e–fold $`N=N_{max}`$. On the other side this scale is the one entering the horizon now, namely $`a_{max}H_{max}=a_0H_0`$ where the subscript $`0`$ indicates the contemporary epoch. This implies that:
$$N_{max}=\mathrm{ln}\frac{a_{end}H_{end}}{a_0H_0}\mathrm{ln}\frac{H_{end}}{H_{max}}$$
(5)
the subscript $`end`$ denotes the epoch at the end of inflation. The slow-roll paradigm tells us that the last term of (5) is usually $`1`$. The first term depends on the evolution of scale factor $`a`$ between the end of slow-roll inflation and the present epoch. Assuming that inflation ends by short matter dominated period, which is followed by RD stage lasting until the present matter dominated era begins, one has
$$N_{max}=62\mathrm{ln}\frac{10^{16}GeV}{\sqrt{H_{end}M_p}}\frac{1}{3}\mathrm{ln}\frac{\sqrt{H_{end}M_p}}{\rho _{reh}^{1/4}},$$
(6)
where $`\rho _{reh}^{1/4}`$ is the reheating temperature when the RD stage is established. With $`H_{end}10^{13}GeV`$ and instant reheating this gives $`N_{max}62`$, the largest possible value. However, if one has to invoke supersymmetry to prevent the flatness of the inflation potential, for example like as in the case of AD baryogenesis, the $`\rho _{reh}^{1/4}`$ should not exceed then $`10^{10}GeV`$ to avoid too many gravitino overproduction , and one have $`N_{max}=58`$, perhaps the biggest reasonable value. Through the paper we will use $`N_{max}=60`$. The smallest cosmological scale of antizillah that is survived after annihilation is $`k_c^1=l_c8h^2kpc`$ . It is $`9`$ order of magnitude smaller then $`H_0^1`$, that corresponds to
$$N_cN_{max}133\mathrm{ln}h45$$
(7)
Thus the $`l_c`$ should left horizon at 45–folds before the end of inflation.
Let us assume that the phase value $`\theta =0`$ corresponds to South Pole of NG field circle valley, and $`\theta =\pi `$ corresponds to the opposite pole. The positive gradient of phase in this picture is routed as anticlockwise direction, and the dish of PNG potential would locate at the South Pole of circle (see fig.1). It will be shown below (see section III) that the antibaryon production corresponds to the regions that would contain phase values caused anticlockwise rolling of PNG field $`\alpha `$ during the first half period of oscillation. If the field $`\alpha `$ rolls clockwise towards the dish of tilted potential just after the start of first oscillation then baryon production will take place.
Now we are in the position to estimate the fraction of the Universe containing antizillahs. To ensure that the Universe would be baryon asymmetric as a whole it is necessary to suppose that the phase average value $`\theta =\theta _{60}`$ within biggest cosmological scale of interest emerging at the $`N_{max}=60`$ e–folds before the end of inflation is located in the range $`[0,\pi ]`$. The $`\theta _{60}`$ is the starting point for Brownian motion of the phase value along the circle valley during inflation. As it has been mentioned above, the phase makes Brownian step $`\delta \theta =\frac{H}{2\pi f}`$ at each e–fold. Because the typical wavelength of the fluctuation $`\delta \theta `$ generated during such timescale is equal to $`H^1`$, the whole domain $`H^1`$, containing $`\theta _{60}`$, after one e–fold effectively becomes divided into $`e^3`$ separate, causal disconnected domains of radius $`H^1`$. Each domain contains almost homogeneous phase value $`\theta _{601}=\theta _{60}\pm \delta \theta `$. In half of these domains the phase evolves towards $`\pi `$ (the North Pole) and in the other domains it moves towards zero (the South Pole). To have antizillah with appropriate sizes to avoid full annihilation one should require that the phase value crosses $`\pi `$ or zero not later then after $`15`$ steps. Only in this case the antizillahs would have the sizes larger than $`l_c`$ and are conserved up to the modern era. This means that one of the two following inequality must be satisfied
$$\pi \frac{15H}{2\pi f}\theta _{60}\frac{15H}{2\pi f}$$
(8)
Consider initially the case of exact equalities in expression (8) when the main part of antimatter is contained in the antizillahs of size $`l_c`$. The number of domains containing the equal values of phase at the $`45`$ e–folds before the end of inflation is given by the following expression
$$n_{45}(e^3/2)^{15}10^{15}.$$
(9)
Then the probability that every domain of size $`l_c`$ would not be separated into $`e^3`$ domains with size one order of magnitude less then $`l_c`$ at the next e–fold is given by $`P_s(1/2)^{e^3}10^6`$. Thus the number of domains serving as the prototypes for antizillahs of size $`l_c`$ looks like
$$\overline{n}=n_{45}P_s10^9$$
(10)
There are about $`10^{11}`$ galaxies in the Universe. Thus, according to such simple consideration, we reveal that $`1\%`$ of volume boxes corresponding to each galaxy contains the region of size $`l_c`$ filled with antimatter of highest possible antibaryonic density if the $`\theta _{60}`$ coincides with left side of inequality (8) or lowest one in the case if the opposite equality is held.
We are able also to find the size distribution for antizillahs. For this purpose it is necessary to study the inhomogeneities of phase induced by (3). It has been well established that for any given scale $`l=k^1`$ large scale component of the phase value $`\theta `$ is distributed in accordance with Gauss’s law . The quantity which will be especially interesting for us is the dispersion (4) for quantum fluctuations of phase with moments from $`k=H^1`$ to $`k_{min}=l_{max}^1`$ (where the $`l_{max}`$ is the biggest cosmological scale that corresponds to $`60`$ e–folds). This quantity can be expressed in the following manner
$$\sigma _l^2=\frac{H^2}{4\pi ^2}\underset{k_{min}}{\overset{k}{}}d\mathrm{ln}k=\frac{H^2}{4\pi ^2}\mathrm{ln}\frac{l_{max}}{l}=\frac{H^2}{4\pi ^2f^2}(60N_l),$$
(11)
where $`N_l`$ is the number of e–folds which relates the biggest cosmological scale to the given scale $`l`$. This means that the distribution of phase has the Gaussian form
$$P(\theta _l,l)=\frac{1}{\sqrt{2\pi }\sigma _l}\mathrm{exp}\left\{\frac{(\theta _{60}\theta _l)^2}{2\sigma _l^2}\right\}$$
(12)
Suppose that at e–fold $`N_t`$ before the end of inflation the volume $`V(\overline{\theta },N_t)`$ has been filled with phase value $`\overline{\theta }`$. Then at the e–fold $`N_{t+\mathrm{\Delta }t}=N_t\mathrm{\Delta }N`$ the volume filled with phase $`\overline{\theta }`$ will follow iterative expression
$`V(\overline{\theta },N_{t+\mathrm{\Delta }t})=e^3V(\overline{\theta },N_t)+`$ (13)
$`+(V_U(N_t)e^3V(\overline{\theta },N_t)P(\overline{\theta },N_{t+\mathrm{\Delta }t})h.`$ (14)
Here the $`V_U(N_t)e^{3N_t}H^3`$ is the volume of the Universe at $`N_t`$ e–fold. Expression (13) makes it possible to calculate the size distributions of domains filled with appropriate value of phase numerically. In order to illustrate quantitatively the number distribution of domains, we present here the numerical results for definite values of $`\theta _{60}`$, $`\overline{\theta }`$ and $`h=\frac{H}{2\pi f}`$. The table contains the results concerning to number of domains with average phase $`\overline{\theta }`$ at e-fold number $`N`$,
The fraction of the Universe filled with phase $`\overline{\theta }`$ appears to be equal to $`7.694\times 10^9`$. Thus we see that the distribution of domains with size is very abrupt and should be peaked at smallest value of size. Adjusting the free parameters $`\theta _{60}`$ and $`h`$ we are able to achieve the situation that volume box corresponding to each galaxy contains ($`1÷10`$) regions with appropriate phase $`\overline{\theta }`$. The sizes of such regions are larger or equal to critical surviving size. In spite of the sufficiently large total number of antizillahs only the small part of our Universe will be occupied by antizillahs (see the last line in the presented table).
The nontrivial question on the actual forms of astrophysical objects antizillahs can have in the modern Universe needs spacial analysis, which, in general, strongly depends on the assumed form of the nonbaryonic dark matter, dominating in the period of galaxy formation. However, based on the early analysis the two extreme cases can be specified, when the evolution of antizillahs is not strongly influenced by the dark matter content. In the first case, the antibaryon density within the antizillah is by an order of magnitude higher than the average baryon density, so that the over-density inside this region can exceed the dark matter density and rapid evolution of such an antizillah with the size exceeding the surviving scale can provide formation of compact antimatter stellar system (globular cluster (see for review )) which can survive in galaxy . The other extreme case is antizillah with extremely low internal antibaryon density $`\mathrm{\Omega }_{\overline{B}}<10^5`$. Then the diffused antiworld is realized, when no compact antimatter objects are formed and antizillahs evolve into low density antiproton-positron plasma regions in voids outside the galaxies .
## III Spontaneous Baryogenesis Mechanism
The following element of our scenario of inhomogeneous baryogenesis should contain a conversion of the phase $`\theta `$ into baryon/antibaryon excess. We consider the ansatz of spontaneous baryogenesis mechanism . The basic feature of this mechanism is that the sign of baryon charge created by relaxation of energy of PNG field critically depends on the direction that the phase is rotated on the bottom of Mexican heat potential. It provides us to convert the domains containing the appropriate phase value, caused by fluctuations, to the antizillahs at the period when the NG potential gets the tilt.
The one of reasonable issue to the spontaneous baryogenesis has been considered in the work . Let us briefly discuss it. It was assumed that in the early Universe a complex scalar field $`\chi `$ coexists with inflaton $`\varphi `$ responsible for inflation. This field $`\chi `$ has non vanishing baryon number. The possible interaction of $`\chi `$ that violates lepton number can be described by following Lagrangian density (see e.g. )
$`L=_\mu \chi ^{}^\mu \chi V(\chi )+i\overline{Q}\gamma ^\mu _\mu Q+i\overline{L}\gamma ^\mu _\mu L`$ (15)
$`m_Q\overline{Q}Qm_L\overline{L}L+(g\chi \overline{Q}L+h.c.)`$ (16)
The fields $`Q`$ and $`L`$ could represent heavy quark and lepton, coupled to the ordinary quark and lepton matter fields. Since fields $`\chi `$ and $`Q`$ possess baryon number while the field $`L`$ does not, the couplings in the (15) violate lepton number . The $`U(1)`$ symmetry that corresponds to baryon number is expressed by following transformations
$$\chi \mathrm{exp}(i\beta )\chi ,Q\mathrm{exp}(i\beta )Q,LL$$
(17)
The effective Lagrangian density for $`\theta `$, $`Q`$ and $`L`$ eventually has the following form after symmetry breaking
$`L={\displaystyle \frac{f^2}{2}}_\mu \theta ^\mu \theta +i\overline{Q}\gamma ^\mu _\mu Q+i\overline{L}\gamma ^\mu _\mu L`$ (18)
$`m_Q\overline{Q}Qm_L\overline{L}L+({\displaystyle \frac{g}{\sqrt{2}}}f\overline{Q}L+h.c.)+_\mu \theta \overline{Q}\gamma ^\mu Q`$ (19)
At the energy scale $`\mathrm{\Lambda }<<f`$, the symmetry (17) is explicitly broken and the Mexican–hat circle gets a little pseudo NG tilt described by the potential
$$V(\alpha )=\mathrm{\Lambda }^4(1\mathrm{cos}\theta )$$
(20)
This potential, of high $`2\mathrm{\Lambda }^4`$, has a unique minimum at $`\theta =0`$. Of course, in the most cases, the potential (20) is the lowest–order approximation to a more complicated expressions emerged from particle physics models (see e.g. and Refs. therein).
The important parameter for spontaneous baryogenesis is the curvature of (20) in the vicinity of its minimum, which is determined by the mass of PNG field
$$m_\theta ^2=\frac{\mathrm{\Lambda }^4}{f^2}$$
(21)
As it was mentioned above the field $`\chi `$ is an additional field with nondominant energy density contribution to the Habble constant deriving by de Sitter stage. Suppose that the tilt was formed during inflation. Then the order of magnitude estimation for fluctuations induced by large– scale inhomogeneity of oscillations of the field $`\chi `$ gives $`\frac{\delta T}{T}=\frac{1}{3}\frac{\delta \rho }{\rho }(\mathrm{\Lambda }/T)^4`$. Thus, for $`T=H/2\pi `$ and reasonable value $`\mathrm{\Lambda }10^5H`$ (see the end of this section) the thermal electromagnetic background fluctuations are within the observational limits.
Also we assume that the field $`\theta `$ behaves as massless NG field during inflation implying that the condition
$$m_\theta <<H$$
(22)
is valid, where the $`H`$ is the Hubble constant during the inflation. After the end of inflation condition (22) is violated and the oscillations of field $`\theta `$ around the minimum of potential (20) are started. The energy density $`\rho _\theta \theta _i^2m_\theta ^2f^2`$ of the PNG field which has been created by quantum fluctuations of $`\theta `$ during the inflation converts to baryons and antibaryons . The sign of baryon charge depends on the initial value of phase from which the oscillations are started.
Let us estimate the number of baryons and antibaryons produced by classical oscillations of field $`\theta `$ with an arbitrary initial phase $`\theta _i`$. The appropriate expression for the density of produced baryons (antibaryons) $`n_{B(\overline{B})}`$ is represented in
$$n_{B(\overline{B})}=\frac{g^2}{\pi ^2}\underset{m_Q+m_L}{\overset{\mathrm{}}{}}\omega 𝑑\omega \left|_{\mathrm{}}^{\mathrm{}}𝑑t\chi (t)e^{\pm 2i\omega t}\right|^2,$$
(23)
that is valid if $`\chi (t\mathrm{})=\chi (t+\mathrm{})=0.`$ General case can be obtained in the limits $`\chi (t\mathrm{})0;\chi (t+\mathrm{})=0`$ without loss of generality. After integration by part expression (23) has the form
$$N_{B(\overline{B})}=\frac{g^2}{4\pi ^2}\mathrm{\Omega }_{\theta _i}\underset{m_Q+m_L}{\overset{\mathrm{}}{}}𝑑\omega \left|\underset{\mathrm{}}{\overset{\mathrm{}}{}}𝑑\tau \dot{\chi }(\tau )e^{\pm 2i\omega \tau }\right|^2,$$
(24)
where the $`\mathrm{\Omega }_{\theta _i}`$ is the volume containing the phase value $`\theta _i`$. Here the surface terms appear to be zero at $`t=\mathrm{}`$ due to asymptotic of field $`\chi `$ and at $`t=\mathrm{}`$ due to Feynman radiation conditions.
For our estimations it is enough to accept that the phase changes as
$$\theta (t)\theta _i(1m_\theta t)$$
(25)
during first oscillation. We also put $`m_Q=m_L=0`$ that is reasonable for estimations. Substituting (25) and (2) into (24) we come to
$$N_{B(\overline{B})}\frac{g^2f^2m_\theta }{8\pi ^2}\mathrm{\Omega }_{\theta _i}\theta _i^2\underset{\frac{\theta _i}{2}}{\overset{\mathrm{}}{}}𝑑\stackrel{~}{\omega }\frac{\mathrm{sin}^2\stackrel{~}{\omega }}{\stackrel{~}{\omega }^2},$$
(26)
where the sign in the lower limit of integral corresponds to baryon or antibaryon net excess generation respectively. The reasonability of our approximation follows from comparison of (26) at small $`\theta _i<<1`$
$$N_BN_{\overline{B}}=\frac{g^2f^2m_\theta }{8\pi ^2}\mathrm{\Omega }_{\theta _i}\theta _i^3$$
(27)
with the result of .
Using for spatially homogeneous field $`\chi =\frac{f}{\sqrt{2}}e^{i\theta }`$ the expression for baryon charge
$$Q=i(\chi ^{}d\chi /dtd\chi ^{}/dt\chi )=fd\theta /dt,$$
(28)
one can easily conclude that $`Q>0`$ if $`\theta >0`$ during classical movement of phase $`\theta `$ to zero. Thus the anticlockwise rotation gives rise to antibaryon excess while the clockwise rotation to the baryon excess one.
During reheating, the inflaton energy converts into the radiation. It is assumed that reheating takes place when the Mexican–hat potential is not sensitive to the PNG tilt yet. This implies that the total decay width of inflaton $`\mathrm{\Gamma }_{tot}`$ into light degrees of freedom exceeds the mass $`m_\theta `$. In the other words the reheating is going on under the condition (22). The relaxation of $`\theta `$ field starts when $`Hm_\theta `$ and converts to the baryons or antibaryons. Baryonic charge is converted inside a comoving volume after reheating owing to very effective decay during the cosmological time. This means that the baryon–to–entropy ratio in $`n_{B(\overline{B})}/s=Const`$ in the course of expansion. The entropy density after thermalization is given by
$$s=\frac{2\pi ^2}{45}g_{}T^3$$
(29)
where $`g_{}`$ is the total effective massless degrees of freedom. Here we concern with the temperature above the electroweak symmetry breaking scale. At this temperature all the degrees of freedom of the standard model are in equilibrium and $`g_{}`$ is at least equal to $`106.75`$. The temperature is connected with expansion rate as follow
$$T=\sqrt{\frac{m_pH}{1.66g_{}^{1/2}}}=\frac{\sqrt{m_pm_\theta }}{g_{}^{1/4}}$$
(30)
The last part of expression (30) takes into account that the relaxation starts under the condition $`Hm_\theta `$. Using the formulas (26), (29), (30) we are able to get the baryon/antibaryon asymmetry
$$\frac{n_{B(\overline{B})}}{s}=\frac{45g^2}{16\pi ^4g_{}^{1/4}}\left(\frac{f}{m_p}\right)^{3/2}\frac{f}{\mathrm{\Lambda }}Y(\theta _i)$$
(31)
The function $`Y(\theta )=\theta ^2\underset{\theta /2}{\overset{\theta /2}{}}𝑑\omega \frac{\mathrm{sin}^2\omega }{\omega ^2}`$ takes into account the dependence of amplitude of baryon asymmetry and its sign on the initial phase value in the different space regions during inflation.
The expression (31) allows us to get the observable baryon asymmetry of the Universe as a whole $`n_B/s310^{10}`$. In the model under consideration we have supposed initially that $`fH10^6m_p`$. The natural value of coupling constant is $`g10^2`$. We are coming to observable baryon asymmetry at quite reasonable condition $`f/\mathrm{\Lambda }10^5`$ (see e.g. ).
## IV Discussion
In this paper we have proposed a model for inhomogeneous baryosynthesis on the base of the spontaneous baryogenesis mechanism . The model predicts the generation of antizillahs with sizes exceeding the critical surviving size. The antibaryon number density relative to background baryon density in the resulting antizillahs and its number depends on the value of phase established at the beginning and on the parameters of PNG field potential. It is possible to have one or several antizillahs the volume box corresponding to every galaxy depending on the parameter values. The observational consequences of existence of antizillahs and the restrictions on their number and sizes have been analyzed in papers
Of course we may in general expect that some region with size exceeding $`l_c`$ would contain antibaryon excess after the annihilation of small primordial domains and antidomains contained in this region is completed. However the probability to have such region is suppressed exponentially. Therefore to have observational acceptable number of antimatter regions with the size exceeding the critical survival size, a superluminous cosmological expansion in the formation of primordial antimatter proto–domain seems necessary.
As we have mentioned, the additional problem for the most models of inhomogeneous baryogenesis invoking phase transitions at the inflation epoch is prediction of the large scale unwanted topological defects. Our scheme contains the premise for existence of domain walls too. Such walls are not formed when the only minimum of PNG potential exists, what corresponds in the considered model to the fluctuations around $`\theta =0`$, when the North pole ($`\theta =\pi `$) is not crossed. But in the case, when such crossing takes place the multiple degeneracy of vacua appears (e.g. vacua with $`\theta =0`$ and $`\theta =2\pi `$). The equation of motion that correspondes to potential (20) admits kink–like, domain wall solution, which interpolates between two adjacent vacua. Thus, when the PNG tilt is significant, domain wall is formed along the closed surface (e.g. $`\theta =\pi `$) . In the other words every antizillah with high relative antibaryon density will be encompassed by domain wall bag. The wall stress energy $`\mathrm{\Delta }8f\mathrm{\Lambda }^2`$ leads to the oscillation of wall bag after the whole bag enters the cosmological horizon. During the oscillations the energy stored in the walls is released in the form of quanta of NG field and gravitational waves. As we are taken $`0<\theta _{60}<\pi `$, the wall’s bag will have the scale of the order of modern horizon, if the dispersion $`\sigma _{l_{max}}`$ is large as $`\pi \theta _{60}`$. Owing to very large oscillation period such big wall bag would presumably survive to the present time, which would be cosmological disaster . Thus the upper limit on the dispersion will be $`\sigma _{60}<\pi `$. From the other hand this condition should be valued if we want to have parameters of antizillah population that do not contradict to direct and indirect observational constraints . It means that we will have wall bags with the sizes less then cosmological horizon and that walls had to decay until present time. The mechanisms of their decay is a subject of separate paper, in which we also plan to obtain additional constraints on the model, which follow from the condition that walls do not dominate within the cosmological horizon before the bag decays. If the energy density of walls is sufficiently high to give local wall dominance in the border region before the bag enters the horizon, it is of interest to analyze the role of superluminous expansion in the border regions in the bag evolution (see e.g. ). The interesting question on the wall interaction with baryons in the course of wall contraction and decay will be also studied separately.
In general all baryogenesis models that are able to generate some amount of antimatter regions look like radical limit of models with local baryon number density fluctuations so called isocurvature fluctuations . It is known that the contribution of isocurvature fluctuations to the cosmic microwave background (CMB) anisotropy obeys to $`\frac{\delta T}{T}=\frac{1}{3}\frac{\mathrm{\Omega }_B}{\mathrm{\Omega }_0}\delta _{B_i}`$, where $`\delta _{B_i}`$ is the amplitude of initial baryon number fluctuations and $`\mathrm{\Omega }_0`$ ($`\mathrm{\Omega }_B`$) are the total (baryon) density (in units of critical density). As it follows from our numerical illustration (see II and expression (26)) we must have quite large amplitude of initial baryon number fluctuations $`\delta _{B_i}h/\theta _{60}10^2`$ at the biggest cosmological scales, and consequently we would have large amplitude of isocurvature fluctuations at large scales that contradicts with COBE measurements .
To be keeping away of the problem of large–scale isocurvature fluctuations, we can, for example, prevent the fluctuations of phase at largest cosmological scales. The point is that to have antizillah with size exceeding few kpc. we do not need to start phase fluctuations at the e–folds that correspond to the biggest cosmological scales. It is sufficiently to start fluctuations of phase from the moment, for instance, when the scale $`8h^1Mpc`$ leaves Habble horizon during inflation, namely after the $`6.2`$ e–folds from the beginning of inflation. We took this scale, because it is known that at the scale less then $`8h^1Mpc`$ we could be generated initial baryon number fluctuations at the level $`\delta _{B_i}10^2÷10^3`$ without any contradictions with observations.
One of the natural way to prevent the phase fluctuations at the early inflation is to keep $`U(1)`$ symmetry restored during first $`7`$ e–folds. The mechanism that is able to restore symmetry during inflation has been consider in the works . According to that works we can introduce interaction between inflaton field $`\varphi `$ and field $`\chi `$. The simple potential of such kind may be chosen as $`V(\varphi ,\chi )=\frac{1}{4}\lambda _\varphi \varphi ^4+V(\chi )+\nu \varphi ^2\chi ^{}\chi `$, where $`\nu =m_\chi ^2/cM_p^2`$, and $`c1`$. The effective mass of the field $`\chi `$ depends on $`\varphi `$ directly $`m_\chi ^2(\varphi )=m_\chi ^2+\nu \varphi ^2`$. One considers here for simplicity the case $`\nu =m_\chi ^2/cM_p^2`$. This implies that the effective value of mass $`m_\chi ^2(\varphi )`$ during inflation is given by $`\nu (\varphi ^2cM_p^2)`$ and is positive because of very large value of the inflation field. It means that our $`U(1)`$ symmetry is restored during the period when the amplitude of the inflaton field exceeds $`\varphi _c=\sqrt{c}M_P`$, and the field $`\chi `$ settles into the minimum of its symmetric potential. During this period there was no NG boson valley and phase fluctuations. After the moment that inflaton field turns to be less then $`\varphi _c`$ the symmetry breaking takes place and the NG potential has the radius $`f_{eff}=\sqrt{\nu (cM_p^2\varphi ^2)/\lambda _\chi }`$ and fluctuations are started. To keep symmetry restored during first $`7`$ e–folds we should have $`\varphi _c=4M_p`$. After the moment of symmetry breaking it is allowed to start the fluctuations of phase with appropriate dispersion to create antizillahs, without any contradictions with observed CMB anisotropy. Of course to evaluate the distribution of antizillahs by sizes we have to take another parameters then we have used in our numerical example, but it does not change the main result of this paper.
Another story will take place if we would like to consider the AD baryogenesis as a basis for generation of antizillahs.
As it was discussed in the introduction the dynamics of the AD field is more complicated that in the case of spontaneous baryogenesis. Moreover it depends on the fact, D– or F– term inflation takes place. Also some details depend on the dimension ($`d=4,6..`$) of non–renormalizable term lifting the flat direction , but it is enough for the brief discussion to circumscribe ourself with the minimal AD baryogenesis , where $`d=4`$. Thus in the case of D– term inflation, when the coherent slow rolling of AD field and inflaton are already established, the maximal radius $`f_{eff}^{AD(D)}10^{16}GeV`$ of effectively massless angular direction can be obtained from the requirement that radial de Sitter fluctuations of AD field would not disturb significantly the spectral index of primordial adiabatic density perturbations measured by COBE. Thereby, it is possible to get dispersion of phase fluctuations at the level $`h10^2`$ that is required for successful generation of antizillahs. The similar situation we could have in the case of F– term inflation because the AD potential gets an order of $`H^2`$ negative mass squared term during inflation, which causes the effective minimum at $`f_{eff}^{AD(F)}C_F\sqrt{Hm_p}10^{16}GeV`$ (the $`C_F`$ is a constant of order of one).
The isocurvature fluctuations in the model of inhomogeneous AD baryogenesis with dispersion of phase fluctuations appropriate for antizillahs generation should be already observed by COBE . Moreover this fluctuations can get some amplification owing to possible transformation of fluctuations of AD condensate into the isocurvature fluctuations of neutralinos . The exact solution of the problem of isocurvature fluctuations for the AD based antimatter generation is the subject of separate investigation. Here we can only present some speculations, how to avoid the large isocurvature fluctuations at large cosmological scales, which are based on the similar strategy that has been chosen in the case of spontaneous baryogenesis.
As it has been mentioned in the Introduction, to organize the angular effectively massless direction in the AD potential we should accept the condition of the absence of order $`H`$ correction to the A– term both during and after inflation . This condition gets automatically satisfied in the case of D– term inflation , while it is not true if the inflation is F– term dominated (see for example ). According to this observation we can hope to find the such kind of trajectory of inflaton in field space that corresponds to the F– term dominated inflation in the beginning and then goes into D– term dominated regime. It implies that during the F– term dominated inflation the angular direction gets a mass of order $`H`$ and imaginary component of AD field is dumped and exponentially close to the minimum caused by this effective mass term. In such situation there are no de Sitter fluctuations of the phase. The fluctuations start only at the moment when the inflation goes to the D– term dominated regime and the angular direction turns to be effectively massless, because there is no correction of order $`H`$ to the A– term anymore. As we estimated before, to put the maximal scale of isocurvature fluctuations far below the modern cosmological horizon the transition from F– term to D– term inflation should take place 5–10 e–folds after the beginning of inflation. How to organize such transition is the subject of separate publication, but it seems that it could appear, for example, in the context of a realistic supergravity theory deriven from the weak coupled supestring , which is already beyond the MSSM. There is some possibility to generate the F– term from a Fayet–Iliopoulos D– term . It could preserve the flatness of F– term direction during the first 5–10 e–folds of inflation causing the F– term domination firstly and subsequent trasformation of the vacuum energy into the D– term domination mode when it is allowed to begin phase fluctuations of AD field with dispersion appropriate for generation of antizillahs and without contradictions with COBE measurements.
We would like to notice in conclusion that the regions with antimatter in matter–dominated Universe could arise naturally in the variety of models. The main issue, that is needed, is a valley of potential. It is the valleys that are responsible for formation of causally separated regions with different values of field which in its turn give rise to antimatter domains. Many extensions of standard model based on supersymmetry possess this property, what strongly extends the physical basis for cosmic antimatter searches.
## V Acknowledgments
This work was partially performed in the framework of Section ”Cosmoparticle physics” of Russian State Scientific Technological Program ”Astronomy. Fundamental Space Research”, International project Astrodamus, Cosmion–ETHZ and AMS–EPCOS. MYuK and ASS acknowledge supporte from Khalatnikov–Starobinsky school (grant 00–15–96699). We thank R.Konoplich and A.Sudarikov for interesting discussions and suggestions. We are also grateful to J.Ulbricht, A.D.Linde and I.Tkachev for useful comments.
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# Twisted Traces of Quantum Intertwiners and Quantum Dynamical R-Matrices Corresponding to Generalized Belavin-Drinfeld Triples
## 1 Introduction
This paper is a continuation of \[ES1\] and \[EV2\].
In \[EV2\], A.Varchenko and the first author considered weighted traces of products of intertwining operators for quantum groups $`U_q(𝔤)`$, where $`𝔤`$ is a simple Lie algebra. They showed that the generating function $`F_{V_1,\mathrm{}V_N}(\lambda ,\mu )`$ of such traces (where $`\lambda ,\mu `$ are complex weights for $`𝔤`$) satisfies four commuting systems of difference equations – the Macdonald-Ruijsenaars (MR) system, the quantum Knizhnik-Zamolodchikov-Bernard (qKZB) system, the dual MR system, and the dual qKZB system. The first two systems are systems of difference equations with respect to $`\lambda `$, which involve Felder’s trigonometric dynamical R-matrix depending of $`\lambda `$. The second two systems are systems of difference equations respect to $`\mu `$, which are obtained from the first two by the transformation $`\lambda \mu ,V_iV_{Ni+1}^{}`$. Such a symmetry is explained by the fact that the function $`F_{V_1,\mathrm{}V_N}(\lambda ,\mu )`$ is invariant under this transformation.
If the quantum group $`U_q(𝔤)`$ is replaced with the Lie algebra $`𝔤`$, these results are replaced with their classical analogs (\[EV2\]). Namely, the MR and qKZB equations are replaced by the classical MR and KZB equations, which are differential equations involving Felder’s classical trigonometric dynamical r-matrix. The dual MR and KZB equations retain roughly the same form, but involve the rational quantum dynamical R-matrix rather than the trigonometric one. Thus, the symmetry between $`\lambda `$ and $`\mu `$ is destroyed.
In \[ES1\], we generalized the classical MR and KZB equations to the case when the trace is twisted using a ”generalized Belavin-Drinfeld triple”, i.e. a pair of subdiagrams $`\mathrm{\Gamma }_1,\mathrm{\Gamma }_2`$ of the Dynkin diagram of $`𝔤`$ together with an isomorphism $`T:\mathrm{\Gamma }_1\mathrm{\Gamma }_2`$ between them. It turned out that such twisted traces also satisfy differential equations which involve a dynamical r-matrix, namely the one attached to the triple $`(\mathrm{\Gamma }_1,\mathrm{\Gamma }_2,T)`$ by the second author in \[S\].
After \[ES1\] was finished, we wanted to generalize its results to the quantum case. It was clear to us that to express the result we would need an explicit quantization of classical dynamical r-matrices from \[S\]. Therefore, we hoped that attempts to quantize the results of \[ES1\] using the approach of \[EV2\] could help us obtain such a quantization (which was unknown even for the usual Belavin-Drinfeld classical r-matrices). This program did, in fact, succeed, and the quantization of dynamical r-matrices from \[S\] was recently obtained in \[ESS\].
In this paper, using the results of \[ESS\] and methods of \[EV2\], we generalize the difference equations from \[EV2\] to the twisted case; this also provides a quantum generalization of \[ES1\]. Namely, we deduce difference equations with respect to the weight $`\lambda `$ for the generating function of the twisted traces for $`U_q(𝔤)`$, - the twisted MR and qKZB equations. Not surprisingly, they involve the dynamical R-matrix constructed in \[ESS\]. In the case when $`T`$ is an automorphism of the whole Dynkin diagram of $`𝔤`$, we also deduce the twisted dual MR and qKZB equations, i.e. the difference equations with respect to the other weight $`\mu `$. These equations involve the usual (Felder’s) dynamical R-matrix, but differ from the equations of \[EV2\] by explicit occurrence of $`T`$. Thus, we see that for $`T1`$, there is no symmetry between $`\lambda `$ and $`\mu `$.
If $`T`$ is not an automorphism, we do not expect the existence of the dual equations. This is explained at the end of Section 2.
Replacing $`U_q(𝔤)`$ with $`𝔤`$, we obtain the classical limit of these results. The twisted MR and qKZB equations become their classical analogs from \[ES1\]. The dual equations retain their form, but the trigonometric R-matrices are replaced by their rational limits.
Finally, we adapt the construction of the quantum dynamical R-matrices from \[ESS\] to the case when $`𝔤`$ is an arbitrary symmetrizable Kac-Moody algebra. This yields quantizations of the classical dynamical r-matrices from \[ES1\] in the case of Kac-Moody algebras. Again, the generating functions for twisted traces of intertwiners for $`U_q(𝔤)`$ satisfy a set of difference equations involving these quantum dynamical R-matrices, and a set of dual difference equations if in addition $`T`$ is an automorphism of the Dynkin diagram.
In the next paper, we plan to generalize these results to the case of affine algebras, when traces take values in finite-dimensional representations. This involves dynamical R-matrices with spectral parameters. In particular, we plan to obtain a trace representation of solutions of the elliptic qKZ equation (with Belavin’s elliptic R-matrix), and compute its monodromy.
#### Remark.
The elliptic qKZ equation is important in statistical mechanics (see \[JM\]). For its classical version, the trace representation of solutions and monodromy are obtained in \[E1\]. The problem of quantizing the results of \[E1\] was suggested to the first author by his advisor I. Frenkel as a topic for his PhD thesis in 1992. After this the first author tried to quantize the results of \[E1\] (see \[E2\]) but obtained only partial results.
### 1.1 Notations
####
Let $`𝔤`$ be a simple complex Lie algebra with a fixed polarization $`𝔤=𝔫_{}𝔥𝔫_+`$. Let $`\mathrm{\Gamma }`$ (resp. $`\mathrm{\Delta }`$) be the Dynkin diagram (resp. the root system) of $`𝔤`$. Denote by $`(a_{ij})`$ the Cartan matrix of $`𝔤`$ and let $`d_i`$ be relatively prime positive integers such that $`(d_ia_{ij})`$ is a symmetric matrix. Let $`(,)`$ be the nondegenerate invariant symmetric form for which $`(\alpha ,\alpha )=2`$ if $`\alpha `$ is a short root. Let $`\{e_\alpha ,f_\alpha \}_{\alpha \mathrm{\Delta }}`$ be a Chevalley basis of $`𝔫_{}𝔫_+`$, normalized in such a way that $`(e_\alpha ,f_\alpha )=1`$ for all $`\alpha `$ and set $`h_\alpha =[e_\alpha ,f_\alpha ]`$. We also let $`\mathrm{\Omega }𝔤𝔤`$ and $`\mathrm{\Omega }_𝔥𝔥𝔥`$ be the inverse elements of the restriction of $`()`$ to $`𝔤`$ and $`𝔥`$ respectively.
####
Let $`q=e^{\frac{1}{2}\mathrm{}}`$ where $`\mathrm{}`$ is a formal variable. For any operator $`A`$ we set $`q^A=e^{\mathrm{}\frac{A}{2}}`$. Let $`U_q(𝔤)`$ be the Drinfeld-Jimbo quantized enveloping algebra of $`𝔤`$. It is a $`[[\mathrm{}]]`$-Hopf algebra with generators $`E_\alpha ,F_\alpha `$, $`\alpha \mathrm{\Gamma }`$ and $`q^h`$, $`h𝔥`$ subject to the following set of relations :
$$q^{x+y}=q^xq^y,x,y𝔥q^hE_{\alpha _j}q^h=q^{\alpha _j(h)}E_{\alpha _j},q^hF_{\alpha _j}q^h=q^{\alpha _j(h)}F_{\alpha _j}$$
$$E_{\alpha _i}F_{\alpha _j}F_{\alpha _j}E_{\alpha _i}=\delta _{ij}\frac{q^{d_ih_{\alpha _i}}q^{d_ih_{\alpha _i}}}{q^{d_i}q^{d_i}},$$
$$\underset{k=0}{\overset{1a_{ij}}{}}(1)^k\left[\begin{array}{c}1a_{ij}\\ k\end{array}\right]_{q^{d_i}}E_{\alpha _i}^{1a_{ij}k}E_{\alpha _j}E_{\alpha _i}^k=0,ij,$$
$$\underset{k=0}{\overset{1a_{ij}}{}}(1)^k\left[\begin{array}{c}1a_{ij}\\ k\end{array}\right]_{q^{d_i}}F_{\alpha _i}^{1a_{ij}k}F_{\alpha _j}F_{\alpha _i}^k=0,ij.$$
where as usual
$$\left[\begin{array}{c}n\\ k\end{array}\right]_q=\frac{[n]_q!}{[k]_q![nk]_q!},[n]_q!=[1]_q[2]_q\mathrm{}[n]_q,[n]_q=\frac{q^nq^n}{qq^1}$$
####
Comultiplication $`\mathrm{\Delta },`$ antipode $`S`$ and counit $`ϵ`$ in $`U_q(𝔤)`$ are given by
$$\mathrm{\Delta }(E_{\alpha _i})=E_{\alpha _i}q^{d_ih_{\alpha _i}}+1E_{\alpha _i},\mathrm{\Delta }(F_{\alpha _i})=F_{\alpha _i}1+q^{d_ih_{\alpha _i}}F_{\alpha _i},\mathrm{\Delta }(q^h)=q^hq^h$$
$$S(E_{\alpha _i})=E_{\alpha _i}q^{d_ih_{\alpha _i}},S(F_{\alpha _i})=q^{d_ih_{\alpha _i}}F_i,S(q^h)=q^h$$
$$ϵ(E_{\alpha _i})=ϵ(F_{\alpha _i})=0,ϵ(q^h)=1.$$
####
Let $`U_q(𝔫_\pm )`$ be the subalgebra generated by $`(E_\alpha )_{\alpha \mathrm{\Gamma }}`$ and $`(F_\alpha )_{\alpha \mathrm{\Gamma }}`$ respectively. It is known that $`U_q(𝔤)`$ is quasitriangular, with R-matrix $`q^{\mathrm{\Omega }_𝔥}U_q(𝔫_+)\widehat{}U_q(𝔫_{})`$. Here $`\widehat{}`$ denotes the completion with respect to the principal grading of $`U_q(𝔫_\pm )`$.
### 1.2 Generalized Belavin-Drinfeld triples and classical dynamical r-matrices
####
Let $`𝔩𝔥`$ be a subalgebra on which $`(,)`$ is nondegenerate. Let $`(x_i)_{iI}`$ be an orthonormal basis of $`𝔩`$ and let $`(x^i)_{iI}`$ be the dual basis of $`𝔩^{}`$. The classical dynamical Yang-Baxter equation with respect to $`𝔩`$ is the following equation :
$$\begin{array}{cc}\hfill \underset{i}{}& \left(x_i^{(1)}\frac{r^{23}(\lambda )}{x^i}x_i^{(2)}\frac{r^{13}(\lambda )}{x^i}+x_i^{(3)}\frac{r^{12}(\lambda )}{x^i}\right)\hfill \\ & +[r^{12}(\lambda ),r^{13}(\lambda )]+[r^{12}(\lambda ),r^{23}(\lambda )]+[r^{13}(\lambda ),r^{23}(\lambda )]=0\hfill \end{array}$$
(1.1)
where $`r(\lambda ):𝔩^{}(𝔤𝔤)^𝔩`$ is a meromorphic function. Solutions of (1.1) relevant to the theory of Poisson-Lie groupoids (see \[EV1\], \[ES2\], \[Xu\]) are those satisfying the generalized unitarity condition, i.e $`r(\lambda )+r^{21}(\lambda )=\mathrm{\Xi }`$ is constant and $`\mathrm{\Xi }`$ belongs to $`(S^2𝔤)^𝔤`$. In \[S\] the second author classified all such solutions $`r(\lambda )`$ which are non skewsymmetric (that is, $`\mathrm{\Xi }0`$). Up to isomorphism and gauge transformations, they are labeled by the following combinatorial data called generalized Belavin-Drinfeld triples.
#### Definition.
A generalized Belavin-Drinfeld triple is a triple $`(\mathrm{\Gamma }_1,\mathrm{\Gamma }_2,T)`$ where $`\mathrm{\Gamma }_1,\mathrm{\Gamma }_2\mathrm{\Gamma }`$ and $`T:\mathrm{\Gamma }_1\stackrel{}{}\mathrm{\Gamma }_2`$ is an orthogonal isomorphism.
####
Let $`(\mathrm{\Gamma }_1,\mathrm{\Gamma }_2,T)`$ be a generalized Belavin-Drinfeld triple. Set
$$𝔩=\left(\underset{\alpha }{}(\alpha T(\alpha ))\right)^{}𝔥.$$
Note that $`𝔩`$ is spanned by real elements so that the restriction of $`(,)`$ to $`𝔩`$ is nondegenerate. Let $`𝔥_0𝔥`$ be the orthogonal complement to $`𝔩`$ in $`𝔥`$ and let $`\mathrm{\Omega }_{𝔥_0}𝔥_0𝔥_0`$ be the element inverse to the form $`(,)`$. The following lemmas are proved in \[ES1\].
###### Lemma 1.1.
There exists a unique Lie algebra homomorphism $`B:𝔟_{}𝔟_{}`$ (resp. $`B^1:𝔟_+𝔟_+`$) such that $`B(f_\alpha )=f_{T(\alpha )}`$, $`B(h_\alpha )=h_{T(\alpha )}`$ if $`\alpha \mathrm{\Gamma }_1`$, $`B(f_\alpha )=0`$ if $`\alpha \mathrm{\Gamma }_1`$ (resp. $`B^1(e_\alpha )=e_{T^1(\alpha )}`$, $`B^1(h_\alpha )=h_{T^1(\alpha )}`$ if $`\alpha \mathrm{\Gamma }_2`$, $`B^1(e_\alpha )=0`$ if $`\alpha \mathrm{\Gamma }_2`$), and $`B^{\pm 1}(h)=h`$ if $`h𝔩`$. Moreover the restriction of $`B`$ to $`𝔥`$ is an orthogonal operator.
#### Remark.
We use the symbol $`B^1`$ for notational convenience only. The operators $`B`$ and $`B^1`$ are only inverse to each other when restricted to $`𝔥`$.
###### Lemma 1.2 (Cayley transform).
For any $`x𝔥_0`$, there exists a unique element $`C_T(x)𝔥_0`$ such that for all $`\alpha \mathrm{\Gamma }_1`$ one has $`(\alpha T\alpha ,C_T(x))=(\alpha +T\alpha ,x)`$. The linear operator $`C_T:𝔥_0𝔥_0`$ is skew-symmetric.
The classical dynamical r-matrix associated to $`(\mathrm{\Gamma }_1,\mathrm{\Gamma }_2,T)`$ is
$$r_T(\lambda )=r_0^{21}+\underset{\alpha ,l1}{}e^{l(\alpha ,\lambda )}e_\alpha B^lf_\alpha +\frac{1}{2}(C_T1)\mathrm{\Omega }_{𝔥_0}$$
(1.2)
where $`r_0=\frac{1}{2}\mathrm{\Omega }_𝔥+_\alpha e_\alpha f_\alpha `$ is the standard classical r-matrix.
### 1.3 Quantum dynamical R-matrices
####
In our joint work with Travis Schedler \[ESS\] we obtain an explicit quantization of the r-matrices $`r_T(\lambda )`$. Namely, we construct a trigonometric function $`\stackrel{~}{R}_T(\lambda ):𝔩^{}U_q(𝔤)U_q(𝔤)`$ (tensor product in the category of topologically free $`[[\mathrm{}]]`$-modules) such that $`\stackrel{~}{R}_T(\lambda )1\mathrm{}r_T(\lambda )\mathrm{mod}(\mathrm{}^2)`$ which satisfies the quantum dynamical Yang-Baxter equation
$$\begin{array}{cc}\hfill \stackrel{~}{R}_T^{12}(\lambda \frac{1}{2}\mathrm{}h^{(3)})\stackrel{~}{R}_T^{13}(\lambda & +\frac{1}{2}\mathrm{}h^{(2)})\stackrel{~}{R}_T^{23}(\lambda \frac{1}{2}\mathrm{}h^{(1)})\hfill \\ & =\stackrel{~}{R}_T^{23}(\lambda +\frac{1}{2}\mathrm{}h^{(1)})\stackrel{~}{R}_T^{13}(\lambda \frac{1}{2}\mathrm{}h^{(2)})\stackrel{~}{R}_T^{12}(\lambda +\frac{1}{2}\mathrm{}h^{(3)}).\hfill \end{array}$$
(1.3)
In the above equation we used the usual notation for shifts in the dynamical variable: for instance, if $`S(\lambda )`$ is any meromorphic function $`𝔩^{}U_q(𝔤)^2`$ we set $`S(\lambda \frac{1}{2}\mathrm{}h^{(3)})=S(\lambda )\frac{1}{2}\mathrm{}_i\frac{S}{y^i}y_i^{(3)}+\mathrm{}`$ (the Taylor expansion), where $`y_1,\mathrm{},y_r`$ is a basis of $`𝔩`$ and $`y^1,\mathrm{}y^r`$ is the dual basis of $`𝔩^{}`$.
####
The construction is based on the following result. Let $`I_\pm U_q(𝔟_\pm )`$ be the kernels of the projections $`U_q(𝔟_\pm )U_q(𝔥)`$. Also set $`Z=\frac{1}{2}((1C_T)1)\mathrm{\Omega }_{𝔥_0}`$. The maps $`B:U_q(𝔟_{})U_q(𝔟_{})`$ and $`B^1:U_q(𝔟_+)U_q(𝔟_+)`$ are defined in the same fashion as in the classical case (see Lemma 1.1).
####
To simplify notations we will write $`q_i^A`$ for $`(q^A)_i`$ for any operator $`A`$ (the operator $`q^A`$ acting on the $`i`$-th component of a tensor product).
###### Theorem 1.1 (\[ESS\]).
There exists a unique trigonometric rational function $`𝒥_T:𝔩^{}(U_q(𝔟_{})U_q(𝔟_+))^𝔩`$ such that
1. $`𝒥_Tq^ZI^{}I^+`$,
2. $`𝒥_T(\lambda )`$ satisfies the modified ABRR equation :
$$^{21}q_1^{2\lambda }B_1(𝒥_T(\lambda ))=𝒥_T(\lambda )q_1^{2\lambda }q^{\mathrm{\Omega }_𝔩}.$$
(1.4)
Moreover $`𝒥_T(\lambda )`$ satisfies the shifted 2-cocycle condition :
$$𝒥_T(\lambda )^{12,3}(\lambda )𝒥_T^{12}(\lambda +\frac{1}{2}h^{(3)})=𝒥_T^{1,23}(\lambda )𝒥_T^{23}(\lambda \frac{1}{2}h^{(1)}).$$
(1.5)
The quantum dynamical R-matrix $`\stackrel{~}{R}_T(\lambda )`$ is obtained by twisting $``$ by $`𝒥_T(\lambda )`$ :
$$_T(\lambda )=𝒥_T^1(\lambda )^{21}𝒥_T^{21}(\lambda ),$$
$$\stackrel{~}{R}_T(\lambda )=_T^{21}(\frac{\lambda }{\mathrm{}}).$$
Note that the polarization we use here is the opposite to the polarization used in \[ESS\], where the twist $`𝒥_T(\lambda )`$ was an element of $`U_q(𝔟^+)U_q(𝔟^{})`$.
####
One aim of this paper is to provide a representation-theoretic interpretation of the quantum dynamical R-matrix $`_T(\lambda )`$. This is done in terms of twisted traces of quantum intertwiners and of the systems of difference equations satisfied by them.
## 2 Twisted traces of quantum intertwiners
### 2.1 Definition
####
Let $`M_\mu `$ be the Verma module over $`U_q(𝔤)`$ with highest weight $`\mu 𝔥^{}`$ and let $`v_\mu `$ be a highest weight vector. We will also consider the graded dual Verma module $`M_\mu ^{}`$ and let $`v_\mu ^{}`$ be its lowest weight vector satisfying $`v_\mu ^{},v_\mu =1`$. Let $`V`$ be a finite-dimensional $`U_q(𝔤)`$-module and let $`V=_\nu V[\nu ]`$ be its weight space decomposition. The following result is well-known (see e.g \[ES2\]) :
###### Lemma 2.1.
Suppose that $`M_\mu ^{}`$ is irreducible. Then the map
$`Hom_{U_q(𝔤)}(M_\mu ,M_\lambda V)`$ $`V[\mu \lambda ]`$
$`\mathrm{\Phi }`$ $`v_\lambda ^{},\mathrm{\Phi }v_\mu `$
is an isomorphism.
Conversely, for every weight $`\nu `$ and every homogeneous vector $`vV[\nu ]`$ we will denote by $`\mathrm{\Phi }_\mu ^v:M_\mu M_{\mu \nu }V`$ the unique intertwiner satisfying $`v_{\mu \nu }^{},\mathrm{\Phi }_\mu ^vv_\mu =v`$. It will be convenient to consider all the operators $`\mathrm{\Phi }_\mu ^v`$ simultaneously by setting
$$\mathrm{\Phi }_\mu ^V=\underset{v}{}\mathrm{\Phi }_\mu ^vv^{}Hom_{}(M_\mu ,\underset{\nu }{}M_{\mu \nu }VV^{}),$$
where $``$ is a homogeneous basis of $`V`$.
####
Let $`(\mathrm{\Gamma }_1,\mathrm{\Gamma }_2,T)`$ be a generalized Belavin-Drinfeld triple. Let $`𝔩,𝔥_0,C_T,\mathrm{}`$ have the same meanings as in Section 1. Finally, let $`\mu ,\mu ^{}`$ be weights satisfying the following relation :
$$(\mu ,\alpha )=(\mu ^{},T(\alpha ))\mathrm{for}\mathrm{all}\alpha \mathrm{\Gamma }_1.$$
(2.1)
We define a linear map $`B:M_\mu M_\mu ^{}`$ by $`uv_\mu B(u)v_\mu ^{}`$ for all $`uU_q(𝔫_{})`$.
####
Now consider finite-dimensional $`U_q(𝔤)`$-modules $`V_1\mathrm{},V_N`$ and let $`v_1V_1[\mu _1],`$ $`\mathrm{},v_NV_N[\mu _N]`$ be homogeneous vectors such that $`\overline{\mu }:=_i\mu _i𝔩^{}`$. The set of pairs of weights $`(\mu ,\mu ^{})`$ satisfying (2.1) and such that $`\mu ^{}\mu =\overline{\mu }`$ is an $`𝔩^{}`$-torsor $`\stackrel{~}{𝔩}^{}`$. For any such pair $`(\mu ,\mu ^{})`$ and for $`\lambda 𝔩^{}`$, we define the following formal power series in $`(V_1\mathrm{}V_N)^𝔩q^{2(\lambda ,\mu )}[[q^{2(\lambda ,\alpha _1)},\mathrm{},q^{2(\lambda ,\alpha _r)}]]`$ by analogy with \[EV2\]:
$$\mathrm{\Psi }_{v_1,\mathrm{},v_N}^T(\lambda ,\mu )=\mathrm{Tr}_{|M_\mu }(\mathrm{\Phi }_{\mu ^{}_{i=2}^N\mu _i}^{v_1}\mathrm{}\mathrm{\Phi }_\mu ^{}^{v_N}Be^\lambda )$$
and
$$\mathrm{\Psi }_{V_1,\mathrm{},V_N}^T(\lambda ,\mu )=\underset{v_i_i}{}\mathrm{\Psi }_{v_1,\mathrm{},v_N}^T(\lambda ,\mu )v_N^{}\mathrm{}v_1^{},$$
where $`_i`$ is a homogeneous basis of $`V_i`$. It is clear that $`\mathrm{\Psi }_{V_1,\mathrm{},V_N}^T(\lambda ,\mu )`$ takes values in $`(V_1\mathrm{}V_N)^𝔩(V_N^{}\mathrm{}V_1^{})^𝔩`$.
### 2.2 The main results
####
Our main result is that the functions $`\mathrm{\Psi }_{V_1,\mathrm{},V_N}^T(\lambda ,\mu )`$ satisfy some interesting difference equations. These difference equations are more conveniently expressed after some renormalizations. Set
$$J_T(\lambda )=𝒥_T(\lambda \rho +\frac{1}{2}(h^{(1)}+h^{(2)})),𝕁_T(\lambda )=𝒥_T(\lambda +\frac{1}{2}(h^{(1)}+h^{(2)})).$$
Put $`_T(\lambda )=m_{21}(1S^1)(𝕁_T(\lambda ))=m_{21}(1S^1)(J_T(\lambda ))`$. We will denote simply by $`(\lambda )`$ the element corresponding to the trivial triple $`(\mathrm{\Gamma },\mathrm{\Gamma },Id)`$. Also set $`R_T(\lambda )=J_T^1(\lambda )^{21}J_T^{21}(\lambda )`$ and $`_T(\lambda )=𝕁_T^1(\lambda )^{21}𝕁_T^{21}(\lambda )`$. Define
$$𝕁_T^{1\mathrm{}N}(\lambda )=𝕁_T^{1,2\mathrm{}N}(\lambda )\mathrm{}𝕁_T^{N1,N}(\lambda ).$$
Finally, let
$$\delta _q^T(\lambda )=\left(\mathrm{Tr}_{|M_\rho }(Bq^{2\lambda })\right)^1$$
be the twisted Weyl denominator. The explicit expression for $`\delta _q^T(\lambda )`$ is as follows. Let $`\mathrm{\Gamma }_3`$ be the subset of $`\mathrm{\Gamma }_1\mathrm{\Gamma }_2`$ consisting of roots which return to their original position after applying $`T`$ several times, and let $`\mathrm{\Gamma }_3`$ be the set of positive roots which are linear combinations of roots from $`\mathrm{\Gamma }_3`$. For each $`\alpha \mathrm{\Gamma }_3`$ let $`N_\alpha `$ be the order of the action of $`B`$ on $`\alpha `$. Consider the Lie algebra $`𝔤`$ and define $`\theta _\alpha `$ by $`B^{N_\alpha }u_\alpha =\theta _\alpha u_\alpha `$ for any $`u_\alpha 𝔤[\alpha ]`$. Then (see \[ES1\])
$$\delta _q^T(\lambda )=q^{2(\rho ,\lambda )}\underset{\overline{\alpha }\mathrm{\Gamma }_3/B}{}(1\theta _\alpha q^{2N_\alpha (\alpha ,\lambda )}).$$
Define the renormalized trace function by
$$F_{V_1,\mathrm{},V_N}^T(\lambda ,\mu )=[^1(\mu +h^{(1\mathrm{}N)})^{(N)}\mathrm{}^1(\mu +h^{(1)})^{(1)}]\phi _{V_1,\mathrm{},V_N}^T(\lambda ,\mu \rho )$$
where
$$\phi _{V_1,\mathrm{},V_N}^T(\lambda ,\mu )=𝕁_T^{1\mathrm{},N}(\lambda )^1\mathrm{\Psi }_{V_1,\mathrm{},V_N}^T(\lambda ,\mu )\delta _q^T(\lambda ).$$
####
Let $`W`$ be a finite-dimensional $`U_q(𝔤)`$-module. Consider the following difference operator acting on functions $`𝔩^{}(V_1\mathrm{}V_N)^𝔩`$ :
$$𝒟_W^T=\underset{\nu }{}\mathrm{Tr}_{|W[\nu ]}\left((_T)^{WV_1}(\lambda +h^{(2\mathrm{}N)})\mathrm{}(_T)^{WV_N}(\lambda )\right)𝕋_\nu $$
(2.2)
where $`𝕋_\nu f(\lambda )=f(\lambda +\nu )`$. In the above, we only consider the trace of the “diagonal block” of $`(_T)^{WV_1}(\lambda +h^{(2\mathrm{}N)})\mathrm{}(_T)^{WV_N}(\lambda )`$, i.e the part that preserves $`W[\nu ]`$.
###### Theorem 2.1 (Twisted Macdonald-Ruijsenaars equations).
$$𝒟_W^TF_{V_1,\mathrm{},V_N}^T(\lambda ,\mu )=\chi _W(q^{2\mu })F_{V_1,\mathrm{},V_N}^T(\lambda ,\mu ),$$
(2.3)
where $`\chi _W(x)=\mathrm{dim}W[\nu ]x^\nu `$ is the character of $`W`$ and where $`𝒟_W^T`$ acts on the variable $`\lambda `$.
This theorem is proved in Section 3.
####
For each $`j\{1,\mathrm{},N\}`$ consider the two following operators :
$`D_j^T`$ $`=q_j^{2\mu C_𝔥}q_{j,1}^{2\mathrm{\Omega }_𝔥}\mathrm{}q_{j,j1}^{2\mathrm{\Omega }_𝔥},`$ (2.4)
$`K_j^T`$ $`=_T^{j+1,j}(\lambda +h^{(j+2,\mathrm{},N)})^1\mathrm{}_T^{Nj}(\lambda )^1\mathrm{\Gamma }_j_T^{j1}(\lambda +h^{(2\mathrm{},j1)}+h^{(j+1\mathrm{},N)})\times `$
$`\times \mathrm{}_T^{j,j1}(\lambda +h^{(j+1\mathrm{},N)})`$ (2.5)
where $`C_𝔥=m_{12}(\mathrm{\Omega }_𝔥)U(𝔥)`$ is the quadratic Casimir element for $`𝔥`$, and where $`\mathrm{\Gamma }_jf(\lambda )=f(\lambda +h^{(j)})`$.
###### Theorem 2.2 (Twisted qKZB equations).
The function $`F_{V_1,\mathrm{},V_N}^T(\lambda ,\mu )`$ satisfies the following difference equation for all $`j=1,\mathrm{},N`$ :
$$F_{V_1,\mathrm{},V_N}^T(\lambda ,\mu )=(D_j^TK_j^T)F_{V_1,\mathrm{},V_N}^T(\lambda ,\mu ).$$
(2.6)
This theorem is proved in Section 4.
####
Now suppose that $`(\mathrm{\Gamma }_1,\mathrm{\Gamma }_2,T)`$ is a complete triple, i.e $`\mathrm{\Gamma }_1=\mathrm{\Gamma }_2=\mathrm{\Gamma }`$ and $`T`$ is an automorphism. In this case, the functions $`F_V^T(\lambda ,\mu )`$ satisfy in addition some dual difference equations, with respect to the variable $`\mu `$.
####
In a complete triple the maps $`B:U_q(𝔟_{})U_q(𝔟_{})`$ and $`B^1:U_q(𝔟_+)U_q(𝔟_+)`$ come from an automorphism $`B:U_q(𝔤)U_q(𝔤)`$. Let $`d`$ be the order of $`B`$ and let $`B\mathrm{Aut}(U_q(𝔤))`$ be the subgroup generated by $`B`$.
####
Let $`W`$ be any finite-dimensional $`U_q(𝔤)`$-module. We denote by $`W^B`$ the twist of $`W`$ by $`B`$: as a vector space $`W=W^B`$ and the $`U_q(𝔤)`$-action is given by $`uw=B^1(g)w`$. Now suppose that $`WW^B`$ as $`U_q(𝔤)`$-modules and let us fix an intertwiner in $`\mathrm{Hom}_{U_q(𝔤)}(W,W^B)\mathrm{Aut}_{}(W)`$ of order $`d`$. This endows $`W`$ with the structure of a module over $`BU_q(𝔤)`$.
####
Consider the following difference operator acting on functions with values in $`(V_N^{}\mathrm{}V_1^{})^𝔩`$ :
$$𝒟_W^{,T}=\underset{\nu }{}\mathrm{Tr}_{|W[\nu ]}\left(^{WV_N^{}}(\mu +h^{(1\mathrm{}N1)})\mathrm{}^{WV_1^{}}(\mu )B_W\right)𝕋_\nu ^{}$$
(2.7)
where $`𝕋_\nu ^{}f(\mu )=f(\mu +\nu )`$.
###### Theorem 2.3 (Dual twisted Macdonald-Ruijsenaars equations).
$$𝒟_W^{,T}F_{V_1,\mathrm{},V_N}^T(\lambda ,\mu )=\mathrm{Tr}_{|W^{𝔥_0}}(q^{2\lambda }B)F_{V_1,\mathrm{},V_N}^T(\lambda ,\mu ).$$
(2.8)
#### Remark.
Any $`B`$-invariant finite-dimensional $`U_q(𝔤)`$ module is a direct sum of modules $`\overline{V}_{\nu _0}:=_{\nu B\nu _0}V_\nu `$ where $`V_\nu `$ is the irreducible highest weight module of highest weight $`\nu `$ and where $`\nu _0`$ is dominant integral. It is easy to see that both sides of (2.8) identically vanish when $`W=\overline{V}_{\nu _0}`$ and $`\nu _0𝔩^{}`$ (i.e when $`B(\nu _0)\nu _0`$).
####
The twisted character $`\mathrm{Tr}_{|W^{𝔥_0}}(q^{2\lambda }B)`$ can be expressed explicitly when $`W=V_\nu `$ with $`\nu 𝔩^{}`$. Consider the quotient Dynkin diagram $`\overline{\mathrm{\Gamma }}=\{\overline{\alpha }=_{\alpha \overline{\alpha }}\alpha |\overline{\alpha }\mathrm{\Gamma }/T\}`$ which forms a set of simple roots in $`𝔩^{}`$ with respect to the restriction of $`(,)`$ to $`𝔩^{}`$. Namely, if $`\mathrm{\Gamma }=\{\alpha _1,\mathrm{},\alpha _n\}`$ is of type $`A_n`$ and if $`T`$ is the “flip” $`T:\alpha _i\alpha _{n+1i}`$ then $`\mathrm{\Gamma }/T`$ is $`B_k`$ where $`n=2k1`$ or $`n=2k`$; if $`\mathrm{\Gamma }=D_4`$ and $`T`$ is the rotation of order three around the trivalent root then $`\mathrm{\Gamma }/T=G_2`$; if $`\mathrm{\Gamma }=D_k`$ and $`T`$ is the symmetry of order two around the trivalent root then $`\overline{\mathrm{\Gamma }}=C_{k1}`$; and if $`\mathrm{\Gamma }=E_6`$ and $`T`$ is the symmetry around the trivalent root then $`\mathrm{\Gamma }/T=F_4`$. Let $`\overline{\mathrm{\Delta }}`$ be the root system of $`\overline{\mathrm{\Gamma }}`$. Let $`\overline{𝔤}`$ be the simple complex Lie algebra associated to $`\overline{\mathrm{\Gamma }}`$. Note that any weight $`\nu 𝔩^{}`$ is naturally a weight for $`\overline{\mathrm{\Gamma }}`$. However, the scalar product on $`𝔩^{}`$ is not the usual one corresponding to the root system $`\overline{\mathrm{\Delta }}`$. For instance we have $`(\overline{\alpha },\overline{\alpha })=2N_\alpha `$ if $`N_\alpha `$ is the number of elements in the $`T`$-orbit $`\overline{\alpha }`$ and if any two elements of that orbit are orthogonal. For $`\lambda =_\alpha c_{\overline{\alpha }}\overline{\alpha }𝔩^{}`$ set $`\overline{\lambda }=_{\overline{\alpha }}\frac{2N_\alpha }{(\overline{\alpha },\overline{\alpha })}c_{\overline{\alpha }}\overline{\alpha }`$.
###### Proposition 2.1.
For any $`\nu 𝔩^{}`$ we have
$$\mathrm{Tr}_{|W^{𝔥_0}}(q^{2\lambda }B)=\chi _{\overline{V}_\nu }(q^{2\overline{\lambda }})$$
where $`\overline{V}_\nu `$ is the irreducible $`\overline{𝔤}`$-module of highest weight $`\nu `$.
Theorem 2.3 and Proposition 2.1 are proved in Section 5.
####
Now, define for each $`j\{1,\mathrm{},N\}`$ the following operators
$$\begin{array}{cc}\hfill D_j^{,T}& =q_j^{2\lambda C_𝔩}q_{j,j+1}^{\mathrm{\Omega }_𝔩}\mathrm{}q_{j,N}^{\mathrm{\Omega }_𝔩},\hfill \\ \hfill K_j^{,T}& =_{j1,j}(\mu +h^{(1\mathrm{}j2)})^1\mathrm{}_{1,j}(\mu )^1\mathrm{\Gamma }_{B^1(j)}^{}\times \hfill \\ \hfill _{j,N}(\mu & +h^{(j+1\mathrm{}N1)}+h^{(1\mathrm{}j1)})\mathrm{}_{j,j+1}(\mu +h^{(1\mathrm{}j1)}),\hfill \end{array}$$
(2.9)
where $`C_𝔩=m_{12}(\mathrm{\Omega }_𝔩)U(𝔩)`$ and where $`\mathrm{\Gamma }_{B^1(j)}^{}f(\mu )=f(\mu +B^1(h^{(j)}))`$.
###### Theorem 2.4 (Dual twisted qKZB equations).
The functions $`F_{V_1,\mathrm{},V_N}^T(\lambda ,\mu )`$ satisfy the following difference equation for each $`j=1\mathrm{},N`$ :
$$B_{V_j}B_{V_j^{}}^{}F_{V_1,\mathrm{},V_j,\mathrm{},V_N}^T(\lambda ,\mu )=(D_j^{,T}K_j^{,T})F_{V_1,\mathrm{},V_j^B,\mathrm{},V_N}^T(\lambda ,\mu ).$$
(2.10)
This theorem is proved in Section 6.
#### Remark 1.
For $`T=Id`$, Theorems 2.1-2.4 appear in \[EV2\].
#### Remark 2.
We do not expect the dual equations to exist for non-complete triples. This can be explained in the following way. Suppose that $`𝔤`$ is an affine Lie algebra and that $`T=Id`$, so that $`r_T(\lambda ,z)`$ is the Felder elliptic dynamical r-matrix. In that case it is known that the the dual trigonometric qKZB equations without spectral parameter can be interpreted as monodromy of the flat connection on the torus defined by the classical (elliptic) KZB equations (see \[Ki\]). One can show that this is true for any elliptic dynamical r-matrix. On the other hand, it was proved in \[ES1\] Proposition 4.2 that the classical dynamical r-matrix with spectral parameter $`r_T(\lambda ,z)`$ associated to an affine Lie algebra and a triple $`(\mathrm{\Gamma }_1,\mathrm{\Gamma }_2,T)`$ is elliptic only when $`T`$ is an automorphism; for general triples, it is partially elliptic and partially trigonometric (for instance, it is purely trigonometric when $`T`$ is nilpotent). This shows that the monodromy of these KZB equations should be defined only for complete triples, and hence the existence of the dual equations should be expected only for them.
#### Remark 3.
The above theorems are also valid for the specialized quantum group $`U_q(𝔤)`$, which is obtained from the formal quantum group when we take $`\mathrm{}^{}\{i\}`$ to be a complex number. In that case, it is more convenient to consider the twist $`𝒥_T(\lambda )`$ as an endomorphism of the functor
$$F:Rep(U_q(𝔤))\times Rep(U_q(𝔤))Vec$$
which assigns to any two finite-dimensional $`U_q(𝔤)`$-modules $`V`$ and $`W`$ the vector space $`VW`$. Here $`Rep(U_q(𝔤))`$ is the category of finite-dimensional $`U_q(𝔤)`$-modules and $`Vec`$ is the category of finite-dimensional $``$-vector spaces. For instance, equation (1.5) means that for every three representations $`U,V,W`$ and vectors $`uU`$, $`vV`$ and $`wW`$ with respective weights $`\lambda _u,\lambda _v`$ and $`\lambda _w`$ we have
$$𝒥_T(\lambda )^{12,3}(\lambda )𝒥_T^{12}(\lambda +\frac{1}{2}\lambda _w)(uvw)=𝒥_T^{1,23}(\lambda )𝒥_T^{23}(\lambda \frac{1}{2}\lambda _u)(uvw).$$
## 3 The twisted Macdonald-Ruijsenaars equations
####
The proof of Theorem 2.1 is an extension of the proof of Theorem 1.1 of \[EV2\] to the case of an arbitrary generalized Belavin-Drinfeld triple. From now on we fix such a triple $`(\mathrm{\Gamma }_1,\mathrm{\Gamma }_2,T)`$. We first note that the notion of radial part generalizes straightforwardly to the twisted setting :
###### Proposition 3.1.
Let $`V`$ be a finite-dimensional $`U_q(𝔤)`$-module. For any $`XU_q(𝔤)`$ there exists a unique difference operator $`𝒟_X^T`$ (with respect to the variable $`\lambda `$) acting on formal power series in $`V^𝔩q^{2(\lambda ,\mu )}[[q^{2(\lambda ,\alpha _1)},\mathrm{},,q^{2(\lambda ,\alpha _r)}]]`$, $`\lambda 𝔩^{}`$ such that we have
$$\mathrm{Tr}_{|M_\mu }(\mathrm{\Phi }_\mu ^{}^VXq^{2\lambda }B)=𝒟_X^T\mathrm{Tr}_{|M_\mu }(\mathrm{\Phi }_\mu ^{}^Vq^{2\lambda }B).$$
The operator $`𝒟_X^T`$ is called the twisted radial part of $`X`$.
####
For any finite-dimensional $`U_q(𝔤)`$-module $`W`$ set
$$C_W=\mathrm{Tr}_{|W}(1\pi _W)(^{21}(1q^{2\rho })).$$
It is well-known (see \[D\], \[R\]) that the map $`WC_W`$ defines a homomorphism from the Grothendieck ring of the category of finite-dimensional $`U_q(𝔤)`$-modules to the center of $`U_q(𝔤)`$. Set $`_W^T=𝒟_{C_W}^T`$.
###### Proposition 3.2.
We have
1. $`_W^T_V^T=_V^T_W^T`$ for all $`V,W`$,
2. $`_W^T\mathrm{\Psi }_V^T(\lambda ,\mu )=\chi _W(q^{2(\mu +\rho )})\mathrm{\Psi }_V^T(\lambda ,\mu )`$ where $`\chi _W(x)=_\nu \mathrm{dim}W[\nu ]x^\nu `$ is the character of $`W`$.
Proof. See \[EK\], \[EV2\].
####
Let us now proceed to explicitly compute the operator $`_W`$.
####
Let $`V`$ be a finite-dimensional $`U_q(𝔤)`$-module. Introduce the following function with values in $`VV^{}U_q(𝔤)`$, with components labeled as $`1,1`$ and $`2`$ respectively :
$$Z_V(\lambda ,\mu )=\mathrm{Tr}_{|M_\mu }(\mathrm{\Phi }_\mu ^{}^V^{20}B_0q_0^{2\lambda }).$$
###### Lemma 3.1.
We have
$$Z_V(\lambda ,\mu )=𝒥_T^{12}(\lambda )\mathrm{\Psi }_V^T(\lambda +\frac{1}{2}h^{(2)},\mu ).$$
(3.1)
Proof. First we note that, by pulling the R-matrix around the trace and using the intertwining property together with the fact that $`B_1^1=B_2`$ we obtain
$$\begin{array}{cc}\hfill Z_V(\lambda ,\mu )& =^{21}q_1^{2\lambda }\mathrm{Tr}_{|M_\mu }(\mathrm{\Phi }_\mu ^{}^V(B_2^1^{20})B_0q_0^{2\lambda })\hfill \\ & =^{21}q_1^{2\lambda }B_2^1\left(\mathrm{Tr}_{|M_\mu }(\mathrm{\Phi }_\mu ^{}^V^{20}B_0q_0^{2\lambda })\right)\hfill \\ & =^{21}q_1^{2\lambda }B_2^1Z_V(\lambda ,\mu )\hfill \end{array}$$
On the other hand, using the defining equation for $`𝒥_T(\lambda )`$, the relation $`B_2^1𝒥_T(\lambda )=B_1𝒥_T(\lambda )`$ and the $`𝔩`$-invariance of $`\mathrm{\Psi }_V^T(\lambda ,\mu )`$ we have
$$\begin{array}{cc}\hfill ^{21}q_1^{2\lambda }B_2^1\left(𝒥_T^{12}(\lambda )\mathrm{\Psi }_V^T(\lambda +\frac{1}{2}h^{(2)},\mu )\right)=& ^{21}q_1^{2\lambda }\left(B_1𝒥_T^{12}(\lambda )\right)\mathrm{\Psi }_V^T(\lambda +\frac{1}{2}h^{(2)},\mu )\hfill \\ \hfill =& 𝒥_T^{12}(\lambda )q_1^{2\lambda }q_{12}^{\mathrm{\Omega }_𝔩}\mathrm{\Psi }_V^T(\lambda +\frac{1}{2}h^{(2)},\mu )\hfill \\ \hfill =& 𝒥_T^{12}(\lambda )\mathrm{\Psi }_V^T(\lambda +\frac{1}{2}h^{(2)},\mu )\hfill \end{array}$$
The lemma now follows from the fact both sides of (3.1) satisfy the relation $`Y=^{21}q_1^{2\lambda }B_2^1Y`$ and are of the form $`Y=q^Z\mathrm{\Psi }_V^T(\lambda +\frac{1}{2}h^{(2)},\mu )+l.o.t`$, where $`l.o.t`$ stands for terms of strictly positive degree in component $`2`$.$`\mathrm{}`$
####
Consider the following function with values in $`VV^{}U_q(𝔤)U_q(𝔤)`$ (with components labeled as $`1,1,2`$ and $`3`$ respectively) :
$$X_V(\lambda ,\mu )=\mathrm{Tr}_{|M_\mu }(\mathrm{\Phi }_\mu ^{}^V^{20}q_0^{2\lambda }B_0(^{03})^1).$$
###### Lemma 3.2.
We have
$$X_V(\lambda ,\mu )=𝒥_T^{3,12}(\lambda )𝒥_T^{12}(\lambda \frac{1}{2}h^{(3)})\mathrm{\Psi }_V^T(\lambda +\frac{1}{2}(h^{(2)}h^{(3)}),\mu )𝒥_T^{32}(\lambda )^1$$
(3.2)
Proof. Moving $`^{03}`$ around the trace, using the quantum Yang-Baxter equation for $``$ and the $`B`$-invariance property of $``$ again, we get
$$\begin{array}{cc}\hfill X_V(\lambda ,\mu )& =^{13}\mathrm{Tr}_{|M_\mu }(\mathrm{\Phi }_\mu ^{}^V(^{03})^1^{20}q_0^{2\lambda }B_0)\hfill \\ & =^{13}^{23}\mathrm{Tr}_{|M_\mu }(\mathrm{\Phi }_\mu ^{}^V^{20}(^{03})^1q_0^{2\lambda }B_0)(^{23})^1\hfill \\ & =^{13}^{23}q_3^{2\lambda }\mathrm{Tr}_{|M_\mu }(\mathrm{\Phi }_\mu ^{}^V^{20}q_0^{2\lambda }(^{03})^1B_0)q_3^{2\lambda }(^{23})^1\hfill \\ & =^{13}^{23}q_3^{2\lambda }B_3\left(\mathrm{Tr}_{|M_\mu }(\mathrm{\Phi }_\mu ^{}^V^{20}q_0^{2\lambda }B_0(^{03})^1)\right)q_3^{2\lambda }(^{23})^1\hfill \\ & =^{13}^{23}q_3^{2\lambda }B_3\left(X_V(\lambda ,\mu )\right)q_3^{2\lambda }(^{23})^1.\hfill \end{array}$$
On the other hand, using the modified ABRR equation (1.4) we have
$$\begin{array}{cc}\hfill ^{13}& ^{23}q_3^{2\lambda }B_3\left(𝒥_T^{3,12}(\lambda )\right)𝒥_T^{12}(\lambda \frac{1}{2}h^{(3)})\mathrm{\Psi }_V^T(\lambda +\frac{1}{2}(h^{(2)}h^{(3)}),\mu )\times \hfill \\ & \times B_3(𝒥_T^{32}(\lambda )^1)q_3^{2\lambda }(^{23})^1\hfill \\ & =\mathrm{\Delta }_1\left(^{13}q_3^{2\lambda }B_3(𝒥_T^{31}(\lambda ))\right)𝒥_T^{12}(\lambda \frac{1}{2}h^{(3)})\mathrm{\Psi }_V^T(\lambda +\frac{1}{2}(h^{(2)}h^{(3)}),\mu )\times \hfill \\ & \times q_3^{2\lambda }q_{23}^{\mathrm{\Omega }_𝔩}𝒥_T^{32}(\lambda )^1\hfill \\ & =\mathrm{\Delta }_1\left(𝒥_T^{31}(\lambda )q_3^{2\lambda }q_{31}^{\mathrm{\Omega }_𝔩}\right)𝒥_T^{12}(\lambda \frac{1}{2}h^{(3)})\mathrm{\Psi }_V^T(\lambda +\frac{1}{2}(h^{(2)}h^{(3)}),\mu )\times \hfill \\ & \times q_3^{2\lambda }q_{23}^{\mathrm{\Omega }_𝔩}𝒥_T^{32}(\lambda )^1\hfill \\ & =𝒥_T^{3,12}q_3^{2\lambda }q_{31}^{\mathrm{\Omega }_𝔩}q_{32}^{\mathrm{\Omega }_𝔩}𝒥_T^{12}(\lambda \frac{1}{2}h^{(3)})\mathrm{\Psi }_V^T(\lambda +\frac{1}{2}(h^{(2)}h^{(3)}),\mu )q_3^{2\lambda }q_{23}^{\mathrm{\Omega }_𝔩}𝒥_T^{32}(\lambda )^1\hfill \\ & =𝒥_T^{3,12}𝒥_T^{12}(\lambda \frac{1}{2}h^{(3)})\mathrm{\Psi }_V^T(\lambda +\frac{1}{2}(h^{(2)}h^{(3)}),\mu )𝒥_T^{32}(\lambda )^1\hfill \end{array}$$
Now set $`X(\lambda )=(𝒥_T^{3,12})^1X_V(\lambda )𝒥_T^{32}(\lambda )`$. By the above and by Lemma 3.1, both $`X(\lambda )`$ and $`Z_V(\lambda \frac{1}{2}h^{(3)})`$ satisfy the equation
$$q_3^{2\lambda }q_{31}^{\mathrm{\Omega }_𝔩}q_{32}^{\mathrm{\Omega }_𝔩}Y=Yq_3^{2\lambda }q_{23}^{\mathrm{\Omega }_𝔩}$$
and are both of the form $`Y=𝒥_T^{12}(\lambda \frac{1}{2}h^{(3)})\mathrm{\Psi }_V^T(\lambda +\frac{1}{2}(h^{(2)}h^{(3)}))+l.o.t`$. Hence $`X(\lambda )=Z_V(\lambda \frac{1}{2}h^{(3)})`$ and the lemma is proved.$`\mathrm{}`$
###### Corollary 3.1.
We have
$$\begin{array}{cc}\hfill \mathrm{Tr}_{|M_\mu }& (\mathrm{\Phi }_\mu ^V^{20}(^{03})^1q^{2\lambda }B_0)\hfill \\ \hfill =& B_3\left(q_3^{2\lambda }𝒥_T^{3,12}(\lambda )𝒥_T^{12}(\lambda \frac{1}{2}h^{(3)})\mathrm{\Psi }_V(\lambda +\frac{1}{2}(h^{(2)}h^{(3)}),\mu )𝒥_T^{32}(\lambda )^1q_3^{2\lambda }\right)\hfill \end{array}$$
####
Now let $`W`$ be any finite-dimensional $`U_q(𝔤)`$-module. By Corollary 3.1, we get
$$\begin{array}{cc}\hfill _W^T\mathrm{\Psi }_V^T(\lambda ,\mu )& =\mathrm{Tr}_{|M_\mu }(\mathrm{\Phi }_\mu ^{}^VC_Wq_0^{2\lambda }B_0)\hfill \\ & =\mathrm{Tr}_{|M_\mu }(\mathrm{Tr}_{|W}(^{W0}^{0W}q_W^{2\rho })\mathrm{\Phi }_\mu ^{}^Vq_0^{2\lambda }B_0)\hfill \\ & =\mathrm{Tr}_{|W}\mathrm{Tr}_{|M_\mu }\left(m_{23}(^{20}S_3(^{03})^1q_3^{2\rho })\mathrm{\Phi }_\mu ^{}^Vq_0^{2\lambda }B_0\right)\hfill \\ & =\mathrm{Tr}_{|W}\left\{m_{23}\left(S_3\mathrm{Tr}_{|M_\mu }(^{20}(^{03})^1\mathrm{\Phi }_\mu ^Vq_0^{2\lambda }B_0)\right)q_2^{2\rho }\right\}\hfill \\ & =\mathrm{Tr}_{|W}\{\underset{ijk}{}[d_k^{(1)}(\lambda )c_iq^{2\lambda }d_k^{(2)}(\lambda )d_i(\lambda +\frac{1}{2}h^{(2)})]\times \hfill \\ & \times \mathrm{\Psi }_V^T(\lambda +h^{(W)})\left(d_j^{}(\lambda )q^{2\lambda }S(B(c_j^{}))S(B(c_k))q^{2\rho }\right)_W\}\hfill \end{array}$$
(3.3)
where $`m:U_q(𝔤)U_q(𝔤)U_q(𝔤)`$ is the multiplication map, $`𝒥_T(\lambda )=_ic_id_i(\lambda )`$, $`𝒥_T^1(\lambda )=_ic_i^{}d_i^{}(\lambda )`$ and where we used Sweedler’s notation for coproducts : $`\mathrm{\Delta }(x)=x^{(1)}x^{(2)}`$.
####
Let us set $`=_ia_ib_i`$, $`^1=_ia_i^{}b_i^{}`$.
###### Lemma 3.3.
We have
$$\underset{j}{}d_j^{}(\lambda )q^{2\lambda }S(B(c_j^{}))=q^{C_𝔩}P(\lambda )S(u)q^{2\lambda }$$
(3.4)
where $`u=_iS(b_i)a_i`$ is the Drinfeld element and $`P(\lambda )=_jd_j^{}(\lambda )S^1(c_j^{})`$.
Proof. This is obtained by applying $`m_{21}(S1)`$ to the relation $`(B_1𝒥_T(\lambda ))^1q_1^{2\lambda }=q^{\mathrm{\Omega }_𝔩}q_1^{2\lambda }𝒥_T^1(\lambda )^{21}`$, which itself follows from (1.4) and the $`B`$-invariance of $`𝒥_T(\lambda )`$.$`\mathrm{}`$
####
Substituting (3.4) in (3.3) yields
$$\begin{array}{cc}\hfill _W^T\mathrm{\Psi }_V^T(\lambda ,\mu )& =\underset{ijk,\nu }{}d_k^{(1)}(\lambda )c_i\mathrm{Tr}_{|W[\nu ]}\{q^{C_𝔩}P(\lambda )q^{2\lambda }S^1(B(c_k))q^{2\lambda }d_k^{(2)}(\lambda )\times \hfill \\ & \times d_i(\lambda +\frac{1}{2}\nu )S(u)q^{2\rho }\}\mathrm{\Psi }_V^T(\lambda +\nu ).\hfill \end{array}$$
(3.5)
###### Lemma 3.4.
We have
$$\begin{array}{cc}\hfill \underset{k}{}& d_k^{(1)}(\lambda )q^{2\lambda }S^1(B(c_k))q^{2\lambda }d_k^{(2)}(\lambda )\hfill \\ & =\underset{jk}{}(a_j^{})^{(1)}d_k^{(1)}(\lambda )q^{\mathrm{\Omega }_𝔩+1C_𝔩}\{S^1(c_k)S^1(b_j^{})(a_j^{})^{(2)}d_k^{(2)}(\lambda )\}_2\hfill \end{array}$$
(3.6)
Proof. From the modified ABRR relation we get
$$q_1^{2\lambda }B_1(𝒥_T(\lambda ))q_1^{2\lambda }=(^{21})^1𝒥_T(\lambda )q^{\mathrm{\Omega }_𝔩}.$$
Applying $`1\mathrm{\Delta }`$ yields
$$q_1^{2\lambda }B_1(𝒥_T^{1,23}(\lambda ))q_1^{2\lambda }=(^{23,1})^1𝒥_T^{1,23}(\lambda )q_{12}^{\mathrm{\Omega }_𝔩}q_{13}^{\mathrm{\Omega }_𝔩},$$
which can be written as
$$\begin{array}{cc}\hfill \underset{k}{}q_1^{2\lambda }B(c_k)q_1^{2\lambda }& d_k^{(1)}(\lambda )d_k^{(2)}(\lambda )\hfill \\ & =\underset{ik}{}(b_i^{}(a_i^{})^{(1)}(a_i^{})^{(2)})\times (c_kd_k^{(1)}(\lambda )d_k^{(2)}(\lambda ))q_{12}^{\mathrm{\Omega }_𝔩}q_{13}^{\mathrm{\Omega }_𝔩}.\hfill \end{array}$$
Equation (3.6) is now obtained by applying $`m_{13}(S^111)`$. $`\mathrm{}`$
####
We introduce the following notation. For any linear operator $`H(\lambda )\mathrm{End}(V_1\mathrm{}V_N)`$ we set
$$H(\lambda +\widehat{h}^{(i)})(v_1\mathrm{}v_N)=\underset{\nu }{}H_\nu (\lambda +\nu )(v_1\mathrm{}v_N)$$
where $`H_\nu (\lambda ):V_1\mathrm{}V_i\mathrm{}V_NV_1\mathrm{}V_i[\nu ]\mathrm{}V_N`$ is the block of $`H(\lambda )`$ with image $`V_i[\nu ]`$ in the $`i`$-th component. In other words, we replace $`\widehat{h}^{(i)}`$ by the weight in the $`i`$-th component after the action of $`H`$.
###### Lemma 3.5.
The following identities hold :
1. $`_j(a_j^{})^{(1)}S^1(b_j^{})(a_j^{})^{(2)}=a_kub_k,`$
2. $`^{23}𝒥_T^{1,23}(\lambda )=𝒥_T^{1,32}(\lambda )^{23},`$
3. $`S(c_i)d_i^{(1)}(\lambda )d_i^{(2)}(\lambda )=S(_T)(\lambda \frac{1}{2}h^{(2)})_1𝒥_T^1(\lambda +\frac{1}{2}\widehat{h}^{(1)}).`$
Proof. Equalities i) and iii) are proved in the same fashion as in \[EV2\]. Equality ii) follows from the relation $`\mathrm{\Delta }=\mathrm{\Delta }^{op}`$.$`\mathrm{}`$
###### Corollary 3.2.
We have
$$\begin{array}{cc}\hfill d_k^{(1)}& (\lambda )q^{2\lambda }S^1(B(c_k))q^{2\lambda }d_k^{(2)}(\lambda )\hfill \\ & =q^{\mathrm{\Omega }_𝔩1C_𝔩}u_2^1S(_T)_2(\lambda \frac{1}{2}h^{(1)})(𝒥_T^{21})^1(\lambda +\frac{1}{2}\widehat{h}^{(2)}).\hfill \end{array}$$
Proof. Use i),ii) and iii) successively, as in \[EV2\], (2.32).$`\mathrm{}`$
####
By Corollary 3.2 and using the relation $`B_1𝒥_T(\lambda )=B_2^1𝒥_T(\lambda )`$, we can rewrite (3.5) as follows :
$$\begin{array}{cc}& _W^T\mathrm{\Psi }_V^T(\lambda ,\mu )\hfill \\ & =\underset{\nu }{}\mathrm{Tr}_{|W[\nu ]}\{P_2(\lambda )q^{\mathrm{\Omega }_𝔩1C_𝔩}u_2^1S(_T)_2(\lambda \frac{1}{2}h^{(1)})(𝒥_T^{21})^1(\lambda +\frac{1}{2}\widehat{h}^{(2)})\hfill \\ & \times (c_i)_1d_i(\lambda +\frac{1}{2}\nu )_2S(u)_2q_2^{2\rho }q_2^{m_{12}\mathrm{\Omega }_𝔩}\}\mathrm{\Psi }_V^T(\lambda +\nu ,\mu )\hfill \\ & =\underset{\nu }{}\mathrm{Tr}_{|W[\nu ]}\left(\stackrel{~}{G}(\lambda )(_T)^{WV}(\lambda )\right)\mathrm{\Psi }_V(\lambda +\nu ,\mu )\hfill \end{array}$$
(3.7)
where $`\stackrel{~}{G}(\lambda )=q^{2\rho }P(\lambda )S(_T)(\lambda )`$. We now proceed to compute $`\stackrel{~}{G}(\lambda )`$.
###### Proposition 3.3.
We have $`\stackrel{~}{G}(\lambda )=\frac{\delta _q^T(\lambda +h)}{\delta _q^T(\lambda )}`$.
Proof. The following lemma is proved as in \[EV2\] :
###### Lemma 3.6.
We have $`P(\lambda )=_T^1(\lambda +h)`$, i.e $`\stackrel{~}{G}(\lambda )=G(\lambda +h)`$ where $`G(\lambda )=q^{2\rho }_T^1(\lambda )S(_T)(\lambda h)`$.
A direct (though lengthy) computation shows that
$$\mathrm{\Delta }(G(\lambda ))=𝕁_T(\lambda )\left(G(\lambda +h^{(2)})G(\lambda )\right)𝕁_T^1(\lambda )$$
(3.8)
(see \[M\] for a detailed proof of this in the nondynamical case; the dynamical case is analogous). In particular, replacing $`\lambda `$ by $`\frac{\lambda }{\mathrm{}}`$, we have
$$\begin{array}{cc}\hfill 𝕁_T(\frac{\lambda }{\mathrm{}})(G(\frac{\lambda +\mathrm{}h^{(2)}}{\mathrm{}})G(\frac{\lambda }{\mathrm{}}))& 𝕁_T^1(\frac{\lambda }{\mathrm{}})\hfill \\ & =𝕁_T^{21}(\frac{\lambda }{\mathrm{}})(G(\frac{\lambda }{\mathrm{}})G(\frac{\lambda +\mathrm{}h^{(1)}}{\mathrm{}}))(𝕁_T^{21})^1(\frac{\lambda }{\mathrm{}})\hfill \end{array}$$
which can be rewritten as
$$_T^{21}(\frac{\lambda }{\mathrm{}})(G(\frac{\lambda +\mathrm{}h^{(2)}}{\mathrm{}})G(\frac{\lambda }{\mathrm{}}))=(G(\frac{\lambda }{\mathrm{}})G(\frac{\lambda +\mathrm{}h^{(1)}}{\mathrm{}}))_T^{21}(\frac{\lambda }{\mathrm{}})$$
(3.9)
Let us now expand $`_T(\frac{\lambda }{\mathrm{}})`$ and $`G(\frac{\lambda }{\mathrm{}})`$ around $`\mathrm{}=0`$ :
$$_T(\frac{\lambda }{\mathrm{}})=1+\mathrm{}r(\lambda )+𝒪(\mathrm{}^2),G(\frac{\lambda }{\mathrm{}})=1+\mathrm{}g_1(\lambda )+𝒪(\mathrm{}^2)$$
where $`g_1(\lambda )=(G(\lambda /\mathrm{})1)/\mathrm{}U_q(𝔤)/\mathrm{}U_q(𝔤)U(𝔤)`$. Note that by (3.8) we have $`\mathrm{\Delta }_0(g_1(\lambda ))=g_1(\lambda )1+1g_1(\lambda )`$ (where $`\mathrm{\Delta }_0`$ is the usual coproduct on $`U𝔤`$), which implies that $`g_1(\lambda )𝔤`$. But since $`G(\lambda )`$ is of $`𝔩`$-weight zero, $`g_1(\lambda )𝔥`$. Now, by (3.9), we have
$$\underset{i}{}x_i\frac{g_1(\lambda )}{x_i}=[r(\lambda ),g_1(\lambda )1+1g_1(\lambda )],$$
where $`(x_i)`$ is a basis of $`𝔩`$. In particular, $`[r(\lambda ),g_1(\lambda )1+1g_1(\lambda )]\mathrm{\Lambda }^2𝔥`$. But this implies that $`[r(\lambda ),g_1(\lambda )1+1g_1(\lambda )]=0`$. Thus $`g_1:𝔩^{}𝔩`$ is a closed 1-form on $`𝔩^{}`$ and there exists functions $`f_1(\lambda )`$ and $`g_2(\lambda )=\frac{1}{\mathrm{}^2}(G(\frac{\lambda }{\mathrm{}})\frac{f_1(\frac{\lambda }{\mathrm{}})}{f_1(\frac{\lambda \mathrm{}h}{\mathrm{}})})U_q(𝔤)/\mathrm{}U_q(𝔤)U(𝔤)`$. By the same argument, $`𝔤_2(\lambda )`$ is a closed 1-form. Continuing in this process, we finally obtain a function $`f`$ defined on $`𝔩^{}`$ such that $`G(\lambda )=\frac{f(\lambda )}{f(\lambda h)}`$. It remains to determine $`f(\lambda )`$ explicitly. For this, apply (3.7) and Proposition 3.2 2. to the case of the trivial representation $`V=`$. Then $`\mathrm{\Psi }_V^T(\lambda ,\mu )=\frac{q^{2(\mu +\rho ,\lambda )}}{\delta _q^T(\lambda )}`$ and $`_{VW}=1`$. We get
$$\underset{\nu }{}(\frac{f(\lambda +\nu )}{f(\lambda )})_{|W[\nu ]}\mathrm{dim}W[\nu ]\frac{q^{2(\mu +\rho ,\lambda +\nu )}}{\delta _q^T(\lambda +\nu )}=\chi _W(q^{2(\mu +\rho )})\frac{q^{2(\mu +\rho ,\lambda )}}{\delta _q^T(\lambda )}.$$
As in \[EV2\], Corollary 2.16 we conclude that one can take $`f(\lambda )=\delta _q^T(\lambda )`$.
####
Theorem 2.1 now follows from (3.7), Proposition 3.2 ii), Proposition 3.3 and from the following easily checked fusion identity :
$$𝕁_T^{1\mathrm{}N}(\lambda )^1\left(_T^{0,1\mathrm{}N}\right)𝕁_T^{1\mathrm{}N}(\lambda +h^{(0)})=(_T^{01}(\lambda +h^{(2\mathrm{}N)}))\mathrm{}(_T^{0N}(\lambda )).$$
(3.10)
## 4 The twisted qKZB equations
####
We will first prove that the twisted qKZB equations hold for two finite-dimensional $`U_q(𝔤)`$-modules $`V`$ and $`W`$. As in the preceding section, we start with several preliminary lemmas.
####
Consider the following function with values in $`WVV^{}W^{}U_q(𝔤)`$, with components labeled as $`1,2,2,1`$ and $`3`$ respectively :
$$Z_{WV}(\lambda ,\mu )=\mathrm{Tr}_{|M_\mu }(\mathrm{\Phi }_{\mu ^{}+h^{(2)}}^W^{30}\mathrm{\Phi }_\mu ^{}^Vq_0^{2\lambda }B_0).$$
###### Lemma 4.1.
We have
$$Z_{WV}(\lambda ,\mu )=(^{32})^1𝒥_T^{12,3}(\lambda )\mathrm{\Psi }_{WV}^T(\lambda +\frac{1}{2}h^{(3)},\mu ).$$
(4.1)
Proof. Moving the R-matrix around the trace and using the cyclicity property, we get
$$\begin{array}{cc}\hfill Z_{WV}(\lambda ,\mu )& =^{31}\mathrm{Tr}_{|M_\mu }(\mathrm{\Phi }_{\mu ^{}+h^{(2)}}^W\mathrm{\Phi }_\mu ^{}^Vq_0^{2\lambda }B_0^{30})\hfill \\ & =^{31}q_{12}^{2\lambda }B_3^1\mathrm{Tr}_{|M_\mu }(\mathrm{\Phi }_{\mu ^{}+h^{(2)}}^W^{30}\mathrm{\Phi }_\mu ^{}^Vq_0^{2\lambda }B_0)\hfill \\ & =^{31}q_{12}^{2\lambda }(B_3^1^{32})B_3^1Z_{WV}(\lambda ,\mu ).\hfill \end{array}$$
On the other hand,
$$\begin{array}{cc}\hfill ^{31}q_{12}^{2\lambda }(B_3^1^{32})B_3^1[(^{32})^1𝒥_T^{12,3}& (\lambda )\mathrm{\Psi }_{WV}^T(\lambda +\frac{1}{2}h^{(3)},\mu )]\hfill \\ & =^{31}q_{12}^{2\lambda }B_{12}𝒥_T^{12,3}(\lambda )\mathrm{\Psi }_{WV}^T(\lambda +\frac{1}{2}h^{(3)},\mu ).\hfill \end{array}$$
From the modified ABRR equation it follows that
$$^{3,12}q_{12}^{2\lambda }B_{12}(𝒥_T^{12,3}(\lambda ))=𝒥_T^{12,3}(\lambda )q_{13}^{\mathrm{\Omega }_𝔩}q_{23}^{\mathrm{\Omega }_𝔩}q_{12}^{2\lambda }.$$
Using the coproduct formula $`^{3,12}=^{32}^{31}`$ and the $`𝔩`$-invariance of $`Z_{VW}(\lambda ,\mu )`$, we see that
$$\begin{array}{cc}\hfill ^{31}q_{12}^{2\lambda }(B_3^1^{32})B_3^1[(^{32})^1𝒥_T^{12,3}(\lambda )& \mathrm{\Psi }_{WV}^T(\lambda +\frac{1}{2}h^{(3)},\mu )]\hfill \\ & =(^{32})^1𝒥_T^{12,3}(\lambda )\mathrm{\Psi }_{WV}^T(\lambda +\frac{1}{2}h^{(3)},\mu ).\hfill \end{array}$$
Thus both sides of (4.1) satisfy the equation $`X=^{31}q_{12}^{2\lambda }B_3^1X`$ and are of the form $`X=q_{32}^{\mathrm{\Omega }_𝔥}q_{12,3}^Z\mathrm{\Psi }_{VW}^T(\lambda +\frac{1}{2}h^{(3)},\mu )+l.o.t`$. But it is easy to see that such an $`X`$ is unique and the lemma follows.$`\mathrm{}`$
####
Now set
$$\stackrel{~}{Z}_{WV}(\lambda ,\mu )=m_{32}S_3(Z_{WV}(\lambda ,\mu ))=\mathrm{Tr}_{|M_\mu }(\mathrm{\Phi }_{\mu ^{}+h^2}^W(^{20})^1\mathrm{\Phi }_\mu ^Vq_0^{2\lambda }B_0).$$
###### Lemma 4.2.
We have
$$\stackrel{~}{Z}_{WV}(\lambda ,\mu )=S(u)_2^1_T(\lambda \frac{1}{2}h^{(1)})_2𝒥_T^1(\lambda \frac{1}{2}\widehat{h}^{(2)})\mathrm{\Psi }_{WV}^T(\lambda \frac{1}{2}\widehat{h}^{(2)},\mu ).$$
(4.2)
Proof. From (4.1) it follows that
$$\begin{array}{cc}\hfill \stackrel{~}{Z}_{VW}(\lambda ,\mu )& =m_{32}S_3\left[(^{32})^1𝒥_T^{12,3}(\lambda )\mathrm{\Psi }_{WV}^T(\lambda +\frac{1}{2}h^{(3)})\right]\hfill \\ & =\underset{ij}{}c_j^{(1)}S(d_j(\lambda ))S(a_i^{})b_i^{}c_j^{(2)}\mathrm{\Psi }_{WV}^T(\lambda \frac{1}{2}\widehat{h}^{(2)},\mu ).\hfill \end{array}$$
To conclude the proof of the lemma we use the following relations : $`_iS(a_i^{})b_i^{}=S(u^1)`$, $`S(u^1)x=S^2(x)S(u^1)`$ for all $`xU_q(𝔤)`$ and
$$\underset{j}{}c_j^{(1)}S^1(d_j(\lambda ))c_j^{(2)}=_T(\lambda \frac{1}{2}h^{(1)})_2𝒥_T^1(\lambda \frac{1}{2}\widehat{h}^{(2)}).$$
This last equation is obtained by applying $`m_{32}(11S^1)`$ to the cocycle identity (1.5).$`\mathrm{}`$
####
Consider the following function with values in $`VV^{}U_q(𝔤)`$ :
$$Y_V(\lambda ,\mu )=\mathrm{Tr}_{|M_\mu }(\mathrm{\Phi }_\mu ^{}^V(^{02})^1q_0^{2\lambda }B_0).$$
###### Lemma 4.3.
We have
$$Y_V(\lambda ,\mu )=q_1^{2\lambda }𝒥_T^{21}(\lambda )\mathrm{\Psi }_V^T(\lambda \frac{1}{2}h^{(2)},\mu ).$$
(4.3)
Proof. A computation similar to the one in Lemma 4.1 shows that both $`Y_V(\lambda ,\mu )`$ and $`q_1^{2\lambda }𝒥_T^{21}(\lambda )\mathrm{\Psi }_V^T(\lambda \frac{1}{2}h^{(2)},\mu )`$ satisfy the equation
$$X=(^{12})^1q_1^{2\lambda }B_2^1X$$
and are of the form $`X=q_1^{2\lambda }q_{12}^Z\mathrm{\Psi }_V^T(\lambda \frac{1}{2}h^{(2)},\mu )+l.o.t`$. It is easy to see that such an $`X`$ is unique.$`\mathrm{}`$
####
We will also need the following two-representations analogue of $`Y_V(\lambda ,\mu )`$ :
$$Y_{WV}(\lambda ,\mu )=\mathrm{Tr}_{|M_\mu }(\mathrm{\Phi }_{\mu ^{}+h^{(2)}}^W\mathrm{\Phi }_\mu ^{}^V(^{03})^1q_0^{2\lambda }B_0).$$
###### Lemma 4.4.
We have
$$Y_{WV}(\lambda ,\mu )=(^{12,3})^1𝒥_T^{3,12}(\lambda )\mathrm{\Psi }_{WV}^T(\lambda \frac{1}{2}h^{(3)},\mu ).$$
(4.4)
Proof. One checks that both sides of (4.4) are solutions of the equation
$$X=(^{12,3})^1q_{12}^{2\lambda }B_3^1X$$
of the form $`X=q_{12}^{\mathrm{\Omega }_𝔥}q_{13}^{\mathrm{\Omega }_𝔥}q_{31}^Zq_{32}^Z\mathrm{\Psi }_{WV}^T(\lambda \frac{1}{2}h^{(3)})+l.o.t`$.$`\mathrm{}`$
####
Finally, we introduce a function :
$$\stackrel{~}{Y}_{WV}(\lambda ,\mu )=m_{32}(q_2^{2\rho }S_3(Y_{WV}(\lambda ,\mu )))=\mathrm{Tr}_{|M_\mu }(\mathrm{\Phi }_{\mu ^{}+h^{(2)}}^Wm_{32}(q_2^{2\rho }\mathrm{\Phi }_\mu ^{}^V^{03})q_0^{2\lambda }B_0).$$
###### Lemma 4.5.
We have
$$\stackrel{~}{Y}_{WV}(\lambda ,\mu )=(q^{2\rho }u^1S(_T)(\lambda \frac{1}{2}h))_2(^{21}𝒥_T^1\mathrm{\Psi }_{WV}^T)(\lambda +\frac{1}{2}\widehat{h}^{(2)},\mu ).$$
(4.5)
Proof. By (4.4) we have
$$\stackrel{~}{Y}_{WV}(\lambda ,\mu )=\left(\underset{ij}{}(a_i^{})^{(1)}d_j^{(1)}(\lambda )S(c_j)S(b_i^{})q^{2\rho }(a_i^{})^{(2)}d_j^{(2)}(\lambda )\right)\mathrm{\Psi }_{WV}^T(\lambda +\frac{1}{2}\widehat{h}^{(2)},\mu ).$$
Now we use the following relations: $`q^{2\rho }x=S^2(x)q^{2\rho }`$ for any $`xU_q(𝔤)`$,
$$\underset{i}{}(a_i^{})^{(1)}S^1(b_i^{})(a_i^{})^{(2)}=\underset{i}{}a_i^{}S^1(b_i^{})u^1$$
(4.6)
and $`(1S)^1=`$. We get
$$\stackrel{~}{Y}_{WV}(\lambda ,\mu )=(q^{2\rho }u^1)_2\left(\underset{ij}{}a_id_j^{(1)}(\lambda )S(c_j)b_id_j^{(2)}(\lambda )\right)\mathrm{\Psi }_{WV}^T(\lambda +\frac{1}{2}\widehat{h}^{(2)},\mu ).$$
Using the identity $`^{12}𝒥_T^{3,12}(\lambda )=𝒥_T^{3,21}(\lambda )^{12}`$ and
$$\underset{i}{}S(c_i)d_i^{(1)}(\lambda )d_i^{(2)}(\lambda )=S(_T)(\lambda \frac{1}{2}h^{(2)})_1𝒥_T^1(\lambda \frac{1}{2}\widehat{h}^{(1)})$$
which is obtained by applying $`m_{12}(S11)`$ to (1.5), we can further simplify $`\stackrel{~}{Y}_{WV}(\lambda ,\mu )`$ :
$$\begin{array}{cc}& \stackrel{~}{Y}_{WV}(\lambda ,\mu )\hfill \\ & =(q^{2\rho }u^1)_2S(_T)(\lambda \frac{1}{2}h^{(1)})_2(𝒥_T^{21})^1(\lambda +\frac{1}{2}\widehat{h}^{(2)})^{12}\mathrm{\Psi }_{WV}^T(\lambda +\frac{1}{2}\widehat{h}^{(2)})\hfill \\ & =(q^{2\rho }u^1)_2S(_T)(\lambda \frac{1}{2}h^{(1)})_2_T^{21}(\lambda +\frac{1}{2}\widehat{h}^{(2)})𝒥_T^1(\lambda +\frac{1}{2}\widehat{h}^{(2)})\mathrm{\Psi }_{WV}^T(\lambda +\frac{1}{2}\widehat{h}^{(2)},\mu ).\hfill \end{array}$$
which proves the Lemma. $`\mathrm{}`$
####
We need one last technical result :
###### Lemma 4.6.
We have
$$m_{21}(q_1^{2\rho }\mathrm{\Phi }_\mu ^{}^V^{02})=q_1^{2(\mu ^{}+\rho )_ix_i^2}(^{10})^1\mathrm{\Phi }_\mu ^{}^V$$
(4.7)
where $`(x_i)_i`$ is an orthonormal basis of $`𝔥`$.
Proof. Let $`Z=u^1q^{2\rho }`$. This is a ribbon element of $`U_q(𝔤)`$ (see \[D\]). Thus
$$\mathrm{\Phi }_\mu ^{}^VZ=\mathrm{\Delta }(Z)\mathrm{\Phi }_\mu ^{}^V=^{21}(ZZ)\mathrm{\Phi }_\mu ^{}^V.$$
But by (4.6) we have
$$\begin{array}{cc}\hfill m_{21}(q_1^{2\rho }\mathrm{\Phi }_\mu ^{}^V^{02})=& m_{21}(q_1^{2\rho }^{02}^{12})\mathrm{\Phi }_\mu ^{}^V\hfill \\ \hfill =& ^{01}m_{21}(q_1^{2\rho }q_1^{2\rho })q_1^{2\rho }\mathrm{\Phi }_\mu ^{}^V\hfill \\ \hfill =& (1Z)\mathrm{\Phi }_\mu ^{}^V.\hfill \end{array}$$
On the other hand, it is easy to see that $`Z_{|M_\nu }=q^{(2\rho +\nu ,\nu )}`$. The Lemma now follows by a direct computation.$`\mathrm{}`$
####
We are now in position to prove Theorem 2.2. From (4.7), we have
$$\begin{array}{cc}\hfill (q^{2(\mu ^{}+\rho )_ix_i^2})_2\mathrm{Tr}_{|M_\mu }(\mathrm{\Phi }_{\mu ^{}+h^2}^W(& ^{20})^1\mathrm{\Phi }_\mu ^{}^Vq_0^{2\lambda }B_0)\hfill \\ & =\mathrm{Tr}_{|M_\mu }(\mathrm{\Phi }_{\mu ^{}+h^2}^Wm_{32}(q_2^{2\rho }\mathrm{\Phi }_\mu ^{}^V^{03})q_0^{2\lambda }B_0).\hfill \end{array}$$
In other words,
$$(q^{2(\mu ^{}+\rho )_ix_i^2})_2\stackrel{~}{Z}_{WV}(\lambda ,\mu )=\stackrel{~}{Y}_{WV}(\lambda ,\mu ).$$
Using (4.2), (4.5), the relation $`uS(u^1)=q^{4\rho }`$, Proposition 3.3 and the definition of $`\phi _{WV}^T(\lambda )`$ we finally obtain
$$(q^{2(\mu ^{}+\rho )_ix_i^2})_2\phi _{WV}^T(\lambda h^{(2)},\mu )=_T^{21}(\lambda )\phi _{WV}^T(\lambda ,\mu ).$$
(4.8)
From this we derive the qKZB equation with N representations in the following way. We start with the easily checked fusion identities
$`𝕁_T^{23}(\lambda )^1_T^{1,23}(\lambda )𝕁_T^{23}(\lambda +h^{(1)})`$ $`=_T^{12}(\lambda +h^{(3)})_T^{13}(\lambda ),`$ (4.9)
$`𝕁_T^{12}(\lambda +h^{(3)})^1_T^{12,3}(\lambda )𝕁_T^{12}(\lambda )`$ $`=_T^{23}(\lambda )_T^{13}(\lambda +h^{(2)}).`$ (4.10)
Now, from (4.8) with $`W=V_1\mathrm{}V_j`$ and $`V=V_{j+1}\mathrm{}V_N`$ we get
$$\begin{array}{cc}\hfill (& q^{2(\mu +\rho )+{\scriptscriptstyle x_i^2}})_{j+1\mathrm{},N}(q^{2{\scriptscriptstyle x_ix_i}})_{j+1,\mathrm{},N,1,\mathrm{},j}𝕁_T^{1\mathrm{},j}(\lambda )𝕁_T^{j+1\mathrm{}N}(\lambda h^{(j+1\mathrm{}N)})\hfill \\ & \times \phi _{V_1,\mathrm{}V_N}^T(\lambda h^{(j+1\mathrm{}N)},\mu )\hfill \\ & =_T^{j+1\mathrm{}N,1\mathrm{}j}(\lambda )𝕁_T^{1\mathrm{}j}(\lambda +h^{(j+1\mathrm{}N)})𝕁_T^{j+1\mathrm{}N}(\lambda )\phi _{V_1,\mathrm{}V_N}^T(\lambda ,\mu ).\hfill \end{array}$$
(4.11)
By (4.9) this implies that
$$\begin{array}{cc}\hfill (& q^{2(\mu +\rho )+{\scriptscriptstyle x_i^2}})_{j+1\mathrm{}N}(q^{2{\scriptscriptstyle x_ix_i}})_{j+1,\mathrm{}N,1,\mathrm{}j}𝕁_T^{1\mathrm{}j1}(\lambda +h^{(j)})𝕁_T^{j+1\mathrm{}N}(\lambda h^{(j+1\mathrm{}N)})\hfill \\ & \times \phi _{V_1,\mathrm{},V_N}^T(\lambda h^{(j+1\mathrm{}N)},\mu )\hfill \\ & =_T^{j+1\mathrm{}N,1\mathrm{}j1}(\lambda +h^{(j)})_T^{j+1\mathrm{}N,j}(\lambda )𝕁_T^{1\mathrm{}j1}(\lambda +h^{(j\mathrm{}N)})𝕁_T^{j+1\mathrm{}N}(\lambda )\phi _{V_1,\mathrm{},V_N}^T(\lambda ,\mu ).\hfill \end{array}$$
(4.12)
On the other hand, by (4.9) with $`W=V_1\mathrm{}V_{j1}`$ and $`V=V_j\mathrm{}V_N`$ and using (4.10) we also have
$$\begin{array}{cc}\hfill (& q^{2(\mu +\rho )+{\scriptscriptstyle x_i^2}})_{j\mathrm{}N}(q^{2{\scriptscriptstyle x_ix_i}})_{j,\mathrm{}N,1,\mathrm{}j1}𝕁_T^{1\mathrm{}j1}(\lambda )𝕁_T^{j+1\mathrm{}N}(\lambda +h^{(1\mathrm{}j1)})\hfill \\ & \times \phi _{V_1,\mathrm{},V_N}^T(\lambda h^{(j\mathrm{}N)},\mu )\hfill \\ & =_T^{j+1\mathrm{}N,1\mathrm{}j1}(\lambda )_T^{j,1\mathrm{}j1}(\lambda +h^{(j+1\mathrm{}N)})𝕁_T^{1\mathrm{}j1}(\lambda +h^{(j\mathrm{}N)})𝕁_T^{j+1\mathrm{}N}(\lambda )\hfill \\ & \times \phi _{V_1,\mathrm{},V_N}^T(\lambda ,\mu ).\hfill \end{array}$$
(4.13)
Applying the operator $`\mathrm{\Gamma }_j`$ to both sides of (4.13) we get
$$\begin{array}{cc}\hfill (& q^{2(\mu +\rho )+{\scriptscriptstyle x_i^2}})_{j\mathrm{}N}(q^{2{\scriptscriptstyle x_ix_i}})_{j,\mathrm{}N,1,\mathrm{}j1}𝕁_T^{1\mathrm{}j1}(\lambda +h^{(j)})𝕁_T^{j+1\mathrm{}N}(\lambda +h^{(1\mathrm{}j)})\hfill \\ & \times \phi _{V_1,\mathrm{},V_N}^T(\lambda h^{(j+1\mathrm{}N)},\mu )\hfill \\ & =_T^{j+1\mathrm{}N,1\mathrm{}j1}(\lambda +h^{(j)})\mathrm{\Gamma }_j_T^{j,1\mathrm{}j1}(\lambda +h^{(j+1\mathrm{}N)})𝕁_T^{1\mathrm{}j1}(\lambda +h^{(j\mathrm{}N)})𝕁_T^{j+1\mathrm{}N}(\lambda )\hfill \\ & \times \phi _{V_1,\mathrm{},V_N}^T(\lambda ,\mu ).\hfill \end{array}$$
(4.14)
Comparing (4.12) with (4.14) and using the $`𝔩`$-invariance of $``$ we obtain
$$\begin{array}{cc}\hfill (& q^{2(\mu +\rho )+{\scriptscriptstyle x_i^2}})_j(q^{2{\scriptscriptstyle x_ix_i}})_{j,1\mathrm{}j1}𝕁_T^{1\mathrm{}j1}(\lambda +h^{(j\mathrm{}N)})_T^{j+1\mathrm{}N,j}(\lambda )𝕁_T^{j+1\mathrm{}N}(\lambda )\hfill \\ & \times \phi _{V_1,\mathrm{},V_N}^T(\lambda ,\mu )\hfill \\ & =𝕁_T^{j+1\mathrm{}N}(\lambda +h^{(j)})\mathrm{\Gamma }_j_T^{j,1\mathrm{}j1}(\lambda +h^{(j+1\mathrm{}N)})𝕁_T^{1,\mathrm{}j1}(\lambda +h^{(j\mathrm{}N)})\phi _{V_1,\mathrm{},V_N}^T(\lambda ,\mu ),\hfill \end{array}$$
(4.15)
which can be rewritten as
$$\begin{array}{cc}\hfill (& q^{2(\mu +\rho )+{\scriptscriptstyle x_i^2}})_j(q^{2{\scriptscriptstyle x_ix_i}})_{j,1\mathrm{}j1}𝕁_T^{j+1\mathrm{}N}(\lambda +h^{(j)})^1_T^{j+1\mathrm{}N,j}(\lambda )𝕁_T^{j+1\mathrm{}N}(\lambda )\hfill \\ & \times \phi _{V_1,\mathrm{},V_N}^T(\lambda ,\mu )\hfill \\ & =\mathrm{\Gamma }_j𝕁_T^{1\mathrm{}j1}(\lambda +h^{(j+1\mathrm{}N)})^1_T^{j,1\mathrm{}j1}(\lambda +h^{(j+1\mathrm{}N)})𝕁_T^{1,\mathrm{}j1}(\lambda +h^{(j\mathrm{}N)})\hfill \\ & \times \phi _{V_1,\mathrm{},V_N}^T(\lambda ,\mu ).\hfill \end{array}$$
(4.16)
Finally, taking into account identities (4.9) and (4.10), we obtain
$$\begin{array}{cc}& (q^{2(\mu +\rho )+{\scriptscriptstyle x_i^2}})_j(q^{2{\scriptscriptstyle x_ix_i}})_{j,1\mathrm{}j1}_T^{Nj}(\lambda )\mathrm{}_T^{j+1,j}(\lambda +h^{(j+2\mathrm{}N)})\phi _{V_1,\mathrm{},V_N}^T(\lambda ,\mu )\hfill \\ & =\mathrm{\Gamma }_j_T^{j1}(\lambda +h^{(2\mathrm{}j1)}+h^{(j+1\mathrm{}N)})\mathrm{}_T^{jj1}(\lambda +h^{(j+1\mathrm{}N)})\phi _{V_1,\mathrm{},V_N}^T(\lambda ,\mu ).\hfill \end{array}$$
(4.17)
The proof of Theorem 2.2 is now obtained by replacing $`\mu `$ by $`\mu \rho `$ and by rewriting (4.17) in terms of $`F_{V_1,\mathrm{},V_N}^T(\lambda ,\mu )`$.
## 5 The twisted dual Macdonald-Ruijsenaars equation
####
In this section we let $`(\mathrm{\Gamma }_1,\mathrm{\Gamma }_2,T)`$ be a complete generalized Belavin-Drinfeld triple. Let $`W`$ be a finite-dimensional $`U_q(𝔤)`$-module such that $`WW^B`$ and let us consider $`W`$ as a $`BU_q(𝔤)`$-module as in Section 2.
####
For generic values of $`\mu `$, the tensor product $`M_\mu W`$ decomposes as a direct sum of Verma modules, and
$`\eta _\nu :W[\nu ]M_{\mu +\nu }`$ $`M_\mu W`$
$`wy`$ $`\mathrm{\Phi }_{\mu +\nu }^w(y)`$
is an isomorphism onto the isotypic component corresponding to $`M_{\mu +\nu }`$. The following lemma is straightforward :
###### Lemma 5.1.
We have $`(BB)\eta _\nu =\eta _{B(\nu )}(BB)`$.
Now let $`V`$ be any finite-dimensional $`U_q(𝔤)`$-module and consider the composition
$$P_{VV^{},W}^{VW}\mathrm{\Phi }_{B(\mu )}^V(BB)\eta _\nu :W[\nu ]M_{\mu +\nu }M_{B(\mu )+h^{(V)}}WVV^{}.$$
By \[EV2\], Proposition 3.1, we have
$$P_{VV^{},W}^{VW}\mathrm{\Phi }_{B(\mu )}^V\eta _\nu =R^{WV}(B(\mu +\nu ))^{t_2}\mathrm{\Phi }_{B(\mu +\nu )}^V.$$
It follows from Lemma 5.1 that
$$P_{VV^{},W}^{VW}\mathrm{\Phi }_{B(\mu )}^V(BB)\eta _\nu =\eta _{h^{(W)}}R^{WV}(B(\mu +\nu ))^{t_2}\mathrm{\Phi }_{B(\mu +\nu )}^V(BB)$$
(5.1)
where $`R(\lambda )=R_{Id}(\lambda )`$ is the quantum dynamical R-matrix corresponding to the trivial triple $`(\mathrm{\Gamma },\mathrm{\Gamma },Id)`$ and where $`t_2`$ means transposition in the second component (so that $`R^{WV}(B(\mu +\nu ))^{t_2}`$ acts on $`WV^{}`$). Now let us multiply both sides of (5.1) by $`q_{M_\mu W}^{2\lambda }`$ and sum over all values of $`\nu `$. This yields
$$P_{VV^{},W}^{VW}\mathrm{\Phi }_{B(\mu )}^V(BB)q_{M_\mu W}^{2\lambda }=\eta R^{WV}(B(\mu +h^{(W)}))(BB)q^{2\lambda }\eta ^1$$
(5.2)
where $`\eta =_\nu \eta _\nu :_\nu W[\nu ]M_{\mu +\nu }\stackrel{}{}M_\mu W`$. Let us take the trace in the Verma modules and in $`W`$. Using the fact that $`q^{\mathrm{\Omega }_𝔥}U_q(𝔫_+)U_q(𝔫_{})`$ and that $`\nu B(\nu )`$ is never a linear combination of negative roots, we obtain
$$\mathrm{Tr}_{|W}(q^{2\lambda +h^{(V)}}B)\phi _V^T(\lambda ,\mu )=\underset{\nu }{}\mathrm{Tr}_{|W[\nu ]}(R^{WV}(B(\mu +\nu ))^{t_2}B)\phi _V^T(\lambda ,\mu +\nu ).$$
It is clear that
$$\mathrm{Tr}_{|W}(q^{2\lambda +h^{(V)}}B)=\underset{\nu ,B(\nu )=\nu }{}\mathrm{Tr}_{|W[\nu ]}(q^{2\lambda +h^{(V)}}B)=\mathrm{Tr}_{|W^{𝔥_0}}(q^{2\lambda }B).$$
Hence from (5.2) we get
$$\mathrm{Tr}_{|(W^{})^{𝔥_0}}(q^{2\lambda }B^{})\phi _V^T(\lambda ,\mu )=\underset{\nu }{}\mathrm{Tr}_{|W^{}[\nu ]}(B_W^{}^{}R^{WV}(B(\mu +\nu ))^{t_1t_2})\phi _V^T(\lambda ,\mu +\nu ),$$
which can be rewritten in terms of $`F_V^T(\lambda ,\mu )`$ as
$$\begin{array}{cc}& \mathrm{Tr}_{|(W^{})^{𝔥_0}}(q^{2\lambda }B)F_V^T(\lambda ,\mu )\hfill \\ & =\underset{\nu 𝔩^{}}{}\mathrm{Tr}_{|W^{}[\nu ]}\left(_{|V^{}}^1(B(\mu ))B_W^{}^{}^{WV}(B(\mu +\nu ))^{t_1t_2}_{|V^{}}(B(\mu +\nu ))\right)F_V^T(\lambda ,\mu +\nu ).\hfill \end{array}$$
(5.3)
Finally, using the formula
$$_{WV}(\lambda )^{t_1t_2}=((\lambda )(\lambda h^{(1)}))_{W^{}V^{}}(\lambda h^{(1)}h^{(2)})(^1(\lambda h^{(2)})^1(\lambda ))$$
(5.4)
(see \[EV2\], (3.12)) and using the fact that $``$ is of weight zero, we simplify (5.3) to
$$\begin{array}{cc}\hfill \mathrm{Tr}& {}_{|(W^{})^{𝔥_0}}{}^{}(q^{2\lambda })F_V^T(\lambda ,\mu )\hfill \\ & =\underset{\nu }{}\mathrm{Tr}_{|W^{}[\nu ]}(B_W^{}^{}_W^{}(B(\mu +\nu ))_{W^{}V^{}}(B(\mu +\nu )\nu h^{(2)})\times \hfill \\ & \times _W^{}^1(B(\mu +\nu )h^{(2)}))F^T_V(\lambda ,\mu +\nu )\hfill \\ & =\underset{\nu }{}\mathrm{Tr}_{|W^{}[\nu ]}\left(_W^{}(B(\mu +\nu ))_{W^{}V^{}}(\mu )B_W^{}^{}_W^{}^1(B(\mu +\nu ))\right)F_V^T(\lambda ,\mu +\nu )\hfill \\ & =\underset{\nu }{}\mathrm{Tr}_{|W^{}[\nu ]}\left(_{W^{}V^{}}(\mu )B_W^{}^{}\right)F_V^T(\lambda ,\mu +\nu ).\hfill \end{array}$$
(5.5)
The twisted dual Macdonald-Ruijsenaars equations for an arbitrary number of modules $`V_1,\mathrm{}V_N`$ can now be deduced from (5.5) and from the fusion identity (3.10). Theorem 2.3 is proved.
####
Proof of Proposition 2.1. Let $`𝐖`$ and $`\overline{𝐖}`$ be the Weyl groups of $`\mathrm{\Gamma }`$ and $`\overline{\mathrm{\Gamma }}`$ respectively. By the Bernstein-Gelfand-Gelfand resolution, we have
$$\mathrm{Tr}_{|V_\nu }(q^{2\lambda }B)=\underset{w𝐖}{}(1)^{l(w)}\mathrm{Tr}_{|M_{w(\nu +\rho )\rho }}(q^{2\lambda }B).$$
Denote by $`s_\alpha `$ the simple reflection corresponding to the simple root $`\alpha \mathrm{\Gamma }`$. The group generated by $`B`$ acts on $`𝐖`$ by $`B(s_\alpha )=s_{T\alpha }`$. It follows from the facts that $`𝐖`$ acts simply transitively on the sets of simple roots and that $`\nu `$ is dominant that $`B(w(\nu +\rho )\rho )=w(\nu +\rho )\rho `$ if and only if $`B(w)=w`$. Moreover, $`𝐖^B`$ is naturally isomorphic to $`\overline{𝐖}`$. Hence,
$$\begin{array}{cc}\hfill \underset{wW}{}(1)^{l(w)}\mathrm{Tr}_{|M_{w(\nu +\rho )\rho }}(q^{2\lambda }B)& =\underset{wW^B}{}(1)^{l(w)}\mathrm{Tr}_{|M_{w(\nu +\rho )\rho }}(q^{2\lambda }B)\hfill \\ & =\underset{w\overline{W}}{}(1)^{l(w)}\frac{q^{2(\lambda ,w(\nu +\rho )\rho )}}{_{\overline{\alpha }\overline{\mathrm{\Delta }}^+}(1\theta _\alpha q^{2(\overline{\alpha },\lambda )})}\hfill \end{array}$$
(5.6)
Let $`\omega _\alpha `$ be the fundamental weight corresponding to $`\alpha \mathrm{\Gamma }`$. It is easy to check that $`\{\overline{\omega }_{\overline{\alpha }}=\frac{(\overline{\alpha },\overline{\alpha })}{2N_\alpha }_{\alpha \overline{\alpha }}\omega _\alpha ,\overline{\alpha }\overline{\mathrm{\Gamma }}\}`$ is the set of fundamental weights of $`\overline{\mathrm{\Gamma }}`$. Thus $`(2\lambda ,w(\nu +\rho )\rho )=(2\overline{\lambda },w(\nu +\overline{\rho })\overline{\rho })`$ where $`\overline{\rho }=_{\overline{\alpha }}\overline{\omega }_{\overline{\alpha }}`$. Hence, by the Weyl character formula for $`\overline{𝔤}`$ we have
$$\mathrm{Tr}_{|V_\nu }(q^{2\lambda }B)=\chi _{\overline{V}_\nu }(q^{2\overline{\lambda }})\frac{\overline{\delta }_q(\overline{\lambda })}{\delta _q^T(\lambda )}$$
where
$$\overline{\delta }_q(\overline{\lambda })=q^{2(\rho ,\overline{\lambda })}\underset{\overline{\alpha }\overline{\mathrm{\Delta }}^+}{}(1q^{2(\overline{\alpha },\overline{\lambda })})$$
is the (usual) Weyl denominator for $`\overline{\mathrm{\Gamma }}`$. Setting $`\nu =0`$ we see that in fact $`\overline{\delta }_q(\overline{\lambda })=\delta _q^T(\lambda )`$ . The Proposition follows.$`\mathrm{}`$
## 6 The twisted dual qKZB equations
####
In this section we prove Theorem 2.4. As in the preceding section, let $`T`$ be an automorphism of $`\mathrm{\Gamma }`$ and let $`V_1,\mathrm{}V_N`$ be finite-dimensional $`U_q(𝔤)`$-modules. We will extensively use the following two identities which are proved in \[EV3\] :
$$\mathrm{\Phi }_{\mu +h^{(V^{})}}^W\mathrm{\Phi }_\mu ^V=^1R_{21}^{}(\mu )\mathrm{\Phi }_{\mu +h^{(W^{})}}^V\mathrm{\Phi }_\mu ^W=_{21}R^{}(\mu )^1\mathrm{\Phi }_{\mu +h^{(W^{})}}^V\mathrm{\Phi }_\mu ^W$$
(6.1)
for any two modules $`V,W`$.
####
Consider
$$\mathrm{\Psi }_{V_1,\mathrm{},V_N}^T(\lambda ,\mu )=\mathrm{Tr}_{|M_\mu }\left(\mathrm{\Phi }_{B(\mu )+h^{(2\mathrm{}N)}}^{V_1}\mathrm{}\mathrm{\Phi }_{B(\mu )}^{V_N}q^{2\lambda }B\right)$$
and move the $`j`$th intertwiner to the right using (6.1). We get
$$\begin{array}{cc}\hfill \mathrm{\Psi }& {}_{V_1,\mathrm{},V_N}{}^{T}(\lambda ,\mu )\hfill \\ & =_{j+1,j}\mathrm{}_{N,j}q_j^{2\lambda }R_{j,j+1}^{}(B(\mu )+h^{(j+2\mathrm{}N)})^1\mathrm{}R_{j,N}^{}(B(\mu ))^1\times \hfill \\ & \times \mathrm{Tr}_{|M_\mu }\left(\mathrm{\Phi }_{B(\mu )+h^{(2\mathrm{}N)}}^{V_1}\mathrm{}\mathrm{\Phi }_{B(\mu )+h^{(j)}}^{V_N}q^{2\lambda }\mathrm{\Phi }_{B(\mu )}^{V_j}B\right).\hfill \end{array}$$
(6.2)
Now, we have
$$\begin{array}{cc}\hfill \mathrm{Tr}_{|M_\mu }(& \mathrm{\Phi }_{B(\mu )+h^{(2\mathrm{}N)}}^{V_1}\mathrm{}\mathrm{\Phi }_{B(\mu )+h^{(j)}}^{V_N}q^{2\lambda }\mathrm{\Phi }_{B(\mu )}^{V_j}B)\hfill \\ & =B_{V_j^{}}B_{V_j^{^{}}}^{}\mathrm{Tr}_{M_{\mu +h^{(j)}}}\left(\mathrm{\Phi }_{B(\mu +h^{(j)})+_{i=1,ij}^Nh^{()i}}^{V_j^{}}\mathrm{}\mathrm{\Phi }_{B(\mu +h^{(j)})}^{V_N}q^{2\lambda }B\right)\hfill \\ & =B_{V_j^{}}B_{V_j^{^{}}}^{}\mathrm{\Gamma }_j^{}\mathrm{Tr}_{M_\mu }\left(\mathrm{\Phi }_{B(\mu )+_{i=1,ij}^Nh^{(i)}}^{V_j^{}}\mathrm{}\mathrm{\Phi }_{B(\mu )}^{V_N}q^{2\lambda }B\right)\hfill \end{array}$$
(6.3)
where we note $`V_j^{}=V_j^{B^1}`$ for simplicity. Finally, we move $`\mathrm{\Phi }^{V_j^{}}`$ to the right back to its original position, thereby completing a cycle. By (6.1) we obtain
$$\begin{array}{cc}\hfill \mathrm{Tr}_{|M_\mu }& \left(\mathrm{\Phi }_{B(\mu )+_{i=1,ij}^Nh^{(i)}}^{V_j^{}}\mathrm{}\mathrm{\Phi }_{B(\mu )}^{V_N}q^{2\lambda }B\right)\hfill \\ \hfill =& _{j,1}^1\mathrm{}_{j,j1}^1R_{1,j}^{}\left(B(\mu )+\underset{i=2,ij}{\overset{N}{}}h^{(i)}\right)\mathrm{}R_{j1,j}^{}\left(B(\mu )+\underset{i=j+1}{\overset{N}{}}h^{(i)}\right)\hfill \\ & \times \mathrm{\Psi }_{V_1,\mathrm{}V_j^{},\mathrm{}V_N}^T(\lambda ,\mu ).\hfill \end{array}$$
(6.4)
Combining (6.2), (6.3) and (6.4) yields the following relation
$$\begin{array}{cc}\hfill \mathrm{\Psi }& {}_{V_1,\mathrm{},V_N}{}^{T}(\lambda ,\mu )\hfill \\ \hfill =& \left[_{j+1,j}\mathrm{}_{N,j}q_j^{2\lambda }(B_j_{j,1}^1)\mathrm{}(B_j_{j,j1}^1)\right]\times \hfill \\ & \times [R_{j,j+1}^{}(B(\mu )+\underset{i=j+2}{\overset{N}{}}h^{(i)})^1\mathrm{}R_{j,N}^{}\left(B(\mu )\right)^1\mathrm{\Gamma }_{B^1(j)}^{}\times \hfill \\ & \times B_j^{}R_{1,j}^{}(B(\mu )+\underset{i=2,ij}{\overset{N}{}}h^{(i)})\mathrm{}B_j^{}R_{j1,j}^{}(B(\mu )+\underset{i=j+1}{\overset{N}{}}h^{(i)})]\times \hfill \\ & \times B_{V_j^{}}B_{V_j^{^{}}}^{}\mathrm{\Psi }_{V_1,\mathrm{},V_j^{},\mathrm{}V_N}^T(\lambda ,\mu ).\hfill \end{array}$$
(6.5)
Let us replace $`\mu `$ by $`\mu \rho `$ and let us rewrite this equation in terms of $`F^T(\lambda ,\mu )`$. We get
$$\begin{array}{cc}\hfill F& {}_{V_1,\mathrm{},V_N}{}^{T}(\lambda ,\mu )\hfill \\ \hfill =& \left[𝕁_T^{1\mathrm{}N}(\lambda )^1_{j+1,j}\mathrm{}_{N,j}q_j^{2\lambda }(B_j_{j,1}^1)\mathrm{}(B_j_{j,j1}^1)(B_j𝕁_T^{1\mathrm{}N}(\lambda ))\right]\times \hfill \\ & \times [\{_N^1(B(\mu ))\mathrm{}_1^1(B(\mu )\underset{i=2}{\overset{N}{}}h^{(i)})\}_{j,j+1}^{}(B(\mu )\underset{i=j+2}{\overset{N}{}}h^{(i)})^1\times \mathrm{}\hfill \\ & \times _{j,N}^{}(B(\mu ))^1\mathrm{\Gamma }_{B^1(j)}^1B_j^{}_{1,j}^{}(B(\mu )\underset{i=2,ij}{\overset{N}{}}h^{(i)})\mathrm{}B_j^{}_{j1,j}^{}(B(\mu )\underset{i=j+1}{\overset{N}{}}h^{(i)})\times \hfill \\ & \times B_j\{_N(B(\mu ))\mathrm{}_1(B(\mu )\underset{i=2,ij}{\overset{N}{}}h^{(i)}B^1(h^{(j)}))\}]\times \hfill \\ & \times B_{V_j^{}}B_{V_j^{^{}}}^{}F_{V_1,\mathrm{},V_j^{},\mathrm{},V_N}^T(\lambda ,\mu ).\hfill \end{array}$$
(6.6)
Inverting, we obtain
$$\begin{array}{cc}\hfill B& {}_{V_j^{}}{}^{}B_{V_j^{^{}}}^{}F_{V_1,\mathrm{},V_j^{},\mathrm{},V_N}^T(\lambda ,\mu )\hfill \\ \hfill =& \left[(B_j𝕁_T^{1\mathrm{}N}(\lambda ))^1(B_j_{j,1\mathrm{}j1})q_j^{2\lambda }_{j+1\mathrm{}N,j}^1𝕁_T^{1\mathrm{}N}(\lambda )\right]\times \hfill \\ & \times B_j\{_N^1(B(\mu ))\mathrm{}_1^1(B(\mu )\underset{i=2,ij}{\overset{N}{}}h^{(i)}B^1(h^{(j)}))\}\times \hfill \\ & \times B_j^{}_{j1,j}^1(B(\mu )\underset{i=j+1}{\overset{N}{}}h^{(i)})\mathrm{}\times B_j^{}_{1,j}^1(B(\mu )\underset{i=2,ij}{\overset{N}{}}h^{(i)})\times \hfill \\ & \times \mathrm{\Gamma }_{B^1(j)}^{}_{j,N}^{}(B(\mu ))\mathrm{}_{j,j+1}^{}(B(\mu )\underset{i=j+2}{\overset{N}{}}h^{(i)})\times \hfill \\ & \times \left\{_N(B(\mu ))\mathrm{}_1\left(B(\mu )\underset{i=2}{\overset{N}{}}h^{(i)}\right)\right\}\times F_{V_1,\mathrm{},V_N}^T(\lambda ,\mu ).\hfill \end{array}$$
(6.7)
Using (5.4) and using the fact that $`\mu =B(\mu )h^{(i)}`$ it is easy to check that
$$\begin{array}{cc}\hfill K_j^{,T}& =\hfill \\ & B_j\{_N^1(B(\mu ))\mathrm{}_1^1(B(\mu )\underset{i=2,ij}{\overset{N}{}}h^{(i)}B^1(h^{(j)}))\}\times \hfill \\ & \times B_j^{}_{j1,j}^1(B(\mu )\underset{i=j+1}{\overset{N}{}}h^{(i)})\times \mathrm{}\times B_j^{}_{1,j}^1(B(\mu )\underset{i=2,ij}{\overset{N}{}}h^{(i)})\times \hfill \\ & \mathrm{\Gamma }_{B^1(j)}^{}_{j,N}^{}(B(\mu ))\mathrm{}_{j,j+1}^{}(B(\mu )\underset{i=j+2}{\overset{N}{}}h^{(i)})\times \hfill \\ & \times \left\{_N(B(\mu ))\mathrm{}_1\left(B(\mu )\underset{i=2}{\overset{N}{}}h^{(i)}\right)\right\}.\hfill \end{array}$$
Finally, we have
$$D_j^{,T}=(B_j𝕁_T^{1\mathrm{}N}(\lambda ))^1(B_j_{j,1\mathrm{}j1})q_j^{2\lambda }_{j+1\mathrm{}N,j}^1𝕁_T^{1\mathrm{}N}(\lambda )$$
(6.8)
when applied to $`(V_1\mathrm{}V_N)^𝔩`$. The proof of this last equality is similar to the proof of \[EV2\] (4.10): it is enough to check (6.8) for $`N=3`$, for which it follows from the modified ABRR equation (1.4). This concludes the proof of Theorem 2.4.
## 7 The classical limits
####
Let us now examine the classical limits of Theorems 2.1-4, that is, the behavior of (a suitable renormalization of) the functions $`F_{V_1,\mathrm{},V_N}^T(\lambda ,\mu )`$ when $`\mathrm{}0`$. In that limit, the quantum group $`U_q(𝔤)`$ becomes the usual enveloping algebra $`U(𝔤)`$. We will denote by $`\mathrm{\Phi },_T^c,_T^c,𝕁_T^c,\mathrm{}`$ the classical limits of the operators constructed in Section 2.1, obtained when we replace $`U_q(𝔤)`$ by $`U(𝔤)`$.
####
Let $`V_1,\mathrm{}V_N`$ be finite-dimensional $`U_q(𝔤)`$-modules and let $`V_1^c,\mathrm{}V_N^c`$ be the corresponding $`U(𝔤)`$ modules. Let us fix a generalized Belavin-Drinfeld triple $`(\mathrm{\Gamma }_1,\mathrm{\Gamma }_2,T)`$ and set
$$\mathrm{\Psi }_{V_1,\mathrm{},V_N}^{T,c}=\mathrm{Tr}_{|M_\mu ^c}\left(\mathrm{\Phi }_{\mu ^{}+_{i=2}^Nh^{(i)}}^{V_1^c}\mathrm{}\mathrm{\Phi }_\mu ^{}^{V_N^c}Be^\lambda \right).$$
Also set $`\delta ^T(\lambda )=(\mathrm{Tr}_{|M_\rho ^c}(Be^\lambda ))^1`$. We define the classical limit of the function $`F_{V_1,\mathrm{},V_N}^T(\lambda ,\mu )`$ as
$$F_{V_1,\mathrm{},V_N}^{T,c}(\lambda ,\mu ):=\underset{\mathrm{}0}{\mathrm{lim}}F_{V_1,\mathrm{},V_N}^T(\frac{\lambda }{\mathrm{}},\mu ).$$
The following result is clear from the definitions.
###### Lemma 7.1.
We have
$$\begin{array}{cc}& F_{V_1,\mathrm{},V_N}^{T,c}(\lambda ,\mu )\hfill \\ & =\delta ^T(\lambda )[_N^c(\mu +h^{(1\mathrm{}N)})^1\mathrm{}_1^c(\mu +h^{(1)})^1]\mathrm{\Psi }_{V_1,\mathrm{},V_N}^{T,c}(\lambda ,\mu \rho ).\hfill \end{array}$$
####
The classical analogue of Proposition 3.1 is as follows.
###### Proposition 7.1.
Let $`V`$ be any finite-dimensional $`U(𝔤)`$-module and let $`XU(𝔤)`$.
1. There exists a unique differential operator $`d_X^T`$ acting on functions $`𝔩^{}V^𝔩`$ such that
$$\mathrm{Tr}_{|M_\mu }(\mathrm{\Phi }_\mu ^{}^VXBe^\lambda )=d_X^T\mathrm{Tr}_{|M_\mu }(\mathrm{\Phi }_\mu ^{}^VBe^\lambda ).$$
2. If $`X,Y`$ belong to the center of $`U(𝔤)`$ then $`d_X^Td_Y^T=d_Y^Td_X^T`$.
Unfortunately, there is no convenient classical analogue of the Drinfeld-Reshetikhin construction of central elements in $`U_q(𝔤)`$, and therefore no convenient explicit computation of the operator $`d_X^T`$ in general. However, this can be done when $`X=m_{12}(\mathrm{\Omega })`$ is the quadratic Casimir, which yields the following classical analogue of Theorem 2.1 (which is proved directly in \[ES1\], Theorem 3.2).
###### Theorem 7.1 (\[ES1\]).
The function $`F_{V_1,\mathrm{},V_N}^{T,c}(\lambda ,\mu )`$ satisfies the following second order differential equation :
$$\left(\underset{iI_1}{}\frac{^2}{x_i^2}\underset{l,n=1}{\overset{r}{}}S_T(\lambda )_{|V_lV_n}\right)F_{V_1,\mathrm{}V_N}^{T,c}(\lambda ,\mu )=(\mu ,\mu )F_{V_1,\mathrm{},V_N}^{T,c}(\lambda ,\mu )$$
(7.1)
where $`(x_i)_{iI_1}`$ (resp. $`(x_i)_{iI_2}`$) is an orthonormal basis of $`𝔩`$ (resp. of $`𝔥_0`$) and where
$$\begin{array}{c}\hfill S_T(\lambda )=\underset{\alpha }{}\underset{k=0}{\overset{\mathrm{}}{}}\underset{v=1}{\overset{\mathrm{}}{}}e^{(s+v)(\alpha ,\lambda )}(B^sf_\alpha B^ve_\alpha +B^ve_\alpha B^sf_\alpha )\\ \hfill \underset{iI_2}{}\frac{1C_T}{2}x_i\frac{1C_T}{2}x_i.\end{array}$$
Theorem 7.1 can also be deduced from Theorem 2.1 by expanding powers of $`\mathrm{}`$.
####
The classical limit of Theorem 2.2 are the twisted (trigonometric) KZB equations.
###### Theorem 7.2 (\[ES1\]).
The function $`F_{V_1,\mathrm{},V_N}^{T,c}(\lambda ,\mu )`$ satisfies the following system of differential equations, for $`j=1,\mathrm{}N`$:
$$\begin{array}{cc}\hfill (\underset{iI_1}{}x_{i|V_j}\frac{}{x_i}+& \underset{l>j}{}r_T(\lambda )_{|V_jV_l}\underset{l<j}{}r_T(\lambda )_{|V_lV_j})F^{T,c}_{V_1,\mathrm{},V_N}(\lambda ,\mu )\hfill \\ & =\left((\mu +\frac{1}{2}_𝔥)_{|V_j^{}}+\underset{l=1}{\overset{j1}{}}(\mathrm{\Omega }_𝔥)_{|V_i^{}V_j^{}}\right)F_{V_1,\mathrm{},V_N}^{T,c}(\lambda ,\mu ).\hfill \end{array}$$
(7.2)
This theorem is proved in \[ES1\] but can also be deduced from Theorem 2.2 by expanding in powers of $`\mathrm{}`$.
####
Finally, when $`T`$ is an automorphism of the Dynkin diagram $`\mathrm{\Gamma }`$ we consider classical limits of the dual Macdonald-Ruijsenaars and dual qKZB equations. Let $`W`$ be a $`B`$-invariant finite-dimensional $`𝔤`$-module and let $`𝒟_W^{,T,c}`$ denote the difference operator given by formula (2.7) when $`U_q(𝔤)`$ is replaced by $`U(𝔤)`$.
###### Theorem 7.3.
We have
$$𝒟_W^{,T,c}F_{V_1,\mathrm{},V_N}^{T,c}=\mathrm{Tr}_{|W^{𝔥_0}}(e^\lambda B)F_{V_1,\mathrm{},V_N}^{T,c}.$$
Similarly, let $`K_j^{,T,c}`$ be the classical limit of $`K_j^{,T}`$, i.e the difference operator given by formula (2.9) when $`q=1`$ (and hence $`(\mu )`$ is just the classical exchange matrix evaluated at $`\mu \rho `$, see \[EV3\]).
###### Theorem 7.4.
For $`j=1,\mathrm{}N`$ we have
$$B_{V_j}B_{V_j^{}}^{}F_{V_1,\mathrm{},V_N}^{T,c}=(e^\lambda )_{|V_j}K_j^{,T,c}F_{V_1,\mathrm{},V_j^B,\mathrm{},V_N}^{T,c}.$$
## 8 Extension to Kac-Moody algebras
####
In this section we briefly explain how to adapt the construction of \[ESS\] to Kac-Moody algebras and how to generalize Theorems 2.1-2.4 to this setting.
####
Let $`A=(a_{ij})`$ be a symmetrizable generalized Cartan matrix of size $`n`$ and rank $`l`$. Let $`(𝔥,\mathrm{\Gamma },\stackrel{ˇ}{\mathrm{\Gamma }})`$ be a realization of $`A`$, i.e $`𝔥`$ is a complex vector space of dimension $`2nl`$, $`\mathrm{\Gamma }=\{\alpha _1,\mathrm{}\alpha _n\}𝔥^{}`$ and $`\stackrel{ˇ}{\mathrm{\Gamma }}=\{h_1,\mathrm{}h_n\}𝔥`$ are linearly independent sets and $`\alpha _j,h_i=a_{ij}`$. Let $`𝔤=𝔫_{}𝔥𝔫_+`$ be the Kac-Moody algebra associated to $`A`$, i.e $`𝔤`$ is generated by elements $`e_i`$, $`f_i`$, $`i=1,\mathrm{}n`$ and $`𝔥`$ with relations
$$[e_i,f_j]=\delta _{ij}h_i,[𝔥,𝔥]=0,[h,e_i]=\alpha _i,he_i,[h,f_i]=\alpha _i,hf_i,$$
together with the Serre relations (see \[K\]). Let $`(,)`$ be a nondegenerate invariant bilinear form on $`𝔤`$. Let $`\mathrm{\Omega }_𝔥`$ be the inverse element to the restriction $`(,)`$ to $`𝔥`$. For every root $`\alpha 𝔥^{}`$ we set $`\alpha ^{}=(1\alpha )\mathrm{\Omega }_𝔥`$.
####
Let $`U_q(𝔤)`$ be the quantum Kac-Moody algebra. It is defined by the same relations as in Section 1, where now $`(a_{ij})`$ is the generalized Cartan matrix $`A`$.
#### Construction of the twist.
Let $`(\mathrm{\Gamma }_1,\mathrm{\Gamma }_2,T)`$ be a generalized Belavin-Drinfeld triple. As before we set $`𝔩=\left(_{\alpha \mathrm{\Gamma }_1}(\alpha T\alpha )\right)^{}`$ and $`𝔥_0=𝔩^{}𝔥`$. We will say that $`(\mathrm{\Gamma }_1,\mathrm{\Gamma }_2,T)`$ is nondegenerate if the restriction of $`(,)`$ to $`𝔩`$ is, and we make this assumption from now on. Let $`𝔥_1𝔥`$ (resp. $`𝔥_2𝔥`$) be the subspace spanned by simple roots $`\alpha \mathrm{\Gamma }_1`$ (resp. $`\alpha \mathrm{\Gamma }_2`$).
####
In \[ESS\] we obtained an explicit construction of a twist $`𝒥_T(\lambda )`$ for simple complex Lie algebras. An important observation there was that $`𝔥=𝔥_1+𝔩`$, which makes it is possible to extend $`B`$ to an orthogonal automorphism of $`𝔥`$, and to define maps $`B^{\pm 1}:U_q(𝔟_{})U_q(𝔟_{})`$. However, in general we only have $`𝔥_1+𝔩𝔥`$ but $`𝔥_1+𝔩𝔥`$, and thus it is necessary to modify the construction in \[ESS\], which is done below.
The following lemma is obvious.
###### Lemma 8.1.
There exist unique algebra morphism $`B:U_q(𝔫_{}𝔥_1)U_q(𝔫_{}𝔥_2)`$ and $`B:U_q(𝔫_+𝔥_2)U_q(𝔫_{}𝔥_1)`$ such that $`B(F_\alpha )=F_{T\alpha }`$, $`B(h_\alpha )=h_{T\alpha }`$ if $`\alpha \mathrm{\Gamma }_1`$, $`B(F_\alpha )=0`$ if $`\alpha \mathrm{\Gamma }\mathrm{\Gamma }_1`$, and $`B^1(E_\alpha )=E_{T^1\alpha }`$, $`B^1(h_\alpha )=h_{T^1\alpha }`$ if $`\alpha \mathrm{\Gamma }_2`$, $`B^1(E_\alpha )=0`$ if $`\alpha \mathrm{\Gamma }\mathrm{\Gamma }_2`$.
####
Let $`\alpha ,\beta \mathrm{\Gamma }`$. Write $`\alpha \beta `$ if there exists $`l0`$ such that $`T^l(\alpha )=\beta `$. We extend this relation to $`^+\mathrm{\Gamma }`$ by setting $`\alpha \beta `$ if there exists $`\alpha _1,\mathrm{}\alpha _r,\beta _1,\mathrm{}\beta _r\mathrm{\Gamma }`$ such that $`\alpha _i\beta _i`$ for $`i=1,\mathrm{}r`$ and $`\alpha =_i\alpha _i,\beta =_i\beta _i`$. It is easy to see that this relation is transitive, i.e if $`\alpha \beta `$ and $`\beta \gamma `$ then $`\alpha \gamma `$. Set
$$^+\mathrm{\Gamma }_\alpha =\{\sigma ^+\mathrm{\Gamma },\sigma \alpha \},^+\mathrm{\Gamma }_\alpha =\{\sigma ^+\mathrm{\Gamma },\alpha \sigma \}.$$
Now let us consider the space
$$I_T=\underset{\beta \alpha }{}\left(U_q(𝔫_{})[\alpha ]q^{(^+\mathrm{\Gamma }_\alpha )^{}}U_q(𝔫_+)[\beta ]q^{(^+\mathrm{\Gamma }_\beta )^{}}\right)U_q(𝔟_{})U_q(𝔟_+).$$
###### Lemma 8.2.
The space $`I_T`$ is stable under the actions of $`B1`$, $`1B^1`$ and $`Ad(q^{\mathrm{\Omega }_𝔥})`$.
Proof. Note that the actions of $`(B1)`$ and $`(1B^1)`$ are well-defined on $`I_T`$ as $`B(U_q(𝔫_{})[\alpha ])=0`$ if $`\alpha ^+\mathrm{\Gamma }_1`$ and $`B^1(U_q(𝔫_+)[\beta ])=0`$ if $`\beta ^+\mathrm{\Gamma }_2`$. It is clear that $`(B1)I_TI_T`$ and $`(1B^1)I_TI_T`$. The last claim in the Lemma follows easily from the formula
$$Ad(q^{\mathrm{\Omega }_𝔥})(u_\alpha v_\beta )=u_\alpha q^\beta ^{}q^\alpha ^{}v_\beta $$
if $`u_\alpha U_q(𝔫_{})[\alpha ]`$ and $`v_\beta U_q(𝔫_+)[\beta ]`$.$`\mathrm{}`$
####
Note that the Cayley transform $`C_T:𝔥_0𝔥_0`$ is still well-defined in the Kac-Moody setting. Set $`Z=\frac{1}{2}((1C_T)1)\mathrm{\Omega }_{𝔥_0}`$. Let $`\overline{I}_T`$ be the completion of $`I_T`$ with respect to the principal gradings in $`U_q(𝔟_\pm )`$ and let $`\overline{I}_T^{}`$ be the subspace consisting of elements of strictly negative degree in the first component and strictly positive degree in the second component.
###### Theorem 8.1.
There exists a unique element $`𝒥_T^0(\lambda ):𝔩^{}(1+\overline{I}_T^{})^𝔩`$ such that
$$^{21}q_1^{2\lambda }B_1𝒥_T^0(\lambda )=𝒥_T^0(\lambda )q_1^{2\lambda }q^{\mathrm{\Omega }_𝔥}.$$
(8.1)
Moreover $`𝒥_T(\lambda ):=𝒥_T^0(\lambda )q^Z`$ satisfies the 2-cocycle relation
$$𝒥_T^{12,3}(\lambda )𝒥_T^{12}(\lambda +\frac{1}{2}h^{(3)})=𝒥_T^{1,23}(\lambda )𝒥_T^{23}(\lambda \frac{1}{2}h^{(1)}).$$
Proof. The first statement is proved exactly as in \[ESS\]. We write $`𝒥_T^0(\lambda )=1+_{ji}𝒥_T^{0,i}(\lambda )`$ where $`𝒥_T^{0,i}(\lambda )`$ has degree $`i`$ in the first component. Then (8.1) is equivalent to a system of equations labelled by $`j1`$
$$Ad(q^{\mathrm{\Omega }_𝔥}q_1^{2\lambda })B_1𝒥_T^{0,j}(\lambda )=𝒥_T^{0,j}(\lambda )+\mathrm{}$$
where $`\mathrm{}`$ stands for terms involving $`𝒥_T^{0,i}(\lambda )`$ with $`i<j`$. But the operator $`Ad(q^{\mathrm{\Omega }_𝔥}q_1^{2\lambda })B_11`$ is invertible on $`I_T^𝔩`$ for generic $`\lambda `$ and $`𝒥_T^{0,j}(\lambda )`$ can be computed recursively.
The second claim is proved as \[ESS\], Section 4. We consider the three components versions of (8.1)
$$^{21}^{31}q_1^{2\lambda }B_1X_T^0(\lambda )=X_T^0(\lambda )q_{12}^{\mathrm{\Omega }_𝔥}q_{13}^{\mathrm{\Omega }_𝔥},$$
(8.2)
$$^{32}^{21}q_3^{2\lambda }B_3^1X_T^0(\lambda )=X_T^0(\lambda )q_{12}^{\mathrm{\Omega }_𝔥}q_{13}^{\mathrm{\Omega }_𝔥},$$
(8.3)
acting on (a suitable completion of) the space
$$\underset{\alpha ,\beta ,\gamma }{}\left(U_q(𝔫_{})[\alpha ]q^{(^+\mathrm{\Gamma }_\alpha )^{}}U_q(𝔤)[\beta ]U_q(𝔫_+)[\gamma ]q^{(^+\mathrm{\Gamma }_\gamma )^{}}\right)$$
where the sum runs over all triples $`(\alpha ,\beta ,\gamma )`$ such that $`\beta `$ can be written as $`\beta =\beta ^+\beta ^{}`$ where $`\beta ^++\gamma \beta ^{}+\alpha `$. It is not difficult to show that $`(𝒥_T^0(\lambda ))^{1,23}\mathrm{Ad}q_{1,23}^Z(𝒥_T^0(\lambda +\frac{1}{2}h^{(1)}))^{23}`$ and $`(𝒥_T^0(\lambda ))^{12,3}\mathrm{Ad}q_{12,3}^Z(𝒥_T^0(\lambda +\frac{1}{2}h^{(3)}))^{12}`$ are two solutions of (8.2) and (8.3) with the same degree zero terms (in component $`1`$ or in component $`3`$). This implies that they are equal (see \[ESS\], Lemma 4.3).$`\mathrm{}`$
####
Now let $`V_1,\mathrm{},V_N`$ be $`U_q(𝔤)`$-modules from the category $`𝒪`$. Define the renormalized twisted traces functions $`F_{V_1,\mathrm{},V_N}^T(\lambda ,\mu )`$ in the same way as in Section 2. Note that all the operators $`𝒥_T(\lambda )`$, $`_T(\lambda )`$, $`_T(\lambda )`$,… are well-defined on any module from the category $`𝒪`$ when considered as formal powers series in $`q^{2(\lambda ,\mu )}[[q^{(\lambda ,\alpha _i)},q^{(\mu ,\alpha _i)}]]`$, $`\alpha _i\mathrm{\Gamma }`$. Operators $`𝒟_W`$ for affine algebras $`𝔤`$ are defined in some particular situation in \[E3\].
###### Theorem 8.2.
The function $`F_{V_1,\mathrm{},V_N}^T(\lambda ,\mu )`$ satisfies the following difference equation for all $`j=1,\mathrm{},N`$ :
$$F_{V_1,\mathrm{},V_N}^T(\lambda ,\mu )=(D_j^TK_j^T)F_{V_1,\mathrm{},V_N}^T(\lambda ,\mu )$$
(8.4)
where $`D_j^T`$ and $`K_j^T`$ are defined by (2.4) and (2.5).
###### Theorem 8.3.
Let $`T`$ be an automorphism of $`\mathrm{\Gamma }`$. The functions $`F_{V_1,\mathrm{},V_N}^T(\lambda ,\mu )`$ satisfy the following difference equation for each $`j=1\mathrm{},N`$ :
$$B_{V_j}B_{V_j^{}}^{}F_{V_1,\mathrm{},V_j,\mathrm{},V_N}^T(\lambda ,\mu )=(D_j^{,T}K_j^{,T})F_{V_1,\mathrm{},V_j^B,\mathrm{},V_N}^T(\lambda ,\mu ),$$
(8.5)
where $`D_j^{,T}`$ and $`K_j^{,T}`$ are defined by (2.9).
The above two theorems are proved in the same way as Theorems 2.2 and 2.4 respectively.
####
Similarly, let $`W`$ be an integrable highest weight $`U_q(𝔤)`$-module (resp. a $`B`$-invariant integrable highest weight $`U_q(𝔤)`$-module) and let $`V_1,\mathrm{},V_N`$ be $`U_q(𝔤)`$-modules from the category $`𝒪`$.
###### Theorem 8.4.
$$𝒟_W^TF_{V_1,\mathrm{},V_N}^T(\lambda ,\mu )=\chi _W(q^{2\mu })F_{V_1,\mathrm{},V_N}^T(\lambda ,\mu ),$$
(8.6)
where $`𝒟_W^T`$ is defined by (2.2).
###### Theorem 8.5.
Let $`T`$ be an automorphism of $`\mathrm{\Gamma }`$. Then
$$𝒟_W^{,T}F_{V_1,\mathrm{},V_N}^T(\lambda ,\mu )=\mathrm{Tr}_{|W^{𝔥_0}}(q^{2\lambda }B)F_{V_1,\mathrm{},V_N}^T(\lambda ,\mu ),$$
(8.7)
where $`𝒟_W^{,T}`$ is defined by (2.7).
The proof of the above two theorems is the same as in the finite-dimensional case.
#### Remark.
The integrability condition on the module $`W`$ is not essential.
####
The classical limits of Theorems 8.2-8.5 are analogous to the the corresponding classical limits of Theorems 2.1-2.4 in Section 7.
Acknowledgments
The first author was partially supported by the NSF grant DMS-9700477. The work of both authors was partly done when they were employed by the Clay Mathematics Institute as CMI Prize Fellows. O.S. would like to thank the MIT mathematics department for its hospitality.
Pavel Etingof, MIT Mathematics Dept., 77 Massachusetts Ave., CAMBRIDGE 02139 MA., USA
etingof@math.mit.edu
Olivier Schiffmann, MIT Mathematics Dept., 77 Massachusetts Ave., CAMBRIDGE 02139 MA., USA
schiffma@math.mit.edu.
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# The Evolution of X–ray Clusters and the Entropy of the Intra–Cluster Medium
## 1 INTRODUCTION
Clusters of galaxies are the largest virialized objects in the Universe, and are usually considered a canonical data set for testing cosmology. They are the largest collections of diffuse, highly ionized baryons that are directly observable in X–rays mostly through thermal bremsstrahlung emission. The strong dependence of X–ray emission on density $`L\rho ^2`$ allows one to select clusters and define complete samples much better than in the optical band.
X-ray observations of cluster number counts, luminosity functions and temperature distributions indicate little apparent evolution in clusters back to redshifts as high as $`0.7`$ (e.g., Henry 1997, 2000; Rosati et al. 1998; Schindler 1999), with the exception of very high luminosity objects or very high redshifts (Gioia et al. 1990; Rosati et al. 2000). This set of results provides one of the strongest challenges to high–density cosmological models in which cluster evolution is expected to be detectable even at redshifts as low as $`z0.3`$. However, these tests are highly dependent on the thermodynamic evolution of the ICM (e.g. see Borgani et al. 1999 and references therein; Bower 1997). The best–fit cosmological parameters are degenerate with the phenomenological parameters used to describe the evolutionary properties of the ICM. In fact, the diffuse baryons in clusters do not simply follow the dark matter, as would be the case if they were driven only by gravity as in self–similar models (Kaiser 1986). Significant efforts have been devoted recently to building a physical model for the ICM including an energy scale at which baryons and dark matter effectively decouple and the self–similarity is broken.
The presence of a minimum entropy in the pre-collapse IGM has been advocated for some time as a way to naturally break the self-similar behaviour (Kaiser 1991, Evrard & Henry 1991). Such an extra entropy is the key ingredient in reproducing the observed luminosity–temperature relation $`LT^n`$ with $`n3`$ (David et al. 1993, Mushotzky & Scharf 1997, Allen & Fabian 1998; Arnaud & Evrard 1999; Markevitch 1998), which is at variance with the self–similar prediction $`LT^2`$. Such an entropy minimum bends the relation from self–similar $`LT^2`$ behaviour at very large scales ($`10^{15}M_{}`$) towards a steeper slope on the scale of groups ($`10^{13}10^{14}M_{}`$) which is actually observed (Ponman et al. 1996; Helsdon & Ponman 2000). The average $`LT^3`$ relationship is essentially produced by the flattening of the density distribution in the cores of the X–ray halos; such cores grow larger as the mass scale decreases, and the luminosity steepens further on the scale of groups, where the gas is only adiabatically compressed (see Balogh, Babul & Patton 1999; Cavaliere, Menci & Tozzi 1997, hereafter CMT97; Cavaliere, Menci & Tozzi 1999).
The picture has been reinforced by the net change observed in the chemical properties and the spatial distribution of the ICM on the scale of groups, below the observed temperature of 1 keV (Renzini 1997, 1999) where the effects of the entropy excess are expected to be strongest. Another piece of evidence can be obtained from the observed mass–temperature relation (see Horner, Mushotzky & Scharf 1999). Recently, an excess of entropy (with respect to the self similar scaling) has been directly detected in the central regions of small clusters with temperatures between $`1`$ and $`3`$ keV (Ponman, Cannon & Navarro 1999, hereafter PCN; see also Lloyd–Davies, Ponman & Cannon 2000), pointing to the role of the entropy as the key ingredient determining the different properties of clusters and groups.
Independent hints come from the extragalactic X–ray background: without a substantial entropy injection at early epochs, its level and correlation function would exceed the observed limits, due to the widespread cooling phenomena that would radiate most of the gravitational energy of the collapsing baryons in the soft X–ray band (Pen 1999; Wu, Fabian & Nulsen 1999).
However, even if there are many hints pointing towards a comprehensive picture, there is a large uncertainty on the amount of extra–energy that effectively generates the entropy excess. It can be shown that it is the final entropy distribution that determines both the spatial distribution of the ICM and its evolutionary properties, irrespective of the total energy released in the past. A given entropy level can be reached through different thermodynamic histories, so that it is not possible to relate the ICM properties directly to a given energy excess without knowing the detailed physics of the heating processes. As we will show in this paper, the first question to answer is not: how much energy has been released in the ICM? but rather: what is the sequence of adiabats through which the baryons evolve?
It is difficult to predict a priori the entropy excess of the cosmic baryons, since most of the processes regulating nuclear activity, star and galaxy formation, and the transfer of energy to the surrounding baryons, are out of reach of present–day techniques. Thus, at present there is no general consensus on the production mechanism of such extra–entropy. For example, it is not clear whether the entropy minimum has been established in the IGM before it has been accreted –the external scenario, or in the high density ICM after accretion–the internal scenario. A different energy budget is required in the two different scenarios: a few tenths of a keV per particle are needed if the entropy is generated early enough to keep the baryons on a high adiabat, which prevents them from reaching high densities and cooling massively; much higher energy excess ($`>1`$ keV per particle) is required if the entropy is generated later, when the cooling process is eventually already widespread and most of the gas is already at high densities (Tozzi, Scharf & Norman 2000, hereafter TSN00).
The external scenario, which we will assume as a reference model, is provided by a ubiquitous entropy floor in the diffuse gas, which is entirely due to non–gravitational processes and is assumed to be in place before the onset of gravitational collapse of massive halos. The initial extra entropy is ineffective in large mass systems, where most of the entropy is due to strong shocks, but is more important in smaller mass systems, where the entropy production via shocks is strongly reduced. Eventually a large part of the baryons are merely adiabatically compressed and retain full memory of the initial entropy level. The non–gravitational origin of the excess entropy is crucial, since its level is independent of the mass scale and it breaks the self–similarity, while gravitational processes always scale self–similarly with mass.
We present a detailed model to relate the thermodynamic properties of the ICM in groups and clusters of galaxies to an initial entropy excess in the IGM, taking into account the transition between the adiabatic and the shock regime in the growth of X–ray emitting halos. The effect of radiative cooling is also included. We show that, despite the many complexities involved, the entropy is always a convenient synthetic quantity to describe the thermodynamic history of the cosmic baryons at least on the scale of groups and clusters. In particular, we show that in many circumstances the entropy track of a shell of baryons being accreted onto dark matter halos goes through three major regimes: (1) adiabatic compression, during which both heating and cooling are negligible and the entropy is constant; (2) step–like discontinuities due to gravitationally induced shocks; and (3) slow decrease when cooling becomes efficient for baryons in the inner regions of large halos. The entropy jump, the onset of cooling, and the final spatial distribution of the ICM, depend on the initial entropy. Such an external, initial entropy level can be reconstructed from the observation of a large number of distant clusters, or from the spatially and spectrally resolved profiles of nearby halos (see TSN00). Even if the knowledge of the entropy does not resolve the details of the underlying heating history and determine unambiguously the energy budget, the combination of data in the X-ray band with data in the optical and infrared bands can help to identify the major source of heating. In principle, this allows a detailed reconstruction of the energetic processes that affect the cosmic baryons over a wide range of scales and cosmic epochs.
The paper is organized as follows. In §2 we establish a one–to–one correspondence between the entropy level and the distribution of the ICM in halos in equilibrium. In §3 we present a generalized spherical infall model to follow the entropy track of each shell. In §4 we derive the average density and temperature profiles and the related global properties such as luminosity, emission weighted temperature and core radius, as a function of mass scale, cosmology, epoch and dark matter profile. In §5 we widen the parameter space, and investigate a time–dependent background entropy to show how the evolution in the entropy reflects in the X–ray properties of clusters of galaxies. In §6 we discuss the limitation of the present approach. Finally, our conclusions and future perspectives are presented in §7.
## 2 ICM THERMODYNAMICS: ENTROPY
The position, density and temperature of each shell in hydrostatic equilibrium in a given dark matter halo (whose average properties are determined by its total virialized mass $`M_0`$ at the epoch of observation $`z_0`$), can be unambiguously recovered once the final entropy profile is known. Assuming a spherical mass distribution, the equation of hydrostatic equilibrium for diffuse baryons in the potential well is:
$$\frac{1}{\rho }\frac{dp}{dx}=C\frac{m(<x)}{x^2},$$
(1)
where the radius $`x`$, the pressure $`p`$ and the density $`\rho `$ refer to the baryons and are normalized to the respective values at the last accreted shell at $`z=z_0`$, while $`m`$ is the total mass profile normalized to the total virialized mass. Explicitly, $`xR/R_s`$, $`pP/P_s`$, $`\rho =\rho _B/\rho _s`$, and $`m(<x)M(<x)/M_0`$. Since dark matter and baryons are distributed differently, we write $`M(<x)=M_{DM}(<x)+M_B(<x)`$. The constant is $`C=GM_0\mu m_p/R_sk_BT_s`$, where $`m_p`$ is the proton mass, $`G`$ is the gravitational constant, $`k_B`$ is the Boltzmann constant and $`\mu `$ is the molecular weight of the plasma (we will assume $`\mu 0.59`$ for a primordial IGM). $`T_s`$ is the temperature of the last accreted baryonic shell. In the following we will refer to the values of the last accreted shell as the shock value, even in the limit of a vanishingly small shock. We assume that hydrodynamic equilibrium is instantaneously established after each accretion event.
We define the adiabat $`Kk_BT/\mu m_p\rho _B^{\gamma 1}`$ (following the notation of Balogh, Babul & Patton 1999), where $`S\mathrm{ln}(K)`$ is the entropy and $`\gamma `$ is the microscopic adiabatic index which is $`\gamma =5/3`$ for a monoatomic gas. Using the perfect gas equation, we can write the density in terms of pressure and entropy normalized to the value at the last accreted shell, with each shell scaled to the corresponding adiabat: $`\rho =p^{1/\gamma }k^{1/\gamma }`$, where $`kK(x)/K_s`$. Substituting in Equation (1), the equilibrium pressure profile is re-written as:
$$\frac{dp}{dx}=Cp^{1/\gamma }k^{1/\gamma }\frac{m(<x)}{x^2}.$$
(2)
The above expression allows us to calculate the thermodynamic properties of a hydrostatic distribution of gas when the adiabat profile $`K(x)`$ is known. The main difference from the usual solutions of the hydrostatic equilibrium equation is that there is no need to assume a polytropic index, since each shell already sits on its adiabat which is determined by its previous history, and the correspondence between density and temperature is unambiguous.
The problem reduces to finding the proper adiabat of each infalling shell, or the entropy as a function of the accreted baryons, since the baryonic mass included in a given shell is constant with time. This procedure is convenient when applied to clusters of galaxies, because the entropy is conserved for the majority of the time. In fact, the dynamic history of a shell of gas can be described in three steps: 1) adiabatic compression during the infall; 2) shock heating at the accretion; 3) compression within the potential well due to further growth of the halo. The entropy is therefore constant during the first and third phase, and the jump at the shock is the most important feature needed to reconstruct the final profile. Cooling introduces further complexity, because for the inner, higher density shells, the radiative loss becomes important, changing substantially the final adiabats with respect to the initial value. However, as we will see later, the cooling can be included in the above picture, as long as the initial adiabat is not too low.
To begin with, we focus on the most important event in the entropy history of each shell: the accretion epoch. To calculate the value of $`K`$ immediately after the accretion shock, we need to estimate both the density and the temperature of each shell after shock heating eventually raised the adiabat from the external value to the post–accretion value $`K_i(x)`$. If a shock does not occur, the baryons are only adiabatically compressed and are accreted with the same adiabat. To determine whether a shell is shocked or not during accretion, we build a spherical infall model for the baryons, generalized for different cosmologies and epochs.
## 3 A GENERALIZED SPHERICAL MODEL
In the framework of the hierarchical clustering scenario, the baryons are accreted along with the dark matter during the process of gravitational collapse. An expanding accretion shock at the interface of the inner hydrostatic gas with a cooler, adiabatically–compressed, external medium, located approximately at the virial radius of the cluster, is a longstanding prediction from such gravitationally–driven models (see the 1D models of Bertschinger 1985, Ryu & Kang 1997, Knight & Ponman 1997, Takizawa & Mineshige 1998, and the 3D numerical simulations of Evrard 1990, Roettiger et al. 1993, Metzler & Evrard 1994, Bryan & Norman 1998, Abadi, Bower & Navarro 2000). Due to the growth of the total virialized mass, the baryons accreted later experience larger shocks, and the resulting entropy profile is always growing outwards. Such gravitationally–driven models predict X-ray properties which scale self–similarly with mass and fail to reproduce the X-ray observations of clusters.
A non–negligible value of the background entropy is needed in order to break the self–similarity. In fact, an initial adiabat will prevent shocks occurring below a given mass scale. We now discuss the external scenario in which an initial adiabat $`K_{}`$ is imprinted on all the diffuse IGM at some epoch prior to the formation of the dark matter potential wells. We refer to $`K_{}`$ as to the background entropy established in the IGM by non–gravitational processes before the baryons are accreted.
### 3.1 Accretion and Shock Conditions
The most prominent feature of the entropy history of each shell is the discontinuity at the accretion shock. To calculate the discontinuity we need to know the pre–shock density and the temperature that the infalling gas reaches moving along the initial adiabat $`K_{}`$ before accretion. Then we calculate the postshock temperature and density using mass, momentum and energy conservation, in the limit of complete thermalization of the kinetic energy of the gas.
The first important quantity is the infall velocity $`v_i`$. The dependence of $`v_i`$ on the total mass enclosed by the shell can be written as:
$$\frac{v_i^2}{2}=\frac{v_{ff}^2}{2}+\mathrm{\Delta }W\frac{c_s^2}{\gamma 1}+\frac{c_s^2}{\gamma 1}\left(\frac{\rho _{ta}}{\rho _e}\right)^{\gamma 1},$$
(3)
where $`\rho _{ta}`$ is the density at turnaround, $`\rho _e`$ is the gas external density, $`c_s=\sqrt{\gamma K_{}\rho _e^{\gamma 1}}`$ is the sound speed (hereafter $`\gamma =5/3`$), both calculated at the accretion radius $`R_s`$, and $`v_{ff}`$ is the free–fall velocity of a particle containing always the same amount of mass during the infall. Equation 3 is a generalized version of the Bernoulli equation for an adiabatic, spherically simmetrical accretion (Bondi 1952). The last quantity can be written as:
$$\frac{v_{ff}^2}{2}\frac{GM}{R_s}\frac{GM}{R_{ta}},$$
(4)
where $`M`$ is the total mass initially included by the baryonic shell. The term $`\mathrm{\Delta }W`$ is the contribution added to $`v_{ff}^2/2`$ to obtain the total work done by the gravitational potential on the baryonic shell, from the turnaround radius, $`R_{ta}`$, to the accretion radius $`R_s`$, including the effect of the time–varying enclosed mass. To evaluate this term it is strictly necessary to solve the trajectory of each baryonic shell. However we can make the simplifying assumption that the amount of dark matter enclosed by each shell, is a monotonically growing function of time, from the mass enclosed at turn around, to the final mass enclosed at the shock radius. The term $`\mathrm{\Delta }W`$ is estimated in §A, and the uncertainty on it turns out to be approximately $`10`$$`30`$ %. We show later that this error is not important in determining the transition scale between the shock and the adiabatic regime.
The other two terms proportional to $`c_s^2`$ describe the energy needed to compress the gas. In fact, due to the non–negligible value of $`K_{}`$ in the infalling IGM, part of the gravitational energy goes into internal energy in an amount proportional to the square of the sound speed in the external IGM at the epoch of accretion, so that in general $`v_i<v_{ff}`$. The compression term carries an increasing fraction of the potential energy when the mass of the system is lower, or, since the sound speed is proportional to $`K_{}^{1/2}`$, when the entropy is higher. The fourth term on the right hand side of Equation (3) results from the initial condition $`v_i=0`$ for a gas shell at the turnaround radius, when the gas had a density $`\rho _{ta}`$ and it is assumed to be at the same contrast of the dark matter. The epoch of turnaround is assumed to be half of the infall epoch.
Of course to solve Equation (3) we need to evaluate $`\rho _e`$. To do this, we first note that the knowledge of both the external density and the infall velocity gives the net infall accretion rate of baryonic matter through the surface defined by the shock radius. Then, we make the assumption that the growth rate of the total virialized mass $`\dot{M}`$ is proportional to the growth rate of the thermalized baryonic mass $`\dot{M}_B`$. Here $`\dot{M}`$ is the average total mass accretion rate as predicted in the hierarchical clustering scenario. This means that all the baryons, that initially were in the same lagrangian volume of the mass that is currently virialized, have been accreted. The proportionality constant is simply the average mass fraction of baryons in diffuse form $`f_B`$, so that at each epoch the fraction of accreted baryons (with respect to the total baryons accreted at $`z=z_0`$), is equal to the fraction of the accreted matter to the total virialized mass at the same final epoch. This does not imply that the baryons are in the same volume; they are distributed in a volume typically larger than that of the accreted dark matter. This occurs especially in the adiabatic regime, when the baryons have too high a temperature to sink into the potential well and thus the accretion radius is significantly larger than the virial one. The constraint on the mass accretion rate translates into the relation:
$$\dot{M}_B=f_B\dot{M}=\rho _e4\pi R_S^2\left(v_i+\frac{dR_S}{dt}\right),$$
(5)
where $`\dot{M}`$ is given for a particular cosmological model (see §3.3). We can derive $`\rho _e`$ as a function of $`v_i`$, and then the external temperature is $`k_BT_e=\mu m_pK_{}\rho _e^{2/3}`$.
The condition $`v_i>c_s`$ determines if the shell is shocked. In the frame of the infalling gas the shock expands with a velocity $`v_i+dR_S/dt`$. In the case of a shock, we assume that all the kinetic energy of the infalling gas is thermalized (i.e., the post-shock velocity $`v_{ps}=0`$ in the rest–frame of the cluster), and obtain for the postshock temperature (Landau & Lifshitz 1957; Cavaliere, Menci & Tozzi 1998):
$$k_BT_i=\frac{\mu m_pv_i^2}{3}\left[\frac{(1+\sqrt{1+ϵ})^2}{4}+\frac{7}{10}ϵ\frac{3}{20}\frac{ϵ^2}{(1+\sqrt{1+ϵ})^2}\right],$$
(6)
where $`ϵ15k_BT_e/4\mu m_pv_i^2`$.
The postshock density is then $`\rho _i=g\rho _e`$, where $`g`$ is the shock compression factor which depends on the postshock temperature, $`T_i`$, and the external temperature, $`T_e`$, and is given by (see CMT97):
$$g=2\left(1\frac{T_e}{T_i}\right)+\left[4\left(1\frac{T_e}{T_i}\right)^2+\frac{T_e}{T_i}\right]^{1/2}.$$
(7)
If the gas is shocked, we calculate the new adiabat $`K_i=k_BT_i/\mu m_p\rho _i^{2/3}`$ of the baryonic shell after accretion. If the infalling velocity is smaller than the sound speed in the external IGM and the shock does not occur, the gas is accreted adiabatically, and therefore the post–accretion adiabat is the inital one $`K_i=K_{}`$, which is all we need to solve for the final equilibrium.
Thus, using Equations (3), (6) and (7), we are able to associate with each shell, including a mass $`M_B`$ of baryons, its postshock adiabat $`K_i(M_B)`$. For a given object, the adiabat of the infalling shells initially will be $`K_i=K_{}`$, since for sufficiently low velocities the shocks are suppressed. As the total mass grows, the velocities of the infalling shells rise approximately as $`v_iM^{1/3}`$, more rapidly than the sound speed (which in general decreases with epoch, since $`c_s\rho _e^{1/3}1+z`$), and eventually a shock regime begins. In Figure 1 the transition between the two regimes is shown as a function of the accreted mass for a given initial adiabat $`K_{}`$. As it is shown in the first panel, the maximum uncertainty in the infalling velocity, $`v_i`$, grows toward the adiabatic regime, but it does not introduce a large error in the transition scale, since the infall velocity falls steeply below $`c_s`$. The rapid increase of both the infall and the free–fall velocity at the transition, occurs because the gravitational energy becomes sufficient to overcome the pressure support, and the accretion radius moves from a relatively distant position to a position very close to the virial radius. Clearly, the presence of a larger $`K_{}`$ further delays the onset of the shock heating regime, inhibiting adiabatic accretion for the majority of the baryons, especially in small mass systems.
At this stage, if we neglect further changes in the entropy, the adiabat in the final position is simply $`K(x)=K_i`$ and the Equation (2) can be solved easily without any further steps. However, for the inner shells, radiative cooling becomes important and the calculation of the final adiabat requires solving Equation (2) at different epochs, as explained in the following subsection.
### 3.2 The Effect of Radiative Cooling
Each shell of gas is continuously changing its adiabat due to cooling and heating processes. In particular, the first baryonic shells that are accreted drain into the inner, higher density regions of halos as the total virialized mass grows, and their cooling times become small enough to start cooling processes. As a result, the final adiabat of these baryonic shells will be lower than that at the accretion epoch, and eventually part of the gas leaves the diffuse, emitting phase and sinks into the center.
We can model the cooling assuming a homogeneous, single temperature distribution (Fabian & Nulsen 1977, Mathews & Bregman 1978); in this case the energy equation can be formally written as:
$$\frac{d}{dt}[\mathrm{ln}(K)]=\frac{1}{\tau _{cool}(K)},$$
(8)
where the cooling time $`\tau _{cool}`$ is defined as:
$$\tau _{cool}\frac{3}{2}\frac{k_BT}{\mathrm{\Lambda }_{net}}\frac{\rho _B}{\mu m_p},$$
(9)
and therefore it depends on $`K`$ through $`T`$, $`\rho _B`$ and $`\mathrm{\Lambda }_{net}`$. Here $`\mathrm{\Lambda }_{net}`$ is the cooling rate including free-free and line emission (see Sutherland & Dopita 1993).
It is well known that cooling is a runaway process, and the solution of Equation (8) would require the computation of the equilibrium profile at many different epochs. Since we still want to have the benefit of a relatively fast computation, much faster than a full hydrodynamic simulation, we tackle the problem choosing a medium resolution in time ($`\mathrm{\Delta }t0.3`$ Gyr) and solving Equation (8) within $`\mathrm{\Delta }t`$ for every shell with an analytic approximation. This is possible if we assume that the cooling process is isobaric within $`\mathrm{\Delta }t`$, in order to express both density and temperature as a function of the adiabat $`K`$ only. The pressure is updated at each time step, following the new equilibrium configuration. An intermediate step is to approximate the cooling function, $`\mathrm{\Lambda }_{net}`$, with an analytic function of the temperature. In this way the change in the adiabat within $`\mathrm{\Delta }t`$ can be derived as the integral of an analytic function, as described in Appendix §B.
When the cooling times become very short in the center of the halo, part of the gas may eventually cool in a single time step $`\mathrm{\Delta }t`$ (i.e., its entropy drops to zero). In this case, the gas is removed from the diffuse, emitting phase, and is included in a gravitational term as if it is all accumulated in the very center. At this level, we do not implement more sophisticated multiphase models which can be important for the detailed emissivity distribution in cooling flows. However, we can follow the steepening of the baryonic density in the center as the radiative cooling becomes efficient, and compute the corresponding amount of baryons which drop out from the diffuse phase. We stress the fact that we are able to follow the complex cooling processes with good accuracy by virtue of the initial entropy level. The background entropy, in fact, delays and possibly inhibits the onset of strong cooling flows. Our model breaks down in the limit of small initial entropy, where the cooling catastrophe occurs.
The evolution of the adiabat as a function of cosmic epoch for some given shells is plotted in Figure 2. The outermost shells are accreted at later epochs. They are strongly shocked and reach a high adiabat, and find equilibrium at large radii and low densities. Consequently, the cooling times are always large and the adiabat $`K`$ stays almost constant after the accretion. Conversely, inner shells are more affected by cooling for two reasons: they reach much higher densities (being in the central regions), and they have more time to cool since they are accreted much earlier. Eventually, the very inner shells reach very low entropy, corresponding to extremely high densities and very short cooling times, and they rapidly cool and drop out of the diffuse phase.
The calculation without the inclusion of cooling would be much simpler, since the final adiabat would be the accretion value $`K_i`$ for all the shells, and the hydrostatic equilibrium would be solved only once (at the final epoch $`z_0`$). However, solving the equilibrium at several epochs allows us to follow the evolution of the X–ray properties for each (average) dark matter halo. In Figure 3 the evolution of temperature and luminosity for three objects of $`10^{15}h^1`$ (continuous line), $`10^{14}h^1`$ (dashed line), and $`10^{13}h^1M_{}`$ (dotted line), is shown for a constant $`K_{}=0.3\times 10^{34}`$ erg cm<sup>2</sup> g<sup>-5/3</sup> in a $`\mathrm{\Lambda }`$CDM cosmology. In the third panel, the time evolution of the shock radius is plotted for the same objects. The shock radius is normalized to the virial radius at each epoch. It is possible to see how the shock radius is close to the virial one for the largest halo and relaxes in the last few Gyr when the mass accretion slows down and the external pressure term correspondingly decreases. The effect is more pronounced at lower masses, where the internal pressure support is strong enough to dominate the gravitational potential and the external pressure term of the infalling gas.
In Figure 4 we plotted, for the same three final masses, some relevant quantities averaged over the adiabatic cores, defined as regions including the gas accreted during the adiabatic regime. It is possible to see how the initial entropy $`K_{}=0.3\times 10^{34}`$ erg cm<sup>2</sup> g<sup>-5/3</sup> introduces a large difference in the central core as a function of mass. Central densities are much higher for deeper potential wells. In addition, the baryons in the center of massive clusters suffer radiative losses and the baryonic cores shrink to smaller sizes and higher densities. In the second panel the average entropy of the inner cores is shown. The decrease due to the radiative cooling can reduce the initial entropy especially in the most massive halo. The decrease in the entropy is driven by the decrease in the average cooling time shown in the third panel. While the entropy is decreasing, the internal energy of the gas is still rising due to the compressional work done by the gravitational potential. However, the trend of stronger cooling for larger masses, is reversed in the case of very small $`K_{}`$. In fact, as long as there is nothing to prevent the baryons from cooling, the amount of radiative losses is mainly set by the age of the halos.
### 3.3 Dark Matter Properties
In this section we briefly review the properties of the dark matter halos which drive gravitationally the evolution of the diffuse baryons. In particular we describe the mass profiles and the mass accretion rates in the framework of the hierarchical clustering scenario in universes dominated by cold dark matter (CDM). However, the model can be generalized to other cosmologies.
The boundary of a halo is the virial radius, defined as the radius within which the average overdensity with respect to the critical density is $`\mathrm{\Delta }_c`$, where $`\mathrm{\Delta }_c=178`$ for $`\mathrm{\Omega }_0=1`$ with a mild dependence on $`\mathrm{\Omega }_0`$ (see, e.g., Eke et al. 1998). Analytical studies indicate simple power law profiles for the dark matter, of the kind $`\rho x^\xi `$, with $`\xi =9/4`$ (Gunn & Gott 1972, Bertschinger 1985). Numerical works show a more complex behaviour, with a characteristic internal scale radius that depends on the epoch and on the final mass (Navarro, Frenk and White 1997, hereafter NFW; Moore et al. 1998). A very general expression for the universal profile is:
$$\rho =\rho _{c0}\frac{\delta _c}{(cx/x_v)^\nu [1+(cx/x_v)^\zeta ]^\eta },$$
(10)
where $`c`$ is the mass dependent concentration parameter of the dark matter, and $`\delta _c`$ is defined by requiring the average density within $`R_V`$ with respect to the critical density to be $`\mathrm{\Delta }_c`$. Here we used $`x_vR_v/R_s`$ to be consistent with Equation (2) where the radius is normalized to $`R_s`$. Present calculations differ mainly in the inner regions, where NFW predict $`\nu =1`$, $`\zeta =1`$ and $`\eta =2`$, while Moore et al. (1998) have a steeper inner profile with $`\nu =1.5`$, $`\zeta =1.5`$ and $`\eta =1`$. From Equation (10) the mass profiles $`m(<r)`$ entering Equation (2) follow directly.
We will approximate the concentration parameter with power laws, which turn out to be good approximations (Navarro et al. 1997). The expressions used are described in appendix §C. In general, the concentration parameter $`c`$ depends on the characteristic epoch of formation of the halo, which in turn depends on cosmology, perturbation spectrum, $`M_0`$ and $`z_0`$ (see NFW). This is because the dark matter remembers the epoch when each shell was accreted, even if the shell–crossing tends to erase such dynamical memory. For example, in a standard CDM universe groups tend to have a larger concentration ($`c8`$) being formed at higher epochs when the average density was higher, while clusters, being younger, have a lower concentration ($`c6`$). At higher redshifts the concentration parameters are generally lower, since the difference in epoch (and thus in typical density) between formation and the observation epochs $`z_0`$ is reduced. These trends will be included in our calculations.
The accretion processes in groups and clusters show considerable scatter, as observed in numerical simulations and Monte Carlo realizations of hierarchical clustering based on the extended Press & Schechter formula (hereafter PS, Press & Schechter 1974, Bond et al. 1991, Bower 1991, Lacey & Cole 1993). However we are interested in the mass history of typical halos, each of them labeled by the final mass $`M_0`$ and the final (observation) epoch $`z_0`$, for a given cosmology. The natural way to proceed is to average over many realizations of the mass history of the main progenitor, defined as the most massive halo participating in every mass accretion event along the merger tree of a single object. We run $`1000`$ Monte Carlo simulations of the mass history of the main progenitor for different final masses $`M_0`$ and different final redshifts $`z_0`$ ($`z_0=0`$ and $`z_0=1`$) in the two cosmologies discussed below (tCDM and $`\mathrm{\Lambda }`$CDM). We find that the average mass growth of the main progenitor can be approximated within few percent by a parabola in the $`\mathrm{log}(m)`$$`\mathrm{log}(1+z)`$ space:
$$m(z)=\left(\frac{1+z}{1+z_0}\right)^{[B+A\mathrm{log}(\frac{1+z}{1+z_0})]},$$
(11)
where $`A`$ and $`B`$ depend on cosmology, $`M_0`$ and $`z_0`$. The relation (11) is used to determine the accretion epoch of each baryonic shell after Equation (5), and thus to compute its density, $`\rho _e`$, at the accretion shock.
The dispersion in the profiles and in the accretion process is likely to introduce some dispersion in the resulting X–ray properties, and is expected to explain partially the intrinsic scatter observed in the $`L`$$`T`$ relation. The intrinsic scatter in the emission is certainly due also to the presence of cooling flows (Allen & Fabian 1998; Arnaud & Evrard 1999), which in turn can be affected by both the dynamical and the heating history of the gas. For these reasons, we focus on typical halos averaging over many different realizations, and considering the accretion of baryons as a smooth and continuous process. These assumptions clearly break down in the case of massive merger events (see discussion in §6).
## 4 RESULTS
Here we present the X–ray properties of groups and clusters of galaxies in the case of a constant and homogeneous $`K_{}`$ in the external IGM. Our reference calculation will be a flat, low density cold dark matter universe ($`\mathrm{\Lambda }`$CDM), which is currently preferred on the basis of the measurements of the expansion rate of the universe from high $`z`$ SNe (Riess et al. 1998), from the Cosmic Microwave Background (see Lange et al. 2000; Balbi et al. 2000) and of the observation of a high baryonic fraction $`f_{obs}>0.06h^{1.5}`$ in clusters (see Ettori & Fabian 1999), which is consistent with standard nucleosynthesis if $`\mathrm{\Omega }_0<0.3`$ (White et al. 1993). The baryonic density is assumed to be $`\mathrm{\Omega }_B=0.02h^2`$ (Burles & Tytler 1999a, 1999b), consistent with the standard primordial nucleosynthesis scenario. From the diffuse, X–ray emitting component, we exclude a fraction which is assumed to be locked in stars since the beginning, and is choosen to be $`20`$ % of the total baryons in halos (independent of the mass scales and epoch, i.e., we assume a constant efficiency of star formation). The fraction of baryons cooled in the center, instead, is computed at each epoch and subtracted from the diffuse, X–ray emitting phase. For comparison, we will also discuss a tilted cold dark matter universe (tCDM), where we are forced to adopt a baryonic density $`\mathrm{\Omega }_B=0.04h^2`$, larger than the standard value, in order to be consistent with the observed baryonic fraction. The details for the two universes are shown in Table 1. The values for $`A`$ and $`B`$ in the two universes are determined with a $`\chi ^2`$ fitting of the average mass histories with the relation (11), and are reported in Table 2.
### 4.1 Density and Temperature Profiles
First, we discuss a simple case where the cooling is not included, so that the final adiabat $`K(x)`$ is equal to the value at the accretion, $`K_i`$. This case shows the effects of the entropy excess alone without the intervention of cooling processes. In Figure 5 we show the resulting profiles for $`\mathrm{\Lambda }`$CDM at redshift $`z_0=0`$, for an initial $`K_{34}=0.3`$, where $`K_{34}`$ is $`K_{}`$ in units of $`10^{34}`$ erg cm<sup>2</sup> g<sup>-5/3</sup>. This value corresponds to a temperature $`k_BT_{}1.5\times \mathrm{\hspace{0.17em}10}^2(1+z)^2`$ keV at the ambient density. The dark matter is distributed according to the NFW profile. Three final masses are shown: $`M_0=10^{15}10^{14}10^{13}h^1M_{}`$. The plotted profiles are all normalized at the corresponding shock values in order to show how the scaling behaviour departs from self–similarity. Note, however, that the density and temperatures values at the shock in physical units are very different in the three cases.
A characteristic feature is the flat density profile of isentropic gas in the core, which is relatively larger at smaller masses (dashed lines in Figure 5, panel a). Such cores are built in the initial, high redshift stages of the accretion process, when the accretion is adiabatic since the infall velocities are small and shocks do not occur. This regime is relatively more extended going to lower masses. The pressure is more effective in pushing the baryons over a region larger than that of the dark matter (panel c). All this information is synthesized in the entropy profiles: at larger radii the entropy rises since the outer shells experience stronger shocks (panel d). Since the entropy is normalized to the value at the shock radius, the constant entropy floor in the center appears different at different masses. The slope of the entropy profile in the shock dominated regime is almost independent of the initial value $`K_{}`$, yielding $`d\mathrm{ln}(K)/d\mathrm{ln}x1.1`$; this value is close to the value $`1.3`$ expected for the simple case of an isothermal profile where the entropy is due only to shock heating and $`M_Br`$. The sharp knee in the entropy profile is due to the fact that the transition from adiabatic accretion to strong shocks is very fast, and the intermediate shock regime virtually does not exist, so that during the shock regime the entropy is always dominated by shock heating. In contrast, in the center, isentropic cores are clearly emerging. The ratio of mass accreted adiabatically to the total baryonic mass, is correspondingly larger at lower mass scales (panel e).
Here we note that a departure from a power–law behaviour for the entropy profile has been observed in hydrodynamical simulations where neither radiative cooling nor extra-entropy were included (see Frenk et al. 1999). This may suggest that departures from a power–law behaviour in the entropy profile can also be originated by asphericity.
The temperature profiles (panel b) do show mild gradients in the regions where the gas has been shocked (variation less than a factor of $`3`$ between $`R_s`$ and $`0.1R_s`$), while they show considerable gradients when the entropy is constant, following $`TK_{}\rho ^{2/3}`$. Part of the large gradient in the smallest system corresponds to very low luminosity regions, where the gas is relaxed due to the very small pressure term. These regions, and the corresponding large temperature gradients, have never been observed. In fact, if we compute the temperature gradient in the inner regions of halos with $`M_0=10^{13}h^1M_{}`$, we find an increase of about 2 within a radius of $`0.1R_S100h^1`$ kpc, an effect hardly visible, e.g., in the data by ROSAT.
A good quantity to characterize the properties of the temperature profile is the effective polytropic index defined by the relation $`p\rho ^{\gamma _p}`$. In general, a family of polytropic relations can be used to describe the ICM and investigate the energy budget underlying each polytropic family (Loewenstein 2000). As a result of the combined action of shock heating and adiabatic compression, the index $`\gamma _p`$ is found to be approximately $`\gamma _p0.81.2`$ between the adiabatic core and the shock radius, roughly consistent with an isothermal temperature profile (in the Figure we show the value of $`\gamma _p`$ averaged over $`\mathrm{\Delta }\mathrm{log}(x)=0.3`$). In the adiabatic cores the polytropic index is simply $`\gamma _p=\gamma =5/3`$, since all the gas is on the same adiabat.
To elucidate how the breaking of self–similarity occurs, in Figure 6 we show the same profiles for a negligible value of the external entropy, but without the inclusion of cooling. This is what we call the self–similar case, which is different from the more realistic case of negligible entropy and the inclusion of cooling, since cooling also alters the entropy profile, as shown in §4.4. In the absence of an entropy floor, the profile $`K(x)`$ always decreases at smaller radii, and exhibits a power law behaviour without any particular scale. The only differences between groups and clusters are now driven by the dark matter distributions. Despite the pressure support, the gas essentially follows the dark matter, and groups appear more concentrated than clusters, reversing the trend of Figure 5.
In the $`K_{}0`$ case the cooling starts very early and deeply affects the profiles of massive clusters. The majority of the initially diffuse baryons cool in the center of small halos, where, without an effective background entropy, nothing prevents the baryons from cooling and the luminosity is dominated by the central regions. The cooling selectively removes the lower entropy gas in the center of lower mass objects, helping to create an entropy plateau at the very center, but with the entropy entirely produced by gravitational processes. This mechanism to create an entropy plateau has been advocated by PCN but a large amount of cooled baryons need to be accomodated in the center. The strongest evidence for the presence of a background entropy at high $`z`$, is given by the low fraction of baryons in stars with respect to the total baryons available, which implies a strong suppression of the cooling processes especially in low mass halos (see Prunet & Blanchard 1999).
Figure 7 shows the case with $`K_{34}=0.3`$ and with the inclusion of cooling. This case can be considered a realistic, complete scenario. As we shall see later, this value of the background entropy gives a good fit to the $`L`$$`T`$ relation. The inclusion of cooling introduced some change with respect to Figure 5, especially in the very inner regions, where the entropy evolved towards lower values. However the entropy excess in the center is still present (panel d). Cores with constant density appear more peaked, but small groups still show much flatter density profiles with respect to large clusters. The temperature profiles are lower, and the polytropic index $`\gamma _p`$ is rapidly decreasing in the center.
For a more comprehensive view, the differences in the density profiles can be expressed in terms of fitting parameters $`\beta `$ and $`r_c`$ after adopting a beta model (Cavaliere & Fusco Femiano 1976). The results are shown in Figure 8. The $`\beta _{fit}`$ parameter is about $`0.8`$ in $`\mathrm{\Lambda }`$CDM at $`z=0`$, and about $`0.6`$ at $`z=1`$. The density profiles are slightly steeper in the outer regions at smaller masses. However the most prominent feature is the core radius, whose scaling departs from the self–similar behaviour $`RM^{1/3}`$ (dotted line) below $`1`$ keV. No significant differences are predicted in the tCDM universe. The flattening of the $`RM`$ or the $`RT`$ relation has been clearly detected in the data, and related to heating processes, by Mohr & Evrard (1997); note, however, that they plotted an isophotal radius, which is a much better defined quantity from the observational point of view. Smaller cores are found at higher $`z`$, since all the linear dimensions are reduced approximately by a factor $`(1+z)`$.
We note that our results differ from those found by Fujita & Takahara (2000). In fact, their assumption of isothermality allows to relate the $`\beta `$ parameter directly to the temperature of the external gas. This is no longer valid in our model, where the entropy of the external gas affects the dimension of the adiabatic core $`r_c`$, breaking the self–similar scaling, while yielding a $`\beta `$ parameter independent of the mass.
Finally, we note that our values of $`\beta `$ are somewhat larger than that observed in clusters (see Mohr, Mathiesen & Evrard 1999). This may be due to our fitting procedure, that extends up to the shock radius. In fact, our profiles are steeper than a $`\beta `$–model at large radii, and the best fits usually give larger $`\beta `$ for larger cores in order to reproduce the rapid steepening of the profiles out of the core. The outer regions are generally too weak to be detected in ROSAT data, since their surface brightness is below $`10^{15}`$ erg s<sup>-1</sup> cm<sup>-2</sup> arcmin<sup>-2</sup> (see TSN00). On the other hand, such a regions are expected to be efficiently detected in the future Chandra and XMM data.
### 4.2 The effect of Cosmology and Dark Matter
From the considerations above, it is clear that the level of the initial adiabat strongly affects the final properties of the ICM, and that, in principle, it is not necessary to invoke substantial heating after the collapse, provided that $`K_{34}0.3`$. The profiles are affected also by changing the cosmological background, the epoch of observations, or the dark matter profile. To show these variations, not directly related to the entropy, in Figure 9 we plot the density and temperature profiles, along with the polytropic index, for a typical massive cluster ($`0.6\times 10^{15}h^1M_{}`$, corresponding to a virial temperature of $`k_BT5`$ keV) changing in turn cosmology, epoch and dark matter profile and comparing them to the case with $`\mathrm{\Lambda }`$CDM at $`z=0`$, NFW profile, $`K_{34}=0.4`$. The cooling is included for all the cases.
A steeper dark matter profile (Moore et al. 1998) gives higher gas densities in the center (dashed line). The temperature gradient is correspondingly larger. In any case $`\gamma _p`$ is always bounded between $`0.9`$ and $`1.2`$ outside the adiabatic core. In principle, observations can discriminate between different dark matter profiles, and the observed temperature profiles (see Markevitch 1998) would favour profiles steeper than NFW, but we recall that this minor effect can be overwhelmed by changes in the entropy or by the presence of substructure.
At higher redshifts (here we focus on a typical value of $`z=1`$ which is the nominal goal of the future X–ray surveys) the adiabatic accretion is relatively more extended in time during the lifetime of the object, and, for the same value of $`K_{}`$, the imprint of the background entropy is more evident. This is because virialized objects form at a total density contrast which is almost constant with respect to the critical density, and the baryons will consequently reach larger densities before being accreted. These larger densities translate into pre–shock temperatures larger approximately by a factor $`(1+z)^2`$, and thus in a larger sound speed $`c_s(1+z)`$. On average, the shock condition is harder to satisfy since the infalling velocities scale only as $`v_i\sqrt{1+z}`$, and, consequently, a larger number of baryons are accreted adiabatically. The resulting density and temperature profiles are flatter (dot–short dashed line). This effect adds to the flattening of the total dark matter profile at high redshift, as envisaged by NFW. As we will see, this mechanism is responsible for keeping the $`L`$$`T`$ relation approximately constant with redshift.
The case for a tCDM cosmology at $`z=0`$ (dotted line) shows flatter profiles. This is easily understood if we recall that the external density is proportional to the mass accretion rate, and that the mass accretion rates are higher at $`z=0`$ in tCDM with respect to $`\mathrm{\Lambda }`$CDM (similar to the rates at $`z=1`$ in $`\mathrm{\Lambda }`$CDM for objects of the same mass). In general, the cosmology does not have a large effect on the evolution of the $`L`$$`T`$.
### 4.3 The Shock Radius and the Baryonic Fraction
The boundaries of the emitting gas are given by the shock radius of the last accreted shell, where there is a discontinuity between the inner hot gas and the outer cooler gas. The outer unshocked gas gives also a contribution to the emission, and can be detected in the outskirts of rich clusters giving important information on the entropy level of the external baryons (TSN00). It always gives a small contribution if compared to the total emission from the cluster, and here it is neglected. For very small mass objects, the last accretion radius is quite distant from the virial radius, in a region of very low density and very low infall velocity. The shocks are typically very weak, and the gravitational entropy production is negligible. In such low mass objects the X–ray emission is expected to fade outwards without discontinuity.
The position of the last accreting shell is calculated simply using mass conservation. In fact, following Equation (5), the total mass of diffuse baryons involved in the cluster collapse, is equal to the mass included in the initial comoving region, $`M_B=f_BM`$, after subtraction of the baryons in stars and the cooled baryons in the center which depend on epoch and mass. Due to the different distribution of the baryons with respect to the dark matter, the ratio of the shock to the virial radius is a function of epoch and of the total mass accreted, as shown in Figure 3. In Figure 10(a) we show the position of the final shock radius with respect to the virial one at redshift $`z_0=0`$. At small masses, where the gas distribution is flatter and more extended, the shock radius can be approximately $`2`$ times larger than the virial radius. In other words, the external gas does not fall into the potential well, but is accumulated at large radii. At very high masses the accretion rates are larger, and the pressure term can be important, giving, for high density universes, a shock radius slightly smaller than the virial radius. The same happens at higher redshift when the accretion rates are correspondingly larger. In any case, for large mass systems the shock radius is expected to remain close to the virial radius of the cluster, as was predicted in numerical simulations (see, e.g., Takizawa and Mineshige 1998). Slightly larger shock radii are predicted for higher values of the background entropy.
Here, the ratio of the mass in baryons within the shock radius to the total mass within the virial radius, is, by definition, always equal to the universal average baryonic fraction. However, since the two radii are generally different, the observed baryonic fraction within the virial radius will be a growing fraction of the mass scales. In Figure 10(b) the baryonic fraction within the virial radius $`R_v`$ is shown as a function of the total virialized mass. The largest variations are between masses $`10^{13}`$ and $`10^{14}h^1M_{}`$, roughly corresponding to temperatures below $`1`$ keV at $`z=0`$. In any case, any entropy excess, irrespective of the origin (external or internal) always tends to puff up the baryons with respect to the dark matter (as observed in numerical simulations, see Pearce et al. 1994, Tittley & Couchman 2000). This reinforces the case for a low density universe derived from the observed high ratio of baryonic to total mass.
### 4.4 The Energy Budget
An important quantity is the amount of non–gravitational energy per particle corresponding to $`K_{}`$. The temperature corresponding to a given adiabat is $`k_BT3.2\times 10^2K_{34}(1+z_h)^2\delta ^{2/3}`$ keV, where $`\delta `$ is the overdensity with respect to the ambient density at $`z_h`$, and $`z_h`$ is the epoch of the heating. The assumption of an initial and homogeneous $`K_{}`$, implies that the entropy of each shell must be in place at turn–around. At this epoch the density of the shell is assumed to be the background value . Thus, the minimum energy released in the gas can be computed as
$$k_BT_{min}=3.2\times 10^2K_{34}\frac{1}{M_0}_0^{M_0}[1+z_{ta}(M)]^2𝑑M.$$
(12)
In the case of a $`\mathrm{\Lambda }`$CDM universe we have $`k_BT_{min}0.1(K_{34}/0.4)`$ keV with a small dependence on the final mass $`M_0`$.
As we have already discussed, starting from a high adiabat is not the only way to prevent massive cooling, since non–gravitational heating in the center of the clusters could help in re-establishing the entropy floor. However, the energy needed to re–establish the entropy floor after accretion is much higher than the energy needed to put the baryons initially on the right adiabat. If the baryons are heated preferentially at higher density the excess energy is higher by a factor $`\delta ^{2/3}`$. However, this is not the only reason for a larger energetic budget. In fact, another advantage in heating the gas at lower densities, is that radiative cooling is not able to re–emit the energy on very short time–scales.
To make a simple example without the cooling, if the baryons are heated at $`z0`$ when they are at an average density contrast equal to $`200`$, typical of virialization, we would obtain $`k_BT0.3`$ keV. However, this value underestimates the real energy budget, since the density in the center, where the entropy excess is expected, is much higher than the average contrast, and $`z0`$ is in any case too late to inject the extra energy. A more realistic calculation for the center of rich clusters can require more than $`2`$ keV per particle (see §7 and TSN00) to establish a density core and eventually halt the cooling in the center. This arguments show clearly how the same entropy level, which determines all the X–ray emission properties, can be due to very different heating balances. In this respect, the distribution of metallicity in the ICM may be useful in calculating the actual amount of excess energy dumped into the baryons.
### 4.5 The Luminosity–Temperature–Mass relations and the Entropy–Temperature plot
We can derive the average relation between the bolometric luminosity, the emission weighted temperature and the total virialized mass. The bolometric luminosity over the whole emitting volume defined by $`R_s`$ is:
$$L_x=_0^{R_s}ϵ(r)𝑑V\mathrm{erg}\mathrm{s}^1,$$
(13)
where $`ϵ(r)`$ is the emissivity per unit volume, including free-free and line emission, expressed by:
$$ϵ=n_en_i\mathrm{\Lambda }_N\mathrm{erg}\mathrm{s}^1\mathrm{cm}^3$$
(14)
where $`n_e`$ and $`n_i`$ are the electron and ion density respectively, and $`\mathrm{\Lambda }_N`$ is the normalized cooling function depending on temperature and metallicity (from Sutherland & Dopita 1993). We adopt a value of $`Z=0.3Z_{}`$, as observed on the scale of clusters ($`k_BT>2`$ keV). Such a value is currently observed on the scale of groups with large uncertainties, due to difficult line diagnostic and poor temperature resolution (Renzini 1997; Buote 2000). However, since the cooling function includes emission over a range of energies wider than the usual X–ray bands, we cut the emission at energies lower than $`0.1`$ keV.
The emission weighted temperature defined over the entire emitting volume is:
$$k_BT_{ew}\frac{_0^{R_s}k_BT(r)ϵ(r)𝑑V}{_0^{R_s}ϵ(r)𝑑V}\mathrm{keV}.$$
(15)
The results are shown in the Figures 11 and 12 for $`\mathrm{\Lambda }`$CDM and tCDM respectively for $`K_{34}=0.3`$, with the inclusion of cooling. The self–similar case is shown for comparison (dashed line). Data are taken from Arnaud & Evrard (1999) and Allen & Fabian (1998) for the clusters, and from Ponman et al. (1996) for the groups. An important issue here is that the total luminosity emitted by all the accreted gas (light curves in Figures 11 and 12), overestimates the luminosities found by Ponman et al. (1996) at temperatures below $`1`$ keV. This is because the luminosities of the observed groups are defined within the fixed projected radius of $`100h^1`$ kpc. Therefore we also calculated the luminosity and the emission weighted temperature performing the integrals of Equations (13) and (15) over the cylindrical volume defined by the projected radius of $`100h^1`$ kpc. We show both the total luminsity, including all the gas even at $`R_s>>R_v`$, and the luminosity within $`100h^1`$ kpc. The lower values with respect to the global $`L`$$`T`$ relation is due to a factor of $`1/3`$ in luminosity due to the exclusion of the low surface–brightness gas at radii larger than $`100h^1`$ kpc, and by the factor of $`2`$ gained in the emission weighted temperature since only the inner regions, with strong temperature gradients, are included. Thus, in the simple scenario of an external $`K_{}`$, the groups are expected to be surrounded by a large halo of surface–brightness $`10^{16}`$ erg s<sup>-1</sup> cm<sup>-2</sup> arcmin<sup>-2</sup>. Its detection would constitute an important test for the external entropy scenario (TSN00).
In Figure 11 we also show the prediction for the luminosity within $`100h^1`$ kpc in the cases $`K_{34}=0.2`$, which turns out to give better fits for the groups. Thus, even if clusters with $`k_BT>2`$ keV seems to require $`K_{34}0.30.4`$, a lower value $`K_{}0.2`$ gives a better fit to the low end of the $`L`$-$`T`$ relation. As we will see, this is confirmed by the entropy–temperature relation (see Figure 17).
It is clear how the presence of the background entropy bends the $`L`$$`T`$ relation from the self–similar slope to an average $`LT^3`$. However, with this simple model it is difficult to reproduce the steepening below $`1`$ keV. This is partially due to inclusion of line emission, that prevents the $`L`$$`T`$ relation from reaching the adiabatic slope $`LT^5`$. In fact, for a metallicity $`Z>0.1Z_{}`$ the slope of the emission curve between $`0.31`$ keV is virtually zero, or even negative. In this case the asymptotic slope will be flatter than $`T^4`$.
The $`M`$$`T`$ relation at small masses is lower with respect to the relation between mass and virial temperature (dashed lines, see Equation (2.2) in Eke et al. 1998) which is reproduced by our self–similar case. The predicted $`M`$$`T`$ relation in $`\mathrm{\Lambda }`$CDM with $`K_{34}=0.4\pm 0.2`$ is consistent with the recent finding of Nevalainen, Markevitch & Forman (2000). Note that the values plotted in Figures 13 and 14 are re-scaled to the virial mass from the mass quoted in the paper, using the corresponding NFW profile. The steepening of the temperature profiles in the adiabatic cores gives higher emission weighted temperatures, about $`25`$% larger than the corresponding virial temperatures for $`k_BT<2`$ keV. This translates into an uncertainty of less than a factor of $`2`$ in the total mass (using the self–similar relation). The evolution is similar to that of the self–similar case, and the difference in slope is preserved. In the tCDM case, the $`M`$$`T`$ relation is higher and gives a poor fit to the data of Nevalainen, Markevitch & Forman (2000).
The slope of the $`L`$$`T`$ relation is affected by different values of $`K_{}`$ as shown in Figure 15 at $`z=0`$ in a $`\mathrm{\Lambda }`$CDM universe. Lower values gradually approach the self–similar relation $`LT^2`$. However, the self–similar scaling is never reached in the limit $`K_{}0`$, due to the cooling catastrophe. We recall that our ability to include the cooling processes in the cases presented here, is due to the non–negligible initial entropy level. If $`K_{34}<<0.05`$, the cooling processes are too strong and our computation scheme becomes inadequate. The $`M`$$`T`$ relation is less affected by changes in $`K_{}`$ (see Figure 16).
All the above physics influences the relation between the central entropy (measured at a radius $`r=0.1R_v`$) and the temperature, as shown in PCN. The emergence of the entropy floor at small scales (low temperatures) is directly seen as a departure from the self–similar expectations, shown as a dashed line in Figure 17<sup>1</sup><sup>1</sup>1The entropy is computed using the predicted local value of the temperature at $`r=0.1R_v`$; very similar values are obtained using the emission weighted temperature as effectively used in PCN. Note that in this case the adiabat is defined differently, using the electron density instead of the mass density: $`K_Pk_BT/n_e^{2/3}`$ keV cm<sup>2</sup>. The relation between the two definitions is $`K_P=0.95\times 10^3K_{34}`$. In this respect, the value observed should be considered indicative of the average entropy in the center of the halos. The entropy floor is clearly matched at $`k_BT<2`$ for $`0.1<K_{34}<0.4`$. In particular, $`K_{34}=0.20.3`$ reproduce both the $`L`$$`T`$ and $`K`$$`T`$ relations over the whole temperature range.
## 5 THE ENTROPY HISTORY OF THE UNIVERSE AND THE X–RAY EVOLUTION OF CLUSTERS
From the above results, it is clear that a significant background entropy, $`S_{}`$, present in the IGM before the formation of large dark matter halos affects the X–ray properties of groups and clusters and can explain many scaling properties. However, the assumption of a uniform floor of entropy for all the baryons could be too simplistic. As we showed, the data seems to require a growing value of $`K_{}`$ at larger mass scales: $`K_{34}0.2`$ for $`k_BT<2`$ keV, and $`K_{34}0.4`$ for $`k_BT>2`$ keV. In terms of physical mechanisms, it is reasonable to expect that $`S_{}`$ is correlated with higher density regions where star formation or nuclear activity preferentially occurs. For example, if the excess entropy is linked to star formation processes, an entropy excess should be observed in the diffuse baryons expelled by galaxies at high redshift. The distribution of entropy should follow the light distribution, and should show a dependence on cosmic time that parallels the birth of the first stars and QSOs. This topic can be addressed not only with X–ray observations, but also with the UV and optical investigation of the low density baryons detected, e.g., as Ly$`\alpha `$ clouds. Here we will discuss in greater detail the scenario with a uniform external entropy, but relaxing the assumption of a constant $`K_{}`$.
We already know that the IGM which is observed in high–$`z`$ Ly$`\alpha `$ clouds generally shows an entropy level lower than that observed in the centers of groups. An approximate relation derived from the observations is $`K_{Ly\alpha }=(1.2\pm 0.5)\mathrm{\hspace{0.17em}10}^2(1+z)^1\times 10^{34}`$ erg g<sup>-5/3</sup> cm<sup>2</sup> (extrapolated from Figure 10(b) in Ricotti et al. 2000, see also Schaye et al. 1999). Thus, the ratio of the value $`K_{gr}`$ observed in the center of the groups to that observed in Ly$`\alpha `$ is about $`K_{gr}/K_{Ly\alpha }10(1+z)`$. This may indicate that the ICM baryons undergo substantial heating with respect to the baryons observed in Ly$`\alpha `$ or, possibly, that the baryons seen in Ly$`\alpha `$ clouds are not the same baryons that will be later accreted in clusters. Furthermore, the chemical properties of the IGM seen in the Ly$`\alpha `$ forest are clearly different from those of the ICM in clusters, showing that the ICM was affected by star formation processes and chemical enrichment to a larger extent with respect to the Ly$`\alpha `$ clouds, with a commensurate amount of entropy production. In this respect, it will be interesting to observe the tenuous gas being accreted in the outskirts of nearby, large clusters, but not yet shocked, or at large radii in small groups, and compare it with the gas observed in different enviroments at different cosmic epochs. Such observations would complement the investigation of the entropy excess as observed in nearby and distant clusters.
As expected, the evolution of the background entropy affects both the evolution and the shape of the $`L`$$`T`$ relation. We already emphasized the fact that the uncertainty in the evolution of the $`L`$$`T`$$`M`$ relations reflects on the uncertainty in the derivation of cosmological parameters from the cluster abundance evolution. The $`L`$$`M`$ relation is, in fact, the link between the cluster mass function (predictable for a given cosmology with numerical or analytical calculation) and the observed X–ray luminosity distribution. The complexities due to the evolution in the luminosity are only partially avoided when directly using the temperature. In fact, selection effects for flux limited samples add to the evolution of the emission weighted temperatures (see Eke et al. 1998).
If the minimum background entropy $`S_{}`$ is kept constant at every epoch, the evolution of the $`L`$$`T`$ relation is essentially frozen, or mildly negative, even at redshifts as high as $`z=1`$, as already shown in Figures 11 and 12. The evolution of $`L`$ at fixed $`T_{ew}`$ is negative especially at small temperature. This global behaviour is in agreement with the claim for null evolution of the $`L`$$`T`$ at redshift $`z0.4`$ (Mushotzky & Scharf 1997). A non evolving $`L`$$`T`$ relation, suggested also by the present data on the luminosity function at high redshifts ($`z>0.5`$), would strengthen a low density, eventually flat, universe (c.f. Borgani et al. 1999).
We can investigate how the evolution of the $`L`$$`T`$ is affected if the background entropy evolves substantially with epoch. In the Figure 18 and 19 we assumed $`K_{34}(z)=0.8(1+z)^1`$, which is an evolution that parallels the one observed in the Ly$`\alpha `$ clouds. In this scenario, objects observed at redshift $`z=1`$ have accreted most of their baryonic mass when the entropy was lower, and thus mostly in the shock regime. This allows the cooling to start earlier, and be more efficient. As a net result, the $`L`$$`T`$ and the $`K`$$`T`$ relations at $`z=1`$ are higher with respect to the predictions of the constant $`K_{}`$ scenario. However the positive evolution is about a factor of two, much less than the intrinsic scatter, and very difficult to observe. Such a positive evolution is too small to reconcile a critical universe with the observed high redshift luminosity function. As a further comment, we recall that the large discrepancy between the average level of entropy seen in Ly$`\alpha `$ clouds and that observed in the center of groups, implies that the Ly$`\alpha `$ gas is not the same or the heating rate is much steeper than this. We therefore adopt this entropy evolution as a reference case.
Assuming $`K_{34}=0.8(1+z)^1`$, gives a good fit to the whole temperature range without requiring further dependence on the mass scale. This is because the evolution $`(1+z)^1`$ introduces by itself such a dependence. The core of intermediate mass halos are assembled at $`z1`$, for an effective $`K_{34}(z=1)=0.4`$, while low mass objects build their cores at redshifts $`z23`$, for $`K_{34}(z=3)=0.2`$. Also an evolution as strong as $`K_{34}=3(1+z)^2`$ provides a good fit to the data.
## 6 DISCUSSION
The main limitation of this model is clearly the adopted spherical symmetry and also the assumptions of isotropic and continuous infall. In the real world, some of the baryons are accreted in the form of smaller clumps and substructure, and flow along sheets and filaments. The spherical infall model used here does not include the effects of larger and smaller scale perturbations. Moreover, there are missing ingredients in the physics of baryons. We shall briefly discuss them in turn.
The presence of large–scale structure is not expected to affect strongly the accretion rates and in general the statistical properties of dark matter halos. In fact, the rates used in this work are derived from the PS formalism, which proves accurate within few percent when compared to N–body simulations that include large–scale structure (see, e.g., Governato et al. 1998). However, an effect of the large–scale structure which is of interest here, could be the eventual contribution to the initial entropy in the IGM due to shocks occurring on large scales related to the formation of filaments. Hierarchical gravitational processes do not break the self–similarity, but the anisotropic collapse can produce widespread shocks that raise the average entropy level in the IGM everywhere without being associated with the formation of halos. The baryons that fall in the isotropic potential wells at the intersection of sheets and filaments could be already heated by an amount which depends on the power spectrum on large scales. This can break the self–similarity of the baryons, assuming that the large–scale heating is effective almost uniformly in the IGM.
Focusing on smaller scales, the presence of substructure in the infalling matter necessarily introduces some stochasticity in the accreting processes. The intrinsic scatter in the density and the temperature of the accreted baryons translates into a scatter in the observational quantities (see, e.g., CMT97). The presence of substructures implies some gravitational energy is transferred to the baryons before they are accreted into the main potential well and shocked for the last time. However, the gravitationally–produced entropy on small scales is very different from the above mentioned large scale production. In fact, the mass distribution of satellite halos scales self–similarly with the total mass of the final halo. Thus the amount of entropy given to the baryons in substructures scales with the final mass, and does not produce any break of self–similarity. This entropy contribution can be included in the external entropy, $`K_{}`$, without any distinguishing effect with respect to the mass scale.
Another point related to the dark matter is the case of very massive merger events, where a massive, disruptive event is defined by the mass ratio of the merging halos being larger than about $`0.3`$ (see Roettiger et al. 1998). In these cases it is likely that the ICM is strongly stirred, and, if the lookback time of the event is less than $`1`$ Gyr, the ICM is not even in hydrodynamical equilibrium at the epoch of observation. Massive mergers can also create situations of non–equilibrium ionization (see Ettori & Fabian 1998; Takizawa 2000). It is clear that the model cannot describe the population of such disturbed clusters. In the PS formalism, the fraction of objects that are subject to large merger events is a sensitive function of both the total virial mass $`M_0`$ and the observation epoch $`z_0`$. We calculate that the expected number of major mergers in the last Gyr is between $`0.1`$ and $`0.2`$ in tCDM, and a factor of 2 lower in $`\mathrm{\Lambda }`$CDM, at $`z_0=0`$ (for a mass range between $`10^{15}`$ and $`10^{13}`$ $`h^1M_{}`$). However, such numbers grow to $`1`$$`0.5`$ at $`z_0=1`$ in tCDM and $`0.60.3`$ in $`\mathrm{\Lambda }`$CDM. In this framework, it is reasonable to expect that at $`z1`$, a fraction between $`1/3`$ and $`1/2`$ of the population of clusters has undergone a massive merger event with a lookback time less than 1 Gyr. This has to be regarded as an intrinsic limitation to statistical analyses of the population of high redshift clusters.
Other limitations come from the more complex physics of baryons. An important issue is that the entropy in the center may increase because shocks propagate in the inner part of the halos due to infalling gas along philaments (A. Klypin 2000, private communication). We stress however, that in order to survive the outer shock and propagate in the very central part of the halo, the infalling baryons should be compressed already. The presence of an initial entropy level will inhibit the formation of dense knots of gas at least on small scales, and thus inner shocks are probably limited only to very massive mergers.
Another important component, which is not included in the present model, is the momentum gained from the heated gas, that can push part of the baryons out of the halos without contributing to the average heating. This effect is very difficult to model a priori. Its effect on the X–ray emission can be computed by including semianalytical models of galaxy formation (see Menci & Cavaliere 2000).
Finally, gravitational effects of the baryons on the dark matter profile are neglected. These can be important in the very center, where the baryons can concentrate in the form of cooled gas and contribute to density peaks which may affect the X–ray emission (see Pearce et al. 2000, Lewis et al. 2000).
## 7 CONCLUSIONS
We have presented a detailed model to relate the X–ray properties of diffuse baryons in clusters of galaxies to the entropy history of the cosmic baryons, after including adiabatic compression, shock heating and cooling. Our aim is to build a useful tool to reconstruct the entropy history of the universe from the observations of local and distant clusters. In particular, a major goal is to identify and follow in time the processes that generate the entropy excess. This entropy excess is now probed by many observations and it is connected with many scaling properties of X–ray halos. Even if a given entropy excess does not translate into a unique heating history, the comparison of X–ray data with observations in other bands may allow identification of the major heating sources. Favoured candidates are star formation processes and nuclear activity. At present, however, neither the epoch, nor the source of the related heating process have been identified.
In this paper we have limited the investigation to a scenario in which the excess entropy is present since very high $`z`$ and is uniform throughout the IGM. A case with an external entropy decreasing with redshift, mimicking the rise of a population of heating sources, is also presented. In both the constant and time–evolving case, the scaling properties of local clusters of galaxies are reproduced on a large range of scales, with an appropriate choice of the free parameter $`K_{}`$. The properties of distant X–ray halos are predicted to be generally similar to properties of the local population, but significative differences can be actually observed by the present–day X–ray satellites, shedding light on the thermodynamics history of the ICM. We recall here the general results on density and temperature profiles, together with the results on the evolution of the global X–ray properties, especially luminosity and emission weighted temperatures.
The bending of the $`L`$$`T`$ relation with respect to the self–similar case $`LT^2`$, is due to the flatter profiles of the ICM going from large mass to small mass halos. Good fits are obtained for a background entropy in the range $`K_{}=(0.2\pm 0.1)\times 10^{34}`$ erg cm<sup>2</sup> g<sup>-5/3</sup> for $`k_BT_{ew}<2`$ keV, and $`K_{}=(0.4\pm 0.1)\times 10^{34}`$ erg cm<sup>2</sup> g<sup>-5/3</sup> for $`k_BT_{ew}>2`$ keV. This scale dependence can be introduced by an evolution in the effective value of $`K_{}`$. In particular, $`K_{34}=0.8(1+z)^1`$ gives a good fit over the whole range of observed temperatures.
The central regions of groups and clusters, which dominate the X–ray emission, are formed during the initial stages of accretion. In these early phases, if a significant background entropy is present, the accretion is adiabatic, and the gas is compressed in a flat, low density profile with steep temperature gradients. This is relevant for the smallest halos, where the gravity does not overcome the pressure support of the baryons for the majority of the subsequent accretion of gas. In clusters the infall velocities rapidly become larger than the sound speed, and the shock regime takes over. In the outer regions of clusters the entropy is entirely due to gravitational processes, and the entropy profile is a featurless power law approaching $`Kr^{1.1}`$.
This mechanism is particularly efficient if cooling is neglected. However, it is known that the cooling is an important ingredient in the history of the ICM. The main effect is that the isentropic cores expected in the constant entropy scenario, are partially erased by the process of cooling. Still, if $`K_{}>0.1\times 10^{34}`$ erg cm<sup>2</sup> g<sup>-5/3</sup>, the cooling processes are significantly suppressed and the inner regions of the halos keep the imprint of the initial entropy level. Cooling processes appear again only in massive halos, where the gravity dominates the energy of the system and the excess entropy is no longer able to keep the gas at low density. In the extreme case of negligible $`K_{}`$, it is worth noting that the cooling processes alone would have a dramatic effect on both small and large mass halos. In small mass halos ($`10^{13}M_{}`$) most of the gas is expected to cool and recombine, causing a central baryonic catastrophe.
Other important characteristics are found in the temperature structure especially of smaller halos. Temperature gradients are commonly expected both in clusters and in groups. The polytropic index is predicted to be $`\gamma _p0.91.2`$ in the region where the gas is shock heated. The polytropic index can be higher if the dark matter profile is centrally peaked (e.g., with a power law with index $`1.4`$, see Moore et al. 1998). Another relevant observable (for local halos) is the position of the final shock radius, which is expected to be close to the virial one at large mass scale, while it migrates to larger radii in small groups. In the smallest halos, in fact, the shock is vanishingly small. As a function of epoch, for a given object, the shock/accretion radius is initially quite distant from the virial radius. It is very close to the virial radius when the mass accretion rate reaches its maximum and the shock regime is well developed. Eventually, the mass accretion rate decreases (especially in the $`\mathrm{\Lambda }`$CDM universe) and the shock radius relaxes again to larger positions. A consequence of the above picture is that the ratio of the baryonic mass included in the virial radius to the total mass, is always lower but still close to unity; it can be significantly lower ($`1/3`$) only for small mass halos (corresponding to emission–weighted temperatures of $`0.31`$ keV).
It is remarkable that the simple presence of an initial excess entropy in the diffuse IGM can reproduce many of the scaling properties of the observed X–ray halos, without the contribution of any internal heating. It is interesting to discuss the implications of this simple scenario for the energetic budget and the past cosmic star formation history. The minimum excess energy associated with an initial background entropy $`K_{34}`$ is about
$$k_BT0.1\left(\frac{K_{34}}{0.4}\right)\mathrm{keV},$$
(16)
where the gas is assumed to be at the background density at the epoch of the heating. However, we can speculate on the energy budget when the entropy excess is generated after the collapse, at much larger densities (the internal scenario).
Following PCN, we can establish a relation between the epoch of heating and the energy released. Under the assumption that the heating process can be described with a single epoch and a typical overdensity, we have:
$$1+z_h=\left(\frac{k_BT_h}{3.2\times 10^2K_{34}}\right)^{1/2}\delta ^{1/3},$$
(17)
where $`k_BT_h`$ is the average energy per particle released in the IGM by non–gravitational processes. If we adopt the conservative scenario in which the gas is heated at a typical virial density ($`\delta 200`$), to have an entropy level in the range $`K_{}=(0.4\pm 0.2)\times 10^{34}`$ erg cm<sup>2</sup> g<sup>-5/3</sup>, we obtain:
$$1+z_h(1.5\pm 0.3)\left(\frac{k_BT_h}{1\mathrm{k}\mathrm{e}\mathrm{V}}\right)^{1/2}.$$
(18)
Thus, if we want heating at $`z>1`$ in order to avoid the overcooling catastrophe, the energy budget must be larger than $`1`$ keV per particle. The above estimate would give even larger values after the inclusion of cooling. In fact, if the gas is heated at high densities, most of the extra energy is likely to be re–emitted soon, and this would raise the energetic budget for a given final entropy excess. In this respect, the relation between the epoch of heating and the energy released is strongly dependent on the physical process. Of course, a scenario in which the extra entropy is provided by the contributions of several different sources, active at different epochs, is a likely possibility. In this perspective, the measure of metallicities as a function of the entropy of the baryons in different systems, from Ly$`\alpha `$ clouds to rich clusters, may be useful in determining whether the excess entropy is linked to star formation processes.
The assumption of an initial excess entropy uniformly diffused in the IGM, offers new perspectives in the approach to cluster formation, but also galaxy formation. Such an entropy background, once established, may affect the star formation itself, since the cooling processes on all scales are virtually inhibited. This is the mechanism which is expected to solve the cooling catastrophe (see White & Rees 1978, Blanchard, Valls Gabaud & Mamon 1992; Prunet & Blanchard 1999) and in this view X–ray clusters and galaxy formation processes are intimately related. Current attempts to model ab initio the physics of the heating process, and then link the entropy history of the cosmic baryons to galaxy formation, must include the well known plethora of ingredients that has been already mentioned several times: feedback from star formation processes and SNe explosions, radiative and mechanical heating from active galactic nuclei, radiative heating from hard X–ray background, gravitational heating on large scale filaments (see Menci & Cavaliere 2000; Valageas & Silk 2000; Wu, Fabian & Nulsen 2000; Madau & Efstathiou 1999; Cen & Ostriker 1999). Such different scenarios allow for different entropy histories of the universe, determining both the spatial distribution and the evolution of the entropy in the diffuse gas.
A promising strategy for the near future is to look directly for the distribution of the entropy in the ICM (TSN00). A direct consequence of assuming a uniform entropy everywhere in the gas, is that the groups are expected to be surrounded by large halos of low surface brightness gas, spread out over radii much larger than the virial radius of dark matter halos. This low–density gas may have been missed by observations with the ROSAT satellite, but can be detected by the XMM satellite. Its emission can enhance the total luminosity of the groups by more than a factor of $`3`$, including the lowest energy bins of $`0.1`$ keV. Another promising observational channel is the absorption from metals in the gas seen against bright X or UV sources. If the source of the background entropy is star formation, significant pollution by metals is expected.
The model presented here is to be considered a useful tool to interpret the observations of high redshift clusters, that will be provided especially by the Chandra and XMM satellites. Our aim is to build a solid link between the thermodynamics of the diffuse cosmic baryons and the emitting properties of X–ray halos, in order to be able to reconstruct the entropy history of the universe, at high and low redshifts, from spectral and imaging X–ray observations. This will help in understanding the source of the entropy excess and the time evolution of the corresponding heating process.
We acknowledge discussions with S. Borgani, N. Menci, and P. Rosati. We thank T.J. Ponman for discussions and for providing the data in Figure 17. We thank R. Giacconi for discussions and continuous encouragement. We thank the referee, Greg Bryan, for detailed comments. PT thanks ESO Garching for hospitality during the completion of this work. This work has been supported by NASA grant NAG 8-1133.
## Appendix A The infall velocity
We find upper and lower limits for the infall velocity of the accreted baryonic shells computed with Equation (3). The last two terms of Equation (3) depend on the densities of the shell at accretion ($`\rho _e`$) and at turnaround ($`\rho _{ta}`$). The values of the two densities are derived requiring conservation of mass, and assuming that baryons and dark matter are still not decoupled at turnaround. In particular, the exact value of $`\rho _e`$ depends on the validity of the Equation (5), which is based on the assumption that the total baryonic mass accreted at every epoch is $`f_BM_v`$, where $`f_B`$ is the universal baryonic fraction. Such an assumption can be tested with numerical simulations, and will not be discussed here.
We focus on the numerical uncertainty in the estimate of the term $`\mathrm{\Delta }W`$. The total work per unit baryonic mass done by the gravitational potential on the baryonic shell is:
$$W_{R_{ta}}^{R_s}\frac{GM(<r)}{r^2}𝑑r,$$
(A1)
where the integral is computed along the trajectory $`r(t)`$. If the total mass within the shell were constant, the solution would be simply:
$$W=\frac{v_{ff}^2}{2}\left(\frac{GM_{ta}}{R_S}\frac{GM_{ta}}{R_{ta}}\right),$$
(A2)
where $`M_{ta}`$ is the total mass initially contained in the turn–around radius $`R_{ta}`$. The free fall velocity $`v_{ff}`$ refers to a test particle falling from turnaround to the shock radius, experiencing a gravitational force always from the same amount of matter. However, the actual mass enclosed by a given baryonic shell will depend on time. We can write:
$$M(<r)=[f_B+(1f_B)Y(t)]M_{ta},$$
(A3)
since the baryonic mass inside the shell is constant, but the amount of dark matter can change by a time dependent factor $`Y(t)`$. The complete solution can now be formally written as:
$$W=\frac{v_{ff}^2}{2}+\mathrm{\Delta }W=\frac{v_{ff}^2}{2}\frac{GM_{ta}}{r_{ta}}_{R_{ta}}^{R_s}\frac{(f_B+(1f_B)Y(t)1)}{r^2}𝑑r.$$
(A4)
At this point we note that the amount of mass that is included in a given baryonic shell along its trajectory is always larger than the initial mass $`M_{ta}`$, since the collisionless shells of dark matter fall faster than the baryonic shells, which, instead, are pressure supported. Here, we neglect the shell crossing and the detailed behaviour in time, but we recall that we want the solution only at the accretion radius, which usually occurs just inside the most external caustic of the dark matter (see the self–similar model of Bertschinger 1985). Thus, we can safely assume that the total mass can only grow inside the baryonic shell. The mass excess $`\mathrm{\Delta }M/M_{ta}=(f_B+(1f_B)Y(t)1)`$ can be described with a generic power law dependence on the actual position $`r(t)`$ of the kind:
$$f_B+(1f_B)Y(t)1=(f_B+(1f_B)Y_s1)[(1\left(\frac{r(t)}{r_{ta}}\right)^\alpha ][1\left(\frac{r_s}{r_{ta}}\right)^\alpha ]^1$$
(A5)
where $`\alpha >0`$, and $`Y_s`$ is the value at the accretion. To calculate $`Y_s`$ we must know the dark matter density profile at radii larger than the virial radius. We do not propose a specific model here, instead we simply use the density profile as computed in Bertschinger (1985) as a reasonable approximation at radii larger than the virial one. We can substitute Equation (A5) in Equation (A4) and integrate, obtaining an estimate of $`W`$ as a function of $`\alpha `$. To eliminate the dependence on $`\alpha `$, we take the limit for small and large values of $`\alpha `$, to obtain the upper and lower values for $`\mathrm{\Delta }W`$:
$$\mathrm{\Delta }W=[f_B+(1f_B)(Y_s1)]\frac{GM_{ta}}{r_{ta}}\left(\left[\frac{1x_s}{x_s}+\frac{\mathrm{ln}(x_s)+1x_s}{\mathrm{ln}(x_s)x_s}\right]\pm \left[\frac{1x_s}{x_s}\frac{\mathrm{ln}(x_s)+1x_s}{\mathrm{ln}(x_s)x_s}\right]\right).$$
(A6)
The last term in Equation (A6) bounds the possible values for $`\mathrm{\Delta }W`$, assuming a monotonic increase of the total mass enclosed by the infalling shell. The upper and lower values turn out to be between $`10`$ % and $`30`$ % during the mass history of a given halo, and are plotted in Figure 1 as dotted lines. This reflects our error in computing the infall velocities of the baryonic shells. The uncertainty in the infall velocities does not strongly affect the mass scale at which the adiabatic/shock transition occurs, since the dependence of $`v_i`$ on the accreted mass is very steep when shocks begin to appear. This effect is related to the fast migration of the accretion radius from $`2R_v`$ to $`R_v`$ (see Figure 3).
## Appendix B Cooling processes
Here we discuss how to compute the effect of the radiative cooling on each baryonic shell. The treatment of the cooling is complex and constitutes the largest uncertainty in modelling the X–ray emission from clusters in present–day numerical simulations, since the predicted luminosity of the central region can heavily depend on the adopted resolution (see, e.g., Suginohara & Ostriker 1998).
There is of course no difficulty in solving Equation (8) as long as $`\tau _{cool}>\mathrm{\Delta }t`$ where $`\mathrm{\Delta }t`$ is the time resolution. However, the time resolution needed increases dramatically when the density increases and $`\tau _{cool}\mathrm{\Delta }t`$, since the cooling is a runaway process. Since our calculation is based on a sequence of hydrostatic equilibria, and we do not want to end up with an heavy computation effort, we propose to use a reasonable time step (of the order of few tenths of Gyr) and solve analytically the energy Equation (8) for each shell within each time step. To do this we first assume that the cooling proceeds isobarically within $`\mathrm{\Delta }t`$, and compute the new value of the pressure after each step to take into account the new equilibrium positions of each shell.
If the pressure is constant for each shell within $`\mathrm{\Delta }t`$, the density can be expressed as a function of the adiabat $`K`$ only, to give:
$$\rho =p^{1/\gamma }k^{1/\gamma },$$
(B1)
where $`\gamma =5/3`$ and the variables are assumed to be normalized to the shock values as usual. The temperature is then:
$$t=k^{1/\gamma }p^{(\gamma 1)/\gamma }.$$
(B2)
Following Sutherland & Dopita (1993), we define the normalized cooling function $`\mathrm{\Lambda }_N\mathrm{\Lambda }_{net}n_en_i`$, where $`n_e`$ is the electron number density and $`n_i`$ is the ion number density. For an average metallicity $`Z=0.3Z_{}`$ we can approximate $`n_en_i=0.704(\rho _B/m_p)^2`$. The cooling time now can be expressed as a function of the adiabat $`K`$ and the normalized cooling function $`\mathrm{\Lambda }_N`$:
$$\tau _{cool}2.13\frac{kT_{s0}m_p}{\mu \rho _{s0}}k^{2/\gamma }p^{(\gamma 2)/\gamma }\mathrm{\Lambda }_N^1,$$
(B3)
where the subscripts “$`s0`$” refer to the value at the shock. To write an analytic expression, we approximate $`\mathrm{\Lambda }_N`$ with a polynomial form:
$$\mathrm{\Lambda }_N=C_1(kT)^\alpha +C_2(kT)^\beta +C_3,$$
(B4)
where the exponents take the values $`\alpha =1.7`$ and $`\beta =0.5`$. The constants depend on the assumed metallicity, and are chosen as in Table 2 in order to reproduce the cooling function of Sutherland & Dopita (1993) within few percent in the energy range $`k_BT>0.03`$ keV.
Thus, using the canonical value $`\gamma =5/3`$, the cooling time can be written as:
$$\tau _{cool}=C_\tau T_{s0}\frac{k^{6/5}p^{1/5}}{C_1T_{so}^\alpha k^{\frac{3}{5}\alpha }p^{\frac{2}{5}\alpha }+C_2T_{so}^\beta k^{\frac{3}{5}\beta }p^{\frac{2}{5}\beta }+C_3}.$$
(B5)
The constant $`C_\tau `$ factorizes out the terms that depend on the shock condition, and can be written as:
$$C_\tau =1.62\times 10^2[f_B(1f_{cool}f_{})h^2(z)\delta _{so}g_{s0}]^1\mathrm{Gyr}/\mathrm{keV},$$
(B6)
where $`\delta _{s0}`$ is the overdensity with respect to the critical density at redshift $`z`$, $`T_{s0}`$ is the temperature at the shock and $`g_{s0}`$ is the compression factor at the shock; $`f_{cool}`$ and $`f_{}`$ are respectively the fraction of baryons cooled in the center and the fraction of baryons locked into stars.
Equation (8) can be recast in term of the adiabat $`K`$ only, and the final adiabat $`k_f`$ can be recovered implicitly from the solution in the finite time step $`\mathrm{\Delta }t`$ (expressed in Gyr):
$$\mathrm{\Delta }t=C_\tau T_{s0}_{k_i}^{k_f}𝑑k\frac{k^{1/5}p^{1/5}}{C_1T_{so}^\alpha k^{\frac{3}{5}\alpha }p^{\frac{2}{5}\alpha }+C_2T_{so}^\beta k^{\frac{3}{5}\beta }p^{\frac{2}{5}\beta }+C_3}F(k_i,k_f)$$
(B7)
In particular, the condition $`F(k_i,0)<\mathrm{\Delta }t`$ determine if a shell with initial entropy $`k_i`$ cools completely within $`\mathrm{\Delta }t`$. At each epoch, the region comprised within the largest shell for which $`F(k_i,0)<\mathrm{\Delta }t`$ is included in the cooled fraction $`f_{cool}`$ and excluded from the diffuse, emitting phase.
## Appendix C Concentration parameters
The concentration parameters of the dark matter profiles depend on epoch and cosmology, as shown in the numerical works of Navarro et al. (1997) or analytical models (see, e.g., Lokas 2000). A general trend is that lower mass halos are more centrally concentrated than high mass halos by virtue of the higher redshift of formation. For the same reason, halos of the same virial mass, but observed at higher redshifts, are less concentrated, since the difference in the average density at the formation and at the observation is smaller with respect to low redshift halos. The mass dependence of the concentration parameter, however, can be well approximated with power laws which change slightly as a function of epoch and cosmology. In this paper we used the following approximations:
$`c=8.5M_{15}^{0.086}\mathrm{\Lambda }\mathrm{CDM},z=0`$ (C1)
$`c=5.4M_{15}^{0.070}\mathrm{\Lambda }\mathrm{CDM},z=1`$ (C2)
$`c=5.5M_{15}^{0.070}\mathrm{tCDM},z=0`$ (C3)
$`c=4.4M_{15}^{0.046}\mathrm{tCDM},z=1`$ (C4)
where $`M_{15}M/(10^{15}h^1M_{})`$.
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# Compact singularities of meromorphic mappings between complex 3-dimensional manifolds
## Introduction
The study of the extendibility of holomorphic and meromorphic mappings began with the classical theorem of Hartogs \[Ha\] (see \[Si\]).
> Let $`K`$ be a compact subset of a domain $`M^n`$, $`n2`$, such that $`MK`$ is connected, and let $`f:MK`$ be a holomorphic function. Then there exists a holomorphic function $`\widehat{f}:M`$ extending $`f`$, i.e., $`\widehat{f}_{MK}=f`$.
Shortly after Hartogs proved his theorem, E. E. Levi \[Le\] discovered that this extension result holds true also for meromorphic functions.
A natural problem is to understand under what conditions Hartogs’ Theorem (respectively Levi’s Theorem) holds when the mapping $`f`$ takes values in a general complex manifold $`X`$ rather than $``$ (respectively $`^1`$). Of course it is immediate that Hartogs’ Theorem remains valid for holomorphic mappings with values in a Stein manifold $`X`$, since such a manifold $`X`$ can be embedded into $`^N`$. It similarly follows that Levi’s Theorem also remains valid for meromorphic mappings with values in compact projective manifolds.
In 1971, Griffiths \[Gr\] and the second author \[Sh\] independently showed that Hartogs’ Theorem is valid for holomorphic mappings into manifolds $`X`$ carrying a complete Hermitian metric with non-positive holomorphic scalar curvature, answering a question was asked by Chern in \[Che\]. Concerning the meromorphic mapping problem, the first author \[Iv1\] proved that Hartogs extension holds for meromorphic maps into compact Kähler manifolds.
We recall two more results here due to K. Stein and M. Chazal. Stein proved in \[St\] that Hartogs’ Theorem holds for holomorphic maps if $`dimXn2`$. Recently Chazal \[Cha\] relaxed this condition to $`dimXn1`$ and more generally $`f`$ can be meromorphic. The next case of interest is the equidimensional case $`dimX=n`$. It is well known that one doesn’t always have meromorphic extension in this case, as is illustrated by the (holomorphic) projection $`f:^nX=^n/`$ to the Hopf manifold. (The $``$-action is given by $`z\stackrel{n}{}2^nz`$.) The goal of this paper is to show that, at least for dimension $`3`$, the singularity at 0 of the Hopf map $`f`$ is the only type of singularity that can occur for equidimensional meromorphic maps:
###### Theorem 1.
Let $`K`$ be a compact set with connected complement in a Stein manifold $`M`$ of dimension $`3`$, let $`X`$ be a compact complex manifold of the same dimension and let $`f:MKX`$ be a meromorphic map. Then there exists a finite set $`\{a_1,\mathrm{},a_d\}K`$ such that $`f`$ has a meromorphic extension $`\widehat{f}:M\{a_1,\mathrm{},a_d\}X`$, and if $`B(a_j)`$ are disjoint coordinate balls centered at $`a_j`$, then $`\widehat{f}(B(a_j))`$ is not homologous to zero in $`X`$ ($`1jd`$).
The same result when both $`M`$ and $`X`$ have dimension two was proved by the first author in \[Iv3, Corollary 4(b)\]. It is open whether this result if valid for equidimensional maps of dimension greater than 3. Of course, one cannot expect to obtain such results when the dimension of $`X`$ is greater then the dimension of $`M`$; see the remark in §1 below.
In the case of the Hopf map $`f:^3\{0\}^3/`$ mentioned above, $`A=\{0\}`$. Of course, the elements $`\widehat{f}(B(a_j))`$ of the fifth homology group are very special; they are often called spherical shells. If, for example, $`\widehat{f}`$ is a holomorphic embedding in a neighborhood of $`B(a_j)`$, then $`X`$ is of a very restricted type: it is a deformation of the Hopf 3-fold (see \[Ka1\]).
In particular, if the singular set $`A`$ is nonempty, then $`H^5(X,)0`$. Poincaré duality then implies the following:
###### Corollary 2.
If under the conditions of Theorem 1, $`H^1(X,)=0`$, then $`f`$ extends meromorphically to all of $`M`$.
## 1. Reductions
For degenerate mappings, the result is known and is due to F. Chazal \[Cha\]. Hence, in the sequel, we always suppose that $`f`$ is nondegenerate; i.e., $`\mathrm{rank}f=3`$.
We let $`\mathrm{\Delta }(r)=\{z:|z|<r\}`$ denote the disk or radius $`r`$ about $`0`$, and we write $`\mathrm{\Delta }=\mathrm{\Delta }(r)`$. We consider the polydisk $`\mathrm{\Delta }^n(r)=\mathrm{\Delta }(r)^n`$ and “annulus” $`A^n(r,1)=\mathrm{\Delta }^n\overline{\mathrm{\Delta }^n(r)}`$. We shall make frequent use of the following Hartogs figure in $`^3`$:
$$H_1^2(r)=\left[\mathrm{\Delta }(1r)\times \mathrm{\Delta }^2\right]\left[\mathrm{\Delta }\times A^2(r,1)\right].$$
By the standard method of extending analytic objects (see for example \[Si\]), it suffices to prove either of the following two equivalent results:
###### Proposition 3.
Let $`X`$ be a compact complex $`3`$-fold and let $`f:H_1^2(r)X`$ be a nondegenerate meromorphic map. Then there is a discrete set $`\{a_j\}\mathrm{\Delta }^3H_1^2(r)`$ and a meromorphic extension $`\widehat{f}:\mathrm{\Delta }^3\{a_j\}X`$ such that if $`B(a_j)`$ are disjoint balls in $`\mathrm{\Delta }^3`$ centered at $`a_j`$, then $`\widehat{f}(B(a_j))`$ is not homologous to zero in $`X`$.
###### Proposition 4.
Let $`M,W`$ be open sets in $`^3`$, and suppose $`pMW`$ such that $`W`$ is smooth and strictly pseudoconvex at $`p`$. Let $`U=M\overline{W}`$. Suppose $`f:UX`$ is a nondegenerate meromorphic map to a compact $`3`$-fold $`X`$. Then there is an open set $`\stackrel{~}{U}U\{p\}`$ such that $`f`$ has a meromorphic extension $`\widehat{f}`$ to $`\stackrel{~}{U}\{p\}`$, and either $`\widehat{f}`$ is meromorphic at $`p`$ or $`\widehat{f}(B(p))`$ is not homologous to zero in $`X`$, where $`B(p)`$ is a ball about $`p`$ contained in $`\stackrel{~}{U}`$.
We note that Proposition 4 follows from Proposition 3 with the additional simplifying assumption that $`f`$ is holomorphic on $`\mathrm{\Delta }\times A^2(r,1)`$. To see this, we observe that the set of points of indeterminacy $`I_f`$ of our meromorphic map $`f`$ has codimension at least two, i.e., is a curve together with a discrete set of points. Let $`M,W,p`$ be as in Proposition 4 and let $`\mathrm{\Delta }_{p_1}^2=\{p_1\}\times \mathrm{\Delta }^2`$ denote the vertical bidisk passing through $`p=(p_1,p_2,p_3)`$. We can assume, after making a quadratic change of coordinates, that $`\mathrm{\Delta }_{p_1}^2I_f`$ contains no curves and $`\mathrm{\Delta }_{p_1}^2\overline{W}=\{p\}`$. After translating and stretching coordinates, we then obtain a Hartogs figure $`H_1^2(r)`$ contained in $`U`$, with $`p`$ in the corresponding polydisk $`\mathrm{\Delta }^3`$, such that $`[\mathrm{\Delta }\times A^2(r,1)]I_f=\mathrm{}`$, so we can apply the modified Proposition 3 to obtain the conclusion of Proposition 4 with $`\stackrel{~}{U}=\mathrm{\Delta }^3`$.
Proposition 3 follows from Proposition 4, since the Hartogs figure can be exhausted by a family of strictly pseudoconcave hypersurfaces and this family can be continued to exhaust $`\mathrm{\Delta }^3`$. Thus, when proving Proposition 3, we may assume that $`f`$ is holomorphic on $`\mathrm{\Delta }\times A^2(r,1)`$.
Remark: The reader may observe that these results involves extension from a “$`1`$-concave” 3-dimensional domain. It is worthwhile to note that in general there is no extension of meromorphic maps with values in compact 3-dimensional manifolds from $`2`$-concave domains, such as the classical Hartogs figure
$$H_2^1(r):=\left[\mathrm{\Delta }^2(1r)\times \mathrm{\Delta }\right]\left[\mathrm{\Delta }^2\times A^1(r)\right].$$
Namely, M. Kato constructed in \[Ka2\] an example of a compact complex three-fold $`X`$ and a holomorphic mapping $`f:^2\overline{B}X`$ defined on the complement of a ball $`B^2`$, such that every point of the sphere $`B`$ is an essential singular point of $`f`$.
We shall prove Proposition 3 in §2 after we make the following reductions:
1. First of all, as was already explained, we can assume that $`f:H_1^2(r)X`$ is nondegenerate and holomorphic on $`\mathrm{\Delta }\times A^2(r,1)`$.
2. We can further assume that there is no hypersurface in $`H_1^2(r)`$ which $`f`$ sends to a point. If such a hypersurface exists, then by shrinking $`H_1^2(r)`$ a little bit, we can suppose that there are finitely many of them. Then by blowing up the image points sufficiently many times, we obtain a modification $`\widehat{X}`$ of $`X`$ together with a lift $`\widehat{f}`$ of $`f`$ to a meromorphic map $`\widehat{f}:U\widehat{X}`$ having the desired property. After extending $`\widehat{f}`$, we can push it down to extend $`f`$ itself.
3. We write
$$A_s^2=\{s\}\times A^2(r,1),s\mathrm{\Delta }.$$
After again shrinking $`H_1^2(r)`$ a little, we can suppose that $`A_s^2`$ contains no curves contracted by $`f`$ to a point, for all $`s\mathrm{\Delta }`$. Indeed, since $`\mathrm{rank}f=3`$, there are at most 1-parameter families of contracted curves. We consider small quadratic changes of the $`z_1`$ coordinate: $`\stackrel{~}{z}_1=z_1+Q(z_2,z_3)`$, where $`Q`$ is a polynomial of degree 2 with small coefficients. The set of such Q such that $`\stackrel{~}{z}_1`$ is constant on a fixed holomorphic curve is of codimension at least 2. Whence, for an open dense set of such $`Q`$, the coordinate function $`\stackrel{~}{z}_1`$ is nonconstant on each contracted curve; i.e., for all $`s\mathrm{\Delta }`$, $`\stackrel{~}{z}_1^1(s)`$ contains no contracted curves.
4. By the above argument, we can also assume that none of the $`A_s^2`$ are contained in the critical set $`C_f`$ of $`f`$.
5. By the argument below, we can also assume that for all $`s\mathrm{\Delta }`$, there do not exist nonempty disjoint open subsets $`V_1,V_2`$ of $`A_s^2`$ with $`f(V_1)=f(V_2)`$.
To show that we can realize property (e) after a change of coordinates, we let
$$U=\left[\mathrm{\Delta }\times A^2(r,1)\right](I_fC_f)$$
and we consider the set
$$D=\{(z,w)U\times U:zw,f(z)=f(w)\},$$
which is an analytic subvariety of $`U\times U`$ minus the diagonal. Note that $`D`$ is locally given as the graph of a biholomorphic map (and thus is a smooth 3-dimensional submanifold). It suffices to show that we can make a small perturbation of coordinates so that
$$dimD(\mathrm{\Delta }_s^2\times \mathrm{\Delta }_s^2)1$$
(1)
for all $`s\mathrm{\Delta }`$.
To show (1), we let $`𝒫_n^l`$ denote the vector space of polynomials of degree $`l`$ on $`^n`$. Note that
$$dim𝒫_n^l=(\begin{array}{c}l+n\\ n\end{array}).$$
(2)
We also let $`𝒥_a^l(g)𝒫_n^l`$ denote the $`l`$-jet of a germ $`g{}_{n}{}^{}𝒪_{a}^{}`$ ($`a^n`$). We shall use the following lemma:
###### Lemma 5.
Let $`\phi :\mathrm{\Delta }^m^n`$ be a holomorphic map such that $`\phi (0)=a`$ and $`\mathrm{rank}_0\phi =m`$. Then
$$\mathrm{codim}_{𝒫_n^l}\{f𝒫_n^l:𝒥_0^l(f\phi )=0\}=(\begin{array}{c}l+m\\ m\end{array}).$$
###### Proof.
By a change of coordinates, we can assume without loss of generality that $`\phi (z_1,\mathrm{},z_m)=(z_1,\mathrm{},z_m,0,\mathrm{},0)`$. The result then follows immediately from (1).∎
We let $`𝒬^l`$ denote the set of polynomials $`g`$ in $`𝒫_3^l`$ such that $`dg`$ does not vanish on $`\mathrm{\Delta }^3(2)`$. (Note that small polynomial perturbations of the coordinate function $`z_1`$ are in $`𝒬^l`$.) Let $`a=(z_0,w_0)D\mathrm{\Delta }^6(2)`$ be arbitrary, and let $`_a`$ denote the set of polynomials $`g`$ in $`𝒬^5`$ with
$$dim_a\{(z,w)D:g(z)=g(z_0),g(w)=g(w_0)\}>1.$$
We shall show that
$$\mathrm{codim}_{𝒬^5}_a6.$$
(3)
Since $`dimD=3`$, (3) implies that we can choose $`g𝒫_3^5`$ such that $`g`$ is a small deformation of the coordinate function $`z_1`$ and
$$dimD\left(g^1(s)\times g^1(t)\right)1\text{for all }s,t\mathrm{\Delta }.$$
If we then replace $`z_1`$ with $`\stackrel{~}{z}_1=g`$, (1) will be satisfied.
To verify (3), we first consider an arbitrary quadratic polynomial $`g_1𝒬^2`$, and we let
$$E=\{(z,w)D:g_1(z)=g_1(z_0)\}.$$
Since $`D`$ is locally given as a graph and $`dg_1(z_0)0`$, $`E`$ is smooth at $`z_0`$. Now let $`\phi =(\phi _1,\phi _2):\mathrm{\Delta }^2E`$ with $`\phi (0)=a`$ and $`\mathrm{rank}_0\phi =2`$. This implies that $`\mathrm{rank}_0\phi _1=\mathrm{rank}_0\phi _2=2`$. Let $`^{}(g_1)`$ denote the set of $`g_2𝒬^2`$ such that $`𝒥_0^2(g_2\phi _2)=0`$. By Lemma 5, $`\mathrm{codim}^{}(g_1)\left(\genfrac{}{}{0pt}{}{4}{2}\right)=6`$. Furthermore, we note that if we replace $`g_1`$ with a germ $`\stackrel{~}{g}_1{}_{3}{}^{}𝒪_{z_0}^{}`$ with the same 2-jet at $`z_0`$, then we can choose $`\stackrel{~}{\phi }_2`$ with the same 2-jet (at 0) as $`\phi _2`$ so that $`(\phi _1,\stackrel{~}{\phi }_2):\mathrm{\Delta }^2E`$ has the same 2-jet (at 0) as $`\phi `$. Thus if $`g_2𝒬^2^{}(g_1)`$, we have $`𝒥_0^2(g_2\stackrel{~}{\phi }_2)=𝒥_0^2(g_2\phi _2)0`$. Furthermore, if we also replace this $`g_2`$ with $`\stackrel{~}{g}_2{}_{3}{}^{}𝒪_{w_0}^{}`$ with the same 2-jet at $`w_0`$, then $`𝒥_0^2(\stackrel{~}{g}_2\stackrel{~}{\phi }_2)0`$, and hence
$$dim_a\{(z,w)D:\stackrel{~}{g}_1(z)=g_1(z_0),\stackrel{~}{g}_2(w)=g_2(w_0)\}=1.$$
Now consider the linear map
$$\tau _a:𝒫_3^5𝒫_3^2\times 𝒫_3^2,g(𝒥_{z_0}^2(g),𝒥_{w_0}^2(g)).$$
By the above discussion, $`_a\tau _a^1(_a^{})`$, where
$$_a^{}=\{(g_1,g_2):g_1𝒬^2,g_2^{}(g_1)\}.$$
Since $`\tau `$ is surjective, it follows that
$$\mathrm{codim}_{𝒬^5}_a\mathrm{codim}_{𝒬^2\times 𝒬^2}_a^{}6.$$
This completes the verification of (1), and hence conditions (a)–(e) above can be satisfied.
## 2. Proof of Proposition 3
We are now prepared to prove the Hartogs extension property. By our construction above, we may assume that the map $`f:H_1^2(r)X`$ of Proposition 3 possesses the following properties:
1. $`f`$ is non-degenerate and holomorphic on a neighborhood of $`\overline{\mathrm{\Delta }\times A^2(r,1)}`$;
2. for all $`s\mathrm{\Delta }`$, the set $`A_s^2`$ contains no curves contracted by $`f`$ to a point;
3. for all $`s\mathrm{\Delta }`$, there do not exist nonempty disjoint open subsets $`V_1,V_2`$ of $`A_s^2`$ with $`f(V_1)=f(V_2)`$;
We must show that $`f`$ extends meromorphically to $`\mathrm{\Delta }^3`$ minus a discrete set of points. Denote by $`W`$ some open subset of $`\mathrm{\Delta }`$ such that $`f`$ can be meromorphically extended onto the Hartogs domain
$$H_W(r):=\left[W\times \mathrm{\Delta }^2\right]\left[\mathrm{\Delta }\times A^2(r,1)\right].$$
Let $`\mathrm{\Omega }`$ be a strictly positive $`(2,2)`$-form on $`X`$ with $`dd^c\mathrm{\Omega }=0`$. Existence of such a form on the compact 3-dimensional manifold $`X`$ follows from the absence of nonconstant plurisubharmonic functions on $`X`$ via duality and the Hahn-Banach theorem. In fact even more is true. Every Hermitian metric on $`X`$ is conformally equivalent to a metric whose associated $`(1,1)`$-form $`\omega `$ is $`dd^c`$-closed, see \[Ga\]. In the sequel, we shall take $`\mathrm{\Omega }=\omega ^2`$, where $`\omega `$ is such a Gauduchon form. Denote by $`T`$ the pull-back of $`\mathrm{\Omega }`$ by $`f`$, i.e. $`T=f^{}\mathrm{\Omega }`$. More accurately, $`f^{}\mathrm{\Omega }`$ is defined in the case of meromorphic $`f`$ as follows. Let $`\stackrel{~}{\mathrm{\Gamma }}_f`$ denote the desingularization of the graph $`\mathrm{\Gamma }_fH_W(r)\times X`$ of $`f`$ and let $`\pi _1:\stackrel{~}{\mathrm{\Gamma }}_fH_W(r)`$ and $`\pi _2:\stackrel{~}{\mathrm{\Gamma }}_fX`$ be the natural projections. Note that $`\pi _1`$ is proper by the very definition of meromorphic map. Define
$$f^{}\mathrm{\Omega }:=\pi _1\pi _2^{}\mathrm{\Omega }.$$
(4)
The current $`T=f^{}\mathrm{\Omega }`$ is a positive bidegree $`(2,2)`$ current on $`H_W(r)`$. Being the push-forward of a smooth form (on a desingularization of $`\mathrm{\Gamma }_f`$), $`T`$ has coefficients in $`_{\mathrm{loc}}^1(H_W(r))`$.
To see that the push-forward of a smooth form $`\eta `$ by a modification $`\pi :\stackrel{~}{M}M`$ has coefficients in $`_{\mathrm{loc}}^1`$, it suffices to show that $`\pi _{}\eta `$ has no mass on the center $`C`$ of $`\pi `$. (In our case $`C=I_f`$.) But for any test form $`\phi `$ on $`M`$ and any sequence $`\rho _n\chi _C`$ with $`0\rho _n1`$, we have $`(\pi _{}\eta ,\rho _n\phi )=_{\stackrel{~}{M}}\eta \pi ^{}(\rho _n\phi )0`$. Hence $`\pi _{}\eta (E)=0`$. (In fact, this holds when $`\pi `$ is any surjective holomorphic map that is proper on $`\mathrm{Supp}\eta `$.)
It follows immediately from (4) that $`dd^cT=0`$. Moreover, $`T`$ is smooth on $`H_W(r)I_f`$, since outside the set $`I_f`$ of indeterminacy points of $`f`$, it is the usual pull-back of the smooth form $`\mathrm{\Omega }`$.
We write $`\mathrm{\Delta }_s^2:=\{s\}\times \mathrm{\Delta }^2`$, for $`s\mathrm{\Delta }`$. The function
$$\mu (s):=_{\mathrm{\Delta }_s^2}f^{}\mathrm{\Omega }$$
is well defined for all $`sW`$, since by the above, $`(f^{}\mathrm{\Omega })|_{\mathrm{\Delta }_s^2}=(f_{\mathrm{\Delta }_s^2})_{}\mathrm{\Omega }`$ is a positive, bidegree $`(2,2)`$-current on a neighborhood of $`\overline{\mathrm{\Delta }_s^2}`$ and is in $`_{\mathrm{loc}}^1`$. We remark that $`\mu (s)`$ is nothing but the volume of $`f(\mathrm{\Delta }_s^2)`$ with respect to $`\mathrm{\Omega }`$ counted with multiplicities.
###### Lemma 6.
The function $`\mu `$ is positive and smooth on $`W`$, and its Laplacian $`\mathrm{\Delta }\mu `$ smoothly extends onto the whole unit disk $`\mathrm{\Delta }`$.
###### Proof.
(We follow the method of proof of \[Iv2, Lemma 3.1\].) The positivity of $`\mu `$ follows from the positivity of $`f^{}\mathrm{\Omega }`$ and property (ii) above. To show that $`\mathrm{\Delta }\mu `$ extends to the unit disk, we begin by writing
$$T=\underset{\alpha ,\beta =1}{\overset{3}{}}t_{\alpha \overline{\beta }}dz^\stackrel{\alpha }{}d\overline{z}^\stackrel{\beta }{},$$
where $`dz^\stackrel{1}{}=dz_2dz_3,dz^\stackrel{2}{}=dz_1dz_3,dz^\stackrel{3}{}=dz_1dz_2`$. The function $`\mu `$ is then given by
$$\mu (z_1)=_{\mathrm{\Delta }^2}t_{1\overline{1}}(z_1,z_2,z_3)𝑑z_2d\overline{z}_2dz_3d\overline{z}_3.$$
Let
$$T^\epsilon =\underset{\alpha ,\beta =1}{\overset{3}{}}t_{\alpha \overline{\beta }}^\epsilon dz^\stackrel{\alpha }{}d\overline{z}^\stackrel{\beta }{}$$
be the smoothing of $`T`$ by convolution; the $`T^\epsilon `$ are smooth forms converging to $`T`$ in $`^1`$ as $`\epsilon 0`$. On $`H_W(r)I_f`$ the convergence is in the $`𝒞^{\mathrm{}}`$ topology. The functions
$$\mu ^\epsilon (z_1):=_{\mathrm{\Delta }_{z_1}^2}T^\epsilon =_{\mathrm{\Delta }_{z_1}^2}t_{1\overline{1}}^\epsilon (z_1,z_2,z_3)𝑑z_2d\overline{z}_2dz_3d\overline{z}_3$$
are smooth in $`W`$.
The condition $`dd^cT=0`$ reads as
$$\underset{\alpha ,\beta }{}\frac{^2t_{\alpha \overline{\beta }}}{z_\alpha \overline{z}_\beta }=0.$$
So,
$`\mathrm{\Delta }\mu ^\epsilon (z_1)`$ $`=`$ $`4{\displaystyle _{\mathrm{\Delta }_{z_1}^2}}{\displaystyle \frac{^2t_{1\overline{1}}^\epsilon }{z_1\overline{z}_1}}𝑑z_2dz_3d\overline{z}_2d\overline{z}_3`$
$`=`$ $`4{\displaystyle _{\mathrm{\Delta }_{z_1}^2}}{\displaystyle \underset{(\alpha ,\beta )(1,1)}{}}{\displaystyle \frac{^2t_{\alpha \overline{\beta }}^\epsilon }{z_\alpha \overline{z}_\beta }}dz_2d\overline{z}_2dz_3d\overline{z}_3`$
$`=`$ $`4{\displaystyle _{\mathrm{\Delta }_{z_1}^2}}\left[{\displaystyle \underset{\alpha =2}{\overset{3}{}}}\pm {\displaystyle \frac{t_{\alpha \overline{1}}^\epsilon }{\overline{z}_1}}dz_{5\alpha }d\overline{z}_2d\overline{z}_3+{\displaystyle \underset{\alpha =1}{\overset{3}{}}}{\displaystyle \underset{\beta =2}{\overset{3}{}}}\pm {\displaystyle \frac{t_{\alpha \overline{\beta }}^\epsilon }{z_\alpha }}dz_2dz_3d\overline{z}_{5\beta }\right].`$
Since $`f`$ is holomorphic on a neighborhood of $`\mathrm{\Delta }\times \mathrm{\Delta }^2`$, the current $`T`$ is smooth on $`\mathrm{\Delta }\times \mathrm{\Delta }^2`$ and thus $`\mathrm{\Delta }\mu ^\epsilon `$ converges smoothly to the function $`\psi `$ given by
$$\psi (z_1)=4_{\mathrm{\Delta }_{z_1}^2}\left[\underset{\alpha =2}{\overset{3}{}}\pm \frac{t_{\alpha \overline{1}}}{\overline{z}_1}dz_{5\alpha }d\overline{z}_2d\overline{z}_3+\underset{\alpha =1}{\overset{3}{}}\underset{\beta =2}{\overset{3}{}}\pm \frac{t_{\alpha \overline{\beta }}}{z_\alpha }dz_2dz_3d\overline{z}_{5\beta }\right].$$
(5)
But (5) defines a smooth function on all of $`\mathrm{\Delta }`$. While $`\mu ^\epsilon \mu `$ in $`^1`$ on $`W`$, so $`\mathrm{\Delta }\mu ^\epsilon \mathrm{\Delta }\mu `$ on $`W`$. This shows that $`\mathrm{\Delta }\mu =\psi `$ is smooth and smoothly extends onto the disk $`\mathrm{\Delta }`$. Thus $`\mu `$ is also smooth on $`W`$. ∎
###### Lemma 7.
Suppose that $`f`$ is non-degenerate and that there exists a sequence $`\{s_n\}W`$ converging to $`s_0\mathrm{\Delta }`$ such that $`\mu (s_n)`$ is bounded. Then:
1) $`f_0:=f|A_{s_0}^2`$ meromorphically extends onto $`\mathrm{\Delta }_{s_0}^2`$;
2) the volumes of the graphs $`\mathrm{\Gamma }_{f_{s_n}}`$ are uniformly bounded in $`n`$;
3) $`f`$ meromorphically extends onto $`U_0\times \mathrm{\Delta }^2`$ for some neighborhood $`U_0`$ of $`s_0`$.
###### Proof.
1) We let $`f_n=f|\mathrm{\Delta }_{s_n}^2`$ and we write $`F_n=f_n(\mathrm{\Delta }_{s_n}^2)`$. We further write $`\mathrm{\Sigma }_n=f_n(\mathrm{\Delta }_{s_n}^2),`$ $`\mathrm{\Sigma }_0=f_0(\mathrm{\Delta }_{s_0}^2)`$. Since the volumes $`\mu (s_n)`$ of the $`F_n`$ are uniformly bounded, by Bishop’s theorem (see for example \[HS\]) we can assume, after passing to a subsequence, that $`F_n`$ converges to a pure 2-dimensional analytic subset $`F`$ of $`X\mathrm{\Sigma }_0`$. Note that $`\overline{F}=F\mathrm{\Sigma }_0`$.
Case 1. $`\overline{F}`$ is a subvariety of $`X`$.
In this case
$$f_0|_{A_{s_0}^2}:A_{s_0}^2F^0$$
is a holomorphic map into an irreducible component $`F^0`$ of $`\overline{F}`$. If $`\widehat{F}^0`$ denotes a desingularization of $`F^0`$, then $`f_0`$ lifts to a meromorphic map $`\widehat{f}_0`$ from $`A_{s_0}^2`$ to $`\widehat{F}^0`$. By the 2-dimensional version of Theorem 1 proved in \[Iv3, Cor. 4(b)\], $`\widehat{f}_0`$ extends meromorphically onto $`\mathrm{\Delta }_{s_0}^2`$ minus a finite set of points. If this set is nonempty, then $`\widehat{f}_0(\mathrm{\Delta }_{s_0}^2)`$ would not be homologous to zero in $`\widehat{F}^0`$, and hence $`\mathrm{\Sigma }_0=f_0(\mathrm{\Delta }_{s_0}^2)`$ would not be homologous to zero in $`F^0`$. But $`\mathrm{\Sigma }_0=lim\mathrm{\Sigma }_n=limF_n=limF_n`$ in the sense of currents. Since $`\mathrm{Supp}limF_n\overline{F}`$, $`\mathrm{\Sigma }_0`$ is homologous to zero in $`\overline{F}`$ and therefore in $`F^0`$, a contradiction. So $`f_0`$ extends onto all of $`\mathrm{\Delta }_{s_0}^2`$.
Case 2. $`\overline{F}`$ is not a subvariety of $`X`$.
Let $`F^0`$ be the irreducible component of $`F`$ containing $`f_0(A_{s_0}^2)\mathrm{\Sigma }_0`$. Define the analytic space $`E=F_0A_{s_0}^2/`$, where the equivalence relation is defined as follows: The points $`aF^0`$ and $`bA_{s_0}^2`$ are equivalent iff $`a=f_0(b)`$, and necessarily $`b^{}A_{s_0}^2`$ is equivalent to $`b`$ iff $`f_0(b)=f_0(b^{})`$. By property (ii) above, this is a proper equivalence relation and hence $`E`$ is a complex space. Let $`\pi :E\overline{F}^0X`$ be the projection defined by $`\pi (a)=a`$ for $`aF^0`$ and $`\pi (b)=f_0(b)`$ for $`bA_{s_0}^2`$. Let $`\widehat{E}\stackrel{\eta }{}E`$ denote the normalization. By property (iii), the map $`f_0:A_{s_0}^2E`$ is generically one-to-one and thus is a normalization of its image. By the uniqueness of the normalization, $`f_0`$ lifts to a map $`\widehat{f}_0:A_{s_0}^2\widehat{E}`$, i.e., $`\eta \widehat{f}_0=f_0`$. The map $`\widehat{f}_0`$ is a biholomorphism onto its image.
The boundary $`\widehat{E}`$, being biholomorphic to $`\mathrm{\Delta }^2`$, is strictly pseudoconvex after shrinking slightly, so by Grauert’s theorem, $`\widehat{E}`$ can be blown down to a normal Stein space. This easily yields an extension of $`f_0`$ onto $`\mathrm{\Delta }_{s_0}^2`$.
2) We denote the extension of $`f_0`$ onto $`\mathrm{\Delta }_{s_0}^2`$ also by $`f_0`$. Let $`F^{}`$ be the maximal compact pure 2-dimensional variety contained in $`\overline{F}`$. (In Case 1 above, $`F^{}=\overline{F}`$, whereas in Case 2, $`F^{}=\overline{FF^0}`$.) We consider the pure two-dimensional analytic set
$$\mathrm{\Gamma }=\mathrm{\Gamma }_{f_0}(I_{f_0}\times F^{}).$$
in $`(\mathrm{\Delta }\times \mathrm{\Delta }^2)\times X`$, where $`I_{f_0}`$ is the (finite) set of points of indeterminacy of $`f_0`$.
Step 1. We claim that for all $`\epsilon >0`$ the graph $`\mathrm{\Gamma }_{f_n}`$ belongs to the $`\epsilon `$-neighborhood of $`\mathrm{\Gamma }`$, for $`n0`$.
Neighborhoods are taken with respect to the Euclidean metric on $`^3`$ and Gauduchon metric on $`X`$. (In fact, any choice of metric works as well as this one.) This claim follows immediately from Lemma 8 below.
We say that a sequence of meromorphic maps $`f_n:UX`$ converges to a holomorphic map $`f_0`$ on a domain $`U`$ if for all compact subsets $`KU`$, $`I_{f_n}K=\mathrm{}`$ for $`n0`$ and $`f_nf_0`$ uniformly on $`K`$.
###### Lemma 8.
Let $`f_n:\overline{\mathrm{\Delta }}^2X`$ be a sequence of meromorphic maps, where $`X`$ is a compact complex manifold. Suppose that $`f_n`$ is holomorphic on $`A`$, where $`A=A^2(r,1)`$. If there exists a meromorphic map $`f_0:\mathrm{\Delta }^2X`$ such that $`f_n|_Af_0|_A`$, then $`f_nf_0`$ on $`\mathrm{\Delta }^2I_{f_0}`$.
Lemma 8 is a special case of Proposition 1.1.1 in \[Iv4\]. (Proposition 1.1.1 in \[Iv4\] is stated in terms of “strong convergence” of meromorphic maps. However, if $`\{f_n\}`$ strongly converges to a holomorphic map, then the sequence converges in the above sense. This is the content of the “Rouché principle” of \[Iv4, Theorem 1\].)
To complete the proof of (2), we consider a point $`p\mathrm{\Gamma }`$ and take any open $`Wp`$ adapted to $`\mathrm{\Gamma }`$, i.e. biholomorphic to $`\mathrm{\Delta }^2\times \mathrm{\Delta }^4=U\times B`$ in such a way that $`(\overline{U}\times B)\mathrm{\Gamma }=\mathrm{}`$. Then for $`n1`$, we have $`\mathrm{\Gamma }_{f_n}(\overline{U}\times B)=\mathrm{}`$ and thus $`p_{\mathrm{\Gamma }_{f_n}}:\mathrm{\Gamma }_{f_n}(U\times B)U`$ is a $`d_n`$-sheeted analytic covering, where $`p:U\times BU`$ is a natural projection.
Step 2. The number $`\{d_n\}`$ of sheets is uniformly bounded.
Consider the following two cases. Case 1. $`p(I_{f_0}\times F^{})\mathrm{\Gamma }_{f_0}.`$ In this case as $`Wp`$ we can take the following neighborhood. Let $`p=(a,b)`$, where $`aI_{f_0}^3`$ and $`bF^{}X`$. Take a neighborhood of $`b`$ in $`X`$ of the form $`\mathrm{\Delta }^2\times \mathrm{\Delta }`$ such that $`F^{}(\overline{\mathrm{\Delta }}^2\times \mathrm{\Delta })=\mathrm{}`$. Then take some small $`\mathrm{\Delta }^3a`$ in $`^3`$ and put $`U=\mathrm{\Delta }^2`$ and $`B=\mathrm{\Delta }\times \mathrm{\Delta }^3`$. If the number $`d_n`$ of sheets of the analytic cover $`\pi _U:\mathrm{\Gamma }_{f_n}(U\times B)U`$ is not bounded, it will contradict the fact that $`f_n(\mathrm{\Delta }_{s_n}^2)(\mathrm{\Delta }^2\times \mathrm{\Delta })=F_n(\mathrm{\Delta }^2\times \mathrm{\Delta })`$ has uniformly bounded volume (counted with multiplicities).
One should remark now that boundedness of the number of sheets does not depend on the particular choice of the adapted neighborhood of $`p`$.
Case 2. $`p\mathrm{\Gamma }_{f_0}`$. Let $`W=U\times Bp`$ be some adapted neighborhood. Find a point $`qU`$ such that all its pre-images $`\{q_1,\mathrm{},q_N\}=\pi _U^1(q)\mathrm{\Gamma }`$ are smooth points of $`\mathrm{\Gamma }`$ and $`\pi _U`$ is a biholomorphism between neighborhoods $`V_jq_j`$ on $`\mathrm{\Gamma }`$ and $`V`$ on $`U`$. Denote by $`b_j`$ the projection of $`q_j`$ into $`B`$. Take mutually disjoint polydisks $`B_jB`$ with centers $`b_j`$. Consider $`W_j:=V_j\times B_j`$ as adapted neighborhoods of $`\mathrm{\Gamma }`$ in $`q_j`$. They are adapted also for $`\mathrm{\Gamma }_{f_n},n>>0.`$ Denote by $`d_n^j`$ the corresponding number of sheets. If $`d_n`$ is not bounded then at least one sequence $`d_n^j`$ is also unbounded. Fix $`j`$ with $`d_j^n`$ unbounded.
If $`q_j(I_{f_0}\times F^{})\mathrm{\Gamma }_{f_0}`$, then everything reduces to Case 1. So let $`q_j\mathrm{\Gamma }_{f_0}`$. Perturbing $`q`$ and thus $`q_j`$ if necessary, we can suppose that $`q_j`$ is a point where our map $`f`$ is holomorphic. More precisely $`q_j=(a,f(a))`$ for some $`a\mathrm{\Delta }\times \mathrm{\Delta }^2^3`$. Now the contradiction is immediate, because the graphs $`\mathrm{\Gamma }_{f_n}`$ uniformly approach $`\mathrm{\Gamma }_{f_0}`$ while $`f_n`$ converges to $`f`$ in a neighborhood of $`a`$.
3) We are exactly under the assumptions of Proposition 1.3 of \[Iv3\], i.e., we can apply the “Continuity Principle.” (The condition of boundedness of the cycle geometry is insured by Proposition 1.4 from \[Iv3\].) This gives us an extension of $`f`$ onto $`U_{s_0}\times \mathrm{\Delta }^2`$. ∎
Let us proceed further with the proof of the theorem. Let $`W`$ be the maximal open subset of the disc $`\mathrm{\Delta }`$ such that $`f`$ meromorphically extends onto $`H_W(r)`$.
###### Lemma 9.
$`\mathrm{\Delta }W`$ is a closed complete polar set in $`\mathrm{\Delta }`$.
The proof is the same as that of Lemma 2.4 from \[Iv3\] and will be omitted.
It suffices to show that there exists $`\widehat{f}:\mathrm{\Delta }^3(1\delta )X`$ satisfying the conclusion of Proposition 3 for arbitrary $`\delta >0`$. We now repeat the above arguments using two slightly deformed coordinate systems $`(z_1^{},z_2,z_3)`$ and $`(z_1^{\prime \prime },z_2,z_3)`$, where
$$z_1^{}=z_1+\epsilon z_2+O(|z|^2),z_1^{\prime \prime }=z_1+\epsilon z_3+O(|z|^2).$$
Here the $`O(|z|^2)`$ terms are chosen so that conditions (i)–(iii) at the beginning of this section are satisfied for each of the two coordinate systems, after shrinking $`r`$ if necessary. (As was shown in §1, these terms can be taken to be polynomials consisting of terms of degrees 2 through 5.) We choose $`\epsilon `$ and the $`O(|z|^2)`$ terms to be small enough so that $`\mathrm{\Delta }^3(1\delta )\mathrm{\Delta }^{}{}_{}{}^{3}\mathrm{\Delta }^3,`$ $`\mathrm{\Delta }^3(1\delta )\mathrm{\Delta }^{\prime \prime }{}_{}{}^{3}\mathrm{\Delta }^3`$, where $`\mathrm{\Delta }^{}^3`$ and $`\mathrm{\Delta }^{\prime \prime }^3`$ are the polydisks of radius $`1\frac{\delta }{2}`$ in the new coordinates.
Applying the above argument to the new coordinate systems, we obtain maximal open $`W^{},W^{\prime \prime }`$ in $`\mathrm{\Delta }\stackrel{~}{}:=\mathrm{\Delta }(1\frac{\delta }{2})`$ such that $`f`$ extends meromorphically to the Hartogs domains $`H_W^{}^{}(r),H_{W^{\prime \prime }}^{\prime \prime }(r)`$. We let $`S_1=\mathrm{\Delta }\stackrel{~}{}W,S_2=\mathrm{\Delta }\stackrel{~}{}W^{},S_3=\mathrm{\Delta }\stackrel{~}{}W^{\prime \prime }`$. Now consider the coordinates
$$w_1=z_1,w_2=z_1^{},w_3=z_1^{\prime \prime }$$
and let $`U`$ denote the image of $`\mathrm{\Delta }^3(1\delta )`$ under the coordinate map $`(w_1,w_2,w_3)`$ . (We may assume that $`z_1^{},z_1^{\prime \prime }`$ are chosen so that the $`w_j`$ indeed provide coordinates on $`\mathrm{\Delta }^3`$.)
In terms of the $`w`$-coordinates, $`f`$ then extends to a meromorphic map $`\widehat{f}`$ on $`U(S_1\times S_2\times S_3)`$. Now let $`s_0`$ be an arbitrary point in $`S:=S_1\times S_2\times S_3`$. We must show that $`s_0`$ is an isolated point of $`S`$ and that $`\widehat{f}(B_{s_0}(r))`$ is not homologous to zero in $`X`$, for any ball $`B_{s_0}(r)`$ centered at $`s_0`$ such that $`B_{s_0}(r)S=\{s_0\}`$.
Since polar sets in $``$ are of Hausdorff dimension zero, we can choose a polydisk $`\mathrm{\Delta }_0^3`$ about $`s_0`$ such that the set $`K:=S\mathrm{\Delta }_0^3`$ is compact. An identical proof to that of Lemma 3.3 from \[Iv2\] now shows that the current $`T=f^{}\mathrm{\Omega }`$ has locally summable coefficients on all of $`\mathrm{\Delta }_0^3`$. Hence $`T`$ extends to a unique current $`\stackrel{~}{T}`$ on $`\mathrm{\Delta }_0^3`$ with $`_{\mathrm{loc}}^1`$ coefficients. The following lemma then tells us that $`dd^c\stackrel{~}{T}`$ is of order 0:
###### Lemma 10.
\[Iv2, Proposition 2.3\] Let $`K`$ be a complete pluripolar, compact set in a strictly pseudoconvex domain $`D^n`$ and $`T`$ a positive, bidimension $`(1,1)`$ current in $`DK`$. Suppose that:
1) $`dd^cT0`$ in $`DK`$,
2) $`T`$ has locally finite mass in a neighborhood of $`K`$,
3) $`dT`$ and $`d^cT`$ have measure coefficients on $`DK`$.
Then the current $`dd^c\stackrel{~}{T}`$ has measure coefficients in $`D`$.
(Condition (3) on $`dT`$ and $`d^cT`$ was omitted in \[Iv2\], but is used in the proof. In our case $`T=\pi _1\pi _2^{}\mathrm{\Omega }`$, so this condition follows from the fact that $`dT=\pi _1\pi _2^{}d\mathrm{\Omega }`$ is the push-forward of a smooth form by a proper map, and similarly for $`d^cT`$.)
Since $`dd^cT=0`$, the support of the current $`dd^c\stackrel{~}{T}`$ must be contained in $`K`$. We also conclude from the Lemma 2.6 in \[Iv3\] that $`dd^c\stackrel{~}{T}0`$. Thus we can write $`dd^c\stackrel{~}{T}=\nu \omega _e^3`$, where $`\omega _e`$ is the Euclidean Kähler form on $`^3`$. Then for any ball $`B_{s_0}(r)\mathrm{\Delta }_0^3`$ about $`s_0`$ with $`B_{s_0}(r)K=\mathrm{}`$, we have that either
$$\nu (KB_{s_0}(r))=0,$$
(6)
or
$$\begin{array}{ccc}0\hfill & >& \nu (KB_{s_0}(r))=_{B_{s_0}(r)}𝑑d^c\stackrel{~}{T}=\underset{\epsilon 0}{lim}_{B_{s_0}(r)}𝑑d^c\stackrel{~}{T}_\epsilon \hfill \\ & =& \underset{\epsilon 0}{lim}_{B_{s_0}(r)}d^c\stackrel{~}{T}_\epsilon =_{B_{s_0}(r)}d^cT_\epsilon =_{f(B_{s_0}(r))}d^c\mathrm{\Omega }.\hfill \end{array}$$
(7)
Case 1. $`\nu (KB_{s_0}(r))=0.`$
In this case the negativity of $`\stackrel{~}{T}`$ implies that $`dd^c\stackrel{~}{T}=0`$. Therefore we can find a polydisk neighborhood $`\mathrm{\Delta }^3B_{s_0}(r)`$ of $`s_0`$ and a $`(2,1)`$-form $`\mathrm{\Gamma }`$ in $`\mathrm{\Delta }^3`$ such that:
1) $`f`$ is holomorphic in a neighborhood of $`\mathrm{\Delta }\times \mathrm{\Delta }^2`$ (and therefore $`\stackrel{~}{T}`$ is smooth there);
2) $`\stackrel{~}{T}=i(\overline{\mathrm{\Gamma }}\overline{}\mathrm{\Gamma })`$ in a neighborhood of $`\overline{\mathrm{\Delta }}^3`$;
3) $`\mathrm{\Gamma }`$ is smooth in a neighborhood of $`\mathrm{\Delta }\times \mathrm{\Delta }^2`$.
The zero-dimensionality of $`K`$ implies that there exists a nonempty open $`W\mathrm{\Delta }`$ such that $`f`$ is defined and meromorphic on $`W\times \mathrm{\Delta }^2`$ and that $`s_0^1W\mathrm{\Delta }`$, where $`s_0^1`$ is the first coordinate of $`s_0`$. As before we let
$$\mu (z_1)=_{\mathrm{\Delta }_{z_1}^2}\stackrel{~}{T}=i_{\mathrm{\Delta }_{z_1}^2}(\overline{\mathrm{\Gamma }}\mathrm{\Gamma }).$$
By the smoothness of $`\mathrm{\Gamma }`$, the function $`\mu `$ is bounded. Therefore by Lemma 7, $`f`$ extends meromorphically to a neighborhood of $`s_0`$.
Case 2. $`\nu (SB_{s_0}(r))<0.`$
By (7), the 5-cycle $`f(B_{s_0}(r))`$ is not homologous to zero in $`X`$. Furthermore, $`_{f(B_{s_0}(r))}d^c\mathrm{\Omega }`$ depends only on the integer homology class of $`f(B_{s_0}(r))`$, since $`dd^c\mathrm{\Omega }=0`$. Hence,
$$_{f(B_{s_0}(r))}d^c\mathrm{\Omega }\delta <0,$$
where $`\delta `$ is independent of $`s_0`$ and $`r`$ (and depends only on $`X`$ and $`\mathrm{\Omega }`$). This shows that $`K`$ is finite, and completes the proof. ∎
Remark that our proof gives more. Namely, if $`\mathrm{\Sigma }M\{a_1,\mathrm{},a_d\}`$ is not homologous to zero in $`M\{a_1,\mathrm{},a_d\}`$ then $`f(\mathrm{\Sigma })`$ is not homologous to zero in $`X`$.
## 3. Generalizations and open questions
In \[Iv3\], the classes $`𝒫_k^{}`$ and $`𝒢_k`$ of complex spaces were introduced. Recall that $`𝒫_k^{}`$ is the class of normal complex spaces which carry a strictly positive $`(k,k)`$-form $`\mathrm{\Omega }^{k,k}`$ with $`dd^c\mathrm{\Omega }^{k,k}0`$, and $`𝒢_k`$ is the subclass of $`𝒫_k^{}`$ which consists of complex spaces carrying a strictly positive $`(k,k)`$-form $`\mathrm{\Omega }^{k,k}`$ with $`dd^c\mathrm{\Omega }^{k,k}=0`$. Note that $`𝒢_k`$ contains all compact complex manifolds of dimension $`k+1`$.
It is easy to observe that our above proof gives the following more general statement of Proposition 3:
###### Proposition 11.
Let $`X`$ be a compact complex manifold in the class $`𝒫_2^{}`$. Then every meromorphic map $`f:H_1^2(r)X`$ extends meromorphically onto $`\mathrm{\Delta }^3A`$, where $`A`$ is a closed, complete pluripolar subset of Hausdorff dimension zero. If moreover $`X𝒢_2`$ then $`A`$ is discrete and for every ball $`B`$ with center $`aA`$ such that $`BA=\mathrm{}`$, $`f(B)`$ is not homologous to zero in $`X`$.
To consider the extension of mappings from higher dimensional domains, we introduce the Hartogs figures
$$H_d^k(r):=\left[\mathrm{\Delta }^d(1r)\times \mathrm{\Delta }^k\right]\left[\mathrm{\Delta }^d\times A^k(r)\right]^{d+k}.$$
We conjecture that the analogous result should hold for meromorphic mappings from $`H_d^k(r)`$ to compact manifolds (and spaces) in the classes $`𝒫_k^{}`$ and $`𝒢_k`$. In particular, Theorem 1 should be true for meromorphic mappings between equidimensional manifolds in all dimensions. The main difficulty lies in the fact that it is impossible in general to make the reductions (a)–(c) of §1. (Note that reductions (d)–(e) can be achieved in all dimensions.) However, these reductions are unnecessary in the case when our map is locally biholomorphic, as we state below.
###### Proposition 12.
Let $`X`$ be a compact complex space of dimension $`k+1`$. Then every holomorphic map $`f:H_1^k(r)X`$ with zero-dimensional fibers extends meromorphically onto $`\mathrm{\Delta }^{k+1}A`$, where $`A`$ is discrete, and for every ball $`B`$ with center $`aA`$ such that $`BA=\mathrm{}`$, $`f(B)`$ is not homologous to zero in $`X`$.
The proof is by induction on the dimension $`n=k+1`$. For the inductive step, the function $`\mu `$ is defined in terms of the push-forward of a $`dd^c`$-closed, positive $`(k,k)`$-form $`\mathrm{\Omega }`$ on a desingularization of $`X`$.
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# Inflation models, spectral index and observational constraints.
## 1 Introduction
It is generally supposed that structure in the Universe originates from a primordial gaussian curvature perturbation, generated by the quantum fluctuations of the inflaton field during slow-roll inflation. Then the spectrum of the curvature perturbation $`\delta _H(k)`$ is determined by the inflaton potential $`V(\varphi )`$. In this paper we will consider the scale–dependence of the primordial spectrum, defined by the spectral index $`n`$:
$$n(k)12\frac{\mathrm{ln}\delta _H}{\mathrm{ln}k}=2M_{Pl}^2(V^{\prime \prime }/V)3M_{Pl}^2(V^{}/V)^2,$$
(1)
where<sup>a</sup><sup>a</sup>a$`M_{Pl}=2.4\times 10^{18}\text{GeV}`$ is the Planck mass, $`a`$ is the scale factor and $`H=\dot{a}/a`$ is the Hubble parameter, and $`k/a`$ is the wavenumber. the potential and its derivatives are evaluated at the epoch of horizon exit $`k=aH`$. The value of $`\varphi `$ at this epoch is given by $`\mathrm{ln}(k_{end}/k)=N(\varphi )=M_{Pl}^2_{\varphi _{end}}^\varphi (V/V^{})𝑑\varphi `$ , where $`k_{end}`$ is the scale leaving the horizon at the end of slow roll inflation and $`N(\varphi )`$ is the number of e-folds. In the majority of the inflation models, $`n`$ is practically scale-independent so that $`\delta _H^2k^{n1}`$, but we shall also discuss an interesting class of models giving significant scale dependence.
## 2 The observational constraints on the $`\mathrm{\Lambda }`$CDM model
In the interest of simplicity and due to present observations $`^\mathrm{?}`$, we adopt the $`\mathrm{\Lambda }`$CDM model, with $`\mathrm{\Omega }_{tot}=1`$ and cold non-baryonic dark matter with negligible interaction. We shall constrain this model, including the spectral index, by performing a least-squares fit to the key observational quantities.
The parameters of the $`\mathrm{\Lambda }`$CDM model are the primordial spectrum $`\delta _H(k)`$, the Hubble constant $`h`$ (in units of $`100\text{km}\text{s}^1\text{Mpc}^1`$), the total matter density $`\mathrm{\Omega }_0`$, the baryon density $`\mathrm{\Omega }_b`$, and the reionization redshift $`z_R`$ (we consider complete and sudden reionization). $`z_R`$ can be estimated in terms of the other parameters because it can be related to the density perturbation and the the fraction of collapsed matter $`f`$ at the epoch of reionization, so we exclude it from the least-squares fit. In the case of the constant $`n`$ models we fix it at a reasonable value ($`z_R=20`$), while in the case of the running mass models we compute it assuming that reionization occurs when a fixed fraction of the matter $`f1`$ collapses. The spectrum is conveniently specified by its value at the COBE scale $`k_{COBE}=6.6H_0`$, and the spectral index $`n(k)`$. Excluding $`z_R`$, the $`\mathrm{\Lambda }`$CDM model is therefore specified by five parameters in the case of a constant spectral index, or by six parameters in the case of running mass inflation models.
Taking as our starting point a study performed three years ago $`^\mathrm{?}`$, we consider seven observational quantities: the cosmological quantities $`h`$, $`\mathrm{\Omega }_0`$, $`\mathrm{\Omega }_B`$, which we are also taking as free parameters, and moreover the shape parameter $`\mathrm{\Gamma }`$, $`\sigma _8`$, the COBE normalization and the first peak height in the cmb anisotropy. The adopted values and errors are given in the second and third line of Table 1. For a discussion of the data, see $`^\mathrm{?}`$. In common with earlier investigations, we assume the errors to be uncorrelated and random errors to dominate over systematic ones.
## 3 Results
We perform the least–squares fit with the CERN Minuit package; the quoted error bars use the parabolic approximation, while the exact errors computed by Minuit agree with the approximated ones to better than 10%.
In order to simplify the numerical procedure, we follow $`^\mathrm{?}`$ and parameterize the predicted value of $`\sqrt{\stackrel{~}{C}_{\mathrm{peak}}}`$ with the analytical formula $`\sqrt{\stackrel{~}{C}_{\mathrm{peak}}}=77.5\mu \text{K}\left(\frac{\delta _H(k_{COBE})}{1.94\times 10^5}\right)\left(\frac{220}{10}\right)^{\nu /2}`$ where
$$\nu 0.88(n_{COBE}1)0.37\mathrm{ln}(h/0.65)0.16\mathrm{ln}(\mathrm{\Omega }_0/0.35)+5.4h^2(\mathrm{\Omega }_b0.019)0.65g(\tau )\tau $$
(2)
and $`\tau =0.035\frac{\mathrm{\Omega }_b}{\mathrm{\Omega }_0}h\left(\sqrt{\mathrm{\Omega }_0(1+z_R)^3+1\mathrm{\Omega }_0}1\right)`$. The formula is fitted to the CMBfast$`^\mathrm{?}`$ results and agrees within 10% for a 1-$`\sigma `$ variation of the cosmological parameters, $`h,\mathrm{\Omega }_0`$ and $`\mathrm{\Omega }_b`$, and $`n=1.0\pm 0.05`$. With the function $`g(\tau )`$ set equal to 1, the formula contains the usual factor $`\mathrm{exp}(\tau )`$. By fitting the output of CMBfast, we introduce also $`g(\tau )=10.165\tau /(0.4+\tau )`$, so that our formula is accurate to a few percent over the interesting range of $`\tau `$.
Constant spectral index. For the case of a constant spectral index our result is given in Table 1 for $`z_R=20`$. In the case of no reionization ($`z_R=0`$) we obtain a slightly smaller spectral index, $`n=0.98\pm .05`$, and cosmological parameters within the observational error bar, in agreement with previous analysis $`^\mathrm{?}`$. This result is not enough yet to exclude completely proposed inflationary models, but a better determination of the peak height could strengthen the bound sufficiently to discriminate between them, especially in the case of new inflation models, which give low values of $`n`$ $`^\mathrm{?}`$.
Running mass models. We have also considered the scale-dependent spectral index, predicted in inflation models with a running inflaton mass $`^\mathrm{?}`$. In these models, one–loop corrections to the potential are taken into account by evaluating the scale dependent inflaton mass $`m^2(Q)`$ at $`Q\varphi `$. Then the spectral index can be parameterized by just two quantities:
$$\frac{n(k)1}{2}=\sigma e^{cN(\varphi )}c,$$
(3)
where $`\sigma `$ is an integration constant and $`c`$ is related to the inflaton coupling responsible of the mass running. The different signs of $`\sigma `$ and $`c`$ give raise to four different models of inflation. In general, without fine tuning, we expect
$`|c|\text{ }<|\sigma |\text{ }<1`$ $`|c|g^2{\displaystyle \frac{\stackrel{~}{m}^2M_{Pl}}{V_0}}`$ (4)
with $`g`$ denoting the gauge or Yukawa coupling of the inflaton, $`\stackrel{~}{m}^2`$ the soft supersymmetry breaking mass of the particles in the loop and $`V_0`$ the value of the potential energy during inflation. With gravity-mediated susy breaking, typical values of the masses are $`\stackrel{~}{m}^2V_0/M_{Pl}^2`$, which makes $`c`$ of order of the coupling strength. For a gauge coupling, or an unsuppressed Yukawa coupling, we expect $`|c|10^1\text{ to }10^2`$.
Extremizing with respect to all other parameters, we have computed the $`\chi ^2`$ in the $`\sigma `$ vs. $`c`$ plane and obtained contour levels for $`\chi ^2`$ corresponding to the 70% and 95% confidence level in two variables. The results are shown in Figure 1.
In the case of Models (ii) and (iv), the allowed region corresponds to $`|c|`$ and $`|\sigma |`$ both small, giving a practically scale-independent spectral index, with a red and blue spectrum respectively.
In contrast, the allowed region for Models (i) and (iii) allows strong scale-dependence. In Model (i), a large departure from a constant spectral index is allowed for large $`\sigma `$; for the theoretically favored value $`\sigma 1`$ the variation between $`k_{COBE}`$ and $`8^1h\text{Mpc}^1`$ can be as large as 0.05, while the maximal change allowed by the data is 0.2. For Model (iii), a much larger departure from a constant spectral index is allowed, but in the theoretically favored regime $`|\sigma |c`$ one again finds a variation of at most 0.05.
## 4 Conclusion
We have evaluated the observational constraints on the spectral index $`n`$, using a range of data, and we find, for constant $`n`$ at 2-$`\sigma `$ level, $`0.88n1.11`$ for $`0z_R20`$.
We have also investigated the running mass models, parameterized by $`c`$ and $`\sigma `$. For $`c`$ and $`\sigma `$ with the same sign, we have found that indeed $`n`$ can vary by about $`0.05`$ between the COBE scale and $`8h^1\text{Mpc}`$. Moreover, if $`c`$ is positive as it would be for a gauge coupling, $`n1`$ can change sign between the COBE and $`8h^1\text{Mpc}`$ scales. It will be very interesting to see how the present situation changes with the advent of better data.
## Acknowledgments
It is a pleasure to thank D. H. Lyth with whom this work has been done. I would also like to thank the organizers of Moriond 2000 and the European Union for the financial support.
## References
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# Exact solution for the critical state in thin superconductor strips with field dependent or anisotropic pinning
## I Introduction
The critical state model for the magnetic behavior of superconductors with flux-line pinning has proven very useful though it originally was applied to the simple (demagnetization-free) longitudinal geometry of long superconductors in parallel magnetic field. It took over 30 years until an analytical solution of the critical state model was obtained for the more realistic transverse geometry of thin superconductors. The solutions were derived for thin disks and strips in a perpendicular magnetic field, extending an earlier work on superconductor strips with transport current, and finally for elliptic-shaped platelets. Recent detailed numerical work for strips and disks of finite thickness shows how the transition from longitudinal to transverse geometry occurs with changing aspect ratio of the specimen.
So far, in the transverse geometry all analytical solutions of the critical state model were restricted to the Bean model of constant critical current density $`j_c=`$ const, but in many experiments $`j_c=j_c(B)`$ depends on the local magnetic induction $`B`$. For example, the simple Kim model $`j_c(B)=j_c(0)/(1+|B|/B_0)`$ was considered in many experimental and theoretical papers, see e.g. the reviews and the partly analytical calculations for thin strips and disks. While numerical computations easily allow us to consider any $`j_c(B)`$ dependence, an exact analytical solution of some model may give deeper insight since it yields explicit dependences of the resulting quantities on the input parameters.
In the highly anisotropic high-$`T_c`$ superconductors the flux-line pinning in general depends on the angle $`\theta `$ between the local direction of the magnetic induction $`𝐁`$ and the $`c`$ axis, which in typical experiments is normal to the plane of the sample. For example, this type of anisotropy occurs when one takes into account the intrinsic pinning exerted by the CuO planes or the pinning by extended defects. It has been shown recently that for thin superconductors of any shape (with thickness $`d`$ much smaller than the lateral extension $`L`$ but larger than the magnetic penetration depth $`\lambda `$) any such out-of-plane-anisotropy of pinning is equivalent to an induction dependence of the critical sheet current $`J_c(B)`$ (the sheet current is defined as the current density integrated over the film thickness). Thus, the description of the two-dimensional critical state, e.g., in an anisotropic strip can be reduced to the analysis of a one-dimensional problem with some $`J_c(B)`$. In this case the characteristic scale $`B_0`$ over which $`J_c(B)`$ changes is of the order of $`\mu _0j_cd`$.
In this paper we present an analytical solution for the critical state in thin superconductor strips in perpendicular field with field dependent critical current density $`j_c(B)`$ or, equivalently, with anisotropic pinning described by a $`j_c(\theta )`$. A three-parameter model $`j_c(B)`$ consisting of two straight lines, an inclined line at small $`B`$ and a horizontal line at larger $`B`$, is considered. This rather general model is equivalent to a piecewise constant angular dependence $`j_c(\theta )=j_{c1}`$ for $`0\theta <\theta _0`$ and $`j_c(\theta )=j_{c2}`$ for $`\theta _0\theta <\pi /2`$ where $`j_{c1}`$, $`j_{c2}`$, and $`\theta _0`$ are the parameters of the model. We shall show below that the steepness of the flux front in the superconductor essentially depends on the anisotropy of pinning. In particular, in the case corresponding to the intrinsic pinning in high-$`T_c`$ superconductors, the front is a very sharp step, which should be taken into account in analyzing data of local magnetic measurements. We shall also show that under certain conditions two penetrating flux fronts can occur in an anisotropic superconductor.
As usual, we consider here the cases when the characteristic magnetic field in the sample is sufficiently large such that the difference between the magnetic induction $`B`$ and the field $`H`$ may be disregarded. This condition is satisfied when $`j_cd`$ is much larger than the lower critical field $`H_{c1}`$ (otherwise, the so-called geometric barrier must be taken into account). We shall thus express all the following equations in terms of the magnetic field $`H`$, related to the current density by the Maxwell equation $`𝐣=\times 𝐇`$.
## II Model and its solution
We consider an infinitely long strip of width $`2w`$ and thickness $`d`$, filling the space $`wxw`$, $`d/2zd/2`$, i.e. we place the $`y`$ axis of the coordinate system along the central line of the strip and the $`z`$ axis along the external magnetic field $`H_a`$ which is applied normal to the plane of the strip. The increasing applied field induces a sheet current $`J`$ along $`y`$, which is related to the $`z`$ component of the magnetic field in the plane $`z=0`$ by the Biot-Savart law,
$$H_z(x)=H_a+\frac{1}{2\pi }_w^w\frac{J(t)dt}{tx}.$$
(1)
Here and below all singular integrals are taken in the sense of the Cauchy principal value. The penetration of the magnetic flux into the superconducting strip is described by the following critical state equations: In the flux-free central region $`|x|b(H_a)`$ one has
$$H_z=0,$$
(2)
while in the region $`b(H_a)|x|w`$, where the flux already exists, one has
$$|J(x)|=J_c[H_z(x)].$$
(3)
The position $`x=b(H_a)`$ of the boundary separating the regions, is found by solving these equations. In Eq. (3) $`J_c(H_z)`$ is the critical value of the sheet current. At present an exact solution of Eqs. (1)–(3) is known only for the Bean critical state model where $`J_c=`$ const. Below we shall obtain the exact solution for the more general case when $`J_c(|H_z|)`$ has the model form (see Fig. 1):
$`J_c(H_z)`$ $`=`$ $`J_{c1}\gamma H_z\mathrm{for}0H_zH_z^0,`$ (4)
$`J_c(H_z)`$ $`=`$ $`J_{c0}\mathrm{for}H_zH_z^0.`$ (5)
Here $`\gamma =(J_{c1}J_{c0})/H_z^0`$; the three parameters $`J_{c1}`$, $`J_{c0}`$, and $`H_z^0`$ may have any positive value.
As was mentioned above, in the case of thin superconductors the dependence of the critical current density $`j_c`$ on the angle $`\theta `$ between the local direction of the magnetic induction and the normal to the strip plane can be taken into account if one considers this superconductor as infinitely thin but with an $`H_z`$ dependent sheet current. The model dependence described by Eqs. (4) corresponds to the following $`\theta `$-dependence of the critical current density shown in Fig. 1:
$`j_c(\theta )`$ $`=`$ $`J_{c0}/d\mathrm{for}0\theta \theta _0,`$ (6)
$`j_c(\theta )`$ $`=`$ $`J_{c1}/d\mathrm{for}\theta _0\theta \pi /2,`$ (7)
where $`\mathrm{tan}\theta _0=J_{c0}/2H_z^0`$. Thus, the case $`\gamma >0`$ describes intrinsic pinning by the CuO planes in high-$`T_c`$ superconductors ($`j_c`$ peaks at $`\theta =\pi /2`$), whereas the case $`\gamma <0`$ can be used to analyze pinning by columnar defects normal to the film ($`j_c`$ peaks at $`\theta =0`$). In both these cases one can find two-dimensional solutions of the critical state equations for strips of small but finite thickness using the results obtained below and Eqs. (5,6,9–11) of Ref. .
Accounting for the symmetry of the sheet current, $`J(x)=J(x)`$, we seek the solution of Eqs. (1)–(4) in the form
$$J(x)=\frac{x}{|x|}[J_0(x)+J_1(x)]$$
(8)
where
$`J_0(x)=J_{c0},b^2x^2w^2,`$ (9)
$`J_0(x)={\displaystyle \frac{2J_{c0}}{\pi }}\mathrm{arctan}\left[{\displaystyle \frac{(w^2b^2)x^2}{w^2(b^2x^2)}}\right]^{1/2},x^2b^2,`$ (10)
while $`J_1(x)`$ is a new unknown function. The parameter $`b`$ defines the position of the flux front, i.e., $`x=b`$ is the point where $`H_z`$ goes to zero. This parameter depends on $`H_a`$ and must be determined together with $`J_1(x)`$. Both $`J_0(x)`$ and $`J_1(x)`$ (and the magnetic field below) are even functions, which depend only on $`x^2`$. The function $`J_0(x)`$ has the form of the exact solution to Eqs. (1)–(3) in the case when $`J_c=J_{c0}`$ and the external magnetic field is equal to
$`H_b=H_{cs}\mathrm{arccosh}(w/b)`$
where $`H_{cs}=J_{c0}/\pi `$. Using Eqs. (1), (6)–(8), the expression for the magnetic field can be rewritten as
$$H_z(x)=H_0(x)\frac{1}{2\pi }_0^{a^2}\frac{J_1(\sqrt{s})ds}{sx^2},$$
(11)
where $`a`$ is defined by the equality $`H_z(a)=H_z^0`$, and $`H_0(x)`$ is the sum of $`H_a`$ and the field generated by the current $`J_0(x)`$,
$`H_0(x)=H_aH_b,0x^2b^2,`$ (12)
$`H_0(x)=H_aH_b+`$ (13)
$`H_{cs}\mathrm{arctanh}\left[{\displaystyle \frac{(x^2b^2)w^2}{x^2(w^2b^2)}}\right]^{1/2},b^2x^2w^2.`$ (14)
In Eq. (9) it was taken into account that $`J_1(x)`$ differs from zero only in the region $`0x^2a^2`$ where $`H_z(x)<H_z^0`$.
With the above formulas, the critical state equations take the following form: In the interval $`0x^2b^2`$ one has
$$H_0(x)=\frac{1}{2\pi }_0^{a^2}\frac{J_1(\sqrt{s})ds}{sx^2},$$
(15)
and in the region $`b^2x^2a^2`$ we arrive at
$$H_0(x)H_z^0=\frac{J_1(x)}{\gamma }+\frac{1}{2\pi }_0^{a^2}\frac{J_1(\sqrt{s})ds}{sx^2}.$$
(16)
In deriving Eq. (13) we have expressed $`H_z(x)`$ for $`b^2x^2a^2`$ in terms of $`J_1(x)`$ using the equality
$$H_z(x)=H_z^0\frac{J_1(x)}{\gamma }$$
(17)
that follows from formulas (3), (4), (6), (7). Eqs. (12), (13) are linear singular integral equations with Cauchy type kernel. The theory of such equations is well elaborated, and hence we can find $`a`$, $`b`$, and $`J_1(x)`$ for any given $`H_a`$.
To do this, we introduce the following notations:
$`\alpha {\displaystyle \frac{1}{\pi }}\mathrm{arctan}{\displaystyle \frac{\gamma }{2}},\beta {\displaystyle \frac{1}{2}}\alpha ,`$ (18)
$`\alpha _+\alpha ,\alpha _{}\alpha +1,`$ (19)
$`F_\pm (t)(a^2t^2)^{\alpha _\pm }|t^2b^2|^\beta `$ (20)
and define the function $`f(t)`$ by the equalities
$`f(t)=2H_0(t),0t<b,`$ (21)
$`f(t)=2\mathrm{sin}\pi \alpha [H_z^0H_0(t)],b<ta,`$ (22)
i.e., $`f(t)`$ is discontinuous at $`t=b`$. Then, the solution of Eqs. (12), (13) can be represented as follows: In the interval $`0x^2b^2`$ one has
$$J_1(x)=\frac{2}{\pi }|x|F_\pm (x)_0^a\frac{f(t)dt}{(t^2x^2)F_\pm (t)},$$
(23)
while in the interval $`b^2x^2a^2`$ we arrive at
$`J_1(x)=\mathrm{cos}\pi \alpha [f(x)+{\displaystyle \frac{\gamma }{\pi }}|x|F_\pm (x){\displaystyle _0^a}{\displaystyle \frac{f(t)dt}{(t^2x^2)F_\pm (t)}}],`$ (24)
and $`J_1(x)=0`$ for $`a^2x^2w^2`$. Here the integrals are taken in the sense of the Cauchy principal value; $`F_+`$ and $`F_{}`$ refer to positive and negative values of $`\gamma `$, respectively. If $`\gamma <0`$, for the above solution to exist it is necessary that
$`{\displaystyle _0^a}{\displaystyle \frac{f(t)}{F_{}(t)}}𝑑t=0,`$ (25)
and
$`{\displaystyle _0^a}{\displaystyle \frac{t^2f(t)}{F_{}(t)}}𝑑t=0.`$ (26)
These two equalities enable us to determine $`b`$ and $`a`$ when $`\gamma <0`$. If $`\gamma >0`$, the necessary condition for the existence of the solution is
$$_0^a\frac{f(t)}{F_+(t)}𝑑t=0.$$
(27)
A second relation between $`a`$ and $`b`$ in this case is obtained from the analysis of the magnetic field near the point $`x^2=a^2`$. It turns out that
$$H_z(x)H_z^0C_\pm \frac{\gamma }{2|\gamma |}(4+\gamma ^2)^{1/2}(x^2a^2)^{\alpha _\pm }$$
(28)
if $`x^2`$ tends to $`a^2`$ from above, and
$$H_z(x)H_z^0C_\pm (a^2x^2)^{\alpha _\pm }$$
(29)
if $`x^2`$ approaches $`a^2`$ from below. Here $`C_\pm `$ are certain integrals independent of $`x`$; the subscripts $`+`$ and $``$ refer to the cases of positive and negative $`\gamma `$, respectively. Since $`H_z(x)H_z^0`$ when $`x^2>a^2`$, we find that $`C_+0`$. On the other hand, one has $`H_z(x)H_z^0`$ when $`x^2<a^2`$, and thus $`C_\pm 0`$. Hence, one concludes that $`C_+=0`$. This is the second equality in the case of positive $`\gamma `$, and it has the form
$`{\displaystyle _0^b}{\displaystyle \frac{f(t)dt}{(a^2t^2)F_+(t)}}{\displaystyle \frac{f(a)}{2\alpha a(a^2b^2)^{1/2}}}+`$ (30)
$`{\displaystyle _b^a}\left[{\displaystyle \frac{f(t)}{t(t^2b^2)^\beta }}{\displaystyle \frac{f(a)}{a(a^2b^2)^\beta }}\right]{\displaystyle \frac{tdt}{(a^2t^2)^{1+\alpha }}}=0.`$ (31)
The Eqs. (19) and (22) determine $`a`$ and $`b`$ when $`\gamma >0`$.
## III Analysis
Let us now analyze the obtained solution. For evaluation of the integrals in Eqs. (9), (15)-(19), (22) we use the method given in Appendix A. Some profiles $`J(x)`$, Eq. (6), and $`H_z(x)`$, Eq. (9), obtained in this way are shown in Figs. 2 to 5.
It should be noted that no restriction on $`C_{}`$ is obtained when $`\gamma <0`$. In this situation the constant $`C_{}`$ is not equal to zero but negative, and thus the derivative of $`H_z`$ with respect to $`x`$ becomes infinite at $`x=a`$. In the same point a sharp bend occurs in $`J(x)`$. In other words, we obtain that two flux fronts exist in the sample, at $`x=b`$ and at $`x=a`$, see Figs. 2, 4, and 5. Of course, the singularities in $`H_z`$ and in $`J`$ at $`x=a`$ result from the sharp bend in our model $`J_c(H_z)`$ at $`H_z=H_z^0`$, see Eqs. (4) and Fig. 1. However, one may expect that in the case $`\gamma <0`$ our qualitative conclusion on the existence of the second flux front in the sample remains valid if $`J_c(H_z)`$ is a smooth function but its behavior changes abruptly over an interval smaller than $`H_{cs}`$. Such changes indeed may occur if the critical current density has sufficiently sharp angular dependence $`j_c(\theta )`$.
We shall now describe $`H_z(x)`$ and $`J(x)`$ in the vicinity of the point $`x=b`$ in which $`H_z=0`$. According to Eq. (14), at this point $`|J(b)|=J_{c1}`$. When $`x^2b^2`$, it follows from the exact solution that
$$|J(x)|J_{c1}C_b^\pm (b^2x^2)^\beta ,$$
(32)
while if $`x^2b^2`$, one has
$$|J(x)|J_{c1}\frac{\gamma }{(4+\gamma ^2)^{1/2}}C_b^\pm (x^2b^2)^\beta .$$
(33)
Here $`C_b^+`$ and $`C_b^{}`$ are certain integrals which do not depend on $`x`$ and have negative values. Formulas (23) and (24) show that in the case $`\gamma >0`$, $`|J(x)|`$ has a sharp peak at $`x=b`$, whereas for $`\gamma <0`$, $`J(x)`$ is a monotonic function and its derivative with respect to $`x`$ becomes infinite at $`x=b`$, see Figs. 2, 4, and 5. Taking into account the above formulas and Eq. (14), one obtains the distribution of $`H_z`$ near $`x=b`$,
$`H_z`$ $`=`$ $`0\mathrm{for}xb,`$ (34)
$`H_z`$ $`=`$ $`{\displaystyle \frac{C_b^\pm }{(4+\gamma ^2)^{1/2}}}(x^2b^2)^\beta \mathrm{for}xb.`$ (35)
When $`\gamma =0`$, we arrive at the well-known result $`H_z(x^2b^2)^{1/2}`$. However, in the general case, taking into account the equality $`\beta =\frac{1}{2}\frac{1}{\pi }\mathrm{arctan}(\gamma /2)`$, one may conclude that the greater $`\gamma `$ is, the sharper is the $`H_z`$ profile, Fig. 2. Interestingly, the dependence $`(xb)^\beta `$ sufficiently well describes $`H_z(x)`$ even if $`x`$ is not too close to $`b`$, see Fig. 6.
Consider now the solution in the limit of small positive values of $`\gamma `$. If $`\gamma 0`$, two cases are possible: $`H_z`$ remains a constant, or it increases as $`\gamma ^1`$ (i.e., $`J_{c1}J_{c0}=`$ const.). In the first case one has $`\alpha \gamma /2\pi `$, $`H_aH_b\gamma `$, and the function $`f(t)`$ tends to zero. Thus, according to Eqs. (15) and (16), $`J_10`$, and the solution goes over to the well-known result for the Bean critical state model with $`J_c=J_{c0}`$. In the second case $`J_0(x)+J_1(x)`$ also tends to the solution corresponding to a constant $`J_c`$ but now $`J_c=J_{c1}`$.
In the limiting case $`\gamma +\mathrm{}`$, this parameter drops out from Eqs. (15), (16), (19), (22), and $`J_1(x)`$ depends only on $`H_z^0`$, $`J_{c0}`$. In other words, if $`J_{c1}J_{c0},H_z^0`$ or $`J_{c1}>J_{c0}H_z^0`$, the solution becomes practically independent of $`J_{c1}`$. The distribution of the magnetic field in this case can be understood using Eq. (26). It turns out that $`C_b^+\gamma H_z^0`$ for $`\gamma 1`$, and hence
$$H_z(x)H_z^0(x^2b^2)^\beta $$
(36)
with $`\beta 0`$. This means we have an abrupt step of height $`H_z^0`$ at $`x=b`$ (see Fig. 2).
It should be emphasized that this limiting case, $`\gamma +\mathrm{}`$, corresponds to intrinsic pinning in high-$`T_c`$ superconductors in which the ratio $`[j_c(\pi /2)/j_c(0)]=J_{c1}/J_{c0}`$ can be sufficiently large (see, e.g., Ref. ). Thus, our solution of this limit can be used for analyzing the critical state in these superconductors. In particular, it follows from Eqs. (19), (22) that the position of the flux front, $`b/w`$, is a function of $`H_a/H_{cs}`$ and of the parameter $`H_z^0/H_{cs}`$, see Fig. 7. In general this function can not be
fitted by scaling the dependence found in the isotropic case,
$$\frac{b_0}{w}=\frac{1}{\mathrm{cosh}(H_a/H_{cs})},$$
(37)
using some effective value of $`H_{cs}`$. Rather, the shape of $`b(H_a)`$ essentially depends on the ratio $`H_z^0/H_{cs}`$. Therefore, measuring $`b(H_a)`$ in principle can give information not only on $`H_{cs}=J_{c0}/\pi `$ but also on $`H_z^0`$, i.e., about the width of the peak in $`j_c(\theta )`$, see Eq. (5). In particular, when $`H_z^0J_{c0}`$, Eqs. (19), (22) lead to the following expression for the front position:
$$\left(\frac{b}{w}\right)^2\frac{1+k\left(H_z^0/H_{cs}\right)^2\mathrm{tanh}^2(H_a/H_{cs})}{\mathrm{cosh}^2(H_a/H_{cs})},$$
(38)
where the constant $`k`$ is determined by the root of the equation
$`{\displaystyle \frac{\pi }{4}}(u^21)=u\mathrm{arctan}u,`$ (39)
$`k={\displaystyle \frac{16}{\pi ^2}}{\displaystyle \frac{u^2}{(1+u^2)^2}}0.394.`$ (40)
Note that the right hand side of Eq. (29) cannot be reduced to the dependence (28) in the whole interval of changes of $`H_a`$ when $`H_z^0`$ is different from zero. The exact values of the front position $`b(H_a)`$ are shown in Fig. 7 for the limit of large $`\gamma 1`$, for $`J_{c1}=11`$ and $`H_z^0=0\mathrm{}1.5`$ in units of $`J_{c0}=1`$.
In the third limiting case when $`\gamma \mathrm{}`$, one has $`\beta 1`$, $`C_b^{}H_z^0`$ and the the induction profile becomes
$$H_z(x)(x^2b^2)$$
(41)
with a small prefactor of the order of $`H_z^0/|\gamma |`$. Thus, for $`\gamma 1`$ the flux front at $`x=b`$ practically disappears while, according to Eq. (20), the second front near $`x=a`$ is well developed, see Fig. 5.
Finally, we consider in some detail the case of small negative values of $`\gamma `$ when $`H_z^0H_{cs}=J_{c0}/\pi `$ while the ratio $`J_{c0}/J_{c1}`$ is not close to unity. This case can give some idea of pinning by columnar defects, which produce a peak in $`j_c(\theta )`$ at $`\theta =0`$. Indeed, if one assumes that the characteristic width of the peak, $`\theta _0`$, is small ($`\theta _01`$), then it follows from the definitions of $`H_z^0`$ and $`\gamma `$ that $`H_z^0J_{c0}/2\theta _0`$ and $`|\gamma |<2\theta _0`$. Since the solution with $`\gamma =0`$ and $`J_c=J_{c1}`$ describes the critical state in the strip before the irradiation \[we assume that the columnar defects do not change $`j_c(\theta )`$ at $`\theta >\theta _0`$\], the difference between the solutions corresponding to $`\gamma 0`$ and $`\gamma =0`$ provides information on pinning by columnar defects. In the considered case this difference is small, and it can be analyzed analytically. In particular, we obtain the following relation between the positions of the flux fronts, $`b`$ and $`b_1`$, obtained at the same $`H_a`$ in the strip with and without columnar defects, respectively:
$$\mathrm{arccosh}\frac{w}{b_1}\mathrm{arccosh}\frac{w}{b}=\frac{|\gamma |}{\pi }g(h),$$
(42)
where $`h\pi H_a/J_{c1}`$, $`w/b_1=\mathrm{cosh}(h)`$, and the function $`g(h)`$ has the form:
$`g(h)={\displaystyle _0^h}\mathrm{ln}(2\mathrm{cosh}t)𝑑t.`$ (43)
Since $`g`$ is a nonlinear function of $`h`$,
$`g(h){\displaystyle \frac{1}{2}}h^2+0.411(1e^{1.8h}),`$ (44)
the exact dependence $`b(H_a)`$ can not be described by Eq. (28) with some effective $`H_{cs}`$. The prefactor
$`{\displaystyle \frac{|\gamma |}{\pi }}{\displaystyle \frac{2\theta _0}{\pi }}{\displaystyle \frac{j_c(0)j_c(\pi /2)}{j_c(0)}}`$
in Eq. (31) is determined by the characteristics of pinning by the columnar defects, i.e., by the width and height of the peak in $`j_c(\theta )`$.
## IV conclusions
An exact solution of the critical state equations for the strip in perpendicular magnetic field is derived for an induction-dependent critical sheet current $`J_c(H_z)`$ described by Eqs. (4). This model dependence may be used to simulate the intrinsic pinning by CuO planes
($`\gamma >0`$) or pinning by extended defects ($`\gamma <0`$) in high-$`T_c`$ superconductors. In the case $`\gamma >0`$, the $`H_z`$ profile in the vicinity of the flux front is sharper than in the isotropic case, and the current density has a sharp peak there. In the limiting case, $`\gamma 1`$, which may describe the intrinsic pinning in high-$`T_c`$ superconductors, the field profile $`H_z(x)`$ has a sharp rectangular step. In the opposite situation, $`\gamma <0`$, two flux fronts can occur in the superconductor; the $`H_z`$ profile near $`x=b`$ is less steep than in the isotropic case, and the current density is a monotonic function of $`x`$. In both cases of positive and negative $`\gamma `$ the profile $`H_z(x)`$ in a sufficiently large vicinity of the flux front is well approximated by the expression $`H_z(x)(xb)^\beta `$ with the exponent $`\beta =0.5\pi ^1\mathrm{arctan}(\gamma /2)`$.
The experimental investigation of flux-density profiles near the flux front and of the $`H_z`$ dependence of the penetration depth can give information on the strength and anisotropy of flux line pinning in superconductors.
###### Acknowledgements.
G.P.M. acknowledges the hospitality of the Max-Planck-Institut für Metallforschung, Stuttgart.
## A Numerical evaluation
The condition that two integrals have to vanish, e.g. Eqs. (17,18) of the form $`I_1(a,b)=0`$ and $`I_2(a,b)=0`$, we satisfy by minimizing the function $`U(a,b)=I_1^2+I_2^2`$ with respect to $`a`$ and $`b`$. After this we calculate the sheet current $`J_1(x)`$ from Eqs. (15,16) and the magnetic field $`H_z(x)`$ from Eqs. (9) and (14).
The integrals (9), (15-19), and (22) over the variable $`t`$ have integrands which possess one or several infinities at the points $`t=0`$, $`t=x`$, $`t=b`$ and $`t=a`$ where the denominators vanish. We evaluate such integrals in the following way.
In the integrals containing a factor $`(tx)^1`$ we subtract the singular part and integrate it analytically, e.g.,
$`{\displaystyle _0^a}{\displaystyle \frac{f(t)dt}{t^2x^2}}𝑑t={\displaystyle _0^a}{\displaystyle \frac{f(t)f(x)}{t^2x^2}}𝑑t{\displaystyle \frac{f(x)}{2x}}\mathrm{ln}{\displaystyle \frac{a+x}{ax}}.`$ (A1)
Then we divide the integration interval into pieces bounded by the remaining singularities, $`0tb`$, $`bta`$, and $`at1`$. In each interval we substitute the integration variable by an appropriate function $`t=t(u)`$ and integrate over $`u`$ such that the new integrand has no infinity and vanishes rapidly at the boundaries. This new integral may thus be evaluated as a sum over an equidistant grid $`u_i`$ with constant weights. For example we write
$`{\displaystyle _0^\tau }g(t)𝑑t={\displaystyle _0^1}g[t(u)]t^{}(u)𝑑u{\displaystyle \underset{i=1}{\overset{N}{}}}g_iw_i`$ (A2)
with $`g_i=g[t(u_i)]`$, $`u_i=(i1/2)/N`$, $`w_i=t^{}(u_i)/N`$, $`t^{}(u)=dt/du`$, and $`i=1,\mathrm{\hspace{0.17em}2},\mathrm{\hspace{0.17em}3},\mathrm{}N`$. This integration method is very accurate if the substitution is chosen such that the weights $`w_i`$ and the products $`g_iw_i`$ vanish rapidly at the integration boundaries, e.g., $`w_iu_i^p`$ and $`w_i(1u_i)^q`$ with $`p1`$ and $`q1`$. Simple choices of this substitution in the example (A2) are
$`t(u)=(3u^22u^3)\tau ,t^{}(u)=6u(1u)\tau ,`$ (A3)
or better,
$`t(u)=(10u^315u^4+6u^5)\tau ,t^{}(u)=30u^2(1u)^2\tau .`$ (A4)
Higher accuracy is achieved by the following substitution. We chose equidistant $`u_i=(i1/2)/N`$ as above and then iterate (A3) $`m`$ times starting with $`s_i=u_i`$ and $`w_i=\tau /N`$ according to
$`w:=6(ss^2)w,s:=3s^22s^3(m\mathrm{times}).`$ (A5)
Finally we write $`t(u_i)=s_i\tau `$. The weights $`w_i=t^{}(u_i)/N`$ of this substitution vanish at the boundaries with exponents $`p=q=2^{(m1)}`$, which can be made arbitrarily large. For example, using $`m=5`$ iterations one gets the exponents $`p=q=2^4=16`$.
An infinity $`g(t)1/t^\eta `$ in the original integral (A2) leads, after this substitution, to a new integrand vanishing at $`t=0`$ as $`g[t(u)]t^{}(u)u^\vartheta `$ with $`\vartheta =p(1\eta )\eta `$. Thus, for the example $`\eta =1/2`$ with $`p=16`$ the new integrand near $`u=0`$ vanishes as $`u^{7.5}`$ and the terms in the sum (A2) as $`(i1/2)^{7.5}`$, in spite of the singular original integrand. For general exponent $`\eta `$, to reach high accuracy one should choose $`m`$ so large that the new exponent is $`\vartheta =(1\eta )2^{m1}\eta 4`$, or approximately $`m3.51.5\mathrm{ln}(1\eta )`$. To avoid spurious results due to rounding errors, one has to add in all vanishing denominators a small $`ϵ10^{15}`$ by writing, e.g., $`(|t^2b^2|+ϵ)^\beta `$.
In the limit of a large negative slope $`\gamma \mathrm{}`$ one has $`\beta 1`$ and the integrals (17,18) containing a factor $`|t^2b^2|^\beta `$ are close to diverging. In this case the singular part in these integrals should be integrated analytically, similar as shown in Eq. (A1). The subtracted terms are conveniently chosen such that the integral which has to be taken analytically is simple, e.g., $`t(b^2t^2)^\beta 𝑑t`$. Note that the numerator $`f(t)`$ in Eqs. (15–19,22) is discontinuous at $`t=b`$.
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# Tsallis’ entropy maximization procedure revisited
## I Introduction
Tsallis’ thermostatistics is by now recognized as a new paradigm for statistical mechanical considerations. One of its crucial ingredients, Tsallis’ normalized probability distribution , is obtained by following the well known MaxEnt route . One maximizes Tsallis’ generalized entropy
$$\frac{S_q}{k}=\frac{1_{i=1}^wp_i^q}{q1},$$
(1)
($`kk(q)`$ tends to the Boltzmann constant $`k_B`$ in the limit $`q1`$ ) subject to the constraints (generalized expectation values)
$`{\displaystyle \underset{i=1}{\overset{w}{}}}p_i`$ $`=`$ $`1`$ (2)
$`{\displaystyle \frac{_{i=1}^wp_i^qO_j^{(i)}}{_{i=1}^wp_i^q}}`$ $`=`$ $`O_j_q,`$ (3)
where $`p_i`$ is the probability assigned to the microscopic configuration $`i`$ ($`i=1,\mathrm{},w`$) and one sums over all possible configurations $`w`$. $`O_j^{(i)}`$ ($`j=1,\mathrm{},n`$) denote the $`n`$ relevant observables (the observation level ), whose generalized expectation values $`O_j_q`$ are (assumedly) a priori known.
The Lagrange multipliers recipe entails maximizing
$$F=\frac{S_q}{k}\alpha _0\left(\underset{i=1}{\overset{w}{}}p_i1\right)\underset{j=1}{\overset{n}{}}\lambda _j\left(\frac{_{i=1}^wp_i^qO_j^{(i)}}{_{i=1}^wp_i^q}O_j_q\right),$$
(4)
yielding
$$p_i=\frac{f_i^{1/(1q)}}{\overline{Z}_q},$$
(5)
where
$$f_i=1\frac{(1q)_j\lambda _j\left(O_j^{(i)}O_j_q\right)}{_jp_j^q},$$
(6)
is the so-called configurational characteristic and
$$\overline{Z}_q=\underset{i}{}f_i^{1/(1q)},$$
(7)
stands for the partition function.
The above procedure, originally employed in , overcomes most of the problems posed by the old, unnormalized way of evaluating Tsallis’generalized mean values . Some hardships remain, though. One of them is that numerical difficulties are sometimes encountered, as the $`p_i`$ expression is explicitly self-referential. An even more serious problem is also faced: a maximum is not necessarily guaranteed. Indeed, analyzing the concomitant Hessian so as to ascertain just what kind of extreme we face, one encounters the unpleasant fact that this Hessian is not diagonal.
In the present effort we introduce an alternative Lagrange route, that overcomes the above mentioned problems.
## II The new Lagrange multipliers’ set
We extremize again (1) subject to the constraints (3), but, and herein lies the central idea, rephrase (4) by recourse to the alternative form
$$\underset{i=1}{\overset{w}{}}p_i^q\left(O_j^{(i)}O_j_q\right)=0,$$
(8)
with $`j=1,\mathrm{},n`$. We have now
$$F=\frac{S_q}{k}\alpha _0\left(\underset{i=1}{\overset{w}{}}p_i1\right)\underset{j=1}{\overset{n}{}}\lambda _j^{}\underset{i=1}{\overset{w}{}}p_i^q\left(O_j^{(i)}O_j_q\right),$$
(9)
so that, following the customary variational procedure and eliminating $`\alpha _0`$ we find that the probabilities are, formally, still given by (5). However, in terms of the new set of Lagrange multipliers, the configurational characteristics do not depend explicitly on the probabilities
$$f_i^{}=1(1q)\underset{j}{}\lambda _j^{}\left(O_j^{(i)}O_j_q\right).$$
(10)
Comparing (10) with (6), it is clear that the Lagrange multipliers $`\lambda _j`$ of the Tsallis-Mendes-Plastino formalism (TMP) and $`\lambda _j^{}`$ (present treatment) can be connected via
$$\lambda _j^{}=\frac{\lambda _j}{_ip_i^q},$$
(11)
which leads to the nice result $`f_i^{}=f_i`$. The probabilities that appear in (11) are those special ones that maximize the entropy, not generic ones. The ensuing, new partition function is also of the form (7), with $`f_i>0`$ the well known Tsallis’ cut-off condition . Notice that now the expression for the MaxEnt probabilities $`p_i`$ is NOT explicitly self-referential.
In order to ascertain the kind of extreme we are here facing we study the Hessian, that now is of diagonal form. The maximum condition simplifies then to the requirement
$$\frac{^2F}{p_i^2}<0.$$
(12)
The above derivatives are trivially performed yielding
$$\frac{^2F}{p_i^2}=qp_i^{q2}f_i,$$
(13)
which formally coincides with the maximum requirement one finds in the case of Tsallis’ unnormalized formalism. Since the $`f_i`$ are positive-definite quantities, for a maximum one should demand that $`q>0`$.
Extremes found by following the celebrated Lagrange procedure depend only on the nature of the constraints, not on the form in which they are expressed. Thus, the two sets of multipliers lead to the same numerical values for the micro-state probabilities. Via (11) one is always able to establish a connection between both treatments.
The present algorithm exhibits the same nice properties of the TMP formalism, namely:
* The MaxEnt probabilities are invariant under uniform shifts of the Hamiltonian’s energy spectrum (see, for instance, the illuminating discussion of Di Sisto et al.).
Indeed, after performing the transformation
$`ϵ_i`$ $``$ $`ϵ_i+ϵ_0`$ (14)
$`U_q`$ $``$ $`U_qϵ_0,`$ (15)
on equation (5), with $`f_i`$ given by (10), we trivially find that the probabilities $`p_i`$ keep their forms invariant if the $`\lambda _j^{}`$ do not change. Due to relation (11), the $`\lambda _j`$ are invariant too.
* The mean value of unity equals unity, i.e., $`1_q=1`$, which is not the case with the unnormalized expectation values .
* One easily finds that, for two independent subsystems $`A,B`$, energies add up: $`U_q(A+B)=U_q(A)+U_q(B)`$.
## III OLM treatment and Rènyi’s measure
The above OLM recipe has the purpose of diagonalizing the Hessian characterizing the kind of extreme one finds after a Lagrange treatment. It is natural to ask whether the concomitant recipe works also for other information measures in addition to the Tsallis’ one. We discuss next the Rènyi’s measure ($`R_q`$) instance in an OLM context. We have
$$\frac{R_q}{k}=\frac{\mathrm{ln}\left[1+(1q)S_q\right]}{1q},$$
(16)
or (Cf. (1))
$$\frac{R_q}{k}=\frac{\mathrm{ln}\left(_ip_i^q\right)}{q1}.$$
(17)
We extremize now $`R_q`$ subject to the constraints (8) with Lagrange multipiers $`\mu _j^{}`$. The ensuing Lagrange functional reads
$$F^{}=\frac{R_q}{k}\alpha _0\left(\underset{i=1}{\overset{w}{}}p_i1\right)\underset{j=1}{\overset{n}{}}\mu _j^{}\underset{i=1}{\overset{w}{}}p_i^q\left(O_j^{(i)}O_j_q\right),$$
(18)
whose first derivatives are
$$\frac{F^{}}{p_i}=\frac{qp_i^{q1}}{(q1)_ip_i^q}\alpha _0qp_i^{q1}\underset{j}{}\mu _j^{}\left(O_j^{(i)}O_j_q\right)=0.$$
(19)
The entropy term will lead to a self referencial probability distribution. We have
$$p_i=\frac{(f_i^{})^{\frac{1}{1q}}}{Z_q^{}},$$
(20)
with a functional characteristic given by
$$f_i^{}=1+(1q)\underset{i}{}p_i^q\underset{j}{}\mu _j^{}\left(O_j^{(i)}O_j_q\right),$$
(21)
and the partition function
$$Z_q^{}=\underset{i}{}(f_i^{})^{\frac{1}{1q}}.$$
(22)
We see that the OLM treatment leads to a self-referential solution that the “Rènyi-TMP” treatament does not exhibit. Indeed, on account of the fact both the constraints and the entropy share the $`_ip_i^q`$ denominator, the TMP-Rènyi distribution has characteristics of the form
$$f_i=1+(1q)\underset{j}{}\mu _j\left(O_j^{(i)}O_j_q\right).$$
(23)
Comparing Eqs. (21) and (23), one observes that the Lagrange multipliers are connected via
$$\mu _j^{}=\frac{\mu _j}{_ip_i^q},$$
(24)
that is the same relation found for Tsallis’ measure (see Equation (11)). The following correspondence can be appreciated
$$\lambda _j^{}=\mu _j,$$
since the configurational characteristics of Equation (23) is identical to the one of Equation (10).
Consider now, for the second derivatives, the non-diagonal terms. One has. in the OLM instance,
$$\frac{^2F^{}}{p_jp_i}=\frac{q^2}{1q}\frac{(p_ip_j)^{q1}}{\left(_kp_k^q\right)^2}.$$
(25)
If one uses the standard TMP approach of Ref. one finds, instead, for these non-diagonal terms,
$$\frac{^2F}{p_jp_i}=\frac{q^2}{1q}\frac{f_i(p_ip_j)^{q1}}{\left(_kp_k^q\right)^2}.$$
(26)
There is an important distinction to be made between (25) and (26), that refers to the origin of the non-diagonal terms. In the first case, they have as a source just the entropy term of the Lagrangian. In the second, the constraints also contribute.
As a consequence we gather that the OLM approach will work as nicely as in the Tsallis case only for those measures that do not give rise to off-diagonal terms in the associated Hessian.
## IV Thermodynamics
We pass now to the question of writing down the basic mathematical relationships of Thermodynamics, as expressed with respect to the new set of Lagrange multipliers $`\lambda _j^{}`$.
In order to do this in the most general quantal fashion we shall work in a basis-independent way. This requires consideration of the statistical operator (or density operator) $`\widehat{\rho }`$ that maximizes Tsallis’ entropy, subject to the foreknowledge of $`M`$ generalized expectation values (corresponding to $`M`$ operators $`\widehat{O}_j`$). These take the form
$$\widehat{O}_j_q=\frac{Tr(\widehat{\rho }^q\widehat{O}_j)}{Tr(\widehat{\rho }^q)},j=1,\mathrm{},M.$$
(27)
To these we must add, of course, the normalization requirement
$$Tr\widehat{\rho }=1.$$
(28)
The TMP formalism, where relations are written in terms of the “old” Lagrange multipliers $`\lambda _j`$, yields the usual thermodynamical relationships , namely
$`{\displaystyle \frac{}{\widehat{O}_j_q}}\left({\displaystyle \frac{S_q}{k}}\right)`$ $`=`$ $`\lambda _j`$ (29)
$`{\displaystyle \frac{}{\lambda _i}}\left(\mathrm{ln}_qZ_q\right)`$ $`=`$ $`\widehat{O}_j_q,`$ (30)
where
$$\mathrm{ln}_q\overline{Z}_q=\frac{\overline{Z}_q^{1q}1}{1q}$$
(31)
and
$$\mathrm{ln}_qZ_q=\mathrm{ln}_q\overline{Z}_q\underset{j}{}\lambda _j\widehat{O}_j_q,$$
(32)
so that the essential mathematical structure of Thermodynamics is preserved.
Following the standard procedure one gets
$$\widehat{\rho }=\overline{Z}_q^1\left[1(1q)\underset{j}{\overset{M}{}}\lambda _j^{}\left(\widehat{O}_j\widehat{O}_j_q\right)\right]^{\frac{1}{1q}},$$
(33)
where $`\overline{Z}_q`$ stands for the partition function
$$\overline{Z}_q=Tr\left[1(1q)\underset{j}{}\lambda _j^{}\left(\widehat{O}_j\widehat{O}_j_q\right)\right]^{\frac{1}{1q}}.$$
(34)
Enters here Tsallis’cut-off condition . The form (33) does not a priori guarantee that we will have a positive-definite operator. Some additional considerations are requested.
Consider the operator
$$\widehat{A}=1(1q)\underset{j}{}\lambda _j^{}\left(\widehat{O}_j\widehat{O}_j_q\right)$$
(35)
enclosed within parentheses in (33). One must ensure its positive-definite character. This entails that the eigenvalues of $`\widehat{A}`$ must be non-negative quantities. This can be achieved by recourse to an heuristic cut-off procedure. We replace (33) by
$$\widehat{\rho }=\overline{Z}_q^1\left[\widehat{A}\mathrm{\Theta }(\widehat{A})\right]^{1/(1q)},$$
(36)
with $`\overline{Z}_q`$ given by
$$\overline{Z}_q=Tr\left[\widehat{A}\mathrm{\Theta }(\widehat{A})\right]^{1/(1q)},$$
(37)
where $`\mathrm{\Theta }(x)`$ is the Heaviside step-function. Equations (36)-(37) are to be re-interpreted as follows. Let $`|i`$ and $`\alpha _i,`$ stand, respectively, for the eigenvectors and eigenvalues of the operator (35), whose spectral decomposition is then
$$\widehat{A}=\underset{i}{}\alpha _i|ii|.$$
(38)
In the special basis used above $`\widehat{\rho }`$ adopts the appearance
$$\widehat{\rho }=\overline{Z}_q^1\underset{i}{}f(\alpha _i)|ii|,$$
(39)
with $`f(x)`$ defined as
$$f(x)=0,forx0,$$
(40)
and
$$f(x)=x^{\frac{1}{1q}},forx>0.$$
(41)
Notice that $`f(x)`$ possesses, for $`0<q<1,`$ a continuous derivative for all $`x.`$ Moreover,
$$\frac{df(x)}{dx}=\left(\frac{1}{q1}\right)\left[x\mathrm{\Theta }(x)\right]^{\frac{q}{1q}}.$$
(42)
In terms of the statistical operator, Tsallis’ entropy $`S_q`$ reads
$`{\displaystyle \frac{S_q}{k}}`$ $`=`$ $`{\displaystyle \frac{1}{q1}}Tr\left[\widehat{\rho }^q\left(\widehat{\rho }^{1q}\widehat{I}\right)\right]`$ (43)
$`=`$ $`{\displaystyle \frac{1}{q1}}Tr\left[\widehat{\rho }^q\left(\overline{Z}_q^{q1}\widehat{A}\mathrm{\Theta }(\widehat{A})\widehat{I}\right)\right]`$ (44)
$`=`$ $`{\displaystyle \frac{\overline{Z}_q^{q1}}{q1}}Tr\left[\widehat{\rho }^q\widehat{A}\mathrm{\Theta }(\widehat{A})\right]{\displaystyle \frac{Tr(\widehat{\rho }^q)}{(q1)}},`$ (45)
where $`\widehat{I}`$ is the unity operator.
Obviously, $`\widehat{\rho }`$ commutes with $`\widehat{A}`$. The product of these two operators can be expressed in the common basis that diagonalizes them
$$\widehat{\rho }^q\widehat{A}\mathrm{\Theta }(\widehat{A})=\overline{Z}_q^q\underset{i}{}[f(\alpha _i)]^q\alpha _i|ii|,$$
(46)
which entails, passing from the special basis $`|i`$ to the general situation, that
$$\widehat{\rho }^q\widehat{A}\mathrm{\Theta }(\widehat{A})=\widehat{\rho }^q\left[1(1q)\underset{j}{}\lambda _j^{}\left(\widehat{O}_j\widehat{O}_j_q\right)\right],$$
(47)
and, consequently
$$\frac{S_q}{k}=\frac{\overline{Z}_q^{q1}1}{q1}Tr\left(\widehat{\rho }^q\right)+\overline{Z}_q^{q1}\underset{j}{}\lambda _j^{}Tr\left[\widehat{\rho }^q\left(\widehat{O}_j\widehat{O}_j_q\right)\right].$$
(48)
Since the last term of the right-hand-side vanishes, by definition (8), we finally arrive at
$$\frac{S_q}{k}=\frac{\overline{Z}_q^{q1}1}{q1}Tr\left(\widehat{\rho }^q\right).$$
(49)
Now, from the very definition (in terms of $`\widehat{\rho }`$) of Tsallis’ entropy $`S_q`$ , we find
$$Tr\left(\widehat{\rho }^q\right)=1+(1q)\frac{S_q}{k},$$
(50)
so that (49) and (50) lead to
$$Tr(\widehat{\rho }^q)=\overline{Z}_q^{1q}$$
(51)
and
$$S_q=k\mathrm{ln}_q\overline{Z}_q,$$
(52)
where $`\mathrm{ln}_q\overline{Z}_q`$ has been introduced in (31).
Using (51), equation (11) can be rewritten as
$$\lambda _j^{}=\frac{\lambda _j}{\overline{Z}_q^{1q}}.$$
(53)
Following we define now
$$\mathrm{ln}_qZ_q^{}=\mathrm{ln}_q\overline{Z}_q\underset{j}{}\lambda _j^{}\widehat{O}_j_q,$$
(54)
which leads finally to (see (29) and (30))
$`{\displaystyle \frac{}{\widehat{O}_j_q}}\left({\displaystyle \frac{S_q}{k}}\right)`$ $`=`$ $`\overline{Z}_q^{1q}\lambda _j^{}=\lambda _j`$ (55)
$`{\displaystyle \frac{}{\lambda _j^{}}}\left(\mathrm{ln}_qZ_q^{}\right)`$ $`=`$ $`\widehat{O}_j_q.`$ (56)
Equations (55) and (56) constitute the basic Information Theory relations on which to build up, à la Jaynes , Statistical Mechanics. Notar que OLM conduce a las mismas relaciones que TMP, y que no ha sido recesario tomar $`k`$ como constante en ninguna parte del desarrollo.
As a special instance of Eqs. (55) and (56) let us discuss the Canonical Ensemble, where they adopt the appearance
$`{\displaystyle \frac{}{U_q}}\left({\displaystyle \frac{S_q}{k}}\right)`$ $`=`$ $`\overline{Z}_q^{1q}\beta ^{}=\beta `$ (57)
$`{\displaystyle \frac{}{\beta ^{}}}\left(\mathrm{ln}_qZ_q^{}\right)`$ $`=`$ $`U_q,`$ (58)
where (see equation (54))
$$\mathrm{ln}_qZ_q^{}=\mathrm{ln}_q\overline{Z}_q\beta ^{}U_q.$$
(59)
Proponemos que $`\beta ^{}`$ Finally, the specific heat reads
$$C_q=\frac{U_q}{T}=\frac{U_q}{\beta ^{}}\frac{d\beta ^{}}{dT}=\frac{U_q}{\beta ^{}}\frac{1}{k^{}T^2}=k^{}\beta ^2\frac{U_q}{\beta ^{}}.$$
(60)
We conclude that the mathematical form of the thermodynamic relations is indeed preserved by the present treatment. Both sets of Lagrange multipliers accomplish this feat and they are connected via (11). The primed one, however, allows for a simpler treatment, as will be illustrated below.
## V Simple applications
We consider now some illustrative examples. They are chosen in such a manner that each of them discusses a different type of situation: classical and quantal systems, the latter in the case of both finite and infinite number of levels.
### A The classical harmonic oscillator
Let us consider the classical harmonic oscillator in the canonical ensemble. We can associate with the classical oscillator a continuous energy spectrum $`ϵ(n)=ϵn`$ with $`ϵ>0`$ and $`n^+`$ compatible with the cut-off condition. The ensuing MaxEnt probabilities adopt the appearance
$$p_q(n,t^{})=\frac{\left[f_q(n,t^{})\right]^{1/(1q)}}{\overline{Z}_q(t^{})},$$
(61)
where
$$f_q(n,t^{})=1(1q)\frac{(nu_q)}{t^{}},$$
(62)
and
$$\overline{Z}_q(t^{})=_0^{n_{max}}\left[f_q(n,t^{})\right]^{1/(1q)}𝑑n,$$
(63)
with $`u_q=U_q/ϵ`$ and $`t^{}=k^{}T/ϵ`$. We have introduced also $`n_{max}`$ as the upper integration limit on account of Tsallis’ cut-off condition. One appreciates the fact that $`n_{max}\mathrm{}`$ if $`q>1`$. $`n_{max}`$ is, of course, the maximum $`n`$-value that keeps $`\left[f_q(n,t^{})\right]^{1/(1q)}>0`$ for $`q<1`$.
The normalization condition reads
$$u_q(t^{})=\frac{_0^{n_{max}}\left[p_q(n,t^{})\right]^qn𝑑n}{_0^{n_{max}}\left[p_q(n,t^{})\right]^q𝑑n},$$
(64)
or, using (61),
$$u_q(t^{})=\frac{_0^{n_{max}}\left[f_q(n,t^{})\right]^{q/(1q)}n𝑑n}{_0^{n_{max}}\left[f_q(n,t^{})\right]^{q/(1q)}𝑑n}.$$
(65)
Due to the form of $`f_q`$, equation (65) constitutes a well-defined expression. By explicitly performing the integrals for $`1<q<2`$ (for $`q2`$ the integrals diverge) we obtain
$$u_q(t^{})=\frac{t^2/(2q)(1+(1q)u_q/t^{})^{(2q)/(1q)}}{t^{}(1+(1q)u_q/t^{})^{1/(1q)}}.$$
(66)
After a little algebra, the above equation leads to the simple result
$$u_q(t^{})=t^{}.$$
(67)
Replacing now $`u_q=U_q/ϵ`$ and $`t^{}=k^{}T/ϵ`$, we obtain $`U_q=k^{}T`$, so that the specific heat reads
$$C_q=k^{}.$$
(68)
It is worthwhile to remark that, in the case of this particular example, we formally regain the usual expressions typical of the $`q=1`$ case. Due to fact that we possess a degree of freedom in the definition of $`k^{}`$, we can set $`k^{}=k_B`$ and thus recover Gibbs’ Thermodynamics. Performing the pertinent integral and using (67), the partition function becomes
$$\overline{Z}_q(t^{})=t^{}(2q)^{1/(1q)}.$$
(69)
According to equation (53), $`t^{}`$ can be written in terms of $`t`$ and $`\overline{Z}_q`$, allowing us to recover
$$\overline{Z}_q(t)=t^{1/q}(2q)^{1/[q(1q)]},$$
(70)
and, consequently,
$`u_q`$ $`=`$ $`t^{1/q}(2q)^{1/q}`$ (71)
$`C_q`$ $`=`$ $`{\displaystyle \frac{k}{2}}(2q)^{1/q}t^{(1q)/q}.`$ (72)
These results are identical to those of , but are here derived in a remarkably simpler fashion.
### B The two-level system and the quantum harmonic oscillator
Let us consider the discrete case of a single particle with an energy spectrum given by $`E_n=ϵn`$, where $`ϵ>0`$ and $`n=0,1,\mathrm{},N`$. If $`N=1`$, we are facing the non degenerate two level system, while, if $`n\mathrm{}`$, the attendant problem is that of the quantum harmonic oscillator.
The micro-state probabilities are of the form, once again
$$p_n=\frac{f_n^{1/(1q)}}{\overline{Z}_q}$$
(73)
with
$$\overline{Z}_q=\underset{n=0}{\overset{N}{}}f_n{}_{}{}^{1/(1q)}.$$
(74)
The configurational characteristics take the form
$$f_n(t^{})=1(1q)(nu_q)/t^{}$$
(75)
where again (see (V A)), $`t^{}=k^{}T/ϵ`$ and $`u_q=U_q/k^{}`$.
Using (73), the mean energy can be written as
$$u_q=\frac{_{n=0}^Nf_n^{q/(1q)}n}{_{n=0}^Nf_n^{q/(1q)}},$$
(76)
which, using the explicit form of $`f_n`$ and rearranging terms, allows one to write down the following equation
$$\underset{n=0}{\overset{N}{}}\left[1\frac{(1q)}{t^{}}(nu_q)\right]^{q/(1q)}(nu_q)=0,$$
(77)
which implicitly defines $`u_q`$. Notice that one does not arrive to a closed expression. However, in order to numerically solve for $`u_q`$, we just face (77). This equation is easily solved by recourse to the so-called “seed” methods (cut-off always taken care of), with quick convergence (seconds). This is to be compared to the TMP instance . In their case, one faces a non-linear coupled system of equations in order to accomplish the same task. This coupled system can be recovered from (77) and (73), writing $`t^{}`$ in terms of $`t`$.
### C Magnetic Systems
Consider now a very simple magnetic model, discussed, for instance, in : a quantum system of $`N`$ spin 1/2 non-interacting atoms in the presence of a uniform, external magnetic field $`\stackrel{}{H}=H\widehat{k}`$ (oriented along the unit vector $`\widehat{k}`$). Each atom is endowed with a magnetic moment $`\widehat{\stackrel{}{\mu }}^{(i)}=g\mu _0\widehat{\stackrel{}{S}}^{(i)},`$ where $`\mu _0=e/(2mc)`$ is Bohr’s magneton and $`\widehat{\stackrel{}{S}}^{(i)}=(\mathrm{}/2)\widehat{\stackrel{}{\sigma }}^{(i)}`$, with $`\widehat{\stackrel{}{\sigma }}^{(i)}`$ standing for the Pauli matrices. The concomitant interaction energy reads
$$\widehat{}=\underset{i=1}{\overset{N}{}}\widehat{\stackrel{}{\mu }}^{(i)}\stackrel{}{H}=\frac{g\mu _0}{\mathrm{}}H\widehat{S}_z,$$
(78)
where $`\widehat{\stackrel{}{S}}=_{i=1}^N\widehat{\stackrel{}{S}}^{(i)}`$ the total (collective) spin operator. The simultaneous eigenvectors of $`\widehat{\stackrel{}{S}}^2`$ and $`\widehat{S}_z`$ constitute a basis of the concomitant $`2^N`$-dimensional space. We have $`|S,M`$, with $`S=\delta ,\mathrm{},N/2,`$ $`M=S,\mathrm{},S,`$ and $`\delta N/2[N/2]=0`$ $`(1/2)`$ if $`N`$ is even (odd). The corresponding multiplicities are $`Y(S,M)=Y(S)=N!(2S+1)/[(N/2S)!(N/2+S+1)!]`$ . We recast the Hamiltonian in the simple form
$$\widehat{}=\frac{x^{}}{\beta ^{}}\widehat{S}_z,$$
(79)
with $`x^{}=g\mu _0H\beta ^{}/\mathrm{}`$. Our statistical operator can be written as
$$\widehat{\rho }=\frac{1}{\overline{Z}_q}\left[1(1q)x^{}\left(\widehat{S}_z\widehat{S}_z_q\right)\right]^{1/(1q)},$$
(80)
where
$$\overline{Z}_q=Tr\left[1(1q)x^{}\left(\widehat{S}_z\widehat{S}_z_q\right)\right]^{1/(1q)}.$$
(81)
Due to the cut-off condition, $`1(1q)x^{}\left(\widehat{S}_z\widehat{S}_z_q\right)>0`$.
The mean value of the spin $`z`$-component is computed according to (27)
$$\widehat{S}_z_q=\frac{Tr\left(\widehat{\rho }^q\widehat{S}_z\right)}{Tr\left(\widehat{\rho }^q\right)},$$
(82)
so that, replacing (80) into (82) and rearranging then terms we arrive at
$$Tr\left\{\left[1+(1q)x^{}\left(\widehat{S}_z\widehat{S}_z_q\right)\right]^{q/(1q)}\left(\widehat{S}_z\widehat{S}_z_q\right)\right\}=0.$$
(83)
More explicitly, one has
$$\underset{S=\delta }{\overset{N/2}{}}Y(S)\underset{M=S}{\overset{S}{}}\left[1+(1q)x^{}\left(M\widehat{S}_z_q\right)\right]^{q/(1q)}\left(M\widehat{S}_z_q\right)=0,$$
(84)
which is the equation to be solved in order to find $`\widehat{S}_z_q`$.
Notice that, once again, one faces just a single equation, that can be easily tackled. If one uses instead the TMP prescription (as discussed in ) one has to solve a coupled, highly non-linear system of equations. Such a system can be recovered from (84) if one replaces $`x^{}`$ by $`x/Tr(\rho ^q)`$ and adds the condition $`Tr(\rho ^q)`$ from (80).
As in , we consider now two asymptotic situations from the present viewpoint.
For $`x^{}0`$ we Taylor-expand (84) around $`x^{}=0`$ and find
$$\widehat{S}_z_q=\frac{qx^{}N}{4},$$
(85)
that leads to an effective particle number
$$N_{eff}^0=qN,$$
(86)
as in . Following the same mechanism and using (85), one finds that
$$Tr(\rho ^q)=2^{N(1q)}.$$
(87)
Remembering that $`x^{}=x/Tr(\rho ^q)`$, it is possible to recover the TMP normalized solution
$$\widehat{S}_z_q=\frac{qxN}{4}2^{N(q1)},$$
(88)
and
$$N_{eff}^{0^{(3)}}=qN2^{N(q1)}.$$
(89)
For $`x^{}\mathrm{},`$ and for $`0<q<1,`$ expression (84) leads to an equation identical to that of
$$\underset{S=\delta }{\overset{N/2}{}}Y(S)\underset{M=S}{\overset{S}{}}\left(M\widehat{S}_z_q\right)^{1/(1q)}=0,$$
(90)
whose solution reads $`\widehat{S}_z_q=N/2`$.
## VI Conclusions
In order to obtain the probability distribution $`p_i`$ that maximizes Tsallis’ entropy subject to appropriate constraints, Tsallis-Mendes-Plastino extremize
$`F=S_q\alpha _0\left({\displaystyle \underset{i=1}{\overset{w}{}}}p_i1\right){\displaystyle \underset{j=1}{\overset{n}{}}}\lambda _j\left({\displaystyle \frac{_{i=1}^wp_i^qO_j^{(i)}}{_{i=1}^wp_i^q}}O_j_q\right),`$
and obtain
$`p_i={\displaystyle \frac{f_i^{1/(1q)}}{\overline{Z}_q}},`$
where
$`f_i=1{\displaystyle \frac{(1q)_j\lambda _j\left(O_j^{(i)}O_j_q\right)}{k_jp_j^q}},`$
and $`\overline{Z}_q`$ is the partition function. Two rather unpleasant facts are thus faced, namely,
* $`p_i`$ explicitly depends upon the probability distribution (self-reference).
* The Hessian of $`F`$ is not diagonal.
In this work we have devised a transformation from the original set of Lagrange multipliers $`\{\lambda _j\}`$ to a new set $`\{\lambda _j^{}\}`$ such that
* Self-reference is avoided.
* The Hessian of $`F`$ becomes diagonal.
As a consequence, all calculations, whether analytical or numerical, become much simpler than in , as illustrated with reference to several simple examples. The primed multipliers
$`\lambda _j^{}={\displaystyle \frac{\lambda _j}{_ip_i^q}}`$
incorporate the $`p_i`$<sup>*</sup><sup>*</sup>*that maximize the entropy in their definition. Since one solves directly for the primed multipliers, such a simple step considerably simplifies the TMP treatment. Finally, we remark on the fact that the two sets of multipliers lead to thermodynamical relationships that involve identical intensive quantities (LABEL:rel).
###### Acknowledgements.
The financial support of the National Research Council (CONICET) of Argentina is gratefully acknowledged. F. Pennini acknowledges financial support from UNLP, Argentina.
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# 1 Introduction
## 1 Introduction
Noncommutative geometry has been applied to string theory in several different aspects -. In particular, Connes, Douglas and Schwarz found that in the matrix theory a constant $`C_{ij}`$ field background results in a gauge field theory living on a noncommutative space. It was also found that the low energy effective theory for D-branes in the NS-NS $`B`$ field background is also given by the same noncommutative gauge theory (NCGT) . We feel that it is important to understand the physics of noncommutative field theories.
The purpose of this paper is to understand better twisted bundles on noncommutative space. A noncommutative space is defined by the algebra $`𝒜`$ of functions on the space, and the Hilbert space on which $`𝒜`$ is realized as operators. According to Connes, vector bundles on a noncommutative space are defined as projective modules of $`𝒜`$ . In other words, a twisted bundle can be understood as a projection $`P`$ from a trivial bundle. Classically we restrict ourselves to bundles with fibers all of the same dimension at every point. This means that the projection $`P`$ is of constant rank. However, due to the existence of projective operators in $`𝒜`$, one can extend the definition of bundles to projective modules for which $`P`$ has non-constant rank. The $`U(1)`$ instanton solution of Nekrasov and Schwarz is found to be such an example .
In sec.2 we review the noncommutative algebra $`𝒜`$ of functions on a quantum plane, and its realization on Hilbert spaces. In sec.3, we find projective operators in $`𝒜`$ constructed out of operators $`U`$ which has right inverse but no left inverse. Formal gauge transformations by $`U`$ bring a bundle of constant dimension to a bundle with non-constant dimension. This can be viewed as a mathematical motivation for us not to exclude bundles with non-constant dimension. We review the notion of twisted bundles in sec.4, with special care to projective modules with non-constant rank . As an example we review the ADHM construction of $`U(1)`$ instantons in sec.5.
## 2 Quantum Plane
In Connes’ formulation of noncommutative geometry , a noncommutative space is specified by the spectral triple $`(𝒜,,𝒟)`$, where $`𝒜`$ is the noncommutative algebra of functions on the space, $``$ is the Hilbert space on which $`𝒜`$ is realized as an operator algebra, and $`𝒟`$ is the so-called Dirac operator which defines the metric and differential calculus. <sup>1</sup><sup>1</sup>1 For a short introduction to the role $`𝒟`$ plays in noncommutative geometry please see . In this paper, we will not use $`𝒟`$, but instead we will directly define the differential calculus for the nonocommutative space under consideration.
For the noncommutative space we are interested in, the coordinates $`x^i`$ satisfy
$$[x_i,x_j]=i\theta _{ij},$$
(1)
where $`þ_{ij}𝐑`$ are constant real numbers. The algebra $`𝒜`$ is defined by this relation and elements in $`𝒜`$ are functions of $`x_i`$.
There are two different Hilbert spaces that are often used in the literature for our $`𝒜`$. The first is the Hilbert space of simple harmonic oscillators (SHO), on which complex coordinates act as creation and annihilation operators. The second is the phase space of one particle on a classical space, on which the noncommutative coordinates are linear combinations of coordinates and momenta in quantum mechanics. We discuss these Hilbert spaces in the following two subsection.
For simplicity, we will consider a two dimensional noncommutative space, which is also called a quantum plane, with coordinates $`\widehat{x}`$ and $`\widehat{y}`$ satisfying
$$[\widehat{x},\widehat{y}]=iþ.$$
(2)
This relation is invariant under translation and rotation on the two dimensional plane. There is an anti-involution denoted by $``$ which is identified with Hermitian conjugation when $`𝒜`$ is realized as operators on a Hilbert space. We assume that $`þ`$ is a real number and
$$\widehat{x}^{}=\widehat{x},\widehat{y}^{}=\widehat{y}.$$
(3)
The algebra $`𝒜`$ is generated by $`\widehat{x}`$ and $`\widehat{y}`$. As a convention we choose the coordinates such that
$$þ0.$$
(4)
It follows that the commutation relation of the complex coordinates
$$\widehat{z}=\frac{\widehat{x}+i\widehat{y}}{2},\widehat{\overline{z}}\widehat{z}^{}=\frac{\widehat{x}i\widehat{y}}{2}$$
(5)
is similar to that of the creation and annihilation operators for a simple harmonic oscillator:
$$[\widehat{z},\widehat{\overline{z}}]=\frac{þ}{2}.$$
(6)
One can therefore use the SHO Hilbert space as a representation of $`𝒜`$.
One can adjoin differential one-forms $`(d\widehat{z})`$ and $`(d\widehat{\overline{z}})`$ to this algebra and generate the differential calculus. By definition, $`(d\widehat{z})`$ and $`(d\widehat{\overline{z}})`$ commute with $`\widehat{z}`$ and $`\widehat{\overline{z}}`$, and satisfy the anticommuting relations
$$(d\widehat{z})(d\widehat{\overline{z}})=(d\widehat{\overline{z}})(d\widehat{z}),(d\widehat{z})^2=(d\widehat{\overline{z}})^2=0.$$
(7)
One can also define the derivatives $`\widehat{}_i`$ ($`i=1,2`$) which satisfy
$`\widehat{}_i\widehat{x}_j=\widehat{x}_j\widehat{}_i+\delta _{ij},`$ (8)
$`[\widehat{}_1,\widehat{}_2]=0.`$ (9)
The action of a derivative on a function is
$$(\widehat{}_if)[\widehat{}_i,f].$$
(10)
When acting on functions, $`\widehat{}_iþ_{ij}^1\widehat{x}_j`$.
### 2.1 SHO Hilbert Space
The algebra $`𝒜`$ can be realized as an operator algebra according to
$$\widehat{z}|n=\sqrt{\frac{þn}{2}}|n1,\widehat{\overline{z}}|n=\sqrt{\frac{þ(n+1)}{2}}|n+1$$
(11)
on the SHO Hilbert space $`_{\text{SHO}}=\{|n,n=0,1,2,\mathrm{}\}`$.
This realization of $`𝒜`$ is not a one-to-one map. Any smooth function of
$$\widehat{n}=\frac{2}{þ}\widehat{\overline{z}}\widehat{z},$$
(12)
such as $`\mathrm{sin}\pi \widehat{n}`$, which vanishes for all $`\widehat{n}𝐙`$ is identical to zero. Similarly, a smooth function of $`\widehat{n}`$ which vanishes at all integers except $`\widehat{n}=n`$ is equivalent to the projective operator $`P_n=|nn|`$.
Intuitively, points on the noncommutative space correspond to states in the Hilbert space. The existence of these projective operators suggests that in some sense the noncommutative space is a disconnected space. Classically, the step function is a projective operator, but it is not a smooth function.
In general, a necessary condition for a Hilbert space $``$ to respect a symmetry $`\widehat{x}\widehat{x}^{}=\widehat{x}^{}(\widehat{x})`$ is that
$$Tr(f(\widehat{x}^{}))=Tr(f(\widehat{x}))$$
(13)
for the trace on $``$. Note that the natural notion for intergration on a quantum space is just the trace over its Hilbert space. <sup>2</sup><sup>2</sup>2To be more precise one should use the Dixmier trace.
Denoting the trace of $`f𝒜`$ on a Hilbert space for the quantum plane by $`f(\widehat{x},\widehat{y})`$, one can check whether it is true that
$$f(\widehat{x}+a,\widehat{y}+b)=f(\widehat{x},\widehat{y})$$
(14)
for arbitrary real constants $`a,b`$ whenever $`f`$ is well defined. This requirement is equivalent to
$$(\widehat{}_if)=0,$$
(15)
which is the quantum analogue of the Stokes’ theorem. It can be checked that the trace over $`_{\text{SHO}}`$ respects both the rotational and the translational symmetry.
### 2.2 GNS Construction
The problem here is generic. Given a noncommutative algebra of functions $`𝒜`$, how do we define its Hilbert space such that certain symmetries of the algebra are respected? The natural solution is the following. First define the integration $``$ on the noncommutative space as a functional invariant under all required symmetries, then use the Gel’fand-Naǐmark-Siegel (GNS) construction to build the Hilbert space from this functional. In the GNS construction, the Hilbert space $`_{\text{GNS}}`$ is roughly speaking the set of all functions $`\{|f(\widehat{x})\}`$, where the inner product is defined by
$$f(\widehat{x})|g(\widehat{x})=f(\widehat{x})^{}g(\widehat{x}).$$
(16)
However, only normalizable states should be considered, and two states are viewed equivalent if they have the same inner product with all other states.
Let $`\widehat{x}_1=\widehat{x},\widehat{x}_2=\widehat{y}`$. The operators $`\widehat{x}_i`$ can be realized in terms of the coordinates $`x_i`$ and derivatives $`_i`$ on a classical plane as
$$\widehat{x}_i=x_i+\frac{i}{2}þ_{ij}_j,$$
(17)
where $`þ_{ij}=þϵ_{ij}`$. Define the “vacuum” $``$ which is annihilated by the derivatives
$$_i=0.$$
(18)
The vacuum correpsonds to the identity function $`1`$ and a classical function $`f`$ acts on it to give a new state $`f(x)`$. The vacuum may not be a state in $``$ and is simply a covenient notation.
For any function $`f`$ of $`\widehat{x}`$, there is a classical function $`f_c`$ of $`x`$ such that
$$f(\widehat{x})=f_c(x)$$
(19)
due to (18). The function $`f_c`$ is also the classical function corresponding to $`f(\widehat{x})`$ in the star product representation (see the next subsection). We can now define the integration $`f(\widehat{x})`$ on the quantum plane by the classical integration of $`f_c(x)`$. Since according to (17), the translation of $`\widehat{x}`$ is identical to the translation of $`x`$, this functional $``$ preserves translational symmetry.
This integral is different from the integral on $`_{\text{SHO}}`$. For instance, any function of $`\widehat{\overline{z}}\widehat{z}`$ which vanishes when $`(þ\widehat{\overline{z}}\widehat{z}/2)`$ is integer is always integrated to zero on $`_{\text{SHO}}`$, but not necessarily zero on $`_{\text{GNS}}`$.
Following the GNS construction with this invariant functional, we obtain the Hilbert space $`_{\text{GNS}}`$ which is the same as the Hilbert space of one-particle quantum mechanics on a classical space.
$`_{\text{GNS}}`$ is not an irreducible representation of $`𝒜`$, and $`_{\text{SHO}}`$ is a subset of this Hilbert space corresponding to picking up only the states $`\{|nH_n(x/\sqrt{þ})e^{x^2/þ}\}`$ where the $`H_n`$’s are the Hermite polynomials.
Due to (18), the requirement (15) is consistent with viewing the functional as the “vacuum expectation value”, with the assumption that
$$\widehat{}_i=0.$$
(20)
It seems that $`_{\text{GNS}}`$ is the correct choice to be used in matrix theory with $`C`$ field background or for D-branes in a $`B`$ field background. In the derivation of the noncommutative relations (1) via open string quantization , the noncommutativity comes from a mixing of classical coordinates and momenta due to mixed boundary conditions for the open string. This is very much like what eq.(17) represents.
### 2.3 Star Product Representation
It is well known that the quantum plane can be described by commutative functions with a star product. Let $`x`$ and $`y`$ be the coordinates over a classical plane. The star product is defined by
$$fg=\left(e^{i\frac{þ}{2}(_x_y_y_x)}(fg)\right),$$
(21)
where $`(AB)`$ is the ordinary commutative product of $`A`$ and $`B`$.
For each classical function $`f(x)`$ we associate a pseudo-differential operator
$$\widehat{f}=\underset{n=0}{\overset{\mathrm{}}{}}i^n\frac{þ_{i_1j_1}}{2}\mathrm{}\frac{þ_{i_nj_n}}{2}\left(_{i_1}\mathrm{}_{i_n}f(x)\right)_{j_1}\mathrm{}_{j_n},$$
(22)
which can be written as
$$\widehat{f}=\left(e^{i\frac{þ}{2}(_x_y_y_x)}(f)\right).$$
(23)
It can be shown that in general
$$\widehat{f}\widehat{g}=\widehat{h},\text{where}h=fg.$$
(24)
Let the Fourier transform of $`f`$ be denoted $`\stackrel{~}{f}_k`$:
$$f(x)=𝑑k\stackrel{~}{f}_ke^{ikx},$$
(25)
then one finds
$$\widehat{f}=𝑑k\stackrel{~}{f}_ke^{ik\widehat{x}},$$
(26)
where $`\widehat{x}`$ is defined in (17). Thus $`\widehat{f}`$ is a function of $`\widehat{x}`$, and the algebra of classical functions with the star product is equivalent to $`𝒜`$. It is straightforward to check from (26) that
$$\widehat{f}(\widehat{x})=f(x).$$
(27)
Comparing this with (19), one sees that $`f=f_c`$.
## 3 An Algebraic Problem
Consider the standard algebraic problem of finding all solutions of $`\psi `$ from an equation like
$$f\psi =0$$
(28)
for given $`f(x)`$. Given a special example of (28), one can obtain other examples by replacing $`f`$ by $`f^{}f`$ and $`\psi `$ by $`\psi g`$ with arbitrary functions $`f^{}`$ and $`g`$.
In noncommutative algebra, an element $`f`$ may not have a left inverse $`f^{}`$ such that $`f^{}f=1`$, or a right inverse $`f^{}`$ such that $`ff^{}=1`$. If $`f`$ has a left inverse, then the unique solution of $`\psi `$ is zero. Conversely, if $`\psi `$ has a trivial solution, it means that $`f`$ can not have a left inverse.
Because the star product involves derivatives, the equation (28) is not really an algebraic relation but rather a differential equation. Differential equations are uniquely solved only if sufficient initial or boundary conditions are specified, so we should expect that there are generically more than one solution to the equation (28).
Consider the example of $`f=z`$:
$$z\psi =z\psi +\frac{þ}{4}\overline{}\psi =0.$$
(29)
It can be easily solved by
$$\psi =e^{\frac{4z\overline{z}}{þ}+h(z)},$$
(30)
where $`h(z)`$ is an arbitrary function of $`z`$. The special case of $`h`$ equal to a constant gives the Gaussian distribution
$$\psi =\psi _0e^{\pm \frac{x^2+y^2}{þ}}.$$
(31)
This function corresponds to the operator $`|00|`$ on $`_{\text{SHO}}`$.
There are also examples of $`f`$ for which the solution of $`\psi `$ in (28) is unique ($`\psi =0`$). An example is given by $`f=e^{ikx}`$.
It is easy to see that if $`f`$ is a polynomial of degree $`k`$ then (28) is a differential equation of order $`k`$, and one needs the appropriate initial or boundary condition for the given differential equation.
In terms of the operator algebra on $`_{\text{SHO}}`$, the fact that (28) has many solutions for $`f=z^k`$ is related to the fact that $`\widehat{z}^k`$ annihilates all states $`|n`$ with $`n<k`$. Instead of using $`\widehat{z}`$ in this discussion, let us normalize it and define
$$u=(\widehat{\overline{z}}\widehat{z}+þ/2)^{1/2}\widehat{z}.$$
(32)
One can check that
$$uu^{}=1,\text{but}u^{}u=1|00|.$$
(33)
Thus
$$p=u^{}u$$
(34)
is a projective operator on $`_{\text{SHO}}`$, i.e., $`pp=p`$. This is because $`u`$ has the annihilation operator $`\widehat{z}`$ on the right to annihilate $`|0`$. Thus $`u`$ has a right inverse $`u^{}`$ and one can prove that it has no left inverse. The existence of this kind of operators $`u`$ is not a special feature of the SHO Hilbert space. The operator $`u`$ has the same properties on $`_{\text{GNS}}`$. In general, $`(u^{})^ku^k`$ for positive integer $`k`$ is a projective operator which annihilates all $`|n`$ with $`n<k`$.
In general, we are interested in operators $`U`$ which satisfy
$$UU^{}=1,\text{but}U^{}U1.$$
(35)
As above, $`U^{}U`$ is always a projective operator. Since $`U`$ does not have a left inverse, a formal gauge transformation of the gauge potential $`A`$ by $`U`$
$$AA^{}=UAU^{}+U(dU^{})$$
(36)
is not really a gauge transformation. We can not transform $`A^{}`$ back to $`A`$ by any transformation. It follows that a transformation by $`U^{}`$ also should not be viewed as a true gauge transformation. However we will show in sec.4.2 that $`U`$ still defines an equivalence relation by modifying the bundle at the same time.
## 4 Projective Module as Twisted Bundle
In noncommutative geometry, vector bundles are defined as projective modules of the algebra $`𝒜`$ of functions on the base space. In this section we review the notion of twisted bundles as projective modules and derive some formulas to be used latter for the ADHM construction.
For a trivial vector bundle of $`N`$ dimensional fibers, a section of the bundle is a vector of $`N`$ entries in $`𝒜`$. The set of sections is a free module $`𝒜^N`$.
To motivate the definition of twisted bundles, we recall the example of the tangent bundle on $`S^2`$. This tangent bundle is nontrivial, but the direct sum of the tangent space and the normal vector space constitute a trivial bundle of three dimensions. One can choose a projective operator acting on this trivial bundle which is the $`3\times 3`$ matrix that projects any 3-vector on a fiber of the trivial bundle to the tangent space. Explicitly, the projection is given by $`P=1|nn|`$, where $`|n`$ is the normal vector on $`S^2`$. Thus sections of the tangent bundle on $`S^2`$ form a projective module, since the definition of a projective module is just that it is a direct summand of a free module $`𝒜^N`$ with finite $`N`$.
Classically all bundles are projective modules and vise versa due to a theorem by Serre and Swan. <sup>3</sup><sup>3</sup>3 The theorem of Serre and Swan states that all locally trivial finite-dimensional complex vector bundles over a compact space $`X`$ are one-to-one corresponding to all finite projective modules over the algebra $`𝒜=C(X)`$. The natural definition of bundles on noncommutative spaces is thus just projective modules . Given a bundle, one can always find an $`N\times N`$ matrix $`P`$ with elements in $`𝒜`$ which satisfies $`PP=P`$, such that the set of sections on the bundle is the projective module $`P𝒜^N`$ for some integer $`N`$. <sup>4</sup><sup>4</sup>4 $`P𝒜^N`$ is a projective module because its direct sum with $`(1P)𝒜^N`$ is a free module. An obvious property of projective modules is that since the projection is taken from the left by convention, it is a right $`𝒜`$ module because one can multiply elements in $`𝒜`$ from the right, to each copy of $`𝒜`$ in $`𝒜^N`$.
In practice it is not always necessary to define a bundle in terms of the projection $`P`$ and free module $`𝒜^N`$. For instance the $`U(N)`$ twisted bundles on a noncommutative torus are constructed in directly without referring to $`P`$ or the trivial bundle.
### 4.1 Connection From Projection
Any connection on the trivial bundle induces a connection on the twisted bundle via this projection. A familiar example is how the tangent bundle of a curved space gets its connection from its embedding in a flat space. We will give explicit expressions for the covariant derivative on the twisted bundle, provided that the projection and free module are given.
A section on a trivial bundle $`_0`$ of dimension $`N_0`$ can be denoted as $`|s(x)=|as_a(x)`$, $`a=1,2,\mathrm{},N_0`$, where $`s_a(x)𝒜`$ are functions on the base space. A section is an element of the free module $`𝒜^{N_0}`$. Let the covariant derivative on $`_0`$ be $`D=d+A`$, where $`A`$ is an $`N_0\times N_0`$ matrix of elements in $`𝒜`$. For simplicity, we first assume that the basis $`\{|a\}`$ is orthonormal and covariantly constant. That is, it satisfies
$`a|b=\delta _{ab},a,b=1,2,\mathrm{},N_0,`$ (37)
$`(D|a)(d|a)+|bA_{ba}(x)=0,`$ (38)
where $`A_{ba}=b|A|a`$. The action of $`D`$ should satisfy the Leibniz rule
$$(D|sf)=(D|s)f+|s(df)$$
(39)
for any section $`|s`$ and function $`f`$. It follows from (38) that
$$(D|s(x))=|a(ds_a(x)).$$
(40)
Consider an arbitrary projection $`P`$ which can be “locally” <sup>5</sup><sup>5</sup>5 What is meant by “locally” is not obvious for noncommutative spaces in general. A special technique is developed in . For the case of instantons it is clear that by this we mean the quantum $`𝐑^4`$, which is a local patch of $`S^4`$. expressed as
$$P=\underset{i=1}{\overset{N}{}}|\psi _i(x)\psi _i(x)|,$$
(41)
where $`N`$ is the rank of the projector ($`NN_0`$), and
$$|\psi _i(x)=|a\psi _{ai}(x)$$
(42)
are sections of $`_0`$. Until the next subsection we assume that
$$\psi _i|\psi _j=\underset{a}{}\psi _{ai}^{}(x)\psi _{aj}(x)=\delta _{ij}$$
(43)
so that
$$PP=P.$$
(44)
The projection $`P`$ defines a projective module $`P𝒜^{N_0}`$ which corresponds to a new bundle $``$. Sections on $``$ are of the form
$$|u(x)=|\psi _iu_i(x)$$
(45)
for $`u_i(x)𝒜`$ ($`i=1,2,\mathrm{},N`$).
The induced covariant derivative on $``$ is
$$D_P=PD.$$
(46)
When acting on $``$,
$`(D_P|u(x))`$ $`=`$ $`(PD|\psi _iu_i)`$ (47)
$`=`$ $`(PD|a\psi _{ai}u_i)`$
$`=`$ $`P|a(d\psi _{ai}u_i)`$
$`=`$ $`|\psi _j\psi _{aj}^{}(d\psi _{ai}u_i)`$
$`=`$ $`|\psi _i(D_Pu)_i,`$
where
$`(D_Pu)_i(x)`$ $`=`$ $`(du_i(x))+(A_P)_{ij}u_j(x),`$ (48)
$`(A_P)_{ij}`$ $`=`$ $`\psi _{ai}^{}(x)(d\psi _{aj}(x))=\psi _i|d|\psi _j.`$ (49)
This is the connection on a twisted bundle obtained as a projection from a free module. Eq.(49) is reminisent of the Berry phase.
If the orthonormal basis $`\{|a\}`$ does not satisfy (38), but rather $`D|a=|bA_{ab}^{}`$, then (49) will be modified by an additional term $`A_{ij}^{}u_j`$ where $`A_{ij}^{}=i|A^{}|j=\psi _{ai}^{}A_{ab}^{}\psi _{bj}`$.
The field strength is by definition
$$F=(dA)+AA$$
(50)
for a given gauge potential $`A`$. Using (41) and (43), the field strength of (49) is
$`(F_P)_{ij}`$ $`=`$ $`(d\psi _{ai}^{})(d\psi _{aj})+\psi _{ai}^{}(d\psi _{ak})\psi _{bk}^{}(d\psi _{bj})`$ (51)
$`=`$ $`(d\psi _{ai}^{})(d\psi _{aj})(d\psi _{ai}^{})\psi _{ak}\psi _{bk}^{}(d\psi _{bj})`$
$`=`$ $`(d\psi _{ai}^{})(\delta _{ab}P_{ab})(d\psi _{bj}),`$
where $`P_{ab}=a|P|b=\psi _{ak}\psi _{bk}^{}`$.
The projective operator $`(1P)`$ in (51) can be written in terms of a basis of vectors as
$$(1P)=\underset{\alpha =1}{\overset{N_0N}{}}|\varphi _\alpha \varphi _\alpha |,$$
(52)
where $`|\varphi _\alpha =|a\varphi _{a\alpha }(x)`$ satisfy
$`\varphi _\alpha |\psi _i=0,`$ (53)
$`\varphi _\alpha |\varphi _\beta =\delta _{\alpha \beta },`$ (54)
for all $`\alpha ,\beta =1,2,\mathrm{},(N_0N)`$, $`i=1,2,\mathrm{},N`$. So $`\{|\psi _i\}`$ is a basis for $``$, and together with $`\{|\varphi _\alpha \}`$ they give a complete basis for $`_0`$. Using (52) and (53), the field strength (51) can be written as
$$(F_P)_{ij}=\psi _{ai}^{}(d\varphi _{a\alpha })(d\varphi _{b\alpha }^{})\psi _{bj}.$$
(55)
In the ADHM construction of instanton solutions, we are given conditions of the form
$$\xi _{a\alpha }^{}\psi _a=0,$$
(56)
where $`\xi _\alpha `$ may not be orthonormal, but can be related to the orthonormal basis $`\varphi _\alpha `$ by a linear transformation
$$\xi _{a\alpha }=\varphi _{a\beta }M_{\beta \alpha },$$
(57)
where $`M`$ satisfies
$$\xi _{a\alpha }^{}\xi _{a\beta }=M_{\alpha \gamma }^{}M_{\gamma \beta }$$
(58)
as a result of (54). The solutions of $`\psi `$ for (56) are labelled as $`\psi _i`$, which defines a projective operator as (41). Thus eq.(56) can be viewed as defining equations for a twisted bundle. Using (56) and (58), we derive the expression for $`F`$ as
$$(F_P)_{ij}=\psi _{ai}^{}(d\xi _{a\alpha })(\xi ^{}\xi )_{\alpha \beta }^1(d\xi _{b\beta }^{})\psi _{bj}.$$
(59)
The expression (59) is particularly useful for the ADHM construction of instanton solutions.
The calculations above are valid even for noncommutative algebras because we have only assumed associativity of the algebra $`𝒜`$.
### 4.2 Projections of Non-Constant Rank
There are two kinds of projective operators depending on whether they have a constant rank through out the space. For instance, for the free module $`𝒜`$ ($`N=1`$) on a classical space which has disconnected pieces, one can define a projection which equals $`1`$ on one piece and equals $`0`$ on another. This projection is a smooth function in $`𝒜`$, so it defines a projective module. However, the corresponding bundle does not have a constant dimension everywhere on the space. It is of one dimension on one piece and zero dimension (no bundle at all) on another. As we have seen in sec.3, for a noncommutative space there can be elements in $`𝒜`$ which are projective even though the classical limit of this space is simply connected. The question is: whether projective modules associated with these projections should be forbidden, or how different they are from others. In this subsection we first give a description of bundles with non-constant dimension, and then show that some of these bundles may be related to bundles with constant dimension via gauge transformations. This suggests that one should not exclude these bundles in noncommutative gauge theories. The appropriate formulation for these bundles is given by Furuuchi , where it is also noted that the $`U(1)`$ instanton solution discovered by Nekrasov and Schwarz in is of this kind. We will have more discussions on it in sec.5.1.
In the previous subsection, we have assumed that the rank of the projection $`P`$ is constant in (43). To be more precise about what we mean, recall that a state in the Hilbert space $``$ can be roughly speaking viewed as a fuzzy point on the noncommutative space. One can evaluate a function in $`𝒜`$ at a fuzzy point by taking the expectation value of it for the corresponding state. If (43) is satisfied, $`|\psi _i`$ is nonzero everywhere, and the rank of $`P`$ is $`N`$ everywhere.
To include the projections with non-constant rank, we have to replace the condition (43) by a weaker one in the derivation above. To simplify our notation, we will omit all indices on matrices. In the following $`\psi `$ represents an $`N_0\times N`$ matrix, and $`\varphi `$ represents an $`N_0\times N^{}`$ matrix. $`N^{}`$ is not smaller than $`(N_0N)`$. Assume that we are given $`\varphi `$ and $`\psi `$ satisfying
$`\varphi \varphi ^{}+\psi \psi ^{}=1,`$ (60)
$`\varphi ^{}\psi =0,`$ (61)
$`\psi ^{}\psi \psi ^{}=\psi ^{},\psi \psi ^{}\psi =\psi .`$ (62)
The first two conditions are the same as (52) and (53) which were used before. The condition (62) insures that
$$p\psi ^{}\psi $$
(63)
is projective, i.e. $`pp=p`$. It is weaker than (43).
Now consider the bundle $``$ composed of sections $`|u=|\psi _iu_i`$ with $`u_i𝒜`$ satisfying
$$pu_i=u_i.$$
(64)
Let $`D=d+A`$ with
$$A=\psi ^{}(d\psi )$$
(65)
as in (49), and define the covariant derivative on $``$ by
$`D_p`$ $``$ $`pDp`$ (66)
$`=`$ $`\psi ^{}d\psi ,`$
where we have used (62). In the above we have used the notation that the action of the exterior derivative is denoted by parentheses, so that $`df=(df)+fd`$ for a function or even differential form $`f`$, and $`dA=(dA)Ad`$ for an odd differential form $`A`$. When acting on sections of the bundle,
$$(D_p|u)=(du)+A_pu,$$
(67)
where
$$A_p=(d\psi ^{})\psi .$$
(68)
Note that we have to define the bundle by (64) in order for the action of $`D_p`$ to be decomposed as $`(d+A_p)`$. This in general will not be possible if one defines the bundle by all sections of the form (45). By imposing (64), the bundle has zero dimension at states which is projected out by $`p`$, so the bundle has non-constant dimension.
The field strength can be calculated as
$`F_p`$ $``$ $`D_p^2`$ (69)
$`=`$ $`\psi ^{}d\psi \psi ^{}d\psi `$
$`=`$ $`\psi ^{}d\psi \psi ^{}((d\psi )+\psi d)`$
$`=`$ $`\psi ^{}((d\psi \psi ^{}(d\psi ))(d\psi )d)+d\psi d`$
$`=`$ $`\psi ^{}((d\psi \psi ^{})(d\psi )+\psi dd)`$
$`=`$ $`((d\psi ^{})(d\psi ^{})\psi \psi ^{})(d\psi )`$
$`=`$ $`\psi ^{}(d\varphi )(d\varphi ^{})\psi ,`$
where we used (60)-(62) and $`dd=0`$. This expression for the field strength is exactly the same as (55). It follows that (59) is also valid when a linear combination of $`\varphi `$ is used.
Suppose we are given $`\varphi `$ and $`\psi `$ satisfying (60)-(62), then we can find new sets of $`\varphi `$ and $`\psi `$ by the formal gauge transformation
$$\psi \psi U,\varphi \varphi ,$$
(70)
where $`U`$ is an $`N\times N`$ matrix of elements in $`𝒜`$ satisfying
$$UU^{}=1.$$
(71)
It follows that all three conditions are still valid, and so by (66) and (69)
$$D_pU^{}D_pU,F_pU^{}F_pU.$$
(72)
If $`U^{}U`$ is not $`1`$, the derivation still goes through, and it means that the $`U`$ transforms sections on the bundle $``$ to the bundle $`p`$. By $`U^{}`$ one can transform sections on $`p`$ back to $``$.
For the special case where $`U`$ is given by the unit $`N\times N`$ matrix times $`u^k`$ with $`u`$ defined in (32), it transforms a bundle of constant dimension to another bundle of dimension zero at $`|n`$ for all $`n<k`$, and dimension $`N`$ everywhere else. This can be viewed as a motivation why one should not exclude bundles of non-constant dimension. However, there are also bundles of non-constant dimension which can not be transformed to those with constant dimension. The $`U(1)`$ instanton discussed below may be such an example. It can be interpreted as a D0-brane on a D4-brane in a constant B field background. Since there are no smooth projective functions on a connected classical space, these bundles can not have a smooth classical limit. For the case of $`U(1)`$ instanton it is suggested that its classical interpretation via the Seiberg-Witten map is a bundle on a Kähler manifold which is a blowup of $`𝐂^2`$ at a finite number of points . It would be interesting to see whether such a correspondence can be established on general grounds, not only for anti-self dual configurations.
For a matter field $`\mathrm{\Phi }`$ in the adjoint representation, $`\mathrm{\Phi }U^{}\mathrm{\Phi }U`$. For any state $`|n`$, it is always possible to find $`U`$ such that after the transformation by $`U`$ the expectation value of $`\mathrm{\Phi }`$ at $`|n`$ is zero. This means that one can choose a gauge such that the field is trivial within any given radius around the origin. This may be a hint of holography in noncommutative gauge theories.
## 5 ADHM Construction
For completeness we review the ADHM construction of instantons on noncommutative space in this section. Explicit expressions for some instanton configurations are given in .
Instantons are defined as anti-self dual configurations of the gauge field. We consider the case where the spacetime coordinates are two commuting copies of (6):
$$[\widehat{z}_A,\widehat{\overline{z}}_B]=\delta _{AB}þ/2,[\widehat{z}_A,\widehat{z}_B]=0,A,B=1,2.$$
(73)
This can always be achieved by linear transformations on the coordinates if $`þ`$ is of maximal rank.
The ADHM construction is a prescription for finding the conditions (56) defining the projective module. For the $`U(N)`$ instanton solution of charge $`k`$, the condition (56) is of the form
$$D_z^{}\psi =0,$$
(74)
where
$$D_z=(\tau _z^{},\sigma _z),$$
(75)
with
$$\tau _z=(B_0\widehat{\overline{z}}_0,B_1+\widehat{\overline{z}}_1,I),\sigma _z^{}=(B_1^{}\widehat{z}_1,B_0^{}\widehat{z}_0,J^{}).$$
(76)
$`D_z`$ plays the role of $`\xi `$ in sec.4. In order to give instanton solutions, the constant matrices $`B_0`$, $`B_1`$, $`I`$ and $`J`$ should be solutions of
$`[B_0,B_0^{}]+[B_1,B_1^{}]+II^{}J^{}J=þ,`$ (77)
$`[B_0,B_1]+IJ=0,`$ (78)
where $`B_0,B_1`$ are $`k\times k`$ matrices and $`I,J^{}`$ are $`k\times N`$ matrices. We can write $`\psi `$ as
$$\psi =\left(\begin{array}{c}\psi _0\\ \psi _1\\ \xi \end{array}\right),$$
(79)
where $`\psi _0,\psi _1`$ are $`k`$-vectors and $`\xi `$ is an $`N`$-vector.
Eq.(74) imposes $`2k`$ constraints on $`\psi `$, which is a $`(2k+N)`$-vector, so there are $`N`$ independent solutions for $`\psi `$, and the projective matrix defined by solutions of $`\psi `$ is a $`(2k+N)\times (2k+N)`$ matrix of rank $`N`$. Hence we can think of the ADHM construction as a recipe to define the $`U(N)`$ instanton bundle as a rank $`N`$ projection from a trivial bundle of dimension $`(2k+N)`$.
In the case of $`U(1)`$ instantons, it is impossible to find $`\psi `$ such that both (60) and (43) hold. In order to use (59), which guarantees that the solution is anti-self dual, one has to choose $`\psi `$ which satisfies (60-62).
### 5.1 $`U(1)`$ Instanton Solution
Following , we consider the $`U(1)`$ instanton solution for the SHO Hilbert space.
For $`N=k=1`$, there is no regular classical instanton solution. On the noncommutative space (73), we can first choose $`B_0=B_1=0`$ by translation, and a solution of (77), (78) is
$$I=\sqrt{þ},J=0.$$
(80)
From (74), one can easily find a solution of this form
$$\psi _0=\widehat{z}_0f,\psi _1=\widehat{z}_1f,\xi =\frac{1}{\sqrt{þ}}(\widehat{\overline{z}}\widehat{z})f,$$
(81)
where the function $`f`$ is arbitrary, and
$$\widehat{\overline{z}}\widehat{z}\widehat{\overline{z}}_0\widehat{z}_0+\widehat{\overline{z}}_1\widehat{z}_1.$$
(82)
The operator $`\widehat{\overline{z}}\widehat{z}`$ has eigenvalue $`\frac{þ}{2}(n_0+n_1)`$ on the state $`|n_0,n_1`$.
In order to use $`\psi `$ to define the projection $`P`$ (41), naively the function $`f`$ should be determined by (43)
$`\psi ^{}\psi `$ $`=`$ $`1`$ (83)
$`=`$ $`f^{}[(\widehat{\overline{z}}\widehat{z}+þ)\widehat{\overline{z}}\widehat{z}/þ]f.`$
In a formal solution of $`f`$ is given as
$$f=f_0\left((\widehat{\overline{z}}\widehat{z}+þ)\widehat{\overline{z}}\widehat{z}/þ\right)^{1/2}.$$
(84)
This expression is in fact ill defined since it diverges on the state $`|0,0`$. However, $`\psi `$ is well defined according to (81) if the factors of $`\widehat{z}_0,\widehat{z}_1`$ are ordered to the right of $`f`$ so that the state $`|0,0`$ is annihilated before it causes any trouble. It is pointed out in that this projective module is the kind that is related to a projection $`p`$ in $`𝒜`$. It is
$$\psi ^{}\psi =p1|0,00,0|$$
(85)
on the SHO Hilbert space, and (83) does not hold. Another way to interpret this is that one needs to make a further projection by $`p`$ to get a new projective module on which (85) is equivalent to (83). As we have mentioned in sec. 4.2, the rank of the gauge group is not a constant on the whole space. It has rank zero at $`|0,0`$ and rank one everywhere else. This special property may be related to its classical interpretation as a bundle on a blowup of $`𝐑^4`$ via the Seiberg-Witten map.
As a generalization of (32), let
$$u_A=(\widehat{\overline{z}}_A\widehat{z}_A+þ/2)^{1/2}\widehat{z}_A,A=0,1,$$
(86)
which are well defined operators satisfying
$$u_Au_A^{}=1.$$
(87)
(Note here that the index $`A`$ is not summed over.) But $`u_A^{}u_A`$ is not equal to $`1`$. The $`þ0`$ limit of $`u_A`$ is the phase factor of $`z_A`$, which is ill-defined on $`𝐂`$. When we try to solve (83), we can also choose $`f`$ to be
$$f=f_0U,U=u_0^mu_1^n$$
(88)
for non-negative integers $`m,n`$. This results in a change of $`\psi `$ by $`\psi \psi U`$, and is thus an example of the situation discussed in sec.4.2.
On the other hand, if we choose $`U`$ to be given by products of $`u_A^{}`$, eq.(43) is satisfied, but the condition (60) is not satisfied for $`\xi =D_z`$. ($`\xi `$ is related to $`\varphi `$ by a linear transformation.) In order to use (59), one has to adjoin new vectors to $`\varphi `$, making it an $`N_0\times (N_0N+1)`$ matrix, and the result is not anti-self dual anymore.
## Acknowledgment
I would like to thank Keun-Young Kim, Bum-Hoon Lee, Kimyeong Lee, Miao Li, Hyun Seok Yang and Piljin Yi for discussions. I also thank APCTP and KIAS where part of this work is done for hospitality. This work is supported in part by the National Science Council, Taiwan, R.O.C. and the Center for Theoretical Physics at National Taiwan University.
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# Study of pseudogap in underdoped cuprate
## Abstract
A mean field SDW analysis of pseudogap in the underdoped cuprates is proposed on the basis of the $`tt^{}U`$ Hubbard model. The prediction of our theory is consistent with the experiment quite well within the uncertainty of the present experimental measurement. Therefore we argue that the pseudogap phenomenon in the underdoped cuprates can be well explained within the mean field approximation.
1. Introduction
A large body of experimental investigations have indicated that underdoped high-temperature superconductors exhibit intriguing properties at temperatures above the superconducting transition temperature $`T_c`$. Most notably, the underdoped cuprates exhibit a pseudogap behavior below a characteristic temperature $`T^{}`$ which can be well above the superconducting transition temperature $`T_c`$. The so-called “pseudogap” means a partial gap. “An example of such a partial gap would be a situation where, within the band theory approximation, some regions of the Fermi surface become gapped while other parts retain their conducting prope rties and with increased doping the gapped portion diminishes and the materials become more metallic”(quoted from Ref. ). What is the origin behind it? A number of scenarios like pair formation well above $`T_c`$ , spin–charge separation , spin-density wave (SDW) or antiferromagnetic fluctuations have been proposed as possible origins of these pseudogap phenomena. However, no consensus has been reached so far, which one is correct in these microscopic theories. It should be noted that these theories of the pseudogap are all beyond the mean-field approximation.
In present paper, we propose a mean field SDW analysis for understanding the pseudogap phenomena in the underdoped cuprates. Our aim is to examine to what extend can we interpret this phenomenon within mean field theory.
2. Electronic band structure
Soon after the discovery of the cuprate superconductors, the electronic band structure of the cuprates has been calculated by the local density approximation (LDA) band calculation . The result of the electronic band structure of the cuprates of the LDA band calculation is consistent with the later angle resolved photoemission experiment . The electronic band structure of the cuprates can be well fitted by a tight-binding model, which is written as
$$\overline{\epsilon }_𝐤=2t(\mathrm{cos}k_x+\mathrm{cos}k_y)4t^{}\mathrm{cos}k_x\mathrm{cos}k_y.$$
(1)
Where $`t`$ is nearest-neighbor, $`t^{}`$ is next-to-nearest-neighbor. In this paper we consider $`t>0`$ and $`t^{}<0`$ only. Energy contour lines for the electronic band structure (1) are shown in Fig. 1. There are two different saddle points locate at the $`\overline{M}`$ points \[($`\pm \pi `$, 0) and (0, $`\pm \pi `$)\] of the Brillouin zone. The energy contour line with energy $`\overline{\epsilon }_s=4t^{}`$ pass through the saddle points.
For convenience, we choice $`\overline{M}`$ as the new origin of the $`𝐤`$-space and take the energy $`\overline{\epsilon }_s=4t^{}`$ at the saddle point $`\overline{M}`$ as zero. Then the dispersion (1) is reexpressed in the form
$`\epsilon _𝐤`$ $`=`$ $`\overline{\epsilon }_𝐤4t^{}`$ (2)
$`=`$ $`2t(\mathrm{cos}k_x+\mathrm{cos}k_y)+4t^{}(\mathrm{cos}k_x\mathrm{cos}k_y1).`$ (3)
If without specific statement, we keep this usage later.
We replot in the period Brillouin zone the energy contour lines passing through the saddle points. As shown in Fig. 2, there are two different regions: $`𝐈+𝐈^{}`$ and $`\mathrm{𝐈𝐈}`$. In the region $`\mathrm{𝐈𝐈}`$, $`\epsilon _𝐤<0`$, in the regions $`𝐈+𝐈^{}`$, $`\epsilon _𝐤>0`$. The area of the region $`𝐈+𝐈^{}`$ is larger than that of the region $`\mathrm{𝐈𝐈}`$. When the region $`\mathrm{𝐈𝐈}`$ is shifted by the vector $`𝐐=(\pi ,\pi )`$, it coincides with the region $`𝐈`$. The region $`𝐈^{}`$ is called as the necklace region, which has following features. Firstly, when $`𝐤`$ locates in a bubble, $`𝐤+𝐐`$ will locate in another one, both $`\epsilon _𝐤`$ and $`\epsilon _{𝐤+𝐐}`$ are larger then zero. On the other hand, in the regions outside the necklace region $`𝐈^{}`$, both sign of the $`\epsilon _𝐤`$ and $`\epsilon _{𝐤+𝐐}`$ are always opposite. For example, when $`𝐤`$ locates in $`𝐈`$, $`\epsilon _𝐤>0`$, then $`𝐤+𝐐`$ will locate in $`\mathrm{𝐈𝐈}`$, $`\epsilon _{𝐤+𝐐}<0`$. Secondly, in the overdoping regime, the Fermi surface entirely lies outside the necklace region (as shown in Fig. 3). But for the underdoping case, only part of the Fermi surface lies outside the necklace region, and further, with decreased doping the portion outside the necklace region increases (as shown in Fig. 3).
It is interest to note the fact that when $`t^{}=0`$, the necklace region and said peculiarity of the band structure of the cuprates disappears.
3. Mean-field theory
The starting point of our calculation is the Hubbard model. In the momentum representation, the $`tt^{}U`$ Hubbard model can be written as
$$H=\underset{𝐤\sigma }{}(\epsilon _𝐤\mu )a_{𝐤\sigma }^{}a_{𝐤\sigma }\frac{U}{2N}\underset{𝐪}{}\underset{𝐤\sigma 𝐤^{}\sigma ^{}}{}a_{𝐤+𝐪\sigma }^{}\sigma a_{𝐤\sigma }a_{𝐤^{}\sigma ^{}}^{}\sigma ^{}a_{𝐤^{}+𝐪\sigma ^{}}.$$
(4)
Here a term $`\frac{1}{2}NU`$ has been omitted. $`U`$ is the local Coulomb repulsion. $`a_{𝐤\sigma }(a_{𝐤\sigma }^{})`$ is the annihilation (creation) operator for the electron with momentum $`𝐤`$ and spin $`\sigma `$. $`\mu `$ is the chemical potential. $`\epsilon _𝐤`$ is given by Eq. (3). All the momentum summations extend over the Brillouin zone. Considering commensurate SDW state and using the mean-field approximation, the Hamiltonian reduces
$$H=\underset{𝐤\sigma }{}^{}(\epsilon _𝐤\mu )a_{𝐤\sigma }^{}a_{𝐤\sigma }+\underset{𝐤\sigma }{}^{}(\epsilon _{𝐤+𝐐}\mu )a_{𝐤+𝐐\sigma }^{}a_{𝐤+𝐐\sigma }\mathrm{\Delta }\underset{𝐤\sigma }{}^{}(a_{𝐤+𝐐\sigma }^{}\sigma a_{𝐤\sigma }+h.c).$$
(5)
Here $`_𝐤^{}`$ means that the sum extends over the magnetic Brillouin zone (shown in Fig. 4 by the thick square). The term $`\frac{N}{2U}\mathrm{\Delta }^2`$ has been omitted. The order parameter $`\mathrm{\Delta }`$ is given by
$$\mathrm{\Delta }=\frac{2U}{N}\underset{𝐤\sigma }{}^{}<a_{𝐤+𝐐\sigma }^{}\sigma a_{𝐤\sigma }>.$$
(6)
By the following canonical transformation
$`\alpha _{𝐤\sigma }`$ $`=`$ $`u_𝐤a_{𝐤\sigma }v_𝐤\sigma a_{𝐤+𝐐\sigma },`$ (7)
$`\gamma _{𝐤\sigma }`$ $`=`$ $`v_𝐤\sigma a_{𝐤\sigma }+u_𝐤\sigma a_{𝐤+𝐐\sigma },`$ (8)
the Hamiltonian (5) is diagonalised as
$$H=\underset{𝐤\sigma }{}^{}(\epsilon _1(𝐤)\alpha _{𝐤\sigma }^{}\alpha _{𝐤\sigma }+\epsilon _2(𝐤)\gamma _{𝐤\sigma }^{}\gamma _{𝐤\sigma }),$$
(9)
in which,
$$\epsilon _1(𝐤)=\frac{\epsilon _𝐤+\epsilon _{𝐤+𝐐}}{2}\mu +\sqrt{(\frac{\epsilon _𝐤\epsilon _{𝐤+𝐐}}{2})^2+\mathrm{\Delta }^2},$$
(10)
$$\epsilon _2(𝐤)=\frac{\epsilon _𝐤+\epsilon _{𝐤+𝐐}}{2}\mu \sqrt{(\frac{\epsilon _𝐤\epsilon _{𝐤+𝐐}}{2})^2+\mathrm{\Delta }^2},$$
(11)
$$\mathrm{\Delta }=\frac{U}{N}\underset{𝐤}{}^{}\frac{\mathrm{\Delta }}{E(𝐤)}(\mathrm{tanh}(\frac{\epsilon _1(𝐤)}{2T})\mathrm{tanh}(\frac{\epsilon _2(𝐤)}{2T}))$$
(12)
and
$$E(𝐤)=\sqrt{(\frac{\epsilon _𝐤\epsilon _{𝐤+𝐐}}{2})^2+\mathrm{\Delta }^2}.$$
(13)
Here $`\epsilon _1(𝐤)`$ and $`\epsilon _2(𝐤)`$ are energy dispersions of the quasiparticles. For the hole doping system, the Fermi surface lies inside the lower band ($`\epsilon _2(𝐤)`$). The pseudogap is given by
$`\mathrm{\Delta }_{PS}(\varphi )`$ $`=`$ $`|\epsilon _2(𝐤)|`$ (14)
$`=`$ $`\mu 4t^{}(\mathrm{cos}k_x\mathrm{cos}k_y1)+\sqrt{4t^2(\mathrm{cos}k_x\mathrm{cos}k_y)^2+\mathrm{\Delta }^2}.`$ (15)
In Fig. 5 we plot the part of the magnetic Brillouin zone of the Fig. 4. The light curve represents the Fermi surface. $`k_x`$\- and $`k_y`$-axis are parallel with $`\overline{M}\mathrm{\Gamma }`$ and $`\overline{M}`$X, respectively. In Eq. (15), $`𝐤=(k_x,k_y)`$ is the wave vector of the Fermi surface, i.e. $`\epsilon _𝐤\mu =0`$. $`\varphi =\mathrm{arctan}(k_x/k_y)`$ is polar angle of the wave vector $`𝐤`$. For convenience, we take $`\varphi ^{}=\mathrm{arctan}(k_x/(\pi k_y))`$ as variable instead of the $`\varphi `$ in the following calculations.
4. Results
In this section, we analyse the angular dependence of the pseudogap $`\mathrm{\Delta }_{PS}(\varphi ^{})`$ along the Fermi surface. Ouing to the symmetry of the energy spectrum $`\epsilon _2(𝐤)`$, our analysis can be limited only in the interval $`0\varphi ^{}45^{}`$.
By solving Eqs. (10), (11), (12), (13) and (15) numerically, we compute $`\mathrm{\Delta }_{PS}(\varphi ^{})`$ at $`T=0`$ K in the underdoping regime ($`\mu >0`$). In the computation, we choose $`t=430`$ meV, $`t^{}/t=0.18`$, $`U/t=0.8`$ and the hole doping concentration x=0.13. The results are plotted as $`\mathrm{\Delta }_{PS}(\varphi ^{})`$ versus $`\varphi ^{}`$ curve in Fig. 6. It shows that there is strong angular dependence of the $`\mathrm{\Delta }_{PS}(\varphi ^{})`$, as one moves along the Fermi surface from $`\varphi ^{}=0`$ (i.e. near the saddle point $`\overline{M}`$, or at the hot spot) to $`\varphi ^{}=45^{}`$ (i.e. cold spot). At first, we see the maximum pseudogap at $`\varphi ^{}=0`$. As we go into the necklace region, the pseudogap drops quickly and, at approximately $`18^{}`$, drops down to 2 meV. And then, the pseudogap decreases monotonously to $`\mathrm{\Delta }_{PS}(45^{})`$. Experimental measurement reveals that only a portion of the Fermi surface near the saddle point $`\overline{M}`$ becomes gapped while in other parts, the pseudogap is equal to zero . However, the error-bar of the pseudogap data is rather larger<sup>*</sup><sup>*</sup>*<sup>)</sup>See, for example, Fig. 8 of paper <sup>)</sup>. It is impossible to say certainly that along the part of the Fermi surface near the cold point, the pseudogap is real zero or only a small quantity. Keeping this fact in mind, we conclude that the general structure of the pseudogap along the Fermi surface, shown in Fig. 6, captures the main feature of experiment <sup>)</sup>Ref. also the review article and the papers listing in it<sup>)</sup>.
The dependence of $`\mathrm{\Delta }_{PS}(0)`$ on the hole doping concentration is shown in Fig. 7. It shows that $`\mathrm{\Delta }_{PS}(0)`$ increase with the decrease of hole doping. In Fig. 8, we plot $`\mathrm{\Delta }\varphi ^{}`$ versus the hole doping concentration curve. Here, $`\mathrm{\Delta }\varphi ^{}`$ is the interval $`\varphi ^{}`$ (measured from $`\varphi ^{}=45^{}`$), defined by the requirement that the value $`\mathrm{\Delta }_{PS}(\varphi ^{})`$ is less than a proper chosen value (say, 2 meV in Fig. 6 and 8). It can be seen from Fig. 8 that the length of the Fermi arc, along which the pseudogap less than 2 meV, increase with the increase of doping. It implies that as doping increase, the portion of the Fermi surface destroyed by the pseudogap decreases. The prediction discribed above is consistent with the experiment <sup>2)</sup>.
5. Concluding remarks
It is of interest to note that the situation is entirely different if $`t^{}=0`$. For in this case, Eq. (15) reduces to
$$\mathrm{\Delta }_{PS}(\varphi ^{})=\mu +\sqrt{\mu ^2+\mathrm{\Delta }^2}.$$
(16)
It is in contradiction with the experiment , for the pseudogap along the Fermi surface, according to (16), is constant.
Now, it is clearly that the peculiarity of the band structure of the cuprate plays an important role in understanding the pseudogap phenomenon in underdoped cuprate. This is the reason why our mean field SDW analysis of the pseudogap, based on the $`tt^{}U`$ Hubbard model, meets with success.
The mean-field solution has an antiferromagnetic long-range order. At sufficient doping concentration, the spin long-range order will be removed by fluctuations but there are still short-range orderings. We assume implicitly in our theory that the pseudogap structure, at least near the saddle point ($`\pi ,0`$), is not sensitive to the long-range order and will survive in underdoped region, leading to the pseudogap phenomenon.
Acknowledgments
One of the authors (H. S. Wu) would like to thank Prof. Z. Y. Weng for very valuable discussion.
Figure Captions:
Fig. 1. The Brillouin zone and energy contour lines: The $`\mathrm{\Gamma }`$ point is at the middle of the Brillouin zone and the $`\overline{M}`$ points \[($`\pm \pi `$, 0) and (0, $`\pm \pi `$)\] are midway along the edges. The curves are the energy contour lines ($`t^{}/t=0.16`$) with energy $`\overline{\epsilon }_𝐤/t=1.59,0.64,0.49`$ and $`0.4`$, which are from inside to outside.
Fig. 2. The period Brillouin zone: The solid curves are the energy contour lines with $`\epsilon _𝐤=0`$. When the region $`\mathrm{𝐈𝐈}`$ is shifted by the vector $`𝐐=(\pi ,\pi )`$, it coincides with the region $`𝐈`$. In regions $`𝐈+𝐈^{}`$, $`\epsilon _𝐤>0`$. In the region $`\mathrm{𝐈𝐈}`$, $`\epsilon _𝐤<0`$. The region $`𝐈^{}`$ is called as the necklace region.
Fig. 3. The Fermi surfaces ($`t^{}/t=0.16`$) in the quarter of Brillouin zone: The light curve represents the Fermi surface. 1 and 2 for the overdoping. 3 and 4 for the underdoping. The heavy dashed and the solid curve represent the necklace region boundary.
Fig. 4. Our choice of the Brillouin zone. The heavy rectangle is our choice of the Brillouin zone boundary and the origin is at the $`\overline{M}`$. The heavy square is the magnetic Brillouin zone boundary.
Fig. 5. This figure is the part of the magnetic Brillouin zone of the Fig. 4. The light curve represents the Fermi surface for the underdoping region ($`t^{}/t=0.16`$). The dashed and the solid curve represent the necklace region boundary. The heavy solid lines are the magnetic Brillouin zone boundary. $`k_x`$\- and $`k_y`$-axis are parallel with $`\overline{M}\mathrm{\Gamma }`$ and $`\overline{M}`$X, respectively. $`\varphi =\mathrm{arctan}(k_x/k_y)`$. $`\varphi ^{}=\mathrm{arctan}(k_x/(\pi k_y))`$.
Fig. 6. The angle dependence of the pseudogap located at $`\varphi ^{}`$, $`\mathrm{\Delta }_{PS}(\varphi ^{})`$, for the hole doping concentration x=0.13 ($`t^{}/t=0.18`$ and $`U/t=0.8`$). The $`\mathrm{\Delta }\varphi ^{}`$ is the region where the values of $`\mathrm{\Delta }_{PS}(\varphi ^{})`$ are all smaller than 2 meV.
Fig. 7. The hole doping concentration dependence of the pseudogap located at $`\varphi ^{}=0`$, $`\mathrm{\Delta }_{PS}(0)`$, for $`t^{}/t=0.18`$ and $`U/t=0.8`$. The x indicates the hole doping concentration. The solid curve represents the pseudogap in the underdoping region.
Fig. 8. The hole doping concentration dependence of the region where the values of $`\mathrm{\Delta }_{PS}(\varphi ^{})`$ are all smaller than 2 meV, $`\mathrm{\Delta }\varphi ^{}`$, for the underdoping region ($`t^{}/t=0.18`$ and $`U/t=0.8`$). The x indicates the hole doping concentration.
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# Neutrinoless double beta decay with and without Majoron-like boson emission in a 3-3-1 model
## I Introduction
The issue of neutrino masses continues to be a golden plate in elementary particle physics. Although data coming from solar , atmospheric , and the accelerator LSND neutrino experiments strongly suggest that neutrinos must be massive particles, direct measurements did not obtain any positive result .
It is a very well known fact that if neutrinos are massive Majorana particles it should exist the neutrinoless double beta $`(\beta \beta )_{0\nu }`$ decay . If the neutrino mass is the main effect that triggers this decay, the decay lifetime is proportional to (for the case of light neutrinos)
$$M_\nu =\underset{i}{}U_{ei}^2m_{\nu _i},$$
(1)
where $`U_{ei};i=1,2,3`$ denote the elements of a mixing matrix that relates symmetry $`\nu _\alpha ;\alpha =e,\mu ,\tau `$ and mass eigenstates $`\nu _i`$ through the relation $`\nu _\alpha =_iU_{\alpha i}\nu _i`$; and $`m_{\nu _i}`$ are the neutrino masses. Experimentally a half life limit $`T_{1/2}^{0\nu }>1.8\times 10^{25}`$ yr implies
$$M_\nu <0.2\text{eV}.$$
(2)
The important point is that the $`(\beta \beta )_{0\nu }`$ decay probes the physics beyond the standard model. In particular the observation of this decay would be an evidence for a massive Majorana neutrino although it could say nothing about the value of the mass. This is because although right-handed currents and/or scalar-bosons may affect the decay rate, it has been shown that whatever the mechanism of this decay is, a nonvanishing neutrino mass is required for the decay to take place . However, this does not mean that the neutrino mass is necessarily the main factor triggering this decay. In some models the $`(\beta \beta )_{0\nu }`$ decay can proceed with arbitrary small neutrino mass via scalar boson exchange . The mechanism involving a trilinear interaction of the scalar bosons was proposed in Ref. in the context of a model with $`SU(2)U(1)`$ symmetry with doublets and a triplet of scalar bosons. However, since in this type of models there is no large mass scale, it was shown in Refs. that the contribution of the trilinear interactions are in fact negligible. In general, in models with that symmetry, a fine tuning is needed if we want that the trilinear terms give important contributions to the $`(\beta \beta )_{0\nu }`$ decay .
Here we will show that in a model with gauge symmetry $`SU(3)_cSU(3)_LU(1)_N`$ (3-3-1 model by short) , which has a rich Higgs bosons sector as in the multi-Higgs extensions of the standard model, there are new contributions to the $`(\beta \beta )_{0\nu }`$ decay. However, unlike the latest sort of models, a fine tuning of the parameters of the 3-3-1 model is not necessary since some trilinear couplings, which have mass dimension, could imply an enhancement of the respective amplitudes (See. Sec. III).
We will use the following strategy: First, we consider the several new contributions to the $`(\beta \beta )_{0\nu }`$ decay introduced by the 3-3-1 model. Next, once this decay has not experimentally been seen, we will consider the usual standard model amplitude (that would arise with massive Majorana neutrinos) as the reference one and make the assumption that all the new amplitudes are at most as contributing as this one. Hence, we can obtain constraints on some typical mass scale 3-3-1 parameters. The new contributions to the $`(\beta \beta )_{0\nu }`$ decay are of the short range type . Since the respective matrix elements are different from those of the long range contributions (the exchange of a light-Majorana neutrino) our results should be considered only as an indication of the possible large contributions to this decay in the context of the 3-3-1 model.
The outline of the paper is the following. In Sec. II we introduce the interactions which are relevant to the present study. The model with $`\sigma _1^00`$, which is some cases it has a Majoron-like Goldstone boson, is also discussed. In Sec. III we consider the more important contributions to the $`(\beta \beta )_{0\nu }`$ decay and the constraints upon some masses of the model. In Sec. IV we show that if we add a neutral scalar singlet to the minimal model a Majoron-like Goldstone boson is consistent with the $`Z^0`$ invisible width and we also discuss briefly the Majoron emission process $`(\beta \beta )_{0\nu M}`$ comparing the relative strength of two amplitudes. Our conclusions remain in the last section.
## II The model
Here we will consider the 3-3-1 model with the leptons belonging to triplets $`(\nu _lll^c)^T;l=e,\mu ,\tau `$ and in which a sextet of scalar bosons
$$S=\left(\begin{array}{ccc}\sigma _1^0& \frac{h_2^{}}{\sqrt{2}}& \frac{h_1^+}{\sqrt{2}}\\ \frac{h_2^{}}{\sqrt{2}}& H_1^{}& \frac{\sigma _2^0}{\sqrt{2}}\\ \frac{h_1^+}{\sqrt{2}}& \frac{\sigma _2^0}{\sqrt{2}}& H_2^{++}\end{array}\right)(\mathrm{𝟔},\mathrm{𝟎}).$$
(3)
is necessary to give to the charged leptons a mass if $`\sigma _2^0v_{\sigma _2}0`$ .
Most of the phenomenological studies of the model has been done by considering $`\sigma _1^0v_{\sigma _1}=0`$. The case when $`\sigma _1^00`$ it was considered in Ref. , where the other scalar multiplets are explicitly given. The main difference in the latter case with respect to the former one is that there is a mixing between the vector bosons $`W^+`$ and $`V^+`$:
$$\left(\begin{array}{cc}W_\mu ^+& V_\mu ^+\end{array}\right)\left(\begin{array}{cc}M_W^2& \delta \\ \delta & M_V^2\end{array}\right)\left(\begin{array}{c}W_\mu ^{}\\ V_\mu ^{}\end{array}\right),$$
(4)
where $`\delta =(g^2/2)(2v_{\sigma _1}v_{\sigma _2})`$, $`M_W^2`$ and $`M_V^2`$ are the mass eigenvalues when $`\delta =0`$; if $`\delta 0`$ (i.e. when $`v_{\sigma _1}0)`$ the mass of the physical fields are now
$$2M_{1,2}=(M_W^2+M_V^2)\pm \left[(M_W^2M_V^2)^2+4\delta ^2\right]^{1/2}$$
(5)
and we have defined $`M_W^2=(g^2/2)(v_\eta ^2+v_\rho ^2+v_{\sigma _2}^2+v_{\sigma _1}^2)`$, $`M_V^2=(g^2/2)(v_\eta ^2+v_\chi ^2+v_{\sigma _2}^2+v_{\sigma _1}^2)`$ and $`g`$ is the $`SU(3)_L`$ coupling constant which is numerically equal to the coupling of $`SU(2)_L`$ i.e., $`g^2=8M_W^2G_F/\sqrt{2}`$. We have denoted by $`v_\eta `$, $`v_\rho `$ and $`v_\chi `$ the vacuum expectation values of the neutral components of the triplets. Notice that $`M_1M_W`$ and $`M_2M_V`$ when $`\delta 0`$. The vector bosons $`W_\mu ^+`$ and $`V_\mu ^+`$ are related to the new mass eigenstates $`W_{1\mu }^+`$ and $`W_{2\mu }^+`$ as
$$\left(\begin{array}{c}W^+\\ V^+\end{array}\right)=\left(\begin{array}{cc}c_\theta & s_\theta \\ s_\theta & c_\theta \end{array}\right)\left(\begin{array}{c}W_1^+\\ W_2^+\end{array}\right)$$
(6)
with $`\mathrm{tan}2\theta =2\delta /(M_W^2M_V^2)`$. We can obtain an upper bound on $`\delta `$ by assuming that the main contribution to the $`M_W^2`$ mass is given by $`v_{\sigma _2}246`$ GeV and that $`v_{\sigma _1}`$ has its maximum value 3.89 GeV allowed by the value of the $`\rho `$-parameter . In fact if $`v_{\sigma _2}`$ were the main contribution to the $`M_W`$ mass we would have $`\delta /M_W^2(2v_{\sigma _1}/v_{\sigma _2})<0.032`$. The constraint on the mixing angle $`\theta `$ is:
$$0s_\theta ^2=\frac{1}{2}\left(1\frac{M_V^2M_W^2}{[4\delta ^2+(M_V^2M_W^2)^2]^{1/2}}\right)<\frac{1}{2}.$$
(7)
Some illustrative values for $`s_\theta `$ are obtained by using typical values for the parameters. For instance, for $`v_{\sigma _1}=3.89`$ GeV; $`v_{\sigma _2}=10`$ GeV, $`M_W=80.41`$ GeV and $`M_V=100(300)`$ GeV we get $`s_\theta ^2=1.9\times 10^5(3.4\times 10^8)`$; or if $`M_V=100`$ GeV and if $`v_{\sigma _2}`$ has its maximal value $`v_{\sigma _2}=246`$ GeV we have $`s_\theta ^2=1.1\times 10^2`$. We see that only for values of $`M_VM_W`$ the $`s_\theta ^2`$ is almost 0.5 but this light vector boson may be not phenomenologically safe. However if $`v_{\sigma _1}`$ is of the same order of magnitude of the neutrino mass smaller values for the mixing angle are obtained. Hence, it may be no relevant for the collider physics and low energy processes like the $`(\beta \beta )_{0\nu }`$ decay at all and in practice $`W_1^+W^+`$, $`W_2^+V^+`$; but this could not be the case in astrophysical processes .
Next, we consider the several interactions that are present in this model. The scalar-quark interactions are
$`_Y^{ud}`$ $`=`$ $`{\displaystyle \frac{\sqrt{2}}{|v_\rho |}}\overline{D}_LV_{\mathrm{CKM}}^{}M^uU_R\rho ^{}+{\displaystyle \frac{\sqrt{2}}{|v_\eta |}}\overline{U}_LV_{\mathrm{CKM}}M^dD_R\eta _1^+`$ (10)
$`+\overline{D}_L\left(V_L^d\right)^T\mathrm{}\text{ }V_L^uM^uU_R\left[{\displaystyle \frac{\sqrt{2}}{|v_\eta |}}\eta _1^{}{\displaystyle \frac{\sqrt{2}}{|v_\rho |}}\rho ^{}\right]`$
$`+\overline{U}_L\left(V_L^u\right)^T\mathrm{}\text{ }V_L^dM^dD_R\left[{\displaystyle \frac{\sqrt{2}}{|v_\rho |}}\rho ^+{\displaystyle \frac{\sqrt{2}}{|v_\eta |}}\eta _1^+\right]`$
$`+`$ $`H.c.,`$ (11)
with
$$\mathrm{\Delta }\left(\begin{array}{ccc}0& 0& 0\\ 0& 0& 0\\ 0& 0& 1\end{array}\right),$$
(12)
and $`V_L^{u,d}`$ are unitary mixing matrices, and $`M^{u,d}`$ are the diagonal mass matrices of the $`u`$-like and $`d`$-like quark sectors, and $`V_{\mathrm{CKM}}`$ denotes the usual mixing matrix of Cabibbo-Kobayashi-Maskawa.
The Yukawa interactions in the lepton sector are
$`_Y^l`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\overline{\nu _L}𝒦_1l_RH_1^++{\displaystyle \frac{1}{\sqrt{2}}}\overline{l}_L𝒦_2\nu _R^cH_2^{}`$ (13)
$`+`$ $`{\displaystyle \frac{1}{2}}\overline{l_L}𝒦_3(l_L)^cH_1^{}+{\displaystyle \frac{1}{2}}\overline{(l^c)_L}𝒦_4l_RH_2^{++}`$ (14)
$`+`$ $`2\overline{\nu _L}𝒦_1^{}l_R\eta _1^+2\overline{l_L}𝒦_2^{}\nu _R^c\eta _2^{}+H.c.,`$ (15)
where $`𝒦_1=E_L^\nu GE_R`$, $`𝒦_2=E_L^{}GE_L^\nu `$; $`𝒦_3=E_L^{}GE_L`$, $`𝒦_4=E_R^TGE_R`$; $`𝒦_1^{}=E_L^\nu G^{}E_R`$, $`𝒦_2^{}=E_L^{}G^{}E_L^\nu `$; $`G`$ and $`G^{}`$ are symmetric and antisymmetric (they can be complex) matrices, respectively. $`E_R,E_L,E_L^\nu `$ are the right- and left-handed mixing unitary matrices in the lepton sector relating symmetry eigenstates (primed fields) with mass eigenstates (unprimed fields) :
$$l_R^{}=E_Rl_R,l_L^{}=E_Ll_L,\nu _L^{}=E_L^\nu \nu _L,$$
(16)
Some of the couplings in Eq. (15) do not depend on the charged lepton masses and since all matrices in Eq. (15) are not unitary, the model breaks the lepton universality but it can be shown that, for the massless neutrino case no strong constraints arise from exotic muon and tau decays .
In the scalar sector we have also mixing angles. In the singly charged sector we have $`\varphi _i^{}=_lO_{ij}H_j^{}`$, where $`\varphi _i^{}=\eta _1^{},\eta _2^{},\rho ^{},\chi ^{},h_1^{},h_2^{}`$ and $`H_j^{},j=1,\mathrm{},6`$ denotes the respective mass eigenstate field; similarly in the doubly charged sector we have $`\mathrm{\Phi }_i^{}=_l𝒪_{ij}\mathrm{\Psi }_j^{}`$, with $`\mathrm{\Psi }_i^{}=\rho ^{},\chi ^{},H_1^{},H_2^{}`$ and $`\mathrm{\Psi }_j^{},j=1,..,4`$ the respective mass eigenstate fields. However, in the following we will use $`H^{}`$ and $`H^{}`$ as typical mass eigenstates of the respective charged fields and omit the scalar mixing parameters.
We recall that the model conserves the $`=L+B`$ quantum number; $`L`$ is the total lepton number and $`B`$ is the baryon number. The assignments are:
$`(U^{})`$ $`=`$ $`(V^{})=(\rho ^{})=(\chi ^{})=(\eta _2^{})`$ (17)
$`=`$ $`(\chi ^{})=(\sigma _1^0)=(h_2^{})=(H_1^{})`$ (18)
$`=`$ $`(H_2^{})=2,`$ (19)
and all other scalar fields with $`=0`$.
The charged currents coupled to the vector bosons are given by
$`^{CC}`$ $`=`$ $`{\displaystyle \frac{g}{\sqrt{2}}}(\overline{U}_L\gamma ^uV_{\mathrm{CKM}}D_LW_\mu ^+\overline{\nu }_L\gamma ^\mu V_Wl_LW_\mu ^+`$ (20)
$`+`$ $`\overline{l_L^c}\gamma ^\mu V_UV_W^{}\nu _LV_\mu ^+\overline{l_L^c}\gamma ^\mu V_Ul_LU_\mu ^{++})+H.c,`$ (21)
with the mixing matrices defined as $`V_{\mathrm{CKM}}=\left(V_L^u\right)^{}V_L^d`$ in the quark sector; and $`V_W=E_{L}^{\nu }{}_{}{}^{}E_L`$, $`V_U=E_R^TE_L^\nu `$ in the leptonic sectors.
We have the trilinear interactions involving one vector- and two scalar-bosons which are of the form (up to a $`ig/\sqrt{2}`$ factor):
$`^{V2S}`$ $`=`$ $`^\mu \chi ^+\chi ^{}W_\mu ^++\chi ^{}^\mu \chi ^{++}W_\mu ^{}+^\mu h_2^+H_1^{}W_\mu ^+`$ (22)
$`+`$ $`^\mu H_1^{++}h_2^{}W_\mu ^{}+^\mu \rho ^{}\rho ^{++}V_\mu ^{}+^\mu \rho ^{}\rho ^+V_\mu ^+`$ (23)
$`+`$ $`h_1^+^\mu H_2^{}V_\mu ^++^\mu h_1^{}H_2^{++}V_\mu ^{}+(\eta _2^+^\mu \eta _1^+`$ (24)
$`+`$ $`\eta _1^+^\mu \eta _2^++h_2^+^\mu h_1^++h_1^+^\mu h_2^+)U^{}_\mu +H.c..`$ (25)
There are also trilinear interactions involving two vector- and one scalar-bosons (the $`\chi ^{}`$ scalar couples to ordinary and exotic quarks and for this reason it is not of our concern here). They are given by (up to a $`g^2/2`$ factor)
$`^{2VS}`$ $`=`$ $`{\displaystyle \frac{v_{\sigma _1}}{2}}H_1^{++}W^{}W^{}+{\displaystyle \frac{v_{\sigma _1}}{2}}H_2^{++}V^{}V^{}`$ (26)
$`+`$ $`\left(v_\rho \rho ^{++}+v_\chi \chi ^{++}+{\displaystyle \frac{v_{\sigma _2}}{2}}(H_1^{++}+H_2^{++})\right)`$ (28)
$`W^{}V^{}.`$
Notice that there is a coupling which is proportional to $`v_\chi `$ and hence it will be the dominant one.
Next, we write down the trilinear interactions among three vector bosons
$`^{3V}`$ $`=`$ $`i{\displaystyle \frac{g}{\sqrt{2}}}(W_{\mu \nu }^+V^{\nu +}U^\nu +V_{\mu \nu }^+V^{\nu +}U^\mu `$ (29)
$`+`$ $`W^{\nu +}V^{\mu +}U_{\mu \nu }^{})`$ (30)
where $`W_{\mu \nu }=_\mu W_\nu _\nu W_\mu `$ and so on.
Finally, we have trilinear interactions among scalar-bosons only
$$^{3S}=\frac{f_1}{2}ϵ^{ijk}\eta _i\rho _j\chi _k+\frac{f_2}{2}\chi ^TS^{}\rho +H.c..$$
(31)
The couplings $`f_{1,2}`$ have dimension of mass but are both of them arbitrary parameters (see next section). Other terms like the trilinears $`f_3\eta ^TS^{}\eta `$ and $`f_4ϵSSS`$ and the quartic interactions $`ϵ\chi (S\eta ^{})\rho `$, $`\chi ^{}\eta \rho ^{}\eta `$ and $`ϵ\chi \rho SS`$ violate the conservation of $``$.
However, as we will show in Sec. IV, when discussing the Majoron emission, the model must be modified by adding a scalar singlet in order to be consistent with the LEP data.
## III The neutrinoless double beta decay
Some of the more relevant diagrams of the $`(\beta \beta )_{0\nu }`$ decay in the present model are shown in Figs. 1-6. Our goal is to analyze the order of magnitude of each diagram and to obtain constraints on some mass scales of the model. We will consider the diagram in Fig. 1 as the reference one, i.e., it is the diagram that already exist in the standard model framework with massive Majorana neutrinos and which is parameterized by two effective four-fermion interactions. The other contributions will be considered as been at most equally important than the standard one.
The strength of the diagram in Fig. 1 is given by
$$A(1)\frac{g^4M_\nu }{M_W^4p^2}c_\theta ^4=\frac{32G_F^2M_\nu }{p^2}c_\theta ^4,$$
(32)
where $`M_\nu `$ is the effective mass defined in Eq. (1) and $`p^2`$ is the average of the four-momentum transfer squared, which is of the order of $`(100\text{MeV})^2`$. Below we will use a small $`\delta `$ so that $`M_1M_W`$ and $`M_2M_V`$.
In Eq. (32) and hereafter we will omit for simplicity the mixing parameters. Only in the vertices we will take care about the mixing between $`W`$ and $`V`$ defined in Eq. (6) but in the propagator we will use the masses of $`W`$ and $`V`$.
Next, let us consider the diagram in Fig. 2 which has the strength given by
$$A(2)32G_F^2\left(\frac{M_W}{M_V}\right)^2\frac{c_\theta ^3s_\theta }{\sqrt{p^2}},$$
(33)
and we have the ratio
$$\frac{A(2)}{A(1)}=\left(\frac{M_W}{M_V}\right)^2\frac{\sqrt{p^2}}{M_\nu }\mathrm{tan}\theta ,$$
(34)
and if $`A(2)/A(1)<1`$ we have that
$$M_V>2.2\times 10^4M_W\sqrt{\mathrm{tan}_\theta }=1.79\times 10^6\sqrt{\mathrm{tan}_\theta }\mathrm{GeV}.$$
(35)
We recall that a lower limit of 440 GeV is obtained for $`M_V`$ from the muon decays but when only the bilepton contributions to those decays are considered . However, in the minimal 3-3-1 model the scalar-boson contributions cannot be negligible since some of the charged scalar-bosons can be lighter than the vector bilepton boson $`V^{}`$. Hence, a lighter vector boson $`V`$ may still be possible but this subject deserves a more detailed study of the muon decay considering both vector and scalar contributions. A contribution similar to that in Fig. 1 but with two $`V^{}`$ bosons instead of two $`W^{}`$ bosons may be not negligible but it does not constraint the mass $`M_V`$ as much as those in Eq. (35) since the condition that its ratio to the $`A(1)`$ amplitude be less than one gives the condition $`M_V>M_W\sqrt{\mathrm{tan}\theta }`$.
All the Lagrangian interactions in Eqs. (11), (15), (21), (25), (28) and (30) are written in terms of symmetry eigenstates. We have assumed Yukawa couplings of the order of unity. As we are not considering the mixing among the scalar fields our constraints are valid only for the main component of the symmetry eigenstate scalar fields. It means that $`H^{}`$ and $`H^{}`$ denote the dominant mass eigenstates of the singly and doubly charged scalar fields, respectively.
The amplitude of the diagram in Fig. 3 is
$$A(3)\frac{M_\nu }{p^2M_H^{}^4}.$$
(36)
The scalar contribution in Fig. 3 can be as important as the standard one in Fig. 1. We have
$$\frac{A(3)}{A(1)}=\frac{1}{32G_F^2M_H^{}^4c_\theta ^4},$$
(37)
and assuming that $`A(3)/A(1)<1`$ and $`c_\theta =1`$ we get
$$M_H^{}>124\text{GeV}.$$
(38)
From Eq. (28) we see that the contribution $`v_\chi \chi ^{++}W_\mu ^{}V^\mu `$ is the dominant one in diagrams like that in Fig. 4. As we said before we will omit the mixing angles, i.e., assuming $`\chi ^{}H^{}`$. Hence we have
$$A(4)\frac{v_\chi }{M_W^2M_V^2M_H^{}^2}c_\theta ^2s_\theta ^2.$$
(39)
Next, we note that
$`{\displaystyle \frac{A(4)}{A(1)}}`$ $``$ $`{\displaystyle \frac{v_\chi }{M_\nu }}{\displaystyle \frac{p^2}{32G_F^2M_W^2M_V^2M_H^{}^2}}\mathrm{tan}^2\theta ,`$ (41)
$`5.33\times 10^{15}\mathrm{tan}^2\theta {\displaystyle \frac{(1\mathrm{G}\mathrm{e}\mathrm{V})^4}{M_V^2M_H^{}^2}},`$
where we used $`M_\nu =0.2`$ eV , $`v_\chi =3`$ TeV and $`p^2=(100\text{MeV})^2`$. If $`A(4)/A(1)<1`$ it implies
$$M_V>7.3\times 10^7\mathrm{tan}\theta \frac{(1\mathrm{G}\mathrm{e}\mathrm{V})^2}{M_H^{}}.$$
(42)
Similar analysis arises by considering Fig. 5, however it is less enhanced than the contribution of Fig. 4 because instead of $`v_\chi `$ it appears the momentum of one of the vector bosons, $`p\sqrt{p^2}`$.
More interesting are the contributions involving trilinear scalar interactions given in Eq. (31) like that of the diagram in Fig. 6. We have in this case
$$A(6)\frac{f}{M_H^{}^4M_H^{}^2},$$
(43)
where $`M_H^{}`$ represents a typical mass of the singly charged scalar bosons, say 124 GeV; $`M_H^{}`$ is the mass of the doubly charged scalar boson and $`f`$ is the trilinear coupling $`f_1`$ or $`f_2`$ in Eq. (31) with dimension of mass. The ratio of these amplitudes is:
$`{\displaystyle \frac{A(6)}{A(1)}}`$ $``$ $`{\displaystyle \frac{fp^2}{32G_F^2M_H^{}^4M_H^{}^2M_\nu c_\theta ^4}}`$ (44)
$``$ $`\left({\displaystyle \frac{f/\text{GeV}}{M_H^{}^2/\text{GeV}^2}}\right){\displaystyle \frac{4.8\times 10^7}{c_\theta ^4}},`$ (45)
If $`A(6)/A(1)<1`$, and assuming $`c_\theta =1`$ and $`M_H^{}=124`$ GeV, we obtain the constraint
$$\frac{f}{M_H^{}^2}<2.1\times 10^8\text{GeV}^1.$$
(46)
For arbitrary $`U(1)_N`$ charge for the scalar multiplets the symmetry of the potential is $`SU(3)_L[U(1)]^2`$. If the triplet $`\eta `$ and the sextet $`S`$ have both $`N=0`$, as it is the case for the present model, the trilinear couplings $`f_{1,2}`$ break the extra $`U(1)`$ symmetry. We have verify that if both $`f_{1,2}=0`$ there is indeed a pseudo-Goldstone boson . It means that $`f_{1,2}`$ are arbitrary parameters and in principle they can be small (say 1 GeV), or large (say 1 TeV) mass scales.
We see that if $`f=1`$ ($`10^3`$) TeV then $`M_H^{}`$ is greater or of the order of 300 (10) TeV. For this value for the mass of the doubly charged scalar field and $`\theta `$ small the constraint given in Eq. (35) is stronger than that of Eq. (42). For instance if $`\mathrm{tan}\theta =10^8`$ we have from Eq. (35) that $`M_V179`$ GeV.
There is also a diagram in which the doubly charged scalar field in Fig. 6 is substituted by a vector boson $`U^{}`$. Although the interactions in Eq. (25) are proportional to $`g`$ they are also derivative and proportional to the momentum $`p\sqrt{p^2}`$; hence it is suppressed with respect to the diagram in Fig. 6.
## IV Majoron emission
If the $``$ quantum number is spontaneously broken as in the present model, it means that a Majoron-like bosons does exist. Since the scalar field that is responsible for the breakdown of that continuous symmetry is $`\sigma _1^0`$, and it belongs to a triplet of the subgroup $`SU(2)U(1)`$, this Majoron-like Goldstone has similar couplings that the triplet Majoron model of Ref. . It is well known that this sort of Majoron model has been ruled out by the LEP data . Apparently, since the Higgs sector of the present model is rather complicated having a neutral scalar singlet (under $`SU(2)U(1)`$), $`\chi ^0`$, it seems that the Majoron-like Goldstone in this case it will be able to avoid the LEP constraints as claimed in Refs. . However, we will show that this is not indeed the case.
The mass matrices of the scalar and pseudoscalar in this model have been given in Ref. . Here we will only give the results of the mass eigenvalues and the respective mixing matrix in the CP-even scalar sector. The argument in Ref. was the following. Let us begin with the relation $`R_4^0=_j𝒪_{4j}^oH_j^0`$, where $`H_j^0,j=1,\mathrm{}5,`$ denotes the mass scalar eigenstates and $`R_4`$ the real component of the scalar field $`\sigma _1^0`$ according to the general shifting of the neutral scalar fields in the scalar potential of the form $`X_i^0(1/\sqrt{2})(v_{X_i}+R_i+iI_i),i=1,2,3,4,5`$ where $`X_i^0=\eta ,\rho ,\chi ,\sigma _1,\sigma _2`$ respectively. In this case if $`H_1^0`$ denotes the lightest scalar boson ($`M_{H_1}<M_Z`$, we are assuming a mass spectrum where $`M_{H_i}<M_{H_j}`$ if $`i<j`$), the contribution to the decay mode $`Z^0H_1^0M^0`$ is $`\mathrm{\Gamma }_{H_1^0M^0}^Z=2|𝒪_{41}^o|^2\mathrm{\Gamma }_{\nu \overline{\nu }}^Z`$. Hence, if $`|𝒪_{41}^o|<10^2`$ the model would be consistent with the LEP data i.e., now $`\mathrm{\Gamma }_{H_1^0M^0}^Z`$ would be reduced to an acceptable level.
First of all recall that as shown in Ref. the Majoron-like boson decouples from the other pseudoscalar fields i.e., $`\mathrm{Im}\sigma _1^0I_4=M^0`$. For instance, using the same values of the VEV and $`f_{1,2}=1`$ TeV and with the dimensionless constant of the scalar potential given in Ref. with $`\lambda _k=0.1`$ for $`k=1,2,3,4,5,6,7,8,9,18,20`$, $`\lambda _m=0.01`$ for $`m=13,14,16,17`$, $`\lambda _n=0.001`$ for $`n=10,11,12,19`$ and $`\lambda _{15}=0.05`$ we obtain the following masses in the scalar sector $`(0.056,102,1342,3626,4325)`$ GeV and the mixing matrix (up to three decimal places)
$$O^o=\left(\begin{array}{ccccc}0.0& 0.081& 0.010& 0.996& 0.021\\ 0.0& 0.995& 0.029& 0.082& 0.039\\ 0.0& 0.030& 0.999& 0.008& 0.004\\ 1.0& 0.0& 0.0& 0.0& 0.0\\ 0.0& 0.040& 0.005& 0.017& 0.999\end{array}\right)$$
(47)
This pattern of mixing matrix remains the same for several values of the parameters provided that $`v_{\sigma _1}`$ is a small VEV restricted to the condition that it has to be smaller than 3.89 GeV . From Eq. (47) it can be seen that the scalar partner of the Majoron is always mainly the lightest scalar i.e., $`|𝒪_{41}^o|1`$ and it would be always produced at LEP. We see that the Majoron in the minimal 3-3-1 model has been also ruled out by the LEP data.
One possibility to recover consistency with the LEP data is to break explicitly the $``$ symmetry by adding trilinear terms like $`f_3\eta ^TS\eta ,f_4\chi ^TS^{}\rho `$ in the scalar potential, see Eq. (31). In this case there is no Majoron at all and although $`v_{\sigma _1}`$ still has a small value, due to $`f_3`$ all scalars are heavy enough to not be produced at the LEP energies . Of course, in this case there is no contribution to Majoron emission in the neutrinoless double beta decay. However, our results in Sec. III are still valid since they depend only on the small value of $`v_{\sigma _1}`$.
Another possibility which we will consider here is to modify the model by introducing a scalar singlet $`\mathrm{\Sigma }^0`$, which carries $`=2`$ (or $`L=2`$), in the same way as considered in Ref. in the context of a $`SU(2)U(1)`$ model. In this case we have to add the following terms to the scalar potential in Ref.
$`V(X_i,\mathrm{\Sigma })`$ $`=`$ $`\mu _5^2\mathrm{\Sigma }^2+\lambda _{21}\mathrm{\Sigma }^4+{\displaystyle \underset{i}{}}[\lambda _{X_i}Tr(X_i^{}X_i)\mathrm{\Sigma }^2`$ (48)
$``$ $`\kappa \eta ^TS^{}\eta \mathrm{\Sigma }+H.c.]`$ (49)
where $`X_i`$ denotes any triplet $`\eta ,\rho ,\chi `$ or the sextet $`S`$, and we will denote $`\lambda _{X_i}`$ as $`\lambda _{22,23,24,25}`$ respectively and $`\kappa >0`$. The neutral Higgs sector contains six CP-even scalars and three massive CP-odd pseudoscalar beside the massless CP-odd Majoron. The neutral scalar singlet also gains a VEV, i.e., $`\mathrm{\Sigma }=(v_\mathrm{\Sigma }+R_6+iI_6)/\sqrt{2}`$, and the mass term is given by $`M^2/2`$, where $`M^2`$ in the pseudoscalar sector in the basis $`I_1,I_2,I_3,I_4,I_5,I_6`$ is given by (the constraints equation appear in the Appendix)
$`M_{11}={\displaystyle \frac{\lambda _{16}}{2\sqrt{2}}}{\displaystyle \frac{v_\rho ^2v_{\sigma _2}}{v_\eta }}+2\lambda _{17}v_{\sigma _2}^2+{\displaystyle \frac{1}{2\sqrt{2}}}(\lambda _{15}v_\chi v_{\sigma _2}`$ (50)
$``$ $`f_1v_\rho ){\displaystyle \frac{v_\chi }{v_\eta }}+2\kappa v_{\sigma _1}v_\mathrm{\Sigma }+{\displaystyle \frac{t_\eta }{v_\eta }},`$ (55)
$`M_{22}={\displaystyle \frac{1}{4}}(\sqrt{2}f_1v_\eta +f_2v_{\sigma _2}){\displaystyle \frac{v_\chi }{v_\rho }}+{\displaystyle \frac{t_\rho }{v_\rho }},`$
$`M_{33}={\displaystyle \frac{1}{4}}(\sqrt{2}f_1v_\eta +f_2v_{\sigma _2}){\displaystyle \frac{v_\rho }{v_\chi }}+{\displaystyle \frac{t_\chi }{v_\chi }},`$
$`M_{44}=\kappa {\displaystyle \frac{v_\eta ^2v_\mathrm{\Sigma }}{2v_{\sigma _1}}}+{\displaystyle \frac{t_{\sigma _1}}{v_{\sigma _1}}},`$
$`M_{55}={\displaystyle \frac{1}{2\sqrt{2}}}(\lambda _{15}v_\chi ^2\lambda _{16}v_\rho ^2){\displaystyle \frac{v_\eta }{v_{\sigma _2}}}+2\lambda _{17}v_\eta ^2{\displaystyle \frac{f_2}{4}}{\displaystyle \frac{v_\rho v_\chi }{v_{\sigma _2}}}`$
$`+`$ $`{\displaystyle \frac{t_{\sigma _2}}{v_{\sigma _2}}},M_{66}=\kappa {\displaystyle \frac{v_\eta ^2v_{\sigma _1}}{2v_\mathrm{\Sigma }}}+{\displaystyle \frac{t_\mathrm{\Sigma }}{v_\mathrm{\Sigma }}},`$ (61)
$`M_{12}={\displaystyle \frac{f_1}{2\sqrt{2}}}v_\chi ,M_{13}={\displaystyle \frac{f_1}{2\sqrt{2}}}v_\rho ,M_{14}=\kappa v_\eta v_\mathrm{\Sigma },`$
$`M_{15}={\displaystyle \frac{1}{2\sqrt{2}}}(\lambda _{15}v_\chi ^2\lambda _{16}v_\rho ^2)+2\lambda _{17}v_\eta v_{\sigma _2},M_{16}=\kappa v_\eta v_{\sigma _1},`$
$`M_{23}={\displaystyle \frac{1}{4}}(\sqrt{2}f_1v_\eta +f_2v_{\sigma _2}),M_{24}=0,M_{25}={\displaystyle \frac{f_2}{4}}v_\chi ,`$
$`M_{26}=0,M_{34}=0,M_{35}={\displaystyle \frac{f_2}{4}}v_\rho ,M_{36}=0,`$
$`M_{45}=0,M_{46}=\kappa {\displaystyle \frac{v_\eta ^2}{2}},M_{56}=0.`$
The mass matrix above has two true Goldstone bosons $`G_{1,2}^0`$ and the Majoron-like one, $`M^0`$, and three massive CP-odd pseudoscalar bosons. The massless ones are given by
$`G_1^0`$ $`=`$ $`(0,v_\rho /v_\chi ,1,0,0,0)/(1+v_\rho ^2/v_\chi ^2)^{1/2}`$ (62)
$`G_2^0`$ $`=`$ $`({\displaystyle \frac{v_\eta }{v_{\sigma _2}}},{\displaystyle \frac{v_\rho v_\chi ^2}{V_1}},{\displaystyle \frac{v_\rho ^2v_\chi }{V_1}},{\displaystyle \frac{2v_{\sigma _1}v_\mathrm{\Sigma }^2}{V_2}},1,{\displaystyle \frac{2v_{\sigma _1}^2v_\mathrm{\Sigma }}{V_2}})/N`$ (63)
$`M^0`$ $`=`$ $`(0,0,0,v_{\sigma _1}/v_\mathrm{\Sigma },0,1)/(1+v_{\sigma _1}^2/v_\mathrm{\Sigma }^2)^{1/2}`$ (64)
where $`V_1=v_{\sigma _2}(v_\rho ^2+v_\chi ^2)`$, $`V_2=v_{\sigma _2}(v_{\sigma _1}^2+v_\mathrm{\Sigma }^2)`$; $`N`$ is the normalization factor that we will omit here. We have verify that $`M^0`$ in Eq. (64) is in fact the Majoron: by adding an explicit $``$-violating term, like $`f_3\eta ^TS\eta `$, it gets a mass while the other two $`G_{1,2}`$ remain massless. The massive pseudoscalars, for the parameters used before have the following masses (in GeV): 174, 3625 and 4325. On the other hand, if $`v_{\sigma _1}=0`$, which forces $`\kappa =0`$, the Majoron is purely singlet and the real and imaginary parts of $`\sigma _1^0`$ are mass degenerate, i.e., form a complex field, with mass
$`m_{\sigma _1}`$ $`=`$ $`\mu _4^2+\lambda _{10}v_{\sigma _2}^2+(\lambda _{12}+\lambda _{19}){\displaystyle \frac{v_\eta ^2}{2}}+\lambda _{13}{\displaystyle \frac{v_\chi ^2}{2}}`$ (65)
$`+`$ $`\lambda _{14}{\displaystyle \frac{v_\rho ^2}{2}}+\lambda _{25}{\displaystyle \frac{v_\mathrm{\Sigma }^2}{2}}.`$ (66)
We see that in this case the Majoron has not doublet components at all and it is mainly singlet. Hence it is possible to keep consistence with LEP data. Although there are astrophysical constraints (the Majoron emission implies a different rate for the stellar cooling) that have to be taken into account , in the basis we have chosen they are less severe since we have avoided the doublet component of the Majoron. Any way, since these constraints have been already considered in Ref. and they imply that $`v_{\sigma _1}<0.33`$ GeV if $`v_\mathrm{\Sigma }=1`$ TeV, we will use these values for $`v_{\sigma _1},v_\mathrm{\Sigma }`$.
Once we have shown in what situation there is a safe Majoron-like boson in the present model we can consider the emission of this Goldstone boson in the neutrinoless double beta decay. In fact, as in the triplet Majoron model, in the present model it is possible to have the neutrinoless double beta decay with Majoron emission: $`2n2p+2e^{}+M^0`$ , denoted here by $`(\beta \beta )_{0\nu M}`$. We will denote the strength of the amplitude of the diagram $`i`$ of the $`(\beta \beta )_{0\nu M}`$ decay by $`B(i)`$. This decay proceeds via the diagram in Fig. 7 and it has an strength proportional to
$$B(7)\frac{m_\nu (v_{\sigma _1}/v_\mathrm{\Sigma })}{M_X^{}^4p^2v_{\sigma _1}},$$
(67)
where $`X^{}`$ can be an scalar or a vector boson, i.e., the diagram in Fig. 7 can be formed with anyone of Figs.1, 2 or 3 with a Majoron attached to the neutrinos. The couplings between neutrinos and the Majoron are diagonal and given by $`m_\nu /v_{\sigma _1}`$. Notice that in Eq. (67) it appears the truly neutrino mass instead of the effective mass $`M_\nu `$ defined in Eq. (1). However we still can assume that neutrinos have small masses and numerically $`m_\nu M_\nu `$. We will assume also that the contribution to the $`(\beta \beta )_{0\nu M}`$-decay in Fig. 7 with $`X^{}=W^{}`$ is the reference one. This diagram depends only on the neutrino masses and mixing angles, and we will compare it with other contributions like the one in Fig. 8. The couplings of the Majoron to the vector bosons are proportional to $`v_{\sigma _1}`$ and so they are negligible. We will consider only the diagram with the Majoron coupled to the scalar $`H^{}`$ since it is proportional to the trilinear $`f_2`$ shown in Eq. (31). We have
$$B(8)\frac{ff_2}{M_H^{}^6M_H^{}^2},$$
(68)
with $`f`$ can be $`f_1`$ or $`f_2`$.
Let us consider the ratio
$$\frac{B(7)}{A(1)}\frac{m_\nu Q}{32G_F^2M_X^4M_\nu c_\theta ^4v_\mathrm{\Sigma }},$$
(69)
where we have introduced the factor $`Q`$ which denotes the available energy. It implies that the diagram in Fig. 7 is a potentially important contribution when $`X`$ is the $`W`$ vector boson since for $`Q3`$ MeV the suppression of $`B(7)`$ will depend mainly on the value of $`v_\mathrm{\Sigma }`$. If $`B(7)/A(1)<1`$ we obtain that $`v_\mathrm{\Sigma }>1.65\times 10^2`$ GeV which is automatically satisfied.
On the other hand, comparing the amplitudes of the diagrams in Figs. 7 and 8 we have
$$\frac{B(8)}{B(7)}\frac{ff_2p^2M_X^4v_\mathrm{\Sigma }}{M_H^{}^6M_H^{}^2m_\nu },$$
(70)
and for $`M_X=M_V=400`$ GeV and $`v_\mathrm{\Sigma }=1`$ TeV, using typical values as $`f=f_1=f_2=1`$ TeV, $`M_H^{}=124`$ GeV and $`M_H^{}=500`$ GeV, and the other parameters in Eq. (70), we have that $`B(8)/B(7)1.3\times 10^9`$ or; $`B(8)/B(7)1.5\times 10^3`$ if $`f_1=f_2=f=10^3`$ TeV. The relative importance of the processes in Fig. 7 and 8 will depend on the values of the trilinear parameters and on the value of $`v_\mathrm{\Sigma }`$.
## V Conclusions
We see that in the 3-3-1 model, like in other models with complicated Higgs sector , besides the well known mechanism of exchanging massive Majorana neutrinos between two standard model $`VA`$ vertices, there are new contributions involving the exchange of scalar bosons. However, unlike similar mechanism in the context of extensions of the standard model there is no need of fine tuning in order to have trilinear scalar couplings giving large contributions to the several neutrinoless double beta decay modes. Notice that effective interactions from diagrams like those in Fig. 3 are still parameterized in the form of two general four-fermion effective interactions (they are point-like at the Fermi scale) exchanging a light neutrino in between . However, contributions involving trilinear interactions like those in Figs. 4, 5 and 6 necessarily need a six-fermion effective interaction parameterization.
Another important point to be stressed here is that in the present model the double Majoron emission: $`2d2p+2e^{}+2M^0`$ may be as contributing as the decay with only one Majoron boson. This decay is expected to be important in supersymmetric models . In the present model it can occur because in diagrams like that in Fig. 8 a second Majoron can be attached to the scalar lines. Since this coupling is proportional to the trilinear $`f_2`$ it is still possible that the suppression coming from the mass square in the denominator do not sufficiently suppresses this process (there is also an important contribution coming from the vertex $`v_\chi \chi ^{}\eta _1^+M^0`$). There are also contributions similar to the one in Fig. 8 but now with the scalar-Majoron vertex above being substituted by the vertex $`W^{}V^+M^0`$ which is proportional to $`g^2v_{\sigma _2}`$ and for this reason it is not necessarily suppressed. It is interesting to note that this sort of contribution to the $`(\beta \beta )_{0\nu MM}`$ decay, coming from adding another trilinear coupling in the diagram in Fig. 8, which is not derivatively suppressed, was not considered in Ref. . Recent experimental data on Majoron emission decays have been constrained only the effective Majoron-neutrino coupling constant . Other processes like the double $`K`$ capture can also be important in the present model.
Some comments are now in order. i) We have not considered possible cancellations among several contributions to each diagram. It means that our constraints are valid, as we said in Sec. III, for the main component of each scalar field of the singly and doubly charged scalar sectors.
ii) Our results were obtained assuming that all new contributions, say to the $`(\beta \beta )_{0\nu }`$ decay, are at most as important as the contribution due to a light massive Majorana neutrino exchange which is proportional to $`M_\nu `$. However, we can wonder what would be the value of the effective mass $`M_\nu `$ if we use the oscillation data and direct measurements on neutrino masses. Recent analysis shown that assuming a normal mass hierarchy the effective mass parameter can take any value from zero to the present upper limit . In fact, if the data on oscillation is put together with that from $`(\beta \beta )_{0\nu }`$ decay and tritium beta decay, it was shown that if the minimum of $`M_\nu `$ with respect to the mixing angles is greater than the present bound of 0.2 eV, then neutrinos are quasi-Dirac particles . As discussed previously, the black box of Ref. may induce a negligible Majorana mass to the neutrinos and in the context of the present model we must interpret this situation as an indication of the fact that the main contribution to the $`(\beta \beta )_{0\nu }`$ decay is not the diagram in Fig. 1. In this case the neutrinos would be almost Dirac particles and the constraints on the several mass scales of the model should be obtained by comparing directly these contributions with the lower bound on the half life $`T_{1/2}^{0\nu }>1.8\times 10^{25}`$ yr .
iii) In the basis we have chosen, see Eq. (64), the Majoron couples at the tree level only to neutrinos, and hence the constraints on the $`\nu `$-$`\nu `$-$`M`$ vertex, coming from muon ($`\mu e\nu \nu M`$), pion and kaon $`\pi ^+(K^+)l\overline{\nu }M`$ decays, are the same as in Ref. . The existence of the vertex $`\nu e^{}V^+`$, which is proportional to $`\mathrm{sin}\theta `$ and for this reason is not relevant for laboratory processes, may have, as we said before, important astrophysical consequences .
iv) phenomenology of non-zero initial electric charge processes, like $`e^{}e^{}`$ and hadronic ones, will furnish constraints on the trilinear vertices appearing in Figs. 4, 5 and 6 but this will be considered elsewhere. The 3-3-1 model has a rich scalar sector indeed. This implies that it may be, in principle, difficult to separate, in a given process, the contributions of all fields belonging to a charged sector. However, it has been shown that in lepton-lepton colliders the left-right asymmetries are not sensible to the scalar contributions. It means that those asymmetries are the appropriate observable for the doubly charged vector bilepton discovery . We see that the opposite occurs in the $`(\beta \beta )_{0\nu }`$ decay: it is possible that the main contribution comes from the doubly charged scalar boson, through the diagram in Fig. 6, while the respective vector boson contribution seems to be negligible.
Finally we would like to compare our Majoron model with that of Schechter and Valle . Firstly, we notice that although our model has two singlet, three doublets and a triplet of scalars under the subgroup $`SU(2)U(1)`$, the respective scalar potential is not reduced to the scalar potential invariant under the standard $`SU(2)U(1)`$ symmetry, involving the same multiplets. For instance in our model there are cubic invariants which are not present in the former. Secondly, we have not introduced right-handed neutrinos and for this reason we have only light neutrinos. It means that the singlet $`\mathrm{\Sigma }`$ does not couple to the leptons and that the coupling of neutrinos to Majoron and $`Z^0`$ are diagonal. Thus, the decays $`\nu _H\nu _L+M^0`$ and $`\nu _H\nu _L+\nu _L^{}+\nu _L^{}`$, are not induced at the tree level, where $`\nu _H`$, although light, is heavier that $`\nu _L`$. The decay $`\nu _H\nu _L^c+M^0`$ is produced at the one loop level due to the mixing between $`W`$ and $`V`$. The vertex is proportional to $`(g^2/\sqrt{2})(v_{\sigma _2}v_{\sigma _1}/v_\mathrm{\Sigma })`$; then even with $`v_{\sigma _2}`$ of the order of 10 GeV and $`v_\mathrm{\Sigma }`$ of the order of 1 TeV the lifetime is of the order of the age of the universe. (The decay $`\nu _H\nu _L+M^0`$ also occurs but the vertex involved is proportional to $`(g^2/\sqrt{2})(v_{\sigma _1}^2/v_\mathrm{\Sigma })`$.) Notice also that in the basis given in Eq. (64) the Majoron does not couple to the charged leptons so there is not the process $`\gamma +eM^0+e`$ at the tree level which imposes severe astrophysical constraints in $`v_{\sigma _1}`$ as has been noted in Ref. .
###### Acknowledgements.
This work was supported by Fundação de Amparo à Pesquisa do Estado de São Paulo (FAPESP), Conselho Nacional de Ciência e Tecnologia (CNPq) and by Programa de Apoio a Núcleos de Excelência (PRONEX).
## A Contraints equation of the scalar potential
Here we show the constraints equation that must be satisfied by the scalar potential
$`t_\eta =\mu _1^2v_\eta +\lambda _1v_\eta ^3+{\displaystyle \frac{\lambda _4}{2}}v_\rho ^2v_\eta +{\displaystyle \frac{\lambda _5}{2}}v_\chi ^2v_\eta +{\displaystyle \frac{\lambda _{12}}{2}}(v_{\sigma _1}^2+v_{\sigma _2}^2)v_\eta `$
$``$ $`{\displaystyle \frac{\lambda _{15}}{2\sqrt{2}}}v_\chi ^2v_{\sigma _2}+{\displaystyle \frac{\lambda _{16}}{2\sqrt{2}}}v_\rho ^2v_{\sigma _2}\lambda _{17}v_{\sigma _2}^2v_\eta +{\displaystyle \frac{\lambda _{19}}{2}}v_{\sigma _1}^2v_\eta `$
$``$ $`\kappa v_\eta v_{\sigma _1}v_\mathrm{\Sigma }+{\displaystyle \frac{f_1}{2\sqrt{2}}}v_\rho v_\chi ,`$
$`t_\rho =\mu _2^2v_\rho +\lambda _2v_\rho ^3+{\displaystyle \frac{\lambda _4}{2}}v_\eta ^2v_\rho +{\displaystyle \frac{\lambda _6}{2}}v_\chi ^2v_\rho +{\displaystyle \frac{\lambda _{14}}{2}}(v_{\sigma _1}^2+v_{\sigma _2}^2)v_\rho `$
$`+`$ $`{\displaystyle \frac{\lambda _{16}}{\sqrt{2}}}v_{\sigma _2}v_\eta v_\rho +{\displaystyle \frac{\lambda _{20}}{4}}v_{\sigma _2}^2v_\rho +\lambda _{23}{\displaystyle \frac{v_\rho v_\mathrm{\Sigma }^2}{2}}`$
$`+`$ $`{\displaystyle \frac{f_1}{2\sqrt{2}}}v_\eta v_\chi +{\displaystyle \frac{f_2}{4}}v_{\sigma _2}v_\chi ,`$
$`t_\chi =\mu _3^2v_\chi +\lambda _3v_\chi ^3+{\displaystyle \frac{\lambda _5}{2}}v_\eta ^2v_\chi +{\displaystyle \frac{\lambda _6}{2}}v_\rho ^2v_\chi +{\displaystyle \frac{\lambda _{13}}{2}}(v_{\sigma _1}^2+v_{\sigma _2}^2)v_\chi `$
$``$ $`{\displaystyle \frac{\lambda _{15}}{\sqrt{2}}}v_{\sigma _2}v_\eta v_\chi +{\displaystyle \frac{\lambda _{18}}{4}}v_{\sigma _2}^2v_\chi +\lambda _{24}{\displaystyle \frac{v_\chi v_\mathrm{\Sigma }^2}{2}}`$
$`+`$ $`{\displaystyle \frac{f_1}{2\sqrt{2}}}v_\eta v_\rho +{\displaystyle \frac{f_2}{4}}v_{\sigma _2}v_\rho ,`$
$`t_{\sigma _1}=\mu _4^2v_{\sigma _1}+\lambda _{10}(v_{\sigma _1}^2+v_{\sigma _2}^2)v_{\sigma _1}+\lambda _{11}v_{\sigma _1}^3+{\displaystyle \frac{\lambda _{12}}{2}}v_\eta ^2v_{\sigma _1}`$
$`+`$ $`{\displaystyle \frac{\lambda _{13}}{2}}v_\chi ^2v_{\sigma _1}+{\displaystyle \frac{\lambda _{14}}{2}}v_\rho ^2v_{\sigma _1}+{\displaystyle \frac{\lambda _{19}}{2}}v_\eta ^2v_{\sigma _1}+{\displaystyle \frac{\lambda _{25}}{2}}v_{\sigma _1}v_\mathrm{\Sigma }^2{\displaystyle \frac{\kappa }{2}}v_\eta ^2v_\mathrm{\Sigma },`$
$`t_{\sigma _2}=\mu _4^2v_{\sigma _2}+\lambda _{10}(v_{\sigma _2}^2+v_{\sigma _1}^2)v_{\sigma _2}+{\displaystyle \frac{\lambda _{11}}{2}}v_{\sigma _2}^3+{\displaystyle \frac{\lambda _{12}}{2}}v_\eta ^2v_{\sigma _2}`$
$`+`$ $`{\displaystyle \frac{\lambda _{13}}{2}}v_\chi ^2v_{\sigma _2}+{\displaystyle \frac{\lambda _{14}}{2}}v_\rho ^2v_{\sigma _2}{\displaystyle \frac{\lambda _{15}}{2\sqrt{2}}}v_\chi ^2v_\eta +{\displaystyle \frac{\lambda _{16}}{2\sqrt{2}}}v_\rho ^2v_\eta \lambda _{17}v_\eta ^2v_{\sigma _2}`$
$`+`$ $`{\displaystyle \frac{\lambda _{18}}{4}}v_\chi ^2v_{\sigma _2}+{\displaystyle \frac{\lambda _{20}}{4}}v_\rho ^2v_{\sigma _2}+{\displaystyle \frac{\lambda _{25}}{2}}v_{\sigma _2}v_\mathrm{\Sigma }^2+{\displaystyle \frac{f_2}{4}}v_\rho v_\chi ,`$
$`t_\mathrm{\Sigma }=\mu _5^2v_\mathrm{\Sigma }+\lambda _{21}v_\mathrm{\Sigma }^3+{\displaystyle \frac{\lambda 22}{2}}v_\eta ^2v_\mathrm{\Sigma }+{\displaystyle \frac{\lambda _{23}}{2}}v_\rho ^2v_\mathrm{\Sigma }+{\displaystyle \frac{\lambda _{24}}{2}}v_\chi ^2v_\mathrm{\Sigma }`$
$`+`$ $`{\displaystyle \frac{\lambda _{25}}{2}}(v_{\sigma _1}^2+v_{\sigma _2}^2)v_\mathrm{\Sigma }{\displaystyle \frac{\kappa }{2}}v_\eta ^2v_\sigma .`$
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# 1 Introduction
## 1 Introduction
Over the past decade information theory has been generalized to include quantum mechanical systems, for example, a two-level quantum system has come to be known as a qubit in this context. The additional freedom introduced with the quantum mechanical superposition principle has opened up a variety of capabilities that go well beyond those of conventional information techniques. There are two distinct directions in which progress is currently being made: quantum computation and error correction or prevention on the one hand, and nonlocality and distillation, on the other hand. In each of those progresses, quantum entanglement that provides a good measure of quantum correlations plays an important role.
There are a number of good measures of the amount of entanglement for two quantum systems in a pure state, a good measure of entanglement for mixed states is also found though it is hard to compute for a general state.
In present paper, considering the entanglement from the other aspect, we prefer to discuss the change of entanglement due to the state changes rather than to compute straightforwardly the entanglement of an arbitrary state.
In the framework of quantum information theory, the state change allowed by quantum mechanics may by treated in terms of quantum operations, a simple example is the unitary evolution experienced by a close quantum system. The final state of the system is related to the initial state by a unitary transformation $`U`$,
$$\rho \epsilon (\rho )=U\rho U^+.$$
Unitary evolution is widely in use of quantum gates and circuits as a quantum operation. In addition to the unitary evolution, environment coupling to a quantum system or a measurement performed on the quantum system changes the state too. The connection of quantum operations to quantum measurements is easy to explain. Standard text book treatments describe quantum measurement in terms of a complete set of orthogonal projection operators for the system being measured. This formalism, however, does not describe many of the measurements that can be performed on a quantum system. The most general type of measurement that can be performed on a quantum system is known as generalized measurement. Generalized measurement can be understood within the framework of unitary evolution, because most generalized measurements can be realized through many dynamical processes, and the state change due to environment may be also treated in terms of the unitary evolution, since an arbitrary open system may be enlarged by including the environment to be a close system. In this sense, unitary evolution is one of the most general types of state change possible in quantum mechanics.
The rest of present paper is organized as follows: In Sec.2, we present a general method to factorize the unitary evolution operator $`U(t)`$ for a close system. The results may be generalized in the treatment of many autonomous dynamical systems. Sec.3 contains our results on the entanglement change occurring in a dynamical process. Finally, in Sec.4, we present two typical examples and some conclusions.
## 2 Factorizing the unitary evolution operator $`U(t)`$
As noted above, the unitary evolution operator is one of the most general types of state change possible in quantum mechanics, the point of this section is to factorize the evolution operator into a set of independent one. For this end, we discuss the following cases.
case A— The Hamiltonian can be written as a finite sum
$$H(t)=\underset{i}{\overset{m}{}}a_i(t)H_i,$$
(2.1)
where $`a_i(t)`$ are a set of linearly independent complex valued functions of time, and $`H_i`$ are constant operators. In addition, the set of operators $`H_i(i=1,\mathrm{},m)`$ may be enlarged by repeated commutation to a Lie algebra $`L`$ with finite dimension $`n(nm)`$. With this presupposition, the unitary evolution operator can be uncoupled into a set of independent operators
$$U(t)=U_1(t)U_2(t)\mathrm{}U_n(t),$$
(2.2)
where each component $`U_i(t)`$ is an operator satisfying
$$\frac{d}{dt}U_i(t)=\dot{g}_i(t)H_iU_i(t),U_i(0)=1.$$
(2.3)
With the scalar function $`g_i(t)`$ being the solution to a set of nonlinear differential equations
$$\frac{d}{dt}g_i(t)=\underset{i=1}{\overset{n}{}}\eta _{ik}a_k(t),g_i(0)=0,$$
(2.4)
where $`\eta _{ik}`$ are nonlinear function of $`g`$’s. Thus we have factorized the unitary evolution operator into the form:
$$U(t)=\underset{i=1}{\overset{n}{}}e^{g_i(t)H_i}.$$
(2.5)
Especially, for a general case of a dynamically closed quantum system which consists of two interacting subsystems $`A`$ and $`B`$, the total Hamiltonian may be written as a sum of three terms
$$H=H_A+H_B+H_{int},$$
(2.6)
the first two terms represent the free Hamiltonian of subsystem $`A`$ and $`B`$, respectively, and the last term describes the interaction between the $`A`$ and $`B`$. Following the procedure stated above, we arrive at
$$U(t)=\underset{i=AB,int,\mathrm{},M}{}e^{g_i(t)}H_i.$$
(2.7)
Here, $`H_i(i=A,B,int,\mathrm{}M)`$ are elements of Lie algebra with finite dimension enlarged by $`H_A,H_B,H_{int}`$ .
case B— In the case of the dimension of the Lie algebra enlarged by $`\{H_i\}`$ is infinite, we can factorized the unitary evolution operator using the general Baker-campbell-Hausdorff formula. To start with, we give the evolution operator of the system under consideration
$$U(t)=e^{itH}=e^{it(H_A+H_B+H_{int})},$$
(2.8)
where $`H_A,H_B,H_{int}`$ are the same as in eq.(2.6), the eq.(2.8) can approximately be written as
$$U(t)=e^{\frac{1}{2}[H_A+H_B,H_{int}]t^2}e^{iH_{int}t}e^{i(H_A+H_B)t}+O(t^3).$$
(2.9)
This splitting formula is hold in the case that $`t`$ has to be safely smaller than a typical energy of the system. Thus, even in the simplest case, a better method is needed. Let $`n`$ be a positive integer. The exponential function satisfies the scaling identity
$$exp(iHt)=[exp(iHt/2^n)]^{2n}.$$
(2.10)
When $`n`$ is sufficiently large, the argument $`\frac{t}{2^n}`$ is in some sense small. Eqs.(2.9) and (2.10) together give
$`U(t)`$ $`=`$ $`e^{\frac{1}{2}[H_A+H_B,H_{int}]\tau ^2}e^{iH_{int}\tau }e^{i(H_A+H_B)\tau }e^{\frac{1}{2}[H_A+H_B,H_{int}]\tau ^2}`$ (2.11)
$``$ $`e^{iH_{int}\tau }e^{i(H_A+H_B)\tau }\mathrm{}e^{\frac{1}{2}[H_A+H_B,H_{int}]\tau ^2}e^{iH_{int}\tau }e^{i(H_A+H_B)\tau }+O(\tau ^3),`$
where $`\tau =\frac{t}{2^n}`$. Still higher-order formulae are known. We would like to point out that, in quantum computation, the $`n=1`$ is widely taken in use and it is large enough to avoid the decoherence during quantum computing.
## 3 Quantification of entanglement
In the previous section we have factorized the time evolution operator $`U(t)`$. The question remains open about how does the entanglement change in a dynamical process. Of course, this question is not entirely well defined unless we state what physical circumstances characterized the amount of entanglement. This suggest that there is no unique measure of entanglement. Before we define the measure of entanglement we expand the density operator for a close system. Suppose that the two interacting subsystems are initially separable, i.e., their initial density operator (state) can be written in a form
$$\rho (0)=\rho _A(0)\rho _B(0),$$
(3.1)
to use the entanglement for quantum information processing, however, we need a inseparable state, more precisely, a state in pure entanglement form. The procedure of converting a separable state to inseparable one can be performed, as stated in Sec.1, through a unitary evolution operator $`U(t)`$(in addition, a partial trace is also needed sometimes.)
$$\rho (t)=U(t)\rho (0)U^+(t).$$
(3.2)
If the interaction between the two subsystems is small, it is natural to attempt some sort of Taylor series expansion of the exponential in eqs. (3.2) (2.8) and (2.11), which give
case A
$`\rho (\lambda ,t)`$ $`=`$ $`\rho _A^0(t)\rho _B^0(t)+\lambda {\displaystyle \underset{iA,B}{}}({\displaystyle \frac{f_i}{\lambda }}H_i\rho _A^0(t)\rho _B^0(t)+\rho _A^0(t)\rho _B^0(t){\displaystyle \frac{f_i^{}}{\lambda }}H_i)`$ (3.3)
$`+`$ $`{\displaystyle \frac{\lambda ^2}{2}}{\displaystyle \underset{i,jA,B}{}}({\displaystyle \frac{f_i}{\lambda }}{\displaystyle \frac{f_j}{\lambda }}H_iH_j\rho _A^0(t)\rho _B^0(t)+\rho _A^0(t)\rho _B^0(t){\displaystyle \frac{f_i}{\lambda }}{\displaystyle \frac{f_j}{\lambda }}H_iH_j`$
$`+`$ $`{\displaystyle \frac{f_i}{\lambda }}H_i\rho _A^0(t)\rho _B^0(t){\displaystyle \frac{f_j^{}}{\lambda }}H_j)+O(\lambda ^2)`$
and case B
$`\rho (\lambda ,t)`$ $`=`$ $`\rho _A^0(t)\rho _B^0(t){\displaystyle \frac{\lambda }{2}}({\displaystyle \frac{t}{2^n}})^2{\displaystyle \underset{i=0}{\overset{2^n1}{}}}\{[H_A+H_B,H_{int}(t_i)]_{},\rho _A^0(t)\rho _B^0(t)\}_+`$ (3.4)
$``$ $`i\lambda ({\displaystyle \frac{t}{2^n}}){\displaystyle \underset{i=0}{\overset{2^n1}{}}}[H_{int}(t_i),\rho _A^0(t)\rho _B^0(t)]_{}+O(\lambda ^2)+O(({\displaystyle \frac{t}{2^n}})^3),`$
where $`\lambda `$ denotes the coupling constant, and $`\rho _i^0(t)`$ represents the state of subsystem $`i`$ at time $`t`$ with $`\lambda =0`$. The results presented above suggest that we may take the form
$$\delta D(\rho )=\rho (t)\rho _A^0(t)\rho _B^0(t)^2=Tr(\rho (t)\rho _A^0(t)\rho _B^0(t))^2$$
(3.5)
as a measure of entanglement change in the time evolution process. Noticing the initial state is separable, the measure of entanglement change (3.5) is a measure of entanglement in reality. Although the definition of the measure for entanglement is not unique, it has to satisfies the three condition stated below:
(1)$`D(\rho )=0`$ if and only if $`\rho `$ is separable.
(2)Local unitary operators leave $`D(\rho )`$ invariant, i.e.
$$D(\rho )=D(U_AU_B\rho U_A^+U_B^+).$$
(3)The expected entanglement cannot increase under Local general measurements+Classical communication+Postselection (LGM+CC+PS) given by $`_iV_i^+V_i=1,`$i.e.,
$$Tr(\rho _i)D(\rho _i/Tr\rho _i)D(\rho ),$$
where$`\rho _i=V_i\rho V_i^+.`$ For the measure of entanglement change proposed above, (1)follows from the fact that $`D(\rho )`$ is a true metric, and (2) is obvious. Property (3) is satisfied too. We believe that there are numerous other nontrivial choices for measure of entanglement, one of the choices could not be said to be more important than any other,the present choice has the advantage that it is easy to compute for any dynamical process.
Our discussion so far has centered on the entanglement change in a dynamical process. To complete it we still need to show that this definition can be generalized for any process that quantum mechanics allowed. For a general process, the quantum operator that change a state of the system should be factorized by the subsystem’s operators. For instance, a control not operation in quantum computation given by (in fact, control not is a unitary evolution operator)
$$O=|0_10|1_2+|1_11|\sigma _2^x,$$
(3.6)
where $`1_2`$ is the unit operator for the second qubit, $`\sigma _2^x`$ stands for the $`x`$ pauli matrix of the second qubit. $`|1_1`$ and $`|0_1`$ represent two state of the first qubit. This control not operator is factorisable, i.e. $`O`$ can be written in the form
$$O=\underset{i}{}O_1^iO_2^i,$$
where $`O_1^i`$ and $`O_2^i`$ denote operators for the first and second qubit, respectively. Hence, according to the definition eq.(3.5), the control not operator in the form (3.6) does not change the entanglement of the system.
We would like to point out that the discussions presented here are for the unitary evolutions, for non-unitary evolution such as a trace over some of the degree of freedom, we should find a auxiliary unitary process instead of the non-unitary one.
## 4 Example and Conclusion
In order to understand how our program for calculating the amount of entanglement change works, we present in this section two examples, one of them consists of two interacting qubits (two identical two-level system) in a laser beam\[5a\], and the another two independent qubits coupling simultaneously to a bath.
example 1:
The Hamiltonian describing the system in this example has the following form(set $`\mathrm{}=1`$):
$`H`$ $`=`$ $`H_A+H_B+H_{int}+H_f,`$
$`H_i`$ $`=`$ $`{\displaystyle \frac{1}{2}}\omega \sigma _z^i,(i=A,B)`$
$`H_{int}`$ $`=`$ $`g(\sigma _A^+\sigma _B^{}+\sigma _A^{}\sigma _B^+)+\lambda ({\displaystyle \underset{i=A,B}{}}\sigma _i^+a+\sigma _i^{}a^+),`$
$`H_f`$ $`=`$ $`\omega _fa^+a,`$ (4.1)
where $`\sigma _i^z,\sigma _i^{},\sigma _i^+`$ describe the pauli operator of the $`i`$ qubit, $`g`$ denotes the coupling constant, and $`H_f`$ stands for the free Hamiltonian of the laser beam. Suppose the state is initially in the form
$`\rho (0)`$ $`=`$ $`\rho _A(0)\rho _B(0)\rho _f(0),`$
$`\rho _A(0)\rho _B(0)`$ $`=`$ $`|e_A,e_Be_A,e_B|,`$
$`\rho _f(0)`$ $`=`$ $`{\displaystyle \underset{n}{}}p(n)|nn|,`$ (4.2)
where $`|e_i`$ denotes the excited state of the qubit $`i`$ and $`|n`$ stands for a Fock state of the laser beam field. In the Schrödinger picture, the density operator that obeys von Neumann equation is given by
$$\rho _{AB}(t)=Tr_f\rho (t)=$$
$$\left(\begin{array}{ccc}_np^2\left(n\right)f_{gg}^2(n,t)\hfill & _np\left(n+1\right)p\left(n\right)f_{gg}(n+1,t)f_{EG}(n,t)\hfill & _np\left(n+2\right)p\left(n\right)f_{ee}\left(n+2\right)f_{gg}\left(n\right)\hfill \\ _np\left(n+1\right)p\left(n\right)f_{gg}(n+1,t)f_{EG}(n,t)\hfill & _np^2\left(n\right)f_{EG}^2(n,t)\hfill & _np\left(n\right)p\left(n+1\right)f_{EG}\left(n+1\right)f_{ee}\left(n\right)\hfill \\ _np\left(n+2\right)p\left(n\right)f_{ee}\left(n+2\right)f_{gg}\left(n\right)\hfill & _np\left(n\right)p\left(n+1\right)f_{EG}\left(n+1\right)f_{ee}\left(n\right)\hfill & _np^2\left(n\right)f_{ee}^2(n,t)\hfill \end{array}\right),$$
(4.3)
where we take $`|g_A,g_B,|E,G,`$ and $`|e_A,e_B`$ as a set of basis, and
$$|g_A,g_B=|g_A|g_B,|e_A,e_B=|e_A|e_B,|E,G=\frac{1}{\sqrt{2}}(|g_A|e_B+|e_A|g_B),$$
$$\rho _A^0(t)\rho _B^0(t)=\rho _A^0(0)\rho _B^0(0),$$
$$f_{gg}(n,t)=\frac{1}{4}\mathrm{sin}2\varphi \mathrm{sin}\theta e^{iE_+t}+\frac{1}{2}\mathrm{sin}2\varphi \mathrm{cos}^2\frac{\theta }{2}e^{iE_{}t}\frac{1}{2}\mathrm{sin}2\varphi e^{iE_0t},$$
$$f_{ee}(n,t)=\mathrm{sin}^2\varphi \mathrm{sin}^2\frac{\theta }{2}e^{iE_+t}+\mathrm{cos}^2\varphi e^{iE_0t}+\mathrm{sin}^2\varphi \mathrm{cos}^2\frac{\theta }{2}e^{iE_{}t},$$
$$f_{EG}(n,t)=\mathrm{sin}\theta \mathrm{sin}\varphi \mathrm{sin}\frac{\mathrm{\Omega }t}{2},$$
$$E_\pm =\mathrm{\Omega }\frac{\mathrm{cos}\theta \pm 1}{2}+\omega _f(n+1),E_0=(n+1)\omega _f,$$
and$`\theta =\frac{\pi }{2},\mathrm{tan}\varphi =\sqrt{\frac{n+2}{n+1}},\mathrm{\Omega }^2=(16n+24)g^2.`$ Eqs.(3.5) and (4.3) together give
$$\delta D(\rho )=\underset{i,j=1,2,3}{}(\rho _{AB}^{ij})^22\rho _{AB}^{33}+1,$$
(4.4)
where $`\rho _{AB}^{ij}`$ denotes the element of matrix $`\rho _{AB}`$ given by eq.(4.3), which represents the entanglement change or entanglement of subsystems $`A`$ and $`B`$ at time $`t`$.
example 2
The Hamiltonian describing dissipation of the two qubits has the following form\[1d\](setting $`\mathrm{}=1`$)
$`H`$ $`=`$ $`\omega _0(\sigma _a^z+\sigma _b^z)+{\displaystyle \underset{l=a,b}{}}{\displaystyle 𝑑\omega [g_{\omega l}A_l(a_{\omega l}^++a_{\omega l})]}`$ (4.5)
$`+`$ $`{\displaystyle 𝑑\omega }_{l=a,b}(\omega a_{\omega l}^+a_{\omega l}),`$
where $`\stackrel{}{\sigma }_i`$ describe the pauli’s matrix of the $`i`$ qubit, $`a_{\omega l}`$ stands for the bath mode $`\omega `$ coupling to the $`l`$ qubit, and $`_{l=a,b}a_{\omega l}^+a_{\omega l}=a_{\omega l}^+a_{\omega l}`$ for $`a_{\omega a}=a_{\omega b}`$, whereas $`_{l=a,b}a_{\omega l}^+a_{\omega l}=a_{\omega a}^+a_{\omega a}+a_{\omega b}^+a_{\omega b}`$ for $`a_{\omega a}a_{\omega b}`$. The coupling coefficients are denoted by $`g_{\omega l}`$, and the qubit operator $`A_l`$ in general is expressed as a linear superposition of three pauli’s operators, i.e. $`A_l=\lambda ^{(1)}\sigma _l^x+\lambda ^{(2)}\sigma _l^y+\lambda ^{(3)}\sigma _l^z`$. The ratio $`\lambda ^{(1)}:\lambda ^{(2)}:\lambda ^{(3)}`$ is determined by the type of the dissipation. For instance, $`\lambda ^{(1)}=\lambda ^{(2)}=0`$ for phase damping and $`\lambda ^{(3)}=0`$ for amplitude damping. Phase damping induces pure dephasing, whereas amplitude damping induces loss and dephasing simultaneously. Many source of decoherence in quantum computers are described by amplitude damping.
Without any loss of generality, we discuss in detail the case with $`\lambda ^{(1)}=\lambda ^{(2)}=0`$ i.e. phase damping. Some words of caution are now in order. As mentioned above, the bath can also cause another unwanted effect in computation process, i.e. amplitude dissipation. It is easy, however, to make system have small loss rate of amplitude dissipation, so a considerable number of operations are allowed to perform.
In the case of phase damping, the unitary evolution operator may be factorized in the following form
$$U(t)=U_0(t)U_I(t),$$
(4.6)
where $`U_0(t)=e^{iH_0t}`$ with $`H_0=\omega _0(\sigma _a^z+\sigma _b^z)+𝑑\omega _{l=a,b}(\omega a_{\omega l}^+a_{\omega l})`$ is the free evolution operator, while $`U_I(t)`$ denotes the evolution operator in the interaction picture. A readily calculation shows that
$`U_I(t)=U_I^a(t)U_I^b(t),`$
$`U_I^i(t)=u_I^i(t)|e_i><e_i|+v_I^i|g_i><g_i|,`$ (4.7)
where $`|e_i>`$ and $`|g_i>`$ are the eigenstates of $`\sigma _i^z`$ with eigenvalues $`+1`$ and $`1`$, respectively, and $`u_I^i`$ and $`v_I^i`$ satisfy($`i=a,b`$):
$`i\mathrm{}{\displaystyle \frac{}{t}}u_I^i`$ $`=`$ $`{\displaystyle d\omega g_{\omega i}(a_{\omega i}^+e^{i\omega t}+a_{\omega i}e^{i\omega t})u_I^i},`$
$`i\mathrm{}{\displaystyle \frac{}{t}}v_I^i`$ $`=`$ $`{\displaystyle d\omega g_{\omega i}(a_{\omega i}^+e^{i\omega t}+a_{\omega i}e^{i\omega t})v_I^i}.`$ (4.8)
The Wei-Norman’s algebraic method that provides a way to factorize the evolution operator gives
$$u_I^i=\underset{\omega }{}e^{f_\omega ^i(t)}e^{A_\omega ^i(t)a_\omega ^+}e^{B_\omega ^i(t)a_\omega }$$
(4.9)
and
$$v_I^i=\underset{\omega }{}e^{h_\omega ^i(t)}e^{C_\omega ^i(t)a_\omega ^+}e^{D_\omega ^i(t)a_\omega }.$$
(4.10)
Here,
$$A_\omega ^i(t)=\frac{g_{\omega i}}{\omega }(e^{i\omega t}1),$$
$$B_\omega ^i(t)=(A_\omega ^i(t))^{},$$
$$f_\omega ^i(t)=i\frac{g_{\omega i}^2}{\omega }t+\frac{g_{\omega i}^2}{\omega ^2}(1e^{i\omega t}),$$
$$C_\omega ^i(t)=A_\omega ^i(t),D_\omega ^i(t)=B_\omega ^i(t),h_\omega ^i(t)=f_\omega ^i(t).$$
Now, we turn our attention to compute the reduced density operator of the two-qubit system, first of all, we calculate the total density operator, which follows straightforwardly from eq.(4.6)
$$\rho (t)=U_0U_I\rho (0)U_I^+U_0^+,$$
(4.11)
where $`\rho (0)`$ denotes the initial density operator (state), which may be written in a separable form
$$\rho (0)=\rho _a(0)\rho _b(0)\rho _B(0).$$
Here, $`\rho _i(0)(i=a,b)`$ represents the initial state of qubit $`i`$, and $`\rho _B(0)`$ stands for the initial state of the bath. In following, we use the notation,$`|e_a,g_b`$ to indicate the eigenstates of $`\sigma _a^z`$ and $`\sigma _b^z`$ with eigenvalues $`1`$ and $`1`$, $`Tr_B`$ indicate a trace over the bath, and $`\rho _i^0(t)`$ to represent the free two-qubit state i.e.
$$\rho ^0(t)=Tr_BU_0(t)\rho (0)U_0^+(t).$$
With this notation, in a subspace spanned by $`\{|11=|e_a,e_b,|12=|e_a,g_b,|21=|g_a,b_2,|22=|g_a,g_b\}`$, the state of the two-qubit system at time t takes the form
$$\rho _{ab}(t)=Tr_B\rho (t)=\left(\begin{array}{cccc}\rho _{1111}\hfill & \rho _{1112}& \rho _{1121}& \hfill \rho _{1122}\\ \rho _{1211}\hfill & \rho _{1212}& \rho _{1221}& \hfill \rho _{1222}\\ \rho _{2111}\hfill & \rho _{2112}& \rho _{2121}& \hfill \rho _{2122}\\ \rho _{2211}\hfill & \rho _{2212}& \rho _{2221}& \hfill \rho _{2222}\end{array}\right),$$
(4.12)
where $`\rho _{ijkl}=\rho _{ijkl}^0F_{ijkl}(i,j,k,l=1,2)`$, $`\rho _{ijkl}^0=`$ $`Tr_Bij|\rho ^0(t)|kl,`$ and $`F_{ijkl}=Tr_B`$ $`ij|_{c,d,e,f=1}^2(U_I^a)_i`$ $`(U_I^b)_j|cd`$ $`ef|((U_I^b)^+)_k((U_I^a)^+)_l|kl.`$ Here, $`(U_I^a)_i=_ai|U_I^a|i_a,`$ and $`|2_a=|e_a`$, $`|1_a=|g_a.`$ In order to get more information about the reduced density operator, we make some discussion on the quantity $`F_{ijkl}`$. It can be easily verified that $`F_{ijkl}=1`$ for $`i=k,j=l`$, while $`F_{ijkl}=F_{klij}^{}`$ for $`ik`$ and $`jl`$. Moreover, the quantity results from the interaction between the two-qubit system and the bath, hence it depends on the states of the bath. Although different bathes result in different results $`F_{ijkl}`$, the physical results discussed here do not rely on the bath. In this sense, we may consider a simple case with zero temperature. In this case, $`F_{ijkl}`$ is given that
$$F_{ijkl}=F_{ijkl}(t)=e^{_0^{\mathrm{}}[\mathrm{\Delta }_{ik}(\omega ,t)+\mathrm{\Delta }_{jl}^{}(\omega ,t)]\rho (\omega )𝑑\omega },$$
(4.13)
where $`\mathrm{\Delta }_{ij}(\omega ,t)=2\frac{(g_{\omega i}g_{\omega j})^2\mathrm{sin}^20.5\omega t}{\omega ^2},`$ $`\rho (\omega )`$ stands for the bath spectrum distribution. The eq.(4.13) suggests that $`F_{ijkl}`$ approaches zero with the passage of time except some moments at which $`_0^{\mathrm{}}[\mathrm{\Delta }_{ij}(\omega ,t)\mathrm{\Delta }_{kl}^{}(\omega ,t)]\rho (\omega )𝑑\omega =0.`$ This attractive results might be used in preventing information loss stored in quantum states. Now we come back to the entanglement change, eqs. (3.5) and (4.12) together give
$$\delta D(\rho _{ab}(t)||\rho _a^0(t)\rho _b^0(t))=\underset{i,j,k,l=1}{\overset{2}{}}(\rho _{ijkl}\rho _{ijkl}^0)(\rho _{klij}\rho _{klij}^0).$$
(4.14)
In summary, we propose a new method to compute the entanglement change in a dynamical process.We see the above treatment in Sec.2 and Sec. 3 does not refer to specific entangled systems. This is a desired property as it makes our measure of entanglement universal. Especially, the results yielded by present paper can be easily generalized to more than two subsystems, this is just the case of many qubits interacting simultaneously with environment. In addition to the measure stated above, the quantum relative entropy defined as
$$D(\rho (t)||\rho _A^0(t)\rho _B^0(t))=Tr[\rho (t)(ln\rho (t)ln\rho _A^0(t)\rho _B^0(t))]$$
and the Bures metric given by
$$D(\rho (t)||\rho _A^0(t)\rho _B^0(t))=22\sqrt{F(\rho ,\rho _A^0\rho _B^0)},$$
with $`F(\rho ,\rho _A^0\rho _B^0)=[Tr[\sqrt{\rho _B^0\rho _A^0}\rho \sqrt{\rho _B^0\rho _A^0}]^{\frac{1}{2}}]^2`$ are other good measures of entanglement. With this modified definitions, the measures of entanglement can be given in easy way.
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# Many pion decays of 𝜌(770) and 𝜔(782) mesons in chiral theory.
## I Introduction
Despite the lack of a straightforward derivation from first principles of QCD, the effective Lagrangians that describe the low-energy interactions of the ground state octet of pseudoscalar mesons $`\pi `$, $`K`$, $`\eta `$ are constructed upon treating these mesons as the Goldstone bosons of the spontaneously broken chiral $`U_L(3)\times U_R(3)`$ symmetry of the massless three-flavored QCD Lagrangian. The key point in this task is that the transformation properties of the Goldstone fields under the nonlinear realization of chiral symmetry are sufficient for the establishing the most general form of the effective Lagrangian . As far as vector mesons are concerned, the situation is not so clear, because the vector mesons, contrary to the pseudoscalar ones, cannot be considered as the Goldstone bosons of the spontaneously broken symmetry. For this reason there exist different schemes of including these mesons into the effective chiral Lagrangians. The scheme of Ref. where the vector mesons are treated as the dynamical gauge bosons of hidden local symmetry (HLS), incorporates these mesons into the effective chiral Lagrangian in a most elegant way. However, a straightforward comparison with the data of the predictions of these models immediately meets with difficulty, since the dominant decay modes of all the ground state vector mesons produce pseudoscalar mesons that are far from being soft. This means that both the higher derivatives in the low energy expansion of the effective Lagrangians and the loop quantum corrections should be taken into account.
We suggest here many pion decay modes of $`\rho (770)`$ and $`\omega (782)`$ mesons as a test ground for the chiral models of vector mesons . Theoretically, pions in the final state in such decays are soft to the extent to be specified below. This fact allows one to neglect the higher derivative and loop corrections in the effective Lagrangian. On the other hand, the unconventional, from the point of view of the chiral pion dynamics, sources of soft pions are feasible. Indeed, the progress in increasing the intensity of low energy $`e^+e^{}`$ colliders ($`\varphi `$ factories) , photon beams , and a huge number of the specific hadronic decays of $`\tau `$ leptons could offer the naturally controlled sources of soft pions, provided the sufficiently low invariant mass regions of hadronic systems are isolated. The yield of pions is considerably enhanced when they are produced through the proper vector resonance states, which will hopefully offer the possibility of testing the mentioned models.
Pions emitted in the decay $`\rho 4\pi `$ are rather soft, because their typical momentum is $`|𝐩|m_\pi `$, where $`m_\pi `$ is the pion mass. By this reason this decay attracts much attention from the point of view of the study of the predictions of the effective chiral Lagrangians for vector mesons. As was found in Refs. , the above decay should be rather strong, B$`(\rho 4\pi )10^4`$. The calculations of Refs. were analyzed in detail in Ref. , where a number of shortcomings of the former related with the actual violation of chiral invariance, in particular, the Adler condition for soft pions, was uncovered. The correct results based on the amplitudes obeying the Adler condition and obtained in Refs. , correspond to B$`(\rho 4\pi )10^5`$. The large magnitude of the branching ratio B$`(\rho 4\pi )10^4`$ obtained in Ref. is related, in all appearance, with a very rough method of calculation. A common drawback of Refs. is that their authors evaluate the partial width at the only energy equal to the mass of the $`\rho `$, as if the latter were a genuine narrow peak. However, the fact that the width of the $`\rho `$ resonance is rather large, and $`\mathrm{\Gamma }(\rho 4\pi ,E)`$ rises rapidly with the energy increase even at energies inside the $`\rho `$ peak, forces one to think that the magnitude of the $`4\pi `$ partial width at the $`\rho `$ mass cannot be an adequate characteristics of the dynamics of the process. In this respect, it is just the resonance excitation curve in the channel $`e^+e^{}\rho ^04\pi `$ which is of much interest, being a test ground of various chiral models of the decay under consideration.
Pions emitted in the decay $`\omega 5\pi `$ are truly soft, because they possess the typical momentum $`|𝐩|0.5m_\pi `$. Hence, the lowest order terms obtained upon neglecting the higher derivatives and loop corrections should give the reliable results.
The aim of the present paper is to consider in detail the many pion decay modes of the vector $`\rho (770)`$ and $`\omega (782)`$ mesons in the framework of the effective chiral Lagrangian approach. The first Lagrangian of this kind incorporating the isotopic triplets of the $`\rho `$ meson field $`𝝆_\mu `$, the pion field $`𝝅`$, and their interaction was proposed by Weinberg under the nonlinear realization of the chiral symmetry
$`^{(\rho +\pi )}`$ $`=`$ $`{\displaystyle \frac{1}{4}}\left(_\mu 𝝆_\nu _\nu 𝝆_\mu +g[𝝆_\mu \times 𝝆_\nu ]\right)^2`$ (3)
$`+{\displaystyle \frac{m_\rho ^2}{2}}\left[𝝆_\mu +{\displaystyle \frac{[𝝅\times _\mu 𝝅]}{2gf_\pi ^2(1+𝝅^2/4f_\pi ^2)}}\right]^2`$
$`+{\displaystyle \frac{(_\mu 𝝅)^2}{2\left(1+𝝅^2/4f_\pi ^2\right)^2}}{\displaystyle \frac{m_\pi ^2𝝅^2}{2(1+𝝅^2/4f_\pi ^2)}},`$
where $`f_\pi =92.4`$ MeV is the pion decay constant, and the cross denotes the vector product in the isovector space. As was shown, the Lagrangian Eq. (3) results from the HLS approach . The $`\rho \rho \rho `$ coupling constant $`g`$ and the $`\rho \pi \pi `$ coupling constant $`g_{\rho \pi \pi }`$ are related to the $`\rho `$ mass $`m_\rho `$ and the pion decay constant $`f_\pi `$ via the parameter of hidden local symmetry $`a`$ as
$`g`$ $`=`$ $`m_\rho /f_\pi \sqrt{a},`$ (4)
$`g_{\rho \pi \pi }`$ $`=`$ $`\sqrt{a}m_\rho /2f_\pi .`$ (5)
Note that $`a=2`$, if one demands the universality condition $`g=g_{\rho \pi \pi }`$ to be satisfied. Then the so called Kawarabayashi-Suzuki-Riazzuddin-Fayyazuddin (KSRF) relation arises
$$2g_{\rho \pi \pi }^2f_\pi ^2/m_\rho ^2=1$$
(6)
which beautifully agrees with experiment. The $`\rho \pi \pi `$ coupling constant resulting from this relation is $`g_{\rho \pi \pi }=5.89`$. The Lagrangian of Eq. (3) should be added with the kinetic and mass terms of the isosinglet $`\omega (782)`$ field $`\omega _\mu `$,
$$^{(\omega )}=\frac{1}{4}(_\mu \omega _\nu _\nu \omega _\mu )^2+\frac{m_\omega ^2}{2}\omega _\mu ^2,$$
(7)
and the term describing the interaction of $`\omega `$ with $`𝝆`$ and $`𝝅`$. The latter comes from the term induced by the anomalous Lagrangian of Wess and Zumino ,
$$^{(\omega \rho \pi )}=\frac{N_cg^2}{8\pi ^2f_\pi }\epsilon _{\mu \nu \lambda \sigma }_\mu \omega _\nu \left(𝝅_\lambda 𝝆_\sigma \right),$$
(8)
where $`N_c=3`$ is the number of colors, and $`\epsilon _{\mu \nu \lambda \sigma }`$ is the antisymmetric unit tensor. In agreement with Ref. , the contribution of the pointlike vertex $`\omega 3\pi `$ is omitted. See, however, the discussion following Eq. (88) in Sec. IV where the role of this vertex in the $`\omega 5\pi `$ decay is briefly discussed.
The HLS approach permits one to include the axial mesons as well . An ideal treatment would consist of following the line of reasoning that under the assumption of $`m_\rho Em_{a_1}`$, the difference between the models with and without $`a_1`$ meson were reduced to the taking into account the higher derivatives expected to be small. In the real life, one has $`m_{a_1}^2m_\rho ^2m_\rho ^2`$, and the correction may appear to be appreciable even at the $`\rho `$ mass. In fact, the calculation shows that the corrections amount to $`20\text{ - }30\%`$ in the width. This means, in particular, that the left shoulder of the $`\rho `$ peak, where the contributions of higher derivatives are vanishing rapidly, is the best place to work. In the present paper, we do not take into account the $`a_1`$ meson contribution in the $`\rho 4\pi `$ decay amplitude. In the meantime, such a contribution is negligible, as will be clear later on, in the $`\omega 5\pi `$ decay amplitude.
The rest of the paper is devoted to the working out the consequences of the above Lagrangians for the many pion decays of the $`\rho (770)`$ and $`\omega (782)`$ mesons. Specifically, the partial widths and resonance excitation curves are calculated for the reactions $`e^+e^{}\rho ^02\pi ^+2\pi ^{}`$ and $`e^+e^{}\rho ^0\pi ^+\pi ^{}2\pi ^0`$. It is shown that the intensities of the above decays change as fast as two times the phase space variation, upon the energy variation inside the $`\rho `$ widths. All this means that $`e^+e^{}`$ offers an ideal tool for the study of such effects. The decay widths of charged $`\rho `$ meson, $`\rho ^\pm \pi ^\pm 3\pi ^0`$ and $`\rho ^\pm 2\pi ^\pm \pi ^{}\pi ^0`$, as well as $`\omega `$ meson, $`\omega 2\pi ^+2\pi ^{}\pi ^0`$ $`\omega \pi ^+\pi ^{}3\pi ^0`$, are also evaluated. We choose these particular modes because the final pions in the decay $`\rho 4\pi `$ are practically soft while in the decay $`\omega 5\pi `$ they are soft. Note that the final pions in the G-parity violating decay $`\rho 3\pi `$ are not sufficiently soft to compare the calculated branching ratio with the predictions of chiral models. The final pions in the decay $`\omega 4\pi `$ are sufficiently soft to make such a comparison meaningful. However, the branching ratios of both above G-parity violating decays $`\rho 3\pi `$ and $`\omega 4\pi `$ are determined completely by the $`\omega \rho `$ transition amplitude and by the well known decay $`\omega 3\pi `$ and the decay $`\rho 4\pi `$, respectively.
The following material is organized as follows. Section II contains the expressions for the $`\rho 4\pi `$ amplitudes. The results of calculation of the excitation curves and partial widths for different isotopic states of four pions are presented in Sec. III. This task is fulfilled in the cases of $`e^+e^{}`$ annihilation, $`\tau `$ decays, and photoproduction. The partial widths of the decays $`\omega 5\pi `$ are discussed in Sec. IV. Although an appreciable part of the material of Sec. IV is contained in Ref. , we include here the basic results from that paper in order to keep the presentation self-contained. Section V contains concluding remarks. The angular distributions of the various combinations of emitted pions in the decays $`\rho 4\pi `$ and $`\omega 5\pi `$ obtained for the case of $`e^+e^{}`$ annihilation are presented in the Appendix.
## II Amplitudes of $`\rho 4\pi `$ decay
The amplitudes of the decays of our interest are obtained from the Weinberg Lagrangian Eq. (3). First, let us obtain the $`\pi 3\pi `$ transition amplitudes necessary for the calculation of the many-pion decays of vector mesons. They are given by the diagrams shown in Fig. 1(a) and look as
$`M(\pi ^+\pi _{q_1}^+\pi _{q_2}^+\pi _{q_3}^{})`$ $`=`$ $`(1+P_{12}){\displaystyle \frac{1}{2f_\pi ^2}}\left\{2(q_1,q_2)+a(q_1,q_2q_3)\left[1{\displaystyle \frac{m_\rho ^2}{D_\rho (q_2+q_3)}}\right]\right\},`$ (9)
$`M(\pi ^+\pi _{q_1}^+\pi _{q_2}^0\pi _{q_3}^0)`$ $`=`$ $`(1+P_{23}){\displaystyle \frac{1}{2f_\pi ^2}}\left\{(q_3,q_12q_2)+a(q_3,q_1q_2)\left[1{\displaystyle \frac{m_\rho ^2}{D_\rho (q_1+q_2)}}\right]\right\},`$ (10)
$`M(\pi ^0\pi _{q_1}^+\pi _{q_2}^{}\pi _{q_3}^0)`$ $`=`$ $`(1+P_{12}){\displaystyle \frac{1}{2f_\pi ^2}}\left\{(q_1,q_22q_3)a(q_1,q_2q_3)\left[1{\displaystyle \frac{m_\rho ^2}{D_\rho (q_2+q_3)}}\right]\right\},`$ (11)
$`M(\pi ^0\pi _{q_1}^0\pi _{q_2}^0\pi _{q_3}^0)`$ $`=`$ $`{\displaystyle \frac{1}{f_\pi ^2}}\left[(q_1,q_2)+(q_1,q_3)+(q_2,q_3)\right].`$ (12)
where $`P_{ij}`$ stands for the operator of the interchange of the pion momenta $`q_i`$ and $`q_j`$, and
$$D_\rho (q)=m_\rho ^2q^2im_\rho ^2\left(\frac{q^24m_\pi ^2}{m_\rho ^24m_\pi ^2}\right)^{3/2}\frac{\mathrm{\Gamma }_{\rho \pi \pi }(m_\rho ^2)}{\sqrt{q^2}}$$
(13)
is the inverse propagator of $`\rho `$ meson. Our notation for the Lorentz invariant scalar product of two four vectors $`a`$ and $`b`$ hereafter is $`(a,b)=a_0b_0𝐚𝐛`$. As it will be clear later on, the nonrelativistic expressions for the above amplitudes are needed. They are obtained upon neglecting the space components of the pion four momenta and look as
$`M(\pi ^+\pi _{q_1}^+\pi _{q_2}^+\pi _{q_3}^{})`$ $`=`$ $`{\displaystyle \frac{2m_\pi ^2}{f_\pi ^2}},`$ (14)
$`M(\pi ^+\pi _{q_1}^+\pi _{q_2}^0\pi _{q_3}^0)`$ $`=`$ $`{\displaystyle \frac{m_\pi ^2}{f_\pi ^2}},`$ (15)
$`M(\pi ^0\pi _{q_1}^+\pi _{q_2}^{}\pi _{q_3}^0)`$ $`=`$ $`{\displaystyle \frac{m_\pi ^2}{f_\pi ^2}},`$ (16)
$`M(\pi ^0\pi _{q_1}^0\pi _{q_2}^0\pi _{q_3}^0)`$ $`=`$ $`{\displaystyle \frac{3m_\pi ^2}{f_\pi ^2}}.`$ (17)
Note that the HLS parameter $`a`$ drops from the expressions in the nonrelativistic limit.
The diagrams representing the amplitudes of the decays $`\rho 4\pi `$ for different combinations of the charges of final pions are shown in Fig. 1(b) and (c). Introducing the four vector of polarization of the decaying $`\rho `$ meson $`\epsilon _\mu `$ one can write the general expression for the amplitude in the form
$$M=\frac{g_{\rho \pi \pi }}{f_\pi ^2}\epsilon _\mu J_\mu ,$$
where
$$g_{\rho \pi \pi }/f_\pi ^2=\sqrt{a}m_\rho /2f_\pi ^3$$
(18)
results from Eq. (5). Let us give the expressions for the current $`J_\mu `$ for all the decay modes considered here.
(1) The decay $`\rho ^0(q)\pi ^+(q_1)\pi ^+(q_2)\pi ^{}(q_3)\pi ^{}(q_4)`$. One has
$`J_\mu `$ $`=`$ $`(1+P_{12})(1+P_{34})\{q_{1\mu }[{\displaystyle \frac{1}{2}}+{\displaystyle \frac{a(q_2,q_3)(a2)(q_3,q_4)}{D_\pi (qq_1)}}]`$ (21)
$`+q_{3\mu }\left[{\displaystyle \frac{1}{2}}+{\displaystyle \frac{a(q_1,q_4)(a2)(q_1,q_2)}{D_\pi (qq_3)}}\right]`$
$`+am_\rho ^2(1+P_{13}){\displaystyle \frac{q_{1\mu }(q_3,q_2q_4)}{D_\pi (qq_1)D_\rho (q_2+q_4)}}\}.`$
Hereafter $`D_\pi (q)=m_\pi ^2q^2`$ is the inverse propagator of pion.
(2) The decay $`\rho ^0(q)\pi ^+(q_1)\pi ^{}(q_2)\pi ^0(q_3)\pi ^0(q_4)`$. In this case one has $`J_\mu =J_\mu ^{\mathrm{nan}}+J_\mu ^{\mathrm{an}}`$, where
$`J_\mu ^{\mathrm{nan}}`$ $`=`$ $`(1P_{12})(1+P_{34})q_{1\mu }\{{\displaystyle \frac{1}{4}}+{\displaystyle \frac{1}{D_\pi (qq_1)}}[(a1)(q_3,q_4)(a2)(q_2,q_3)`$ (25)
$`+am_\rho ^2{\displaystyle \frac{(q_3,q_2q_4)}{D_\rho (q_2+q_4)}}]\}`$
$`+(1+P_{34}){\displaystyle \frac{m_\rho ^2}{2D_\rho (q_1+q_3)D_\rho (q_2+q_4)}}[(q_1+q_3q_2q_4)_\mu (q_1q_3,q_2q_4)`$
$`2(q_1q_3)_\mu (q_1+q_3,q_2q_4)+2(q_2q_4)_\mu (q_2+q_4,q_1q_3)]`$
is obtained from Eq. (3), while the contribution of the term induced by the anomalous Lagrangian of Wess and Zumino manifesting in the process $`\rho ^0\omega \pi ^0\pi ^+\pi ^{}\pi ^0\pi ^0`$, is given by the expression
$`J_\mu ^{\mathrm{an}}`$ $`=`$ $`2\left({\displaystyle \frac{N_cg^2}{8\pi ^2}}\right)^2(1+P_{34})[q_{1\mu }(1P_{23})(q,q_2)(q_3,q_4)`$ (28)
$`+q_{2\mu }(1P_{13})(q,q_3)(q_1,q_4)+q_{3\mu }(1P_{12})(q,q_1)(q_2,q_4)]`$
$`\times \left[{\displaystyle \frac{1}{D_\rho (q_1+q_2)}}+{\displaystyle \frac{1}{D_\rho (q_1+q_3)}}+{\displaystyle \frac{1}{D_\rho (q_2+q_3)}}\right]{\displaystyle \frac{1}{D_\omega (qq_4)}},`$
where $`D_\omega (q)=m_\omega ^2q^2`$ is the inverse $`\omega `$ meson propagator. In general, this term is attributed to the contribution of higher derivatives. Nevertheless, we take it into account to show the effect of the latter and the dynamical effect of the opening of the channel $`\rho \omega \pi 4\pi `$. The following amplitudes of the charged $`\rho `$ decay are necessary for obtaining the $`\omega 5\pi `$ decay amplitude, and are of their own interest when studying the reactions of peripheral $`\rho `$ meson production and $`\tau `$ decays.
(3) The decay $`\rho ^+(q)\pi ^+(q_1)\pi ^0(q_2)\pi ^0(q_3)\pi ^0(q_4)`$. One has
$`J_\mu `$ $`=`$ $`(1+P_{24}+P_{34})\{2q_{1\mu }[{\displaystyle \frac{1}{3}}+{\displaystyle \frac{(q_2,q_3)}{D_\pi (qq_1)}}]{\displaystyle \frac{q_{4\mu }}{D_\pi (qq_4)}}[2(a1)(q_2,q_3)`$ (30)
$`(a2)(q_1,q_2+q_3)+am_\rho ^2(1+P_{23}){\displaystyle \frac{(q_2,q_1q_3)}{D_\rho (q_1+q_3)}}]\}.`$
(4) The decay $`\rho ^+(q)\pi ^+(q_1)\pi ^+(q_2)\pi ^{}(q_3)\pi ^0(q_4)`$. Here, the contribution induced by the anomalous Lagrangian of Wess and Zumino is also possible, hence $`J_\mu =J_\mu ^{\mathrm{nan}}+J_\mu ^{\mathrm{an}}`$, where
$`J_\mu ^{\mathrm{nan}}`$ $`=`$ $`(1+P_{12})\{{\displaystyle \frac{1}{2}}(q_1q_4)\mu +{\displaystyle \frac{q_{1\mu }(1+P_{23})}{D_\pi (qq_1)}}[(a1)(q_2,q_3)(a2)(q_2,q_4)]`$ (35)
$`{\displaystyle \frac{q_{4\mu }}{D_\pi (qq_4)}}\left[a(q_1,q_3)(a2)(q_1,q_2)\right]`$
$`am_\rho ^2\left[{\displaystyle \frac{q_{1\mu }}{D_\pi (qq_1)}}(1+P_{23}){\displaystyle \frac{(q_2,q_3q_4)}{D_\rho (q_3+q_4)}}+{\displaystyle \frac{q_{4\mu }(q_1,q_2q_3)}{D_\pi (qq_4)D_\rho (q_2+q_3)}}\right]`$
$`+{\displaystyle \frac{m_\rho ^2}{2D_\rho (q_1+q_3)D_\rho (q_2+q_4)}}[(q_1+q_3q_2q_4)_\mu (q_1q_3,q_2q_4)`$
$`2(q_1q_3)_\mu (q_1+q_3,q_2q_4)+2(q_2q_4)_\mu (q_1q_3,q_2+q_4)]\}`$
is obtained from Eq. (3), while the term induced by the anomaly looks as
$`J_\mu ^{\mathrm{an}}`$ $`=`$ $`2\left({\displaystyle \frac{N_cg^2}{8\pi ^2}}\right)^2(1+P_{23})[q_{1\mu }(1P_{24})(q,q_4)(q_2,q_4)`$ (38)
$`+q_{2\mu }(1P_{14})(q,q_1)(q_3,q_4)+q_{4\mu }(1P_{12})(q,q_2)(q_1,q_3)]`$
$`\times \left[{\displaystyle \frac{1}{D_\rho (q_1+q_2)}}+{\displaystyle \frac{1}{D_\rho (q_1+q_4)}}+{\displaystyle \frac{1}{D_\rho (q_2+q_4)}}\right]{\displaystyle \frac{1}{D_\omega (qq_3)}}.`$
One can verify that up to the corrections of the order of $`m_\pi ^2/m_\rho ^2`$, the above written amplitudes vanish in the limit of vanishing 4-momentum of each final pion. In other words, they obey the Adler condition.
It is useful to obtain the nonrelativistic expressions for the $`\rho 4\pi `$ decay amplitudes which are relevant for the four pion invariant mass below 700 MeV. This can be made upon neglecting the space components of the pion momenta. One can convince oneself that $`a`$ enters the expressions for amplitudes as an overall factor Eq. (18) in this limit, so that the latter look as
$`M(\rho ^0\pi _{q_1}^+\pi _{q_2}^+\pi _{q_3}^{}\pi _{q_4}^{})`$ $``$ $`{\displaystyle \frac{g_{\rho \pi \pi }}{2f_\pi ^2}}(\epsilon ,q_1+q_2q_3q_4),`$ (39)
$`M(\rho ^0\pi _{q_1}^+\pi _{q_2}^{}\pi _{q_3}^0\pi _{q_4}^0)`$ $``$ $`{\displaystyle \frac{g_{\rho \pi \pi }}{4f_\pi ^2}}(\epsilon ,q_1q_2),`$ (40)
$`M(\rho ^+\pi _{q_1}^+\pi _{q_2}^0\pi _{q_3}^0\pi _{q_4}^0)`$ $``$ $`{\displaystyle \frac{g_{\rho \pi \pi }}{f_\pi ^2}}(\epsilon ,q_1),`$ (41)
$`M(\rho ^+\pi _{q_1}^+\pi _{q_2}^+\pi _{q_3}^{}\pi _{q_4}^0)`$ $``$ $`{\displaystyle \frac{g_{\rho \pi \pi }}{4f_\pi ^2}}(\epsilon ,q_1+q_22q_4).`$ (42)
The amplitudes of the four pion decays of the $`\rho ^{}`$ are obtained from corresponding expressions for the $`\rho ^+`$ by reversing an overall sign. These considerably simplified expressions are especially convenient in the calculation of the $`\omega 5\pi `$ decay amplitude, because the typical invariant masses of the four pion system in the above decay are in the vicinity of 620 MeV (see Sec. IV for more detail).
## III Results for various $`\rho 4\pi `$ decays
When evaluating the partial widths of the $`4\pi `$ decays of $`\rho `$ meson the modulus squared of the matrix element is expressed via the Kumar variables . The idea of speeding up the numerical integration suggested in Ref. is realized in the numerical algorithm. The results of evaluation of the partial widths at $`\sqrt{s}=m_\rho =770`$ MeV are as follows: $`\mathrm{\Gamma }(\rho ^02\pi ^+2\pi ^{},m_\rho )=0.89`$ keV, $`\mathrm{\Gamma }(\rho ^0\pi ^+\pi ^{}2\pi ^0,m_\rho )=0.24`$ and 0.44 keV, respectively, without and with the induced anomaly term being taken into account. This coincides with the results obtained in Ref. . In the case of the charged $`\rho `$ meson decays it is obtained for the first time: $`\mathrm{\Gamma }(\rho ^+\pi ^+3\pi ^0,m_\rho )=0.41`$ keV, $`\mathrm{\Gamma }(\rho ^+2\pi ^+\pi ^{}\pi ^0,m_\rho )=0.71`$ and 0.90 keV respectively, without and with the anomaly induced term being taken into account. When obtaining these figures, the narrow $`\rho `$ width approximation is used. This is equivalent to the setting $`\mathrm{\Gamma }_{\rho \pi \pi }0`$ in Eq. (13). Keeping the physical value of the $`\rho `$ width gives the results deviating from those obtained in the narrow width approximation by a quantity that does not exceed a few percents of the values obtained in the latter. This is true in the case of the invariant mass of the four pion state lying below the $`\rho \pi `$ threshold energy, $`m_{4\pi }<910`$ MeV. Recall that the allowing for the finite widths effects is in fact equivalent to the loop correction being taken into account.
The above results are obtained at $`a=2`$. The variation of $`a`$ within $`20\%`$ around this value implies the variation of the branching ratios within $`20\%`$ around the values cited above. This fact can be easily traced in the nonrelativistic limit where the parameter $`a`$ enters the expressions for the amplitudes as an overall factor $`\sqrt{a}`$, see Eqs. (5) and (42).
### A The decay $`\rho ^04\pi `$ as manifested in $`e^+e^{}`$ annihilation.
The results of the evaluation of the $`4\pi `$ state production cross section in the reaction $`e^+e^{}\rho ^04\pi `$
$$\sigma _{e^+e^{}\rho 4\pi }(s)=\frac{12\pi m_\rho ^3\mathrm{\Gamma }_{\rho e^+e^{}}(m_\rho )\mathrm{\Gamma }_{\rho 4\pi }(E)}{E^3\left|D_\rho (s)\right|^2},$$
(43)
where $`s=E^2`$ is the square of the total center-of mass energy, and $`D_\rho (s)`$ is obtained from Eq. (13) upon the substitution $`q^2s`$, are plotted in Figs. 2 and 3. Note that the values of the vector meson parameters taken from Ref. are used hereafter. The following notations are such that
$$q(m_a,m_b,m_c)=\frac{1}{2m_a}\lambda ^{1/2}(m_a^2,m_b^2,m_c^2),$$
(44)
with the function $`\lambda `$ given by the equation
$$\lambda (x,y,z)=x^2+y^2+x^22(xy+xz+yz),$$
(45)
is the momentum of final particle $`b`$ (or $`c`$) in the rest frame system of decaying particle $`a`$.
To demonstrate the effects of chiral dynamics, also shown is the energy dependence of the cross section evaluated in the model of pure phase space for the four pion decay. In this model, the $`4\pi `$ partial width normalized to the width at the $`\rho `$ mass calculated in the dynamical model, is given by the expression
$$\mathrm{\Gamma }^{\mathrm{LIPS}}(\rho 4\pi ,s)=\mathrm{\Gamma }(\rho 4\pi ,m_\rho ^2)\frac{W_{4\pi }(s)}{W_{4\pi }(m_\rho ^2)},$$
(46)
where the four pion phase space volume is
$`W_{4\pi }(s)`$ $`=`$ $`{\displaystyle \frac{\pi ^3}{16(2\pi )^8s^{3/2}N_s}}{\displaystyle _{(3m_\pi )^2}^{(\sqrt{s}m_\pi )^2}}{\displaystyle \frac{ds_1}{s_1}}\lambda ^{1/2}(s,s_1,m_\pi ^2)`$ (49)
$`\times {\displaystyle _{(2m_\pi )^2}^{(\sqrt{s_1}m_\pi )^2}}{\displaystyle \frac{ds_2}{s_2}}\lambda ^{1/2}(s_1,s_2,m_\pi ^2)`$
$`\times \lambda ^{1/2}(s_2,m_\pi ^2,m_\pi ^2).`$
In the above formula, $`N_s=4\text{ (2)}`$ is the factor that takes into account the identity of final pions in the final state $`2\pi ^+2\pi ^{}`$($`\pi ^+\pi ^{}2\pi ^0`$), respectively. As the evaluation shows, the ratio
$$R(s)=\mathrm{\Gamma }(\rho 2\pi ^+2\pi ^{},s)/\mathrm{\Gamma }^{\mathrm{LIPS}}(\rho 2\pi ^+2\pi ^{},s)$$
changes from 0.4 at $`\sqrt{s}=650`$ MeV to 1 at $`\sqrt{s}=m_\rho `$. As can be observed from the figures, the rise of the $`\rho 4\pi `$ partial width with the energy increase is such fast that it compensates completely the falling of the $`\rho `$ meson propagator and electronic width. Also noticeable is the dynamical effect of the $`\omega \pi ^0`$ production threshold in the decay $`\rho ^0\pi ^+\pi ^{}2\pi ^0`$ at $`\sqrt{s}>850`$ MeV which results from the anomaly induced Lagrangian, see Fig. 3. To quantify the abovementioned effect of vanishing contribution of higher derivatives at the left shoulder of the $`\rho `$ resonance it should be noted that the difference between magnitude of $`\mathrm{\Gamma }(\rho \pi ^+\pi ^{}2\pi ^0,s)`$ with and without the term originating from the anomaly induced Lagrangian, equal to $`100\%`$ at $`\sqrt{s}=m_\rho `$, diminishes rapidly with the energy decrease amounting to $`8\%`$ at $`\sqrt{s}=700`$ MeV and $`0.3\%`$ at $`\sqrt{s}=650`$ MeV. It should be pointed out that the evaluation of the partial widths with the nonrelativistic expressions for the $`\rho 4\pi `$ amplitudes, Eq. (42), gives the values which deviate from those obtained with the exact expressions, by the quantity ranging from 7 to 15 $`\%`$, depending on the energy in the interval from 610 to 770 MeV.
As is seen from Fig. 2, the predictions of chiral symmetry for the $`e^+e^{}2\pi ^+2\pi ^{}`$ reaction cross section do not contradict to the three lowest experimental points of CMD-2 detector . However, at $`\sqrt{s}>800`$ MeV one can observe a substantial deviation between the predictions of the Lagrangian (3) and the data of CMD-2. In all appearance, this is due to the contribution of higher derivatives and chiral loops neglected in the present work. It is expected that the left shoulder of the $`\rho `$ is practically free of such contributions, and by this reason it is preferable for studying the dynamical effects of chiral symmetry. Note that even at $`\sqrt{s}=650`$ MeV, where the contribution of higher derivatives is negligible, one can hope to gather one event of the reaction $`e^+e^{}2\pi ^+2\pi ^{}`$ per day, and up to 10 events of this reaction per day at $`\sqrt{s}=700`$ MeV, provided the luminosity $`L=10^{32}\text{cm}^2\text{s}^1`$ is achieved, i.e. to have a factory for a comprehensive study of the chiral dynamics of many-pion systems.
Due to helicity conservation, $`\rho `$ meson is produced in the states with the spin projections $`\lambda =\pm 1`$ on the $`e^+e^{}`$ beam axes characterized by the unit vector $`𝐧_0`$. The latter is assumed to be directed along the z axes. Then, using the expressions for the total $`\rho 4\pi `$ amplitudes, one can obtain the angular distributions for the final pions. The are expected to be cumbersome. However, the good approximation for these distributions are obtained from the approximate nonrelativistic expression Eq. (42). The specific formulas are collected in the Appendix.
### B $`\rho 4\pi `$ in $`\tau `$ decays
Based on the vector current conservation, the partial width of the decay $`\tau ^{}\nu _\tau (4\pi )^{}`$ can be written as the integral over the invariant mass of the four pion state $`m`$, extended up to some mass $`m_0`$, whose maximal value is $`m_{0\mathrm{m}\mathrm{a}\mathrm{x}}=m_\tau `$:
$`B_{\tau ^{}\nu _\tau (4\pi )^{}}(m_0)`$ $`=`$ $`T_\tau {\displaystyle _{4m_\pi }^{m_0}}𝑑m{\displaystyle \frac{2m^2\mathrm{\Gamma }_{\tau ^{}\nu _\tau \rho ^{}}(m)}{\pi |D_\rho (m^2)|^2}}`$ (51)
$`\times \mathrm{\Gamma }_{\rho ^{}(4\pi )^{}}(m),`$
where $`T_\tau `$ and
$`\mathrm{\Gamma }_{\tau ^{}\nu _\tau \rho ^{}}(m)`$ $`=`$ $`{\displaystyle \frac{G_F^2\mathrm{cos}^2\theta _C}{8\pi f_\rho ^2}}m_\tau ^3m_\rho ^2\left(1{\displaystyle \frac{m^2}{m_\tau ^2}}\right)^2`$ (53)
$`\times \left(1+2{\displaystyle \frac{m^2}{m_\tau ^2}}\right)`$
are, respectively, the lifetime of $`\tau `$ lepton and its partial width of the decay $`\tau ^{}\nu _\tau \rho ^{}`$ , with $`m`$ being the invariant mass of the four pion state. Using the numerical values of the $`\rho 4\pi `$ decay widths, one can evaluate the branching ratios of the four pion $`\tau `$ decays for various values of the upper invariant mass $`m_0`$ of the latter. The results of the evaluation of the branching ratios of the decays $`\tau ^{}\nu _\tau 2\pi ^{}\pi ^+\pi ^0`$ and $`\tau ^{}\nu _\tau \pi ^{}3\pi ^0`$ for the values of the invariant mass of the four pion system from 600 to 850 MeV are plotted in Figs. 4 and 5, respectively. In particular, taking $`m_0=740`$ MeV one obtains
$$B(\tau ^{}\nu _\tau 2\pi ^{}\pi ^+\pi ^0,740\text{MeV})=\{\genfrac{}{}{0pt}{}{7.6\times 10^8(\text{ without anomaly induced term})}{8.4\times 10^8(\text{ with anomaly induced term})},$$
(54)
and
$$B(\tau ^{}\nu _\tau 3\pi ^0\pi ^{},740\text{MeV})=4.6\times 10^8.$$
(55)
Corresponding values for the upper integration mass $`m_0=640`$ MeV are
$$B(\tau ^{}\nu _\tau 2\pi ^{}\pi ^+\pi ^0,640\text{MeV})=\{\genfrac{}{}{0pt}{}{2.895\times 10^{10}(\text{ without anomaly induced term})}{2.900\times 10^{10}(\text{ with anomaly induced term})},$$
(56)
and
$$B(\tau ^{}\nu _\tau 3\pi ^0\pi ^{},640\text{MeV})=1.8\times 10^{10}.$$
(57)
The comparison of Eqs. (54) and (56), and both curves in Fig. 4 again demonstrates that the contributions of higher derivatives represented by the terms induced by the anomalous Lagrangian of Wess and Zumino vanish rapidly with the decreasing of mass. Unfortunately, the domains in the low four pion invariant mass where the effects of chiral dynamics are clean, are hardly accessible with $`\tau `$ factories. Indeed, guided by the expression for the cross section of the $`\tau `$ lepton pair production in $`e^+e^{}`$ annihilation
$$\sigma _{\tau ^+\tau ^{}}(s)=\frac{4\pi \alpha ^2}{3s}\sqrt{1\frac{4m_\tau ^2}{s}}\left(1+2\frac{m_\tau ^2}{s}\right),$$
(58)
one can find that up to $`N=25\times 10^7`$ $`\tau `$ lepton pairs with the total energy $`\sqrt{s}=m_{\psi (2S)}`$ can be produced per season at $`\tau `$-charm factory with expected luminosity $`L=10^{34}`$ cm<sup>-2</sup> s<sup>-1</sup> . This implies that one can detect only from 2 to 4 events per season in the four pion mass range below 700 MeV. Nevertheless, the event counting rate rises rapidly with the increase of the upper integration mass $`m_0`$ in Eq. (51), reaching, at $`m_0=m_\rho `$, the figure about 60 to 120 events per season, depending on the charge combination of the final pions.
### C The decay $`\rho 4\pi `$ in photoproduction, $`\pi N\rho \pi N`$, and so on.
To characterize the possibility of the study of the $`\rho 4\pi `$ decays in photoproduction, we calculate the quantity
$$B_{\rho 4\pi }^{\mathrm{aver}}(m_0)=\frac{2}{\pi }_{4m_\pi }^{m_0}𝑑m\frac{m^2\mathrm{\Gamma }_{\rho 4\pi }(m)}{|D_\rho (m^2)|^2},$$
(59)
which is the average of the branching ratio over the invariant mass of the four pion state. In the limit $`m_0\mathrm{}`$, Eq. (59) serves as the definition of the branching ratio in case of a wide resonance. Equation (59) should be confronted with the familiar definition of the branching ratio at the $`\rho `$ mass
$$B_{\rho 4\pi }(m_\rho )=\mathrm{\Gamma }_{\rho 4\pi }/\mathrm{\Gamma }_\rho ,$$
(60)
which results from Eq. (59) upon the replacement $`m\mathrm{\Gamma }_\rho /\pi |D_\rho (m^2)|^2\delta (m^2m_\rho ^2)`$ valid in the limit of narrow width. With the partial widths evaluated here one finds
$$B_{\rho ^02\pi ^+2\pi ^{}}(m_\rho )=5.9\times 10^6,$$
(61)
and
$$B_{\rho ^0\pi ^+\pi ^{}2\pi ^0}(m_\rho )=\{\genfrac{}{}{0pt}{}{1.6\times 10^6(\text{ without anomaly induced term})}{2.9\times 10^6(\text{ with anomaly induced term})}.$$
(62)
The results of plotting the quantity $`B_{\rho 4\pi }^{\mathrm{aver}}(m_0)`$ are shown in Fig. 6 and 7. In particular, the evaluation gives $`B_{\rho ^02\pi ^+2\pi ^{}}^{\mathrm{aver}}(m_0)=4.4\times 10^6\text{}6.1\times 10^8`$, and $`1.4\times 10^9`$ at $`m_0=850`$, 700, and 640 MeV, respectively. In the case of other four pion decay mode of the $`\rho ^0`$ the results are the following. In the model with the vanishing term induced by the anomalous Lagrangian of Wess and Zumino one obtains $`B_{\rho ^0\pi ^+\pi ^{}2\pi ^0}^{\mathrm{aver}}(m_0)=1.3\times 10^6\text{}1.58\times 10^8`$, and $`3.66\times 10^{10}`$ at $`m_0=850`$, 700, and 640 MeV. In the model that includes the above term, one obtains $`B_{\rho ^0\pi ^+\pi ^{}2\pi ^0}^{\mathrm{aver}}(m_0)=4.9\times 10^6\text{}1.65\times 10^8`$, and $`3.63\times 10^{10}`$ at the same respective values of $`m_0`$. As is expected, the branching ratios in the two mentioned models converge to each other in view of the rapid vanishing of the contributions due to the terms with higher derivatives. The difference between the two definitions of the branching ratio is seen upon comparison of $`B_{\rho 4\pi }^{\mathrm{aver}}(m_0=850\text{ MeV})`$ evaluated for various charge combinations of the final pions, with Eqs. (61) and (62).
With the total number of $`\rho `$ mesons $`N6\times 10^9`$ expected to be produced on nucleon at the Jefferson Laboratory ”photon factory” one may hope to observe about 100, 360 events of the $`\rho `$ decays into the states $`\pi ^+\pi ^{}2\pi ^0`$, $`2\pi ^+2\pi ^{}`$, respectively, in the mass range $`m_0<700`$ MeV where the effects of chiral dynamics are most clean. The photoproduction on heavy nuclei results in increasing the number of produced $`\rho `$ mesons faster than $`A^{2/3}`$, where $`A`$ is the atomic weight. A generally adopted behavior is in accord with the behavior $`A^{0.80.95}`$ . Thus the photoproduction of the four pion states on heavy nuclei would give the possibility of the high statistics study of the effects of chiral dynamics in the four pion decays of the $`\rho (770)`$. It should be recalled once more that the counting rate rises rapidly with the increase of $`m_0`$.
The conclusions about the angular distributions of the final pions with zero net charge in photoproduction are the following. Of course, their general expression should be deduced from the full decay amplitudes which can be found in Sec. II, together with the detailed form of the photoproduction mechanism. The qualitative picture, however, can be obtained upon assuming $`s`$-channel helicity conservation to be a good selection rule for the photoproduction reactions . Then in the helicity reference frame characterized as the frame where the $`\rho `$ is at rest, while its spin quantization axes is directed along the $`\rho `$ momentum in the center-of-mass system, the expressions for the angular distributions coincide with the corresponding expressions for the production of these states in $`e^+e^{}`$ annihilation that can be found in the Appendix. Since, at high energies, the direction of the final $`\rho `$ momentum lies at the scattering angle less than $`0.5^{}`$ in the case of the photoproduction on heavy nuclei, the vector $`𝐧_0`$ can be treated as pointed along the photon beam direction.
Note that another peripheral reactions can provide a sufficiently intense source of $`\rho `$ mesons. For example, the diffractive production of the $`\rho \pi `$ state in $`\pi N`$ collisions are currently under study with the VES detector in Protvino. The regions of the four pion invariant mass spectrum larger than $`m_\rho `$, namely, $`m_0`$ 850 MeV where $`B^{\mathrm{aver}}(\rho 4\pi ,m_0)10^5`$ should be included to measure the $`\rho 4\pi `$ branching ratio reliably. As is explained in Introduction, this would require the inclusion of the contributions of $`a_1`$ meson and higher derivatives to the total amplitude. Nevertheless, the results of the present paper shown in Figs. 8 and 9, obtained upon neglecting the latter contributions can be regarded as a guess in the experimental work in this direction.
## IV The decay $`\omega 5\pi `$
One may convince oneself that the $`\omega \rho \pi 5\pi `$ decay amplitude unambiguously results from the anomaly induced Lagrangian (8). This amplitude is represented by the diagrams shown in Fig. 10. As one can foresee, its general expression looks cumbersome. However, it can be considerably simplified upon noting that due to the low pion momentum, $`|𝐪_\pi |0.5m_\pi `$, the nonrelativistic expressions Eq. (42) for the $`\rho 4\pi `$ decay amplitudes in the diagrams Fig. 10(a) are valid with the accuracy $`5\%`$ in the $`4\pi `$ mass range relevant for the present purpose . This accuracy is estimated from the direct evaluation of the $`\rho 4\pi `$ branching ratios with the exact $`\rho 4\pi `$ decay formulas given in Sec. II, and with the approximate ones, Eq. (42). The evaluation shows that the results differ by approximately 10 $`\%`$ in the width. Likewise, the expression for the combination $`D_\pi ^1M(\pi 3\pi )`$ standing in the expression for the diagrams in Fig. 10(b) can be replaced, with the same accuracy, by $`(8m_\pi ^2)^1`$ times the nonrelativistic $`\pi 3\pi `$ amplitudes in Eq. (17). First, using Eq. (42) one obtains the expression for the sum of the diagrams shown in Fig. 10(a). Second, using Eq. (17), one obtains the expression for the sum of the diagrams shown in Fig. 10(b). Note that the contribution of the diagrams Fig. 10(b) was neglected in Ref. . The final expressions for the $`\omega 5\pi `$ decay amplitudes, upon neglecting the terms of the order of $`O(|𝐪_\pi ^4/m_\pi ^4)`$ or higher, can be represented in the form:
$`M(\omega 2\pi ^+2\pi ^{}\pi ^0)`$ $`=`$ $`{\displaystyle \frac{N_cg_{\rho \pi \pi }g^2}{8(2\pi )^2f_\pi ^3}}\epsilon _{\mu \nu \lambda \sigma }q_\mu ϵ_\nu \{(1+P_{12})q_{1\lambda }[{\displaystyle \frac{(q_2+3q_4)_\sigma }{D_\rho (qq_1)}}{\displaystyle \frac{2q_{4\sigma }}{D_\rho (q_1+q_4)}}]`$ (65)
$`(1+P_{35})q_{3\lambda }\left[{\displaystyle \frac{(q_5+3q_4)_\sigma }{D_\rho (qq_3)}}{\displaystyle \frac{2q_{4\sigma }}{D_\rho (q_3+q_4)}}\right]`$
$`(1+P_{12})(1+P_{35})q_{3\lambda }[{\displaystyle \frac{2q_{4_\sigma }}{D_\rho (qq_4)}}+{\displaystyle \frac{q_{1\sigma }}{D_\rho (q_1+q_3)}}]\},`$
with the final momentum assignment according to $`\pi ^+(q_1)\pi ^+(q_2)\pi ^{}(q_3)\pi ^{}(q_5)\pi ^0(q_4)`$, and
$`M(\omega \pi ^+\pi ^{}3\pi ^0)`$ $`=`$ $`{\displaystyle \frac{N_cg_{\rho \pi \pi }g^2}{8(2\pi )^2f_\pi ^3}}(1P_{12})(1+P_{34}+P_{35})\epsilon _{\mu \nu \lambda \sigma }q_\mu ϵ_\nu q_{1\lambda }`$ (68)
$`\times \{q_{3\sigma }[{\displaystyle \frac{1}{D_\rho (qq_3)}}{\displaystyle \frac{1}{D_\rho (q_1+q_3)}}]`$
$`q_{2\sigma }[{\displaystyle \frac{4}{3D_\rho (qq_1)}}{\displaystyle \frac{1}{2D_\rho (q_1+q_2)}}]\},`$
with the final momentum assignment according to $`\pi ^+(q_1)\pi ^{}(q_2)\pi ^0(q_3)\pi ^0(q_4)\pi ^0(q_5)`$. In both above formulas, $`ϵ_\nu `$, $`q_\mu `$ stand for four-vectors of polarization and momentum of $`\omega `$ meson. Note that the first term in each square bracket refers to the specific diagram shown in Fig. 10(a) while the second one does to the diagram shown in Fig. 10(b).
Yet even in this simplified form the expressions for the $`\omega 5\pi `$ amplitudes are not easy to use for evaluation of the branching ratios. To go further, one should note the following. One can check that the invariant mass of the $`4\pi `$ system on which the contribution of the diagrams shown in Fig. 10(a) depends, changes in the very narrow range $`558\text{ MeV}<m_{4\pi }<642\text{ MeV}`$. Hence, one can set it in all the $`\rho `$ propagators standing as the first terms in all square brackets in Eqs. (68) and (65), with the accuracy $`20\%`$ in width, to the ”equilibrium” value $`\overline{m_{4\pi }^2}^{1/2}=620`$ MeV evaluated for the pion energy $`E_\pi =m_\omega /5`$ which gives the dominant contribution. The same is true for the invariant mass of the pion pairs on which the $`\rho `$ propagators standing as the last terms in square brackets of the above expressions, depend. This invariant mass varies in the narrow range $`280\text{ MeV}<m_{2\pi }<360\text{ MeV}`$. With the same accuracy, one can set it to $`\overline{m_{2\pi }^2}^{1/2}=295`$ MeV in all relevant propagators. On the other hand, the amplitude of the process $`\omega \rho ^0\pi ^0(2\pi ^+2\pi ^{})\pi ^0`$ is
$`M[\omega \rho ^0\pi ^0(2\pi ^+2\pi ^{})\pi ^0]`$ $`=`$ $`4{\displaystyle \frac{N_cg_{\rho \pi \pi }g^2}{8(2\pi )^2f_\pi ^3}}`$ (71)
$`\times \epsilon _{\mu \nu \lambda \sigma }q_\mu ϵ_\nu (q_1+q_2)_\lambda `$
$`\times {\displaystyle \frac{q_{4\sigma }}{D_\rho (qq_4)}},`$
where the momentum assignment is the same as in Eq. (65). The other relevant amplitude corresponding to the first diagram in Fig. 10(b) is
$`M[\omega \rho ^0\pi ^0(\pi ^+\pi ^{})(\pi ^+\pi ^{}\pi ^0)]`$ $`=`$ $`{\displaystyle \frac{N_cg_{\rho \pi \pi }g^2}{8(2\pi )^2f_\pi ^3}}\epsilon _{\mu \nu \lambda \sigma }q_\mu ϵ_\nu `$ (73)
$`\times (1+P_{12})(1+P_{35}){\displaystyle \frac{q_{1\lambda }q_{3\sigma }}{D_\rho (q_1+q_3)}}`$
Taking into account the above consideration concerning the invariant masses, one can replace all the $`\rho `$ propagators standing as the first terms in square parentheses of Eq. (65) with $`1/D_\rho (qq_4)`$ one. In the same manner, all the $`\rho `$ propagators standing as the second terms in square parentheses can be replaced with $`1/D_\rho (q_1+q_3)`$. Then the comparison of Eqs. (65), (71), and (73) shows that
$`M(\omega 2\pi ^+2\pi ^{}\pi ^0)`$ $``$ $`{\displaystyle \frac{5}{2}}M[\omega \rho ^0\pi ^0(2\pi ^+2\pi ^{})\pi ^0]`$ (75)
$`\times \left[1{\displaystyle \frac{D_\rho (\overline{m_{4\pi }^2})}{2D_\rho (\overline{m_{2\pi }^2})}}\right],`$
where we replace the ratio $`D_\rho (qq_4)/D_\rho (q_1+q_3)`$ with the ratio $`D_\rho (\overline{m_{4\pi }^2})/D_\rho (\overline{m_{2\pi }^2})`$ evaluated at the ”equilibrium” point. The same treatment shows that
$`M(\omega \pi ^+\pi ^{}3\pi ^0)`$ $``$ $`{\displaystyle \frac{5}{2}}M[\omega \rho ^+\pi ^{}(\pi ^+3\pi ^0)\pi ^{}]`$ (77)
$`\times \left[1{\displaystyle \frac{D_\rho (\overline{m_{4\pi }^2})}{2D_\rho (\overline{m_{2\pi }^2})}}\right],`$
where
$`M[\omega \rho ^+\pi ^{}(\pi ^+3\pi ^0)\pi ^{}]`$ $`=`$ $`4{\displaystyle \frac{N_cg_{\rho \pi \pi }g^2}{8(2\pi )^2f_\pi ^3}}`$ (79)
$`\times {\displaystyle \frac{\epsilon _{\mu \nu \lambda \sigma }q_\mu ϵ_\nu q_{1\lambda }q_{2\sigma }}{D_\rho (qq_2)}},`$
and the final momenta assignment is the same as in Eq. (68). The numerical values of $`\overline{m_{4\pi }^2}^{1/2}`$ and $`\overline{m_{2\pi }^2}^{1/2}`$ found above are such that the correction factor in parentheses of Eqs. (75) and (77) amounts to $`20\%`$ in magnitude. In what follows, the above correction will be taken into account as an overall factor of 0.64 in front of the branching ratios of the decays $`\omega 5\pi `$. When making this estimate, the imaginary part of the $`\rho `$ propagators in square brackets of Eq. (75) and (77) is neglected. This assumption is valid with the accuracy better than $`1\%`$ in width.
The evaluation of the partial widths valid with the accuracy $`20\%`$ can be obtained upon using the expressions (71)- (79). Such an accuracy is estimated by noting that the numerical magnitude of the $`\rho `$ propagators in the expressions for the $`\omega 5\pi `$ decay amplitudes in Eqs. (65) and (68) evaluated by assuming the invariant mass of the four pion system to be determined by either the ”equilibrium” pion energy or by the center of allowed range of the variation of this mass, respectively, differs by the quantity not exceeding 10 $`\%`$ of the numerical value of the $`\rho `$ propagator. Defining the branching ratio at the $`\omega `$ mass as
$$B_{\omega 5\pi }=\mathrm{\Gamma }_{\omega 5\pi }/\mathrm{\Gamma }_\omega ,$$
(80)
one finds
$`B_{\omega 2\pi ^+2\pi ^{}\pi ^0}`$ $`=`$ $`\left|1{\displaystyle \frac{D_\rho (\overline{m_{4\pi }^2})}{2D_\rho (\overline{m_{2\pi }^2})}}\right|^2\left({\displaystyle \frac{5}{2}}\right)^2{\displaystyle \frac{2}{\pi \mathrm{\Gamma }_\omega }}`$ (84)
$`\times {\displaystyle _{4m_{\pi ^+}}^{m_\omega m_{\pi ^0}}}dm`$
$`\times {\displaystyle \frac{m^2\mathrm{\Gamma }_{\omega \rho ^0\pi ^0}(m)\mathrm{\Gamma }_{\rho 2\pi ^+2\pi ^{}}(m)}{|D_\rho (m^2)|^2}}`$
$`=1.1\times 10^9`$
where
$$\mathrm{\Gamma }_{\omega \rho ^0\pi ^0}(m)=g_{\omega \rho \pi }^2q^3(m_\omega ,m,m_{\pi ^0})/12\pi ,$$
$$g_{\omega \rho \pi }=\frac{N_cg^2}{8\pi ^2f_\pi }=14.3\text{ GeV}^1,$$
and the numerical data for the $`\rho 2\pi ^+2\pi ^{}`$ decay width obtained in Sec. III are used. Note also the $`a^1`$ dependence of the $`\omega 5\pi `$ width on the HLS parameter $`a`$. The branching ratio $`B_{\omega \pi ^+\pi ^{}3\pi ^0}`$ is obtained from Eq. (84) upon changing the lower integration limit to $`m_{\pi ^+}+3m_{\pi ^0}`$, the substitution $`m_{\pi ^0}m_{\pi ^+}`$ in the expression for the momentum $`q`$, and the substitution of the $`\rho ^+\pi ^+3\pi ^0`$ decay width numerically calculated in Sec. III, instead of the $`\rho 2\pi ^+2\pi ^{}`$ one. Note that the former is corrected for the mass difference of charged and neutral pions. Of course, the main correction of this sort comes from the phase space volume of the final 4$`\pi `$ state. One obtains
$`B_{\omega \pi ^+\pi ^{}3\pi ^0}`$ $`=`$ $`\left|1{\displaystyle \frac{D_\rho (\overline{m_{4\pi }^2})}{2D_\rho (\overline{m_{2\pi }^2})}}\right|^2\left({\displaystyle \frac{5}{2}}\right)^2{\displaystyle \frac{2}{\pi \mathrm{\Gamma }_\omega }}`$ (88)
$`\times {\displaystyle _{m_{\pi ^+}+3m_{\pi ^0}}^{m_\omega m_{\pi ^+}}}dm`$
$`\times {\displaystyle \frac{m^2\mathrm{\Gamma }_{\omega \rho ^+\pi ^{}}(m)\mathrm{\Gamma }_{\rho ^+\pi ^+3\pi ^0}(m)}{|D_\rho (m^2)|^2}}`$
$`=8.5\times 10^{10},`$
where
$$\mathrm{\Gamma }_{\omega \rho ^+\pi ^{}}(m)=g_{\omega \rho \pi }^2q^3(m_\omega ,m,m_{\pi ^+})/12\pi .$$
As is pointed out in Ref. , the inclusion of the direct $`\omega \pi ^+\pi ^{}\pi ^0`$ vertex reduces the 3$`\pi `$ decay width of the $`\omega `$ by $`33\%`$. This implies that one should make the following replacement to take into account the effect of the pointlike diagrams in Fig. 10(b) in the expression for the suppression factor:
$`\left|1{\displaystyle \frac{D_\rho (\overline{m_{4\pi }^2})}{2D_\rho (\overline{m_{2\pi }^2})}}\right|^2`$ $``$ $`|1{\displaystyle \frac{D_\rho (\overline{m_{4\pi }^2})}{2}}[{\displaystyle \frac{1}{D_\rho (\overline{m_{2\pi }^2})}}`$ (91)
$`{\displaystyle \frac{1}{3m_\rho ^2}}\left]\right|^2|1{\displaystyle \frac{D_\rho (\overline{m_{4\pi }^2})}{3D_\rho (\overline{m_{2\pi }^2})}}|^2`$
$`0.75,`$
instead of 0.64, which results in the increase of the above branching ratios by the factor of 1.17.
The numerical value of the $`\omega 5\pi `$ decay width changes by the factor of two when varying the energy within $`\pm \mathrm{\Gamma }_\omega /2`$ around the $`\omega `$ mass. In other words, the dependence of this partial width on energy is very strong. This is illustrated by Fig. 11 where the $`\omega 5\pi `$ excitation curves in $`e^+e^{}`$ annihilation,
$`\sigma _{e^+e^{}\omega 5\pi }(s)`$ $`=`$ $`12\pi \left({\displaystyle \frac{m_\omega }{E}}\right)^3\mathrm{\Gamma }_{\omega e^+e^{}}(m_\omega )`$ (93)
$`\times {\displaystyle \frac{\mathrm{\Gamma }_\omega B_{\omega 5\pi }(E)}{\left[(sm_\omega ^2)^2+(m_\omega \mathrm{\Gamma }_\omega )^2\right]}},`$
are plotted. Here $`B_{\omega 2\pi ^+2\pi ^{}\pi ^0}(E)`$ \[$`B_{\omega \pi ^+\pi ^{}3\pi ^0}(E)`$\] is given by Eq. (84) \[(88)\], respectively, with the substitution $`m_\omega E`$. The mentioned strong energy dependence of the partial width results in the asymmetric shape of the $`\omega `$ resonance and the shift of its peak position by +0.7 MeV. As is seen from Fig. 11, the peak value of the $`5\pi `$ state production cross section is about 1.5-2.0 femtobarns. Yet the decays $`\omega 5\pi `$ can be observable on $`e^+e^{}`$ colliders. Indeed, with the luminosity $`L=10^{33}\text{cm}^2\text{s}^1`$ near the $`\omega `$ peak, which seems to be feasible, one may expect about 2 events per week for the considered decays to be detected at these colliders.
The angular distributions of the final pions should be deduced from the full amplitudes Eqs. (65) and (68). However, some qualitative conclusions about the angular distributions can be drawn from the simplified expressions Eqs. (71), (75), (77), (79). Since helicity is conserved, only the states of the $`\omega (782)`$ with the spin projections $`\lambda =\pm 1`$ on the $`e^+e^{}`$ beam axes are populated. Corresponding expressions for the various combination of the final pions can be found in Appendix.
The strong energy dependence of the five pion partial width of the $`\omega `$ implies that the branching ratio at the $`\omega `$ mass, Eq. (80), evaluated above, is slightly different from that determined by the expression
$$B_{\omega 5\pi }^{\mathrm{aver}}(E_1,E_2)=\frac{2}{\pi }_{E_1}^{E_2}𝑑E\frac{E^2\mathrm{\Gamma }_\omega B_{\omega 5\pi }(E)}{(E^2m_\omega ^2)^2+(m_\omega \mathrm{\Gamma }_\omega )^2}.$$
(94)
Taking $`E_1=`$ 772 MeV and $`E_2=`$ 792 MeV, one finds $`B_{\omega 2\pi ^+2\pi ^{}\pi ^0}^{\mathrm{aver}}(E_1,E_2)=9.0\times 10^{10}`$ and $`B_{\omega \pi ^+\pi ^{}3\pi ^0}^{\mathrm{aver}}(E_1,E_2)=6.7\times 10^{10}`$ to be compared to Eqs. (84) and (88), respectively. In particular, the quantity $`B_{\omega 2\pi ^+2\pi ^{}\pi ^0}^{\mathrm{aver}}(E_1,E_2)`$ is the relevant characteristics of this specific decay mode in photoproduction experiments. The Jefferson Lab ”photon factory” could also be suitable for detecting the five pion decays of the $`\omega `$. However, in view of the suppression of the $`\omega `$ photoproduction cross section by the factor of 1/9 as compared with the $`\rho `$ one, the total number of $`\omega `$ mesons will amount to $`7\times 10^8`$ per nucleon. Hence, the increase of intensity of this machine by the factor of 50 is highly desirable, in order to observe the decay $`\omega 5\pi `$ and measure its branching ratio. Evidently, the $`\omega `$ photoproduction on heavy nuclei is preferable in view of the dependence of the cross section on atomic weight $`A`$ growing as $`A^{0.80.95}`$ .
The conclusions about how the angular distributions in the $`\omega `$ photoproduction are related with those in $`e^+e^{}`$ annihilation, are basically similar to those concerning the $`\rho `$ photoproduction discussed in Sec. III C. The expressions for these distributions coincide with the corresponding expressions Eqs. (115), (120), (122), and (126) obtained for the case of $`e^+e^{}`$ annihilation. The vector $`𝐧_0`$ in those formulas can be treated as pointed along the photon beam direction.
## V Conclusion
The results presented in this paper show that, in the $`\rho 4\pi `$ decay amplitude, the contributions of the higher derivatives specified by the term induced by the anomalous Lagrangian of Wess and Zumino, vanish rapidly when decreasing the invariant mass of the four pion system below 700 MeV. The loop corrections are also expected to behave similarly. The decay $`\omega 5\pi `$ is of a special interest, because its kinematics is such that the final pions are essentially nonrelativistic, and the above effects are completely suppressed in its decay amplitude. Hence the approach to the decays of the vector mesons $`\rho `$ and $`\omega `$ presented in this paper is the zero order approximation to the full amplitude, in a close analogy with the Weinberg amplitude in the classical $`\pi \pi `$ scattering. Under the approximation of the present paper, all chiral models of the vector meson interactions with pions are indistinguishable in their predictions concerning the many pion branching ratios. Any difference could manifest upon including the higher derivatives and chiral loops. As the next step in the development of the present study, the inclusion of the $`a_1`$ meson contribution at the tree level approximation would be important. This could help in extending the validity of the treatment up to the invariant masses just below 910 MeV. The task looks meaningful, since the loop corrections, whose particular manifestation is the finite width effects uncovered to be unimportant at $`m_{4\pi }<910`$ MeV ( see Sec. III), are still expected to be small at these invariant masses.
In our opinion, the above considerations show that the left shoulder of the $`\rho `$ peak is the very perspective place to study the effects of chiral dynamics of the vector meson interactions. The $`e^+e^{}`$ colliders with the large enough luminosity at energies below the $`\rho `$ mass could provide the controlled source of soft pions. The role of higher derivatives, loop corrections, and, possibly, the $`\rho ^{}`$, $`\rho ^{\prime \prime }`$ contributions in the low energy effective Lagrangian for the soft pions, as well as various schemes of incorporation of the vector mesons into the chiral approach, can be successfully tested with such machines. The intense beams of photons from the Jefferson Laboratory ”photon factory” are also of great importance in achieving the mentioned theoretical goals. Certainly, the measurements of the branching ratio of the five pion decays of the $`\omega `$ at the level of $`B_{\omega 5\pi }10^9`$ constitute the real challenge to experimenters, but by the reasons specified above the effects of chiral dynamics of the vector meson interactions are manifested in the decay $`\omega 5\pi `$ in the very clean way.
###### Acknowledgements.
We are grateful to G. N. Shestakov and A. M. Zaitsev for discussion. The present work is supported in part by grant No. RFBR-INTAS IR-97-232.
## Angular distributions
Here a number of expressions for the angular distributions of various combinations of the final pions in the decays $`\rho 4\pi `$ and $`\omega 5\pi `$ are given. In what follows, the one photon $`e^+e^{}`$ annihilation production mechanism for these states is assumed, where the $`e^+e^{}`$ beam axes is characterized by the unit vector $`𝐧_0`$ directed along the z axes.
### 1 The angular distributions in the $`\rho 4\pi `$ decay
Taking $`\theta _i,\varphi _i`$ to be the polar and azimuthal angles of the pion three momentum $`𝐪_i`$, where the momentum assignment corresponds to Eq. (42), one finds the following.
(i) The $`\rho ^02\pi ^+2\pi ^{}`$ decay. The probability density of the emission of four charged pions can be found directly from the first Eq. (42):
$`w`$ $``$ $`(𝐪_1+𝐪_2𝐪_3𝐪_4)^2[𝐧_0(𝐪_1+𝐪_2𝐪_3𝐪_4)]^2`$ (99)
$`={\displaystyle \underset{i=1}{\overset{4}{}}}𝐪_i^2\mathrm{sin}^2\theta _i+2|𝐪_1|(1P_{23}P_{24})|𝐪_2|\mathrm{sin}\theta _1\mathrm{sin}\theta _2\mathrm{cos}(\varphi _1\varphi _2)`$
$`2|𝐪_2|(1+P_{34})|𝐪_3|\mathrm{sin}\theta _2\mathrm{sin}\theta _3\mathrm{cos}(\varphi _2\varphi _3)`$
$`+2|𝐪_3||𝐪_4|\mathrm{sin}\theta _3\mathrm{sin}\theta _4\mathrm{cos}(\varphi _3\varphi _4).`$
One may use the relation
$$(\epsilon ,q_1+q_2+q_3+q_4)=0$$
(100)
that expresses the transverse character of the $`\rho `$ polarization four vector $`\epsilon `$, to get rid of the momenta of negatively charged pions $`q_3`$ and $`q_4`$. Then the probability density of the emission of two $`\pi ^+`$’s found from the first Eq. (42) is
$`w`$ $``$ $`(𝐪_1+𝐪_2)^2[𝐧_0(𝐪_1+𝐪_2)]^2`$ (103)
$`=𝐪_1^2\mathrm{sin}^2\theta _1+𝐪_2^2\mathrm{sin}^2\theta _2+2|𝐪_1||𝐪_2|\mathrm{sin}\theta _1\mathrm{sin}\theta _2`$
$`\times \mathrm{cos}(\varphi _1\varphi _2).`$
Allowing for Eq. (100), the angular distribution for the emission of two $`\pi ^{}`$’s is obtained from Eq. (103) upon the replacement $`𝐪_{1,2}𝐪_{3,4}`$.
(ii) The $`\rho ^0\pi ^+\pi ^{}2\pi ^0`$ decay. The probability density of the emission of $`\pi ^+\pi ^{}`$ pair found from the second Eq. (42) in the form
$`w`$ $``$ $`(𝐪_1𝐪_2)^2[𝐧_0(𝐪_1𝐪_2)]^2`$ (106)
$`=𝐪_1^2\mathrm{sin}^2\theta _1+𝐪_2^2\mathrm{sin}^2\theta _22|𝐪_1||𝐪_2|\mathrm{sin}\theta _1\mathrm{sin}\theta _2`$
$`\times \mathrm{cos}(\varphi _1\varphi _2).`$
Getting rid of the momentum $`q_2`$ one finds the corresponding expression for the final state $`\pi ^+2\pi ^0`$:
$`w`$ $``$ $`(2𝐪_1𝐪_3𝐪_4)^2[𝐧_0(2𝐪_1𝐪_3𝐪_4)]^2`$ (109)
$`=4𝐪_1^2\mathrm{sin}^2\theta _1+𝐪_3^2\mathrm{sin}^2\theta _3+𝐪_4^2\mathrm{sin}^2\theta _4(1+P_{34})4|𝐪_1||𝐪_3|\mathrm{sin}\theta _1\mathrm{sin}\theta _3\mathrm{cos}(\varphi _1\varphi _3)`$
$`+2|𝐪_3||𝐪_4|\mathrm{sin}\theta _3\mathrm{sin}\theta _4\mathrm{cos}(\varphi _3\varphi _4),`$
where $`P_{ij}`$ interchanges the pion momenta $`q_i`$ and $`q_j`$. In view of Eq. (100), the angular distribution for the state $`\pi ^{}2\pi ^0`$ is obtained from the above upon the replacement $`𝐪_1𝐪_2`$ and changing the signs in front of the terms containing $`(1+P_{34})`$.
### 2 The angular distributions in the $`\omega 5\pi `$ decay
In what follows the suitable notation for the vector product of the pion momenta are used:
$`[𝐪_i\times 𝐪_j]`$ $`=`$ $`|𝐪_i||𝐪_j|\mathrm{sin}\theta _{ij}`$ (111)
$`\times (\mathrm{sin}\mathrm{\Theta }_{ij}\mathrm{cos}\mathrm{\Phi }_{ij},\mathrm{sin}\mathrm{\Theta }_{ij}\mathrm{sin}\mathrm{\Phi }_{ij},\mathrm{cos}\mathrm{\Theta }_{ij}).`$
In other words, $`\theta _{ij}`$ is the angle between the pion momenta $`𝐪_i`$ and $`𝐪_j`$, $`\mathrm{\Theta }_{ij}`$, $`\mathrm{\Phi }_{ij}`$ being the polar and azimuthal angles of the normal to the plane spanned by the momenta $`𝐪_i`$ and $`𝐪_j`$. Choosing $`𝐧_0`$ to be the unit vector along z axes, the probability density of the emission of two $`\pi ^+`$’s with the momenta $`𝐪_1`$, $`𝐪_2`$, and $`\pi ^0`$ with the momentum $`𝐪_4`$ is represented as
$`w`$ $``$ $`\left[𝐪_4\times \left(𝐪_1+𝐪_2\right)\right]^2\left(𝐧_0\left[𝐪_4\times \left(𝐪_1+𝐪_2\right)\right]\right)^2`$ (115)
$`=𝐪_4^2[𝐪_1^2\mathrm{sin}^2\theta _{41}\mathrm{sin}^2\mathrm{\Theta }_{41}+𝐪_2^2\mathrm{sin}^2\theta _{42}\mathrm{sin}^2\mathrm{\Theta }_{42}`$
$`+2|𝐪_1||𝐪_2|\mathrm{sin}\mathrm{\Theta }_{41}\mathrm{sin}\mathrm{\Theta }_{42}\mathrm{sin}\theta _{41}\mathrm{sin}\theta _{42}`$
$`\times \mathrm{cos}(\mathrm{\Phi }_{41}\mathrm{\Phi }_{42})]`$
in the case of the final state $`2\pi ^+2\pi ^{}\pi ^0`$. Here the momentum assignment is the same as in Eq. (65). The angular distribution of two $`\pi ^{}`$’s with the momenta $`𝐪_3`$, $`𝐪_5`$, and $`\pi ^0`$ is obtained from Eq. (115) upon the replacement $`𝐪_{1,2}𝐪_{3,5}`$, because the identity
$$\epsilon _{\mu \nu \lambda \sigma }q_\mu ϵ_\nu (q_1+q_2)_\lambda q_{4\sigma }=\epsilon _{\mu \nu \lambda \sigma }q_\mu ϵ_\nu (q_3+q_5)_\lambda q_{4\sigma }$$
is valid. Since another identity
$$\epsilon _{\mu \nu \lambda \sigma }q_\mu ϵ_\nu (q_1+q_2)_\lambda q_{4\sigma }=\epsilon _{\mu \nu \lambda \sigma }q_\mu ϵ_\nu (q_1+q_2)_\lambda (q_3+q_5)_\sigma $$
is valid, one can write the angular distribution that includes four charged pions:
$`w`$ $``$ $`[(𝐪_1+𝐪_2)\times (𝐪_3+𝐪_5)]^2\left(𝐧_0[(𝐪_1+𝐪_2)\times (𝐪_3+𝐪_5)]\right)^2`$ (120)
$`=(1+P_{12})(1+P_{35})𝐪_1^2𝐪_3^2\mathrm{sin}^2\theta _{13}\mathrm{sin}^2\mathrm{\Theta }_{13}`$
$`+2|𝐪_1||𝐪_2|(1+P_{35})𝐪_3^2\mathrm{sin}\theta _{13}\mathrm{sin}\theta _{23}\mathrm{sin}\mathrm{\Theta }_{13}\mathrm{sin}\mathrm{\Theta }_{23}\mathrm{cos}(\mathrm{\Phi }_{13}\mathrm{\Phi }_{23})`$
$`+2|𝐪_3||𝐪_5|(1+P_{12})𝐪_1^2\mathrm{sin}\theta _{13}\mathrm{sin}\theta _{15}\mathrm{sin}\mathrm{\Theta }_{13}\mathrm{sin}\mathrm{\Theta }_{15}\mathrm{cos}(\mathrm{\Phi }_{13}\mathrm{\Phi }_{15})`$
$`+2|𝐪_1||𝐪_2||𝐪_3||𝐪_5|(1+P_{35})\mathrm{sin}\theta _{13}\mathrm{sin}\theta _{25}\mathrm{sin}\mathrm{\Theta }_{13}\mathrm{sin}\mathrm{\Theta }_{25}\mathrm{cos}(\mathrm{\Phi }_{13}\mathrm{\Phi }_{25}).`$
Here $`P_{ij}`$ interchanges the indices $`i`$ and $`j`$. In the case of the final state $`\pi ^+\pi ^{}3\pi ^0`$ the corresponding probability density can be obtained from Eqs. (77) and (79) and looks as
$`w`$ $``$ $`[𝐪_1\times 𝐪_2]^2(𝐧_0[𝐪_1\times 𝐪_2])^2`$ (122)
$`=𝐪_1^2𝐪_2^2\mathrm{sin}^2\theta _{21}\mathrm{sin}^2\mathrm{\Theta }_{21}.`$
Here the momentum assignment is the same as in Eq. (68). The corresponding angular distribution of one charged, say $`\pi ^+`$, and three neutral pions can be obtained from Eqs. (77) and (79) upon using the identity
$$\epsilon _{\mu \nu \lambda \sigma }q_\mu ϵ_\nu q_{1\lambda }q_{2\sigma }=\epsilon _{\mu \nu \lambda \sigma }q_\mu ϵ_\nu q_{1\lambda }(q_3+q_4+q_5)_\sigma $$
and looks as
$`w`$ $``$ $`\left[𝐪_1\times {\displaystyle \underset{i}{}}𝐪_i\right]^2\left(𝐧_0\left[𝐪_1\times {\displaystyle \underset{i}{}}𝐪_i\right]\right)^2`$ (126)
$`=𝐪_1^2[{\displaystyle \underset{i}{}}𝐪_i^2\mathrm{sin}^2\theta _{i1}\mathrm{sin}^2\mathrm{\Theta }_{i1}`$
$`+2{\displaystyle \underset{ij}{}}|𝐪_i||𝐪_j|\mathrm{sin}\theta _{i1}\mathrm{sin}\theta _{j1}\mathrm{sin}\mathrm{\Theta }_{i1}\mathrm{sin}\mathrm{\Theta }_{j1}`$
$`\times \mathrm{cos}(\mathrm{\Phi }_{i1}\mathrm{\Phi }_{j1})].`$
Here indices $`i,j`$ run over 3,4,5.
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# Glass transition of a particle in a random potential, front selection in non linear RG and entropic phenomena in Liouville and SinhGordon models
## I introduction
Despite significant progress in the last two decades, disordered systems continue to pose considerable theoretical challenges. Two important questions still largely open, are, respectively, to which extent the (better understood) mean field models are relevant to describe low dimensional physical systems, and, in the special case of two dimension, to which extent the powerful field theoretic treatments developed for pure models can be adapted to treat disordered models.
A celebrated controversy is whether the structure found in the solution of mean field models for spin glasses and other complex disordered systems, both in the statics and in the dynamics , has any counterpart in the world of experimentally relevant low dimensional models. Specifically it has been vigorously questionned whether the breaking of the phase space in “many pure states”, predicted to occur in mean field, may also occur in short range models, and how it can be properly defined . The unusual nature of the technique used to solve the statics, i.e the replica method with a hierarchical breaking of the permutation symmetry between replicas in the limit $`n0`$ (RSB), did not contribute to make the physics transparent. A distinct structure, which remarkably parallels the one in the statics, has been found to occur in the nonequilibrium dynamics. The dynamical problem can be studied by a priori better defined methods and leads to predictions which are in principle directly testable in experiments, such as a non trivial generalization of the fluctuation dissipation relations. Even so, it has been emphasized that mean field models, which usually involve infinite range or infinite number of component limits, may not capture physical processes important in low dimensions. The alternative “droplet picture” in its simplest form postulates the existence of a single ground state with excitations (droplets) of (free) energy $`\mathrm{\Delta }E`$ scaling with their size $`x`$ as $`\mathrm{\Delta }Ex^\theta `$, $`\theta >0`$. It provides a more conventional scaling description of the glass physics, as being controlled by zero temperature RG fixed points where temperature is formally irrelevant (with eigenvalue $`\theta `$).
Another important advance was the exact solution of simpler prototype models, such as the random energy model (REM) , where one consider simply a collection of independently distributed energy levels, as well as its generalization, the GREM , or the Directed Polymer on the Cayley Tree (DPCT) with disorder . These solutions being direct with no use of replica, their results can be fully relied upon. They exhibit a similar physics though, with a glass transition and in the glass phase, an exponential tail for the distribution of the free energy $`P(f)e^{\beta _cf}`$ for negative $`f`$. This feature is crucial to recover the same physics, and indeed many observables were found to be similar . In fact the alternative solution of the REM using replicas, given in or the one of the DPCT do involve RSB. In the REM model the structure of the glass phase is particularly transparent as being dominated by a few states .
It is important to go beyond models defined in mean field or on hierarchical (or ultrametric) structures and to study simple yet non trivial (and non artificial) finite $`d`$ models with full statistical translational invariance. In this paper we study the model of a particle in a gaussian random potential $`V(𝐫)`$ with spatial correlations which are invariant by translation and which grow as the logarithm of the distance. We consider this model in any dimension $`d`$, but in $`d=2`$ it has also been studied recently since it is of direct relevance for several physical systems . One example is a spin model with $`XY`$ symmetry and random gauge quenched disorder, which arises naturally in describing Josephson junction arrays or 2D crystalline structures with smooth disorder, e.g. flux lattices in superconductors , or electrons at the surface of Helium . In this model, a single topological defect (a $`XY`$ vortex), or a single neutral pair, sees precisely a random potential with logarithmic correlations . Another example arises in a model of localization of Dirac fermions in a random magnetic field, motivated by quantum Hall physics. There, the zero energy $`E=0`$ normalized wavefunction is identical to the Boltzmann weight of the particle studied here . This wave function is “critical” in a sense discussed below.
Here we study this model using a renormalization group (RG) approach, bounds, numerical methods and qualitative arguments. We show that it exhibits a transition at $`T=T_c>0`$ in any $`d1`$. We find that in the high temperature phase the particle is essentially delocalized over the whole system, while in the low temperature glass phase the Gibbs measure is concentrated in a few minima. The fact that such a simple (finite $`d`$) model exhibits a genuine glass transition is already surprising. Indeed, as we argue, this transition exists only for such a “marginal” type of correlations (which correspond to $`\theta =0`$ in the glass scaling mentionned above ). It disappears (for gaussian $`V(𝐫)`$) if correlations grow faster (with only a low temperature phase and single ground state dominance) or slower (with only a high temperature phase). Logarithmic growth of correlations thus produces exactly the right balance between the depth of the energy wells and their number (entropy). Note that for slower growing correlations one can recover a transition but only by artificially rescaling the disorder variance with the size of the system: in the extreme case of uncorrelated variables, it is the REM model. Here by contrast, there is a genuine phase transition in the thermodynamic limit, with no need for rescaling. Most interestingly, the glass phase is non trivial. The Gibbs measure is concentrated in a few distant minima which remain in a finite number in the thermodynamic limit. This is because the extrema of random variables with such correlations exhibit an interesting property of “return near the minimum”: there is, with a finite probability in a sample of size $`L+\mathrm{}`$, at least one second minimum far away (at distances of order $`L`$), and with a finite energy difference with the absolute minimum. And there are not too many (a thermodynamic number) of these secondary minima, leading to a zero entropy. As in the REM, this property leads here naturally to a non trivial ground state structure, reminiscent as we discuss of a genuine property of replica symmetry breaking in the replicated theory. The low temperature limit corresponds to a non trivial problem of extremal statistics of correlated variables, studied here.
Another interesting property of this model is its relation to the Liouville model (LM) and the sinhGordon model (SGM) in $`d=2`$ (and their boundary restriction in $`d=1`$): $`V(𝐫)`$ turns out to be the Liouville field while the LM and SGM partition functions arise simply as generating functions of the probability distribution of the partition sum $`Z[V]=_𝐫e^{\beta V(𝐫)}`$ of a single particle. The high temperature phase for the particle corresponds to the weak coupling regime for the LM and SGM, where we find that known exact results compare well with results for the particle. In the SGM we predict here the existence of a transition (more appropriately, a “change of behaviour”). It corresponds to the glass transition for the particule, which also exhibits interesting similarities with the weak to strong coupling transition in the Liouville theory (and the so called $`c=1`$ barrier). The glassy freezing of the particule is associated, in the LM and SGM, to new local operators and non smooth configurations being generated under RG.
To study the model we will introduce a RG approach based on a Coulomb gas renormalization a la Kosterlitz. It leads to a non linear RG equation (of the so-called Kolmogorov KPP type) for the full probability distribution of the “local disorder”. Indeed, a separation between the long range part of the disorder and the local, short range part, arises naturally in our approach. The RG equation has traveling wave types of solutions. The corresponding well known problem, in such non linear equations, of the selection of the velocity of the traveling front, and its freezing for $`TT_c`$ is related to glassy freezing of the particle free energy and, in the LM or SGM, to the “selection” of the anomalous dimensions (and at the transition dimension degeneracy leads to logarithmic operators). When temperature is lowered the local disorder becomes broadly distributed and the freezing occurs when its tails become relevant. Our RG method indicates that the physics depends only weakly of $`d`$. We will take advantage of this fact and check our results using simulations in $`d=1`$.
It is important to compare the present work to previous studies of the model. The existence of a freezing transition in $`d=2`$ has been conjectured previously based on an approximation which completely neglect spatial correlations (REM approximation, see below). Stronger arguments were given in , but did not fully establish the existence of a transition, which is done here (see Appendix (A)). The present work is thus motivated by the need to go beyond the REM approximation to describe this problem. In particular one wants to know what is the precise universality class of the model, which we hope can be determined from the RG method introduced here. This RG method yields some remarkable universal features of the probability distribution of the free energy and of its finite size corrections, different from the REM. It shows that the problem is more closely related to the directed polymer on a Cayley tree. A qualitative analogy between the present model and the DPCT was in fact cleverly guessed recently in Ref. . It is based on the observation that the energy of polymer configuration on a tree also scale logarithmically with the overlap distance defined on the tree (see Fig 1). It is remarkable that this connection naturally emerges here from the Kosterlitz type RG performed on this problem, via the KPP equation. It is all the more surprising, since the model studied here has statistical translational invariance, while a tree has a hierarchical structure. The solution of Derrida and Spohn (and the mapping onto the DPCT proposed in ) would be exact for random variables $`V(𝐫)`$ correlated with a hierarchical (i.e ultrametric) matrix of correlation. Here instead the correlations are translationally invariant and it is thus important to understand the origin of the analogy with the DPCT and to which extent it holds. The RG procedure developed in this paper is a first attempt to address these questions. The result is that we can make the mapping precise: at least for the universal observables studied here (e.g. the tails of the free energy distribution), the mapping is onto a continuum branching process, i.e a continuum limit of a Cayley tree (whereas could not be so specific).
The present model has also been studied in the context of random Dirac problems and localization. An early study of the $`E=0`$ wavefunction established that it was critical (in the sense of corresponding to a “delocalized” wave function, while $`E0`$ has finite localization length). However this study missed the glass transition. Later studies computed the multifractal spectrum based on the REM approximation and noticed the existence of a strong disorder regime. These and other studies however focused on properties of the high temperature phase: it was conjectured that the (conformal) Liouville field theory (LFT) (i.e a continuum limit of the LM) was able to describe all spatial correlations of the model in the high temperature phase. These works call for more investigations. First the glass transition and the peculiar physical properties of the low T (i.e strong disorder) phase have not been addressed, even at the most qualitative level. We thus find it useful to present the problem from a different perspective by comparing with other types of correlated disorder, or by recasting it as a problem of extremal statistics. Although well known properties of extremal statistics of uncorrelated variables were often used to study model disordered systems (see e.g. ) a lot remains to be understood about the (more realistic) case of correlated variables. Second, the question of the universality class is in our opinion far from established. Evidence for the LFT description mostly comes from reproducing the multifractal spectrum as given by the REM approximation and one would like to check it against more detailed predictions. The present RG procedure is a step towards clarifying the connexion between this model and solvable models such as Derrida’s REM and Derrida-Spohn DPCT. In this respect finite size corrections are important to understand, as they are found to exhibit universal prefactors allowing to distinguish between various universality classes. In addition, they determine the anomalous dimensions, and thus control the critical behaviour, in the full disordered XY model as shown in . Since they are found to be very large, they are also crucial in order to analyze the results of numerical simulations. In particular, although we confirm the result of , we also conclude that the sizes used in the numerical study of were in fact vastly insufficient for drawing firm conclusions: we do perform here a more detailed finite size analysis on much larger samples to confirm analytical predictions.
The model studied here is thus related to a surprising number of interesting problems. Let us mention for completeness that it also has connexions to problems such as two dimensional interfaces, or films, confined between two walls (for $`\beta =+\mathrm{}`$ it is the confinement entropy of a film), wetting transitions , extremal statistics of correlated variables useful e.g. for problems of “persistence” in nonequilibrium dynamics , and finally, to the clumping transition of a self gravitating planar gas . We will not explore these connections in details here.
This paper is organized as follows: in section II, the single particle model is defined and in section II B the random energies approximation (REM) is applied, which amounts to neglect the spatial correlations of the random potential. The full problem, with correlations taken into account, is related to the description of extremal statistics in II C, and three different classes of correlations are identified in II D from qualitative arguments. A new renormalization (RG) technique is applied to this problem in section III. The resulting non-linear scaling equation for the distribution of the local disorder is studied in section III C, and found to be related to the Kolmogorov KPP equation, which admits front solutions. This connection between front solutions of non-linear equations and the renormalization of disordered models is pursued in section III D, where a solution to the REM is found via a similar non-linear RG (details in Appendix C). The non trivial nature of the glass phase is discussed in III E together with its relations to replica symetry breaking. In part IV we present a numerical analysis of the problem of the particle in a random potential in $`d=1`$. Section V is devoted to the connexion between the particle model -and its transition- and entropic phenomena in the Liouville and Sinh-Gordon models. A direct RG analysis in Section V C allows to recover the corresponding change of behaviours in these models. Section VI contains the applications to the properties of the critical wave function of a Dirac fermion in a random magnetic field, in particular the multifractal properties and the property of quasi-localization. Appendix A contains an outline of the proof of the existence of a transition, Appendix B is a review of well known (and not so well known) results about extremal statistics, Appendix D contains an extended model which exhibits three phases.
## II Model and qualitative analysis
In this Section we define the model of a single particle in a correlated random potential. Then we describe the random energy model (REM) approximation used in previous studies which consist in neglecting correlations. We then pose the new questions which we want to address here for the true model and present a qualitative analysis showing physically why we expect that logarithmic correlations (as opposed to faster growing or slower growing correlations) is the only case which leads to (i) a glass transition (ii) a low temperature phase with a non trivial structure of quasi degenerate distant minima
### A the model
The equilibrium problem of a single particle in a $`d`$-dimensional random potential is defined by the canonical partition function
$$Z[V]=\underset{𝐫}{}e^{\beta V(𝐫)}$$
(1)
where $`\beta =1/T`$ is the inverse temperature, in a sample of finite size $`L`$ (and total number of sites $`L^d`$) and for a given configuration of the random variables $`V(𝐫)`$. The equilibrium Gibbs measure, or probability distribution for the position of the particle is:
$`p(𝐫)=e^{\beta V(𝐫)}/Z[V]`$ (2)
We are interested here in cases where the random variables $`V(𝐫)`$ can be correlated. As discussed below, the statics (and dynamics) of this problem in the limit of large sizes depends on the type of correlations, the distribution of the disorder and the dimensionality of space $`d`$. Some of these cases and their dynamical aspects (such as the Sinai model) have been extensively studied, e.g. in the context of diffusion in random media . Even logarithmic correlations in $`d=2`$ were studied then , but it was not realized at that time that a static glass transition could exist in that case.
Correlated random potentials $`V(𝐫)`$ are most conveniently studied for Gaussian distributions, on which we focus, parametrized by the correlator $`\mathrm{\Gamma }(𝐫,𝐫^{})=\overline{V(𝐫)V(𝐫^{})}`$ (and we choose $`\overline{V(𝐫)}=0`$). Non gaussian extensions will be mentionned. Unless specified otherwise the correlations will be chosen translationally invariant $`\mathrm{\Gamma }(𝐫,𝐫^{})=\mathrm{\Gamma }_L(𝐫𝐫^{})`$ with cyclic boundary conditions, or in (discrete) Fourier space $`\overline{V(𝐪)V(𝐩)}=\mathrm{\Gamma }(𝐪,𝐩)=\mathrm{\Gamma }(𝐪)\delta _{𝐩,𝐪}`$. We will often denote $`\overline{(V(𝐫)V(𝐫^{}))^2}=\stackrel{~}{\mathrm{\Gamma }}(𝐫𝐫^{})=2_𝐪\mathrm{\Gamma }(𝐪)(1\mathrm{cos}(𝐪.(𝐫𝐫^{})))`$ (with $`_𝐪=\frac{1}{L}_𝐪\frac{d^d𝐪}{(2\pi )^d}`$).
One important quantity is the free energy:
$`F[V]=T\mathrm{ln}Z[V]`$ (3)
and, since it fluctuates from configuration to configuration, as $`F[V]=\overline{F[V]}+\delta F[V]`$ we will be interested in its average $`F=\overline{F[V]}`$ and in its distribution. From the convexity of the logarithm follows the well known exact bound for $`F`$ in terms of the annealed free energy $`F_A`$:
$`T\overline{\mathrm{ln}Z}=FF_A=T\mathrm{ln}\overline{Z}`$ (4)
$`F(Td\mathrm{ln}L+{\displaystyle \frac{1}{2T}}\overline{V(𝐫)^\mathrm{𝟐}})`$ (5)
for the Gaussian case.
In this paper we will mainly focus on the case of correlations growing logarithmically with distance:
$`\overline{(V(𝐫)V(𝐫^{}))^2}4\sigma \mathrm{ln}{\displaystyle \frac{|𝐫𝐫^{}|}{a}}a|𝐫𝐫^{}|L`$ (6)
which also requires a small distance ultraviolet (UV) cutoff $`a`$ (we can set here $`a=1`$ in accordance with the definition 1 of a discrete model, but in the following Sections we will consider a continuum version and vary $`a`$). This behaviour is achieved in $`d`$ dimension by choosing a propagator in Fourier space $`\mathrm{\Gamma }(𝐪)\frac{2\sigma (2\pi )^d}{S_dq^d}`$. The $`d=2`$ case is also of special interest as the propagator is the usual Coulomb one:
$`\mathrm{\Gamma }(𝐪){\displaystyle \frac{4\pi \sigma }{q^2}}`$ (7)
and boundary conditions must be specified later on. It is important to note that for LR correlations the single site variance $`\overline{V(𝐫)^2}=\mathrm{\Gamma }_L(0)`$ diverges with the system size, e.g. for (6) one has $`\mathrm{\Gamma }_L(0)2\sigma \mathrm{ln}(L/a)`$.
For such logarithmic correlations (as well as for weaker correlations ) one will find that $`F`$ scales as $`d\mathrm{ln}L`$ (consistent with the number of degree of freedom being $`L^d`$ in this problem). Thus it is natural to define the intensive free energy:
$`f(\beta )=\underset{L+\mathrm{}}{lim}{\displaystyle \frac{F[V]}{d\mathrm{ln}L}}`$ (8)
which will be found to be self averaging. The above bound gives:
$`f(\beta )\left({\displaystyle \frac{1}{\beta }}+{\displaystyle \frac{\sigma }{d}}\beta \right)`$ (9)
Thus we will find that $`F[V]f(\beta )d\mathrm{ln}L`$ with subdominant corrections. These corrections have a non fluctuating universal $`O(\mathrm{ln}(\mathrm{ln}L))`$ piece, as well as an $`O(1)`$ fluctuating part $`\delta F[V]`$ which we will both study.
### B the REM approximation
A useful approximation to the problem studied here, which can be called the REM approximation, consists in neglecting all correlations but keeping the on site variance exact :
$`\mathrm{\Gamma }_L(𝐫)\mathrm{\Gamma }_L^{\mathrm{REM}}(𝐫)=\mathrm{\Gamma }_L(0)\delta _{𝐫,𝐫^{}}=2\sigma \mathrm{ln}\left({\displaystyle \frac{L}{a}}\right)\delta _{𝐫,𝐫^{}}`$ (10)
The corresponding Gaussian REM model can then be solved, being identical to , and one finds that it exhibits a transition at $`\beta _c=\sqrt{d/\sigma }`$ with:
$`f(\beta )=\left({\displaystyle \frac{1}{\beta }}+{\displaystyle \frac{\sigma }{d}}\beta \right)\beta <\beta _c`$ (11a)
$`f(\beta )={\displaystyle \frac{2}{\beta _c}}\beta >\beta _c`$ (11b)
Most previous studies of the original model (all in $`d=2`$) amount to study the REM approximation and argue that it is a good approximation. Indeed, as we will also find here, this REM approximation appears to give the exact result for some observables (e.g. for $`f(\beta )`$). In particular, it does seem to give correctly the transition temperature $`\beta _c`$.
### C beyond the REM approximation: extremal statistics of correlated variables
Since it is not obvious a priori why logarithmic correlations can be considered so weak as to be neglected, one would like to go beyond the REM approximation and describe the effect of the neglected correlations . One would like to understand why this approximation works for some observables (and for which ones) and whether it gives exactly the universality class of the model (i.e all universal behaviour of observables). The answer to the latter is negative: as our analysis will reveal, the correlations do matter for the more detailed behaviour and the original model (1,6) is not in the same universality class as the REM model.
In fact, the problem at hand is related to describing universal features of the extremal value statistics for a set of correlated random variables. Indeed, the zero temperature limit ($`T=0`$ for fixed $`L`$) of the problem defined by (1) amounts to finding the distribution of the minimum $`lim_{T0}T\mathrm{ln}Z_L=V_{min}=min_r(\{V_r\})`$ of a set of correlated random variables. In the case of uncorrelated (or short range correlated) variables a lot is known in probability theory on this problem (see e.g. ), some of which is summarized in Appendix B. For the type of distributions considered here (gaussian and some extensions) the distribution of the minimum $`V_{min}`$ has a strong universality property, being given, up to - non trivial - rescaling and shift (see Appendix B and below), by the Gumbell distribution:
$`\text{Proba}(y<x)=𝒢(x)=\mathrm{exp}(e^x)`$ (12)
The Gumbell distribution thus appears as the distribution of the zero temperature free energy in the REM. For the case of a Gaussian distribution the standard probability theory results are usually given in terms of a variable $`X_r`$ such that $`\overline{X_r^2}=1`$. One can simply rescale $`V_r=\sqrt{2\sigma \mathrm{ln}L}X_r`$ from Appendix B and get:
$`V_{min}=2\sqrt{\sigma d}\mathrm{ln}L{\displaystyle \frac{1}{2}}\sqrt{{\displaystyle \frac{\sigma }{d}}}\mathrm{ln}(4\pi d\mathrm{ln}L)+\sqrt{{\displaystyle \frac{\sigma }{d}}}y`$ (13)
where $`y`$ is distributed as in (12).
Much less is known in the case of variables with stronger correlations studied here, though it is more important in practice. The statistics of $`V_{min}`$ in the logarithmically correlated case is thus one of the open issues discussed here. One key question is to determine what is universal in the distribution of the minimum of correlated variables. Here, we can formulate the question as follows: given Gaussian random variables satisfying (6), what in the distribution of the minimum (i.e of the ground state energy for fixed large $`L`$) is universal, i.e depends only on $`\sigma `$ and not on the details of the correlator $`\mathrm{\Gamma }(𝐫)`$ at short scale. Writing
$`V_{min}\overline{V_{min}}+\delta V_{\mathrm{min}}`$ (14)
one finds, for the logarithmic correlator, that the averaged ground state energy must satisfy:
$`\overline{V_{\mathrm{min}}}2\sqrt{\sigma d}\mathrm{ln}L`$ (15)
which follows from the above annealed bound, together with the fact that $`f/T=S0`$. Furthermore one will find here that $`V_{\mathrm{min}}e_{min}d\mathrm{ln}L`$ up to a positive subdominant - universal - piece and that $`e_{min}=2\sqrt{\sigma /d}`$ saturates the bound. In the distribution of $`\delta V_{\mathrm{min}}O(1)`$ we can clearly expect less universality that in the problem of random variables with short range correlations .
### D Qualitative study of a particle in a random potential
Before describing the RG method which allows to go beyond the REM approximation, let us give some simple qualitative arguments and numerical results which illustrate the main physics of the thermodynamics of a particle in a correlated random potential. To put things in context we discuss several types of correlations (short range, long range, and marginal). We focus on $`d=1`$ for simplicity but the arguments extend to any finite $`d`$.
Whether there is a single phase or not here comes simply from whether the entropy of typical sites wins or not over the energy of the low energy sites. When there is a low temperature phase, to decide its structure one must pay special attention to distant secondary local minima.
Indeed, when there is a low temperature phase, it is controlled by the regions with most negative potential. To investigate its structure one can start, for a given system of size $`L`$, with the $`T=0`$ state which is determined by the absolute minimum over the system, denoted $`V_{min}`$ and located at $`𝐫_{min}`$. At $`T`$ very small but strictly positive, each (low lying) secondary local minima $`V`$ will also be occupied with a probability $`e^{(VV_{\mathrm{min}})/T}`$ which is very small except when $`VV_{min}O(T)`$. Thus to characterize the low temperature phase we need to know how many of these secondary minima exist and where they are located. For a smooth enough disorder (see e.g. Fig 3) there will always be “trivial” secondary local minima in the vicinity of $`r_{min}`$. To eliminate these, we define $`V_{min2}(R)`$ as the next lowest minimum constrained to be at a distance at least $`R`$ of the absolute minimum. An interesting quantity to study is then the distribution $`P_{R,L}(\mathrm{\Delta }E)`$ of $`\mathrm{\Delta }E(R,L)=V_{min2}(R)V_{min}`$ over environments (which a priori depends on $`R`$ and $`L`$)
We now distinguish three main cases, according to the behaviour of the correlator $`\overline{(V(𝐫)V(𝐫^{}))^2}=\stackrel{~}{\mathrm{\Gamma }}(𝐫𝐫^{})`$ at large scale (we restrict to gaussian potentials ). In these three cases the distribution $`P_{R,L}(\mathrm{\Delta }E)`$ has markedly different behaviours as illustrated in Fig. 2:
(i) short range correlations. i.e $`\stackrel{~}{\mathrm{\Gamma }}(𝐫)Cst`$ at large $`𝐫`$, equivalently $`\mathrm{\Gamma }_L(𝐫)0`$ at large $`𝐫`$ (or e.g. $`\mathrm{\Gamma }_𝐪q^{d+\delta }`$ with $`\delta >0`$). In this case it is clear that there is only a high temperature phase in any finite $`d`$ and no phase transition. The entropy $`Td\mathrm{ln}L`$ of typical sites (of energy typically $`O(1)`$) always wins over the energy of optimal sites ($`V_{\mathrm{min}}\sqrt{2\sigma d\mathrm{ln}L}`$ for gaussian distributions with on-site variance $`\sigma `$ ). The optimal energy $`V_{min}`$ can be estimated using $`1/L^d=_{\mathrm{}}^{V_{min}}P_1(V)𝑑V`$ in terms of the single site distribution $`P_1(V)`$, which yields the exact leading behavior for uncorrelated disorder (and also for weak enough correlations - see Appendix B). Thus, the particle is delocalized over the system for all $`T>0`$. One estimates the number of states within $`\mathrm{\Delta }E`$ of the minimum as $`N(\mathrm{\Delta }E)L^d_{V_{\mathrm{min}}}^{V_{\mathrm{min}}+\mathrm{\Delta }E}P_1(V)𝑑V\mathrm{exp}(\mathrm{\Delta }E\sqrt{2d\mathrm{ln}L}/\sqrt{\sigma })`$ for a Gaussian distribution. Thus there is a large number of sites almost degenerate with the absolute minimum $`V_{min}`$, separated by finite barriers, and $`\mathrm{\Delta }E_{typ}`$ decays to 0 as a power of $`1/\mathrm{ln}L`$ ($`1/\sqrt{\mathrm{ln}L}`$ for a Gaussian) . These minima however are irrelevant for the thermodynamics of the system at a fixed finite temperature.
For these minima to play a role and to obtain a transition even for SR disorder one needs to perform some artificial rescaling, as in the REM model : either, at fixed size, to concentrate on the very low $`T`$ region, (e.g. take $`\beta \mathrm{ln}L`$ in the Gaussian case), or equivalently, to rescale disorder with the system size. By making disorder larger as the system increases, for instance using $`P_1(V)e^{|V/V_{typ}|^\alpha }`$ with $`V_{typ}(\mathrm{ln}L)^{11/\alpha }`$ one recovers artificially a transition . For $`\alpha =2`$ and uncorrelated $`V(𝐫)`$ this is exactly the REM studied in . There, the simple argument for the transition is that the averaged density of sites at energy $`E=V`$ is $`\overline{\mathrm{\Omega }(E)}=L^de^{E^2/(2\sigma _L)}/\sqrt{2\pi \sigma _L}`$ (related to the annealed partition sum via $`\overline{Z}=_Ee^{\beta E}\overline{\mathrm{\Omega }(E)}`$). If $`\sigma _L=\sigma `$ is not rescaled the average energy is $`O(1)`$ and the huge entropy of these states always wins. If $`\sigma `$ scales with $`L`$ as $`\sigma _L2\sigma \mathrm{ln}L`$ then there is a transition at $`\beta _c=\sqrt{d/\sigma }`$. Indeed, $`\overline{\mathrm{\Omega }(E)}\mathrm{exp}(d\mathrm{ln}L(1(e/e_{min})^2))`$ where $`e=E/(d\mathrm{ln}L)`$ and $`e_{min}=E_{min}/(d\mathrm{ln}L)=2\sqrt{\sigma /d}`$ and there is a saddle point in $`\overline{Z}`$ at $`E/(d\mathrm{ln}L)=e_{sp}=\beta e_{min}^{}{}_{}{}^{2}/2`$: since $`e_{sp}`$ must be larger than $`e_{min}`$ (as $`\overline{\mathrm{\Omega }(E)}`$ cannot become smaller than $`1`$) the saddle point cannot be valid below $`T_c=1/\beta _c=e_{min}/2=\sqrt{\sigma /d}`$ and the system freezes in low lying states. Although this argument implicitly relies on using $`\mathrm{ln}\overline{\mathrm{\Omega }(E)}`$ instead of $`\overline{\mathrm{ln}\mathrm{\Omega }(E)}`$ it does give the correct picture for the REM, as shown in .
This picture generalizes to correlated potentials provided $`\mathrm{\Gamma }_L(r)`$ decreases fast enough at large $`r`$. The decay must be faster than $`1/\mathrm{ln}r`$ (which is a rather slow decay) as indicated by the theorems recalled in Appendix B or also by a simple argument given in Appendix B 2 c. Finally, let us point out also that another way to obtain a transition for SR disorder is to take the $`d=\mathrm{}`$ limit before taking the large $`L`$ limit: there the model (even without rescaling) always exhibit a transition (in the statics and in the dynamics).
(ii) long range correlations: when the typical $`V(𝐫)V(𝐫^{})`$ grows with distance as a power law $`\stackrel{~}{\mathrm{\Gamma }}(𝐫)|𝐫|^\delta `$, there is only a low temperature phase and no transition. The particle is now always localized near the absolute minimum of the potential in the system at $`𝐫_{min}`$. The typical minimum energy $`V_{\mathrm{min}}`$ grows as $`L^{\delta /2}`$ and thus overcomes the entropy $`Td\mathrm{ln}L`$ which is never sufficient to delocalize the particle. The structure of this single low temperature phase is simple: there are no quasi degenerate minima separated by infinite distance (and thus also by infinite barriers) in the thermodynamic limit. As can be seen on Fig. 3 there is typically a single minimum, with many secondary ones near it, but none far away. More precisely, as $`L\mathrm{}`$, the probability that the lowest energy excitation $`\mathrm{\Delta }E(R,L)`$ above the ground state (a distance at least $`RL^c`$ from $`𝐫_{min}`$) be smaller than a fixed finite (arbitrary) value decays algebraically to $`0`$ with $`L`$ (and $`\mathrm{\Delta }E_{typ}`$ and $`\overline{\mathrm{\Delta }E}`$ increase algebraically with $`L`$). This is the scenario familiar from the droplet picture , with $`Proba(\mathrm{\Delta }E<T)TL^{\delta /2}`$ (i.e in some configurations which become more and more rare as $`L+\mathrm{}`$, there are two far away quasidegenerate ground states). In some cases, e.g. in $`d=1`$ Sinai’s model ($`\delta =1`$) the distribution of rare events with quasi degenerate minima has been studied extensively . For instance it has been shown that there is a well defined limit distribution $`Q(R)dR`$ (when $`L+\mathrm{}`$) to find quasi degenerate minima at fixed distance between $`R`$ and $`R+dR`$, with $`Q(R)R^{3/2}`$ at large $`R`$.
(iii) marginal case, logarithmic correlations: the most interesting case is when correlations grow as $`\stackrel{~}{\mathrm{\Gamma }}(𝐫)4\sigma \mathrm{ln}|𝐫|`$. A typical logarithmically correlated landscape is illustrated in Fig. 4. One can already see that, contrarily to Fig. 3, it has states with similar energies far away.
Given the growth of correlations one sees that the typical energy differences over a distance $`L`$ scale as $`(V(0)V(L))_{typ}\pm \sqrt{4\sigma \mathrm{ln}L}`$. Computing the minimum energy is a harder task here, but if one estimates it as in through the REM approximation $`1/L^d=_{\mathrm{}}^{V_{min}}P_1(V)𝑑V`$ (which neglects correlations), one finds that it behaves as $`V_{min}2\sqrt{\sigma d}\mathrm{ln}L`$ (for Gaussian disorder). This estimate appears rather uncontrolled here since correlations grow with distance, while the theorems for uncorrelated random variables apply a priori only for correlations decaying slower than $`1/\mathrm{ln}r`$. In fact the situation is a bit more complex, and as we will find below from the RG and our numerics, the leading behaviour of $`V_{min}`$ with $`\mathrm{ln}L`$ is still correctly given by the REM approximation, although the next subleading -universal- correction is not. Thus the energy of the minimum $`2\sqrt{\sigma d}\mathrm{ln}L`$ can now balance the entropy of typical sites $`Td\mathrm{ln}L`$ which yields the possibility of a transition. The REM approximation of the model indeed yields a transition at $`T_c=\sqrt{\sigma /d}`$ between a high temperature phase for $`\beta <\beta _c=\sqrt{\sigma /d}`$ and a frozen phase $`\beta >\beta _c`$. This scenario is confirmed by various approaches in the following sections.
An interesting feature of this model is that the low temperature phase exhibits a non trivial structure. Unlike long range disorder discussed above, for logarithmic correlations we find that the low temperature phase is dominated not by one, but by a few states in the thermodynamic limit. This is in stark contrast with the standard droplet picture and is reminiscent of the replica symmetry breaking phenomenology, even though we are dealing here with a very simple finite dimensional system.
One can visualize the transition, and the peculiar nature of the low temperature phase in Fig. 5,6, where a typical Gibbs measure $`p(𝐫)`$ is shown in both phases : is fairly delocalized at $`T>T_c`$ (Fig. 5) but peaks around a few states when $`T<T_c`$ (Fig. 6) separated by a distance of the order of the system size.
This peculiar nature of the frozen phase can be tested by showing that distant secondary local minima with a finite $`\mathrm{\Delta }E`$ exist with finite probability in the thermodynamic limit. Thus we have investigated numerically the distribution $`P_{R,L}(\mathrm{\Delta }E)`$ of the lowest excitation. As illustrated in (Fig. 2), if the phase is non trivial, we expect that this distribution has a well defined limit for e.g. $`R=L/3`$ when $`L\mathrm{}`$ with a finite typical $`\mathrm{\Delta }E`$. Contrarily to the LR disorder, we expect the probability that e.g. $`\mathrm{\Delta }E(L/3,L)`$ be smaller than a fixed number to saturate (not to decrease) as $`L\mathrm{}`$, i.e that there is a fixed probability that a second state within $`\mathrm{\Delta }E`$ exists far away (as was already apparent in Fig. 4). We show in Fig. 7, Fig. 8 and Fig. 9 numerical evidence that this distribution has a well defined limit (the details of the simulation are discussed in Section IV). Finite size effects are clearly important in this system, but their magnitude appears compatible with the predictions of our RG approach, as discussed below. Thus we conclude that the numerics are consistent with the existence of such a limit distribution and hence with a frozen phase with a non trivial structure.
## III renormalization group approach
### A Idea of the method
We now study the model (1, 6) using a renormalization approach introduced by us to study $`d=2`$ disordered XY models . There, one is led to study a neutral collection of interacting $`\pm 1`$ charges (XY vortices) in a random potential $`\pm V(𝐫)`$ with (6). The single particle problem studied here amounts to restrict the Coulomb gas RG of to the sector of a single $`+1`$ charge. Here however there is no charge neutrality and one must be careful to study a system of finite size $`L`$, as some quantities (such as $`\overline{V(𝐫)^2}`$) explicitly depend on $`L`$, while appropriately defined quantities have a well defined thermodynamic limit.
The idea is first to formulate the problem in the continuum, with a short distance cutoff $`a`$:
$$Z=\frac{d^d𝐫}{a^d}e^{\beta V(𝐫)}$$
(16)
and an appropriately defined cutoff-dependent distribution for $`V(𝐫)`$, and second, by coarse graining infinitesimally, to relate the problem defined with a cutoff $`a^{}=ae^{dl}`$ to the problem with a cutoff $`a`$. In general, this implies to be able to follow under this transformation the full probability measure of the potential $`V(𝐫)`$, which is quite difficult, as complicated correlations can be generated under coarse graining. In some very favorable cases, for instance in the $`d=1`$ Sinai landscape (where $`V(r)`$ performs a random walk as a function of $`r`$ \- case $`\delta =1`$), it is possible to follow analytically an asymptotically exact RG transformation (in the statics and in the dynamics ). There a very specific real space decimation procedure is required, which can in principle be extended here, although it may not be tractable beyond numerics. The present case of the logarithmically correlated potential is thus a priori less favorable but still, thanks to some known properties of the Coulomb potential, a RG method a la Kosterlitz can be constructed which, we argue, should describe correctly all the universal properties of the model. There are two possible derivations, one which uses replicas and is more precise, and the other one without. We start with the latter, which is physically more transparent.
The key observation is that before (and also after) coarse graining, the logarithmically correlated disorder studied here can naturally be decomposed into two parts as:
$`V(𝐫)=V^>(𝐫)+v(𝐫)`$ (17)
where $`V^>(𝐫)`$ is a smooth gaussian disorder with the same LR correlations as the initial $`V(𝐫)`$ which represents the contribution of the scales larger than the cutoff $`ae^l`$, and $`v(𝐫)`$ is a local short range random potential which represents the contribution of scales smaller than, or of the order of, the cutoff $`ae^l`$. In the starting model $`v(𝐫)`$ appears naturally as a gaussian variable (see below). After coarse graining, $`v(𝐫)`$ does not remain gaussian, but it does remain uncorrelated in space (i.e correlations of short range $`a`$). The decomposition (17) allows to follow the distribution of the $`V(𝐫)`$ under coarse graining in a tractable way.
The precise way of decomposing the disorder in (17) depends on the details of the cutoff procedure, but should not matter as far as universal properties are concerned. For illustration let us indicate a simple way to do it, a more detailed discussion is given in . It starts with the well known continuum approximation in $`d=2`$ of the lattice Coulomb potential $`\stackrel{~}{\mathrm{\Gamma }}(𝐫𝐫^{})4\sigma \left(\mathrm{ln}\left(\frac{|𝐫𝐫^{}|}{a}\right)+\gamma \right)(1\delta ^{(a)}(𝐫𝐫^{}))`$ where $`\delta ^{(a)}(𝐫𝐫^{})=1`$ for $`|𝐫𝐫^{}|<a`$ and $`0`$ otherwise ($`\gamma =\mathrm{ln}(2\sqrt{2}e^C)`$ and $`C=0.5772`$ is the Euler constant). This decomposition can be performed more generally, e.g. with other short-distance regularization of the potential $`\stackrel{~}{\mathrm{\Gamma }}(𝐫)`$ (which preserve the large distance logarithmic behaviour) and in any $`d`$, which amounts to modify the value of $`\gamma `$. Using this approximation the bare disorder (6) can indeed be rewritten equivalently as a sum (17) of two gaussian disorder $`V^>(𝐫)`$ and $`v(𝐫)`$ with no cross correlations and with respective correlators:
$`\overline{(V^>(𝐫)V^>(𝐫^{}))^2}`$ $`=`$ $`4\sigma \mathrm{ln}{\displaystyle \frac{|𝐫𝐫^{}|}{a}}(1\delta ^{(a)}(𝐫𝐫^{}))`$ (18)
$`\overline{v(𝐫)v(𝐫^{})}`$ $`=`$ $`2\sigma \gamma \delta ^{(a)}(𝐫𝐫^{})`$ (19)
With this definition, the problem to be studied is rewritten as:
$$Z=\frac{d^d𝐫}{a^d}z(𝐫)e^{\beta V^>(𝐫)}z(𝐫)=e^{\beta v(𝐫)}$$
(20)
We can now study the behaviour of the model under a change of cut-off. There are two main contributions from eliminated short length scales variables. The first one can be seen most simply by rewriting the correlator in (18):
$`\overline{(V^>(𝐫)V^>(𝐫^{}))^2}`$ $`=`$ $`4\sigma \mathrm{ln}{\displaystyle \frac{|𝐫𝐫^{}|}{a^{}}}(1\delta ^{(a^{})}(𝐫𝐫^{}))`$ (21)
$`+`$ $`4\sigma dl(1\delta ^{(a^{})}(𝐫𝐫^{}))`$ (22)
explicitly as the sum of a new LR disorder correlator with cutoff $`a^{}=ae^{dl}`$ and a SR disorder correlator (we have discarded terms of order $`O(dl^2)`$). Thus the original problem with cutoff $`a`$ can be rewritten as one with cutoff $`a^{}`$ with (i) a new gaussian LR disorder with identical form of the correlator (18) with $`a`$ replaced by $`a^{}`$ (ii) a new short range disorder $`v(𝐫)v(𝐫)+dv(𝐫)`$ with $`\overline{dv(𝐫)dv(𝐫^{})}=2\sigma dl\delta ^{(a)}(𝐫𝐫^{})`$ since it is clear from (22) that when $`aae^{dl}`$ the LR disorder produces an additive gaussian contribution $`dv`$ to the SR disorder.
The second contribution resulting from a change of cutoff is that neighboring regions will merge. Points $`𝐫_1`$ and $`𝐫_2`$ previously separated as $`a<|𝐫_1𝐫_2|<ae^{dl}`$ should now be considered as within the same region. The second important observation is that the resulting transformation can only affect the SR part $`v(𝐫)`$ of the disorder. Indeed, in the region $`a<|𝐫_1𝐫_2|<ae^{dl}`$ the LR part $`V^>(𝐫)`$ can be considered as constant up to higher order terms of order $`dl`$. One must view this coarse graining as resulting in a ”fusion of local environments” : the two local partition sum variables $`z(𝐫_1)`$ and $`z(𝐫_2)`$ combine into a single one $`z(𝐫)`$ according to a rule which we will write as $`z(𝐫)=z(𝐫_1)+z(𝐫_2)`$. The exact choice of the form of this fusion rule is again dependent of the cutoff procedure and thus to a large extent arbitrary.
Putting together these two contributions we obtain the following RG equation for the distribution $`P_l(z)`$ of the local disorder $`z=e^{\beta v}`$ variable (also called ”fugacity” in the Coulomb Gas context).
$`_lP(z)=\beta ^2\sigma \left(1+z_z\right)^2PdP(z)+d{\displaystyle _{z^{}z^{\prime \prime }}}P(z^{})P(z^{\prime \prime })\delta (z(z^{}+z^{\prime \prime }))`$ (23)
This equation also describes the evolution of the universal part of the total free energy distribution with the system size. Indeed, the total partition function can be written at any scale as:
$$Z(\beta )=\frac{d^d𝐫}{a^d}e^{\beta V(𝐫)}\frac{d^d𝐫}{(ae^l)^d}z_l(𝐫)e^{\beta V_l^>(𝐫)}z_{(l^{})}$$
(24)
where the $`z_l(𝐫)`$ are independent variables distributed with $`P_l(z)`$ and the $`V_l^>(𝐫)`$ are gaussian distributed as (18). In the last equality we have coarse grained up to the system size : $`L=ae^l^{}`$. At this scale, there remains a single site of (random) fugacity $`z_l^{}`$. Thus the distribution function of the partition function $`Z(\beta )`$ can be deduced from the distribution of the random fugacities at scale $`l^{}`$. The distribution of the free energy $`F=T\mathrm{ln}Z`$ is thus given by $`\stackrel{~}{P}_l^{}(v=F)`$ (where $`\stackrel{~}{P}(v)dv=P(z)dz`$ from the change of variable from $`z`$ to $`v=T\mathrm{ln}z`$). Note that the $``$ in (24) means that these distributions are the same a priori only up to subdominant non universal terms (multiplicative for $`Z`$ and additive for $`\mathrm{ln}Z`$).
For a fixed system size $`L`$, the above RG equation describes the evolution with the scale $`l`$ smaller that $`l^{}`$ of the distribution of $`z(𝐫)`$, which is the local partition sum over scales around $`𝐫`$ smaller or equal to $`ae^l`$ (i.e of a “local free energy” $`T\mathrm{ln}z_l(𝐫)=v_l(𝐫)`$). The remaining long wavelength disorder at that scale, $`V_l^>(𝐫)`$ should still be taken into account when computing the total partition sum.
It is striking that the equation (35) is identical to the RG equation for the partition function of a continuum version of a directed polymer on a Cayley tree (a so called branching process ). We note that it has been derived here for a problem with complete (statistical) translational invariance, with no ad-hoc assumption about an underlying tree structure and simply adapting to the present problem the Coulomb gas renormalization a la Kosterlitz. That the correspondence between the two problems naturally appears within the RG with no additional assumptions, is even more apparent on the derivation using replicas of the next section. Thus we consider that this establishes on a firm footing the strong connection between the two problems.
Before analyzing the consequences of the above RG equation let us sketch the more precise derivation using replicas. Other derivations without replicas are also possible and we refer the reader to for more details.
### B derivation of the RG equation using replicas
Let us consider the whole set of moments $`\overline{Z^m}`$ which encode for the distribution function $`P[Z]`$. They can be written as:
$`\overline{Z^m}={\displaystyle \frac{d^d𝐫_1}{a^d}\mathrm{}\frac{d^d𝐫_m}{a^d}e^{\frac{\beta ^2}{2}\overline{\left(_{i=1,..m}V(𝐫_i)\right)^2}}}`$ (25)
This can be rewritten as:
$$\begin{array}{c}\overline{Z^m}=\frac{d^d𝐫_1}{a^d}\mathrm{}\frac{d^d𝐫_m}{a^d}\hfill \\ \hfill e^{\frac{\beta ^2}{4}_{ij=1,..m}\stackrel{~}{\mathrm{\Gamma }}(𝐫_i𝐫_j)}e^{m^2\sigma \beta ^2\mathrm{ln}\frac{L}{a}}\end{array}$$
(26)
We have used that $`\mathrm{\Gamma }(𝐫,𝐫)=\mathrm{\Gamma }_L(0)=2\sigma \mathrm{ln}\frac{L}{a}`$. One can choose a regularisation, e.g. $`\mathrm{\Gamma }(𝐫𝐫^{})=\overline{V(𝐫)V(𝐫^{})}=\sigma \mathrm{ln}\frac{|𝐫𝐫^{}|^2+a^2}{L^2}`$. Notice that only the large distance behaviour of the above correlator is important for the following renormalization.
We now switch to another representation of the replica partition sum. (26) is a partition sum of $`m`$ particles located at $`𝐫_1,..𝐫_m`$ corresponding to $`m`$ replicas. Now instead we will index the configurations using (vector) columnar replicated charges. To each point $`𝐫`$, within a hard core size $`a`$, we associate a $`m`$-component vector $`𝐧`$ whose components $`n^i(𝐫)`$ are either $`1`$ or $`0`$ depending on whether the particle corresponding to the $`i`$-th replica is present within $`a`$ of $`𝐫`$ ($`|𝐫𝐫_i|<a`$) or not. These charges thus correspond to $`𝐧=(0,1,0,\mathrm{..0},1,1)`$ since several replicas can be present near a given point. Choosing a columnar hard core for the vector charges corresponds to a choice of cutoff, which is arbitrary, but the universal features of the renormalization should not depend on it .
The $`m`$th moment of $`P[Z]`$ then read
$`\overline{Z^m}=\left({\displaystyle \frac{L}{a}}\right)^{\beta ^2\sigma m^2}{\displaystyle \underset{\{n_\alpha ^i\}}{\overset{}{}}}{\displaystyle \underset{\alpha }{}}Y[𝐧_\alpha ]`$ (27)
$`{\displaystyle _{|𝐫_\alpha 𝐫_\alpha ^{}|a}}{\displaystyle \frac{d^d𝐫_\alpha }{a^d}}\mathrm{exp}\left(2\beta ^2\sigma {\displaystyle \underset{\alpha <\alpha ^{}}{}}n_\alpha n_\alpha ^{}\mathrm{ln}\left({\displaystyle \frac{|𝐫_\alpha 𝐫_\alpha ^{}|}{a}}\right)\right)`$ (28)
where the primed sum correspond to a sum over all distinct non zero configurations of replica charges $`𝐧_\alpha `$ at sites $`𝐫_\alpha `$. We have defined $`n_\alpha =_in_\alpha ^i`$ as the total number of replicas present in a given charge ($`n_\alpha ^i=1`$). The quantities $`Y[𝐧]`$ are functions of the local vector charge and are the so-called vector charge fugacities. In the bare model they appear as soon as the continuum approximation to the lattice Green function is used and read $`Y[𝐧]=e^{2\sigma \gamma n^2}`$. Since we are studying a single particle problem, there is also an important global constraint on the configuration sum that only one particle in any replica $`i`$ is present in the system, i.e:
$`{\displaystyle \underset{\alpha }{}}n_\alpha ^i=1`$ (29)
which is preserved by the RG.
The RG equations for this model read:
$`_lY[𝐧]`$ $`=`$ $`\left(d+\beta ^2\sigma n^2\right)Y[𝐧]`$ (31)
$`+{\displaystyle \frac{S_{d1}}{2}}{\displaystyle \underset{𝐧^{}+𝐧^{\prime \prime }=𝐧}{}}Y[𝐧^{}]Y[𝐧^{\prime \prime }]`$
where the sum is over $`𝐧^{}`$ and $`𝐧^{\prime \prime }`$ non zero vector charges (also $`𝐧`$ is non zero) and $`S_{d1}`$ is the volume of the unit sphere in dimension $`d`$. We recall that $`n=_{i=1}^mn_i`$. These equations are obtained by a generalization of the Kosterlitz procedure as follows. The first term comes from explicit cutoff dependence in (27). Upon increasing the cutoff infinitesimally $`aa^{}=ae^{dl}`$ the integration measure and the $`a`$ dependence in all logarithms combine to give $`Y[𝐧_\alpha ]Y[𝐧_\alpha ]e^{dl(d+\sigma \beta ^2n_\alpha ^2)}`$. We have used that $`2_{\alpha <\alpha ^{}}n_\alpha n_\alpha ^{}=m^2_\alpha n_\alpha ^2`$ which holds due to (29). The last term in the above equation (31) comes from the fusion of replica charges upon increase of the cutoff. The above RG equations hold for any $`m`$.
We should now look for solutions of this set of equations analytically continued to $`m0`$. One way to do that is to find a convenient parametrization for the set of $`Y[𝐧]`$. Here we preserve replica permutation symmetry within the RG and we can thus choose $`Y[𝐧]`$ to be a function of $`n=_in_i`$ only. Then we define the parametrization $`Y[n]=𝑑z\mathrm{\Phi }_l(z)z^n=𝑑u\stackrel{~}{\mathrm{\Phi }}_l(v)e^{\beta nv}`$. The different terms in the above equation then translate into
$`n^2Y[𝐧]`$ $`=`$ $`{\displaystyle 𝑑ve^{\beta nv}(\beta ^1_v)^2\stackrel{~}{\mathrm{\Phi }}_l(v)}`$ (32)
$`{\displaystyle \underset{𝐧^{(1)}+𝐧^{(2)}=𝐧}{}}Y[𝐧_1]Y[𝐧_2]`$ $`=`$ $`{\displaystyle _{z^{},z^{\prime \prime }}}\mathrm{\Phi }_l(z^{})\mathrm{\Phi }_l(z^{\prime \prime })\delta (zz^{}z^{\prime \prime })`$ (34)
$`2𝒩\mathrm{\Phi }_l(z)+\delta (z)𝒩^2`$
where $`𝒩=_z\mathrm{\Phi }_l(z)`$. One then easily converts the equation for $`\mathrm{\Phi }_l(z)`$ into an equation for a normed function $`P_l(z)=\mathrm{\Phi }_l(z)/𝒩_>`$ defined only for $`z>0`$, with $`𝒩_>=_{z>0}\mathrm{\Phi }_l(z)`$ (as in ) by noting that $`𝒩_>`$ converges quickly to $`𝒩_>=2d/S_{d1}`$. The resulting equation for $`P_l(z)`$ is exactly the one (23) given above, and its physical interpretation in terms of the probability distribution of the fugacity (i.e the local partition sum) was given in the previous section.
What is the small parameter which controls the validity of the above RG equations (with and without replicas)? In conventional Coulomb gas context, these RG equations are known to become exact in the dilute limit of non zero (vector) charges . It is easy to see that this corresponds to the tail of the distribution $`P(z)`$ for large $`z`$ (or equivalently small $`v`$). This is further confirmed, a posteriori, by the remarkable universality properties of the resulting non linear RG equation (23), analyzed in the following section, which arises precisely in this region of $`z`$. So to obtain the universal behaviour (e.g. of the distribution of free energy) we are working with sufficient accuracy. On the other hand the bulk of the distribution $`P_l(z)`$ seem to be sensitive to details of the cutoff procedure (e.g. details in the fusion rule) and as discussed below, is thus likely (unless proven otherwise) to be non universal.
### C analysis of RG equation and results
#### 1 KPP front propagation equation and velocity selection
Let us analyze the solutions to the RG equation (23). In terms of the (local free) energy variable $`v(𝐫)=T\mathrm{ln}z(𝐫)`$ (from (20) and its distribution $`P_l(v)=P_l(z=e^{\beta v})\beta e^{\beta v}`$) it has a well defined zero temperature limit, since then the fusion rule simply becomes the extremal rule $`v^{}=\mathrm{min}(v_1,v_2)`$ leading to :
$`_lP(v)=\sigma _v^2P+dP(v)\left(1+2{\displaystyle _v^+\mathrm{}}P(v^{})𝑑v^{}\right)`$ (35)
To be able to work at all temperatures, it is in fact useful to trade the distributions $`P_l(z)`$ or $`P_l(v)`$ for the generating function :
$$G_{l;\beta }(x)=e^{ze^{\beta x}}_{P_l(z)}=e^{e^{\beta (xv)}}_{P_l(v)}$$
(36)
We will sometimes drop the index $`\beta `$. At zero temperature, the double exponential becomes a theta function and $`G_l(x)`$ simply identifies with the distribution function:
$$G_{l;\beta =+\mathrm{}}(x)=_x^+\mathrm{}P_l(v)𝑑v=\text{Proba}(v>x)$$
(37)
and for all $`\beta `$ it is a decreasing function of $`x`$ with $`G_l(x\mathrm{})=1`$ and $`G_l(x+\mathrm{})=0`$. Note the asymptotic behaviour at very large negative $`x`$, $`1G_l(x)<z>_{P_l}e^{\beta x}`$. The temperature appears only via the initial condition and the problem at hand is thus to determine the large $`l`$ behaviour of $`G_l(x)`$ for a given initial condition.
The equation (35) is easily transformed, at all temperatures, into the Kolmogorov (KPP) non linear equation
$`{\displaystyle \frac{1}{d}}_lG(x)={\displaystyle \frac{\sigma }{d}}_x^2G+F[G]`$ (38)
$`F[G]=G(1G)`$ (39)
which describes the diffusive invasion of a stable state $`G=0`$ into an unstable one $`G=1`$. This class of equations admits a family of traveling wave solutions $`G_l(x)=g(x+m(l))`$ which describe a front moving towards negative $`x`$ and located around $`xm(l)`$. This is readily seen by plugging this form in (38) and assuming that $`_lm_\beta (l)c`$ one obtains the equation for the front shape:
$`{\displaystyle \frac{1}{d}}cg^{}(x)={\displaystyle \frac{\sigma }{d}}g^{\prime \prime }(x)+F[g(x)]`$ (40)
The family of such traveling wave solutions $`g_c(x)`$ can thus be parametrized by the velocity $`c`$. (40) simplifies for large negative $`x`$ when $`g1`$. Denoting $`\stackrel{~}{g}=1g`$ and using that $`F[g]\stackrel{~}{g}`$ for $`g1`$, one finds the linearized front equation for $`\stackrel{~}{g}_c`$ :
$`{\displaystyle \frac{1}{d}}c\stackrel{~}{g}^{}={\displaystyle \frac{\sigma }{d}}\stackrel{~}{g}^{\prime \prime }+\stackrel{~}{g}`$ (41)
This equation allows to relate the speed of the front $`c`$ to the asymptotic decay of the front, since if $`\stackrel{~}{g}(x)e^{\alpha x}`$ for large negative $`x`$ one finds:
$`{\displaystyle \frac{c}{d}}={\displaystyle \frac{\sigma }{d}}\alpha +{\displaystyle \frac{1}{\alpha }}`$ (42)
The problem at hand now is to determine toward which of these front solutions $`g_c(x)`$ will $`G_l(x)`$ converge at large $`l`$, and thus what will be the asymptotic front velocity. This velocity will determine the intensive free energy of the original problem. Indeed, the convergence at large $`l`$ of the solutions of non linear equations of the type (38) (with a general $`F[G]`$) towards one of such front solutions, and the corresponding problem of the selection of the front velocity $`c`$, is a famous problem, still under current interest in nonlinear physics .
The simplest argument is to use the fact that for very large negative $`x`$, one must have $`\stackrel{~}{g}(x)e^{\beta x}`$ and thus $`\alpha =\beta `$. This seems to imply that the front velocity is:
$`c=c(\beta )=\left({\displaystyle \frac{\sigma }{d}}\beta +{\displaystyle \frac{1}{\beta }}\right)d`$ (43)
This however is not always true. First note that the curve $`c(\beta )`$ has two branches, i.e that in this naive estimate two different $`\beta `$ would correspond to the same velocity. The special point $`\beta _c=\sqrt{d/\sigma }`$ corresponds to $`c=c^{}=2d\sqrt{\sigma /d}`$. For more general non linear equations one usually relies on the so called marginal stability criterion (e.g. which shows that the large $`\beta `$ branch is unstable and can be eliminated) Here there are rigorous results available : the Bramson theorem ensures the following results, which are independent of the precise form of $`F[G]`$ (up to some rather weak conditions on $`F[G]`$ ):
(i) At high temperature, $`\beta <\beta _c=\sqrt{d/\sigma }`$ the asymptotic front is indeed an exponential for large negative $`x`$ and $`G_l(x)`$ uniformly converges towards the traveling wave solution $`g_{c(\beta )}(x+m(l))`$ where the velocity is given by (43), thus continuously dependent on temperature.
(ii) At low temperature $`\beta \beta _c`$ the velocity freezes to the value $`c=c^{}`$ and the front decays as:
$`\stackrel{~}{g}(x)xe^{\beta _cx}`$ (44)
for large negative $`x`$, thus independent of the temperature. The solution $`G_l(x)`$ uniformly converges towards the traveling wave solution $`g_c^{}(x+m(l))`$. Thus in that regime, one must then distinguish two regions in $`G_l(x)`$ at large $`l`$, the front region and the region very far ahead of the front ($`x+m(l)\sqrt{l}`$) where the decay is again as $`G_l(x)\mathrm{exp}(\beta x)`$ as it should: this will be discussed again below.
There are additional rigorous results from and in particular the remarkable fact that not only the velocity but also the corrections to the velocity are universal (independent of $`F[G]`$) i.e one has for the position of the traveling wave $`m_\beta (l)`$ at “time” $`l`$:
$`m(l)=\left({\displaystyle \frac{\sigma }{d}}\beta +\beta ^1\right)dl+Cst\beta <\beta _c=\sqrt{{\displaystyle \frac{d}{\sigma }}}`$ (46)
$`m(l)=\sqrt{{\displaystyle \frac{\sigma }{d}}}\left(2dl{\displaystyle \frac{1}{2}}\mathrm{ln}l\right)\beta =\beta _c`$ (47)
$`m(l)=\sqrt{{\displaystyle \frac{\sigma }{d}}}\left(2dl{\displaystyle \frac{3}{2}}\mathrm{ln}l\right)\beta >\beta _c`$ (48)
#### 2 Results for the fugacity and free energy distribution and extremal statistics
These results on the KPP equation (38) can now be translated (via (36)) into results for the fugacity distribution $`P_l(z)`$ and for the distribution of free energy $`\stackrel{~}{P}_l(v)`$. One finds that $`P_l(z)`$ and $`\stackrel{~}{P}_l(v)`$ also take the form of a front at large $`l`$, e.g.:
$`\stackrel{~}{P}_l(v)p(v+m(l))`$ (49)
with $`p(v^{})`$ related to $`g(x)`$ by $`g(x)=_v^{}p(v^{})e^{e^{\beta (xv^{})}}`$. Thus we obtain that the local free energy is:
$`\beta ^1\mathrm{ln}zm_\beta (l)`$ (50)
up to a finite constant, where the position of the front $`m_\beta (l)`$ is given above in (46). Using the result (24), $`N=d\mathrm{ln}(L/a)=dl^{}`$ we obtain using the RG that the free energy $`F[V]`$ of the system of size $`L`$ reads:
$`F[V]=f_L(\beta )d\mathrm{ln}L+\delta F`$ (51)
where $`\delta F`$ is a fluctuating part of $`O(1)`$ of probability distribution $`p(\delta F)`$ and the intensive free energy reads:
$`f_L(\beta )=\left({\displaystyle \frac{\beta }{\beta _c^2}}+{\displaystyle \frac{1}{\beta }}\right)+O\left({\displaystyle \frac{1}{\mathrm{ln}L}}\right)\beta <\beta _c=\sqrt{{\displaystyle \frac{d}{\sigma }}}`$ (52a)
$`f_L(\beta )={\displaystyle \frac{1}{\beta _c}}\left(2{\displaystyle \frac{1}{2}}{\displaystyle \frac{\mathrm{ln}(\mathrm{ln}L)}{d\mathrm{ln}L}}\right)+O\left({\displaystyle \frac{1}{\mathrm{ln}L}}\right)\beta =\beta _c`$ (52b)
$`f_L(\beta )={\displaystyle \frac{1}{\beta _c}}\left(2{\displaystyle \frac{3}{2}}{\displaystyle \frac{\mathrm{ln}(\mathrm{ln}L)}{d\mathrm{ln}L}}\right)+O\left({\displaystyle \frac{1}{\mathrm{ln}L}}\right)\beta >\beta _c`$ (52c)
where the factors $`1/2`$ and $`3/2`$ which arise in the finite size corrections are universal.
Thus we have found using our RG method that in any dimension $`d1`$ the original model (1,6) exhibits a phase transition at $`\beta =\beta _c(d)`$. This transition is very similar to the freezing transition of the continuous version of the random directed polymer on the Cayley tree. Our RG thus confirms that the REM approximation (10) to the model does give the transition at the same $`\beta _c`$, and with same asymptotic intensive free energies (11b) as (52c). It allows however for a more detailed study and shows that the universal finite size corrections differ in the two model. In the REM the above formula with the factor $`1/2`$ holds in all the low temperature phase, which is not the case for the present model. Thus the present model is in a different universality class than the REM. The physics that we find here is much closer to the one of the directed polymer on the Cayley tree: it remains to be seen whether this can be extended to other observables.
The RG method also yields the distribution of the $`O(1)`$ fluctuating part $`\delta F`$ of the free energy, and in particular at $`T=0`$ it gives a result for the extremal statistics of the correlated variables. We must now carefully distinguish between what is clearly universal (and thus for which we can be confident that the RG approach gives the exact result) and what may not be (as it depends on the details of the cutoff procedure, yielding e.g. a different KPP non linearity $`F[G]`$).
Let us start with $`T=0`$. We find (cf. (51,52c)) that the minimum $`V_{min}`$ of $`L^d`$ logarithmically correlated variables behaves as:
$`V_{min}=2\sqrt{\sigma d}\mathrm{ln}L+{\displaystyle \frac{3}{2}}\sqrt{{\displaystyle \frac{\sigma }{d}}}\mathrm{ln}(\mathrm{ln}L)+\delta V`$ (53)
and $`\delta V`$ is a fluctuating part of order $`O(1)`$. Since at $`T=0`$ one has $`p(v)=\stackrel{~}{g}^{}(v)`$, from the result (44) we get that the tail of the distribution of $`u=\delta V<\delta V>`$ for $`u\mathrm{}`$ is universal and behaves as:
$`p(u)ue^{\beta _cu}`$ (54)
with $`\beta _c=\sqrt{d/\sigma }`$. Thus we find a distribution different from the Gumbell distribution, and thus correlations do matter.
The question of what is universal in this distribution is non trivial. We find from our method that the full distribution of $`P(u)`$ depends on the detailed form of the front (and thus on $`F[G]`$ and a priori on the cutoff procedure) and is thus less likely to be universal (although this remains to be investigated). Hence we believe that universal features include at least the tail of the distribution (54).
The above result (54) carries through the tail of the distribution of the free energy $`u=F<F>`$ for $`u\mathrm{}`$ for $`T<T_c`$ and it was shown in that for $`T>T_c`$ one has:
$`p(u)e^{u\beta _c^2/\beta }\beta <\beta _c`$ (55)
### D More on fronts, REM via nonlinear RG and extremal statistics
To illustrate how the previous results fit in a broader context, let us show how the simpler properties of extremal statistics of uncorrelated variables and of the random energy model can be recovered within the same RG framework. This provides, en passant, yet another solution of the REM.
#### 1 uncorrelated variables with fixed distribution: Gumbell via RG
Let us consider $`N=e^{ld}=(L/a)^d`$ independent random variables $`V(r)`$ $`r=1,..N`$ with a fixed distribution $`P(V)`$ ($`d`$ here does not play any role as the true variable is $`ld`$ but we keep it for the sake of comparison). The generating function of the distribution of the partition function $`Z[V]=_re^{\beta V(r)}`$ of model (1) reads:
$`G_l(x)=<\mathrm{exp}(Z[V]e^{\beta x})>_{P(V)}`$ (56)
$`=\left({\displaystyle 𝑑VP(V)\mathrm{exp}\left(e^{\beta (xV)}\right)}\right)^{e^{ld}}`$ (57)
It satisfies the equation:
$`{\displaystyle \frac{1}{d}}_l\mathrm{ln}\mathrm{ln}{\displaystyle \frac{1}{G}}=1`$ (58)
Or, interestingly enough, it obeys a KPP type equation with no diffusion term:
$`{\displaystyle \frac{1}{d}}_lG=F[G]`$ (59)
$`F[G]=G\mathrm{ln}G`$ (60)
The Gumbell distribution now emerges naturally from the front solutions of this equation. Writing $`G_l(x)g(\alpha _l(x+m_l))`$ and assuming $`_l(\alpha _lm_l)c`$ yields $`cg^{}=g\mathrm{ln}g`$ whose solutions with the above boundary conditions are $`g(y)=\mathrm{exp}(\gamma e^{y/c})`$ ($`\gamma `$ being a positive constant). We have assumed $`_l\alpha _l0`$. Since there is some freedom of choice for $`\alpha _l`$ and $`m_l`$, one can always set $`c=\gamma =1`$. The determination of the rescaling factors $`\alpha _l`$ and $`m(l)`$ is performed in Appendix C. At $`T=0`$ one has $`P(V_{min})=G^{}(V_{min})`$ and one recovers the known results from probability theory for the convergence to the Gumbell distribution detailed in Appendix B, but the generating function $`G_l(x)`$ takes a Gumbell form also at finite $`T`$.
#### 2 REM via RG
We now turn to an alternative derivation of the solution to the Gaussian REM model using a RG approach and a traveling wave analysis. This allows to make some connections with the correlated case studied previously. Let $`l=\mathrm{ln}L`$ and $`\mathrm{ln}N=ld`$.
We want to write a RG equation for:
$`G_l(x)=\left(e^{e^{\beta (xV)}}_{P_l(V)}\right)^{e^{ld}}`$ (61)
where the single site distribution $`P_l(V)`$ now is scaled with $`l`$. We introduce
$`\stackrel{~}{G}_l(x)=e^{e^{\beta (xV)}}_{P_l(V)}=\mathrm{exp}\left(e^{ld}\mathrm{ln}G_l(x)\right)`$ (62)
Let us choose the single site distribution $`P_l(V)`$ which corresponds to the REM approximation (10) of the model studied here (1,6), i.e the gaussian:
$`P_l(V)={\displaystyle \frac{1}{\sqrt{4\pi \sigma l}}}e^{\frac{V^2}{4\sigma l}}`$ (63)
It satisfies:
$`_lP_l(V)=\sigma _V^2P_l(V)`$ (64)
One easily checks that it implies that:
$`_l\stackrel{~}{G}_l(x)=\sigma _x^2\stackrel{~}{G}_l(x)`$ (65)
This leads to the equation for $`G_l(x)`$:
$`_lG=\sigma _x^2G+dG\mathrm{ln}G\sigma (1e^{ld}){\displaystyle \frac{1}{G}}(_xG)^2`$ (66)
Thus the RG equation of the REM, for large $`l`$ reads:
$`_lG=\sigma _x^2G+dG\mathrm{ln}G\sigma {\displaystyle \frac{1}{G}}(_xG)^2`$ (67)
and is almost a KPP equation, except that it has an additional gradient (KPZ type) term. This term here plays an important role and yields a different universality class than KPP. We now search for the front solutions.
Let us rewrite the exact equation (66) using the function $`h=\mathrm{ln}G`$ (remember that $`0<G<1`$):
$`_lh=dh+\sigma h^{\prime \prime }+\sigma e^{ld}h_{}^{}{}_{}{}^{2}`$ (68)
For large $`l`$ we can neglect the decaying nonlinear part, and we now look for a solution of the linear equation. The only front solution of the form $`h(x)=\stackrel{~}{h}(x+m(l))`$ with $`_lm(l)c`$ which satisfies the boundary conditions $`h(\mathrm{})=0`$ and $`h(+\mathrm{})=+\mathrm{}`$ is the exponential:
$`h_l(x)=e^{\alpha (x+m(l))}`$ (69)
$`_lm(l)=c={\displaystyle \frac{d}{\alpha }}+\sigma \alpha `$ (70)
By using again the $`h_l(x)e^{\beta x}`$ boundary condition at $`x\mathrm{}`$, we find $`\alpha =\beta `$ and:
$`c(\beta )={\displaystyle \frac{d}{\beta }}+\sigma \beta `$ (71)
as in (43). This is correct in the high $`T`$ phase and yields the correct REM value for the intensive free energy $`f(\beta )=c(\beta )/d+O(1/\mathrm{ln}L)`$ as in (11b) (and also correctly yields the absence of non trivial finite size corrections). Thus for the REM in the high $`T`$ phase we find:
$`G_l(x)\mathrm{exp}(e^{\beta (x+m(l))})`$ (72)
thus again a Gumbell form, with $`\alpha _l=\beta `$ and $`m(l)=(\frac{d}{\beta }+\sigma \beta )l`$.
To see the transition to a low T phase for $`\beta \beta _c=\sqrt{d/\sigma }`$ and the freezing of the velocity at $`c=c^{}=2\sqrt{d\sigma }`$, one needs to carry a slightly more detailed analysis (discarding again the decaying nonlinear part). The general solution of the linear part of the equation (68) is:
$`h_l(x)={\displaystyle 𝑑x^{}\frac{1}{\sqrt{4\pi \sigma l}}e^{ld\frac{(xx^{})^2}{4\sigma l}}h_0(x^{})}`$ (73)
where $`h_0(x^{})`$ can be interpreted as the $`h_l(x^{})`$ at earlier time $`l_0`$ such that the nonlinear terms can already be neglected and decays as $`h_0(x^{})e^{\beta x^{}}`$ for $`x^{}\mathrm{}`$.
This formula nicely exhibits the REM transition. In the high $`T`$ phase, using the asymptotic form $`h_0(x^{})e^{\beta x^{}}`$ we find that there is a saddle point at $`x^{}=x+2\sigma \beta l`$. This gives $`h_l(x)e^{\beta (x+c(\beta )l)}`$ with $`c(\beta )`$ given in (71). The front $`h_l(x)`$ is centered at $`x^{}=c(\beta )l`$ and consistency requires that the corresponding saddle point $`x_{}^{}{}_{}{}^{}`$ moves to $`\mathrm{}`$ so that the asymptotic form of $`h_0(x^{})`$ can indeed be used. Hence we have $`x_{}^{}{}_{}{}^{}(\sigma \beta \frac{d}{\beta })l`$. Thus the saddle point become inconsistent and the high $`T`$ solution ceases to hold, for $`\beta \beta _c=\sqrt{d/\sigma }`$.
The solution in the low $`T`$ phase is easy to find. Setting $`x=m(l)+y`$ one finds for large $`l`$:
$`h_l(y)e^{ld\frac{1}{4\sigma l}m(l)^2\frac{1}{2}\mathrm{ln}(4\pi \sigma l)}e^{\frac{c^{}}{2\sigma }y}{\displaystyle 𝑑x^{}e^{\frac{c^{}}{2\sigma }x^{}}h_0(x^{})}`$ (74)
where we have denoted $`c^{}=lim_{l+\mathrm{}}m(l)/l`$ and neglected the additional factor $`e^{x_{}^{}{}_{}{}^{2}/(4\sigma l)}`$ in the integral. This is correct provided the integral:
$`{\displaystyle 𝑑x^{}e^{\frac{c^{}}{2\sigma }x^{}}h_0(x^{})}`$ (75)
is convergent, i.e $`c^{}<2\beta \sigma `$. The consistent choice for $`c^{}`$ and $`m(l)`$ must be:
$`c^{}=2\sqrt{\sigma d}m(l)=\sqrt{{\displaystyle \frac{\sigma }{d}}}\left(2ld{\displaystyle \frac{1}{2}}\mathrm{ln}(4\pi \sigma l)\right)+O(1)`$ (76)
which ensures that (74) has a proper limit $`h_l(y)Ae^{\frac{c^{}}{2\sigma }y}=Ae^{\beta _cy}`$ which is again a Gumbell form for $`G_l(x)`$ but now is temperature independent. This holds for $`\beta \beta _c=\sqrt{d/\sigma }`$.
From this method of solving the REM we have recovered the result of namely that for $`\beta \beta _c`$ the free energy behaves as:
$`f_L(\beta )={\displaystyle \frac{m(l)}{dl}}={\displaystyle \frac{1}{\beta _c}}\left(2{\displaystyle \frac{1}{2}}{\displaystyle \frac{\mathrm{ln}(\mathrm{ln}L)}{d\mathrm{ln}L}}\right)+O\left({\displaystyle \frac{1}{\mathrm{ln}L}}\right)`$ (77)
In addition we recover, for $`T=0`$ the result for the minimum $`V_{min}`$ in the REM approximation:
$`V_{min}=2\sqrt{\sigma d}\mathrm{ln}L+{\displaystyle \frac{1}{2}}\sqrt{{\displaystyle \frac{\sigma }{d}}}\mathrm{ln}(\mathrm{ln}L)+\delta V`$ (78)
with $`u=\delta V<\delta V>`$ distributed with a Gumbell distribution:
$`\text{Proba}(u>x)=\mathrm{exp}(Ae^{\beta _cx})`$ (79)
where $`A`$ is a constant.
#### 3 conclusion on RG fronts and extremal statistics
Thus we have seen on two examples that extremal statistics problems (and their $`T>0`$ thermodynamic model counterpart) can be studied using non linear RG equation with traveling wave solutions. In one example (uncorrelated rescaled variables, i.e the REM) the RG equation is exact, while in the second (logarithmically correlated variables) we only know it presumably in the tails. The front position represents the typical value of the minimum $`V_{min}`$ as a function of $`l=d^1\mathrm{ln}N`$ while the shape of the front gives the distribution of the $`V_{min}`$ (resp. of the free energy $`F`$). This suggests that a broader class of such models can be approached by these methods, and raises the question of universality.
Studies of such non linear equations usually distinguish between pushed fronts where the velocity relaxes exponentially in $`l`$ (velocity selection by non linear terms) and pulled fronts (velocity selection by the marginal stability criterion). The extremal statistics (and the glassy phase) correspond to the pulled fronts. There one expects a very broad universality as stressed in : not only is the asymptotic front universal but also the velocity and its corrections. In a nutshell, the argument for the universal $`\frac{3}{2}\mathrm{ln}l`$ corrections to the front position comes from matching of the universal tail of the front $`g(y)(Ay+B)e^{\beta _c\xi }`$ with $`y=x+m(l)`$ with the far tail region, so far ahead of the front that one can linearize the KPP equation and get:
$`1G_l(x)e^{\beta _cy}\psi (y)_l\psi _l(y)=\sigma _y^2\psi _l(y)`$ (80)
The only matching solution is $`\psi _l(y)=y/l^{3/2}e^{y^2/(4\sigma l)}`$. Inserting $`y=x+m(l)=x+c^{}l+C\mathrm{ln}l`$ immediately yields $`C=3/2`$ for proper matching. As discussed in this universality extend for pulled fronts in a very broad class of non linear (or coupled non linear) equations and holds for steep enough initial conditions (i.e in the glass phase in our language).
This argument fails in some cases, such as at the bifurcation between pushed and pulled fronts (e.g. at the glass transition $`\beta =\beta _c`$ or equivalently when the initial condition has slow decay $`\mathrm{exp}(\beta _cx)`$) (see e.g analysis in ). Interestingly, it clearly fails also for the non linear equation corresponding to the REM model, which is thus in a different universality class (this may be related to the fact that fronts are here unbounded ). Presumably what happens there is that the coefficient $`A`$ vanishes, and the solution is exactly $`e^{\beta _cy}`$, hence the $`\frac{1}{2}\mathrm{ln}l`$ (since the above matching function is now $`\psi (y)l^{1/2}e^{y^2/(4\sigma l)}`$).
Next is the question of universality. We will address it only for our model of gaussian variables with logarithmic correlations. We have recast the RG equation (23) into a KPP equation with a specific non linear term $`F[G]`$. From our RG we have obtained $`F[G]=G(1G)`$. The structure of the RG derivation suggests that we have obtained correctly the two lowest orders of $`F[G]`$. From the above discussion this is enough for the universality. Thus, and we call it the restricted universality scenario, it is likely that higher order terms $`F[G]=G(1G)+O((1G)^3)`$ are non universal and thus that only the tail of the distribution of the minimum of log-correlated variables is universal.
Let us mention however that we were not able to rule out another scenario, the broad universality scenario, such that the true distribution of the minimum of log-correlated variables is indeed universal. If this was true the following conjecture would be tempting: since we know that for uncorrelated variables the KPP RG equation is exact with $`F[G]=G\mathrm{ln}G`$ and $`\sigma =0`$ (and see Appendix B is asymptotically exact even for weakly correlated ones), one could conjecture an interpolating KPP equation (38) with $`F[G]=G\mathrm{ln}G`$ and $`\sigma >0`$ which would gives exactly the distribution of the minimum of log-correlated variables. Unfortunately we have been unable to confirm (but also to strictly rule out) numerically this conjecture, due to the very large finite size corrections, as discussed in Section IV.
### E structure of low temperature phase and replica symmetry breaking
Let us now return to the structure of the low temperature phase for the particle in the $`d`$-dimensional random potential with logarithmic correlations. We argue that (i) it has a non trivial structure, with a few states (ii) this structure is reminiscent of the so-called “replica symmetry breaking” This non trivial structure can be caracterised more precisely here as the various states of the model correspond to the different positions of the particle, and have thus a natural meaning in real space. In particular, the minima of the ”energy landscape” (or metastable states) are nothing but the local minima in the sample of the random potential for our problem. A precise caracterisation of these ’local minima’ is given below. Also, approximate replica solutions of our model are shown in the following to exhibit RSB at low T.
#### 1 Spatial distribution of secondary minima
Let us start with a simple argument: for a given realization of disorder, we divide our system into two subsystems of size $`L^d/2`$, and call $`V_{min1}`$ and $`V_{min2}`$ the two corresponding minima in each subsystem.
Within the REM approximation, we know from (13) that $`V_{min1}V_{min2}(y_1y_2)\sqrt{\sigma /d}O(1)`$ where $`y_1`$ and $`y_2`$ have independent Gumbell distributions. Thus clearly in that case there is a non trivial structure: the secondary minimum (defined as being constrained to lie within the other subsystem) is typically within $`\mathrm{\Delta }E=O(1)`$ in energy of the absolute minimum (and within this approximation the distribution is also easily computed).
The RG analysis performed in this paper indicates that adding correlations will not change this conclusion. Indeed, one first coarse grains up to scale $`l_0=\mathrm{ln}(L)\frac{1}{d}\mathrm{ln}2`$. At this scale, the system can be described by two local energies (one for each half) of minima $`v_1`$ and $`v_2`$ distributed according to $`P_{l_0}(v)`$, to which should be added a term $`\delta V`$ which correlates the two halfes and is gaussian of variance $`\frac{2\sigma }{d}\mathrm{ln}2`$. This however does not change the fact that the difference $`V_{min1}V_{min2}O(1)`$. Thus one still finds that there exist secondary minima of $`O(1)`$ in energy from the minimum, and a typical distance $`L`$ away from the absolute minimum. As discussed in Section II D this property was also confirmed by numerical simulations.
It is natural, in view of the analogy with the directed polymer on the Cayley tree, to introduce the ”overlap” between two different states (i.e positions of particles) $`𝐫_1`$ and $`𝐫_2`$ as:
$`q(𝐫_1,𝐫_2)=1{\displaystyle \frac{\mathrm{ln}(a+|𝐫_1𝐫_2|)}{\mathrm{ln}L}}`$ (81)
We expect it to be non self averaging and characterized by the “overlap distribution”:
$`P_2(q)=\overline{{\displaystyle \underset{𝐫_1,𝐫_2}{}}p(𝐫_1)p(𝐫_2)\delta (qq(𝐫_1,𝐫_2))}`$ (82)
Although we have not attempted to compute this function directly using our RG it is natural to expect that, as in the REM and the DPCT, it is non trivial for $`T<T_c`$ and reads:
$`P_2(q)={\displaystyle \frac{T}{T_c}}\delta (q)+(1{\displaystyle \frac{T}{T_c}})\delta (1q)`$ (83)
Similarly one expects that in a given disorder environment, the probability of finding an overlap $`q`$ between two thermal realizations becomes in the large $`L`$ limit:
$`\stackrel{~}{Y}(q)dq=(1Y)\delta (q)+Y\delta (q1)`$ (84)
with $`\overline{Y}=1\frac{T}{T_c}`$ and $`Y`$ has the same distribution as in the REM. Thus the natural expectation, from the DPCT analogy, is that the overlap in the low $`T`$ phase will be either $`1`$ or $`0`$ (i.e secondary minima - of energy difference of order $`T`$ \- will be either near the absolute one $`\mathrm{ln}r_{12}/\mathrm{ln}L0`$, or a distance $`r_{12}O(L)`$ typically a fraction of the system size away) It would however be of interest to investigate further these properties in the present model, in particular to obtain more detailed information at intermediate scales, e.g. correlations probing the whole range $`\mathrm{ln}r_{12}(\mathrm{ln}L)^a`$ with $`0a1`$.
#### 2 Approximate replica symetry breaking solutions of the model
Let us now turn to the replica representation and discuss how the present model exhibits a form of “replica symmetry breaking”. The replicated partition sum reads:
$`\overline{Z^m}={\displaystyle \frac{d^d𝐫_1}{a^d}\mathrm{}\frac{d^d𝐫_m}{a^d}e^{2\sigma \beta ^2_{i<j}\mathrm{ln}\frac{|𝐫_i𝐫_j|}{a}}e^{m^2\sigma \beta ^2\mathrm{ln}\frac{L}{a}}}`$ (85)
It turns out that various approximations of this partition function (specifically the REM and the DPCT approximations) are dominated, in the limit $`m0`$ by replica symmetry breaking configurations.
In the context of 2d Dirac fermions with random vector potential (see Section VI) an estimate of (85) was given in . For small $`\beta `$ it is clear that the exponential containing the logarithmic attraction between replicas does not decay fast enough and thus the integral is dominated by the configurations where the replicas are all far away $`O(L)`$ apart thus:
$`\overline{Z^m}\left({\displaystyle \frac{L}{a}}\right)^{\beta ^2\sigma m^2+dm\beta ^2\sigma m(m1)}=\left({\displaystyle \frac{L}{a}}\right)^{md(1+\frac{\sigma }{d}\beta ^2)}`$ (86)
This estimate of Ref. is in fact incorrect as it misses the glass transition. Indeed, one can redo this argument using configurations where $`m/p`$ packets of $`p`$ replicas are $`O(L)`$ far apart (while in each packet the replica (independent particles) are close to each other). This estimate was performed in Ref. and gives instead:
$`\overline{Z^m}\left({\displaystyle \frac{L}{a}}\right)^{\beta ^2\sigma m^2+d\frac{m}{p}\beta ^2\sigma m(mp)}`$ (87)
The interaction term being proportional to the number of pairs of replicas in different packets, which is $`m(mp)/2`$. In the limit $`m0`$ one can then optimise over $`0<u=p<1`$, i.e:
$`\overline{Z^m}\mathrm{exp}\left(d\mathrm{ln}{\displaystyle \frac{L}{a}}\underset{0<u<1}{\mathrm{max}}\left({\displaystyle \frac{1}{u}}+{\displaystyle \frac{\sigma }{d}}\beta ^2u\right)\right)`$ (88)
For $`\beta <\beta _c=\sqrt{\frac{d}{\sigma }}`$ the saddle is for $`p=1`$ and one recovers the above expression. For $`\beta >\beta _c=\frac{d}{\sigma }`$ one finds that the saddle is for $`u=\beta _c/\beta =T/T_c`$ which gives:
$`\overline{Z^m}e^{d\mathrm{ln}\frac{L}{a}2\frac{\beta }{\beta _c}}`$ (89)
Thus this calculation yields a transition. In Ref. it was claimed that it does not correspond to replica symmetry breaking. We believe that this is incorrect and that this is a (one step) RSB estimate of the above partition sum. This is clear since this calculation exactly amounts to the corresponding one for the REM approximation of the model, i.e replacing in (86) $`_{i<j}\mathrm{ln}\frac{|𝐫_i𝐫_j|}{a}`$ by $`_{i<j}(1\delta _{r_i,r_j})\mathrm{ln}(L/a)`$. In the REM we know from Ref. that the correct solution for $`T<T_c`$ can be obtained by performing the analytical continuation to $`m0`$ on a RSB saddle point (note that the REM finite size correction $`\frac{1}{2}\mathrm{ln}\mathrm{ln}L`$ is also obtained from the saddle integration).
One can go one step further and use an argument based on universality, which puts the present problem in the DPCT universality class (for some observables such as the free energy distribution). For the DPCT, it was shown in Ref. that one can also recover the correct result for the averaged free energy by considering directed polymer configurations which break replica symmetry as $`m0`$. It remains to be demonstrated how to obtain other universal quantities, e.g. the $`\frac{3}{2}\mathrm{ln}\mathrm{ln}L`$ finite size corrections, via a RSB saddle point calculation.
It is interesting to see how the features associated to RSB arise from the RG developed here, despite the fact that it is explicitly replica symmetric. Quite generally, if one can find independent local free energy variables with an exponential distribution $`P(f)e^{\beta _cf}`$ one naturally obtains a RSB picture. This is the case here, up to some more detailed universal preexponential structure in $`P(f)`$. The important feature of our RG is thus that it follows the full distribution $`P_l(z)`$ of local disorder (i.e. of local Boltzman weights $`z`$) which becomes algebraically broad as $`l+\mathrm{}`$. Here this property is sufficient to show that the low $`T`$ phase has a structure reminiscent of RSB. Indeed, let us again coarse grain the system up to an already large scale $`L_0=ae^{l_0}`$ but still much smaller than $`L`$, the ratio $`L/L_0=e^{l_1}=M`$ being large but fixed as $`L+\mathrm{}`$. Assuming that $`L_0`$ is so large that $`P_{l_0}(z)`$ has reached its fixed point already (except in a remote tail region corresponding to very rare events). Since one has the decomposition (17) the RG tells us that the sample is divided in $`M`$ subsystems with free energies $`F_i=v_i+V_i^>`$, $`i=1,..M`$ where the variables $`z_{i=1,..M}=e^{\beta v_i}`$ are independently drawn from the common distribution $`P_{l_0}(z)`$ and the $`V_i^>`$ are still correlated but gaussian. Neglecting first the $`V_i^>`$ we are left with a system of $`M`$ subsystems of Gibbs measure:
$`{\displaystyle \frac{z_i}{_jz_j}}`$ (90)
Since the $`z_i`$ are drawn from a distribution with algebraic tails $`P(z)1/z^{1+\mu }`$ with $`\mu =T/T_c`$ one has $`<z>=+\mathrm{}`$ for $`T<T_c`$ and, as is well known, the partition sum (90) is dominated by a few of the $`z_i`$ variables (which in essence is the physics associated to RSB). Since the correlated $`V^>`$ variables are in finite number and with gaussian tails they cannot change the exponential tails of the $`F_i`$ and thus adding them back should not change the above conclusions.
Thus here, although the RG is replica symmetric, since it allows for generation of broad tails it can capture features usually associated with RSB.
## IV numerical study
Since we found via the RG and other arguments that there should be a transition in any dimension $`d1`$ it is particularly convenient to perform numerical simulations in the ”extreme case” of $`d=1`$ (i.e the further away from mean field). However, even in $`d=1`$ numerical simulations are delicate because the finite size corrections are very large (and interesting to study, in order to distinguish various universality class). Indeed we have found that the main numerical uncertainties come from the finite size effects and not from the number of averages. In most of the numerical work averaging over $`10^4`$ realizations of disorder was sufficient, while a simulation of a system of size $`2^{21}210^6`$ leads to important corrections to the thermodynamic behaviour of the model. In view of this, we believe that the previous numerical investigation was at best approximate.
We have considered a lattice model in $`d=1`$ with $`L=2^n`$ sites. The potential $`V(r)`$ on each site ($`r=1,..L`$) was computed from its Fourier components $`V(r)=w_{L/2}(1)^r+_{k=1}^{L/21}w_k\mathrm{cos}(2\pi kr/L\varphi _k)`$, eliminating the $`k=0`$ mode, with $`w_k`$ independent gaussian variables $`\overline{w_kw_k^{}}=\mathrm{\Delta }(k)\delta _{k,k^{}}`$ ($`k,k^{}=1,..L/2`$) and each $`\varphi _k`$ independently distributed uniformly in $`[0,2\pi ]`$. We choose $`\mathrm{\Delta }(k)`$ such that:
$`\mathrm{\Gamma }(rr^{})`$ $`=`$ $`\overline{V(r)V(r^{})}`$ (91)
$`=`$ $`\sigma {\displaystyle \frac{2\pi }{L}}{\displaystyle \underset{k=1}{\overset{L1}{}}}{\displaystyle \frac{\mathrm{cos}(\frac{2\pi k}{L}(rr^{}))}{|\mathrm{sin}(\frac{\pi k}{L})|\sqrt{62\mathrm{cos}(\frac{2\pi k}{L})}}}`$ (92)
so that $`\overline{(V(r)V(r^{}))^2}=4\sigma \mathrm{ln}(rr^{})`$ for $`1rr^{}L/2`$. This is the choice which also corresponds to correlations along the axis $`y=0`$ on a 2d square lattice.
The behaviour of the model has been studied, without loss of generality, at zero and at finite temperature for a disorder strength $`\sigma =1`$ (others value of $`\sigma `$ can be incorporated in the definition of the temperature scale). We have first computed the average minimum $`e_{min}=\overline{V_{min}}/\mathrm{ln}N`$ (with $`N=L`$) for system sizes ranging from $`L=2^7=128`$ to $`L=2^{21}10^6`$ and for each size we have taken the average over $`10^4`$ realizations of disorder. An estimate of the uncertainty on the disorder average was made by measuring the variance of a series of average over $`10^4`$ realizations. This variance was found to be of the order of $`10^3`$ for all the value of $`\overline{V_{min}}`$. The results are plotted in Fig. 10. We recall that the RG prediction reads for $`\sigma =1`$:
$`{\displaystyle \frac{1}{\mathrm{ln}N}}\overline{V_{min}}=2\mathrm{ln}N{\displaystyle \frac{3}{2}}\mathrm{ln}(\mathrm{ln}N)+O(1)`$ (93)
We should first note that if one does not assume anything about the finite size corrections, the resulting uncertainty on the ratio $`e_{min}=\overline{V_{min}}/\mathrm{ln}N`$ is very large even for sizes $`L=2^{21}`$ since the ratio $`\frac{3}{2}\mathrm{ln}(\mathrm{ln}N)/\mathrm{ln}N0.3`$. Hence with no assumption it is hard to estimate $`e_{min}`$ to better than $`10`$ per cent accuracy.
However, if one assumes that $`e_{min}=2`$, the plot in Fig. 10 shows the existence of the $`\mathrm{ln}(\mathrm{ln}N)`$ corrections with a slope definitely larger than $`1`$ and consistent with $`3/2`$ (although the accuracy is not excellent). It is however sufficient to rule out a REM type behaviour and is consistent with the RG prediction (93).
Next, we have plotted the distribution of $`V_{min}`$ in Fig. 11 and compared with the prediction of the RG for the tails. Here also the agreement is satisfactory.
Finally, we have plotted the “glass order parameter” $`Y_2=\overline{_rp_r^2}`$ which is non zero when the system is dominated by a few states. It is consistent with a very slow convergence towards $`Y_2=(1T/T_c)\theta (T_cT)`$ but clearly other forms cannot be ruled out.
## V Relations with Liouville and Sinhgordon models
In this Section we describe the relation between the problem of the particle in the log-correlated random potential and the Liouville and sinh-Gordon models. Exact results on the sinh-gordon model are compatible with (and also point out towards) the existence of the transition at $`\beta =\beta _c`$.
### A Relations with the sinh-Gordon model in $`d=2`$ and $`d=1`$
Let us start with the correspondence with the sinh-Gordon model. Although less direct, it is also simpler to analyze, as the model does not contain subtle boundary conditions problems. The interest of the connexion is that the sinh-Gordon model is integrable in $`d=2`$ and $`d=1`$ (Boundary sinh-Gordon) .
The connexion requires introducing a slightly different version of the initial problem, defined by the partition function:
$`Z_{sh}[V]=Z[V]+Z[V]={\displaystyle \underset{𝐫}{}}(e^{\beta V(𝐫)}+e^{\beta V(𝐫)})`$ (94)
which corresponds to a particle in a random potential which can explore both $`V(𝐫)`$ and $`V(𝐫)`$. A physical realization would be a particle with an Ising spin in a random field. As it turns out the physics of this disordered model is very similar to the original problem. At low temperature, it is now related to the distribution of the minimum of $`|V(𝐫)|`$.
We define the generating function of this model $`G_{sh}(x)=\mathrm{exp}(\mu Z_{sh}[V])`$, with $`\mu =e^{\beta x}`$, which is related to the distribution of the free energy of the particle. In the continuum limit and in $`d=2`$, it can be rewritten as:
$`G_{sh}(x)=H_{sh}[\mu ]={\displaystyle DVe^{S_{sh}[V]}}`$ (95)
$`S_{sh}[V]={\displaystyle _{\mathrm{}}^+\mathrm{}}{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑x𝑑y\left({\displaystyle \frac{1}{8\pi \sigma }}(V)^2+2\mu \mathrm{cosh}(\beta V(𝐫))\right)`$ (96)
i.e the partition function of the sinh-Gordon model in $`d=2`$. Similarly, the $`d=1`$ version of our model is related to the well studied boundary sinh-Gordon model defined as:
$`G_{shB}(x)=H_{shB}[\mu ]={\displaystyle DVe^{S_{shB}[V]}}`$ (97)
$`S_{shB}[V]={\displaystyle _0^+\mathrm{}}dy{\displaystyle _0^L}dx({\displaystyle \frac{1}{4\pi \sigma }}(V)^2`$ (98)
$`+2\mu \mathrm{cosh}(\beta V(x,0)))`$ (99)
Indeed one has, as required, that $`(V(x,0)V(x^{},0))^24\sigma \mathrm{ln}\frac{|xx^{}|}{a}`$ at large $`|xx^{}|`$, and one only studies (boundary) observables defined at $`y=0`$.
In the limit of $`\beta =+\mathrm{}`$ one has in both case:
$`G_{sh}(x)`$ $`=`$ $`\text{Proba}(x<\text{min}(V_r,V_r))`$ (100)
$`=`$ $`\text{Proba}(x<\text{max}_r|V_r|)`$ (101)
and thus the (properly discretized) partition function of the (boundary) sinh-Gordon model becomes related, in that limit, to the distribution function of the maximum of the set of positive random variables $`|V(𝐫)|`$. The results described in the previous Sections about the statistics of extrema of such variables imply that some transition must occur as a function of $`\beta `$ corresponding to a related “change of behaviour” in the sinh-Gordon and boundary sinh-Gordon models as well. This is a prediction, as we are not aware of such a change of behaviour at $`\beta =\beta _c`$ being mentionned in the literature. As we now discuss, examination of known results is perfectly compatible with the transition at $`\beta =\beta _c`$.
Let us first describe the known exact results both in $`d=2`$ and $`d=1`$. The extensive free energy of the bulk sinh-Gordon model is defined as:
$`f_{sh}=\underset{L+\mathrm{}}{lim}L^2\mathrm{ln}G_{sh}`$ (102)
where the model defined in (96) is considered in finite size $`L`$. The model is studied usually using the field $`\varphi =V\sqrt{2/\sigma }`$, the nonlinear term being $`2\mu \mathrm{cosh}(\beta V)=2\mu \mathrm{cosh}(b\varphi )`$ and its free energy depends on the single variable $`b=\beta \sqrt{\sigma /2}=\beta /\beta _c`$, where $`\beta _c=\sqrt{d/\sigma }`$ is dimension dependent. Using the variable $`b`$, its exact expression, proposed in Ref. , reads when explicited :
$`f_{sh}(\mu )=C_2(b)\mu ^{\frac{1}{1+b^2}}`$ (103a)
$`C_2(b)={\displaystyle \frac{2\pi }{\left(\mathrm{\Gamma }[\frac{1}{2+2b^2}]\right)^2\left(\mathrm{\Gamma }[1+\frac{b^2}{2+2b^2}]\right)^2\mathrm{sin}(\frac{b^2\pi }{1+b^2})}}`$ (103b)
$`\times \left({\displaystyle \frac{\pi \mathrm{\Gamma }[1+b^2]}{\mathrm{\Gamma }[b^2]}}\right)^{\frac{1}{1+b^2}}`$ (103c)
These results are a priori only valid for $`b<1`$ ($`|b|<1`$), as they were obtained in from an analytical continuation of the sine-Gordon model (performing $`\mu \mu `$ and $`b^2b^2`$, $`M`$ being the soliton mass). The constant $`\mu `$ was defined in the continuum model by fixing the normalization of the field $`<\mathrm{cos}(b\varphi (𝐫))\mathrm{cos}(b\varphi 𝐫^{})>=\frac{1}{2}|xy|^{4b^2}`$ of the sine-Gordon model.
The $`d=1`$ version corresponds to the boundary sinh-Gordon model usually studied using the $`\varphi =V/\sqrt{\sigma }`$ and $`2\mu \mathrm{cosh}(\beta V)=2\mu \mathrm{cosh}(b\varphi )`$, with again $`b=\beta /\beta _c`$ ($`\beta _c=1/\sqrt{\sigma }`$). The analogous expression for the free energy reads, from :
$`f_{shB}(\mu )=\underset{L+\mathrm{}}{lim}L^1\mathrm{ln}G_{shB}=C_1(b)\mu ^{\frac{1}{1+b^2}}`$ (104)
$`C_1(b)={\displaystyle \frac{1}{8\pi ^{3/2}}}\mathrm{\Gamma }\left[{\displaystyle \frac{1+2b^2}{2+2b^2}}\right]\mathrm{\Gamma }\left[{\displaystyle \frac{b^2}{2+2b^2}}\right]\left({\displaystyle \frac{2\pi }{\mathrm{\Gamma }[b^2]}}\right)^{\frac{1}{1+b^2}}`$ (105)
Let us now comment on these results. The power law dependence in $`\mu `$ of the free energy is just the naive dimensional result $`\mu ^{\frac{1}{1+b^2}}`$ in both cases. This result should hold for $`\beta <\beta _c`$. However, there is clearly, in both $`d=2`$ and $`d=1`$ cases, a singularity as $`\beta \beta _c^{}`$ as the amplitude $`C(b)`$ diverges as $`b=\beta /\beta _c1^{}`$. This is thus in perfect agreement with the existence of a phase transition in the particle model. In the sinh-Gordon model itself, we do not expect strictly speaking a phase transition, as the model is massive both below and above $`b=1`$, however we do expect some “change of behaviour”, which may be related to a change of nature of the excitations around the ground state. This is not ruled out by exact results as it clearly comes here from the physical mass acquiring a nontrivial dependence in the bare mass parameter $`\mu `$ (contrarily to sine-Gordon model, for sinh-Gordon model there is no presently known exact solution of a lattice version).
Let us now interpret these results for our model. They mean that the generating function $`G_{sh}(x)`$ of the free energy distribution, with $`\mu =e^{\beta x}`$, takes indeed the form of a traveling wave:
$`G_{sh}\mathrm{exp}\left(L^2C_d\left({\displaystyle \frac{\beta }{\beta _c}}\right)\mu ^{\frac{2}{(1+\frac{\beta }{\beta _c})^2}}\right)=g(x+cl+\gamma )`$ (106)
with $`l=\mathrm{ln}L`$ and a velocity:
$`c={\displaystyle \frac{d}{\beta }}+\sigma \beta `$ (107)
This is exactly the velocity given by the KPP equation for the particle model, in the high temperature phase. It also yields a front $`g(y)=\mathrm{exp}(e^{ay})`$ with $`a=\frac{\beta }{1+(\beta /\beta _c)^2}`$ and $`\gamma =\frac{\beta }{1+(\beta /\beta _c)^2}\mathrm{ln}C_d(\beta /\beta _c)`$. This form however should be taken with caution as strictly speaking formula (106) is valid only in the limit where $`L`$ goes to infinity first (at fixed $`\mu =e^{\beta x}`$). It should be compared with the asymptotic behaviour of the front in the region of large positive $`y`$. We expect universality in the other region of the front (of very negative $`y`$ i.e. $`x<<cl`$) and exact knowledge about this region would be equivalent to exact knowledge of the sinh-Gordon model at finite size, which is not yet available.
The physics of the problem of the particle in the random potential leads us to conjecture that the 2d sinh-Gordon model (as well as the boundary sinh-Gordon model) will exhibit a change of behaviour, the algebraic $`\mu `$-dependance of its free energy will freeze for $`\beta \beta _c`$, which corresponds to the low temperature glassy phase of the particle model. We thus expect:
$`f_{sh}(\mu )\mu ^\alpha `$ (108a)
$`\alpha ={\displaystyle \frac{1}{1+(\beta /\beta _c)^2}}`$ $`\beta <\beta _c=\sqrt{d/\sigma }`$ (108b)
$`\alpha ={\displaystyle \frac{1}{2}}`$ $`\beta >\beta _c=\sqrt{d/\sigma }`$ (108c)
and presumably log corrections (at least at $`\beta =\beta _c`$, and maybe for all $`\beta >\beta _c`$).
This is confirmed by a renormalization group analysis directly on the Sinh-Gordon and Liouville models discussed below.
### B Relation with the Liouville model in $`d=2`$
The relation between our original model (1) of the particle in the random potential and the Liouville model proceeds via the generating function:
$`G(x)`$ $`=`$ $`\mathrm{exp}(e^{\beta x}Z[V])_V`$ (109)
$`=`$ $`\mathrm{exp}({\displaystyle \underset{𝐫}{}}e^{\beta (xV(𝐫))})_V`$ (110)
which encodes the full probability distribution of the free energy of the particle. In the case of the $`d=2`$ potential with logarithmic correlations it is identical to the partition function of a Liouville model, which one can write either on the original lattice or in the continuum (with UV and IR cutoffs $`a`$ and $`L`$) as ($`\mu =e^{\beta x}`$) :
$`G(x)=H[\mu ]={\displaystyle DVe^{S[V]}}`$ (111)
$`S[V]={\displaystyle d^2𝐫\left(\frac{1}{8\pi \sigma }(V(𝐫))^2+\mu e^{\beta V(𝐫)}\right)}`$ (112)
where the functional integral is normalized such that $`H[\mu =0]=1`$ (equivalently one redefines $`H[\mu ]H[\mu ]/H[\mu =0]`$). We call it the Liouville Model (LM) to distinguish it from the continuum Liouville field theory LFT whose definition is recalled below. A relation also exists between the correlation functions of the Gibbs measure and some correlation functions in the Liouville model:
$`<p(𝐫_1)..p(𝐫_n)>={\displaystyle _{\mu >0}}\mu ^{n1}e^{\beta (V(𝐫_1)+\mathrm{}V(𝐫_n))}e^{S[V]}`$ (113)
Strictly speaking the model (112) above is not well defined because of the zero mode $`V(𝐫)V(𝐫)+w`$ and must be complemented with boundary conditions. In the particle problem studied here we have chosen periodic boundary conditions with the additional constraint $`_𝐫V(𝐫)=0`$ to pin the zero mode.
On the other hand, many results are known for the (related) continuum Liouville field theory (LFT), of great interest in quantum gravity \[64, 65, 66, 67, kazakov\]. It is usually defined on a arbitrary genus $`h`$ manifold with background metric $`g`$ and associated curvature $`R`$ by the action :
$`S_{LFT}={\displaystyle d^2x\left(\frac{1}{4\pi }(_a\varphi )^2+\mu e^{2b\varphi }+\frac{Q}{4\pi }R\sqrt{g}\varphi \right)}`$ (114)
in conventional notations, with $`\varphi =V/\sqrt{2\sigma }`$ and $`b=\beta /\beta _c`$. The (standard) choice is $`Q=b+1/b`$ for which the theory is critical and has local conformal invariance (with a central charge $`c_L=1+6Q^2=25+c`$) . It can also be formulated as the theory of (liquid) random surfaces e.g. as random triangulations. There one defines the total area $`A=_𝐫e^{2b\varphi (r)}`$ which is nothing but the partition function $`A=Z[V]`$ of the particle problem, and studies the distribution $`𝒵(A)e^{\mu _cA}A^{\gamma _{\text{string}}3}`$ which is nothing but $`P(Z)`$.
The particle model allows us to make precise statements on the Liouville model defined above. The LFT allows for exact calculations (e.g. of correlation functions) and in principle one could hope to translate those in the particule model (Gibbs measure correlations). The relation between the two however is rather subtle. For instance, the boundary conditions chosen in the particle problem would correspond to Liouville on a torus $`h=0`$, except that the additional pinning condition spoils it. We will thus not explore here all these intricacies but give a few general remarks, mostly about the behaviour of the Liouville model under coarse graining.
First we know that $`G_l(x)`$ satisfies a RG equation of the KPP type. Thus upon coarse graining (i.e as a function of $`l=\mathrm{ln}(L/a)`$ see Section III) the Liouville model partition function satisfies a KPP non linear RG equation. The corresponding front velocity gives the scaling of the partition function with $`\mu `$. The glass transition, with freezing of the front velocity, corresponds exactly in the Liouville model to the transition between two regimes (the so-called $`c=1`$ barrier in the LFT):
(i) weak coupling Liouville: $`b=\beta /\beta _c<1`$. In that regime there is no problem to define a continuum limit. The KPP non linear RG front solution of velocity $`c(\beta )=2(\frac{1}{\beta }+\frac{\beta }{\beta _c^2})c(b)`$ yields $`G_l(\mu )\mu ^{1/(1+b^2)}`$ and in the Liouville model the scaling dimensions are the one obtained by naive dimensional counting. e.g. :
$`\mu {\displaystyle d^2𝐫e^{\beta V(𝐫)}}=\mu e^{\mathrm{ln}Z[V]}\mu L^{\beta c(\beta )}=\mu L^{2(1+b^2)}`$ (115)
In the LM only one of the two branches (deduced by $`b1/b`$) is selected as $`l+\mathrm{}`$.
The regime $`b1`$ is the one where the continuum LFT is well defined: there the role of the $`Q`$ term in (114) is to shift the conformal dimension of the fields $`e^{2\alpha \varphi }`$ to $`\mathrm{\Delta }(\alpha )=\alpha (Q\alpha )`$ (while the naive power counting dimensions in the LM are $`\mathrm{\Delta }(\alpha )=\alpha ^2`$, see above) which renders the Liouville term $`e^{2b\varphi }`$ exactly marginal (and thus the theory critical). It was argued in that the LFT gives the correlations of the Gibbs measure, the operator $`p(𝐫)`$ (eq. (2)) corresponding to the Liouville field $`e^{2b\varphi }`$ (and is thus of conformal dimension $`\mathrm{\Delta }(b)=1`$ (i.e $`p(𝐫)L^2`$). A hint in favor of this conjecture was that the corresponding LFT conformal dimensions of the composite fields $`p(𝐫)^qe^{2qb\varphi }`$ is simply $`\mathrm{\Delta }(q)=q(1+b^2b^2q)`$ which correctly reproduces the multifractal spectrum:
$`{\displaystyle d^2𝐫p(𝐫)^q}L^{2(1\mathrm{\Delta }(q))}`$ (116)
given in and Section VI (in the weak disorder regime $`q<q_c`$). This is not a very strong test since the same multifractal spectrum can also be obtained within the LM model by considering the dimension of the normalized Gibbs measure (rather than the unnormalized one $`e^{\beta V}`$). Indeed the effect of the $`Q`$ term is to shift:
$`e^{2bq\varphi }L^{2bqQ}e^{2bq\varphi }Z[V]^qe^{2bq\varphi }`$ (117)
To convincingly establish the conjecture of the effect of the additional $`Q`$ term should be checked on the many point correlations , where it is rather more subtle, and further investigation is needed. In particular, the RG described here suggests by extension that the Gibbs measure correlations (or at least some limits of them) should also be computable within the DPCT model. This suggests a direct relation between LFT and DPCT, a check of which would be of great interest. Note also that a critical model, which mimics the effect of the $`Q`$ term (as adding an average value to $`\varphi `$ ) is studied in Appendix D.
(ii) strong coupling Liouville: $`b=\beta /\beta _c>1`$. This corresponds to the glass phase for the particle which, interestingly, has a non trivial structure. As is well known there are serious difficulties in defining the continuum LFT in that regime. Using what we know from the particle problem we can gain some idea of what happens in the Liouville theory. First let us note that since $`G(x)`$ is such that for $`\beta =+\mathrm{}`$ and fixed $`L`$ it is equal to the distribution function of the minimum of the set of $`V(𝐫)`$:
$`G(x)=\text{Proba}(x<\underset{𝐫}{\mathrm{min}}V(𝐫))`$ (118)
The (infinitely strong coupling) Liouville model can be recast as an extremal statistics problem in that limit. The partition sum $`Z[V]`$ of the particle model being dominated, for $`b>1`$ by a few regions of space where $`V(𝐫)V_{min}`$ (with little dependence in $`\beta `$, i.e the Liouville wall become a hard wall for all $`\beta <\beta _c`$ with thickness of order $`O(1)`$) we expect this spatial heterogeneity to show up in LM as well. From what we have learned in previous Sections we know that upon coarse graining the following effective Liouville model action $`S_{eff}`$ is generated:
$`G_l(x)=H_l[\mu ]={\displaystyle D\stackrel{~}{V}<e^{S[\stackrel{~}{V},z]}>_{P(z)}}`$ (119)
$`S_{eff,l}[V]={\displaystyle d^2𝐫\left(\frac{1}{8\pi \sigma }(\stackrel{~}{V}(𝐫))^2+z(𝐫)e^{\beta \stackrel{~}{V}(𝐫)}\right)}`$ (120)
i.e a new field $`z(𝐫)`$ is dynamically generated, and has short range correlations but has a broad power law distribution:
$`P(z)dzz^{1+\frac{1}{b}}`$ (121)
while $`\stackrel{~}{V}V^>`$ is the smooth field introduced in (17). For $`b<1`$ this dynamically generated local field can be averaged out without changing significantly the action (note that even for $`b<1`$ it changes properties of operators $`e^{q\beta V}`$ for $`q>q_c`$) while for $`b>1`$ it changes crucially the physics. One can define the effective Liouville potential $`U[V]`$ for the smooth field $`\stackrel{~}{V}`$ after averaging over the $`z`$ field as:
$`U_l[\stackrel{~}{V}]=\mathrm{ln}<\mathrm{exp}(ze^{\beta \stackrel{~}{V}})>_{P_l(z)}=\mathrm{ln}G_l(x=\stackrel{~}{V})`$ (122)
the bare Liouville potential being $`U[V]=\mu \mathrm{exp}(\beta V)`$. We can now use the front solution of the KPP equation (i.e the scaling region in the large $`L/a`$ limit) described in previous Sections. For $`b<1`$, since $`<z>_{P(z)}<+\mathrm{}`$ we have that for large $`V`$
$`U_l[V]c\mu e^{2(1+b^2)l}\mathrm{exp}(\beta V)`$ (123)
and thus the coarse grained potential is similar to the bare one. However, for $`b>1`$ one has for large $`V`$
$`U_l[V]c\mu e^{4l}V\mathrm{exp}(\beta V)`$ (124)
because of the broad distribution of the $`z`$ field.
Since the $`z(𝐫)`$ are highly heterogeneous on short scales it is not surprising that a continuum limit is hard to obtain for $`b>1`$. These heterogeneities are linked to the structure of the glass phase reminiscent of replica symmetry breaking. It is tempting to conjecture that it may also be related to the branched polymer structure which appear in LFT for $`b>1`$, i.e beyond the $`c=1`$ barrier , or to the spike instability of fluid membranes.
Furthermore, let us notice that the LFT theory at $`b=1`$ is known to have two marginal operators whose dimensions are degenerate $`e^{2b\varphi }`$ and $`\varphi e^{2b\varphi }`$. This is in exact parallel with the behaviour of the KPP front solution, which develops at $`b=1`$ two degenerate linear eigenmodes $`\mathrm{exp}(\beta V)`$ and $`V\mathrm{exp}(\beta V)`$.
Thus we have seen that the Coulomb gas RG can be used to understand the behaviour of the Liouville model. A scenario is obtained where for $`b1`$ new short scale degree of freedom are generated (short scale instability). Averaging over these changes the effective Liouville potential. The parallel with the particle model suggests that the short scale instability in Liouville may be related to the generation of strong inhomogeneities in the Gibbs measure $`p(𝐫)`$, analogous to structures discussed in the context of replica symetry breaking. Thus, if the mapping onto the LFT is confirmed it suggests to also investigate RSB type effects in strong coupling LFT.
### C Direct renormalisation group analysis of sinh-Gordon and Liouville models and traveling waves
Let us now illustrate how one can see explicitly the freezing of the free energy exponent in the strong coupling phase from renormalisation group approaches directly on the sinh-Gordon and Liouville models. Such functional RG methods have been applied to study the analogous problem of the wetting of an interface of height $`V`$. Related exact RG methods, with various truncation schemes, have also been applied to the Liouville model, and in the context of quantum gravity to the LFT . In all cases we will illustrate how the main physics lies in the selection mechanism for the traveling wave solutions of the non linear RG equation.
The study proceeds as follows. We consider
$`G(x)=H[\mu =e^{\beta x}]={\displaystyle DVe^{S[V]}}`$ (125)
with $`S[V]={\displaystyle d^2𝐫\left(\frac{1}{8\pi \sigma }(V(𝐫))^2+U[V]\right)}`$ (126)
One can perform a Wilson RG analysis (or if one prefers a suitably truncated exact RG analysis), and one finds in $`d=2`$ the flow for the local part $`U_l[V]`$ as:
$`_lU=2U+\sigma U^{\prime \prime }+O(U^2)`$ (127)
There may also be corrections to $`\sigma `$ to $`O(U^2)`$ (in the Sinh-Gordon model), but we focus for now on the RG to linear order. Let us recall that the initial condition is $`U_{l=0}[V]=\mu e^{\beta V}`$ for Liouville and $`U_{l=0}[V]=\mu e^{\beta V}+\mu e^{\beta V}`$ for sinh-Gordon, and that we are interested in the small $`\mu `$ limit. In this limit the initial condition corresponds to a a very wide well $`U[V]`$ (e.g. in the sinh-Gordon model) with a very small curvature $`U^{\prime \prime }[0]`$. To obtain the free energy exponent as $`\mu 0`$, one simply iterates the RG until $`U_l^{\prime \prime }[0]O(1)`$ at a scale $`l^{}=\mathrm{ln}(L^{}/a_0)`$ (more precisely $`U_l^{\prime \prime }[0]1/(\sigma a_0^2)`$ where $`a_0`$ is the bare UV cutoff of the model). At this scale, the free energy is $`O(1)`$, as can be estimated from gaussian fluctuations (straightforwardly at least in the SG model) and thus the initial free energy is:
$`FA_d(\beta )\left({\displaystyle \frac{L}{L^{}}}\right)^2`$ (128)
Remarkably, it is now possible to use what we learned in the previous Sections and demonstrate the “freezing” transition at $`\beta =\beta _c`$ (corresponding to the glass transition for the particle) simply from the RG to this order. Indeed the solution of the truncated equation is:
$`U_l[V]=e^{2l}{\displaystyle \frac{1}{\sqrt{4\pi \sigma l}}}{\displaystyle 𝑑V^{}\mathrm{exp}\left(\frac{(VV^{})^2}{4\sigma l}\right)U_{l=0}[V]}`$ (129)
A straigthforward conclusion would then be that the exact solution corresponding to Liouville is:
$`U_l[V]=\mu e^{(2+\sigma \beta ^2)l}e^{\beta V}`$ (130)
and similarly for the sinh-Gordon:
$`U_l[V]=2\mu e^{(2+\sigma \beta ^2)l}\left(e^{\beta V}+e^{\beta V}\right)`$ (131)
since $`\mathrm{exp}(\pm \beta V)`$ are exact eigenvectors of the linear RG equation for any $`\beta `$. From (128) this immediately yields the “naive dimensional” result for the free energy
$`FL^2\mu ^{\frac{1}{1+(\beta /\beta _c)^2}}`$ (132)
with $`\beta _c=\sqrt{2/\sigma }`$. As we know from the above exact result this is correct for $`\beta <\beta _c`$. Note how the potential $`U_l[V]`$ evolves. Using the notation $`\mu =e^{\beta x}`$ (natural from our extremal statistics interpretation) it forms a “Liouville wall”, which can be seen as a “front solution” moving as $`\mathrm{exp}(\beta (Vxcl))`$ towards $`U=0`$ (and in the sinh-Gordon model there are two symmetric walls moving towards $`U=0`$ and reaching it at $`l=l^{}`$). The Liouville front velocity is:
$`c={\displaystyle \frac{2}{\beta }}+\sigma \beta `$ (133)
which plotted as a function of $`\beta `$ is the famous parabola, such that two values of $`\beta `$ corresponds to the same $`c`$, which is also a well known property of Liouville theory.
As we now show, (130) is incorrect for $`\beta \beta _c`$. This is so for a subtle reason, as apparently the statement that $`\mathrm{exp}(\beta V)`$ is an exact eigenvector of the linear RG (and of (129)) cannot fail ! However, by now we are well used to fronts: in fact we have encountered exactly the same equation in our previous solution of the REM model via RG ($`h_l(x)`$ in (73) is identical to $`U_l(V)`$ in (129)). To describe correctly the bare Liouville (or equivalently the Sinh-Gordon) model one should generalize the initial condition $`U_{l=0}[V]`$, still assuming that $`U_{l=0}[V]\mathrm{exp}(\beta (Vx))`$ for $`Vx`$ ($`x`$ here is very negative corresponding to a small $`\mu `$). Then one can use the saddle point method to estimate (129) as was done in (73) to evaluate $`h_l(x)`$ and one discovers that for $`\beta >\beta _c`$ the velocity freezes into:
$`c=2\sqrt{2\sigma }`$ (134)
which yields a free energy
$`FL^2\mu ^{\frac{1}{2}}`$ (135)
instead of the naive dimensional estimate, thus in agreement with our expectation for the SG model (108c). In addition we find that the decay of the renormalized potential $`U_l[V]e^{\alpha V}`$ is frozen at $`\alpha =\beta _c`$ for all $`\beta >\beta _c`$ consistent with (124).
What has happened is that although $`U_{l=0}[V]=\mathrm{exp}(\beta (Vx))`$ is indeed formally an exact eigenvector, it is dynamically unstable, i.e if one chooses another function with the same large positive $`Vx`$ behaviour one gets a different velocity (which is not the case for $`\beta <\beta _c`$). It is easy to see that the choice $`U_{l=0}[V]=\mathrm{exp}(\beta (Vx))`$ exactly for all $`V`$ does not make sense for $`V\mathrm{}`$. Indeed it is immediately spoiled by the slightest amount of coarse graining (as would appear also by considering the non linearities in the RG equation). The simplest way to see it is to notice that the coarse grained potential:
$`\stackrel{~}{U}[V]=\mathrm{ln}\left({\displaystyle 𝑑v\mathrm{exp}(\left[\mu e^{\beta (V+v)}\right]\frac{v^2}{2s})}\right)`$ (136)
does not grow as $`\mathrm{exp}(\beta V)`$ for large negative $`V`$ but much slower as $`V^2`$. To illustrate further the point let us consider the initial condition:
$`U_{l=0}[V]={\displaystyle \frac{e^{\beta (Vx)}}{1+e^{\beta (Vx)}}}`$ (137)
It behaves as $`e^{\beta (Vx)}`$ for large positive $`Vx`$ (and thus corresponds to the Liouville model) but goes to $`1`$ on the other side. For $`\beta =+\mathrm{}`$ it is easy to compute $`U_l^{\prime \prime }[V=0]`$ from (129) since $`U_{l=0}[V]=\theta (xV)`$. One finds:
$`U_l^{\prime \prime }[V=0]e^{2l\frac{x^2}{4\sigma l}}`$ (138)
and thus one has that $`l^{}`$ defined above is such that:
$`cl^{}=xc=2\sqrt{2\sigma }`$ (139)
This is in fact valid for all $`\beta >\beta _c`$ as was shown in detail in previous sections.
Thus the freezing transition can be obtained from the linearized (i.e lowest order) RG equations, using only elementary insight from coarse graining or the existence of higher order non linear terms. It provides an interesting example where the naive dimensions hold in some regime but are modified in another. Of course, as we have seen in Section III D from the study of fronts, to really establish the transition and determine the universality class one needs to consider higher order non linearities in (127) which goes beyond this paper. For the LFT in quantum gravity the reader can find some exact functional RG studies in Ref. . Although not discussed in this reference, the non linear RG there seems to also exhibit traveling front solutions, whose physics may be important in understanding the problem of the $`c=1`$ barrier.
## VI Critical Dirac fermions in a random gauge field
In this section we relate our RG study of the previous section to the study of the critical wave functions of 2D Dirac fermions in a random magnetic field. We first confirm the results of for the multifractal spectrum, and obtain their finite size corrections. Then we study the transition from the weak disorder to the strong disorder phase, related to the glass transition for the particle, and find that the strong disorder phase has a new and non trivial structure, leading to what we call quasi-localized eigenstates.
### A Critical wave function of 2D random Dirac
Let us first recall the problem of a massless two dimensional Dirac fermion in a static random magnetic field . This model, and its non abelian generalizations, has received a lot of attention in connection with the integer quantum Hall effect transitions with disorder. As discussed in two dimensional Dirac fermions can experience three generic types of disorder: random gauge, random mass and random potential. Random gauge disorder is believed to be a line of fixed points in this general model and is still not yet fully understood. Here we address only the random gauge disorder model of hamiltonian:
$$H=\sigma _\mu \left(iv_F_\mu A_\mu (𝐫)\right)$$
(140)
where the $`\sigma _{1,2}`$ are the $`2`$x$`2`$ Pauli matrices and $`\mu =1,2`$ (we set the Fermi velocity $`v_F=1`$ from now on). The random magnetic field $`𝐁`$ corresponding to the gauge potential $`𝐀`$ is chosen to be gaussian with mean value $`\overline{𝐁(𝐫)}=0`$. The type of correlations studied here correspond to the most interesting case where the gauge potential has short range correlations. In the Coulomb gauge, we can introduce the scalar potential $`\varphi `$ such that $`A_\mu =ϵ_{\mu \nu }_\nu \varphi ,B(𝐫)=_\mu ^2\varphi (𝐫)`$. The gaussian distribution of $`\varphi (𝐫)`$ is thus given by
$$P[\varphi ]=cte\times e^{\frac{1}{4\pi g}_𝐫(_\mu \varphi (𝐫))^2}$$
(141)
where $`g`$ parametrises the strength of the random magnetic field $`B`$. The correlator of the function $`\varphi (𝐫)`$ is thus:
$`\overline{(\varphi (𝐫)\varphi (𝐫^{}))^2}2g\mathrm{ln}{\displaystyle \frac{|𝐫𝐫^{}|}{a}}`$ (142)
In this model, the wave functions at energy $`E`$ are localized for all energies other than the critical energy $`E=0`$. We restrict our study to the $`E=0`$ critical eigenstate, which satisfy:
$$H\mathrm{\Psi }_0(𝐫)=0$$
(143)
For a system of finite size $`L`$ with appropriate boundary conditions there are two independent normalized solutions of (143): the first one can be written $`\mathrm{\Psi }_{0,1}(𝐫)=(\mathrm{\Psi }_0(𝐫),0)`$ with:
$$\mathrm{\Psi }_0(𝐫)^2=\frac{e^{2\varphi (𝐫)}}{_𝐫e^{2\varphi (𝐫^{})}}$$
(144)
the second one being $`\mathrm{\Psi }_{0,2}(𝐫)=(0,\stackrel{~}{\mathrm{\Psi }}_0(𝐫))`$ where $`\stackrel{~}{\mathrm{\Psi }}_0(𝐫))`$ is given by (144) changing $`\varphi (𝐫)\varphi (𝐫)`$. We denote $`_𝐫`$ having in mind either a discrete problem, or a continuous problem with some short scale cutoff $`a`$.
### B participation ratios and multifractal spectrum
Thus in a given configuration of disorder $`\varphi (𝐫)`$ the quantum probability $`|\mathrm{\Psi }_0(𝐫)|^2`$ is identical to the Gibbs probability $`p(𝐫)`$ defined in (2) for the particle in the logarithmically correlated random potential $`V(𝐫)`$ with the correspondence:
$`|\mathrm{\Psi }_0(𝐫)|^2=p(𝐫)`$ (145)
$`2\varphi (𝐫)=\beta V(𝐫)`$ (146)
and thus the model depends on a single parameter $`g=\frac{1}{2}\beta ^2\sigma `$. As we have discussed in previous Sections the particle in the logarithmically correlated random potential undergoes a transition at $`\beta _c=\sqrt{2/\sigma }`$ at which its Gibbs measure changes from being dominated by many sites (high $`T`$ phase) to being dominated to a few sites (low $`T`$ phase). Thus in the quantum problem we expect a transition at:
$`g=g_c=1`$ (147)
with a weak disorder phase for $`g<1`$ and a strong disorder phase for $`g>1`$. In the weak disorder phase the quantum probability (and thus observables such as the mean squared position fluctuations $`<r^2><r>^2`$) is delocalized ($`<..>`$ means averages over $`\mathrm{\Psi }_0`$). In the strong disorder phase the quantum probability is more concentrated, but it cannot be called localized in the usual sense (of an exponential decay around a single center) and in fact both phases have rather peculiar properties.
Properties of wave functions can be described by the inverse participation ratios defined from the normalized wave function $`\mathrm{\Psi }_0(𝐫)`$ in a system of size $`L`$ by
$$R_q(L)=d^2𝐫|\mathrm{\Psi }_0(𝐫)|^{2q}=d^2𝐫\left(p(𝐫)\right)^q$$
(148)
At a very qualitative level, the nature of the eigenfunction can be inferred from the scaling behaviour of the inverse participation ratio with the system size $`L`$ : for an exponentially localized state $`R_q(L)`$ scales as $`R_q(L)\text{const}`$ for all $`q>0`$, while for a plane wave delocalized state we get $`R_q(L)L^{2(q1)}`$. In addition to the localized and delocalized states, there exist states such that $`\tau (q)=\mathrm{ln}R_q(L)/\mathrm{ln}L`$ is a non linear function of $`q`$: they correspond to multifractal wave functions whose moments cannot be described by a single length as usual but rather by a spectrum of exponents. Here, as in , we also find intermediate multifractal behaviour.
To compute the finite size inverse participation ratios we can use the information of previous Sections since:
$$s_q(L)=\mathrm{ln}R_q(L)=\mathrm{ln}Z_{q^2g}+q\mathrm{ln}Z_g$$
(149)
where we have defined $`Z_g=Z(\beta =\sqrt{2g/\sigma })`$ where $`Z(\beta )`$ is the partition function for the particle at inverse temperature $`\beta `$. In particular we will be interested in the multifractal asymptotic scaling exponent $`\tau (q)`$ defined by
$$\tau (q)=\underset{L\mathrm{}}{lim}\frac{s_q(L)}{\mathrm{ln}L}$$
(150)
These exponents were computed previously in using the REM approximation. Here we use our RG results and also obtain finite size corrections. Note that these correspond to properties of $`\mathrm{\Psi }_0`$ defined above and could be changed if other boundary conditions were used.
From the previous Sections we obtain:
$`\mathrm{ln}Z_g=2(1+g)\mathrm{ln}L+\mathrm{\Delta }_g`$ $`,g<1`$ (151a)
$`\mathrm{ln}Z_g=(\sqrt{g}(4\mathrm{ln}L{\displaystyle \frac{1}{2}}\mathrm{ln}(\mathrm{ln}L))+\mathrm{\Delta }_g`$ $`,g=1`$ (151b)
$`\mathrm{ln}Z_g=(\sqrt{g}(4\mathrm{ln}L{\displaystyle \frac{3}{2}}\mathrm{ln}(\mathrm{ln}L))+\mathrm{\Delta }_g`$ $`,g>1`$ (151c)
where $`\mathrm{\Delta }_g`$ is a sample dependent variable of order $`O(1)`$ with a $`g`$ dependent distribution (whose tails we have characterized previously). From there we obtain $`s_q(L)`$, which have different behaviours in the two phases.
(i) weak disorder phase
For $`g<1`$ we find, denoting $`q_c=\frac{1}{\sqrt{g}}`$ :
$`s_q(L)=2(q1)(1gq)\mathrm{ln}L+A_{q,g}`$
$`\text{for }|q|<q_c`$
$`s_q(L)={\displaystyle \frac{2}{\sqrt{g}}}(1\text{sgn}(q)\sqrt{g})^2\mathrm{ln}L+{\displaystyle \frac{1}{2}}\mathrm{ln}\mathrm{ln}L+A_{q,g}`$
$`\text{for }|q|=q_c`$
$`s_q(L)=2q(1\text{sgn}(q)\sqrt{g})^2\mathrm{ln}L+{\displaystyle \frac{3}{2}}|q|\sqrt{g}\mathrm{ln}\mathrm{ln}L+A_{q,g}`$
$`\text{for }|q|>q_c`$
where $`A_{q,g}`$ is a fluctuating part of order $`O(1)`$.
(ii) strong disorder phase
For $`g>1`$ we find:
$`s_q(L)=2(q\sqrt{g}1)^2\mathrm{ln}L{\displaystyle \frac{3}{2}}q\sqrt{g}\mathrm{ln}\mathrm{ln}L+A_{q,g}`$ (153a)
$`\text{for }|q|<q_c`$ (153b)
$`s_q(L)={\displaystyle \frac{1}{2}}\mathrm{ln}\mathrm{ln}L+A_{q,g}\text{for }q=q_c`$ (153c)
$`s_q(L)=A_{q,g}\text{for }q>q_c`$ (153d)
$`s_q(L)=2|q|\sqrt{g}\left(4\mathrm{ln}L{\displaystyle \frac{3}{2}}\mathrm{ln}\mathrm{ln}L\right)+A_{q,g}`$ (153e)
$`\text{for }q<q_c`$ (153f)
$`s_q(L)=2\left(4\mathrm{ln}L{\displaystyle \frac{1}{2}}\mathrm{ln}\mathrm{ln}L\right)+A_{q,g}\text{for }q=q_c`$ (153g)
where $`A_{q,g}`$ is a fluctuating part of order $`O(1)`$.
The corresponding scaling exponents $`\tau (q)`$ are thus identical to the one found in and in addition we have obtained their finite size corrections. In the weak disorder phase one for $`q>0`$:
$$\tau (q)=\{\begin{array}{ccc}2(q1)\left(1\frac{q}{q_c^2}\right)\hfill & \text{ for }\hfill & qq_c=\sqrt{\frac{1}{g}}\hfill \\ 2q\left(1\frac{1}{q_c}\right)^2\hfill & \text{ for }\hfill & qq_c\hfill \end{array}$$
(154)
which means a parabolic form with a termination point at $`q=q_c`$ as represented in the Fig. 14.
In the strong disorder phase $`g>1`$, i.e when $`q_c1`$, the above expression becomes (for $`q>0`$):
$$\tau (q)=\{\begin{array}{ccc}2\left(1\frac{q}{q_c}\right)^2\hfill & \text{ for }\hfill & qq_c=\sqrt{\frac{1}{g}}\hfill \\ 0\hfill & \text{ for }\hfill & qq_c\hfill \end{array}$$
(155)
Since the inverse participation ratio does not scale with the system size $`L`$ for each integer $`q`$, one could naively conclude that it is the sign of a localized state (see however below).
As was discussed in these results can be translated into spectrum for exponent $`\alpha `$. If one assumes that $`p(𝐫)`$ is of order $`L^\alpha `$ in a number $`L^{f(\alpha )}`$ sites then the above spectrum is recovered if:
$`f(\alpha )=8{\displaystyle \frac{(d_+\alpha )(\alpha d_{})}{(d_+d_{})^2}}`$ (156)
with $`d_\pm =2(1\pm \sqrt{g})^2`$ for $`g>1`$ and $`d_+=8\sqrt{g}`$, $`d_{}=0`$ for $`g>1`$. It is easy to see that:
$`<(r<r>)^{2k}>L^{\mathrm{max}_\alpha ((k+1)f(\alpha )\alpha )}`$ (157)
showing that the eigenstate is never localized in the usual sense (exponential decay around a single center) since the exponent is always positive for large enough $`k`$. Since $`lim_{q\pm \mathrm{}}s_q(L)/q=\mathrm{ln}p_{max,min}`$ one obtains that the maximum of the Gibbs measure $`p_{max}=\mathrm{max}_𝐫p(𝐫)`$ and the minimum behave for large $`L`$ as:
$`p_{max}L^{2(1\sqrt{g})}(\mathrm{ln}L)^{\frac{3}{2}\sqrt{g}}`$ (158)
$`p_{min}L^{2(1+\sqrt{g})}(\mathrm{ln}L)^{+\frac{3}{2}\sqrt{g}}`$ (159)
in the weak disorder phase.
### C nature of the strong disorder phase: quasi-localization
Let us now concentrate on the case $`g>1`$. There, we know from previous Sections that the Gibbs measure of the particle is concentrated in a few sites. Thus from (146) the quantum probability $`|\mathrm{\Phi }_0(𝐫)|^2`$ is also concentrated in a few sites, analogous to the RSB picture. This is a very peculiar type of eigenstate. Indeed if one computes the quantum average $`<r^2><r^2>`$ in a given sample, it has a finite probability to be of order $`O(L^2)`$. Thus the eigenstate cannot be considered as localized in the usual sense. Since it is peaked around a few sites we call it “quasi-localized”. Around these centers the wavefunction decays fast enough to be normalizable. It would be interesting to investigate further the typical spatial decay of such eigenstates around their (multiple) centers, which we expect to be slower than exponential.
## VII conclusion
In this paper we have studied the equilibrium problem of a particle in a random potential with logarithmic correlations, through exact bounds, numerical simulation, qualitative arguments and a renormalization group method that we have developed specifically for this problem. We have shown that it exhibits a glass transition at finite temperature $`T_c=1/\beta _c>0`$ in any dimension. This confirms earlier conjectures and allows for a more detailed study of the problem. The RG method allowed to obtain the universal features of the free energy distribution at low temperature. The relation to the problem of extremal statistic of correlated variables was investigated. It has been found that it exhibits universal finite size corrections, consistent with our numerical calculations.
Most interestingly, we found that this logarithmic model provides a particularly simple example (maybe the simplest) of a finite dimensional model - i.e with translationally invariant disorder correlations\- such that the low temperature phase is non trivial. It is non trivial in the sense that in the thermodynamic limit $`L+\mathrm{}`$, there are, with a finite probability, several low lying states (i.e possible positions of the particle) with energy differences of order one, and separated in space by distances of order $`L`$. Thus the Gibbs measure at low temperature is dominated by “a few” spatially well separated states. Interestingly, this transition and this type of glass phase occurs only for logarithmically growing correlations, faster growth (e.g. as in Sinai model) yielding only a glass phase with single ground state dominance, while slower growth yielding only a high temperature phase.
Although oversimplified in some respect (it has no internal space) it does provide one example of a model where the usual droplet picture (which assumes dominance of a single ground state - or several related by a symmetry) does not apply. Rather, it provides one example where some features of the physics usually associated to RSB, namely dominance by a few states with exponential free energy distributions, can be explicitly exhibited. In fact, due to the finite dimensional correlations, there are some departures from the behaviour observed in the simplest prototype mean field models (such as the REM), as can be seen for instance from the free energy distribution which has more structure than a simple exponential. It would of course be interesting to explore further the additional features specific to finite dimensions.
Although the present model is already of obvious physical interest (in 2D it describes e.g. a single vortex in a random gauge XY model) its non trivial properties provide a motivation to search for models with more degrees of freedom and with similar features. One way to proceed would be to search for interface models via an internal dimensional expansion around the present model. The key feature however appears to be the marginality of the model, i.e the logarithmic growth of typical energy fluctuations. This corresponds to a fluctuation energy exponent $`\theta =0`$, i.e the situation where the temperature (i.e the entropy) is marginal in the RG sense. The droplet arguments indeed assume that $`\theta >0`$, consistent with the single ground state dominance (and activated behaviour typical of a zero temperature fized point where $`T`$ is formally irrelevant). In the situation $`\theta =0`$ one does expect more generally power laws with $`T`$ dependent exponents, reminiscent of mean field. It would thus be of great interest to similarly exhibit other non trivial marginal models (e.g. spin models with $`\theta =0`$) with similar features . Spin models where (domain wall) excitations (in root mean square and in average) also scale logarithmically (as vortices in the random gauge XY model) are presumably good candidates.
On one hand we have developed a specific (Coulomb gas) renormalization group (RG) approach to describe the model. From the study of the resulting nonlinear (KPP) RG equation, we found explicitly that a freezing phenomenon occurs at the glass transition temperature, and that in the glass phase a broad (power law) distribution of fugacities develop - or equivalently exponential distribution of local free energy. It is different from more conventional perturbative RG (e.g. the one which was used to study the dynamics of this model) in the sense that the full distribution of probability is followed. This turns out to be crucial to describe the low temperature phase.
On the other hand, as we have discussed, two approximations of the present model, the REM approximation and the DPCT hierarchical version can both be solved using replica and do require considering the analytical continuation to $`m0`$ of contributions of replica symmetry breaking saddle points .
This shows that a RG approach which is explicitly replica symmetric but allows to treat broad disorder distributions can be consistent with (approximate) approaches based on RSB saddle points . We have illustrated this on the REM which can be recast in terms of non linear RG equations, with a freezing transition. In fact one of the striking property of the model is that the RG equations derived here are similar - to the order we have been working - to the one which hold for a continuous version of the DPCT problem, the branching process. In particular it indicates that both problem share the same universal finite size corrections.
We have also analyzed some connections in 2D (and via boundaries in 1D) between the model of the particle and the Liouville and Sinh Gordon models. The intensive free energy of the particle corresponds to the scaling dimension in these models with $`b=\beta /\beta _c`$. The glass transition corresponds to the weak to strong coupling transition at $`b=1`$. Beyond, corresponding to the glass phase, the scaling dimension freezes as we have also shown via a direct RG approach on these models. We have seen that under coarse graining an additional local field appears in the LM and SGM, with broad distribution, and corresponds to inhomogeneous configurations being generated (and broad fluctuations of the local area since the local partition function corresponds to local area).
The present study raises interesting issues to be explored concerning the relations with the continuum Liouville Field Theory (LFT). An outstanding question is whether the conjecture of is correct for the correlations. Since we have obtained another result linking the problem to the DPCT, the direct comparison of the LFT and the DPCT remains to be studied. If it holds it means that the conformally invariant many point correlations can be related in some limits (large separations with fixed ratios $`\mathrm{ln}r_{ij}/\mathrm{ln}r_{kl}`$) to the results from the tree problem. It would also raise interesting issues about the continuation of the LFT beyond $`b=1`$, and its relation with the non trivial structure of the glass phase (with RSB features) in the equivalent particle model.
We have also extracted from our approach some consequences for the problem of the $`E=0`$ critical eigenstate of 2D Dirac fermions in a random magnetic field. We have confirmed, via our RG method, previous results concerning the multifractal spectrum and extracted their finite size corrections. We have found that the non trivial low $`T`$ phase of the particle translates into peculiar quasi-localized eigentstates for the quantum problem, peaked around a few distant centers. It raises the question of whether this property can be present in other quantum systems.
Another interesting question is whether the transition studied here has a signature in the dynamics as well. Note that a similar non trivial structure at low temperature is also present in the the Sinai model with a bias, which renormalizes onto a random walk with algebraic waiting times distribution . However this is a driven system and it would be interesting to see whether non driven systems in low dimension can exhibit similar features.
Finally, an outstanding question is how the present model can be studied using 2D conformal field theory (CFT). In particular one wonders what is the signature in this context, of the physics which was unveiled here, reminiscent of RSB, using RG with broad distributions. The freezing phenomenon within the non linear RG, which transforms the naive scaling dimensions into non trivial ones, should correspond to a similar mechanism in CFT. Recent progress on CFT classification of disordered models where supersymmetry can be used allows to hope that such progress is within sight. We hope that the present RG method will apply to study other two dimensional models with similar features and shed light on the more formal field theoretic methods.
We acknowledge useful discussions with B. Derrida. This work was supported in part by NSF grant DMR-9528578 (D.C.).
## A Existence of a transition
We use the same method as Derrida and Cook for the directed polymer problem . It is easy to compute the first two moments pf $`P[Z]`$, using translational invariance and periodic boundary conditions:
$`\overline{Z}=L^de^{\frac{\beta ^2}{2}\mathrm{\Gamma }_L(0)}L^{d+\beta ^2\sigma }`$ (A1)
and
$`\overline{Z^2}=L^de^{\frac{\beta ^2}{2}4\mathrm{\Gamma }_L(0)}{\displaystyle \underset{𝐫}{}}e^{\frac{\beta ^2}{2}\stackrel{~}{\mathrm{\Gamma }}_L(𝐫)}`$ (A2)
$`BL^{2d+2\beta ^2\sigma }\beta <\beta _2`$ (A3)
the last estimate being valid as long as the sum over $`𝐫`$ is divergent, i.e $`\beta <\beta _2=\sqrt{d/(2\sigma )}`$. The constant $`B>1`$ depends on the details of the model, e.g. for $`d=1`$ one can write $`B=lim_{L+\mathrm{}}_0^1𝑑y\mathrm{exp}(\frac{\beta ^2}{2}(\stackrel{~}{\mathrm{\Gamma }}_L(Ly)4\sigma \mathrm{ln}L))`$. Thus for $`\beta <\beta _2`$ the ratio
$`{\displaystyle \frac{\overline{Z^2}}{\overline{Z}^2}}B`$ (A4)
as $`L+\mathrm{}`$. In it is shown that the property (A4) implies that:
$`\text{prob}({\displaystyle \frac{1}{d\mathrm{ln}L}}\mathrm{ln}Z={\displaystyle \frac{1}{d\mathrm{ln}L}}\mathrm{ln}\overline{Z})1/B`$ (A5)
as $`L+\mathrm{}`$. If we take for granted that the free energy is self averaging, it implies that for $`T>T_2=\sqrt{2\sigma /d}`$ the quenched and annealed (intensive) free energies coincide exactly $`f(T)=f_A(T)`$. Thus for $`T>T_2`$ the (intensive) entropy is $`s(T)=s_A(T)=_Tf_A(T)`$ and thus one has:
$`s(T)=1{\displaystyle \frac{\sigma }{dT^2}}T>T_2`$ (A6)
Since $`s_A(T)`$ becomes negative below $`T=T_g=\sqrt{d/\sigma }`$ it implies that there must be a temperature $`T_c<T_2`$ at which (A6) breaks down and thus a phase transition. Although this is harder to prove, it seems that here (A6) holds down to $`T_c=T_g`$.
Awaiting a rigorous mathematical proof, we have not attempted to prove self averaging of $`f`$. Not only is it highly reasonable in view of our other results but in fact if it were not the case, the above argument would imply a rather curious - and unphysical - distribution for $`f`$ (with a delta peak of non zero weight smaller than one). In addition, as noted in , by adjusting the small scale details of the model, the constant $`B`$ can be chosen as close to $`1`$ as wanted.
## B extremal statistics of correlated variables
In this Section we summarize some results on the extremal statistics of a set of random variables. We selected the ones which are useful in putting the problem studied here in a broader context. We recall some of the classic results from probability theory and we have chosen to illustrate them by adding a few simple arguments which emphasize the importance of some of these results to the physics of disordered systems. We denote the $`N`$ random variables either $`X_r`$, $`r=1,..N`$ when they are normalized in a particular way, or $`V_r`$ when they can be readily interpreted as the random potential variables studied here (the two differing by a trivial uniform rescaling $`V(r)V_rX_r`$). They apply directly to describe $`d=1`$ ($`N=L`$) and can be usually extended to $`d>1`$ ($`V(𝐫)`$ and $`N=L^d`$).
### 1 uncorrelated variables
It is natural to start with the case of $`N`$ uncorrelated variables of identical probability distribution $`P(V)`$. The distribution $`P(V)`$ can belong to three classes of extremal statistics, but we will recall only the Gumbell class. Schematically for this class, a well known theorem states that there exist constants $`a_N`$ and $`b_N`$ such that for a fixed $`\stackrel{~}{y}`$
$`Prob(V_{min}>b_N\stackrel{~}{y}a_N)\mathrm{exp}(e^{\stackrel{~}{y}})`$ (B1)
The constants $`a_N`$ and $`b_N`$ are determined as:
$`\mathrm{ln}{\displaystyle _{\mathrm{}}^{a_N}}𝑑VP(V)={\displaystyle \frac{1}{N}}`$ (B2)
$`b_N=N{\displaystyle _{\mathrm{}}^{a_N}}𝑑y{\displaystyle _{\mathrm{}}^y}𝑑VP(V)`$ (B3)
For variables $`X_r`$ chosen from a centered Gaussian of unit variance $`P(X)=\frac{1}{\sqrt{2\pi }}e^{X^2/2}`$, one can choose $`a_N`$ and $`b_N`$ as:
$`b_N={\displaystyle \frac{1}{\sqrt{2\mathrm{ln}N}}}`$ (B4)
$`a_N=\sqrt{2\mathrm{ln}N}{\displaystyle \frac{1}{\sqrt{2\mathrm{ln}N}}}{\displaystyle \frac{1}{2}}\mathrm{ln}(4\pi \mathrm{ln}N)`$ (B5)
and thus one can then write schematically that:
$`X_{min,N}\sqrt{2\mathrm{ln}N}+{\displaystyle \frac{1}{\sqrt{2\mathrm{ln}N}}}\left({\displaystyle \frac{1}{2}}\mathrm{ln}(4\pi \mathrm{ln}N)+\stackrel{~}{y}\right)`$ (B6)
where $`\stackrel{~}{y}`$ is distributed with the Gumbell distribution $`p(\stackrel{~}{y})=e^{\stackrel{~}{y}}\mathrm{exp}(e^{\stackrel{~}{y}})`$.
It is useful to note the property of reparametrization associated to a monotonous function $`\psi (V)`$. If one has (B1) for the minimum $`V_{min}`$ of the variables $`V_r`$ with the constants $`a_N`$ and $`b_N`$, one also has (under some weak conditions) that (B1) for the minimum $`\psi (V_{min})`$ of the variables $`\psi (V_r)`$ with the constants $`a_N^{}=\psi (a_N)`$ and $`b_N^{}=b_N/\psi ^{}(a_N)`$. Note also that we have illustrated how to show convergence to Gumbell (and generalized it to finite temperature) in the text.
For completeness we recall the necessary conditions for the convergence to Gumbell (i.e $`P(V)`$ belonging to the Gumbell class). First $`P(V)`$ must decay fast enough at $`V\mathrm{}`$ so that there exists $`y_0`$ such that:
$`{\displaystyle _{\mathrm{}}^{y_0}}𝑑y{\displaystyle _{\mathrm{}}^y}P(V)𝑑V<+\mathrm{}`$ (B7)
and second, defining
$`R(t)={\displaystyle \frac{1}{_{\mathrm{}}^tP(V^{})𝑑V^{}}}{\displaystyle _{\mathrm{}}^t}𝑑y{\displaystyle _{\mathrm{}}^y}P(V)𝑑V`$ (B8)
one must have for all $`x<y_0`$:
$`\underset{t\mathrm{}}{lim}{\displaystyle \frac{_{\mathrm{}}^{t+xR(t)}P(V)𝑑V}{_{\mathrm{}}^tP(V^{})𝑑V^{}}}=e^x`$ (B9)
These conditions are in fact rather broad. Finally, note also the very powerful theorem 2.10.1 of for the rate of convergence to the Gumbell fixed point.
### 2 correlated variables
#### a general lower bound
We now consider correlated variables with distribution $`P(V_1,\mathrm{}V_N)`$. Let us start with a simple but very general bound and extract the consequences. One has:
$`G(x)=\text{Proba}(V_{min}<x){\displaystyle \underset{r=1,N}{}}\text{Proba}(V_r<x)`$ (B10)
since the reunion of all events $`V_r<x`$ implies the event $`V_{min}<x`$ and that $`Prob(AUB)Prob(A)+Prob(B)`$ (the bound is exactly saturated e.g. when there are strong correlations such that $`V_rV_r^{}>x`$ for all $`rr^{}`$). For variables which have identical one particle distribution $`P(V_r)=_{r^{}r}dV_r^{}P(V_1,\mathrm{}V_N)`$ one has:
$`G(x)N{\displaystyle _{\mathrm{}}^x}P(V)𝑑V`$ (B11)
Let us illustrate the consequences for correlated variables $`X_1,..X_N`$ such that the one particle distribution is a unit centered Gaussian. Then it implies for $`x\mathrm{}`$:
$`G(x){\displaystyle \frac{N}{\sqrt{2\pi }x}}e^{x^2/2}`$ (B12)
from which one immediately sees that it implies:
$`\text{Proba}\left(X_{min}<x_N=\sqrt{2\mathrm{ln}N}+\alpha {\displaystyle \frac{\mathrm{ln}(4\pi \mathrm{ln}N)}{\sqrt{2\mathrm{ln}N}}}\right)`$ (B13)
$`{\displaystyle \frac{1}{(4\pi \mathrm{ln}N)^{(\frac{1}{2}\alpha )}}}\underset{N+\mathrm{}}{\overset{}{}}0`$ (B14)
by choosing $`x=x_N`$, for any $`\alpha <1/2`$. Thus one has a general lower bound for the minimum of correlated variables. In particular for gaussian variables such that $`\overline{V_r^2}=2\sigma \mathrm{ln}N=2d\sigma \mathrm{ln}L`$ one gets:
$`V_{min}>2d\sqrt{\sigma }\mathrm{ln}L+\sqrt{\sigma }\alpha \mathrm{ln}(4\pi d\mathrm{ln}L)`$ (B15)
with probability $`1`$ in the large $`L`$ limit for any $`\alpha <1/2`$. Moreover choosing $`\alpha =1/2`$ and writing:
$`V_{min}=2d\sqrt{\sigma }\mathrm{ln}L+\sqrt{\sigma }({\displaystyle \frac{1}{2}}\mathrm{ln}(4\pi d\mathrm{ln}L)+\stackrel{~}{y})`$ (B16)
one finds that:
$`Prob(\stackrel{~}{y}<y)e^y`$ (B17)
This yields a lower bound which can be compared with the REM approximation defined in the text. Note that the above upper bound is the exact behaviour of the Gumbell distribution at large negative $`y`$, so in a sense the REM approximation saturates the bound in the tails. Consequently, to allow for a larger tail (such as $`ye^y`$ one needs at least a coefficient of $`\mathrm{ln}\mathrm{ln}N`$ strictly larger than $`1/2`$).
#### b short range correlations and convergence to Gumbell
Let us now consider $`N`$ centered gaussian variables $`X_r`$ with a fixed correlation matrix $`\mathrm{\Gamma }_{rr^{}}=\overline{X_rX_r^{}}`$, normalized so that $`\mathrm{\Gamma }_{rr}=1`$. A powerful bound, which refines (B10) above, allows to easily demonstrate convergence to the Gumbell distribution for a large class of “short enough range” correlations. It compares two arbitrary correlators $`\mathrm{\Gamma }_{rr^{}}^{(1)}`$ and $`\mathrm{\Gamma }_{rr^{}}^{(2)}`$ with $`\mathrm{\Gamma }_{rr}^{(1)}=\mathrm{\Gamma }_{rr}^{(2)}=1`$. Their associated $`G(x)`$ functions satisfy :
$`|G_1(x)G_2(x)|{\displaystyle \underset{rr^{}}{}}{\displaystyle \frac{|\mathrm{\Gamma }_{rr^{}}^{(1)}\mathrm{\Gamma }_{rr^{}}^{(2)}|}{2\pi (1m_{rr^{}}^2)^{1/2}}}e^{\frac{x^2}{1+m_{rr^{}}}}`$ (B18)
with $`m_{rr^{}}=max(|\mathrm{\Gamma }_{rr^{}}^{(1)}|,|\mathrm{\Gamma }_{rr^{}}^{(2)})|`$. It is obtained by bounding $`G(x)/\mathrm{\Gamma }_{rr^{}}`$ and integrating between $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_2`$. It will be used to compare $`\mathrm{\Gamma }_{rr^{}}^{(1)}=\mathrm{\Gamma }_{rr^{}}`$ with the uncorrelated case $`\mathrm{\Gamma }_{rr^{}}^{(2)}=\delta _{rr^{}}`$.
To adress the question of the universality of the Gumbell distribution, let us now consider a ($`d=1`$) translationally invariant correlator $`\mathrm{\Gamma }_{rr^{}}=\mathrm{\Gamma }(rr^{})`$ with $`\mathrm{\Gamma }(0)=1`$ where $`\mathrm{\Gamma }(r)`$ is a N-independent function which decays to zero as $`rr^{}+\mathrm{}`$.
Inserting $`x=a_N\stackrel{~}{y}b_N`$ of (B1) and (B5) into (B18) one easily gets that if $`\mathrm{\Gamma }(r)`$ decreases fast enough, one has $`G(x=a_N\stackrel{~}{y}b_N)=\mathrm{exp}(e^{\stackrel{~}{y}})`$ at large $`N`$, i.e one has convergence to the Gumbell distribution with exactly the same coefficients $`a_N`$ and $`b_N`$ as in the uncorrelated case, so that (B6) still holds. As one sees by studying the bound this result holds as long as $`\mathrm{\Gamma }(r)`$ decreases faster than $`1/\mathrm{ln}(r)`$ (this is theorem 3.8.2. of ). The limiting case (which does not satisfy Gumbell, as discussed below) being $`\mathrm{\Gamma }(r)\tau /\mathrm{ln}(r)`$ at large $`r`$.
Let us give a simple self consistency argument, more enlightening than the bounds, which explains why $`\mathrm{\Gamma }(r)1/\mathrm{ln}r`$ should be the limiting case between the short range (Gumbell) universality class and other behaviours. Let us split a set of $`2N`$ correlated variables $`X_1,..X_{2N}`$ into subsystem 1, $`X_1,..X_N`$, and subsystem 2, $`X_{N+1},..X_{2N}`$. If correlations are very short ranged (e.g. exponentially decaying) it seems reasonable to first neglect correlations between 1 and 2 and find the minimum in each subsystem, which read respectively:
$`\stackrel{~}{X}_{\mathrm{min},i}\sqrt{2\mathrm{ln}N}{\displaystyle \frac{1}{2}}{\displaystyle \frac{\mathrm{ln}(4\pi \mathrm{ln}N)}{\sqrt{2\mathrm{ln}N}}}+{\displaystyle \frac{x_i}{\sqrt{2\mathrm{ln}N}}}`$ (B19)
with $`i=1,2`$ and where $`x_1`$, $`x_2`$ are independently distributed with the Gumbell distribution. The symbol $`\stackrel{~}{V}`$ indicates that the minimum (in each subsystem) is with respect to a slightly different distribution than the original one, since all cross correlations between the two different subsystems have been set to zero. The second stage is to add the correlations between the two subsystems. Typically, the minima $`1`$ and $`2`$ will be a distance $`N`$ apart and thus their original cross-correlation is $`\mathrm{\Gamma }(N)`$, and thus, for short range correlations, very small compared to the fluctuating part $`x_i/\sqrt{2\mathrm{ln}N}`$. Thus the distribution of the mininum $`X_{min}^{(2N)}`$ of the original $`2N`$ variables should be given with better and better accuracy at large $`N`$, as $`X_{min}^{(2N)}=\mathrm{min}(\stackrel{~}{X}_{\mathrm{min},1},\stackrel{~}{X}_{\mathrm{min},2})`$ (which is automatically satisfied by the approximation (B19)). The corrections are irrelevant at large scale provided the typical root mean cross correlation between the subsystems remain smaller than the typical fluctuations of the minimum in each subsystem, a condition which reads:
$`\sqrt{\mathrm{\Gamma }(N)}1/\sqrt{\mathrm{ln}N}`$ (B20)
which indeed gives correctly the basin of attraction of the Gumbell distribution. Furthermore, in the limiting case $`\mathrm{\Gamma }(r)\tau /\mathrm{ln}r`$ the above argument shows that the distribution of the $`x_i`$ should be changed, which is also the case, as we now examine.
So, to summarize, if correlations are short ranged with $`\mathrm{\Gamma }(r)`$ decreasing faster than $`1/\mathrm{ln}(r)`$ this is the “SR universality class”. It includes the REM, and one can check that the finite size corrections in are reproduced (at $`T=0`$).
#### c long range correlations and absence of convergence to Gumbell
There is a simple but instructive model of correlated variables which can be easily solved and illustrate cases where Gumbell does not hold. If one takes:
$`V_r^{}=V_r+U`$ (B21)
with $`V_r`$ a set of uncorrelated gaussian variables, $`U`$ a gaussian variable uncorrelated with the $`V_r`$. Then clearly, if one chooses the variance of $`U`$ big enough (B6) cannot hold. To keep using normalized variables ($`\mathrm{\Gamma }_{rr}=1`$) one defines:
$`X_r^{}={\displaystyle \frac{1}{\sqrt{1+w_N}}}X_r+u\sqrt{{\displaystyle \frac{w_N}{1+w_N}}}`$ (B22)
where $`u`$ is a centered gaussian random variable with unit variance. The correlation matrix is then $`\mathrm{\Gamma }_{rr^{}}^{}=\frac{1}{1+w_N}(\delta _{rr^{}}+w_N)`$. Clearly one has:
$`X_{min}^{}={\displaystyle \frac{1}{\sqrt{1+w_N}}}X_{min}+u\sqrt{{\displaystyle \frac{w_N}{1+w_N}}}`$ (B23)
Using the expression (B6) for $`X_{min}`$ one sees that for deviations from Gumbell to arise one needs that $`w_N\tau /\mathrm{ln}N`$. In that case one gets from (B6) that
$`X_{min}^{}\sqrt{2\mathrm{ln}N}{\displaystyle \frac{1}{2}}{\displaystyle \frac{\mathrm{ln}(4\pi \mathrm{ln}N)}{\sqrt{2\mathrm{ln}N}}}+{\displaystyle \frac{x_i+\sqrt{2\tau }u+\tau }{\sqrt{2\mathrm{ln}N}}}`$ (B24)
These simple considerations thus allow to understand simply the limiting case, that if $`\mathrm{\Gamma }(r)`$ decreases as $`\tau /\mathrm{ln}(r)`$ one has that (B6) still hold (with the same constants) but the distribution of $`\stackrel{~}{y}\tau `$ now converges instead to the convolution of the Gumbell distribution and the gaussian of variance $`2\tau `$ (see e.g. theorem 3.8.2. of ).
Increasing the range of correlations even further, one gets into a regime where the fluctuating part (in the $`X`$ variables) is larger than $`1/\sqrt{\mathrm{ln}N}`$ (and thus in the $`VX\sqrt{\mathrm{ln}N}`$ variables the dominant finite size corrections are non selfaveraging). The fluctuations become then entirely gaussian, being controlled by the $`U`$ part, i.e the $`q=0`$ mode. For instance, if $`\mathrm{\Gamma }(r)`$ decreases as $`1/(\mathrm{ln}(r))^\alpha `$ with $`1/3<\alpha <1`$ then (theorem 3.8.4. of ) one has:
$`P(V_{\mathrm{min}}>\mathrm{\Gamma }(N)^{1/2}x(1\mathrm{\Gamma }(N))^{1/2}\sqrt{2\mathrm{ln}N}`$ (B25)
$`(2\mathrm{ln}N{\displaystyle \frac{1}{2}}\mathrm{ln}(4\pi \mathrm{ln}N))){\displaystyle }_{\mathrm{}}^x2\pi ^{1/2}e^{x^2/2}`$ (B26)
As illustrated below, this behaviour (entirely controlled by the $`q=0`$ mode) is in a sense more long range, and further away from Gumbell than the problem of log-correlated variables that we are interested in and that we now discuss.
#### d log-correlated variables
The case of log-correlated variables is difficult and little is known. We just make a few comments.
Let us first discuss the form of the correlator. The correlator (for the normalized variables $`X_r=V_r/\sqrt{2\sigma \mathrm{ln}N}`$ in $`d=1`$) is of the form:
$`\mathrm{\Gamma }(r)=1{\displaystyle \frac{\mathrm{ln}r}{\mathrm{ln}N}}`$ (B27)
for $`1rN`$. One must then distinguish the two other regions. For small $`r`$, the precise form could vary by adding a short range correlated noise. This is what we term short scale details, and an important question is the extent of universality of the results (scaling of minima, distribution) with respect to the small $`r`$ form. For $`rL`$ the behaviour depends on boundary conditions, which may also be important (see below). For the periodic system in the simulation $`\mathrm{\Gamma }(r)=\mathrm{\Gamma }(Nr)`$ and $`\mathrm{\Gamma }(r)`$ actually becomes negative at $`r=N/2`$ and of order $`c/\mathrm{ln}N`$ (see Section IV). Adding a small uniform $`U`$ noise, as described above in B 2 c, could make $`\mathrm{\Gamma }(N/2)=0`$, so generally speaking one can discuss forms such that $`\mathrm{\Gamma }(N/2)=0`$. Seen as a scaling function of $`z=\frac{\mathrm{ln}r}{\mathrm{ln}N}`$, $`\mathrm{\Gamma }(r)`$ then converges for large $`N`$ towards (B27), but it does have boundary layers at $`z=0`$ and $`z=1`$.
It is useful to plot on the same graph the various cases studied in this Section. This is represented in Figure 15. We have represented schematically $`\mathrm{\Gamma }(r)\mathrm{ln}N`$ versus $`\mathrm{ln}r`$, for the log-correlated form (B27) above, and for the various cases $`\mathrm{\Gamma }(r)\frac{\mathrm{ln}N}{(\mathrm{ln}r)^\alpha }`$ with $`\alpha >1`$ (Gumbell behaviour), $`\alpha =1`$ and $`\alpha <1`$.
As discussed above in the log-correlated case the behaviour of $`\mathrm{\Gamma }(r)`$ near $`\mathrm{ln}r=\mathrm{ln}N`$ can be considered as uncertain to order $`1/\mathrm{ln}N`$. This can be seen either from the $`q=0`$ mode, which depending on boundary conditions one may adjust by this amount, as discussed above, or even looking at the first non trivial mode, $`q=2\pi /L`$ which has a contribution of the same order. We know from the previous paragraph that these contributions can shift the $`x`$ variable by a gaussian, so it makes it unlikely that the Gumbell distribution would hold in that case.
To conclude, we have given the various behaviours as a function of the range of correlations. The presence of the $`\mathrm{ln}\mathrm{ln}N`$ corrections seems to be more robust than the distribution of $`x`$. For the marginal case with $`q=0`$ LR disorder, the same $`\frac{1}{2}\mathrm{ln}\mathrm{ln}N`$ corrections hold as for the REM, while the distribution is changed. On the other hand, for log-correlated variables, we expect a different coefficient $`\frac{3}{2}\mathrm{ln}\mathrm{ln}N`$ as discussed in the text (and we do not expect the Gumbell distribution to hold).
## C Gumbell via RG
In a more detailed analysis Eq. (58) yields $`\mathrm{ln}\mathrm{ln}(1/G_l(x))=l+\mathrm{ln}\mathrm{ln}(1/G_0(x))`$ which can be rewritten in a front-like form:
$`G_l(x)=\mathrm{exp}(e^{l+\varphi (x)})`$ (C1)
where $`\varphi (x)=\mathrm{ln}\mathrm{ln}(1/G_0(x))`$. In this Appendix we set $`dll`$. The center of the front is at $`x=m(l)`$ solution of $`\varphi (m(l))=l`$. One can Taylor expand $`\varphi (x)=l+y+\frac{1}{2}\delta _ly^2+..`$ with $`y=\alpha _l(x+m(l))`$, $`\alpha _l=\varphi ^{}(m(l))`$ and $`\delta _l=\varphi ^{\prime \prime }(m(l))/\varphi ^{}(m(l))^2`$. Thus in the variable $`y`$, $`G_l`$ converges to a Gumbell limit distribution $`𝒢(y)=\mathrm{exp}(e^y)`$. It holds provided higher terms in the Taylor expansion are irrelevant (a necessary, and in simplest cases sufficient, condition being that the second one $`\delta _l0`$).
If no rescaling of disorder is performed, in the relevant large negative $`x`$ region one has $`G_0(x)1`$ and thus $`\varphi (x)\mathrm{ln}(1G_0(x))`$. Two cases must be distinguished because the limit $`T0`$ and $`N+\mathrm{}`$ do not commute:
(i) finite fixed temperature $`T>0`$: then for $`x\mathrm{}`$ one has $`1G_0(x)C_1(\beta )e^{\beta x}(1+O(e^{\beta x}))`$ with $`C_k=_VP(V)e^{k\beta V}`$ and we assume that $`C_1,C_2<+\mathrm{}`$ exists (distributions falling faster than exponentials). Then the situation is simple as $`\varphi (x)=\beta x+\mathrm{ln}C_1(\beta )+O(e^{\beta x})`$, $`m(l)l/\beta +1/\beta \mathrm{ln}C_1(\beta )`$, $`\alpha _l=\beta `$ and $`\varphi ^{\prime \prime }(x)/\varphi ^{}(x)^20`$ exponentially fast. For a Gaussian distribution:
$`m(l){\displaystyle \frac{1}{\beta }}l+{\displaystyle \frac{1}{2}}\sigma \beta `$ (C2)
There is no transition to a glass phase.
(ii) zero temperature: it is an extremal statistics problem. Then clearly $`1G_0(x)`$ does not decay as an exponential. Let us consider a class of distributions such that $`1G_0(x)(A|x|)^\gamma \mathrm{exp}((B|x|)^\alpha )`$ with $`\alpha >1`$ (plus exponentially small corrections). This contains the Gaussian (of variance $`\sigma `$), of most interest here, for $`\alpha =2`$, $`B=1/\sqrt{2\sigma }`$, $`\gamma =1`$ and $`A=\sqrt{2\pi /\sigma }`$. Then one easily finds from above that
$`m(l){\displaystyle \frac{1}{B}}(l{\displaystyle \frac{\gamma }{\alpha }}\mathrm{ln}l\gamma \mathrm{ln}(A/B))^{\frac{1}{\alpha }}`$ (C3)
$`\alpha _lB\alpha (l{\displaystyle \frac{\gamma }{\alpha }}\mathrm{ln}l\gamma \mathrm{ln}(A/B))^{1\frac{1}{\alpha }}`$ (C4)
and that $`\varphi ^{\prime \prime }/\varphi _{}^{}{}_{}{}^{2}1/|x|^\alpha `$ thus the convergence to the Gumbell front holds. Note that the quantity $`\alpha _lm(l)\alpha l\gamma \mathrm{ln}l+O(1)`$ exhibits some universality.
One thus recovers the standard theorems for extremal value statistics reviewed in Appendix B, and the relation to the normalizing constants defined there as:
$`m(l)=a_N\alpha _l={\displaystyle \frac{1}{b_N}}l=\mathrm{ln}N`$ (C5)
In the Gaussian case, using the values given above one finds that (C4) indeed yields Eq. B5 in Appendix B (up to subdominant terms). Although the distribution is universal, the normalizing constants obviously depend on the details of the tail of the distribution. Note in all cases the presence of finite size corrections involving a logarithm.
There is a very small temperature ($`\beta _L\sqrt{\mathrm{ln}L}`$ for Gaussian) where the behaviour of $`G_0(x)`$ changes from (i) to (ii). It can be seen seen by rescaling temperature or equivalently disorder, with system size as in the REM.
Let us examine the case where the constants $`A_l`$ and $`B_l`$ are rescaled and now $`l`$-dependent (see also e.g. ). One can still use formulae (C4). Let us choose $`B_l=bl^{1+1/\alpha }`$ and $`A_l/B_l=cst`$ (which includes the Gaussian REM). One finds at $`T=0`$ that $`m(l)\frac{1}{b}(l\frac{\gamma }{\alpha ^2}\mathrm{ln}l\frac{\gamma }{\alpha }\mathrm{ln}(A/B))`$ and $`\alpha _lb\alpha `$. In the gaussian case $`\sigma _l=2\sigma l`$ one recovers the REM result:
$`m(l)\sqrt{\sigma }(2l{\displaystyle \frac{1}{2}}\mathrm{ln}(4\pi l))`$ (C6)
$`\alpha _l{\displaystyle \frac{1}{\sqrt{\sigma }}}`$ (C7)
at $`T=0`$ (i.e (46) setting $`ldl`$ and $`\sigma \sigma /d`$). The analysis can be performed at any $`T`$ and now yields a transition temperature when the behaviour of $`G_{0,l}(x)`$ at large $`x`$ changes.
## D An extended model
A richer phase diagram can be obtained by adding a logarithmic background potential $`V_0(𝐫)=J\mathrm{ln}\frac{|𝐫|}{a}`$ to the previous random potential $`\overline{(V_d(𝐫)V_d(𝐫^{}))^2}4\sigma \mathrm{ln}\frac{|𝐫𝐫^{}|}{a}`$ for $`a|𝐫𝐫^{}|L`$ and $`\overline{V_d(r)}=0`$ (i.e writing $`V(𝐫)=V_0(𝐫)+V_d(𝐫)`$) in (1). The choice of the origin breaks translational invariance. The competition between the disorder and the binding background potential (which if strong enough, tends to favor localizing the particle in wells far from $`𝐫=0`$) yields the phase diagram of Fig. 16. Another closely related model (model II) which preserves statistical translational invariance and has the same phase diagram is:
$`Z_v[V]=1+\left({\displaystyle \frac{L}{a}}\right)^{\beta J}{\displaystyle \underset{𝐫}{}}e^{\beta V_d(𝐫)}`$ (D1)
describes a problem with either zero or one particle (vortex) present, the energy cost of the vortex being $`J\mathrm{ln}(L/a)`$. It is thus a one vortex toy model of the recently studied XY model with random phase shifts .
In the absence of disorder the model with a background potential (model I) trivially exhibits a transition at $`\beta =\beta ^{}=d/J`$. At low temperature $`\beta >\beta ^{}`$ the Gibbs measure is $`p(𝐫)C(\frac{a}{r})^{\beta J}`$ with $`C=Z_{L=\mathrm{}}`$ a finite constant and the particle is bound to $`𝐫=0`$ (it has a finite probability to be within a fixed distance of $`𝐫=0`$). At high temperature $`\beta <\beta ^{}`$ the Gibbs measure becomes $`p(𝐫)(\frac{a}{L})^{d\beta J}(\frac{a}{r})^{\beta J}`$ and the particle is delocalized. This transition can be seen in the free energy density $`f=F/\mathrm{ln}L=T\mathrm{ln}Z/\mathrm{ln}L`$ since:
$`f=0\beta >\beta ^{}`$ (D2)
$`f=(J\beta ^{}\beta )\beta <\beta ^{}`$ (D3)
for $`\beta <\beta ^{}`$. This first order transition occurs as $`f`$ reaches its bound (since $`Z>1`$ due to the lattice cutoff, one has that $`f0`$). The model II has the same $`f`$ and a similar transition with either one vortex present $`\beta >\beta ^{}`$ or zero $`\beta <\beta ^{}`$.
In the presence of disorder the RG developed in this paper can be extended straightforwardly and leads to:
$`{\displaystyle \frac{1}{d}}_lG(x)={\displaystyle \frac{J}{d}}_xG+{\displaystyle \frac{\sigma }{d}}_x^2G+F[G]`$ (D4)
$`F[G]=G(1G)`$ (D5)
The additional term thus results in a simple shift in the front velocity. The position of the front $`m(l)`$ thus leads to the free energy $`f=m(l)/(dl)`$ which can have three distinct analytical forms:
$`\beta m(l)/ld\beta f(\beta )=`$ (D6)
$`(d+\sigma \beta ^2J\beta )\text{high T phase I}`$ (D7)
$`(2d{\displaystyle \frac{\beta }{\beta _c}}J\beta )\text{localized phase II}`$ (D8)
$`0\text{bound phase III}`$ (D9)
The phase diagram is represented in Fig 16 using the reduced temperature $`T/J`$ and the dimensionless disorder parameter $`\widehat{\sigma }=\sigma /J^2`$. For $`4d\sigma <J^2`$ one defines $`\beta ^{}(\sigma )=\frac{1}{2\sigma }(J\sqrt{J^24d\sigma })`$. The RG analysis yields three phases. In the model with the background potential (model I) they are, respectively:
the high temperature phase (for $`\beta <\beta _c`$ when $`4d\sigma >J^2`$ and for $`\beta <\beta ^{}(\sigma )`$ for $`4d\sigma <J^2`$): Entropy wins and the particle is delocalized over the system.
the localized phase (for $`\widehat{\sigma }>\widehat{\sigma }_c=1/(4d)`$ and $`\beta >\beta _c=\sqrt{d/\sigma }`$): The KPP velocity is frozen. The disorder wins and the particle freezes in wells away from the origin.
the bound phase (for $`\widehat{\sigma }<\widehat{\sigma }_c=1/(4d)`$ and $`\beta <\beta ^{}(\sigma )`$): The particule is bound to the origin. Within this phase near the phase boundaries (where the bound state length is large) a a crossover can be distinguished as a remnant of the freezing transition. The bound phase arises because of the bound $`f0`$ (or equivalently the velocity of the KPP equation must remain positive).
In model II the bound phase correspond to no vortex present. When one vortex is present it can be either localized in a few wells or in a high-T phase (as studied in the text of this paper).
Both transitions away from the bound phase are first order while the transition between high temperature phase and localized phase is continuous. An interesting feature is the multicritical point where the transition becomes continuous.
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# Contents
## 1 Introduction and Summary
BPS D-brane - anti-D-brane system of type IIA and IIB string theories admit tachyonic modes. It has been conjectured that at the minimum of the tachyon potential the tension of the brane-antibrane system is exactly canceled by the negative value of the tachyon potential, so that at this point the brane-antibrane system is indistinguishible from the vacuum. It has been further conjectured that various solitonic solutions on the brane-antibrane pair, where the tachyon approaches the minimum of the potential asymptotically, represent various lower dimensional branes. Thus for example, on a single D$`p`$-$`\overline{\mathrm{D}}p`$ brane system, the kink solution represents a non-BPS D-$`(p1)`$ brane of type II string theory, whereas a vortex solution represents a BPS D-$`(p2)`$ brane of the same theory. There are generalizations of this conjecture in which solitons of higher codimension on more than one pair of D-brane - $`\overline{\mathrm{D}}`$-brane system correspond to BPS and non-BPS D-branes of codimension $`>2`$.<sup>3</sup><sup>3</sup>3Recently evidence for some of these conjectures, and similar conjectures involving D-branes in bosonic string theories, have been found using string field theory. This approach uses the level truncation scheme developed by Kostelecky and Samuel. These conjectures have also been analysed using renormalization group flow on the world-sheet theory following earlier work of ref..
The codimension one case, i.e. the identification of the kink solution on the D$`p`$-brane $`\overline{\mathrm{D}}p`$-brane pair with a non-BPS D-$`(p1)`$-brane, has been demonstrated explicitly (although indirectly) by identifying a series of marginal deformations which interpolate between the D-brane - $`\overline{\mathrm{D}}`$-brane pair and the kink solution, and showing that this series of deformations take the original D$`p`$-brane $`\overline{\mathrm{D}}p`$-brane system to a non-BPS D-$`(p1)`$-brane. This is done by compactifying one of the directions tangential to the brane-antibrane system on a circle, switching on half unit of Wilson line on the antibrane, and reducing the radius of the circle to a critical radius where the lowest mode of the tachyon becomes marginal. One then uses this marginal deformation to create the kink and then takes the radius of the circle back to infinity.
It is natural to ask if this procedure can be generalised to the case of vortex and higher codimension solitons on the brane-antibrane pair. One faces the following problem for the vortex. The original D$`p`$-brane $`\overline{\mathrm{D}}p`$-brane system is neutral under Ramond-Ramond (RR) gauge fields since the RR charge of the brane and the anti-brane cancel. But a vortex, being identified to a BPS D-$`(p2)`$ brane, carries RR charge. Since RR charge is quantized, one cannot hope to have a continuous interpolation between these two configurations.<sup>4</sup><sup>4</sup>4It has been suggested by Maldacena that by switching on background magnetic field as in ref. one may be able to resolve this problem. The approach followed in might also be helpful. What one can hope to do however is to find a marginal deformation which interpolates between the original D$`p`$-brane $`\overline{\mathrm{D}}p`$-brane system and a vortex - antivortex pair on this system. If we can show that this marginal deformation converts the boundary conformal field theory (BCFT) associated with the D$`p`$-brane $`\overline{\mathrm{D}}p`$-brane system to the BCFT associated with the D-$`(p2)`$-brane $`\overline{\mathrm{D}}`$-$`(p2)`$-brane system, then we would establish the equivalence between a vortex solution and a D-$`(p2)`$-brane.
This is the problem we address in this paper. The steps involved in the analysis, which have been summarised in section 2, are more or less the same as the ones used for showing the equivalence of the kink solution on the brane-antibrane pair with a codimension one non-BPS brane. For convenience of notation we study the case of a vortex solution on a D2-brane $`\overline{\mathrm{D}}`$2-brane system. We compactify both directions tangential to the brane, switch on appropriate Wilson lines, and reduce the radii of the torus to certain critical values where the tachyonic deformation corresponding to the creation of the vortex-antivortex pair becomes a marginal deformation. We then switch on this marginal deformation and study the fate of the BCFT under this marginal deformation. These steps have been discussed in detail in section 3. In section 4 we study the effect of increasing the radii of the compact directions back to large values. We show that under this series of deformations the original BCFT gets deformed to a new BCFT describing the dynamics of open strings on a D0-$`\overline{\mathrm{D}}`$0 pair. This establishes the equivalence of a vortex solution on a D2-brane $`\overline{\mathrm{D}}`$2-brane pair, and a D0-brane. In section 5 we discuss generalization of this analysis to solutions of higher codimension.
## 2 The General Strategy
In this section we shall outline the general strategy that we shall follow in order to establish the equivalence between a vortex-antivortex pair on a BPS D$`p`$-brane - $`\overline{\mathrm{D}}p`$ brane system, and a D-$`(p2)`$ \- $`\overline{\mathrm{D}}`$-$`(p2)`$ brane pair. For definiteness we shall take our starting point to be a 2-brane - $`\overline{2}`$-brane pair of type IIA string theory, but the analysis can clearly be generalised to any $`p`$.
The steps are as follows:
1. We take a parallel 2-brane - $`\overline{2}`$-brane pair of type IIA string theory along $`x^1x^2`$ plane, and compactify the $`x^1`$ and $`x^2`$ directions on circles of radii $`R_1`$ and $`R_2`$ respectively. There are four different Chan-Paton (CP) sectors of open string states. States with CP factors proportional to the $`2\times 2`$ identity matrix $`I`$ and the Pauli matrix $`\sigma _3`$, representing open string states with both ends on the same brane, have the standard GSO projection under which the Neveu-Schwarz (NS) sector ground state is taken to be odd. On the other hand states with CP factors $`\sigma _1`$ and $`\sigma _2`$, representing open strings with two ends on two different branes, have opposite GSO projection, and in particular contain tachyonic modes. We shall identify the tachyonic modes associated with the sectors $`\sigma _1`$ and $`\sigma _2`$ with the real and the imaginary parts respectively of a complex tachyon field $`T`$.
2. We switch on half a unit of Wilson line along each of these circles. This makes open strings with CP factors $`\sigma _1`$ and $`\sigma _2`$, including the tachyon field $`T`$, antiperiodic along each of the two circles. Thus we can expand the tachyon field as
$$T(x^1,x^2,t)=\underset{m,nZ}{}T_{m+\frac{1}{2},n+\frac{1}{2}}(t)e^{i(m+\frac{1}{2})\frac{x^1}{R_1}+i(n+\frac{1}{2})\frac{x^2}{R_2}}.$$
(2.1)
The mass of the mode $`T_{m+\frac{1}{2},n+\frac{1}{2}}`$ is given by
$$(M_{m+\frac{1}{2},n+\frac{1}{2}})^2=\frac{(m+\frac{1}{2})^2}{(R_1)^2}+\frac{(n+\frac{1}{2})^2}{(R_2)^2}\frac{1}{2},$$
(2.2)
in the $`\alpha ^{}=1`$ unit that we shall be using.
3. From eq.(2.2) we see that for $`(R_1)^2+(R_2)^2>2`$ there are no tachyonic modes. For $`(R_1)^2+(R_2)^2=2`$ the modes $`T_{\pm \frac{1}{2},\pm \frac{1}{2}}`$ become marginal. We shall see in section 3 that at $`R_1=R_2=1`$ the deformation corresponding to
$`T(x^1,x^2)`$ $`=`$ $`i\alpha [e^{\frac{i}{2}(x^1+x^2)}e^{\frac{i}{2}(x^1+x^2)}+i(e^{\frac{i}{2}(x^1x^2)}e^{\frac{i}{2}(x^1x^2)})]`$ (2.3)
$`=`$ $`2\alpha (\mathrm{sin}{\displaystyle \frac{x^1+x^2}{2}}+i\mathrm{sin}{\displaystyle \frac{x^1x^2}{2}}),`$
for arbitrary constant $`\alpha `$ becomes exactly marginal, i.e. switching on this vev of the tachyon does not cost any energy. $`T(x^1,x^2)`$ has zeroes at $`x^1=x^2=0`$ and at $`x^1=x^2=\pi `$. Near $`x^1=x^2=0`$,
$$T(1+i)\alpha (x^1ix^2).$$
(2.4)
Thus it looks like a vortex.<sup>5</sup><sup>5</sup>5Of course whether we call this a vortex or an anti-vortex is a matter of convention. Once we fix the convention for the vortex, then a configuration with opposite orientation can be identified as an anti-vortex. On the other hand, near $`x^1=x^2=\pi `$,
$$T(1+i)\alpha ((x^1\pi )+i(x^2\pi )).$$
(2.5)
This looks like an anti-vortex. Thus we see that switching on a tachyon background of the form (2.3) corresponds to creating a vortex-antivortex pair.
4. After switching on the deformation (2.3) we take the radii $`R_1`$ and $`R_2`$ back to infinity, since we want to describe a vortex-antivortex pair on infinite D2-brane $`\overline{\mathrm{D}}`$2-brane system. Here we find that for a generic value of $`\alpha `$, once we increase $`R_1`$, $`R_2`$ beyond 1, there is a one point function of the modes $`T_{\pm \frac{1}{2},\pm \frac{1}{2}}`$. This is not surprising, since for $`R_i>1`$, the deformation (2.3) is a relevant perturbation, and hence switching on this background breaks conformal invariance of the world-sheet theory. However, we find that besides the trivial background $`\alpha =0`$, there is another inequivalent point (which we shall choose to be $`\alpha =1`$ by suitably normalizing $`\alpha `$) where the one point function of $`T_{\pm \frac{1}{2},\pm \frac{1}{2}}`$ vanish, and hence the theory is conformally invariant. We shall identify the $`\alpha =1`$ point as the vortex-antivortex pair.
As we shall show in section 4, increasing the $`x^1`$ and $`x^2`$ radii to arbitrary values $`R_1`$ and $`R_2`$ in the presence of $`\alpha =1`$ background can be described equivalently as increasing the radii of a pair of different coordinates $`\widehat{x}^1`$ and $`\widehat{x}^2`$. But in these coordinates, the original D2-$`\overline{\mathrm{D}}`$2 brane system, after being deformed by the tachyon background (2.3) with $`\alpha =1`$, appears as a D0-$`\overline{\mathrm{D}}`$0 brane system situated at diametrically opposite points of the torus spanned by $`\widehat{x}^1`$ and $`\widehat{x}^2`$. This shows that the spectrum of open strings in the background of a vortex-antivortex pair on the D2-$`\overline{\mathrm{D}}`$2 system wrapped on a torus with radii $`(R_1,R_2)`$ is identical to that on a D0-$`\overline{\mathrm{D}}`$0 brane system situated on a torus with the same radii. This establishes the equivalence between the vortex-antivortex pair on a D2-$`\overline{\mathrm{D}}`$2 brane system and D0-$`\overline{\mathrm{D}}`$0 brane pair.
## 3 Conformal Field Theory at the Critical Radii
In this section we shall study the marginal deformation of the BCFT on the upper half plane at the critical radii $`R_1=R_2=1`$ by (2.3). The relevant part of the BCFT at the critical radii before we switch on the perturbation (2.3) is described by a pair of scalar fields $`X^iX_L^i+X_R^i`$ for $`i=1,2`$ and their right- and left-moving fermionic superpartners $`\chi _R^i`$, $`\chi _L^i`$. The Neumann boundary condition satisfied by these fields on the real line is given by
$$X_L^i=X_R^i\frac{1}{2}X_B^i,\chi _L^i=\chi _R^i\chi _B^i.$$
(3.1)
We are considering here the NS sector open string states. In the Ramond (R) sector we have a different boundary condition on $`\chi ^i`$, but it can be handled in a manner similar to the one discussed in ref. and will not be discussed here.
Besides these fields we also have a time coordinate and its superpartners satisfying Neumann boundary condition and 7 other space-like coordinates and their superpartners satisfying Dirichlet boundary condition. We shall refer to these fields as spectator fields as they do not play a major role in the dynamics of the problem. We also have fermionic ghost fields $`b_L,c_L,b_R,c_R`$ and bosonic ghost fields $`\beta _L,\gamma _L,\beta _R,\gamma _R`$ satisfying Neumann boundary condition. We shall denote by $`\mathrm{\Phi }_L`$, $`\mathrm{\Phi }_R`$ the left- and right-moving bosonized ghost field of the $`\beta ,\gamma `$ system, satisfying the boundary condition
$$\mathrm{\Phi }_L=\mathrm{\Phi }_R\mathrm{\Phi }_B.$$
(3.2)
The tachyon vertex operator in the $`(1)`$ picture corresponding to the deformation (2.3) is given by:
$`V_T^{(1)}`$ $``$ $`e^{\mathrm{\Phi }_B}[(e^{\frac{i}{2}(X_B^1+X_B^2)}e^{\frac{i}{2}(X_B^1+X_B^2)})\sigma _1`$ (3.3)
$`+(e^{\frac{i}{2}(X_B^1X_B^2)}e^{\frac{i}{2}(X_B^1X_B^2)})\sigma _2].`$
In ‘zero’ picture this vertex operator takes the form:
$`V_T^{(0)}`$ $``$ $`[(\chi _B^1+\chi _B^2)(e^{\frac{i}{2}(X_B^1+X_B^2)}+e^{\frac{i}{2}(X_B^1+X_B^2)})\sigma _1`$ (3.4)
$`+(\chi _B^1\chi _B^2)(e^{\frac{i}{2}(X_B^1X_B^2)}+e^{\frac{i}{2}(X_B^1X_B^2)})\sigma _2].`$
Let us now define a new set of variables:
$`Y^1`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(X^1+X^2)`$
$`Y^2`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(X^1X^2)`$
$`\psi _R^1`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(\chi _R^1+\chi _R^2)`$
$`\psi _R^2`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(\chi _R^1\chi _R^2)`$
$`\psi _L^1`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(\chi _L^1+\chi _L^2)`$
$`\psi _L^2`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(\chi _L^1\chi _L^2).`$ (3.5)
These fields satisfy the boundary conditions:
$$Y_L^i=Y_R^i\frac{1}{2}Y_B^i,\psi _L^i=\psi _R^i\psi _B^i,(i=1,2)$$
(3.6)
on the real line. In terms of these fields,
$$V_T^{(0)}\left[\psi _B^1(e^{\frac{i}{\sqrt{2}}Y_B^1}+e^{\frac{i}{\sqrt{2}}Y_B^1})\sigma _1+\psi _B^2(e^{\frac{i}{\sqrt{2}}Y_B^2}+e^{\frac{i}{\sqrt{2}}Y_B^2})\sigma _2\right].$$
(3.7)
We now fermionize the scalar fields $`Y^i`$ as follows:
$$e^{i\sqrt{2}Y_R^i}=\frac{1}{\sqrt{2}}(\xi _R^i+i\eta _R^i)\tau _i,e^{i\sqrt{2}Y_L^i}=\frac{1}{\sqrt{2}}(\xi _L^i+i\eta _L^i)\tau _i,$$
(3.8)
where $`\xi _R^i,\eta _R^i`$ ($`\xi _L^i`$, $`\eta _L^i`$) are right- (left-) moving Majorana-Weyl fermions, and the Pauli matrices $`\tau _i`$ denote cocycle factors which must be put in to guarantee correct (anti-)commutation relations between various fields.<sup>6</sup><sup>6</sup>6As we shall be using these bosonization formulae to manipulate vertex operators in the BCFT on the upper half plane, we only need to require that correct (anti-)commutation relations are satisfied by the vertex operators subject to the boundary condition (3.6). We also attach a cocycle factor $`\tau _3`$ to $`\psi _{L,R}^i`$ and all the spectator fermions. This guarantees for example that $`\psi ^i`$ and the spectator fermions commute with both sides of eq.(3.8). The cocycle factors should be taken to commute with CP factors.
We can find another representation of the same conformal field theory by rebosonising the fermions as follows:
$$\frac{1}{\sqrt{2}}(\xi _R^i+i\psi _R^i)=e^{i\sqrt{2}\varphi _R^i}\stackrel{~}{\tau }_i,\frac{1}{\sqrt{2}}(\xi _L^i+i\psi _L^i)=e^{i\sqrt{2}\varphi _L^i}\stackrel{~}{\tau }_i.$$
(3.9)
$`\varphi ^1`$, $`\varphi ^2`$ represent a pair of free scalar fields, and $`\stackrel{~}{\tau }_i`$ are a new set of cocycle factors. In this convention $`\eta _{L,R}^1`$, $`\eta _{L,R}^2`$ and spectator fermions carry the new cocycle factor $`\stackrel{~}{\tau }_3`$. There is a third representation in which we use a slightly different rebosonization:
$$\frac{1}{\sqrt{2}}(\eta _R^i+i\psi _R^i)=e^{i\sqrt{2}\varphi _R^i}\widehat{\tau }_i,\frac{1}{\sqrt{2}}(\eta _L^i+i\psi _L^i)=e^{i\sqrt{2}\varphi _L^i}\widehat{\tau }_i,$$
(3.10)
where $`\varphi ^i`$ is another pair of scalar fields, and $`\widehat{\tau }_i`$ denote another set of cocycle factors. $`\xi _{L,R}^1`$, $`\xi _{L,R}^2`$ and the spectator fermions will carry the cocycle factor $`\widehat{\tau }_3`$ when we use the set of variables $`\varphi ^i`$, $`\xi ^i`$ and the spectator fields to describe the BCFT. We also define
$$\varphi ^i=\varphi _L^i+\varphi _R^i,\varphi ^i=\varphi _L^i+\varphi _R^i.$$
(3.11)
For later use we list here the operator product expansions, and the relations between the currents of free fermions and bosons:
$$\psi _R^i(z)\psi _R^j(w)\xi _R^i(z)\xi _R^j(w)\eta _R^i(z)\eta _R^j(w)\frac{i}{zw}\delta _{ij},$$
(3.12)
$$Y_R^i(z)Y_R^j(w)\varphi _R^i(z)\varphi _R^j(w)\varphi _R^i(z)\varphi _R^j(w)\frac{1}{2(zw)^2}\delta _{ij},$$
(3.13)
$$\psi _R^i\xi _R^i=i\sqrt{2}\varphi _R^i,\eta _R^i\xi _R^i=i\sqrt{2}Y_R^i,\psi _R^i\eta _R^i=i\sqrt{2}\varphi _R^i,$$
(3.14)
with no summation over $`i`$ in the last equation. Here $``$ denotes equality up to non-singular terms. There are also similar relations involving the left-moving currents.
Using eq.(3.8) the boundary condition (3.6) on $`Y^i`$ can be translated to the following boundary condition on the fermions:
$$\xi _L^i=\xi _R^i\xi _B^i,\eta _L^i=\eta _R^i\eta _B^i.$$
(3.15)
Combining these with the boundary condition (3.6) on $`\psi ^i`$, we see from (3.9), (3.10) that $`\varphi ^i`$ and $`\varphi ^i`$ both satisfy Neumann boundary condition on the real line
$$\varphi _L^i=\varphi _R^i\varphi _B^i/2,\varphi _L^i=\varphi _R^i\varphi _B^i/2.$$
(3.16)
Using the bosonization relations (3.8)-(3.10), (3.14), the vertex operator $`V_T^{(0)}`$ given in eq.(3.7) can be expressed as
$$V_T^{(0)}[\psi _B^1\xi _B^1\sigma _1\tau _2\psi _B^2\xi _B^2\sigma _2\tau _1][\varphi _B^1\sigma _1\tau _2\varphi _B^2\sigma _2\tau _1].$$
(3.17)
$``$ denotes tangential derivative along the boundary. The two operators appearing in the right hand side of eq.(3.17) correspond to vertex operators of constant U(1) gauge field along $`\varphi ^1`$ and $`\varphi ^2`$ directions, with generators $`\sigma _1\times \tau _2`$ and $`\sigma _2\times \tau _1`$ respectively. Since these two matrices commute, we see that switching on the tachyon vev corresponds to switching on a pair of commuting U(1) Wilson lines $`A_{\varphi _1}`$ and $`A_{\varphi _2}`$ along the bosonic directions $`\varphi ^1`$ and $`\varphi ^2`$. This represents an exactly marginal deformation. Switching on finite vacuum expectation value (vev) of the tachyon then corresponds to inserting the operator
$$\mathrm{exp}(\frac{i\alpha }{2\sqrt{2}}\varphi _B^1\sigma _1\tau _2\frac{i\alpha }{2\sqrt{2}}\varphi _B^2\sigma _2\tau _1),$$
(3.18)
where $``$ denotes integration along the boundary. Note that we have fixed the normalization of $`\alpha `$ in a specific manner. This is the same normalization convention as in ref., and so we shall be able to use the results of ref. directly. With this normalization, $`\alpha `$ is a periodic variable with period 2.
Before we study the effect of switching on such a tachyon vev on the spectrum of open strings, let us study the open string spectrum before switching on the tachyon vev. Let us denote by $``$ the Hilbert space of states<sup>7</sup><sup>7</sup>7Here we are discussing open string states; so all states are created from $`|0`$ by vertex operators inserted at the boundary of the upper half plane. created by the half-integer moded $`\psi ^i`$, $`\xi ^i`$, $`\eta ^i`$ oscillators and the spectator fields on the $`SL(2,R)`$ invariant vacuum $`|0`$. Alternatively we can also view $``$ as the Hilbert space of states created from $`|0`$ by the spectator fields, together with vertex operators involving $`\psi ^i`$ and $`Y^i`$ with $`Y^i`$ momenta quantized in units of $`1/\sqrt{2}`$, or vertex operators involving $`\eta ^i`$ and $`\varphi ^i`$ with $`\varphi ^i`$ momenta quantized in units of $`1/\sqrt{2}`$, or vertex operators involving $`\xi ^i`$ and $`\varphi ^i`$ with $`\varphi ^i`$ momenta quantized in units of $`1/\sqrt{2}`$. On $``$, we denote by $`(1)^F`$ the world-sheet fermion number which changes the sign of $`\psi ^1`$, $`\psi ^2`$, and the spectator fermions, leaving $`\xi ^i`$, $`\eta ^i`$ (and hence $`Y^i`$) and the spectator bosons unchanged. The SL(2,R) invariant ground state $`|0`$ is taken to be odd under $`(1)^F`$. We also define the transformations $`h_1`$ and $`h_2`$ as follows. Both $`h_1`$ and $`h_2`$ leave $`\psi ^1`$, $`\psi ^2`$ and all the spectator fermions unchanged, but $`h_1`$ changes the sign of $`\xi ^1,\eta ^1`$ and $`h_2`$ changes the sign of $`\xi ^2,\eta ^2`$. In terms of the bosonic variables $`X^i`$ or $`Y^i`$,
$`h_1`$ $`:`$ $`X_{L,R}^1X_{L,R}^1+{\displaystyle \frac{\pi }{2}},X_{L,R}^2X_{L,R}^2+{\displaystyle \frac{\pi }{2}},`$
$`Y_{L,R}^1Y_{L,R}^1+{\displaystyle \frac{\pi }{\sqrt{2}}},Y_{L,R}^2Y_{L,R}^2`$
$`h_2`$ $`:`$ $`X_{L,R}^1X_{L,R}^1+{\displaystyle \frac{\pi }{2}},X_{L,R}^2X_{L,R}^2{\displaystyle \frac{\pi }{2}},`$
$`Y_{L,R}^1Y_{L,R}^1,Y_{L,R}^2Y_{L,R}^2+{\displaystyle \frac{\pi }{\sqrt{2}}}.`$
It is easy to verify that with this definition, $`(h_1)^2`$ and $`(h_2)^2`$ act as identity on all states. States with CP factor $`I`$ and $`\sigma _3`$ are $`h_1h_2`$ even and $`(1)^F`$ even, whereas states with CP factors $`\sigma _1`$ and $`\sigma _2`$ are $`h_1h_2`$ odd and $`(1)^F`$ odd. Taking into account the assignment of the cocycle factors ($`\tau _1`$ for $`h_1`$ odd states, $`\tau _2`$ for $`h_2`$ odd states and $`\tau _3`$ for $`(1)^F`$ odd states) the quantum numbers carried by various open string states, when expressed in terms of $`\psi ^i,\eta ^i,\xi ^i`$ and the spectator fields, are as given in table 1.<sup>8</sup><sup>8</sup>8There will be extra cocycle factors involving $`\stackrel{~}{\tau }_i`$ when we express the vertex operators in terms of the fields $`\varphi ^i`$, $`\eta ^i`$ and the spectator fields. But since these cocycle factors commute with those carried by the gauge fields $`A_{\varphi ^1}`$ and $`A_{\varphi ^2}`$, we can ignore them for the purpose of studying how the spectrum changes when we switch on the gauge fields. The complete spectrum of open string states carrying a given CP$``$cocycle factor is obtained by keeping all states in $``$ carrying the quantum numbers mentioned in the table. The last two columns describe if the open string state is charged or neutral under the gauge fields $`A_{\varphi ^1}`$ or $`A_{\varphi ^2}`$. This is easily determined by examining if the cocycles factors commute with $`\sigma _1\tau _2`$ and $`\sigma _2\tau _1`$, since these represent the U(1) generators associated with the background gauge fields.
From table 1 we note that before we switch on the background (3.18), each state is invariant under $`(1)^Fh_1h_2`$, and that all combinations of the quantum numbers subject to this condition appear exactly twice in this table. Thus the combined spectrum from all the sectors contain two copies of $``$ with $`(1)^Fh_1h_2`$ projection. We also note that open string states carrying CP$``$cocycle factors $`II`$, $`\sigma _3\tau _3`$, $`\sigma _1\tau _2`$ and $`\sigma _2\tau _1`$ are all neutral under both gauge fields. Combining the spectrum from all four neutral sectors we see that each combination of quantum numbers (subject to the condition $`(1)^Fh_1h_2=+1`$) appears exactly once; thus the combined spectrum of charge neutral states contain a single copy of $``$ with $`(1)^Fh_1h_2`$ projection. The combined spectrum of the charged states also contains a single copy of $``$ with $`(1)^Fh_1h_2`$ projection.
Let us now define a new set of symmetry generators $`(1)^{F_\varphi }`$, $`h_{\varphi ^1}`$ and $`h_{\varphi ^2}`$ on $``$ by using the representation where $``$ is generated by action on $`|0`$ by the vertex operator involving $`\varphi ^i`$, $`\eta ^i`$ and the spectator fields. $`(1)^{F_\varphi }`$ changes the signs of $`\eta ^i`$ and the spectator fermions, leaving unchanged $`\varphi ^i`$ and the spectator bosons, and has eigenvalue $`1`$ acting on the SL(2,R) invariant vacuum $`|0`$. $`h_{\varphi ^i}`$ leaves unchanged $`\eta ^1`$, $`\eta ^2`$ and all the spectator fields, and transform $`\varphi ^j`$ as follows:
$`h_{\varphi ^1}`$ $`:`$ $`\varphi _{L,R}^1\varphi _{L,R}^1+{\displaystyle \frac{\pi }{\sqrt{2}}},\varphi _{L,R}^2\varphi _{L,R}^2,`$
$`h_{\varphi ^2}`$ $`:`$ $`\varphi _{L,R}^2\varphi _{L,R}^2+{\displaystyle \frac{\pi }{\sqrt{2}}},\varphi _{L,R}^1\varphi _{L,R}^1.`$ (3.21)
Thus $`h_{\varphi ^i}`$ changes the sign of $`\xi ^i`$ and $`\psi ^i`$. With this definition, and using the bosonization relations (3.8), (3.9), it is easy to see that
$$(1)^Fh_1h_2=(1)^{F_\varphi }h_{\varphi ^1}h_{\varphi ^2}.$$
(3.22)
(Both sides change the sign of $`\xi ^{1,2}`$, $`\eta ^{1,2}`$ and $`\psi ^{1,2}`$, and the spectator fermions, leaving the spectator bosons unchanged.) Thus the combined spectrum of the charge neutral states as well as the combined spectrum of the charged states may be identified as a copy of $``$, subject to the $`(1)^{F_\varphi }h_{\varphi ^1}h_{\varphi ^2}`$ projection.
Upon switching on the gauge fields, the spectrum in the charge neutral sector does not change, but the spectrum in the charged sector changes since the quantization laws for the $`\varphi ^i`$ momenta $`p_{\varphi ^i}`$ change. By combining the fields into charge eigenstates we can follow the change in $`p_{\varphi _1}`$, $`p_{\varphi _2}`$ as a function of $`\alpha `$ and thus completely determine the spectrum at every value of $`\alpha `$.<sup>9</sup><sup>9</sup>9With the normalization of $`\alpha `$ chosen here, $`\alpha =2`$ corresponds to a shift in $`p_{\varphi ^i}`$ by $`\pm \sqrt{2}`$. This corresponds to states carrying same $`h_{\varphi ^1}`$ and $`h_{\varphi ^2}`$ quantum numbers. Hence $`\alpha `$ and $`\alpha +2`$ describes the same BCFT. However as we shall see in the next section, once we deform the radius away from $`R_1=R_2=1`$ point, $`\alpha =0`$ and $`1`$ are the only inequivalent points which give conformally invariant theories. Since $`\alpha =0`$ represents the trivial tachyon background, we shall be interested in the $`\alpha =1`$ point. The spectrum at $`\alpha =1`$ simplifies enormously if we notice that this corresponds to shifting the $`p_{\varphi ^1}`$ and $`p_{\varphi ^2}`$ quantization laws by $`\pm \frac{1}{\sqrt{2}}`$, so that its effect is to simply reverse the sign of the $`h_{\varphi ^1}`$ and $`h_{\varphi ^2}`$ quantum numbers. But the initial spectrum did not contain $`h_{\varphi ^1}`$ and $`h_{\varphi ^2}`$ projections individually. It only contained $`(1)^{F_\varphi }h_{\varphi ^1}h_{\varphi ^2}`$ projection, and this does not change. As a result the spectrum at $`\alpha =1`$ is identical to the spectrum at $`\alpha =0`$!
It may appear from this that at the end of the deformation the system has come back to the original system! However, as we shall see in section 4, the response of this system to a change in the radii $`R_1`$ and $`R_2`$ is very different from that of the original system. In order to facilitate the analysis of section 4, we introduce dual coordinates
$$\widehat{X}_L^i=X_L^i,\widehat{X}_R^i=X_R^i.$$
(3.23)
In terms of the new coordinates $`\widehat{X}^i`$, the BCFT at $`\alpha =1`$ corresponds to a D0-brane $`\overline{\mathrm{D}}`$0-brane pair, situated at diametrically opposite points of a square torus with unit radii.
## 4 Deforming Away from the Critical Radius
In this section we shall consider the effect of switching on the perturbation that deforms the radii away from their critical values. The procedure followed here will be similar to the one used in ref., so our discussion will be brief. There are four possible marginal deformations of the bulk conformal field theory, three of which correspond to deformation of the shape and size of the torus labelled by the $`X^1`$-$`X^2`$ coordinates, and the fourth one corresponds to switching on the anti-symmetric tensor field in the $`X^1`$-$`X^2`$ plane. Using eqs.(3) we can express a general perturbation of this kind in the (0,0) picture as
$$K_{ij}X_L^iX_R^j=a_{ij}Y_L^iY_R^j,$$
(4.1)
where
$$a=\left(\begin{array}{cc}\frac{1}{\sqrt{2}}& \frac{1}{\sqrt{2}}\\ \frac{1}{\sqrt{2}}& \frac{1}{\sqrt{2}}\end{array}\right)K\left(\begin{array}{cc}\frac{1}{\sqrt{2}}& \frac{1}{\sqrt{2}}\\ \frac{1}{\sqrt{2}}& \frac{1}{\sqrt{2}}\end{array}\right).$$
(4.2)
First we shall consider the effect of first order perturbation, and show that in the presence of this perturbation the tachyon vertex operator $`V_T`$ develops a one point function unless $`\alpha =0`$ or $`\alpha =1`$. The procedure for doing this is similar to that discussed in ref.. We insert a tachyon vertex operator $`V_T^{(0)}`$ given in eq.(3.17) at a point on the boundary of the disk (or upper half plane), the background (3.18) at the boundary of the disk, and a closed string vertex operator corresponding to the perturbation (4.1) in the $`(1,1)`$ picture at the center of the disk. This is proportional to
$$a_{ij}e^{\mathrm{\Phi }_L}e^{\mathrm{\Phi }_R}\psi _L^i\psi _R^j,$$
(4.3)
where $`\mathrm{\Phi }_L`$ and $`\mathrm{\Phi }_R`$ are the left and right-moving bosonized ghost fields. Computation of this amplitude is straightforward using the description of the BCFT in terms of the $`(\eta ^i,\varphi ^i)`$ fields. The $`\varphi ^i`$ momentum conservation laws tell us that only the $`i=j`$ terms in (4.3) contribute. The answer is proportional to
$$(a_{11}+a_{22})\mathrm{sin}(\pi \alpha ).$$
(4.4)
This shows that this one point function vanishes only for $`\alpha =0`$ and $`1`$ (mod 2). It is also easy to check that at these two values of $`\alpha `$ the one point function of all other open string vertex operators also vanish, and hence these configurations describe consistent boundary conformal field theories.
Next we need to study what happens when we go beyond the lowest order perturbation in the deformation parameters $`a_{ij}`$. For this let us consider an amplitude with an arbitrary number of insertions of the operator (4.1) at various points in the interior of the disk, insertion of the background (3.18) with $`\alpha =1`$ at the boundary of the disk, and a set of open string vertex operators, whose correlation function we want to calculate, at the boundary of the disk. We take two of these open string vertex operators to be in the $`1`$ picture and the rest in the zero picture so that the total picture number of all the vertex operators add up to $`2`$. Again the calculation proceeds as in ref.. We express all the operators in terms of $`\eta ^i`$, $`\varphi ^i`$, and the spectator fields. The vertex operator (4.1) can be expressed as products of $`\eta _L^i`$, $`\eta _R^i`$, $`e^{\pm i\sqrt{2}\varphi _L^i}`$ and $`e^{\pm i\sqrt{2}\varphi _R^i}`$ fields. The integrated vertex operator (3.18) inserted at the boundary has two effects. It effectively shifts the $`\varphi ^i`$ momenta of the open string vertex operators as discussed in section 3, and $`\varphi _B^i`$ picks up the $`\varphi ^i`$ winding number of all the closed string vertex operators inserted in the interior of the disk. Using the operator product expansion (3.13) it is easy to verify that $`\frac{1}{2\sqrt{2}}\varphi _B^i=\frac{1}{\sqrt{2}}\varphi _R^i`$, acting on any combination of closed string vertex operators, gives an integral multiple of $`\pi `$. Now, since
$$\mathrm{exp}(i\pi n\sigma _1\tau _2)=\mathrm{exp}(i\pi n),\mathrm{exp}(i\pi n\sigma _2\tau _1)=\mathrm{exp}(i\pi n),\text{for integer}n$$
(4.5)
we see that in taking into account the effect of (3.18) for $`\alpha =1`$ on the closed string vertex operators we can replace it by
$$\mathrm{exp}(\frac{i}{2\sqrt{2}}\varphi _B^1\frac{i}{2\sqrt{2}}\varphi _B^2)=\mathrm{exp}(\frac{i}{\sqrt{2}}\varphi _R^1\frac{i}{\sqrt{2}}\varphi _R^2).$$
(4.6)
In going from the left hand side to the right hand side of eq.(4.6) we have used the boundary condition (3.16). Since $`\varphi _R^i`$ are holomorphic fields, we can now deform the integration contour into the interior of the disk, picking up residues from the location of the closed string vertex operators. The net effect is to replace each insertion of (4.1) by
$$\mathrm{exp}\left(\frac{i}{\sqrt{2}}\varphi _R^1\frac{i}{\sqrt{2}}\varphi _R^2\right)a_{ij}Y_L^iY_R^j,$$
(4.7)
where the contours of integration in the exponent are around the location of the closed string vertex operator. This can be easily evaluated using eqs.(3.9)-(3.14), and the result is<sup>10</sup><sup>10</sup>10Here we have differed somewhat from the analysis of ref.. In $`\varphi _B`$ in the exponent was replaced by $`(\varphi _L+\varphi _R)`$ instead of $`2\varphi _R`$. If we follow this procedure here, then eq.(4.8) will be replaced by $`a_{ij}\varphi _L^i\varphi _R^j`$. This would be analogous to the corresponding result in . The final result for the spectrum and correlation functions however does not depend on which procedure we follow, since the two procedures are related by the symmetry transformation $`(\varphi _L^i,\varphi _R^i)(\varphi _L^i+\frac{\pi }{2\sqrt{2}},\varphi _R^i+\frac{\pi }{2\sqrt{2}})`$.
$$a_{ij}Y_L^iY_R^j=K_{ij}X_L^iX_R^j.$$
(4.8)
In terms of the coordinates $`\widehat{X}^i`$ introduced in eq.(3.23), the right hand side of eq.(4.8) can be expressed as
$$K_{ij}\widehat{X}_L^i\widehat{X}_R^j.$$
(4.9)
This has the same form as the left hand side of (4.1) with $`X^i`$ replaced by $`\widehat{X}^i`$. Thus deforming the $`X^i`$ radius from 1 to $`R_i`$ in the presence of the tachyon background (3.18) with $`\alpha =1`$ corresponds to deforming the $`\widehat{X}^i`$ radius from 1 to $`R_i`$. The D-brane system now describes a D0-brane $`\overline{\mathrm{D}}`$0-brane pair situated at diametrically opposite points of this torus.
This shows that the spectrum of open string states living on a vortex antivortex pair at the diametrically opposite points on a D2-$`\overline{\mathrm{D}}`$2 system wrapped on a torus is identical to the spectrum of open string states living on a D0-$`\overline{\mathrm{D}}`$0 brane pair at the diametrically opposite points of the same torus. This leads to the identification of the vortex (antivortex) on a D2-$`\overline{\mathrm{D}}`$2 system with a D0 ($`\overline{\mathrm{D}}`$0) brane.
## 5 Generalizations
The method used here can easily be generalized to prove the identification of a vortex solution on a D$`p`$-$`\overline{\mathrm{D}}p`$ brane pair with a D-$`(p2)`$ brane. The analysis is exactly identical; all that is required is to change the Dirichlet boundary condition on $`(p2)`$ of the spectator superfields to Neumann boundary condition.
It is also possible to generalize this analysis to show that a codimention $`2n`$ soliton on $`2^{n1}`$ pairs of D$`p`$-$`\overline{\mathrm{D}}p`$ brane system represents a D$`(p2n)`$ brane. Again for simplicity let us assume that all the spectator fields except the time direction has Dirichlet boundary condition, i..e. take $`p=2n`$. We compactify each of the $`2n`$ directions tangential to the D-branes on a circle of radius 1. Let us label these coordinates by $`x^1,\mathrm{}x^{2n}`$. The open string states carry $`2^n\times 2^n`$ CP factors; with the tachyon states carrying off diagonal CP factors of the form:
$$\left(\begin{array}{cc}0& A\\ A^{}& 0\end{array}\right),$$
(5.1)
where $`A`$ is a $`2^{n1}\times 2^{n1}`$ complex matrix. Under the $`U(2^{n1})\times U(2^{n1})`$ gauge transformation on the brane-antibrane system, generated by $`U(2^{n1})`$ matrices $`U`$ and $`V`$,
$$AUAV^{}.$$
(5.2)
Let us now, following ref., choose a $`2^n\times 2^n`$ dimensional representation of the $`2n`$ dimensional Clifford algebra in which each of the $`\mathrm{\Gamma }`$-matrices has the form given in (5.1). Let $`X^i`$ ($`1i2n`$) denote the coordinate fields tangential to the D-brane, and $`\chi _R^i`$, $`\chi _L^i`$ be their right and left-moving superpartners. Let us define
$`Y^{2k+1}={\displaystyle \frac{1}{\sqrt{2}}}(X^{2k+1}+X^{2k+2}),Y^{2k+2}={\displaystyle \frac{1}{\sqrt{2}}}(X^{2k+1}X^{2k+2}),`$
$`\psi _R^{2k+1}={\displaystyle \frac{1}{\sqrt{2}}}(\chi _R^{2k+1}+\chi _R^{2k+2}),\psi _R^{2k+2}={\displaystyle \frac{1}{\sqrt{2}}}(\chi _R^{2k+1}\chi _R^{2k+2}),`$
$`\psi _L^{2k+1}={\displaystyle \frac{1}{\sqrt{2}}}(\chi _L^{2k+1}+\chi _L^{2k+2}),\psi _L^{2k+2}={\displaystyle \frac{1}{\sqrt{2}}}(\chi _L^{2k+1}\chi _L^{2k+2}),`$
for $`0k(n1)`$. Consider now the following vertex operator in the $`(1)`$ picture
$$V_T^{(1)}e^{\mathrm{\Phi }_B}\underset{i=1}{\overset{2n}{}}\mathrm{sin}\frac{Y_B^i}{\sqrt{2}}\mathrm{\Gamma }_i,$$
(5.4)
where the subscript $`B`$ denotes the boundary value as usual. In order that this is an allowed vertex operator, we need to switch on appropriate Wilson lines on the branes; we assume that this has been done. (We can, for example, take the Wilson line along the $`X^{2k+1}`$ and the $`X^{2k+2}`$ direction to be $`\mathrm{\Gamma }^{2k+1}\mathrm{\Gamma }^{2k+2}`$. Conjugation by this matrix changes the sign of $`\mathrm{\Gamma }^{2k+1}`$ and $`\mathrm{\Gamma }^{2k+2}`$, keeping the other $`\mathrm{\Gamma }^i`$’s invariant. Unlike in the case analysed earlier, this amounts to switching on Wilson lines both on the D-branes and the anti-D-branes.) In the zero picture the vertex operator (5.4) takes the form:
$$V_T^{(0)}\underset{i=1}{\overset{2n}{}}\psi _B^i\mathrm{cos}\frac{Y_B^i}{\sqrt{2}}\mathrm{\Gamma }_i.$$
(5.5)
It is easy to see via the bosonization procedure that each term in the sum represents switching on a constant gauge field (Wilson line) along a new bosonic coordinate and hence is a marginal deformation. Furthermore, since all the $`\mathrm{\Gamma }`$ matrices anti-commute with each other, and the fermions $`\psi _B^i`$ also anti-commute with each other, we see that the different terms in the sum commute with each other. After bosonization this will imply that the CP$``$cocycle factors carried by the different gauge fields commute with each other. Thus (5.5) corresponds to switching on constant, commuting gauge fields along different directions, and represents a marginal deformation.
In analogy with the analysis of sections 3, 4, one expects that the BCFT at the end of the marginal deformation (the $`\alpha =1`$ point) is more naturally described in terms of a set of dual coordinates $`\widehat{x}^i`$ in which the system appears as $`2^{n1}`$ D0-$`\overline{\mathrm{D}}`$0 brane pair. If we increase the radii of the original torus to arbitrary values $`R_i`$ after the marginal deformation, then it would correspond to increasing the radii of the torus described by the $`\widehat{x}^i`$ coordinates to the same values $`R_i`$ as in section 4. Thus the end result will be $`2^{n1}`$ D0-$`\overline{\mathrm{D}}`$0 brane pairs situated at different point on the torus with radii $`R_i`$.<sup>11</sup><sup>11</sup>11A consistency check for this scenario is that at the critical radii the total mass of the initial configuration is given by $`2^n/g`$, where $`g`$ denotes the closed string coupling constant, since each of the $`2^n`$ D-$`2n`$-branes / $`\overline{\mathrm{D}}`$-$`2n`$-branes has area equal to $`(2\pi )^{2n}`$ and tension equal to $`1/((2\pi )^{2n}g)`$. This agrees with the total mass of $`2^{n1}`$ D0-$`\overline{\mathrm{D}}`$0 brane pair. Also both the initial and the final state has vanishing RR charge.
On the other hand, we can examine the tachyon background (5.4) and try to identify it as a collection of solitons. For this we need to identify the soliton cores as the places where the tachyon field vanishes. This requires each $`y^i`$ to be an integral multiple of $`\pi \sqrt{2}`$. Using eqs.(LABEL:e5.3), and that each $`x^i`$ describes a circle of radius 1, we see that the inequivalent points are given by all possible combination of the configurations
$$(x^{2k+1},x^{2k+2})=(0,0)\text{or}(\pi ,\pi ).$$
(5.6)
Since there are $`n`$ pairs of coordinates, this gives $`2^n`$ possibilities. To study the nature of the solution near the core, let us consider one of them, e.g. $`(x^1,\mathrm{}x^{2n})=(0,\mathrm{}0)`$. The tachyon background near this point is proportional to
$$y^i\mathrm{\Gamma }_i.$$
(5.7)
This is precisely the form suggested in ref. for a codimension $`2n`$ soliton on $`2^{n1}`$ brane-antibrane pairs. On the other hand, expanding the tachyon field near the point $`(x^1,x^2\mathrm{}x^{2n})=(\pi ,\pi ,0,0,\mathrm{}0)`$, we get the tachyon background to be proportional to:
$$(y^1\pi \sqrt{2})\mathrm{\Gamma }^1+\underset{i=2}{\overset{2n}{}}y^i\mathrm{\Gamma }^i.$$
(5.8)
The $``$ sign in front of the first term indicates that it has opposite orientation compared to (5.7) and represents an anti-soliton. In general the solution around a point (5.6) represents a soliton (anti-soliton) if the number of coordinate pairs taking the value $`(\pi ,\pi )`$ is even (odd). Thus the tachyon background (5.5) represents creation of $`2^{n1}`$ soliton - antisoliton pairs. From this we can conclude that $`2^{n1}`$ soliton - antisoliton pairs represent $`2^{n1}`$ D0-$`\overline{\mathrm{D}}`$0 pairs, and hence a single D0-brane should be identified with the single codimension $`2n`$ soliton on the D-$`2n`$-brane $`\overline{\mathrm{D}}`$-$`2n`$-brane pair.
Generalization of this result to solitons of odd codimension, which are expected to describe non-BPS D-branes is also straightforward. In this case we take $`2^{n1}`$ D-$`(2n1)`$-brane $`\overline{\mathrm{D}}`$-$`(2n1)`$-brane pairs in type IIB string theory along $`x^1,\mathrm{}x^{2n1}`$, with $`x^1,\mathrm{}x^{2n2}`$ directions compactified on circles of radii 1 and $`x^{2n1}`$ direction compactified on a circle of radius $`\frac{1}{\sqrt{2}}`$. We define $`Y^i`$’s for $`1i(2n2)`$ as in eq.(LABEL:e5.3) and consider tachyon background associated with the vertex operator<sup>12</sup><sup>12</sup>12This requires choosing the Wilson lines along $`x^1,\mathrm{}x^{2n1}`$ as before. Since the Wilson line along $`x^{2n1}`$ requires a $`\mathrm{\Gamma }_{2n}`$ we see that we need a representation of the $`2n`$ dimensional Clifford algebra even though we have only $`(2n1)`$ coordinates.
$$V_T^{(1)}e^{\mathrm{\Phi }_B}\left[\underset{i=1}{\overset{2n2}{}}\mathrm{sin}\frac{Y_B^i}{\sqrt{2}}\mathrm{\Gamma }_i+\mathrm{sin}\frac{X_B}{\sqrt{2}}\mathrm{\Gamma }_{2n1}\right],$$
(5.9)
This corresponds to creation of $`2^{n1}`$ solitons at
$`(x^{2k+1},x^{2k+2})=(0,0)\text{or}(\pi ,\pi )\text{for}0k(n2),`$
$`x^{2n1}=0.`$ (5.10)
On the other hand it is easy to see that the vertex operator (5.9), in the zero picture, represents a marginal deformation. We expect this to convert the original brane configuration to a configuration of non-BPS branes as in ref.. The number of such D0-branes can easily be seen to be $`2^{n1}`$ by comparing the masses of the initial and the final configurations. This shows the equivalence between the non-BPS D0-brane and a soliton on $`2^{n1}`$ D-$`(2n1)`$ $``$ $`\overline{\mathrm{D}}`$-$`(2n1)`$ brane pairs.
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# Deep X-Ray Observations of Supernova Remnants G 359.1-0.5 and G 359.0-0.9 with ASCA
## 1. Introduction
Hard X-rays have opened a new window to see the Galactic center (GC) region. One of the remarkable discoveries is that a high-temperature ($``$ 10 keV) thin-thermal plasma is prevailing over an $``$ 100 pc-radius region around the GC (Koyama et al. 1986, 1989, 1996; Yamauchi et al. 1990). The presence of the large-scale hot plasma suggests violent activities in the GC region in the past, although the origin is not yet clear. One of the relevant processes to produce such a high-temperature diffuse plasma would be multiple supernova explosions. In this context, X-ray observations of individual supernova remnants (SNRs) near to the GC may provide useful information not only on the SNR physics, but also on the origin of the GC plasma.
Since X-rays from SNRs are usually dominated in the low-energy band ($`2`$ keV), and are easily absorbed by interstellar gas, an X-ray study of the SNRs near to the GC region is rather limited. High sensitivity and imaging capability in the hard X-ray band are required to detect SNRs behind a large absorbing medium, and to separate individual SNRs from the GC X-ray emissions, a complex diffuse plasma, binary sources, and the other stellar objects.
X-ray imaging spectroscopy also provides direct information on the nuclear synthesis, the total explosion energy and the age of SNRs, physical parameters of the surrounding interstellar environment, such as the density, its chemical compositions, and other related subjects: the star forming rate, the structure and the evolution in the central region of the Galaxy.
The ASCA satellite, having high sensitivity in the hard X-ray band ($`2`$ keV) and high energy resolution, enables us to study more elaborate imaging spectroscopy than was possible with previous detectors. We conducted survey and pointing observations near to the GC region with the ASCA satellite, and found X-rays from radio SNRs, including new candidates. Among them, this paper reports on the first detailed X-ray information and analyses of two radio SNRs, G 359.1$``$0.5 and G 359.0$``$0.9.
G 359.1$``$0.5 was first identified as an SNR by a 4.9 GHz observation (Downes et al. 1979) and by a 10.55 GHz observation (Sofue et al. 1984). Uchida et al. (1992b) found a shell-like structure at 1.4 GHz surrounded by a <sup>12</sup>CO ring. Comparing the 21 cm absorption feature of the <sup>12</sup>CO ring with the Galactic rotation curve, they concluded the location of this SNR to be near to the GC. Although and Egger, Sun (1998) discovered X-rays from G 359.1$``$0.5 with ROSAT, the spectral parameters, such as the temperature and the chemical composition, are not well constrained, due to the poor statistics and limited energy resolution. Preliminary results of the ASCA observation on G 359.1$``$0.5 are found in Yokogawa et al. (1999).
G 359.0$``$0.9 was first identified as an SNR by a 10.55 GHz observation (Sofue et al. 1984) and by a 2.7 GHz observation (Reich et al. 1990), and was later found to have an incomplete shell at the 843 MHz (Gray 1994). Leahy (1989) first detected a partial shell of soft X-rays from G 359.0$``$0.9 with the Einstein satellite, but no spectral information was reported.
This paper presents more comprehensive ASCA results and analyses of these two SNRs. Particular care concerning the background subtraction was made to exclude any possible contamination of near-by bright X-ray sources and the GC plasma’s contribution, of which the X-ray flux differs from position to position.
We describe the observations and the method of data reduction in section 2, and the analyses in subsection 3.1 and subsection 3.2, for G 359.1$``$0.5 and G 359.0$``$0.9, respectively. Section 4 is devoted to results and discussions on the distances, chemical compositions and morphology of these SNRs, and also on a relevant subject, the origin of the GC plasma.
## 2. Observations and Data Reduction
The GC region was observed with the X-ray astronomy satellite ASCA (Tanaka et al. 1994). X-ray photons were collected with four XRTs (X-ray Telescopes; Serlemitsos et al. 1995) and simultaneously detected with four detectors on the foci, which are two GISs (Gas Imaging Spectrometers; Ohashi et al. 1996) and two SIS (Solid-state Imaging Spectrometers; Burke et al. 1991) cameras.
During observation of the most prominent radio filament (the Snake) made on 1997 March 20–22, G 359.1$``$0.5 was located in the GIS field (see figure 1). However, it was only partially covered with the SIS field (see figure 3), in which 2-CCD chips for each SIS were operated with the faint or bright modes, depending on the high or medium bit rates. Since the 2-CCD data were significantly degraded by the Residual Dark Distribution (RDD) noise (T.Dotani et al. 1997, ASCA News 5, 14), we applied the RDD correction technique, which is only possible for faint-mode data.
G 359.0$``$0.9 was in the GIS field when the ASCA was pointing at the region of an unusual radio source the “Mouse” (Yusef-Zadeh, Bally 1987) and X-ray bursters SLX 1744$``$299/300 on September 12–14 in 1998 (see figure 4). In this observation, because SIS was operated in the 1-CCD mode, G 359.0$``$0.9 was totally out of the SIS field.
In both observations, the GISs were always operated in the normal PH mode. We excluded high-background data and non-X-ray events with the standard method according to the user guide by NASA Goddard Space Flight Center (GSFC). In total, the available exposures are $``$ 80 ks of the GIS, and $``$35 ks of the SIS for G 359.1$``$0.5, and $``$ 57 ks of GIS for G 359.0$``$0.9.
## 3. Analyses
### 3.1. G 359.1$``$0.5
We found X-rays from the position of the radio SNR G 359.1$``$0.5 only in soft X-ray bands below about 3–4 keV. Figure 1 shows the GIS contour map in the energy of 1.6–2.1 keV, in which band diffuse X-ray excess inside the radio shell is most clearly seen. In the hard X-ray band above 3 keV, on the other hand, we found no significant X-rays, neither along the radio shell nor in the center of the radio SNR. We also note that no significant X-rays were found from the Snake, the most prominent radio non-thermal filament (Uchida et al. 1992b).
Since we found soft X-rays only from the center region surrounded by the radio shell, we made the GIS spectrum of the SNR center, accumulating X-ray photons from GIS 2 and 3 in the X-ray bright region of a 6-radius circle, as shown by the solid line in figure 1.
As the background, we used the dotted line region in figure 1. This background region was carefully selected so that (1) possible contamination from the two nearby bright sources, A 1742$``$294 and 1E 1740.7$``$2942, would be precisely subtracted and (2) the contribution of the Galactic center plasma, distributed symmetrically around the GC, could be properly subtracted. For the two requirements, the angular distances of the background region from the GC and the other two bright sources were taken to be nearly the same as those of the source region. The background-subtracted spectrum of GIS 2+3 is given in figure 2. We have checked the GIS energy scale by using the 1.86 keV Si line in the Galactic diffuse spectrum, since a long-term variation of a few percent level has been reported by Miyata (1996) and Yamauchi et al. (1999). In this observation, we found that the GIS energy scale was very accurate with errors of less than 1%.
In the spectrum of G 359.1$``$0.5 (figure 2), we can notice two prominent emission lines at about 1.9 and 2.6 keV. To determine the accurate energy of the lines, we first fitted the spectrum to a thermal bremsstrahlung (for continuum) and two Gaussian lines with an interstellar absorption. The abundance for the interstellar gas was assumed to be solar, and the absorption cross sections were taken from Morrison and McCammon (1983) (hereafter, we refer to these absorption cross sections unless otherwise mentioned).
The best-fit center energies for the two lines were determined to be 1.86$`{}_{0.04}{}^{}{}_{}{}^{+0.03}`$ keV and 2.61$`{}_{0.07}{}^{}{}_{}{}^{+0.07}`$ keV (here and after, the errors are 90% confidence, unless otherwise mentioned). These line-center energies are consistent with those of K$`\alpha `$ emission from helium-like silicon (He-like Si; 1.86 keV) and hydrogen-like sulfur (H-like S; 2.63 keV); hence, the observed line structures are attributable to these highly ionized atoms. The best-fit line energies, on the other hand, are different from those of H-like Si (2.00 keV) and He-like S (2.45 keV).
Although the presence of K-shell lines of He-like Si and H-like S supports that the SNR X-rays are due to a thin thermal plasma, these two lines can hardly coexist in a single-temperature plasma, because ionization of sulfur atoms requires a higher temperature than that to ionize silicon atoms. In fact, we can reject any single-temperature plasma model, even though allowance is made for a non-ionization equilibrium or a non-solar abundance plasma. Therefore, we applied a two-temperature model (MEKAL: the plasma code established by Mewe et al. (1985) and Kaastra (1992)) with a common interstellar absorption. The abundances of Si and S in each plasma were treated to be free parameters, whereas those of the other elements were fixed to the solar values (Anders, Grevesse 1989).
This model is statistically accepted within the 90% confidence level. The best-fit model and parameters are given in figure 2 and table 1, respectively. A remarkable result is that the sulfur abundance in the higher temperature plasma (here component 2) is larger than $``$40 of solar.
Since SIS has a better energy resolution than GIS, a high-quality spectrum should be provided with SIS. However, as shown in figure 3, the SIS contour map in the 3.2–10.0 keV band, the relevant SIS field (2-CCD mode) covers only limited parts of the SNR. Furthermore, most of the SIS field is occupied by either G 359.1$``$0.5 or 1E 1740.7$``$2942. Therefore, a proper background region, which is free from contamination of the bright source 1E 1740.7$``$2942 and contains the same diffuse GC emission as that in the source region, is not available in the present SIS field. Consequently, we concentrate on analyses of the emission-line structure, which is less sensitive to the background flux. Fortunately, we found that the contamination source, 1E 1740.7$``$2942, exhibits no significant emission line (Sakano et al. 1999a). We then selected the source and the background regions as indicated by the solid and dotted lines in figure 3. As already noted, the surface brightness of the Galactic center diffuse X-rays decreases as the distance from the GC increases; hence, the simple background subtraction in the present case may cause an over-subtraction of the diffuse background. We therefore derived the iron-line fluxes from both the source and the background regions. The spectrum of the background region was subtracted from that in the source region, after normalizing the effective exposure by the iron-line flux ratio. These procedures are justified, because strong iron lines with a nearly constant equivalent width have been reported from the GC plasma (Koyama et al. 1989; Maeda 1998).
To the background-subtracted spectrum, we fitted a model of 4 Gaussians and a power-law spectra with an absorption, of which the latter is simply a phenomenological model to represent the continuous spectrum. The center energies of the Gaussian profiles were fixed to the theoretical values of He-like Si, H-like Si, He-like S, and H-like S. The absorption column density was fixed to the best-fit value of the GIS analyses (see table 1). The best-fit fluxes, the ratios and most probable plasma temperatures under the assumption of ionization equilibrium are listed in table 2 (Mewe et al. 1985).
From this table, we can see that the plasma temperatures determined with Si and S are different from each other, and are consistent with those found in the GIS spectrum.
To examine whether these two temperature plasmas have different spatial structures or not, we made the GIS and SIS images in the 1.6–2.1 keV and 2.1–3.2 keV bands; the former may represent the image of component 1, and the latter is for component 2. However, no significant difference between these two energy band images was found.
### 3.2. G 359.0$``$0.9
We made the X-ray images in several energy bands and found that X-rays from G 359.0$``$0.9 were detected only below $``$ 3 keV. The GIS contour map of G 359.0$``$0.9 in the soft X-ray band (0.7–1.5 keV) is shown in figure 4. Diffuse X-rays were clearly detected from the radio incomplete shell, although the X-ray shape is distorted by the bright X-ray sources SLX 1744$``$299/300. For a spectral study, we selected the source and background regions as given in figure 4. The selection criteria of the background region are the same as those of G 359.1$``$0.5; the angular distances of the background region from the GC and from SLX 1744$``$299/300 are the same as those of the source region.
We also examined any possible error of the GIS energy scale, as was done in the G 359.1$``$0.5 analyses, and found that the GIS energy is smaller than the proper value by about 2%, and hence made a fine tuning of the energy scale to the spectrum. The background-subtracted (and energy-scale tuned) spectrum is shown in figure 5. As expected from the multi-band X-ray images, most of the X-rays fall below $``$3 keV. Because the spectrum exhibits two clear lines, we applied a model of a thermal bremsstrahlung and two Gaussian profiles with an absorption, and determined the center energies to be 1.36$`{}_{0.03}{}^{}{}_{}{}^{+0.03}`$ keV and 1.86$`{}_{0.03}{}^{}{}_{}{}^{+0.03}`$ keV. Since these best-fit center energies are consistent with those of the K$`\alpha `$ lines from He-like Mg at 1.35 keV and He-like Si at 1.86 keV, the origin of the X-rays is a thin thermal plasma. We therefore applied a thin thermal plasma model, MEKAL, with an absorption, and fixing the abundances of all elements to be solar \[“model (a)”\]. The best-fit model spectrum and parameters are shown in figure 5 and table 3, respectively.
This model is, however, rejected with a reduced chi-square of 49.8/38, leaving bump-like and dip-like residuals around 1.3 keV and 1.0 keV, respectively (figure 5). Since these energies correspond to the emission lines of a He-like Mg and Fe-L line complex, we separately treated the abundances of Mg and Fe to be free and fitted the spectrum again \[“model (b)” and “(c)”\]. The best-fit spectra and parameters are given in figure 6, figure 7, and table 3 for models (b) and (c), respectively.
Although both the models (b) and (c) are acceptable, we further tried to fit the spectrum by treating the abundances of both Mg and Fe to be independent free parameters. However, the best-fit abundances have large errors, due to the limited statistics and the coupling of the He-like Mg line and Fe-L lines within the energy resolution of GIS. Therefore, we conclude that the plasma of G 359.0$``$0.9 is non-solar abundance, meaning that either Mg is over-abundant or Fe is under-abundant.
## 4. Results and Discussions
### 4.1. G 359.1$``$0.5
G 359.1$``$0.5 is found to exhibit a large absorption column of $`5.9\times 10^{22}`$ H cm<sup>-2</sup>. Since G 359.1$``$0.5 is reported to be surrounded by the <sup>12</sup>CO ring for a total mass of about $`2.5\times 10^6M_{}`$ (Uchida et al. 1992b), local absorption due to the <sup>12</sup>CO ring may not be ignored. Assuming that the <sup>12</sup>CO ring is a homogeneous shell with nearly the same shape of the G 359.1$``$0.5 radio shell, the absorption column due to the <sup>12</sup>CO ring is estimated to be $`3\times 10^{22}`$ H cm<sup>-2</sup>. Therefore, we infer that the column density of the foreground interstellar matter is about $`3\times 10^{22}`$ H cm<sup>-2</sup>. This value is equal to that of other X-ray sources near to the GC with the same Galactic coordinate of this SNR (Sakano et al 1999b); hence, this SNR would really be located near to the GC region with a distance of about 8.5 kpc. Thus, the diameter of the radio shell is estimated to be $``$ 57 pc, while that of the X-ray emitting central sphere is $``$ 28 pc.
We found that G 359.1$``$0.5 has at least two temperature plasmas; the cooler plasma (component 1) is abundant in Si, whereas the hotter one (component 2) is extremely over abundant in S. The center-filled thermal X-rays imply that these plasmas originated from the ejecta. Assuming a $``$28 pc-dimameter of spherical plasma of uniform density with a filing factor of 0.1, we estimate the total mass of Si and S to be about $`0.1`$ $`M_{}`$ and $`0.3`$ $`M_{}`$, respectively. However, no current theory of nucleosynthesis in supernova explosions predicts such a large mass of S compared to Si (see e.g. Thielemann et al. 1996). This problem can be solved by assuming a smaller filling factor of S than that of Si; a smaller filling factor of S, less than 0.1, reduces the total mass of S to be acceptable for the model of Thielemann et al. (1996).
The context of the very small filling factor and the extreme richness of S lead us to suspect that the S-rich plasma is a “shrapnel” ejected from the massive progenitor of G 359.1$``$0.5, in analogy of the Vela SNR (Aschenbach et al. 1995). However, the narrow band image including only the S-line (2.1–3.2 keV) shows no spatial structure like a “shrapnel”, mainly due to a lack of photon statistics.
Uchida et al. (1992a) argued that the <sup>12</sup>CO ring surrounding the SNR shell was created by stellar winds and/or multiple supernovae of O-type stars, and that several radio sources clustered at the center of the SNR are possibly O-type stars. The X-ray spectrum of G 359.1$``$0.5 contains no clear emission line of Fe (see figure 2), and is thus consistent with the proposed O-star origin. However, we could not quantitatively predict on the mass of the progenitor because of a lack of photon statistics.
We found no shell-like X-ray from G 359.1$``$0.5, although the radio morphology shows an almost complete shell. Shell-like X-rays may originate either by (1) a synchrotron mechanism in a shell, or (2) a thermal plasma made by the shock wave.
The lifetime of X-ray emission in case (1) depends on the magnetic field. Generally, the GC region is known to have a stronger magnetic field than the other regions of our Galaxy. In particular, Robinson et al. (1996) observed Zeeman effects of three OH masers near to the shell of G 359.1$``$0.5, then directly estimated the magnetic field to be 0.4–0.6 mG, which is more than one order of magnitude larger than that in usual shell-likes SNRs ($`10\mu `$G). High-energy electrons to emit synchrotron X-rays have a lifetime of $`10^3`$ yr in an $`10\mu `$G field (e.g. Reynolds 1996); hence, those in the shell of G 359.1$``$0.5 should be much shorter than $`10^3`$ yr. From the large diameter of the G 359.1$``$0.5 shell of 60 pc, this SNR would be middle age, or typically $`>`$$``$$`10^4`$ yr. Thus, no synchrotron X-ray from the shell would be expected from this SNR.
The evolved age of G 359.1$``$0.5 makes a shock-heated shell rather cool, with a typical temperature of less than a few 100 eV. Furthermore, if the shell interacts with the dense <sup>12</sup>CO ring, the X-ray emitting shell density becomes larger, and hence the cooling time is much shorter. Therefore, X-rays should be very soft, and should be entirely absorbed by the large interstellar gas column. Thus, the apparent lack of an X-ray shell of G 359.1$``$0.5 would be reasonable.
Rho and Petre (1998) proposed that SNRs with center-filled X-rays and a shell-like radio structure should be called “mixed morphology” (MM) supernova remnants. They show that the scale height of the MM SNRs’ distribution from the Galactic plane is smaller than that of shell-like SNRs in both radio and X-ray bands, and that many MM SNRs are located near to molecular clouds or H I clouds. G 359.1$``$0.5 is located near to the Galactic plane, and is surrounded by the <sup>12</sup>CO ring (Uchida et al. 1992b). Thus, G 359.1$``$0.5 shares the common features of MM SNRs.
### 4.2. G 359.0$``$0.9
The best-fit column density of $`1.8\times 10^{22}`$ H cm<sup>-12</sup> is smaller than that of the GC region (Sakano et al. 1999b). In fact, using the Galactic interstellar mass model by Olling and Merrifield (1998), and by Dame et al. (1987), we can estimate the distance to be 6 kpc. Hence, G 359.0$``$0.9 would not be in the GC region, but would be a foreground SNR.
The shell-like morphology and a thin-thermal spectrum in X-rays suggest that the shell is a shock-heated plasma. The plasma, however, shows a partial shell extending only 1/4 of a full circle. With the reasonable assumption that the supernova explosion was spherically symmetric, the apparent partial shell is attributable to an inhomogeneity of the interstellar medium; the direction to the X-ray emitting shell would have a higher density than the other directions, and would hence have a larger surface brightness than the others. For simplicity, we assume that the dense region covers $`\pi /2`$ str toward the partial shell; thus, 1/8 of the total explosion energy would be given in this direction. Using the 6-kpc distance, the diameter and X-ray luminosity of the partial shell are estimated to be 38 pc and $`1\times 10^{34}`$erg s<sup>-1</sup>, respectively. Together with the observational temperature of 0.4 keV, we could solve the Sedov equation, and found that the total explosion energy ($`E`$), density ($`n`$) of the interstellar medium (dense region) and the age ($`t`$) are $`1.2\times 10^{51}`$ erg, 0.5 cm<sup>-3</sup>, and 1.8$`\times 10^4`$ yr, respectively. Because the total mass of the X-ray emitting shell is about 40$`M_{}`$, most of them are attributable to the swept-up interstellar matter. Therefore, the result of the spectral fitting (in table 3) implies that the interstellar matter near to the G 359.0$``$0.9 is either over-abundant in Mg, or under-abundant in Fe.
### 4.3. Comments on the Galactic Diffuse Plasma
The GC region is surrounded by the large scale hot plasma, which emits fairly strong X-rays with many emission lines from He-like and H-like Si, S, Ar, Ca and Fe. The co-existence of highly ionized light elements (such as Si and S) and heavier elements (such as Ca and Fe) implies that the plasma is not a single temperature. Kaneda (1996) and Kaneda et al. (1997) confirmed the two-temperature structure of the Galactic ridge plasma, which shows a very similar spectrum to that of the GC region. In fact, Maeda(1998) found that the GC plasma has two-temperature components.
The low-temperature component ( $`<`$$``$ 1 keV) would be the same as that found with ROSAT (Snowden et al. 1997). The plasma has a large-scale height of 1.9 kpc and a temperature of 0.3$``$0.4 keV. Kaneda et al. (1997) suggested that the soft components of the Galactic ridge originated from a multiple supernova explosion. From the spectral similarity, this scenario may be applied to the GC soft component. Then, can we find many individual SNRs which have a similar spectral shape? Because G 359.0$``$0.9 has a 0.4 keV temperature and strong emission line of Mg and Si, it is a possible candidate. The GIS flux at 0.5–10.0 keV of G 359.0$``$0.9 is $`2.4\times 10^{12}`$ erg s<sup>-1</sup>cm<sup>-2</sup>, while that of the GC plasma in a 25 diameter field is $`8.3\times 10^{12}`$ erg s<sup>-1</sup>cm<sup>-2</sup> (Kaneda 1996). We thus require about four G 359.0$``$0.9-like SNRs in this field. However, at this moment, the number density of the resolved X-ray SNRs or its candidates is far less than the requirement.
A more difficult issue is the origin of the high-temperature component. The temperature of 10 keV and the size of about 100-pc diameter are inferred by the observation of Ginga and ASCA GC surveys (Koyama et al. 1989; 1996). As far as we know, no SNR exhibits such a high temperature. In this sense, the high-temperature sulfur-rich plasma of G 359.1$``$0.5 is suggestive. If there are many G 359.1$``$0.5-like SNRs, with an enriched abundance of not only S, but also other elements in high-temperature plasmas, the integrated emission may contribute, at least, some fractions of the GC hot plasma. At present, we are still lacking any quantitative information of individual X-ray sources, like SNRs, near to the GC region. Thus, further systematic study is highly encouraged.
We would like to thank Dr.Y. Maeda for useful comments and discussion about G 359.1$``$0.5. We also thank the members of the ASCA team. This work was supported by the Research Fellowship of the Japan Society for the Promotion of Science for Young Scientists.
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# High energy physics and the very early Universe with LISA
## I Introduction
The Universe became ”thin” to gravitational waves (GWs) at the Plank epoch, corresponding to the cosmic time $`10^{43}`$ sec.; the gravitons decoupled from the surrounding plasma at a temperature of the order of the Planck mass $`10^{19}`$ GeV, and gravitational radiation produced at that epoch or later – including the electro-weak and the Grand Unification (GUT) scale – has travelled undisturbed to us, carrying full information about the state of the Universe, and the physical processes from which it took origin. Indeed GW experiments will open radically new frontiers for cosmology and high energy physics (see and reference therein for an extensive discussion).
In the time frame $``$ 2002-2010 a large portion of the GW spectrum will progressively become accessible, mainly through large-scale laser interferometers. On the ground, the world-wide network of km-size interferometers – LIGO, GEO600, VIRGO and TAMA – sensitive in the frequency band $`10\mathrm{Hz}1\mathrm{kHz}`$, will start carrying out ”science runs” at the beginning of 2002, with the realistic goal of directly detecting GW’s. Several instrumental upgrades, starting around 2005, will drive the sensitivity of the instruments to a GW stochastic background from $`h_{100}^2\mathrm{\Omega }10^6`$ (for the so-called initial generation) to $`h_{100}^2\mathrm{\Omega }10^{10}`$ (for the so-called advanced configuration). In space, a collaboration between ESA and NASA is carrying out the project called LISA (Laser Interferometer Space Antenna). This is a space-borne laser interferometer with arms of length $`5\times 10^6`$ km, planned to fly by 2010 . This instrument guarantees the detection of GW’s at low frequencies ($`10^5\mathrm{Hz}10^2\mathrm{Hz}`$).
The purpose of this paper is to show the central role of the experiments in the low-frequency window $`10^6`$ Hz – 1 Hz, with emphasis on instruments of the LISA class, in the search for the primordial GW background. Our aim is to identify the fundamental issues regarding the achievement of a sensitivity in the range $`h_{100}^2\mathrm{\Omega }10^{16}`$$`10^{15}`$, which is set by the theoretical prediction of ”slow-roll” inflationary models.
### A The stochastic background spectrum
A stochastic GW background is a random process that can be described only in terms of its statistical properties. Without loss of generality, for the issues discussed in this paper, we assume it to be isotropic, stationary, Gaussian and unpolarized. The energy and spectral content of a stochastic background are described by the dimensionless function
$$\mathrm{\Omega }(f)\frac{1}{\rho _\mathrm{c}}\frac{d\rho _{\mathrm{gw}}(f)}{d\mathrm{ln}f};$$
(1)
$`\rho _{\mathrm{gw}}`$ is the energy density carried by the background radiation, and
$`\rho _\mathrm{c}={\displaystyle \frac{3H_0^2c^2}{8\pi G_N}}`$ $``$ $`1.6\times 10^8h_{100}^2\mathrm{erg}/\mathrm{cm}^3,`$ (2)
$``$ $`1.2\times 10^{36}h_{100}^2\mathrm{sec}^2`$ (3)
is the critical energy density required today to close the Universe. The value of the Hubble constant (today) is
$`H_0`$ $`=`$ $`100h_{100}\mathrm{km}\mathrm{sec}^1\mathrm{Mpc}^1,`$ (4)
$``$ $`3.2\times 10^{18}h_{100}\mathrm{sec}^1,`$ (5)
where $`h_{100}`$ is known from observations to be in the range $`0.4h_{100}0.85`$. $`\mathrm{\Omega }(f)`$ is therefore the ratio of the GW energy density to the critical energy density per unit logarithmic frequency interval; one usually refers to $`h_{100}^2\mathrm{\Omega }(f)`$, which is independent of the unknown value of the Hubble constant.
It is useful to introduce the characteristic amplitude $`h_\mathrm{c}(f)`$ of the GW background: it is the dimensionless characteristic value of the total GW background-induced fluctuation $`h(t)`$ at the output of an interferometer per unit logarithmic frequency interval:
$$h^2(t)=2_0^{\mathrm{}}d(\mathrm{ln}f)h_\mathrm{c}^2(f);$$
(6)
here $``$ denotes the expectation value. The spectral density $`S(f)`$ of the background is related to $`h_\mathrm{c}(f)`$ by
$$h_\mathrm{c}^2(f)=2fS(f),$$
(7)
and $`\mathrm{\Omega }(f)`$, $`h_\mathrm{c}(f)`$, and $`S(f)`$ satisfy the relation
$$\mathrm{\Omega }(f)=\frac{2\pi ^2}{3H_0^2}f^2h_\mathrm{c}^2(f)=\frac{4\pi ^2}{3H_0^2}f^3S(f).$$
(8)
The characteristic amplitude over a frequency band $`\mathrm{\Delta }f`$ is therefore:
$`h_\mathrm{c}(f,\mathrm{\Delta }f)`$ $`=`$ $`h_\mathrm{c}(f)\left({\displaystyle \frac{\mathrm{\Delta }f}{f}}\right)^{1/2}`$ (9)
$``$ $`7.1\times 10^{22}\left[{\displaystyle \frac{h_{100}^2\mathrm{\Omega }(f)}{10^8}}\right]^{1/2}\left({\displaystyle \frac{f}{1\mathrm{mHz}}}\right)^{3/2}\left({\displaystyle \frac{\mathrm{\Delta }_\mathrm{b}f}{3.2\times 10^8\mathrm{Hz}}}\right)^{1/2},`$ (10)
where $`\mathrm{\Delta }_\mathrm{b}f3.2\times 10^8(1\mathrm{yr}/T)`$ Hz is the width of the frequency bin for an observation time $`T`$. For comparison, the relevant characteristic amplitude of the LISA noise is $`10^{24}`$.
### B Sources of stochastic backgrounds
The stochastic GW background can be divided into two broad classes, based on its origin: (i) the primordial GW background (PGB), produced by physical processes in the early Universe, and (ii) the astrophysically-generated GW background (GGB), generated by the incoherent superposition of gravitational radiation produced, at much later cosmic times, by a large number of astrophysical sources that can not be resolved individually. The emphasis of this paper is on the detectability of the PGB.
In this paper we will use the following conventions: $`\mathrm{\Omega }_p(f)`$ and $`\mathrm{\Omega }_g(f)`$ identify the fractional energy density in GW’s, Eq. (1), carried by the primordial and the generated component of the GW background, respectively. If no index is used, we refer to a general GW stochastic signal, with no assumption about its production mechanism.
At present, there are three observational constraints on the PGB contribution to $`\mathrm{\Omega }(f)`$:
1. The high degree of isotropy of the cosmic microwave background radiation sets a limit at ultra-low frequencies :
$$h_{100}^2\mathrm{\Omega }_p(f)<7\times 10^{11}\left(\frac{f}{H_0}\right)^2,3\times 10^{18}h_{100}\mathrm{Hz}\stackrel{<}{}f\stackrel{<}{}10^{16}h_{100}\mathrm{Hz};$$
(11)
2. The very accurate timing of milli-second radio-pulsars constrains $`\mathrm{\Omega }_p(f)`$ in a frequency range of the order of the inverse of the observation time, typically of order of a few years :
$$h_{100}^2\mathrm{\Omega }_p(f)<10^8f10^8\mathrm{Hz};$$
(12)
3. The standard model of big-bang nucleosynthesis constrains the total energy content in GWs over a wide frequency range :
$$_{f=10^8\mathrm{Hz}}^{\mathrm{}}h_{100}^2\mathrm{\Omega }_p(f)d(\mathrm{ln}f)<6\times 10^6.$$
(13)
To foresee what physical processes could have produced a detectable GW background is an almost impossible challenge; nonetheless, it is enlighting to discuss some general principles and possible generation mechanisms to show the typical sensitivity that experiments should achieve in order to test different models.
The main mechanisms that produce a PGB can be divided into two broad categories (for a recent detailed review see ): (i) Parametric amplifications of metric tensor perturbations that occurred during an inflationary epoch, and (ii) Some causal processes – such as phase transitions – that took place in the early Universe.
Stochastic backgrounds produced by the parametric amplification of metric tensor perturbations that occurs during an inflationary epoch extend over a huge range of frequencies, from $`3\times 10^{18}`$ Hz up to a cutoff frequency in the GHz range. In the window $`10^{16}\text{Hz}1\text{GHz}`$, slow-roll inflationary models predict a quasi scale-invariant spectrum whose typical magnitude – in order to satisfy the COBE bound – cannot exceed $`h_{100}^2\mathrm{\Omega }_p10^{14}`$ in the LISA frequency band, as well as in the Earth-based detectors observational window ; a more refined analysis yields a more conservative upper limit: $`h_{100}^2\mathrm{\Omega }_p10^{16}10^{15}`$. Superstring-inspired cosmological models predict a spectrum that, for suitable choices of the free parameters of the model, could reach $`h_{100}^2\mathrm{\Omega }_p10^7`$ at the frequencies accessible either to Earth-based or to space-borne experiments, while satisfying the existing observational bound .
Stochastic backgrounds can also be produced by some classical causal processes that took place in the early Universe; for this class of signals, the characteristic frequency is related both to the time of emission and the corresponding temperature $`𝒯`$.
Non-equilibrium processes that occur at the reheating that takes place after inflation could provide a stochastic background with cutoff frequency in the range $`0.1\mathrm{mHz}1\mathrm{kHz}`$, corresponding to reheating temperatures between $`1\text{T}eV`$ and $`10^9\text{G}eV`$. As an example, in hybrid and extended inflationary models, the exit towards a radiation-dominated era is characterized by a first-order phase transition, which, if strongly of the first order, generates a stochastic background with $`h_{100}^2\mathrm{\Omega }_p10^6`$ at frequencies that can vary from the LISA observational window up to the sensitivity band of Earth-based interferometers .
Phase transitions that inevitably occur at $`𝒯10^2\text{M}eV`$ (the QCD phase transition) and $`𝒯10^2\text{G}eV`$ (the electroweak phase transition) produce GWs. In particular, if the electroweak phase transition is strongly of the first order, the spectrum is approximately $`h_{100}^2\mathrm{\Omega }_p10^{11}10^9`$ at $`f1\mathrm{mHz}`$ ; the requirement of a strong first order phase transition, which is necessary in order to have baryogenesis at the electroweak scale (see and references therein for a recent review), is directly related, in a minimal supersymmetric extension of the standard model, to the mass of the super-partner of the top quark .
Cosmic strings, which are topological defects formed during phase transitions, produce GW’s whose typical frequency ranges from $`f10^8`$ Hz up to $`f10^{10}`$ Hz with $`h_{100}^2\mathrm{\Omega }_p10^910^8`$, see and references therein for a review.
Global phase transitions associated with some scalar field which acquires a non-zero vacuum expectation value (VEV) below a critical temperature would produce, via a quite general relaxation process, GW’s whose energy content is very significant, $`h_{100}^2\mathrm{\Omega }_p10^6`$, for VEV’s near the Planck/string scale .
Recently there has been a great amount of theoretical activity investigating higher dimensional ”brane-world” scenarios, where gravity begins to probe the extra dimensions at energies as low as $`10^3`$ GeV; and an estimate of possible GW backgrounds in such models was presented recently in .
These examples clearly show that the investigation of the primordial GW stochastic background in the low-frequency regime would provide us key information about the physics beyond the standard model and/or could allow us to discriminate between different inflationary cosmological models.
### C Detecting a stochastic background
A stochastic background is a random process which is intrinsically indistinguishable from the detector noise. In order to detect such a signal, the optimal signal processing strategy calls for correlations between two (or more pairs of) instruments, possibly widely separated in order to minimize the effects of common noise sources. The relevant data analysis issues have been thoroughly addressed in ; here we simply review the main concepts, and refer to , and references therein, for more details.
The statistical analysis is based on the following assumptions: the signal and the detector noise are uncorrelated; the noise in each detector is stationary and Gaussian, and possible noise correlations between two detectors are negligible.
We define the output (signal + noise) of the two instruments as $`o_1(t)`$ and $`o_2(t)`$; the cross-correlation signal $`C`$ that one constructs is therefore of the form:
$$C_{T/2}^{T/2}𝑑t_{T/2}^{T/2}𝑑t^{}o_1(t)o_2(t^{})Q(tt^{}),$$
(14)
where $`Q(tt^{})`$ is a suitable filter function. In the general case, the filter function depends on $`t`$ and $`t^{}`$ independently, that is $`Q=Q(t,t^{})`$; here we have used the property of the signal of being stationary, and therefore $`Q(t,t^{})=Q(tt^{})`$. The SNR is defined as:
$$\mathrm{SNR}=\frac{\mu }{\sigma },$$
(15)
where $`\mu `$ and $`\sigma `$ are the mean value and the variance of the observable $`C`$:
$$\mu C=T\left(\frac{3H_0^2}{20\pi ^2}\right)(\stackrel{~}{Q},\stackrel{~}{A}),$$
(16)
$$\sigma ^2C^2C^2=\frac{T}{4}(\stackrel{~}{Q},\stackrel{~}{Q}).$$
(17)
Eqs. (16) and (17) are written in terms of the usual inner product
$$(a,b)_{\mathrm{}}^+\mathrm{}𝑑f\stackrel{~}{a}^{}(f)\stackrel{~}{b}(f)R(f),$$
(18)
where $`\stackrel{~}{Q}(f)`$ is the Fourier transform of $`Q(tt^{})`$. The functions $`R(f)`$ and $`\stackrel{~}{A}(f)`$ are defined as follows:
$`R(f)`$ $``$ $`S_n^{(1)}(f)S_n^{(2)}(f)\times `$ (20)
$`\left\{1+\left({\displaystyle \frac{3H_0^2}{10\pi ^2}}\right){\displaystyle \frac{\mathrm{\Omega }(f)}{f^3}}\left[{\displaystyle \frac{S_n^{(1)}(f)+S_n^{(2)}(f)}{S_n^{(1)}(f)S_n^{(2)}(f)}}\right]+\left({\displaystyle \frac{3H_0^2}{10\pi ^2}}\right)^2{\displaystyle \frac{\mathrm{\Omega }^2(f)\left[1+\gamma (f)^2\right]}{f^6S_n^{(1)}(f)S_n^{(2)}(f)}}\right\},`$
$$\stackrel{~}{A}(f)\frac{\gamma (f)\mathrm{\Omega }(f)}{f^3R(f)}.$$
(21)
In Eq. (20), $`S_n^{(k)}(f),k=1,2`$ is the one-sided noise power spectral density of the $`k`$th detector, and $`\gamma (f)`$ is the so-called overlap reduction function, which depends entirely on the relative orientation and location of the two detectors; it accounts for SNR losses that occur when the instruments are not optimally located and oriented, cf. Eq. (24) and Sec. II.
Using Eqs. (16) and (17), one can cast Eq. (15) in the form:
$$\mathrm{SNR}^2=T\left(\frac{3H_0^2}{10\pi ^2}\right)^2\frac{(\stackrel{~}{Q},\stackrel{~}{A})^2}{(\stackrel{~}{Q},\stackrel{~}{Q})}.$$
(22)
The optimal choice of the filter $`\stackrel{~}{Q}`$, is thus based on the maximizing the SNR, Eq. (22), and is given by:
$$\stackrel{~}{Q}(f)=(\mathrm{const}.)\times \stackrel{~}{A}(f),$$
(23)
where the overall normalization factor is arbitrary. Note that Eqs. (14)-(23) are valid for a background of arbitrary energy density $`\mathrm{\Omega }(f)`$. In the case of a signal much weaker than the noise, $`H_0^2\mathrm{\Omega }(f)/f^3S_n^{(k)}(f)`$, one can Taylor expand Eqs. (20) and (21), retaining only the leading order term. As a consequence, Eq. (22) reduces to:
$$\mathrm{SNR}\frac{3H_0^2}{\sqrt{50}\pi ^2}T^{1/2}\left[_0^{\mathrm{}}𝑑f\frac{\gamma (f)^2\mathrm{\Omega }^2(f)}{f^6S_n^{(1)}(f)S_n^{(2)}(f)}\right]^{1/2}(\mathrm{signal}\mathrm{noise}).$$
(24)
It is convenient to introduce the noise characteristic amplitude $`h_{\mathrm{rms}}`$, equivalent to $`h_\mathrm{c}`$, as follows:
$`n^2(t)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑fS_n(f)`$ (25)
$`=`$ $`2{\displaystyle _0^{\mathrm{}}}d(\mathrm{ln}f)h_{\mathrm{rms}}^2(f).`$ (26)
It is enlightening to write Eq. (24), using Eqs. (6) and (26), in the form:
$$\mathrm{SNR}\gamma (f_c)\left(\mathrm{\Delta }fT\right)^{1/2}\left[\frac{h_\mathrm{c}(f_c)}{h_{\mathrm{rms}}(f_c)}\right]^2;$$
(27)
here we have assumed that the frequency band $`\mathrm{\Delta }f`$, which contains most of the SNR, is centered on the characteristic frequency $`f_c`$, and is sufficiently small; the noise spectral density of the two instruments, that for simplicity we assume identical, and the overlap reduction function can be therefore treated as roughly constant. If only one instrument is in operation, one could in principle detect a stochastic background with SNR$`\stackrel{>}{}1`$ when $`h_\mathrm{c}\stackrel{>}{}h_{\mathrm{rms}}`$; with two instruments one can detect the signal when $`h_\mathrm{c}\stackrel{>}{}h_{\mathrm{rms}}/[\gamma (f_c)(\mathrm{\Delta }fT)^{1/4}]`$. Cross-correlation experiments are therefore highly desirable for both detection confidence and sensitivity. In fact, one can isolate the stochastic GW signal from all the spurious contributions which are uncorrelated between the two instruments. Common noise sources, which correlate on the same light-travel time scale, might, however, be present, degrading the overall sensitivity. Moreover, GW signals are expected to be very weak, well buried into the noise; using cross-correlations, through optimal filtering one increases the sensitivity by a factor $`10`$ $`(\mathrm{\Delta }f/1\mathrm{mHz})^{1/4}`$ $`(T/10^7\mathrm{sec})^{1/4}`$, with respect to the single detector case.
### D Summary of the results
The goal of this paper is to explore the capability of space-borne laser interferometers, such as LISA and its successors, in searching for the primordial GW stochastic backgrounds. The analysis of the LISA technology leads us to the following conclusions:
* A PGB of fractional energy density $`h_{100}^2\mathrm{\Omega }_p\stackrel{>}{}10^{10}`$ definitely shows up as an excess power component in the data of a single LISA interferometer over a large frequency window; it might be detectable by calibrating the noise-only response of the instrument, but the issue of decisively assigning this contribution to a real primordial signal remains open.
* Cross-correlations between the data streams of two identical LISAs, characterized by the presently estimated instrumental noise, allow us to reach a minimum value of the fractional energy density carried by GWs in the range $`5\times 10^{14}\stackrel{<}{}h_{100}^2\mathrm{\Omega }^{(\mathrm{min})}\stackrel{<}{}10^{12}`$ for an integration time $`T=10^7`$ sec, depending on the location and orientation of the two detectors.
* Such remarkable sensitivity, however, does not apply to primordial GW backgrounds; in fact, the copious number of short period solar-mass binary systems in the Universe produce a GGB that overwhelms the PGB in the key mHz region; we estimate that the minimum detectable value of the primordial GW background is $`h_{100}^2\mathrm{\Omega }_p^{(\mathrm{min})}\stackrel{>}{}5\times 10^{13}`$.
* The cross-correlation of the data streams from two LISA detectors provides, therefore, a powerful tool to extract information about populations of binary systems of short-period solar-mass compact objects in the Universe: the GGB is detectable at signal-to-noise ratio $`100`$.
We would like to emphasize that the third generation Earth-based laser interferometers, instruments such as advanced LIGO (LIGO III) and EUGO, a new European detector currently under study, will be able to achieve a sensitivity $`h_{100}^2\mathrm{\Omega }^{(\mathrm{min})}10^{10}`$; space-based interferometers will therefore play a primary role in GW cosmology.
Although the former results are very encouraging, our analysis leads to the rather obvious, but somewhat disappointing outcome that experiments in the frequency band $`10^6\mathrm{Hz}0.1\mathrm{Hz}`$ will be limited to a sensitivity of the order $`h_{100}^2\mathrm{\Omega }_p10^{13}10^{12}`$; this limit can not be improved by reducing the instrumental noise and/or increasing the integration time; in fact, GGB’s produce a residual correlation in the filter output designed to detect the PGB, that can not be eliminated: the mHz frequency window is not suitable to search for a primordial gravitational wave background characterized by $`h_{100}^2\mathrm{\Omega }_p<5\times 10^{13}`$.
It is therefore worth asking whether future technological developments, and more ambitious missions will enable the detection of a very weak PGB. Our present understanding suggests that cross-correlation experiments carried out with a pair of space-based interferometers with optimal sensitivity in the band $`0.1\mathrm{Hz}1\mathrm{Hz}`$ could be able to meet the target $`h_{100}^2\mathrm{\Omega }_p10^{16}`$. In this frequency band, in fact, astrophysically generated backgrounds are not present, and an experiment is limited in principle only by the instrumental noise. The design of a detector aimed at the window $`0.1\mathrm{Hz}1\mathrm{Hz}`$ imposes stringent requirements on the power and frequency of the laser, as well as on the dimensions of the ”optics” and on several other components of the instrument. In order to test slow-roll inflation an effective – ı.e. after residual radiation from individual binary systems has been removed – rms noise level $`10^{24}`$, and rather long integration times ($`3`$ yrs) are required. Such technological and data analysis issues have been little investigated so far, and deserve careful consideration.
The paper is organized as follows. In Sec. II we derive a closed form expression of the overlap reduction function for a pair of space-based interferometers. In Sec. III we review our present astrophysical understanding of the GW background generated by the incoherent superposition of radiation emitted by galactic and extra-galactic short period solar-mass binary systems. In Sec. IV we present the key results of the paper: we estimate the sensitivity that can be achieved by the present LISA technology, and possible follow-up missions, by cross-correlating the outputs of two identical detectors with uncorrelated noise. Sec. V contains our conclusions and pointers to future work.
## II The overlap reduction function for LISA detectors
In this section we compute a closed form expression for the overlap reduction function $`\gamma (f)`$, see Eqs. (20)-(23), for space-borne interferometers, and in particular for instruments characterized by a LISA-like orbital configuration. Indeed, our analysis provides explicit formulae that can be directly applied, with little changes, to any orbital configuration.
The overlap reduction function is a dimensionless function of the frequency $`f`$, which measures the degradation of the SNR when the detectors are not optimally oriented and located. At any given frequency, $`\gamma (f)`$ depends entirely on the relative separation and orientation of the instruments, and its behavior has a very intuitive physical explanation. It is maximum when the detectors are co-located and co-aligned. As the detectors (parallel to each others) are shifted apart, at a distance $`D`$, the signal drives oscillations in the two instruments that are progressively out of phase; there is an effective correlation only if the separation is smaller than approximately half of the characteristic wavelength of the gravitational radiation; equivalently, the frequency band over which one accumulates SNR is
$$f\stackrel{<}{}1\left(\frac{D}{1\mathrm{AU}}\right)^1\mathrm{mHz},$$
(28)
where one Astronomical Unit (AU) corresponds to $`1.4959787\times 10^{13}`$ cm. For two coincident detectors, $`\gamma (f)`$ decreases, at any frequency, as one instrument is rotated with respect to the other, because the detectors are excited in different ways by the different polarizations; for a rotation of $`\pi /4`$, the overlap reduction function is identically zero. In the general case, $`\gamma (f)`$ shows a complex behavior which is the superposition of the two effects that we have just described.
### A The LISA mission
LISA is an all-sky monitor with a quadrupolar antenna pattern. Its orbital configuration was conceived in order to keep the geometry of the interferometer as stable as possible during the mission, as well as to provide an optimal coverage of the sky: a constellation of three drag-free spacecraft (containing the free-falling test masses) is placed at the vertices of an ideal equilateral triangle with sides $`5\times 10^6\mathrm{km}`$; it forms a three-arm interferometer, with a $`60^{}`$ angle between two adjacent laser beams. The LISA orbital motion is such that the barycentre of the instrument is inserted in a heliocentric (essentially circular) orbit, following by $`20^{}`$ the Earth; the detector plane is tilted by $`60^{}`$ with respect to the Ecliptic and the instrument counter-rotates around the normal to the detector plane with the same period of $`1\mathrm{yr}`$.
We introduce a Cartesian reference frame $`(x,y,z)`$ tied to the Ecliptic, with $`\widehat{𝐳}`$ perpendicular to the Ecliptic plane, and $`\widehat{𝐱}`$ and $`\widehat{𝐲}`$ in the plane itself, and oriented in such a way to define a left-hand tern. In this frame, LISA’s centre-of-mass is described by the polar angles
$$\mathrm{\Theta }=\frac{\pi }{2},\mathrm{\Phi }(t)=\mathrm{\Phi }_0+n_{}t,(n_{}\frac{2\pi }{1\mathrm{yr}});$$
(29)
$`\mathrm{\Phi }_0`$ sets the position of the detector barycentre at some arbitrary reference time. The time evolution of the unit vectors $`\widehat{𝐥}_j`$ ($`j=1,2,3`$) along each arm are described by the following expression :
$`\widehat{𝐥}_j`$ $`=`$ $`\left[{\displaystyle \frac{1}{2}}\mathrm{sin}\alpha _j(t)\mathrm{cos}\mathrm{\Phi }(t)\mathrm{cos}\alpha _j(t)\mathrm{sin}\mathrm{\Phi }(t)\right]\widehat{𝐱}`$ (31)
$`+\left[{\displaystyle \frac{1}{2}}\mathrm{sin}\alpha _j(t)\mathrm{sin}\mathrm{\Phi }(t)+\mathrm{cos}\alpha _j(t)\mathrm{cos}\mathrm{\Phi }(t)\right]\widehat{𝐲}+\left[{\displaystyle \frac{\sqrt{3}}{2}}\mathrm{sin}\alpha _j(t)\right]\widehat{𝐳},`$
where $`\alpha _j(t)`$ increases linearly with time, according to
$$\alpha _j(t)=n_{}t(j1)\pi /3+\alpha _0,$$
(32)
and $`\alpha _0`$ is just a constant specifying the orientation of the arms at the arbitrary reference time $`t=0`$.
In the next section we will derive the expression of the overlap reduction function for two detectors characterized by a LISA-like motion. The time dependence of $`\widehat{𝐥}_j`$ and the centre-of-mass are described by Eqs. (29), (31) and (32) for both instruments, just with different initial conditions $`\alpha _0`$ and $`\mathrm{\Phi }_0`$. We will use the notation $`\alpha _{01}`$ and $`\mathrm{\Phi }_{01}`$, and $`\alpha _{02}`$ and $`\mathrm{\Phi }_{02}`$ to indicate the corresponding values of the detector ”1” and ”2”, respectively. For future convenience, we also introduce the notation:
$$\mathrm{\Delta }\mathrm{\Phi }_0=\mathrm{\Phi }_{02}\mathrm{\Phi }_{01},\mathrm{\Delta }\alpha _0=\alpha _{02}\alpha _{01}.$$
(33)
The main noise sources that affect the mission have been addressed by the LISA Science Team and yield the following expression for the expected noise spectral density
$$S_n(f)10^{48}\left[1.22\times 10^3f^4+16.74+9.7\times 10^4f^2\right]\mathrm{Hz}^1;$$
(34)
the rms noise amplitude that is inferred from $`S_n(f)`$ is shown in Fig. 1.
### B The overlap reduction function
To derive a closed form expression of the overlap reduction function $`\gamma (f)`$, we follow the formalism developed by Allen and Romano , which in turn was based on the analysis done by Flanagan . $`\gamma (f)`$ is formally given by
$$\gamma (f)=\rho _1(x)d_{ab}^{(1)}d^{(2)ab}+\rho _2(x)d_{ab}^{(1)}S^bd^{(2)ac}S_c+\rho _3(x)d_{ab}^{(1)}d_{cd}^{(2)}S^aS^bS^cS^d.$$
(35)
In the previous expression, $`d_{ab}^{(k)}`$ are the detector response tensors, defined as:
$$d_{ab}^{(k)}=\frac{1}{2}\left[m_a^{(k)}m_b^{(k)}n_a^{(k)}n_b^{(k)}\right],$$
(36)
where $`k=1,2`$ labels the instrument, and $`\widehat{𝐦}^{(k)}`$ and $`\widehat{𝐧}^{(k)}`$ are the unit vectors along the arms of the interferometers; our (arbitrary) choice corresponds to $`\widehat{𝐥}_1`$ and $`\widehat{𝐥}_2`$, respectively, given by Eqs. (31) and (32), with the appropriate initial conditions. The unit vector along the direction that connects the centers of mass of the two detectors is
$$\widehat{𝐒}=\frac{\mathrm{cos}\mathrm{\Phi }_2\mathrm{cos}\mathrm{\Phi }_1}{\sqrt{2(1\mathrm{cos}(\mathrm{\Phi }_2\mathrm{\Phi }_1))}}\widehat{𝐱}+\frac{\mathrm{sin}\mathrm{\Phi }_2\mathrm{sin}\mathrm{\Phi }_1}{\sqrt{2(1\mathrm{cos}(\mathrm{\Phi }_2\mathrm{\Phi }_1))}}\widehat{𝐲};$$
(37)
$`x`$ is a dimensionless parameter defined as
$$x\frac{2\pi fD}{c},$$
(38)
where
$$D=R_{}\sqrt{2\left(1\mathrm{cos}\mathrm{\Delta }\mathrm{\Phi }_0\right)},$$
(39)
and the functions $`\rho _j(x)`$ are given by
$`\rho _1(x)`$ $`=`$ $`5j_0(x){\displaystyle \frac{10}{x}}j_1(x)+{\displaystyle \frac{5}{x^2}}j_2(x),`$ (40)
$`\rho _2(x)`$ $`=`$ $`10j_0(x)+{\displaystyle \frac{40}{x}}j_1(x){\displaystyle \frac{50}{x^2}}j_2(x),`$ (41)
$`\rho _3(x)`$ $`=`$ $`{\displaystyle \frac{5}{2}}j_0(x){\displaystyle \frac{25}{x}}j_1(x)+{\displaystyle \frac{175}{2x^2}}j_2(x),`$ (42)
where
$`j_0(x)`$ $`=`$ $`{\displaystyle \frac{\mathrm{sin}x}{x}},`$ (43)
$`j_1(x)`$ $`=`$ $`{\displaystyle \frac{\mathrm{sin}x}{x^2}}{\displaystyle \frac{\mathrm{cos}x}{x}},`$ (44)
$`j_2(x)`$ $`=`$ $`3{\displaystyle \frac{\mathrm{sin}x}{x^3}}3{\displaystyle \frac{\mathrm{cos}x}{x^2}}{\displaystyle \frac{\mathrm{sin}x}{x}},`$ (45)
are the standard spherical Bessel functions.
By substituting Eq. (31) into Eq. (36), and combining it with Eqs. (37)-(45), one can derive the following expressions:
$`d_{ab}^{(1)}d^{(2)ab}`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{4}{}}}[A_k\mathrm{cos}(2\mathrm{\Delta }\alpha _0+k\mathrm{\Delta }\mathrm{\Phi }_0)+B_k\mathrm{sin}(2\mathrm{\Delta }\alpha _0+k\mathrm{\Delta }\mathrm{\Phi }_0),`$ (46)
$`d_{ab}^{(1)}S^bd^{(2)ac}S_c`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{4}{}}}[C_k\mathrm{cos}(2\mathrm{\Delta }\alpha _0+k\mathrm{\Delta }\mathrm{\Phi }_0)+D_k\mathrm{sin}(2\mathrm{\Delta }\alpha _0+k\mathrm{\Delta }\mathrm{\Phi }_0)],`$ (47)
$`d_{ab}^{(1)}d_{cd}^{(2)}S^aS^bS^cS^d`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{4}{}}}[E_k\mathrm{cos}(2\mathrm{\Delta }\alpha _0+k\mathrm{\Delta }\mathrm{\Phi }_0)+F_k\mathrm{sin}(2\mathrm{\Delta }\alpha _0+k\mathrm{\Delta }\mathrm{\Phi }_0)],`$ (48)
where $`\mathrm{\Delta }\alpha _0`$ and $`\mathrm{\Delta }\mathrm{\Phi }_0`$ are given by Eq. (33), and $`A_k`$, $`B_k`$, $`C_k`$, $`D_k`$, $`E_k`$ and $`F_k`$ (for $`k=0,..,4`$) are numerical coefficients, that are given in Appendix A. Inserting Eqs. (47)-(48) in Eq. (35), the overlap reduction function becomes:
$$\gamma (f)=\underset{k=0}{\overset{4}{}}\left[P_k(x)\mathrm{cos}\phi _k(\mathrm{\Delta }\alpha _0,\mathrm{\Delta }\mathrm{\Phi })+Q_k(x)\mathrm{sin}\phi _k(\mathrm{\Delta }\alpha _0,\mathrm{\Delta }\mathrm{\Phi })\right]$$
(49)
where
$`P_k(x)`$ $`=`$ $`\rho _1(x)A_k+\rho _2(x)C_k+\rho _3(x)E_k,`$ (50)
$`Q_k(x)`$ $`=`$ $`\rho _1(x)B_k+\rho _2(x)D_k+\rho _3(x)F_k,`$ (51)
depend only on the detector separation and the radiation frequency, and
$$\phi _k(\mathrm{\Delta }\alpha _0,\mathrm{\Delta }\mathrm{\Phi }_0)2\mathrm{\Delta }\alpha _0+k\mathrm{\Delta }\mathrm{\Phi }_0,$$
(52)
is a function of the relative orientation of the instruments.
We would like to stress that our definition is such that $`\gamma (f)=1`$ $`f`$, for co-aligned and co-located interferometers with arms perpendicular to each others, that is the angle between $`\widehat{𝐦}`$ and $`\widehat{𝐧}`$ is $`\pi /2`$. However, the LISA opening angle is $`\pi /3`$, and the detector response is reduced by the factor $`\sqrt{3}/2`$: as a consequence, the maximum value that $`\gamma (f)`$ can attain is $`3/4`$.
Notice also that, as pointed out by Cutler , the read-outs from the three arms of LISA can be combined in such a way to form the outputs, say $`o_I`$ and $`o_{II}`$, of two co-located interferometers rotated by $`\pi /4`$ one with respect to the other, whose noise is uncorrelated at all frequencies. Unfortunately, the cross-correlation of $`o_I`$ with $`o_{II}`$ is useless for searching for stochastic backgrounds, as the overlap reduction function is identically zero over the whole frequency range.
Figs. 2 and 3 show the behavior of $`\gamma (f)`$ as a function of the frequency and the separation angle $`\mathrm{\Delta }\mathrm{\Phi }_0`$, which is equivalent to the distance $`D`$, see Eq. (39). Here we assume that the instruments are inserted into their orbits so that $`\mathrm{\Delta }\alpha _0=0`$. In this case, at fixed frequency, the overlap reduction function depends only on $`\mathrm{\Delta }\mathrm{\Phi }_0`$, as the relative orientation of the detectors is determined only by $`\mathrm{\Phi }(t)`$, cf. Eqs. (31) and (32). This also implies that as the two instruments are moved apart – $`\mathrm{\Delta }\mathrm{\Phi }_0`$ increases – their orientation changes too. In order to achieve the maximum $`\gamma (f)`$ for a given separation, one should therefore suitably tune the initial orientation of the detector arms. The plots clearly show that at low frequencies, $`f\stackrel{<}{}5\times 10^4`$ Hz, $`\gamma (f)`$ is fairly flat and close to its maximum value (which is set by the separation and orientation); in fact the radiation wavelength is $`\lambda _{\mathrm{gw}}2(f/1\mathrm{mHz})^1`$ AU, and for separations smaller than 1 AU, the degradation of SNR at very low-frequencies is less than a factor $`2`$. At high frequencies, say $`f\stackrel{>}{}10^2`$ Hz, placing two interferometers at a distance larger than $`10^6`$ km, would severely degrade the SNR, by a factor of 10 or more.
## III The astrophysically generated stochastic background
The so-called astrophysically generated GW stochastic background (GGB) is mainly due to the incoherent superposition of gravitational radiation emitted by short-period solar-mass binary systems. A variety of binary populations contribute to it, but the main contribution, in the mHz region, is due to close white-dwarf binaries (CWDB’s). Present estimates suggest that it is above the LISA instrumental noise in the frequency region $`10^43\times 10^3`$ Hz, right at the heart of the observation window, see Fig. 1. However, sizeable effects are also given by other sources such as W UMa (Ursae Majoris) binary stars, unevolved binaries, cataclysmic binaries, neutron star-neutron star (NS-NS) binary systems, black hole-neutron star (BH-NS) binaries, and possibly BH-BH binary systems.
The GGB is a guaranteed GW source in the low-frequency band; however, it is likely to overwhelm the PGB, degrading the sensitivity of the instruments in searching for a stochastic signal produced in the early Universe. A rigorous analysis of GGB’s goes far beyond the purpose of this paper; here we review the main features and discuss the fundamental theoretical issues. We refer the reader to , and references therein, for a thorough discussion of the astrophysical sources.
Ultimately, the reason that the radiation generated by large populations of binary systems is effectively a stochastic signal is simple: there are too many free parameters that one needs to fit the data in order to resolve all the binary systems that contribute to the signal. In fact, our galaxy contains $`10^7`$ CWDB’s; they evolve, due to radiation reaction, over a time scale $`\stackrel{>}{}10^7`$ years. Therefore, during the typical observation time $`T1`$ yr, they are seen as highly monochromatic, and, in the band $`10^410^3`$ Hz, each frequency bin is ”contaminated” by roughly $`10^3`$ sources. The problem of resolving each individual binary system is actually made worse by the motion of the detector, because, in addition to the frequency, one needs to solve also for the source position in the sky and the orientation of the orbital plane.
In the following we discuss the estimates $`\mathrm{\Omega }_g(f)`$ of the spectral energy density of generated backgrounds, and in particular their isotropic component, and the critical frequency $`f_g`$ up to which GGB’s are indeed present. For $`f\stackrel{>}{}f_g`$, the individual binary sources of astrophysical populations can be resolved, and their radiation subtracted from the data, opening up the band to search for a PGB. The former two properties are essential in addressing the possibility of detecting PGBs with space-based experiments.
In this context, it is useful to divide the sources that contribute to the GGB into two categories: galactic sources and extra-galactic sources. Their key distinguishing feature of relevance here is the degree of isotropy of the GGB that they generate, for an observer on board a LISA-like detector. In fact, the extra-galactic contribution is expected to be isotropic to a rather high degree (the radiation being dominated by binary systems at cosmological distances; for more details see ); it is, therefore, impossible to discriminate it from the PGB. On the contrary, the GGB produced by galactic sources is clearly highly anisotropic. In fact, galactic stars are spatially distributed, approximately, according to $`\mathrm{exp}(r/r_0)\mathrm{exp}[(z/z_0)^2]`$, where $`r`$ is the radial distance to the Galactic centre, $`r_05`$ kpc, and $`z`$ the height above the Galactic plane; $`z_05`$ kpc for neutron star binaries and less than 300 pc for the other types. Due to the peripheral location of the Solar System, and the change of orientation of the LISA arms during the years-long observation time, galactic generated backgrounds appear strongly anisotropic . It is also conceivable that the isotropic portion of the galactic contribution does not exceed the total extra-galactic GGB.
Several astrophysical uncertainties affect the estimates of the GGB that have been carried out so far. A careful analysis of the galactic contribution has been performed by Bender and collaborators , taking into account a wide range of binary populations. A good fit of the galactic GGB spectral density is:
$$S_g(f)=\{\begin{array}{cc}10^{42.685}f^{1.9}\hfill & 10^5f10^{3.15},\hfill \\ 10^{60.325}f^{7.5}\hfill & 10^{3.15}f10^{2.75},\hfill \\ 10^{46.85}f^{2.6}\hfill & 10^{2.75}f,\hfill \end{array}$$
(53)
which is related to $`\mathrm{\Omega }_g(f)`$ by Eq. (8). In Eq. (53) the space density of CWDBs is assumed to be $`10\%`$ of the theoretical value predicted by Webbink , and the radiation from helium cataclysmic variables, likely to contribute significantly in the frequency window $`1\mathrm{mHz}3`$ mHz, is not taken into account. The estimate (53) can be definitely regarded as a solid lower limit, and is likely correct within to a factor $`3`$; see also . We would also like to stress that the dominant contribution from CWDBs switches off at $`10^2`$ Hz, where the binary systems coalesce; for $`f\stackrel{>}{}10^2`$ Hz, NS-NS binary systems are the sources that contribute most to the GGB.
The extragalactic contribution to the GGB has been estimated in , and is weaker than the galactic contribution by a factor $``$ 10-to-3 in the relevant frequency band, where the main uncertainty comes from the star formation rate at high redshifts (see however the optimistic estimates in ). As we have mentioned at the beginning of the section, it is likely that the total (galactic + extra-galactic) isotropic component of the GGB does not exceed the total extra-galactic GGB; indeed, we will assume that the isotropic portion of $`\mathrm{\Omega }_g(f)`$, the only one that affects searches for the PGB, follows the frequency distribution predicted by Eq. (53), but reduced by a factor $`ϵ1`$. In the rest of the paper we will therefore assume
$$\mathrm{\Omega }_{g,is}=ϵ\frac{4\pi ^2}{3H_0^2}f^3S(f).$$
(54)
Note that the determination of the value of $`ϵ`$ is a delicate matter. At frequencies above a few mHz, where the GWs from galactic sources can be subtracted, it is likely that $`ϵ\stackrel{<}{}0.3`$, while below $`1`$ mHz, Ref. suggests $`ϵ0.7`$. However a more detailed analysis is needed in the future to provide a better estimate of this value. In the following we optimistically set $`ϵ=0.1`$, which can be regarded as a solid lower limit; the results that we will present in the next Section can be easily rescaled as a function of $`ϵ`$.
We consider now the critical frequency $`f_g`$ up to which radiation emitted by solar-mass binary systems (in the whole Universe) produce a GGB. For $`f\stackrel{>}{}f_g`$ the observational window becomes ”transparent” to the primordial GW background. Following the discussion at the beginning of the section, we estimate $`f_g`$ by using a very simple argument based on first principles: if the number of independent degrees of freedom of the data set – the number of data points, or, equivalently, the number of frequency bins – is smaller than the total number of independent parameters that describe the radiation, then the superposition of monochromatic GWs must be considered as a stochastic background. If the opposite is true, we have enough information to characterize, at least in principle, each individual source, and the signal is a deterministic one. For the sake of simplicity – although not exactly true, see the discussion at the beginning of the section and – we assume that each binary system is characterized by one parameter. The critical frequency $`f_g`$ is, therefore, formally determined by the condition that the (average) number of sources per frequency bin is less than one:
$$\frac{dN(f)}{df}\mathrm{\Delta }_\mathrm{b}f\stackrel{<}{}1;$$
(55)
here $`dN/df`$ is the number of binary sources, emitting at frequency $`f`$, per unit frequency interval. Assuming that the merger rate is $`R`$, and in the relevant frequency range the binaries evolve only through radiation reaction, we have
$$\frac{dN(f)}{df}\mathrm{\Delta }_\mathrm{b}f=\frac{R}{T}\left(\frac{df}{dt}\right)^1,$$
(56)
where $`df/dt`$ can be estimated using the Newtonian quadrupole formula:
$$\frac{df}{dt}=\frac{96}{5}\pi ^{8/3}^{5/3}f^{11/3},$$
(57)
and $`(m_1m_2)^{3/5}/(m_1+m_2)^{1/5}`$ is the so-called chirp mass ($`m_1`$ and $`m_2`$ are the masses of the two orbiting stars). Substituting Eqs. (56), and (57) into Eq. (55), the frequency $`f_g`$ is easily determined:
$$f_g\left(\frac{5}{96}\right)^{3/11}\pi ^{8/11}^{5/11}\left(\frac{R}{T}\right)^{3/11}.$$
(58)
We are now interested in determining an upper limit to $`f_g`$, considering binary populations from the whole Universe. For $`f\stackrel{>}{}10`$ mHz, the main contribution to the GGB is given by NS-NS binaries . Their merger rate is uncertain; current estimates yield a galactic merger rate in the range :
$$R_{\mathrm{NS}}10^65\times 10^4\mathrm{yr}^1.$$
(59)
We can extrapolate this result to the entire Universe by simply multiplying the galactic rate by the total number of galaxies $`N_G`$:
$$R10^6\left(\frac{R_{\mathrm{NS}}}{10^5\mathrm{yr}^1}\right)\left(\frac{N_G}{10^{11}}\right)\mathrm{yr}^1.$$
(60)
By using this approach, we assume that $`R_{\mathrm{NS}}`$ does not vary with the redshift, which is probably not true. However, even if at high redshift $`R_{\mathrm{NS}}`$ is a factor 10 higher than in our galaxy, the very weak dependence of $`f_gR^{3/11}`$ on the merger rate, see Eq. (58), ensures that this crude estimate is correct within a factor $`2`$. Assuming that the typical chirp mass of the binaries in the population is $`\overline{}=1.2M_{}`$, which corresponds to $`m_1=m_2=1.4M_{}`$, and the rate (60), Eq. (58) yields:
$$f_g1.6\times 10^1\left(\frac{R}{10^6\mathrm{yr}^1}\right)^{3/11}\left(\frac{T}{1\mathrm{yr}}\right)^{3/11}\left(\frac{\overline{}}{1.2M_{}}\right)^{5/11}\mathrm{Hz}.$$
(61)
This simple analysis leads therefore to the conclusion that the window $`f\stackrel{>}{}0.1`$ Hz is likely free from stochastic backgrounds generated by astrophysical sources; radiation from NS-NS binaries is still present, but one can detect each individual source, estimate its parameters, and remove from the data stream the signals. In principle, the search for the PGB becomes limited only by the instrumental noise.
## IV Sensitivity
We can now proceed to discuss the sensitivity of LISA, and possible follow-up missions, to search for the PGB. The major disadvantage of the presently designed LISA mission is the lack of two instruments with uncorrelated noise; in fact, each pair of the three co-located LISA interferometers share one arm, and therefore common noise. One can, in principle, extract two data streams with uncorrelated noise at all frequencies . unfortunately they are equivalent to the outputs of a pair of detectors, one rotated by 45 with respect to the other, and as we have stressed in Sec. II, the response to a stochastic signal is then null. Nonetheless, LISA can play an important role in searching for stochastic backgrounds, as it can possibly achieve a sensitivity $`h_{100}^2\mathrm{\Omega }10^{11}10^{10}`$ by exploiting some intrinsic properties of the signal and/or the unique features of the instrument to operate in a configuration where the GW signal is (almost) absent.
During the observation time, LISA changes orientation with period 1 yr. A background which is anisotropic produces a periodic modulation in the auto-correlation function that can stand above the noise . Unfortunately, one can not exploit this feature to detect a primordial signal that we expect to be intrinsically isotropic to a high degree; the only sizeable anisotropy that one can foresee is a dipolar one, produced by the motion of our local system (and therefore LISA) with respect to the cosmological rest frame with velocity $`v_{\mathrm{prop}}/c10^3`$. Unfortunately, an interferometer has a quadrupolar antenna pattern, and the SNR produced by the quadrupole component is reduced by a factor $`(v_{\mathrm{prop}}/c)^210^6`$. Nonetheless, the use of the signature induced by anisotropies could lead to the detection of galactic generated backgrounds, as suggested in .
Time delay interferometry with multiple readouts provides a way of suppressing by several orders of magnitude the GW contribution from the LISA output at frequencies below a few mHz, providing a shield to GW radiation. By exploiting this feature, one can calibrate the noise-only response of the instrument, and search for an excess power in the data stream that can be assigned to a signal of cosmic origin.
LISA can therefore carry out searches for stochastic backgrounds with a single instrument at a very interesting sensitivity level. LISA is likely to detect the galactic astrophysically-generated background, and could set important upper-limits on the primordial contribution, that can rule out existing models. It also provides invaluable information regarding the astrophysically generated backgrounds that might turn out to be crucial in designing future follow-up missions devoted to GW cosmology.
However, in order to reach the sensitivity $`h_{100}^2\mathrm{\Omega }10^{16}`$ that we have set as a goal, two separated interferometers are essential. It is important to understand whether this target sensitivity is within the reach of near future space technology, and to discuss possible fundamental limitations that might prevent us from achieving a detection at the level $`h_{100}^2\mathrm{\Omega }10^{16}`$. The science goals are of such importance that some suggestions have been already put forward as to how to introduce modifications to the currently envisaged LISA configuration in order to accommodate a pair of independent instruments, possibly with the capability of shifting at-will the centre of the sensitivity window from $`10^3`$ Hz to $`0.1`$ Hz .
We start by considering a pair of identical LISA detectors; we therefore assume present, or near future, technology. Then, we discuss ”second generation” LISA detectors, specifically aimed at PGB searches.
### A The sensitivity of two LISA interferometers
We consider what the present LISA technology allow us to achieve by cross-correlating the outputs of two identical interferometers located at a distance $`D`$, cf. Eq. (39); this is equivalent to determining the minimum value of $`h_{100}^2\mathrm{\Omega }`$ that one is able to detect. No unique answer can be given to this question, as it depends, of course, on the frequency dependence of the true signal $`\mathrm{\Omega }(f)`$. Lacking any solid theoretical prediction, we choose a ”maximum ignorance” approach, and assume that $`\mathrm{\Omega }(f)`$ is constant over the relevant frequency range. This hypothesis is not unreasonable, because the frequency band over which the SNR builds up is fairly small, due to $`S(f)`$ and $`\gamma (f)`$, cf. Eqs. (34), (49), and Figs. 1 and 2. By setting $`\mathrm{\Omega }(f)=\mathrm{const}`$, and $`S_n^{(1)}(f)=S_n^{(2)}(f)=S_n(f)`$, and solving Eq. (24) for $`\mathrm{\Omega }`$, we can estimate the minimum detectable value of the energy density content of the GW background<sup>*</sup><sup>*</sup>*Notice that it is appropriate to use the weak (with respect to the noise) signal approximation, Eq. (24), as we are aiming at the detection of the weakest possible background, which is clearly dominated by the noise.:
$$h_{100}^2\mathrm{\Omega }^{(\mathrm{min})}=\frac{K}{T^{1/2}}\frac{\sqrt{50}\pi ^2}{3H_0^2}\left[_0^{\mathrm{}}\frac{\gamma ^2(f)}{f^6S_n^2(f)}\right]^{1/2};$$
(62)
the constant $`K`$ is related to the false alarm probability and the detection rate associated with the measurement of a background with energy $`\mathrm{\Omega }^{(\mathrm{min})}`$. For a false alarm probability of $`5\%`$, and a detection rate of $`95\%`$, we have $`K3.76`$ (the sum of the false alarm probability and the detection rate need not to be 1, and these two quantities are totally subjective) . We therefore compute Eq. (62), with noise spectral density and overlap reduction function given by Eqs. (34) and (49), respectively. The results presented here are computed, for sake of simplicity, for the case $`\mathrm{\Delta }\alpha _0=0`$, and the integration is carried out over the frequency band $`10^410^2`$ Hz.
It is interesting to analyze first how different frequency regions of the whole sensitivity window contribute to the total SNR. Fig. 4 shows the fraction of the total signal-to-noise ratio that is accumulated per unit logarithmic frequency interval; it indicates clearly that the key frequency band is $`8\times 10^4\mathrm{Hz}5\times 10^3`$ Hz.
Fig. 5 shows the sensitivity to a generic GW stochastic background that one could in principle achieve for a time of observation $`T=10^7`$ sec. It is straightforward to rescale these results for a different observation time, as SNR $`T^{1/2}`$ and $`h_{100}^2\mathrm{\Omega }^{(\mathrm{min})}T^{1/2}`$, see Eqs. (15) and (62). It is remarkable that LISA is capable of measurements in the range $`5\times 10^{14}\stackrel{<}{}h_{100}^2\mathrm{\Omega }^{(\mathrm{min})}\stackrel{<}{}10^{12}`$, which is about three orders of magnitude better than can be achieved (although in a different frequency regime) with Earth-based interferometers operating in the ”advanced” configuration. The latter experimental set-up requires considerable technological and possibly conceptual developments, and is expected to be implemented not before the end of the decade. The one-order-of-magnitude range in $`h_{100}^2\mathrm{\Omega }^{(\mathrm{min})}`$ is due to the effect of the overlap reduction function, where the lower limit is for co-located instruments, and the upper limit for detectors at a distance D = 2 AU with $`\mathrm{\Delta }\alpha _0=0`$, cf. also Fig. 4.
However, such instrumental sensitivity does not correspond to a comparable sensitivity to PGBs; in fact, the peak of the GGB spectrum, of the order $`10^{11}`$ for the isotropic component, is where the instrument is most sensitive, cf. Fig. 1. In order to quantify this effect, and to address the capabilities of space-based interferometers for cosmology, we need a short digression.
Consider a generic two-component stochastic signal:
$$\mathrm{\Omega }(f)=\mathrm{\Omega }_1(f)+\mathrm{\Omega }_2(f),$$
(63)
and assume that $`\mathrm{\Omega }_1(f)`$ and $`\mathrm{\Omega }_2(f)`$ share exactly the same statistical properties (they are isotropic, stationary, Gaussian and unpolarized), and therefore can not be distinguished from each other. In order to search for $`\mathrm{\Omega }_1(f)`$, one constructs the optimal filter
$$\stackrel{~}{Q}_1=\frac{\gamma (f)\mathrm{\Omega }_1(f)}{f^3R(f)},$$
(64)
and correlates it against the data of two instruments, following the scheme described in the Introduction, Sec. I C. The signal-to-noise ratio at the output of the filtering process is
$`\mathrm{SNR}^2`$ $`=`$ $`T\left({\displaystyle \frac{3H_0^2}{10\pi ^2}}\right)^2{\displaystyle \frac{(\stackrel{~}{Q}_1,\left[\frac{\gamma (f)\mathrm{\Omega }(f)}{f^3R(f)}\right])^2}{(\stackrel{~}{Q}_1,\stackrel{~}{Q}_1)}}`$ (65)
$`=`$ $`T\left({\displaystyle \frac{3H_0^2}{10\pi ^2}}\right)^2(\stackrel{~}{Q}_1,\stackrel{~}{Q}_1)\left\{1+{\displaystyle \frac{(\stackrel{~}{Q}_1,\stackrel{~}{Q}_2)}{(\stackrel{~}{Q}_1,\stackrel{~}{Q}_1)}}\right\}^2.`$ (66)
The second term in brackets can be interpreted as the (undesired) residual correlation in the detection filter due to $`\mathrm{\Omega }_2(f)`$, which can not be eliminated. When it exceeds unity, the component $`\mathrm{\Omega }_1`$ can not be detected, regardless of the instrumental sensitivity; observations carried out with smaller noise disturbances, would simply increase both terms $`(\stackrel{~}{Q}_1,\stackrel{~}{Q}_1)`$ and $`(\stackrel{~}{Q}_1,\stackrel{~}{Q}_2)`$ by exactly the same amount, without improving the chances of detecting $`\mathrm{\Omega }_1`$. The same holds for the integration time $`T`$: it has no effect at all on the capability of discriminating the two components. Therefore, the minimum detectable value of $`\mathrm{\Omega }_1(f)`$ is set by the condition
$$(\stackrel{~}{Q}_1,\stackrel{~}{Q}_1)=(\stackrel{~}{Q}_1,\stackrel{~}{Q}_2).$$
(67)
Notice that the frequency dependence of the two components is very important: if $`\mathrm{\Omega }_1(f)`$ and $`\mathrm{\Omega }_2(f)`$ follow a similar frequency behavior, the filter picks up more power from the ”spurious” component $`\mathrm{\Omega }_2(f)`$; if they are drastically different, even if $`\mathrm{\Omega }_2(f)`$ dominates $`\mathrm{\Omega }_1(f)`$, one could achieve a detection.
This example describes exactly the issue that we are considering in this section, by simply identifying $`\mathrm{\Omega }_1`$ with $`\mathrm{\Omega }_p`$, and $`\mathrm{\Omega }_2`$ with $`\mathrm{\Omega }_{g,is}`$. The unresolved radiation from binary systems provides therefore a fundamental sensitivity limit in searching for the primordial GW background. We have computed this limit from Eq. (67), in the case of an experiment carried out with a pair of identical LISA’s, assuming $`\mathrm{\Omega }_p`$ constant, and $`\mathrm{\Omega }_{g,is}(f)`$ given by Eq. (54). The results are summarized in Fig. 5; the key conclusion is that two LISA detectors will be able to detect a PGB (with constant energy spectrum) only if $`h_{100}^2\mathrm{\Omega }_p\stackrel{>}{}5\times 10^{13}`$. The big loss of SNR with respect to the case where the experiment is limited only by instrumental noise is due to the fact that the GGB is very strong in the mHz band, and $`\mathrm{\Omega }_p(f)`$ and $`\mathrm{\Omega }_g(f)`$ have a similar decreasing frequency behavior in the frequency window where most of the SNR is accumulated: the residual correlation at the filter output produced by the GGB is large.
Given the results reported in Fig. 5, it is straightforward to conclude that the GGB is a guaranteed, strong GW signal for space-based detectors. In fact, if one constructs a filter matched to a GGB given by Eq. (54), and performs cross-correlations between two identical LISA instruments which are characterized by the noise curve (34), one can detect such a signal with $`\mathrm{SNR}_\mathrm{g}100`$, for a time of observation $`T=10^7`$ sec; see Fig. 6; two to three days of integration time are sufficient to reach $`\mathrm{SNR}_\mathrm{g}10`$. Two LISA-like detectors would therefore be extremely powerful telescopes to launch deep surveys of populations of binary systems in our Universe with periods between a few hours and a few hundred seconds.
The bottom line of this analysis is therefore clear: the fundamental limiting factor in searching for a primordial GW background in the $`10^5\mathrm{Hz}10^2\mathrm{Hz}`$ frequency window is the stochastic radiation from unresolved binary systems. Searches for PGBs with $`h_{100}^2\mathrm{\Omega }_p\stackrel{<}{}5\times 10^{13}`$ thus call for a change in the observational window. We briefly discuss this issue in the next Section.
### B Towards testing slow-roll inflation
The mHz frequency window is unsuitable to reach the ambitious sensitivity level $`h_{100}^2\mathrm{\Omega }_p^{(\mathrm{min})}10^{16}`$ predicted by slow roll inflation. One needs to design an experiment with optimal sensitivity in a band free from generated backgrounds. Given our present astrophysical understanding, the most promising region seems to be $`0.1\mathrm{Hz}1\mathrm{Hz}`$, which would be optimally accessible through space-borne interferometers with arms shorter than the LISA ones by a factor $`100`$. In fact, the entire frequency window from $`10^7`$ Hz to $`0.1`$ Hz is contaminated by stochastic signals of astrophysical origin with $`h_{100}^2\mathrm{\Omega }10^{16}`$. A background generated by massive black hole binary systems with fractional energy density $`h_{100}^2\mathrm{\Omega }10^{15}10^{14}`$ is present in the $`\mu `$Hz range ; however, our ignorance concerning the formation rate of massive black holes and the merger rate of massive black hole binary systems in the range $`10^5M_{}10^9M_{}`$ prevents us from giving a more solid estimate of the background generated by such objects. For $`10^5\mathrm{Hz}\stackrel{<}{}f\stackrel{<}{}10^2`$ Hz, unevolved binaries and WD-WD binary systems completely swamp the observational window; see Sec. III. Above $`10`$ mHz the only residual sizeable contribution comes from NS-NS binary systems; as we have discussed in Sec. III, above $`0.1`$ Hz, the number of sources per frequency bin becomes less than one, and the sky is “transparent” to a primordial signal. For rather long integration times $`3`$ yrs, the rms instrumental noise level that is required to test the prediction from slow-roll inflationary models is of the order $`10^{24}`$
$$h_{100}^2\mathrm{\Omega }_p^{(\mathrm{min})}8\times 10^{17}\left(\frac{f}{0.1\mathrm{Hz}}\right)^{3/2}\left(\frac{T}{10^8\mathrm{sec}}\right)^{1/2}\left[\frac{h_{\mathrm{rms}}}{10^{24}}\right]^2.$$
(68)
Operating at considerably higher frequency than LISA, the two detectors would have to be closely located, $`D\stackrel{<}{}10^{11}`$ cm, in order to have optimal overlap reduction function; however, this also increases potentially correlated noise sources. Clearly the technological challenge to achieve the mentioned sensitivity is considerable; the main noise sources that would degrade the performance of such a detector are the shot noise, beam pointing fluctuations, and accuracy of the phase measurement technique. This imposes stringent requirements on the power and frequency of the laser, as well as on the dimensions of the ”optics” and on other components of the instrument. We would also like to stress that the value of $`h_{\mathrm{rms}}`$ that we have quoted in (68) is the effective noise fluctuation in the data stream, after the spectral lines from individual NS-NS binaries, that are still copiously present in the observation band, have been removed. How this can be effectively done and what instrument sensitivity is required is still an open question that requires a careful analysis . The main result of this crude analysis is that one could indeed reach the target sensitivity, and the fundamental limitations that make the mHz band unsuitable are removed. However, several questions remain open for future analysis: the discussion of these issues goes far beyond the purpose of the present paper.
## V Conclusions
Gravitational wave experiments in the low-frequency window, together with high-frequency ground-based interferometers, are expected to improve our picture of the very early Universe, and the understanding of the behavior of fundamental fields at high energy, by detecting, or setting stringent upper limits on the primordial background of gravitational radiation.
In this paper we have analyzed the sensitivity of space-borne laser interferometers of the LISA class, and possible succeeding missions. In order to set a reference frame for this discussion, we have regarded the detection of a GW background produced during the early Universe of energy $`h_{100}^2\mathrm{\Omega }_p10^{16}`$, consistent with the prediction of standard slow-roll inflation, as the goal of GW cosmology. We have assumed the operation of two space-detectors, in order to achieve the best sensitivity and detection confidence, and we have shown that the technology available for LISA already ensures the detection of a GW background as weak as $`h_{100}^2\mathrm{\Omega }5\times 10^{14}`$. However, the strong stochastic signal in the mHz band due to short-period solar-mass binary systems that can not be resolved as individual sources prevents us from detecting a primordial background weaker than $`h_{100}^2\mathrm{\Omega }_p5\times 10^{13}`$. Astrophysically generated stochastic backgrounds therefore set a fundamental limit in the mHz band that prevents us from achieving a sensitivity that goes beyond what is already guaranteed by the LISA technology. They also represent a guaranteed strong signal detectable at high signal-to-noise ratio, which enables the study of the distribution and merger rate of populations of binary compact objects in the Universe.
Dedicated missions with optimal sensitivity in the window 0.1 Hz – 1 Hz appear, at present, the only viable option in the search for very weak primordial backgrounds, and we have briefly discussed the technological challenges involved in probing slow-roll inflation. Our order-of-magnitude analysis strengthens the hope that a sensitivity level $`h_{100}^2\mathrm{\Omega }_p10^{16}`$ might be within the capability of future dedicated low-frequency detectors.
Our analysis clearly indicates the key issues that deserve further investigation: a solid estimate of galactic and extra-galactic GW backgrounds produced by astrophysical sources, the investigation of the statistical issues that can lead to the discrimination of the PGB from the GGB, and a more rigorous analysis of the technical and conceptual problems for low-frequency experiments dedicated to GW cosmology. On the observational side, the presently designed single-instrument LISA mission is a fundamental step for the planning of more ambitious, multi-detector experiments: we will be able to measure directly the degree of anisotropy of the generated background, shedding light on the fundamental limiting factor of mHz experiments. In fact, while the present paper deals only with the detection of an isotropic stochastic signal, the remarkable sensitivity of LISA offers the chance of going far beyond: a detailed study of the anisotropy and angular dependence of stochastic signals, both of astrophysical and primordial origin. Such an investigation is currently in progress, and will be reported in a separate publication .
###### Acknowledgements.
We thank C. Cutler and B. Schutz for their help and encouragement throughout this work. We especially thank P. Bender for sharing his latest thoughts on generated stochastic backgrounds, and for several comments on a preliminary version of this paper. We would also like to thank K. Danzmann for illuminating discussions regarding LISA technology and the noise sources at low frequencies. This work has also benefited from interactions with B. Falkner and S. Vitale.
## A
We give here the values of the numerical coefficients that enter expressions (47), (48) and (48) of the overlap reduction function derived in Sec. (II B):
$`A_0`$ $`=`$ $`{\displaystyle \frac{513}{4096}},A_1={\displaystyle \frac{135}{1024}},A_2={\displaystyle \frac{243}{2048}},A_3={\displaystyle \frac{9}{1024}},A_4={\displaystyle \frac{33}{4096}};`$ (A1)
$`B_0`$ $`=`$ $`{\displaystyle \frac{27\sqrt{3}}{4096}},B_1={\displaystyle \frac{27\sqrt{3}}{1024}},B_2={\displaystyle \frac{81\sqrt{3}}{2048}},B_3={\displaystyle \frac{27\sqrt{3}}{1024}},B_4={\displaystyle \frac{27\sqrt{3}}{4096}};`$ (A2)
$`C_0`$ $`=`$ $`{\displaystyle \frac{513}{8192}},C_1={\displaystyle \frac{171}{2048}},C_2={\displaystyle \frac{99}{4096}},C_3={\displaystyle \frac{27}{2048}},C_4={\displaystyle \frac{33}{8192}};`$ (A3)
$`D_0`$ $`=`$ $`{\displaystyle \frac{27\sqrt{3}}{8192}},D_1={\displaystyle \frac{9\sqrt{3}}{2048}},D_2={\displaystyle \frac{63\sqrt{3}}{4096}},D_3={\displaystyle \frac{9\sqrt{3}}{2048}},D_4={\displaystyle \frac{27\sqrt{3}}{8192}};`$ (A4)
$`E_0`$ $`=`$ $`{\displaystyle \frac{513}{16384}},E_1={\displaystyle \frac{207}{4096}},E_2={\displaystyle \frac{339}{8192}},E_3={\displaystyle \frac{63}{4096}},E_4={\displaystyle \frac{33}{16384}};`$ (A5)
$`F_0`$ $`=`$ $`{\displaystyle \frac{27\sqrt{3}}{16384}},F_1={\displaystyle \frac{45\sqrt{3}}{4096}},F_2={\displaystyle \frac{177\sqrt{3}}{8192}},F_3={\displaystyle \frac{45\sqrt{3}}{4096}},F_4={\displaystyle \frac{27\sqrt{3}}{16384}}.`$ (A6)
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# Radio Emission from GRO J1655–40 during the 1994 Jet Ejection Episodes
## 1 Introduction
The X-ray transient GRO J1655$``$40 (Nova Sco 1994) was first detected with the Burst and Transient Source Experiment (BATSE) on board the Compton Gamma-Ray Observatory on 1994 July 27 (Zhang et al. (1994)). Significant flux was observed up to 200 keV and the source reached a luminosity of 1.1 Crab (20–100 keV) on Aug 1. GRO J1655$``$40 remained in outburst until about Aug 15, and after a period of quiescence flared again on Sept 6 (Harmon et al. (1995)). The detection of a steeply-rising radio counterpart was first reported by Campbell-Wilson & Hunstead (1994a) following observations with the Molonglo Observatory Synthesis Telescope (MOST) at 843 MHz on 1994 Aug 6 and 11. The flux density continued to increase, reaching 4.2 Jy on 1994 Aug 14 and 5.5 Jy on 1994 Aug 15 (Campbell-Wilson & Hunstead 1994b ). High resolution radio observations using the Very Large Array (VLA), the Very Long Baseline Array (VLBA) (Hjellming & Rupen 1995, hereafter HR95) and the Southern Hemisphere VLBI Experiment (SHEVE) array (Tingay et al. (1995)) showed repeated episodes of relativistic ejections. Three major ejection events on TJD 9577.5, 9584 and 9596 (TJD defined as JD $``$ 2440000.5) were observed with the VLA and three on TJD 9574$`\pm `$1, 9605$`\pm `$2 and 9668$`\pm `$5 were recorded with the VLBA (HR (95)). The ejection velocity, after correction for inclination to the line of sight, was inferred to be $`v0.92c`$ in the rest frame of GRO J1655$``$40, at a position angle of $`47^{}\pm 1^{}`$ (HR (95)); wiggles about the jet axis were interpreted as precession with a period of $`3.0\pm 0.2`$d. Photometric observations showed GRO J1655$``$40 to be a high-inclination binary system (Bailyn et al. (1995)). Subsequent optical observations led to the determination of a precise mass for the primary, $`M_1=7.02\pm 0.22M_{}`$ (Orosz & Bailyn (1997)), which is well above the theoretical upper limit for a neutron star and direct evidence for it being a black hole. In addition, Soria et al. (1998) found a 95% confidence limit of $`M_1>\mathrm{\hspace{0.17em}5.1}M_{}`$ for the mass of the primary based on measurements of velocity variations in the He ii disk emission lines, thought to reflect the orbital motion of the primary, and hence confirming GRO J1655$``$40 to be a black hole binary. The mass of the secondary, which is classified as F3 IV–F6 IV, is $`M_2=2.34\pm 0.12M_{}`$, and the spectroscopic period of the binary is $`P=2^d.62157\pm 0^d.00015`$ (Orosz & Bailyn (1997)).
During the 1994 ejection events, the radio outbursts from GRO J1655$``$40 were monitored in the southern hemisphere by the Molonglo Observatory Synthesis Telescope (MOST), the Australia Telescope Compact Array (ATCA), and the Hartebeesthoek Radio Astronomy Observatory (HartRAO). We discuss below the results of these observations. Sections 2 and 3 describe the MOST, ATCA and HartRAO observations and the resulting light curves. In Section 4 we present the linear polarization data from the ATCA, and we discuss the evolution of the radio spectra in Section 5. We briefly compare our results with those of HR95 in Section 6. Finally, we interpret the polarization data in terms of a ‘core-lobe’ model in Section 7.
## 2 Observations and data reduction
### 2.1 MOST
The MOST is an east-west synthesis array comprising two colinear cylindrical paraboloids each 11.6 m wide by 778 m long, separated by a 15 m gap (Mills (1981); Robertson (1991)). The telescope operates at 843 MHz, with a detection bandwidth of 3.25 MHz, and measures right circular polarization (as defined by the IEEE standard). The synthesized beamwidth is $`43\mathrm{}\times 43\mathrm{}\mathrm{cosec}|\delta |`$ FWHM (RA $`\times `$ Dec). The telescope forms a comb of 64 real-time fan beams spaced by $`22\mathrm{}`$; for the observations of GRO J1655$``$40 the pointing of the beam set was time shared among three adjoining positions to give a field size of $`70\mathrm{}\times 70\mathrm{}`$ cosec$`|\delta |`$. The background noise in a full 12-hour synthesis image was typically 1.5 mJy rms.
Many GRO J1655$``$40 observations were partial syntheses as the observations were fitted in around a previously scheduled observing program. The choice of field center serendipitously included a strong point source ($``$ 900 mJy, J2000 position $`16^\mathrm{h}53^\mathrm{m}55.9^\mathrm{s},40^\mathrm{°}37^{\mathrm{}}23^{\mathrm{}}`$) that has subsequently been used as an internal reference source. The flux density calibration was based on short fan-beam scan observations of strong point sources from the list of Campbell-Wilson and Hunstead (1994c). For S$`{}_{843\mathrm{MHz}}{}^{}>100`$ mJy the dominant source of error in a synthesis image is the scatter in the pre- and post-observation calibrators; a conservative error of $`\pm 3`$% has been assigned. At lower flux densities the error becomes dominated by noise and, for partial syntheses, residuals arising from incomplete cancellation of strong out-of-field sources. This difficulty can be partly overcome by subtracting a reference image (obtained when GRO J1655$``$40 was quiescent) with the same hour angle coverage.
The data processing used standard MOST reduction software. When the source was bright we were able to search for sample-to-sample (24 s) variations in flux density. Such measurements are affected by confusing sources in MOST’s fan beams, by the effects of ionospheric refraction and by slow drifts in telescope calibration. These effects are eliminated or greatly reduced in a full 12-hour synthesis. For this paper, individual sample flux densities were fitted and averaged in blocks of 200 (80 minutes) and then combined in 4–6-hour blocks for tabulation. Full synthesis imaging was used when the source flux density fell below 500 mJy, with flux densities measured using IMFIT in aips.
In addition to the synthesis-mode observations of GRO J1655$``$40, spot measurements of flux density were obtained using 4-minute scan observations, bracketed by scans of calibrators. This mode of operation measures the target and calibrator source flux densities using the (same) central beams of the 64-beam block. It was used only when the flux density of GRO J1655$``$40 was greater than 1 Jy, with the fit parameters serving as a guide to the presence of confusing sources in the fan-beam response. The final light curve, tabulated in Table 1, contains a mix of flux densities measured from synthesis images, from averaged fits to the synthesis data, and from scan observations, and includes data from a third, weaker outburst which occurred in 1994 Nov (Wu & Hunstead (1997); Hannikainen, Hunstead & Campbell-Wilson (1998)).
### 2.2 ATCA
Following the report of the intense radio outburst by the MOST, GRO J1655$``$40 was monitored as a target-of-opportunity with the ATCA for a period of about twenty days. VLBI observations and H I spectral line observations were also undertaken (McKay & Kesteven (1994); Tingay et al. (1995)).
The ATCA is an earth-rotation aperture synthesis array, comprising six 22-m antennas which can be moved along an east-west railway track to give baselines up to 6 km (Frater, Brooks & Whiteoak (1992)). The observations of GRO J1655$``$40 were made at central frequencies of 1.380, 2.378, 4.800, 5.900, 8.640 and 9.200 GHz with 128 MHz bandwidth in two orthogonal linear polarizations. Each observation was typically 10 minutes duration, with the antenna gain and phase calibration derived from regular observations of the point-source calibrator PKS B1740$``$517. The flux density scale was tied to the primary calibrator PKS B1934$``$638, which we assume has flux densities of 14.96, 11.54, 5.83 and 2.84 Jy at 20, 13, 6 and 3 cm (1.4, 2.3, 4.8, 8.6 GHz) respectively (Reynolds (1994)). The 6.0A configuration was used throughout, giving interferometer spacings from 337–5939 m.
The data were reduced using standard flagging and calibration techniques in the miriad package (Sault, Teuben & Wright (1995)). The polarimetric calibration of ATCA data is described in Sault, Killeen & Kesteven (1991). Polarimetric leakages, which are quite stable with time on the ATCA, were determined simultaneously with the antenna gains using either PKS B1934$``$638 or PKS B1740$``$517, both known to be unpolarized. The resultant polarimetric purity of the calibrated data is believed to be better than 0.1%. At all wavelengths, a single round of phase self-calibration was performed on the GRO J1655$``$40 data (using a point source model) to eliminate residual phase instability. As the observations placed GRO J1655$``$40 at the phase center, we simply summed the real parts of the relevant Stokes visibilities to determine the source flux density; this is equivalent to natural weighting. The errors are dominated by systematic effects. For the total intensity measurements, the errors are smaller than 1% of Stokes I, while for Stokes Q and U they are less than 0.1% of Stokes I. The ATCA total flux densities are tabulated in Table 2.
### 2.3 HartRAO
Continuum measurements of the flux density of GRO J1655$``$40 were made at 3.5 and 6.0 cm using the 26-m telescope of the Hartebeesthoek Radio Astronomy Observatory (HartRAO) situated in Gauteng, South Africa. The observations were made using dual-feed systems operated in a beam-switched Dicke mode. At 3.5 cm (8.58 GHz) the bandwidth was 400 MHz, the beamwidth $`5.7^{}`$ FWHM and the feed recorded right circular polarization. The 6 cm (5 GHz) receiver had a bandwidth of 800 MHz, a beamwidth of $`10^{}`$ FWHM and the feed was linearly polarized east-west.
The flux density at each wavelength was measured using a five-point stepping sequence in declination to correct for pointing errors in this coordinate; the stepping positions in terms of beam parameters were first null, half-power, on-source, half-power, first null. The five points were fitted to a Gaussian to determine declination pointing errors. Right ascension errors were measured at 6 cm by observing midway between the beams. Right ascension pointing corrections at 3.5 cm were assumed to be the same as at 6 cm. The beam separation was sufficiently close ($`15^{}`$ or 1.5 times the beamwidth at 6 cm) that the response of the dual-beam Dicke system varied linearly with right ascension in this region. The observed flux density is an average over both beams, corrected for both declination and right ascension pointing errors. Measurements alternated between the two wavelengths with 10 minutes’ integration on each.
Gain curve corrections, determined by G. Nicolson (private communication), were applied to the data. Flux densities were calibrated to an arbitrary scale using a noise diode, and the calibrators 3C 123, 3C 161, 3C 218 (Hydra A), 3C 274 (Virgo A) and 3C 348 were observed occasionally to put the flux densities on the Ott et al. (1994) scale. The HartRAO data are tabulated in Table 3; we note that on some days the scatter among a closely spaced sequence of measurements is somewhat larger than implied by their formal errors.
## 3 Radio light curves
### 3.1 MOST and HartRAO
The MOST 843 MHz light curve up to the end of 1994 September is plotted in the top panel of Figure 1. The duration of the ATCA observations is also shown, as are the epochs of specific ejection episodes, as reported by HR (95) from their VLA and VLBA data; we labelled these epochs vla1, vla2, vla3 and vlba1, vlba2, corresponding to TJD 9577.5, 9584, 9596 and TJD 9574, 9605, respectively. The BATSE 20–100 keV light curve (provided by S. N. Zhang and W. S. Paciesas) is plotted in the bottom panel, showing two major hard X-ray outbursts, beginning on TJD $`9560`$ and TJD $`9600`$.
Two radio outbursts were observed with the MOST in 1994 August–September, both apparently associated with the hard X-ray bursts detected by BATSE. The MOST data clearly show the rise of the radio flux density lagging behind the rise of the X-ray intensity in the outbursts. During the first X-ray outburst, the radio flux density was still low (S$`_{843\mathrm{MHz}}`$ $``$ 0.36 Jy) on TJD 9570.60, while the X-ray intensity had already reached its maximum. When the X-ray intensity began to decline, the radio flux density then increased rapidly, reaching a local maximum (S<sub>1</sub>) of 5.45 Jy on TJD 9579.22. After a brief decline, the flux density rose again and reached a new peak of 7.62 Jy on TJD 9582.13, about two weeks after the X-ray peak. At this point the X-ray intensity had already fallen to a much lower level, but was still above the background. The radio peak in the second outburst (S<sub>2</sub>) also lagged the X-ray peak, but the delay was shorter. The maximum intensities in the two X-ray bursts were similar, but the second radio burst was obviously weaker.
#### 3.1.1 Short Duration Radio Burst Component
To emphasize some specific features of the temporal evolution of the radio emission, we plot in Figure 2 (upper panel) the MOST 843 MHz and HartRAO 5 GHz flux densities on a logarithmic scale as a function of time. Closely spaced groups of 5 GHz points are plotted as weighted averages. The 843 MHz flux density curve is characterized by short-duration burst components (marked S<sub>1</sub> and S<sub>2</sub>) superimposed on longer-lived components (L<sub>1</sub> and L<sub>2</sub>) which show an approximately exponential decay. HartRAO did not observe GRO J1655$``$40 until after the end of S<sub>1</sub>, but S<sub>2</sub> was clearly recorded at 5 GHz. The confirmation of the short duration burst component is an important new result, since neither the VLA (HR (95)) nor the ATCA collected data during these periods. From the MOST data, the short duration components, S<sub>1</sub> and S<sub>2</sub>, have e-folding rise times of $`1.1`$ days, and S<sub>2</sub> decays with an e-folding time of $`2.2`$ days. The rising time scale for the long duration component cannot be determined from the data, but L<sub>1</sub> and L<sub>2</sub> have e-folding decay times of $`6.5`$ days and $`8.7`$ days respectively. A third short duration event was observed by MOST peaking on TJD $`9666`$ which decayed with an e-folding time of $`2.6`$ days (Table 1; Wu & Hunstead (1997)). In addition, the flare observed in 1996 had an initial e-folding decay time of $`1.4`$ days, similar in scale to the times quoted above (Hunstead, Wu & Campbell-Wilson (1997)). It is tempting to speculate that event S<sub>2</sub> may be a prompt response to the brief hard X-ray event at TJD 9610. A similar association at TJD 9579 can be made for S<sub>1</sub> but seems less plausible.
### 3.2 ATCA
Figure 3 shows the ATCA Stokes I light curves at all six observing frequencies, again plotted on a log-linear scale to emphasize the overall exponential decay; the MOST light curve is included for comparison.
The ATCA observations were initiated during the decay phase of the S1 short duration burst component. It is clear from Figure 3 that the flux density at all six ATCA frequencies declines simultaneously with the MOST flux density between TJD $`95799580`$ suggesting that the short burst component was present at all frequencies. Following the minimum at TJD $`9580`$ the flux density increases for all six ATCA frequencies, reaching a local maximum of 5.49 Jy at 1.4 GHz on TJD 9581.62; note, however, that the overall maximum flux density at frequencies $`4.8`$ GHz occurs prior to the minimum. For TJD $`>9582`$ there is a steady decline, with time constants similar to MOST’s at the lower frequencies but shorter at the higher frequencies (4.8–9.2 GHz). Other features, such as the dip at the higher frequencies around TJD 9583 and the rise between 9596 and 9598 at 8.6 GHz, are discussed in later sections.
## 4 Linear polarization
Linear polarization was detected at all six ATCA frequencies throughout the epoch of the observations. The fractional polarization at each frequency, $`\mathrm{P}/\mathrm{I}=\sqrt{(\mathrm{Q}^2+\mathrm{U}^2)}/\mathrm{I}`$, expressed as a percentage, is plotted together in Figure 4 as a function of time, and the values are listed in Table 4. Because Figure 4 is rather complex, and difficult to interpret on its own, we also include Figure 5 which shows separate overlaid plots of the polarized and total intensity at each frequency.
The frequency-dependent time evolution of the P/I light curves is reminiscent of the ‘expanding synchrotron bubble’ model (e.g., Hjellming & Johnston (1988); Ball & Vlassis (1993)). This simple model is characterized by an optically thick rising phase in which the flux density peaks first at the highest frequencies, followed by delayed peaks of lower intensity at progressively lower frequencies. The decay at all frequencies has a power-law time dependence, with an exponent linked to the electron energy distribution.
We have marked three apparent ‘bubble-type’ events in the linearly polarized ATCA data, and labelled them E1, E2 and E3 in Figures 4 and 5, using the 8.6 and 9.2 GHz curves for reference. However, it is worth noting that the beamsize of the ATCA is larger than the angular size of GRO J1655$``$40 (HR95) at all frequencies, so we are observing effects which encompass both the central object and the ejecta. As a result, we can only hope to establish qualitative comparison with the sychrotron bubble model. In fact, contrary to the model prediction, Figure 3 shows that the source was optically thin at TJD 9582, the time of the major peak in flux density.
Figures 4 and 5 show that the first event, E1, is characterized by a series of declining peaks in the 8.6 and 9.2 GHz data. The peaks at 4.8 and 5.9 GHz are less well defined and their decay is clearly delayed with respect to the higher frequencies, in qualitative agreement with the expanding bubble model. The delays at 1.4 and 2.3 GHz are sufficiently large that the contributions from successive ejecta blend with one another and produce a broad peak in both P and P/I. It is interesting to note that the polarized flux density (Fig. 5) reached essentially the same maximum value ($`\mathrm{P}0.2`$ Jy) at all frequencies.
The early stages of event E2 also show qualitative agreement with the expanding synchrotron bubble model, with the rise in P/I at the higher frequencies preceding that at lower frequencies. However, the four higher frequencies all reach a maximum together and then fall to a minimum at the same time, behavior which is not consistent with the model. Event E3 is also inconsistent with the model, as the fractional polarizations at 4.8–9.2 GHz increase and decrease simultaneously. We discuss possible interpretations of this behavior in Section 7.
### 4.1 Position angle and rotation measure
Faraday rotation of the plane of polarization has been taken into account by fitting the observed position angles (PA) to the usual equation, $`\mathrm{PA}=\mathrm{PA}_0+(\mathrm{RM})\lambda ^2`$, where PA<sub>0</sub> is the intrinsic position angle, RM is the rotation measure (in rad m<sup>-2</sup>) and $`\lambda `$ is the wavelength. The data were first interpolated to a common epoch, namely that of the 9.2 GHz observations. The position angle was determined from the equation $`\mathrm{PA}=\frac{1}{2}[\mathrm{arctan}(\mathrm{U}/\mathrm{Q})+n\pi ]`$, the last term reflecting the ambiguities of $`n\pi `$ that may arise at each frequency, and then plotted against $`\lambda ^2`$. A weighted least-squares linear fit yielded the intrinsic position angle (y-intercept) and the RM (slope). Since the position angle is an average over the beamsize, this means that it is also an average over the central source and the ejecta.
Figure 6 shows the temporal evolution of the intrinsic polarization position angle and rotation measure. The average position angle between TJD 9579 and 9589 was $`47\mathrm{°}`$, which is consistent with the position angle of the jets (HR (95)). This implies that the magnetic field is perpendicular to the jets, a phenomenon also seen in the parsec-scale relativistic jets from powerful radio-loud AGN (e.g. Bridle & Perley (1984)). The overall constancy in position angle suggests that the magnetic field is well ordered, and maintains basically the same orientation. However, on TJD 9589.5 the position angle drops to $`26\mathrm{°}`$, and after a brief increase to $`43\mathrm{°}`$, remains at $`24\mathrm{°}\pm 4\mathrm{°}`$ during the last three observing epochs.
The rotation measure increases rapidly from $`6`$ rad m<sup>-2</sup> on TJD 9579.5 to $`104`$ rad m<sup>-2</sup> on TJD $`9582`$. It then levels off at $`60`$ rad m<sup>-2</sup> for the rest of the observing period, the only exception being an increase to $`80`$ rad m<sup>-2</sup> on TJD 9598. Its relative constancy after TJD 9582 suggests that $`60`$ rad m<sup>-2</sup> arises from the path through the interstellar medium, while the rapid increase prior to TJD 9582 probably reflects local effects arising at or near the source. The low values of rotation measure around TJD 9579.5 imply that the magnetic field at the source is in the opposite direction to that of the interstellar medium, and the increase to $`104`$ rad m<sup>-2</sup> indicates that the orientation of the magnetic field has changed direction and is now more-or-less aligned with that in the interstellar medium.
While Faraday effects predict that the rotation of the plane of polarization is proportional to $`\lambda ^2`$, the validity of the $`\lambda ^2`$ dependence has been questioned by O’Dea (1989). He presents polarization observations of fifteen core-dominated quasars, arguing that poor fits to the PA vs. $`\lambda ^2`$ data for six of the quasars cannot be due to large apparent RM values combined with ambiguities of $`n\pi `$, but rather can be accounted for by wavelength-dependent polarization structure. A similar situation may well apply to GRO J1655$``$40. The upper panel of Figure 7 shows a least-squares fit to the polarization data on days TJD 9588.60 and 9590.53, which resulted in $`\mathrm{PA}_0=47.5^{}`$, $`\mathrm{RM}=56.4`$ rad m<sup>-2</sup> and $`\mathrm{PA}_0=38.2^{}`$, $`\mathrm{RM}=63.5`$ rad m<sup>-2</sup> respectively. Both fits appear secure.
In contrast, the lower panel of Figure 7 shows an example of PA vs. $`\lambda ^2`$ for GRO J1655$``$40 on TJD 9589.53, where $`\pi `$ has been added to both the 1.4 and 2.3 GHz data points (open circles). A least-squares fit to the data (solid line) yielded $`\mathrm{PA}_0=26.1^{}`$, $`\mathrm{RM}=65.5`$ rad m<sup>-2</sup>, and a large reduced $`\stackrel{~}{\chi }^2`$ of 12.8 for 4 d.o.f. We tried to accommodate the high frequency points by adding $`4\pi `$ and $`2\pi `$ to the 1.4 and 2.3 GHz data respectively (filled circles) and obtained $`\mathrm{PA}_0=5.8^{}`$, $`\mathrm{RM}=282.9`$ rad m<sup>-2</sup> (dotted line) and an improved $`\stackrel{~}{\chi }^2`$ of 5.3. However, these fitted parameters are so inconsistent with the values obtained on TJD 9588.60 and 9590.53 (with $`\stackrel{~}{\chi }^2`$ values of 0.31 and 0.05 respectively) that they seem implausible. Hence, we conclude that the discrepancies seen in the high frequency points on TJD 9589.53, especially those at 8.6 and 9.2 GHz, are due to wavelength dependent polarization structure in the source. Considering that GRO J1655$``$40 is a complex time-variable source, and that the observations are integrations over multiple evolving ejecta, a poor fit to a simple $`\lambda ^2`$ dependence is perhaps to be expected.
## 5 Spectra
### 5.1 Spectral indices
To fully investigate the evolution of the radio spectrum, two sets of spectral indices were considered: the first set covers the entire MOST/HartRAO monitoring period, from TJD $`9580`$–9618, while the second set concentrates on the ATCA data which, due to better sampling, show more detail but for a shorter period of time. Two-point spectral indices $`\alpha `$ were calculated, assuming spectra of the form $`S_\nu \nu ^\alpha `$. To examine the spectral evolution more closely, low- and high-frequency spectral indices were calculated separately. For the MOST/HartRAO data these consisted of the 843 MHz–5 GHz ($`\alpha _{0.85}`$) and 5–8.58 GHz ($`\alpha _{58.5}`$) bands, while for the ATCA observations the data were divided into the 1.4–2.3 GHz ($`\alpha _{1.42.3}`$) and 4.8–8.6 GHz ($`\alpha _{4.88.6}`$) bands. In addition, a ‘global’ spectral index was obtained from the MOST 843 MHz and HartRAO 8.58 GHz data ($`\alpha _{0.88.5}`$) and the MOST and ATCA 9.2 GHz data ($`\alpha _{0.89.2}`$).
The single-dish HartRAO flux densities are less accurate than those from the ATCA (and the VLA) and will be affected by confusing sources in the beam; an image of the field of GRO J1655$``$40 is shown in Hunstead et al. (1997). On the other hand, the HartRAO data span a wider time interval than the ATCA and fill in vital gaps in the VLA light curves. Furthermore, while confusion will affect the spectral indices systematically, trends with time will be preserved.
#### 5.1.1 MOST and HartRAO
The lower panel in Figure 2 shows the spectral indices from the MOST and HartRAO data plotted as a function of time. The MOST and HartRAO 8.58 GHz data points were interpolated to the HartRAO 5 GHz epochs. Between TJD 9580–9593 and TJD 9610.5–9616 only the global $`\alpha _{0.88.5}`$ spectral index is plotted, while from TJD 9593 to 9610, $`\alpha _{0.85}`$ and $`\alpha _{58.5}`$ are plotted separately. Prior to TJD 9593, $`\alpha _{0.88.5}`$ ranges between $`0.6`$ and $`0.4`$, whereas after TJD 9610 the spectrum is slightly flatter overall with $`\alpha _{0.88.5}0.3`$. At the time of the vla3 ejection event (TJD $`9596`$), the 5–8.58 GHz spectrum inverts to $`\alpha _{58.5}+0.25`$ (with a large uncertainty); this mini-outburst is also seen in the VLA data (HR (95)), most notably at 15 and 22 GHz. The 5–8.58 GHz spectral index then returns to the pre-outburst value at TJD $`9600`$. A second outburst, also recorded by the VLA, occurs a few days later, with $`\alpha _{58.5}`$ peaking at $`+0.4`$ at TJD $`9607`$ and steepening soon thereafter. The inversion of the spectrum, indicating the emergence of an optically thick component, is consistent with the vlba2 ejection event; the precise time correspondence is, however, uncertain because neither the HartRAO nor VLA dataset covers the critical interval TJD 9601–9606.
#### 5.1.2 MOST and ATCA
Figure 8 shows the spectral indices obtained from the MOST/ATCA datasets plotted as a function of time. To calculate the 843 MHz–9.2 GHz index the MOST data points were interpolated to the 9.2 GHz epochs. The temporal evolution of the three datasets is similar, with $`\alpha _{0.89.2}`$ and $`\alpha _{1.42.3}`$ ranging from $`0.4`$ to $`0.6`$, while $`\alpha _{4.88.6}`$ (excluding the last two points at TJD 9596 and 9598) varies between $`0.5`$ and $`0.8`$, indicating that the high-frequency portion of the spectrum is consistently steeper. There is an overall steepening of the spectrum between TJD 9579 and 9582, as the source becomes optically thin at the lower frequencies. Immediately after this, the high-frequency spectral index begins to flatten and is followed by the lower frequencies on TJD $`9583`$; the light curves in Figure 3 shows a clear dip and rise at this epoch. This behavior is consistent with the ejection of an optically thick synchrotron component and points to an ejection epoch TJD $`9582`$, somewhat earlier than inferred for the vla2 event (HR (95)). The spectra are flattest on TJD $`9586`$, after which the spectrum steepens. Between TJD 9596.33 and 9598.07, $`\alpha _{4.88.6}`$ jumps from $``$0.87 to $``$0.21, coincident with the vla3 ejection episode; the corresponding increase in the 8.6 GHz flux density can be seen in Figure 3.
### 5.2 Linear polarization spectra
In Figure 9 we show a montage of overlaid plots of the total (I) and polarized (P) intensity spectra of GRO J1655$``$40. To construct the spectra, the flux densities at all frequencies, including 843 MHz (but without polarization), were interpolated to the epochs of the ATCA 9.2 GHz observations. The P spectra, in particular, are especially valuable as diagnostics of ejection events and their evolution with time. We now discuss these spectra in conjunction with the overlaid fractional polarization plot in Figure 4 and the light curves in Figure 3.
On the first day, TJD 9579.6 (four observations spanning two hours), the P spectrum is strongly inverted and the I spectrum shows spectral flattening at the low frequencies. On the following day TJD 9580.50, the I spectrum has become essentially a power law. The P spectrum shows a marked increase at 2.3 GHz but is still optically thick at 1.4 GHz. By TJD 9581.64, near the peak in total intensity, the P spectrum has flattened, with signs of an upturn at the highest frequencies signalling another ejection (E1). The following four P spectra, spanning TJD 9582.17–9583.59, trace a steady increase in the polarized flux density at the lower frequencies, a pattern which is qualitatively consistent with the expansion of the synchrotron-emitting region(s).
At TJD 9584.63 the P spectrum is still not completely transparent at the lower frequencies but has steepened markedly at the higher frequencies, with signs of a new component (E2) appearing at 8.6 and 9.2 GHz. The development of this new component over the next three days is again consistent with the expected trend for an expanding, optically-thick synchrotron bubble. Over this same time interval the low frequency P spectrum has become optically thin. Between TJD 9587.62 and 9588.62 we interpret the persistently high level of P for frequencies $`4.8`$ GHz as indicating yet another ejection event, even though there is no perceptible change in the total flux density.
From TJD 9588.62 to 9589.55 there is an interesting transition. At all frequencies the total flux density from TJD 9589.55–9591.37 lies well above an extrapolation of the exponential decay seen over the previous four days (Fig. 3), implying that a flat-spectrum component has been introduced. A simultaneous flux density increase across the spectrum is not a characteristic of the synchrotron bubble model, so an alternative interpretation must be sought.
Accompanying the flux density increase at TJD 9589.55 is a sharp decrease in the high-frequency polarized flux density and fractional polarization which occurs simultaneously at the four frequencies $`4.8`$ GHz. The high-frequency polarized flux ‘reappears’ on TJD 9590.55 and 9591.37 but is quenched again in the same fashion two days later on TJD 9593.36. We put forward a speculative but self-consistent explanation of this behavior in Section 7.
## 6 Comparison with the VLA/VLBA
In Sections 3, 4 and 5 we have described various aspects of the radio behavior of GRO J1655$``$40 as observed with the MOST, ATCA and HartRAO. We now discuss this behavior in the light of the images and light curves recorded with the VLA and VLBA (HR (95)).
The light curves presented here (Figs. 1–3) are completely consistent with those produced at the VLA (HR95), with an outburst peaking on $`\mathrm{TJD}9582`$ followed by a more-or-less exponential decay. Due to sparser temporal coverage the VLA did not record the two short duration outbursts clearly observed with MOST (S<sub>1</sub> and S<sub>2</sub>) and HartRAO (S<sub>2</sub> only).
As discussed in Section 5.2, the ATCA linear polarization spectra (Fig. 9) seem to give the clearest indication of plasmon ejection and expansion. The E1 event is probably composed of several short-lived contributions, beginning at $`\mathrm{TJD}9580`$, and possibly associated with ejection event vla1 on TJD 9578. The vla2 event (HR (95)), originating at TJD 9584, is not well defined because of the 6-day gap in VLA coverage from TJD 9583 to 9589. It was noted previously (Section 5.1.2) that the overall spectral indices pointed towards a major outburst beginning near TJD 9582. This interpretation is supported by our polarization spectra which suggest that events beginning near TJD 9581.6 and 9584.6 may be confusing the interpretation of the component positions recorded later by the VLA.
Both the vla3 and vlba2 events are well traced and supported by our spectral index plots in Figure 2 and Figure 8. It is worth noting, however, that the steepness of the high-frequency P and I spectra in the last panel in Figure 9 suggests that the vla3 event began after TJD 9596.32.
## 7 Physical Interpretation
The radio emission from candidate black hole binaries is usually attributed to optically thin synchrotron emission or thermal free-free emission. The former explains the steep high frequency part of the power-law spectrum, while the latter is invoked if the spectrum is flat, extending to high frequency.
### 7.1 Polarization and the Synchrotron Bubble Model
In the simple synchrotron bubble model (e.g. Hjellming & Johnston 1988; Ball & Vlassis 1993) the emission region is a homogeneous spherical, expanding plasma cloud, composed of relativistic electrons and embedded in a magnetic field. As the cloud (bubble) expands adiabatically, the synchrotron emission from it becomes optically thin at progressively lower frequencies. A natural consequence, therefore, is that the total intensity I and fractional polarization P/I will peak first at high frequency and later, with smaller amplitude, at lower frequencies.
How well does this model describe the observed properties of GRO J1655$``$40? To begin with, the radio emission is strongly polarized, giving direct evidence for the presence of synchrotron emission (Fig. 4). During event E1, P/I peaks at 9.2 and 8.6 GHz ahead of 5.9 GHz and 4.8 GHz. In addition, the maximum values of P/I at the high frequencies are larger than those at the lower frequencies, the only exception being the 1.4 GHz P/I which is explained as a blend of contributions from several outbursts. Event E2 displays a similar pattern. While these characteristics are in qualitative agreement with the synchrotron bubble model, after event E2 the P/I values at 4.8–9.2 GHz reached a minimum at almost the same time, which is not predicted. Thereafter, during event E3, P and P/I rise and fall simultaneously at 4.8–9.2 GHz, again contrary to model predictions. The polarization spectra (Section 5.2, Figure 9) paint a similar picture. Lack of compliance with the synchrotron bubble model is perhaps not surprising considering the complex collimated structures revealed by the VLBA observations (HR95) and the simplicity of the model itself.
### 7.2 “Core-Lobe” Model
As polarization is detected, we cannot dismiss the contribution from synchrotron emission. However, we have shown that the synchrotron bubble model cannot provide a satisfactory explanation for the time evolution of the total intensity spectrum or the polarization in GRO J1655$``$40. As a modification, we propose a new model which takes into account the fact that VLBI observations of GRO J1655$``$40 have shown multiple, time-varying emission regions (HR 95, Tingay et al. 1995). A hybrid model which brings in other emission mechanisms is therefore worth exploring for this very complex system.
There are two examples in the literature where multiple components have been invoked to explain the light curves and/or polarization: V404 Cyg (Han & Hjellming (1992)) and GRS 1915$`+`$105 (Fender et al. 1999).
* The X-ray transient, V404 Cyg (Han & Hjellming (1992)), shows time-dependent polarization properties similar to those seen in GRO J1655$``$40, including little variation in polarization position angle over the 50-day monitoring period. However, in contrast to GRO J1655$``$40, V404 Cyg showed a rapid initial decay, in which the spectrum was optically-thin, followed by a much slower decline during which the radio spectrum remained flat or inverted. Han & Hjellming (1992) interpret this behavior as requiring at least two radiating components.
* During a series of major ejection events in the superluminal jet X-ray binary, GRS 1915$`+`$105, recorded at high angular resolution with MERLIN, Fender et al. (1999) reported that the (stationary) core did not show any significant polarization, even though the approaching component was significantly polarized. They note that the core emission was dominated by flat spectrum oscillations and attribute the lack of polarization to either the superposition of multiple components with different polarization position angles, or to large Faraday depolarization close to the binary system.
In our hybrid model for GRO J1655$``$40 we assume that the emission regions consist of a ‘core’ (an extended plasma cloud surrounding the central source) and ‘lobes’ (the expanding plasmons or ejecta). The core represents the ejecta when they are localized near the binary; it is therefore dense and compact, and unresolved in the VLBA image. For a beamwidth $``$15 mas and an estimated distance of $`3`$ kpc (HR95), we estimate that the corresponding size of the emitting region should be $`10^{14}`$ cm. Synchrotron emission is inefficient in the core, because of strong self-absorption. During the initial stage of ejection, the core is most compact and the electron density is highest, making free-free emission a very efficient process. The electrons in the ejecta are presumably accelerated by some non-thermal (magnetic) processes near the black hole, so their energy distribution will have a non-Maxwellian component giving rise to non-thermal free-free emission. If the density in the core is sufficiently high to allow particle collisions, a sub-population of electrons with a Maxwellian distribution will result. These electrons will emit thermal free-free emission.
When the ejecta leave the immediate environment of the binary and expand to form the lobes, which are then resolved in the radio images, they eventually become transparent to synchrotron emission. Moreover, when the density drops, free-free emission becomes less efficient, leaving synchrotron emission, which is no longer self absorbed, as the dominant mechanism for energy loss.
According to this model, events E1 and E2 can be understood as two consecutive episodes of mass ejection. During the onset of the ejection, the lobes and the core are not actually separable. As the emission region is compact, the synchrotron emission is self-absorbed. There is also a free-free component. Hence the overall spectrum shows a turn-over, with the critical frequency determined jointly by the emissivities of the synchrotron and free-free processes. When the individual plasma bubbles begin to separate from the core and expand, their free-free emission becomes unimportant. Moreover, they become progressively transparent to synchrotron emission, first at the high frequencies and then at the low. The emission from the compact core, on the other hand, may still be dominated by free-free emission, which is weak in comparison with the synchrotron emission from the expanding lobes.
Free-free emission has a flat spectrum, and so any free-free emission from the core does not contribute significantly to the low frequency part of the overall spectrum. For most of the time, therefore, the spectral properties and polarization at low frequency are characterized by the synchrotron emission from the ejecta, especially when that emission peaks. However, if the synchrotron component has declined substantially, or there is a substantial brightening in the emission from the core, the spectral index and the fractional polarization will both be affected, becoming observable first at the high frequencies and then at the lower frequencies. The simultaneous occurrence of the minima in P/I between events E2 and E3, together with the simultaneous rise and fall in P/I at the four higher frequencies during event E3, can be readily explained by dilution of the synchrotron emission by a varying free-free contribution from the core. The magnetic field disruption resulting from core-brightening episodes may also explain the apparently anomalous polarization position angles (Fig. 6) at TJD 9589.55, 9593.36, 9596.32 and 9598.07.
A free-free core component can also explain the polarization spectra for $`\mathrm{TJD}>9588.6`$. Consideration of the timescales evident in Figure 9, and the VLBI structures (HR (95), Tingay et al. (1995)), tells us that the polarized emission at low frequency comes from regions well away from the core, whereas the high-frequency polarized emission is concentrated close to the core. A sudden increase in the size or electron density in the core will therefore have a significant impact on P at high frequency and little or no effect at low frequency, which seems to explain qualitatively what occurs in the transition between TJD 9588.62 and 9589.55 (Fig. 9). We note that a small but significant upturn in P at TJD 9589.55 between 8.6 and 9.2 GHz probably signals a new ejection event. The following days, TJD 9590.55 and 9591.37, see an increase in the polarized flux at high frequency, presumably because the ejecta have moved outside the region of core absorption. Interpretation of the last two panels in Figure 9 is complicated by the poorer time sampling, but it is plausible that a similar series of events follows the minimum in P and P/I at TJD 9593.36.
Although the hybrid ‘core-lobe’ model has provided a more satisfactory explanation for the radio emission from GRO J1655$``$40 during the 1994 ejection episodes, we emphasize that the model is only qualitative at this stage. Further work is necessary to quantify the model, so that fits to the data can be carried out.
## 8 Summary and conclusions
As Figure 1 shows, there is a relationship between the hard X-ray and radio emission from GRO J1655$``$40. The ejection episodes traced by the VLBA both originated during enhanced activity in the hard X-rays, and preceded the radio outbursts recorded with the MOST, suggesting a connection between activity near the event horizon of the black hole and the production of relativistic electrons. One possible explanation for the decline in intensity of the radio outbursts with time could be that the first ejection occurred in an environment that was relatively undisturbed, whereas the subsequent outbursts will have taken place in an environment already disrupted by previous activity (Hjellming et al. (1996)). The implications of the X-ray/radio correlations have been discussed elsewhere (e.g. Harmon et al. (1995)), with a general consensus that the hard X-rays may be indicating enhanced accretion near the black hole which, through processes still not well understood, triggers the formation of relativistic radio jets.
In general the ATCA, MOST and HartRAO flux density light curves agree well with the VLA light curves and ejection epochs reported in HR95. However, the better time sampling in Figures 2 and 3 reveals the presence of short-lived events that were not recorded by the VLA or VLBA, and shows that the light curves do not decay as smooth exponentials. The radio spectra, especially the linear polarization spectra, have proved to be valuable diagnostics of the timing of plasmon ejection events and their subsequent evolution.
The ATCA polarization data show that the magnetic field is aligned at right angles to the radio jets (jet $`\mathrm{PA}=47\pm 1^{}`$), except towards the end of the monitoring period when core contributions may have become important. The rotation measure initially shows a contribution local to the GRO J1655$``$40 system, but after TJD 9582.5, the rotation measure is roughly constant at $``$60 rad m<sup>-2</sup> which must correspond to the interstellar value.
After examining the time evolution of the total and polarized flux density of GRO J1655$``$40 we conclude that there are specific aspects of the behavior which cannot be explained by the simple synchrotron bubble model. We therefore invoke a hybrid ‘core-lobe’ model, with a core which emits by non-thermal (or maybe thermal) free-free emission and lobes which are classical synchrotron emitters. We suggest that a similar model may apply to the other Galactic superluminal jet X-ray binary, GRS 1915+105.
The Australia Telescope Compact Array is funded by the Commonwealth of Australia for operation as a National Facility managed by CSIRO. MOST is operated by the University of Sydney and funded by grants from the Australian Research Council. KW acknowledges the support of the Australian Research Council through an Australian Research Fellowship. DM acknowledges support for his research by the European Union under contract FMGECT950012, and thanks N.P.F. McKay and M.J. Kesteven for assisting with the ATCA observations. DH acknowledges financial support from the Academy of Finland. The director of HartRAO, G. Nicolson, is thanked for the allocation of observing time that made this project possible. We thank T. Ye for assistance with the MOST data reduction, and J-P Macquart for his invaluable help in checking some of the position angle fits.
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# References
The $`\eta \eta `$ , $`\eta \eta ^{}`$ and $`\eta ^{}\eta ^{}`$ channels are important for the determination of the gluonic content of mesons, since it is thought likely that glueballs will decay with the emission of $`\eta `$s and $`\eta ^{}`$. In addition, central production is proposed as a good place to produce glueballs via Double Pomeron Exchange (DPE) . Hence the interest of the present paper which studies the centrally produced $`\eta \eta `$ system.
The $`\eta \eta `$ channel has been studied previously in radiative $`J/\psi `$ decays by the Crystal Ball experiment , in $`\pi ^{}p`$ interactions by the NA12 experiment , in central production by NA12/2 and in $`p\overline{p}`$ annihilations by the Crystal Barrel Collaboration and E760 at Fermilab . In all these reactions evidence for the $`f_0(1500)`$ has emerged. In addition, in $`p\overline{p}`$ annihilation evidence is claimed for a $`f_0(2100)`$ and an $`\eta \eta `$ decay mode of the $`f_2(1950)`$ . In $`\pi ^{}p`$ interactions the NA12 experiment claim a $`f_0(2000)`$ . In radiative $`J/\psi `$ decays evidence was reported for the $`f_J(1710)`$ . In central production the NA12/2 experiment claimed evidence for a $`f_2(2175)`$ .
In this paper a study is presented of the $`\eta \eta `$ final state formed in the reaction
$$ppp_f(\eta \eta )p_s$$
(1)
at 450 GeV/c. It represents more than a factor of 6 increase in statistics over the only other data on the centrally produced $`\eta \eta `$ final state and moreover will present the first partial wave analysis of this channel in central production. The data come from the WA102 experiment which has been performed using the CERN Omega Spectrometer, the layout of which is described in ref. . Reaction (1) has been isolated using the following decay modes:
$$\begin{array}{cc}\eta \gamma \gamma & \eta \gamma \gamma \\ \eta \gamma \gamma & \eta \pi ^+\pi ^{}\pi ^0\\ \eta \pi ^+\pi ^{}\pi ^0& \eta \pi ^+\pi ^{}\pi ^0\end{array}$$
The above decay modes account for 38.8 % of the total.
Fig. 1a) shows a scatter plot of $`M(\gamma \gamma )`$ versus $`M(\gamma \gamma )`$ which has been extracted from the sample of events having two outgoing charged tracks and four $`\gamma `$s reconstructed in the GAMS-4000 calorimeter using momentum and energy balance. A clear signal of the $`\eta \eta `$ channel can be observed. Fig. 1b) shows the $`\gamma \gamma `$ mass spectrum if the other $`\gamma \gamma `$ pair is compatible with being an $`\eta `$ (0.48 $``$ M($`\gamma \gamma `$$``$ 0.62 GeV) where a clear $`\eta `$ signal can be observed. The $`\eta \eta `$ final state has been selected using the mass cuts described above.
Fig. 1c) shows a scatter plot of $`M(\gamma \gamma )`$ versus M($`\pi ^+\pi ^{}\pi ^0)`$ for the sample of events having four outgoing charged tracks and four $`\gamma `$s reconstructed in the GAMS-4000 calorimeter after imposing momentum and energy balance. A clear signal of the $`\eta \eta `$ channel can be observed. Fig. 1d) shows the $`\pi ^+\pi ^{}\pi ^0`$ mass spectrum if the $`\gamma \gamma `$ mass is compatible with being an $`\eta `$ (0.48 $``$ M($`\gamma \gamma `$$``$ 0.62 GeV) and fig 1e) shows the $`\gamma \gamma `$ mass spectrum if the $`\pi ^+\pi ^{}\pi ^0`$ mass is compatible with being an $`\eta `$ (0.52 $``$ M($`\pi ^+\pi ^{}\pi ^0`$$``$ 0.58 GeV). The $`\eta \eta `$ final state has been selected by requiring that 0.48 $``$ M($`\gamma \gamma `$$``$ 0.62 GeV and 0.52 $``$ M($`\pi ^+\pi ^{}\pi ^0`$$``$ 0.58 GeV.
Fig. 1f) shows a scatter plot of $`M(\pi ^+\pi ^{}\pi ^0)`$ versus $`M(\pi ^+\pi ^{}\pi ^0)`$ for the the sample of events having six outgoing charged tracks and four $`\gamma `$s reconstructed in the GAMS-4000 calorimeter after imposing momentum and energy balance. A signal of the $`\eta \eta `$ channel can be observed. Fig. 1g) shows the $`\pi ^+\pi ^{}\pi ^0`$ mass spectrum if the other $`\pi ^+\pi ^{}\pi ^0`$ combination has a mass compatible with being an $`\eta `$ (0.52 $``$ M($`\pi ^+\pi ^{}\pi ^0`$$``$ 0.58 GeV). The $`\eta \eta `$ final state has been selected using the mass cuts described above.
The background below the $`\eta `$ signal has several sources including combinatorics, fake gammas and other channels. The combinatorial background is removed, in part, in the selection procedure. The remaining background varies from 16 % to 32 % dependent on the decay topology. Three methods have been used to determine the effects of this background; studying the side bands around the $`\eta `$ signal, studying events that do not balance momentum and using event mixing. All three methods give a similar background which peaks near threshold in the $`\eta \eta `$ mass spectrum. To illustrate this point, superimposed on figs. 1b), d), e) and g) as a shaded histogram are the respective mass distributions when the central system has a mass greater than 1.3 GeV. As can be seen the background below the $`\eta `$ signal is reduced with respect to the total sample. In the remainder of this paper the method used to determine the background will be the one using events that do not balance momentum. The resulting distributions well represent the background in the data.
The $`\eta \eta `$ mass spectra from each decay mode are very similar and the combined mass spectrum is shown in fig. 2a) and consists of 3351 events. Superimposed on the mass spectrum as a shaded histogram is the estimate of the background. Fig. 2b) shows the background subtracted mass spectrum. The mass spectrum has a threshold enhancement and has peaks at 1.5 and 2.1 GeV and a shoulder in the 1.7 GeV region. Superimposed on fig 2b) as a shaded histogram is the expected contribution from the $`f_2(1270)`$ (47 events) and $`f_2^{}(1525)`$ (3 events) using the observed $`f_2(1270)`$ signal in the $`\pi ^+\pi ^{}`$ final state , the observed $`f_2^{}(1525)`$ signal in the $`K^+K^{}`$ final state and correcting for the experimental acceptance and using the PDG branching ratios. As can be seen the majority of the signal at 1.5 GeV is not due to the $`f_2^{}(1525)`$.
A Partial Wave Analysis (PWA) of the centrally produced $`\eta \eta `$ system has been performed assuming the $`\eta \eta `$ system is produced by the collision of two particles (referred to as exchanged particles) emitted by the scattered protons. The $`z`$ axis is defined by the momentum vector of the exchanged particle with the greatest four-momentum transferred in the $`\eta \eta `$ centre of mass. The $`y`$ axis is defined by the cross product of the momentum vectors of the two exchanged particles in the $`pp`$ centre of mass. The two variables needed to specify the decay process were taken as the polar and azimuthal angles ($`\theta `$, $`\varphi `$) of one of the $`\eta `$s in the $`\eta \eta `$ centre of mass relative to the coordinate system described above.
The acceptance corrected moments $`\sqrt{4\pi }t_{LM}`$, defined by
$$I(\mathrm{\Omega })=\underset{L}{}t_{L0}Y_L^0(\mathrm{\Omega })+2\underset{L,M>0}{}t_{LM}Re\{Y_L^M(\mathrm{\Omega })\}$$
have been rescaled to the total number of observed events and are shown in fig. 3. The moments with $`M>2`$ and all the moments with $`L>4`$ (e.g. the $`t_{44}`$ and $`t_{60}`$ moments shown in fig. 3) are small and hence only partial waves with spin $`l=0`$ and 2 and absolute values of spin $`z`$-projection $`m1`$ have been included in the PWA.
The amplitudes used for the PWA are defined in the reflectivity basis . In this basis the angular distribution is given by a sum of two non-interfering terms corresponding to negative and positive values of reflectivity. The waves used were of the form $`J_m^\epsilon `$ with $`J`$ $`=`$ $`S`$ and $`D`$, $`m`$ $`=`$ $`0,1`$ and reflectivity $`\epsilon `$ = $`\pm 1`$. The expressions relating the moments ($`t_{LM}`$) and the waves ($`J_m^\epsilon `$) are given in ref.. Since the overall phase for each reflectivity is indeterminate, one wave in each reflectivity can be set to be real ($`S_0^{}`$ and $`D_1^+`$ for example) and hence two phases can be set to zero ($`\varphi _{S_0^{}}`$ and $`\varphi _{D_1^+}`$ have been chosen). This results in 6 parameters to be determined from the fit to the angular distributions.
The PWA has been performed independently in 80 MeV intervals of the $`\eta \eta `$ mass spectrum. In each mass an event-by-event maximum likelihood method has been used. The function
$$F=\underset{i=1}{\overset{N}{}}ln\{I(\mathrm{\Omega })\}+\underset{L,M}{}t_{LM}ϵ_{LM}$$
(2)
has been minimised, where $`N`$ is the number of events in a given mass bin, $`ϵ_{LM}`$ are the efficiency corrections calculated in the centre of the bin and $`t_{LM}`$ are the moments of the angular distribution. The moments calculated from the partial amplitudes are shown superimposed as histograms on the experimental moments in fig 3. As can be seen the results of the fit reproduce well the experimental moments.
The system of equations which express the moments via the partial wave amplitudes is non-linear which leads to inherent ambiguities. For a system with S and D waves there are two solutions for each mass bin. In each mass bin one of these solutions is found from the fit to the experimental angular distributions, the other one can then be calculated by the method described in ref. . In the case under study the bootstrapping procedure is trivial because the Barrelet function has only two roots ($`Z_i`$ with $`i`$ = 1,2) and their real and imaginary parts do not cross zero as functions of mass, as seen in fig 2c) and d). In order to link the solutions in adjacent mass bins, the real parts of the roots are sorted in each mass bin in such a way that the real part of the first root should be larger than the real part of the second root (real parts of the two roots have different signs). For the first solution the imaginary parts of both roots are taken positive, the second solution is obtained by complex conjugation of one of the roots.
In this case two different PWA solutions are found. One solution is dominated by the S-wave the other solution has the events split between the different D-waves. By definition both solutions give identical moments and identical values of the likelihood. A Monte Carlo study has shown that if the input distribution is really dominated by a D-wave (i.e. $`D_0^{}`$, $`D_1^{}`$ or $`D_1^+`$) then both solutions of the PWA will be dominated by that D-wave. However, if the input distribution is dominated by the S-wave then two different solutions will result. One of the PWA solutions will be dominated by the S-wave, the other solution will be split equally between the D-waves and hence the physical solution can still be determined. The physical solution is shown in fig. 4. The $`S_0^{}`$-wave for the physical solution is characterised by a broad enhancement at threshold and a peak at 1.5 GeV. A broad enhancement is also seen in the $`D_0^{}`$-wave at 2.1 GeV. Superimposed on the waves as a histogram is the result of running the PWA on the background events.
Fig. 5 shows the background subtracted $`S_0^{}`$ and $`D_0^{}`$ waves. The PWA analysis has also been performed by extending the likelihood function given in equation (2) to include the background subtraction, namely
$$F=\underset{i=1}{\overset{N}{}}ln\{I(\mathrm{\Omega })\}+\underset{L,M}{}t_{LM}ϵ_{LM}+\underset{i=1}{\overset{N_{bg}}{}}ln\{I(\mathrm{\Omega })\}$$
where $`N_{bg}`$ is the number of background events. Since the background is dominantly S-wave the results of this method are similar to those shown in fig. 5.
The $`S_0^{}`$-wave has been fitted using a K-matrix parameterisation similar to that used to fit the $`K^+K^{}`$ spectrum with the addition of an incoherent background term. Poles have been introduced to describe the $`f_0(1370)`$, $`f_0(1500)`$ and $`f_0(1710)`$ with parameters fixed to those found from the coupled channel fit to the $`\pi ^+\pi ^{}`$ and $`K^+K^{}`$ spectra , namely M($`f_0(1370)`$) = 1312 MeV, $`\mathrm{\Gamma }(f_0(1370))`$ = 218 MeV, M($`f_0(1500)`$) = 1502 MeV, $`\mathrm{\Gamma }(f_0(1500))`$ = 98 MeV and M($`f_0(1710)`$) = 1727 MeV, $`\mathrm{\Gamma }(f_0(1710))`$ = 126 MeV. The fit describes well the data but, due to the size of the errors, it is not possible to conclude about the need for a further resonance above 2 GeV. The inclusion of the $`f_0(1710)`$ is essential to describe the spectrum.
The $`D_0^{}`$ wave has been fitted using two spin 2 relativistic Breit-Wigners and a linear background. The first Breit-Wigner is used to describe the $`f_2(1270)`$ with mass and width fixed to the PDG values and the second to describe the peak at 2.1 GeV where we have used M = 2130 MeV, $`\mathrm{\Gamma }`$ = 270 MeV which are the parameters found for $`f_2(2150)`$ observed in the $`K^+K^{}`$ final state . As can be seen the fit well describes the data. This state is compatible with the $`f_2(2175)`$ previously observed by the NA12/2 experiment in the $`\eta \eta `$ final state .
The error bars introduced by the partial wave analysis do not allow the parameters of the resonances to be determined from a free fit to the waves. Instead, we have performed a fit to the total mass spectrum shown in fig. 2b) using an incoherent sum of the expressions used to fit the $`S_0^{}`$ and $`D_0^{}`$ waves. The parameters for the $`f_0(1370)`$ and $`f_2(1270)`$ have been fixed to the values used above. The fit is shown superimposed on fig. 2b) and for the scalar resonances yields sheet II T-Matrix poles at
| $`f_0(1500)`$ | M | = | (1510 | $`\pm `$ | $`\mathrm{\hspace{0.33em}\hspace{0.33em}8}`$) | $`i`$ | ($`\mathrm{\hspace{0.33em}\hspace{0.33em}55}`$ | $`\pm `$ | $`\mathrm{\hspace{0.33em}\hspace{0.33em}8}`$) | MeV |
| --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- |
| $`f_0(1710)`$ | M | = | (1698 | $`\pm `$ | 18) | $`i`$ | ($`\mathrm{\hspace{0.33em}\hspace{0.33em}60}`$ | $`\pm `$ | $`13`$) | MeV |
These parameters are consistent with the PDG values for these resonances and our previous fits . For the $`f_2(2150)`$ the fit gives
| $`f_2(2150)`$ | M | = | 2151 | $`\pm `$ | 16 | MeV, | $`\mathrm{\Gamma }`$ | = | 280 | $`\pm `$ | 70 | MeV. |
| --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- |
The parameters found are compatible with those from the $`K^+K^{}`$ channel .
After correcting for the unseen decay modes and using data from the previously observed $`\pi \pi `$ final state the branching ratio $`\eta \eta `$ /$`\pi \pi `$ has been calculated for the $`f_0(1500)`$ from the fit to the $`S_0^{}`$ wave and gives :
$$\frac{f_0(1500)\eta \eta }{f_0(1500)\pi \pi }=0.18\pm 0.03$$
This value agrees well with the value that can be derived from the PDG of 0.23 $`\pm `$ 0.10.
For the $`f_2(1270)`$ the ratio is
$$\frac{f_2(1270)\eta \eta }{f_2(1270)\pi \pi }=(3\pm 1)\times 10^3$$
which is compatible with the PDG value.
For the $`f_0(1500)`$, $`f_0(1710)`$ and $`f_2(2150)`$ the branching ratio $`\eta \eta `$ /$`K\overline{K}`$ has been calculated from the WA102 data to be:
$$\frac{f_0(1500)\eta \eta }{f_0(1500)K\overline{K}}=0.54\pm 0.12$$
$$\frac{f_0(1710)\eta \eta }{f_0(1710)K\overline{K}}=0.48\pm 0.15$$
$$\frac{f_2(2150)\eta \eta }{f_2(2150)K\overline{K}}=0.78\pm 0.14$$
The branching ratios for the $`f_0(1500)`$, $`f_0(1710)`$ and $`f_2(2150)`$ differ from that of a known $`s\overline{s}`$ state, the $`f_2^{}(1525)`$, which has a ratio of 0.12 $`\pm `$ 0.04 which is consistent with the SU(3) prediction for such a state . No signal is seen for the $`f_2^{}(1525)`$ in the $`\eta \eta `$ final state of WA102 and an upper limit for its decay to this channel has been calculated, which gives
$$\frac{f_2^{}(1525)\eta \eta }{f_2^{}(1525)K\overline{K}}<0.14(90\%CL)$$
which is compatible with the PDG value given above.
For the $`f_0(1370)`$ there is considerable uncertainty due to the background subtraction which mainly affects the $`S_0^{}`$ wave at threshold. The strongest $`f_0(1370)`$ signal has been observed in the $`4\pi `$ final state therefore the branching ratio $`\eta \eta `$ /$`4\pi `$ has been calculated to be
$$\frac{f_0(1370)\eta \eta }{f_0(1370)4\pi }=(4.7\pm 2)\times 10^3$$
which is compatible with the Crystal Barrel measurement of (1.7 $`\pm `$ 0.9)$`\times 10^3`$ .
In the D-waves there is no evidence for the $`f_2(1950)`$ which has been claimed to have been seen in the $`\eta \eta `$ final states formed in $`p\overline{p}`$ annihilations . Since the $`f_2(1950)`$ is clearly seen in the $`4\pi `$ final state of experiment WA102 , an upper limit for its decay to $`\eta \eta `$ has been calculated and gives
$$\frac{f_2(1950)\eta \eta }{f_2(1950)4\pi }<5.0\times 10^3(90\%CL)$$
Hence this would imply that if the observation in the $`\eta \eta `$ final state of $`p\overline{p}`$ annihilations is correct, then a very large signal for the $`f_2(1950)`$ should be seen in the $`4\pi `$ final state of the same experiment.
The $`f_0(1500)`$ has previously been observed in the $`\pi \pi `$ , $`K\overline{K}`$ , $`4\pi `$ and $`\eta \eta ^{}`$ final states of the WA102 experiment. The relative decay rates for the $`f_0(1500)`$ are calculated to be:
$$\pi \pi :K\overline{K}:\eta \eta :\eta \eta ^{}:4\pi =1:\mathrm{\hspace{0.33em}0.33}\pm 0.07:\mathrm{\hspace{0.33em}0.18}\pm 0.03:\mathrm{\hspace{0.33em}0.096}\pm 0.026:\mathrm{\hspace{0.33em}1.36}\pm 0.15$$
The $`f_0(1370)`$ is below $`\eta \eta ^{}`$ threshold. The remaining relative decay rates for the $`f_0(1370)`$ are:
$$\pi \pi :K\overline{K}:\eta \eta :4\pi =1:0.46\pm 0.19:0.16\pm 0.07:34.0_9^{+22}$$
No signal was observed for the $`f_0(1710)`$ in the $`4\pi `$ or the $`\eta \eta ^{}`$ channels . Therefore, an upper limit has been calculated for its decay to these final states. For the $`f_0(1710)`$ the relative decay rates are:
$$\pi \pi :K\overline{K}:\eta \eta :\eta \eta ^{}:4\pi =1:5.0\pm 0.7:2.4\pm 0.6:<\mathrm{\hspace{0.33em}0.18}(90\%CL):<\mathrm{\hspace{0.33em}5.4}(90\%CL)$$
In summary, a partial wave analysis of the $`\eta \eta `$ channel has been performed for the first time in central production. Clear evidence is found for an $`\eta \eta `$ decay mode of the $`f_0(1500)`$, $`f_0(1710)`$ and $`f_2(2150)`$ and the decay branching ratios of these states have been determined.
Acknowledgements
This work is supported, in part, by grants from the British Particle Physics and Astronomy Research Council, the British Royal Society, the Ministry of Education, Science, Sports and Culture of Japan (grants no. 07044098 and 1004100), the French Programme International de Cooperation Scientifique (grant no. 576) and the Russian Foundation for Basic Research (grants 96-15-96633 and 98-02-22032).
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# Possibly interacting Vorontsov-Vel’yaminov galaxies
## 1 Introduction
Among the definitely interacting galaxies in close pairs or groups, many of which were discovered by Vorontsov-Vel’yaminov and included in his Atlas and Catalog of interacting galaxies (Vorontsov-Vel’yaminov, 1959, 1977) there are objects having VV numbers, in which the interaction is not so evident. Such objects were usually classified by Vorontsov-Vel’yaminov as the ”nests” or often similarly looking systems called ”chains” or ”pairs in contact”. He considered these objects as compact fragmenting systems, giving birth to young galaxies, but very soon observations showed that their nature may be different. Such systems can present both galaxies of strange and unusual shapes (either single objects or mergers), and really multiple systems, where it is hard to guess the number of components without detailed studies. In many cases only spectral measurements of their gas velocity distribution enable to prove or to disprove their solitude.
Earlier spectral observations confirmed that some VV systems actually present single dwarf galaxies with clumpy inner structure, resembling dwarf irregulars with multiple regions of active star formation (SF), or blue compact galaxies (BCGs) (Vorontsov-Vel’yaminov 1979a,b; Afanasiev et al. 1980; Arkhipova et al. 1981, 1987a,b,c). Only a few single nearby VV galaxies have been studied in detail (e.g. VV 556 $``$ GR 8, VV 499 $``$ DDO 053).
The important questions on the chemical abundances in the studied VV galaxies left outside the scope of these early observations. New opportunities appeared due to new CCD-detectors, which enable to realize high sensitivity and large dynamical range. It induced authors to return back to the study of ”nest-like” VV galaxies. Besides, many new observational data have appeared on many of these objects.
For the current study we selected more than twenty VV galaxies looking like multiple systems or singular irregular systems in the POSS/DSS <sup>1</sup><sup>1</sup>1DSS is Digital Sky Survey distributed by Hubble Space Telescope Science Institute images, for which such indicators of interaction as well defined tidal tails or bridges are absent. The main objectives of our investigation are:
– to clarify the distances, total luminosities and masses of those systems, systemic velocities of which are not known or badly known;
– to carry out high S/N ratio spectrophotometry in order to address problems of chemical abundances and evolutional status of these objects,
– to analyse the inner gas kinematics and structural properties of the objects.
Since many of VV galaxies in question are dwarfs with recent or current SF burst, it is also important to check possible companions which could exert strong enough tidal action.
In this paper, the first in the series, we present the results of recent long-slit spectroscopy for three VV objects: VV 432, VV 543 and VV 747 and the observation in HI-line of VV 747.
In section 2 we describe observations, data reduction, abundances determination and the measurements of velocity distribution of ionized gas along the slit.
Observations of VV 747 in the HI-line 21 cm and their results are presented in section 3. In section 4 we consider the individual properties of studied galaxies. Discussions and preliminary conclusions are presented in section 5. We adopt throughout the paper H$`{}_{0}{}^{}=`$75 km sec<sup>-1</sup>.
## 2 Spectral observations, data reduction and analysis
### 2.1 Observations
The spectroscopic data were obtained with the 6 m telescope of the Special Astrophysical Observatory of Russian Academy of Science (SAO RAS) during two runs in February and April 1999. The Long-Slit spectrograph (LSS in Table 1) (Afanasiev et al. Afanasiev95 (1995)) at the telescope prime focus was equipped with a Photometrics CCD-detector PM1024 (with $`24\times 24\mu `$m pixel size) (PMCCD in Table 1) installed at Schmidt-Cassegrain camera F/1.5. Most of the long-slit spectra ($`1.2\mathrm{}\times 180\mathrm{}`$) were obtained with the grating of 325 grooves/mm, giving a dispersion of 4.6 Å/pixel. Additional data were obtained with the grating of 1302 grooves/mm and dispersion 1.2 Å/pixel. For the latter set-up the slit of $`2\mathrm{}\times 180\mathrm{}`$ was used. The scale along the slit was 0.39$`\mathrm{}`$/pixel.
The resulting resolution (FWHM) was about $`1415`$ Å for the first set-up, and about 3.7 Å for the second set-up. Reference spectra of an Ar–Ne–He lamp were recorded before or after each observation to provide wavelength calibration. Spectrophotometric standard stars from Massey et al. (Massey88 (1988)) were observed for flux calibration at least twice a night.
Observations and data processing in this set-up have been conducted mainly under the software package NICE in MIDAS, described by Kniazev & Shergin (Kniazev95 (1995)).
### 2.2 Data reduction and abundances determination
The data reduction was performed in SAO RAS, using various packages of MIDAS (see Kniazev et al. (Kniazev2000 (2000)) for details).
In Table 2 we summarize the main observational parameters of the discussed three VV “nests”. They include the names of the objects, their coordinates for the epoch J2000, the apparent blue magnitudes and the corresponding references, the radial heliocentric velocities, measured in this work with their r.m.s. uncertainties, maximal angular sizes, absolute blue magnitudes and the oxygen abundances (12+log(O/H)).
Direct images of studied galaxies, extracted from the DSS and the position of long slit, indicated by bar are presented in Fig. 1a, 2a, 3a. Corresponding 2-D spectra are shown in Fig. 1c, 2c, 3c. The brightness profiles of H$`\alpha `$ line along the slit, and corresponding velocity curves are illustrated in Fig. 1b, 2b, 3b. In Fig. 1d, 2d, and 3d we present 1-D spectra, extracted from 2-D spectra, which were used for the measurements of line intensities, determination of physical conditions and abundances of Hii-regions.
The resulting observed emission line intensities $`F(\lambda )`$ of various ions relative to H$`\beta `$, both uncorrected and corrected for interstellar extinction and underlying stellar absorption $`I(\lambda )`$ (following the procedure described by Izotov et al. Izotov97 (1997)) for the brightest parts of the galaxies are presented in Tables 3 along with the extinction coefficient C(H$`\beta `$), the equivalent width of absorption Balmer hydrogen lines EW(abs), the equivalent width of H$`\beta `$ line EW(H$`\beta `$) and the observed H$`\beta `$ flux.
For the abundances determination we used the scheme, described in detail by Izotov et al. (Izotov94 (1994), Izotov97 (1997)). The electron temperatures and densities in Hii-regions of the observed VV-galaxies and their abundances of O, N and Ne are summarized in Table 4.
For VV 543 the \[Oiii\]-line 4363 Å is not detected. Therefore to estimate its metallicity we employ the empirical method (see e.g. Pagel et al. (Pagel79 (1979)), McGaugh (McGaugh91 (1991)) and Olofsson (Olofsson97 (1997))). Its uncertainty for 12+log(O/H) can be as large as 0.2–0.3 dex. We also applied this empirical method to all four individual knots of VV 432.
The ratio \[Oiii\]/\[Nii\] of extinction corrected intensities of the lines $`\lambda `$ 4959,5007 Å and $`\lambda `$ 6548,6584 Å enable one to avoid well known ambiguity of the empirical method (Alloin et al. Alloin79 (1979)). Taking into account these ratios, one can use $`R_{23}`$ and $`Q`$ parameters to get more reliable estimate of $`O/H`$ through the model curves by Olofsson (Olofsson97 (1997)). By this way it was obtained for VV 543W the abundance 12+log(O/H) = 8.5$`\pm `$0.1, which is about 0.3 lower than the value, derived from McGaugh (McGaugh91 (1991)) tracks, and seems to be in a better agreement with the abundance, expected for underluminous Hii-galaxies.
### 2.3 Velocity curves
To obtain line-of-sight velocity distribution along the slit we use the MIDAS programs, kindly presented to authors by D. Makarov.
To increase the accuracy of derived rotation curves the programs were modified to include additional corrections using close lines in the reference spectrum.
The procedure includes the following steps: 1) a linearisation of 2-D spectrum of the object and of the reference spectrum; 2) a measurement of the position of the H$`\alpha `$ emission for each of the 250 position rows using gaussian fitting; 3) an estimation of the background and S/N ratio for the H$`\alpha `$ emission in each row; 4) a similar measurement of the nearby line Nei $`\lambda `$ 6598.95 Å in the reference spectrum and compiling the table of differences between the laboratory wavelength of this line and the measured one for each row; 5) a polynomial fitting of these differences and a determination of the residual scattering (r.m.s.); 6) the application of this fitting polynomial to correct the measured wavelengths of H$`\alpha `$. The resulting error of wavelengths measurements is determined by quadratic summing of the fitting r.m.s., found in the previous step, and the measurement error of velocity for each point along the slit, estimated from gaussian fitting. The latter varies between 2.5 km s<sup>-1</sup> for the points with the highest signal and 9 km s<sup>-1</sup> for the points with low S/N ratio for dispersion 1.2 Å/pixel; 7) rebinning of H$`\alpha `$ data along the slit, corresponding to the seeing during the observations (4 or 5 pixels in April and February, respectively).
Below we use only those velocity estimates which satisfy the criteria S/N $`>`$ 3 and $`\sigma _V`$ $`<`$ 15 km s<sup>-1</sup>. Resulting high accuracy of corrected observed wavelength enables one to study irregularities of the velocity curve with an amplitude as low as 10 km s<sup>-1</sup> on the angular scale up to 20″.
For VV 432, the rotation curves which were obtained with the low and high resolution spectra are presented in Fig. 1b. Their similarity shows that our low dispersion spectra can be used to derive preliminary dynamical parameters of the studied galaxies.
## 3 Hi data for VV 747
For VV 747, the interpretation of the optical data remains ambiguous regarding its multiplicity. (see Sect. 4.3). In order to settle this ambiguity we investigated the dynamical state of the system and performed single-dish Hi observations.
### 3.1 Observations and data reduction
The 21 cm Hi-line observations were carried out in July 1999 with the Nançay<sup>2</sup><sup>2</sup>2The Nançay Radioastronomy Station is part of the Paris Observatory and is operated by the Ministère de l’Education Nationale and Institut des Sciences de l’Univers of the Centre National de la Recherche Scientifique. 300 m telescope (NRT). The NRT has a half-power beam width of 3.7 (EW) $`\times `$ 22 (NS) at the declination $`\delta =0^{}`$.
A dual polarization receiver has been used, with a system temperature of $``$ 40 K in the horizontal and vertical linear polarizations. The gain of the telescope was 1.1 K Jy<sup>-1</sup> at a declination of $`\delta =0^{}`$. The observations were made in position switching mode with 1-minute on-source and 1-minute off-source integrations.
Since VV 747 had a known optical redshift, we split the autocorrelator into two, each bank covering a bandwidth of 6.4 MHz, and centred at the frequency corresponding to the optical redshift. The velocity coverage was 1 350 km s<sup>-1</sup>. The channel width was 2.6 km s<sup>-1</sup> and after averaging pairs of adjacent channels the effective resolution was 6.3 km s<sup>-1</sup>.
The data has been reduced using the software developed by the telescope’s staff. Both polarization spectra were calibrated and processed independently, and were averaged to improve the signal-to-noise ratio. Errors were calculated following Schneider et al. (Schneider86 (1986)).
### 3.2 Results. Hi-profile and its parameters
For the total integration time of 234 minutes (ON+OFF source) the r.m.s. noise is 4 mJy. The galaxy is detected with S/N ratio of 12. The spectrum is presented in Fig. 4. The line shape is close to gaussian and typical of low-mass galaxies. Its width at 50 % of peak W<sub>0.5</sub> = 90 km s<sup>-1</sup> is close to the median value of the distribution derived for the statistical sample of BCGs from the zone of the Second Byurakan Survey (SBS) by Thuan et al. (Thuan99 (1999)). The line width at 20 % of peak W$`{}_{0.2}{}^{}=117`$ km s<sup>-1</sup> is very close to that of VV 432. This is important in the following discussion about the applicability of the Tully-Fisher relation (Tully77 (1977), hereafter TF) for the distance problem of the latter galaxy.
The Hi-mass of VV 747 of 0.8$`\times `$10<sup>8</sup> M, is typical of low luminosity gas-rich galaxies, like the ratio of the hydrogen mass over the blue luminosity: M(Hi)/L$`{}_{B}{}^{}=1.0`$.
## 4 Properties of individual galaxies
### 4.1 VV 432 = IC 3105 = UGC 7326
The first long-slit spectrum of this chain-like galaxy has been obtained along the major axis by Arkhipova et al. (Arkhipova81 (1981)). They revealed high-excitation spectrum of Hii-region in the NE region. The amplitude of the rotation curve has been found to be relatively small (less than about 60 km s<sup>-1</sup>).
Our long-slit spectrophotometry allowed to detect the NE Hii-region (knot “a” on Fig. 1) \[Oiii\]-line $`\lambda `$ 4363 Å with a good S/N ratio in both spectra. The first spectrum, obtained on February 12, 1999 covered the range 3600–8000 Å and the spectra taken on April 24, 1999 with a dispersion 1.2 Å/pixel, covered two separate ranges 4000–5200 Å and 6000–7200 Å. Table 3 presents the data for all measured emission lines. According to this table, the relative intensities of strong oxygen and hydrogen lines in both spectra are the same within 3%, in accordance with their internal errors. For the abundance derivation we supposed that, for the April spectra, the relative intensity of the \[Oii\]-line $`\lambda `$ 3727 Å and of the H$`\beta `$ line to be equal to those derived from the February spectrum.
Oxygen abundance determination gives consistent values for both data sets, so we accept the weighted mean of the two independent measurements. The resulting value of 12+log(O/H) = 7.58$`\pm `$0.06 obtained for this Hii-region appear to be very low.
Since the \[Oiii\]-line $`\lambda `$ 4363 Å was not detected in other emission-line knots, the oxygen abundances for knots “b”, “c” and “d” were estimated by the same empirical method mentioned in Sect. 2.3. The most uncertain parameter – the mean gas density – was assumed to be close to that found for knot “a”, what is about 1 atom cm<sup>-3</sup> (the value, derived from the consistency of oxygen abundances, determined by classical and empirical methods). Oxygen abundances, found by the empirical method for this value of N<sub>e</sub> for three knots, are consistent with the estimates, obtained for knot “a”, taking into account the uncertainties of parameter R<sub>23</sub> and possible variations of the mean density N<sub>e</sub> within the factor of 2.
This galaxy is elongated (the axial ratio of $`b/a`$ $``$ 1/3 according to LEDA)<sup>3</sup><sup>3</sup>3LEDA is the acronym of Lyon-Meudon Extragalactic Database, http://leda.univ-lyon1.fr, which is not untypical for dwarf galaxies. It looks like an edge-on disk, bent on both NE and SW edges, which probably indicates a tidal action from some other galaxy. The bright Hii-region on the NE edge appears to be outside of the main body of the galaxy, and its connection to the disk is not evident, although a smooth velocity distribution along the galaxy favours its interpretation as a single object. Another possibility which cannot be excluded is that the outer knots are large SF regions at the ends of spiral arms, seen not exactly edge-on. However due to the rather small luminosity of VV 432, the presence of spiral arms is quite improbable.
Our long-slit spectrum with the dispersion of 1.2 Å/pixel allows to derive the rotation curve along the “disk”, using the H$`\alpha `$ line data (see Fig. 1b). A velocity gradient is well seen over the whole disk; there is no sign of flattening of the velocity curve on both NE and SW edges. The mean heliocentric velocity of the galaxy is $`160`$ km s<sup>-1</sup>. The amplitude of the rotation curve within the whole extent of H$`\alpha `$-emission of about 70<sup>′′</sup> corresponds to the maximal rotational velocity of 40$`\pm `$5 km s<sup>-1</sup>. The expected inclination correction is very small, since the disk is seen at the inclination angle $`i>70^{}`$.
The uncertainties of the data points of the Fig. 1b’s velocity curve are small and remain in the range 4 to 15 km s<sup>-1</sup>. The irregularities which are seen on the rotation curve can then be considered as real ones, probably connected with the regions of active SF.
The galaxy is situated in the direction of Virgo cluster (VC), and is also catalogued as VCC 0241. Its negative radial velocity does not contradict its membership to the cluster (Binggeli et al. BPT93 (1993)). In this case, its distance is of $``$ 20.7 Mpc (distance modulus = 31.6 mag, Federspiel et al. (Federspiel98 (1998))). Its maximal diameter on the isophote $`\mu _B`$= 25$`{}_{}{}^{m}/\mathrm{}\mathrm{}`$ is 111$`\mathrm{}`$, according to LEDA, which corresponds to a linear size of about 11 kpc. The extent of the H$`\alpha `$ emission on our long-slit spectrum is of 70$`\mathrm{}`$, which gives the total extent of Hii-emission to be $``$ 6.8 kpc.
If VV 432 is a member of VC it appears to be the most metal-deficient known galaxy of this cluster (see for comparison the metallicity data of the VC BCGs in Izotov & Guseva (Izotov89 (1989))). Such low metallicity implies that the galaxy presumably experienced only one or two major SF episodes during its life. Since galaxy interactions in clusters are very important triggers, it is difficult to understand how VV 432 could sustain its very low metallicity in the dense VC environment. The existence of such metal-poor gas-rich galaxy poses serious questions on its evolution history, and thus gives additional arguments to make independent check of its distance, using, in particular, color-magnitude diagrams for resolved stars.
To estimate physical parameters of VV 432 we accept further that this galaxy is in VC, at the distance 20.7 Mpc.
Since as we already noticed VV 432 resembles by its morphology an edge-on bent disk, it is natural to check probable companion galaxies, acting as strong enough disturbing bodies.
According to NED <sup>4</sup><sup>4</sup>4The NASA/IPAC Extragalactic Database (NED) is operated by the Jet Propulsion Laboratory, California Institute of Technology, under contract with the National Aeronautics and Space Administration., the nearest galaxies in both the projected distance and the relative velocity are VCC 0200 (MCG +02–31–076) at the angular distance 41.4 and V<sub>hel</sub> = 65 km s<sup>-1</sup>, with B-magnitude = $`14\stackrel{m}{.}69`$, and NGC 4216 (UGC 7284) at the angular distance 51.8 and V<sub>hel</sub> = 131 km s<sup>-1</sup>, with B-magnitude = $`10\stackrel{m}{.}99`$. If they are members of VC, their respective projected distances relative to VV 432 are 240 and 312 kpc. NGC 4216 according to LEDA is a giant spiral with M<sub>B</sub>= $`21\stackrel{m}{.}6`$. If it is indeed on the same radial distance as VV 432, it really can exert strong enough tidal action on VV 432 to trigger gravitational instability and subsequent SF bursts. In particular, tidal generation of shocks in gaseous disk according to Icke (Icke85 (1985)) mechanism can be responsible for observed SF bursts (see e.g. Pustilnik et al. (Pustilnik00 (2000)) for the estimates of tidal effect in similar situation).
### 4.2 VV 543 = NGC 5275
This galaxy was considered as a candidate dwarf galaxy (M<sub>B</sub>= $`15\stackrel{m}{.}5`$) due to its catalog radial velocity V = 1 395 km s<sup>-1</sup> (RC3, de Vaucouleurs et al. (deVaucouleurs91 (1991)), cited also in both LEDA and NED, and originally obtained by Arkhipova & Esipov (Arkhipova79 (1979))). However colour indices of this object are rather typical of giant E galaxies than for dIrr systems (Zasov & Arkhipova 2000). According to our observations, it is evident that there was some misprint in the original work, caused the catalog velocity of this galaxy to be in error by the factor of ten. Its real velocity, measured from our spectrum of emission-line region on the western periphery (hereafter VV 543W) is 14 100 km s<sup>-1</sup>.
In fact, the spectrum of the central bright region of this galaxy shows, that we have in this case an optical pair in projection. Indeed the central part shows only the absorption lines, typical of elliptical galaxy (e.g. Pickles (Pickles88 (1988))), but its radial velocity is 1620$`\pm `$120 km s<sup>-1</sup> lower than that for the emission-line galaxy. Its luminosity M<sub>B</sub> = $`20\stackrel{m}{.}9`$ evidences that this galaxy falls to the class of normal ellipticals. The apparent magnitude of VV 543W is $``$ 2$`\stackrel{m}{.}`$5 fainter (B $`17\stackrel{m}{.}7`$) than that of VV 543E (this estimate follows from the comparison of the flux near $`\lambda `$ 4400 Å in both galaxies), what leads to M$`{}_{B}{}^{}`$ $`18\stackrel{m}{.}7`$. Its linear size along the major axis $``$ 13 kpc is quite modest, and, owing to the typical spectrum of Hii-region, this western component can be considered as a bright blue compact/Hii galaxy. The compact object lying at about 20$`\mathrm{}`$ to NE from the absorption-line galaxy is a foreground star.
The extent of H$`\alpha `$ emission along the minor axis is traced within 8$`\mathrm{}`$. Even after the binning in 4 pixels (which corresponds to the seeing of 2$`\mathrm{}`$) the velocity curve does not indicate clear gradient. The full range of the radial velocity is about 80 km s<sup>-1</sup> with the mean value of 14 100$`\pm `$30 km s<sup>-1</sup>.
According to NED, VV 543W has a probable companion galaxy NGP9 F324–0303806 at 5.8 with the radial velocity 13 699$`\pm `$126 km s<sup>-1</sup> and B = $`16\stackrel{m}{.}68`$. Corresponding projection distance $``$ 300 kpc and velocity difference 401$`\pm `$130 km s<sup>-1</sup> are in the range typical of wide pairs of binary galaxies (see e.g. the study by Chengalur et al. (Chengalur93 (1993)) and Nordgren et al. (Nordgren98 (1998))). Some weak tidal action from this galaxy can be responsible for the enhanced SF in VV 543W (see e.g. Reshetnikov & Combes (Reshet97 (1997)), and Rudnick & Rix (Rudnick98 (1998))).
### 4.3 VV 747 = CG 798
This dwarf galaxy consists of two clearly separated regions embedded into the common envelope of low surface brightness. The brighter SW component has full size of about 15$`\mathrm{}`$ and the fainter NE can be traced down to about 8$`\mathrm{}`$ ($``$ 0.8 and 0.4 kpc respectively). Both regions show emission-line spectrum (see Fig. 3b). The high excitation spectrum with the observed \[Oiii\] line $`\lambda `$ 4363 Å of the SW component allows to determine oxygen abundance using the direct measurement of $`T_e`$. Our O abundance is 12+log(O/H) = 7.85$`\pm `$0.05 compared to the abundance 7.97 derived by Izotov & Thuan (Izotov99 (1999)). This low metallicity is not untypical for BCGs.
Smooth velocity curve along the slit shows a small but clear slope across the SW component in the region of high S/N ratio of H$`\alpha `$ with the total extent of about 15$`\mathrm{}`$ and the velocity range from 540 to 630 km s<sup>-1</sup>, consistent with the maximal rotation velocity of about 45 km s<sup>-1</sup> and mean radial velocity of 585 km s<sup>-1</sup>. The galaxy body is well elongated, so the inclination correction does not seem to be larger than 20$`÷`$30%. In the NE component the velocity can be measured only in two independent points, showing significant scattering, consistent with their internal uncertainties. The mean velocity of NE component, taking from these two points, is about 690$`\pm `$42 km s<sup>-1</sup>.
Huchra et al. (Huchra95 (1995)) presented the values of the observed velocities for two components of this galaxy, designated as SW and NE, but their coordinates given in the paper, are the same for both of them. If we accept that they observed the same components, our results are consistent with theirs (V(SW) = 621$`\pm `$32 and V(NE) = 665$`\pm `$44 km s<sup>-1</sup>) within the cited uncertainties.
Since the two separate velocity points for the NE component do not show any clear gradient, which would indicate its independent rotation, and the mean velocity of the NE component matches well the continuation of the velocity curve for the SW component, there is no reason to consider this system as two different galaxies in collision. The current data favour the interpretation of this system as a single galaxy with two super-giant Hii-regions in different excitation stages. In this case the systemic velocity of this single galaxy is about 620 km s<sup>-1</sup> with the full velocity range of 540 to 710 km s<sup>-1</sup>. Both these values are quite well consistent with the parameters of Hi-profile of VV 747 described and shown in section 3.2. The latter full velocity range corresponds to the maximal rotational velocity of 85 km s<sup>-1</sup>.
The better knowledge of the velocity curve in the region of the NE component is necessary to exclude completely the hypothesis of merger of two dwarf galaxies. At least one indirect argument favours this interpretation. If TF relation holds for this galaxy, its Hi mass and blue luminosity are more consistent with the narrower width (W<sub>0.2</sub> = 70–80 km s<sup>-1</sup>), than with the observed one in the integrated Hi-profile. It can indicate that we observe the sum of two “narrow” Hi-profiles displaced at about 40–50 km s<sup>-1</sup> each from other. Only 21 cm line mapping of Hi velocity field and/or very high quality H$`\alpha `$ velocity data can help to make a final interpretation of this system.
The search in the NED resulted in the bright companion galaxy NGC 3432 (= Arp 206 = VV 011) at the angular distance 67 with B = $`11\stackrel{m}{.}67`$ and the radial velocity 616 km s<sup>-1</sup>, very close to that of VV 747. This galaxy, classified as SB(s)m (LINER Hii), has M<sub>B</sub> = $`18\stackrel{m}{.}3`$. At the projected distance $``$ 190 kpc the tidal action of this galaxy to VV 747 may also be strong enough to trigger SF in this tiny dwarf.
## 5 Discussion and conclusions
One of the important questions concerning the nature of the “nest” and “chains” of VV-galaxies is their evolution status. Current study as well as several previous publications demonstrate that many of the low-luminosity VV-galaxies are relatively nearby irregular galaxies with several bright knots of enhanced SF. Their position in two-colour diagram indicates a presence of SF burst in many of them (Zasov & Arkhipova Zasov00 (2000)). They roughly follow so called luminosity – metallicity relation: the least luminous galaxies show in general smaller values of $`O/H`$.
VV 432 has a very low heavy-element abundance. Its $`O/H`$ is among the lowest ten values of the most metal-deficient BCGs out of more than one thousand BCGs/Hii-galaxies known up-to-now. It may be considered as an example of non-evolved galaxy. Analysis of the empirical correlations suggests that dwarf galaxies with 12+log(O/H) $`<`$ 7.6 can currently experience only the first in their history episode of SF (Izotov & Thuan Izotov99 (1999)). Therefore VV 432 with 12+log(O/H) = 7.58 is very good candidate for more detailed study. If it is situated in Virgo cluster, it will be the most metal deficient galaxy of this aggregate, being even less chemically evolved than another well-known metal-poor Hii-galaxy in the direction of Virgo cluster Hi 1225+01 with 12+log(O/H) = 7.66 (Salzer et al. Salzer91 (1991); Chengalur et al. Chengalur95 (1995)). One of the possible ways to resolve the dilemma of radial distance to VV 432 is a detection of brightest stars and construction of their color-magnitude diagram.
The galaxies we discuss have rather close neighbours. Icke (Icke85 (1985)) first has drawn attention to the importance of relatively weak interactions to trigger gravitational instability in gas disks via generation of shocks. Many observational evidences for the important role of weak interactions to trigger SF were obtained since that time, including the detection of low mass Hi-companions of nearby Hii-galaxies (Chengalur et al. Chengalur95 (1995); Taylor et al. Taylor93 (1993), Taylor95 (1995); Taylor Taylor97 (1997)) and optical faint companions of BCGs (Pustilnik et al. Pustilnik97 (1997)). Recent results on late spirals by Reshetnikov & Combes (Reshet97 (1997)) and Rudnick & Rix (Rudnick98 (1998)) also suggest the importance of weak interactions to modulate SF history in these galaxies.
From the observational data discussed above some preliminary conclusions can be drawn:
* Spectrophotometry of VV-galaxies shows that low luminosity representatives of this sample are in general metal-deficient objects, and in this aspect they are similar to dwarf irregular galaxies.
* The extremely metal-deficient galaxy VV 432 (12+log(O/H) = 7.58) is probably the least evolved known member of Virgo cluster.
* The system VV 543 radial velocity cited in the RC3 catalog and other databases is wrong. This object consists of two galaxies with discordant redshifts (a unique example among the galaxies of this type!), and probably presents an optical pair. VV 543W is an Hii-galaxy with the radial velocity 1620 km s<sup>-1</sup> higher than that of absorption-line E-type galaxy VV 543E.
* VV 747 is probably a single dwarf galaxy rich of neutral hydrogen.
* The presence of massive companion galaxies at the distances of few hundred kpc from the studied VV-objects with enhanced SF rate is probably indicative of the important role of weak interactions to trigger SF activity at least in some fraction of low mass VV galaxies.
###### Acknowledgements.
We thank D.Makarov (SAO) for providing us with his MIDAS package to derive velocity curve on long-slit spectra. SAO authors appreciate the partial financial support from the RFBR grant No. 96-02-16398 One of the authors (A.Z.) thanks RFBR (grant 98-02-17102) and Federal program “Astronomy” for financial support. We have made use of the NASA/IPAC Extragalactic Database (NED), which is operated by the Jet Propulsion Laboratory, California Institute of Technology, under contract with the National Aeronautics and Space Administration, and of the Lyon-Meudon Extragalactic Database, http://leda.univ-lyon1.fr.
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# Young Star Clusters: Metallicities, Ages, and Masses
## 1. Why do we want to know GC and YSC masses?
GCs are conventionally believed to be (among) the oldest objects in the Universe, dating back to the times of galaxy formation. They are used to constrain the age of the Universe. LFs of GC systems are believed to be universal enough for determinations of distances to $`>20`$ Mpc and of the Hubble constant. On the other hand, “present-day GCs are the hardiest survivors of a larger original population” (Harris 1991). Hence, their observed LF is not only their LF at formation shifted by stellar evolutionary fading, but might be additionally modified by cluster destruction processes. Dynamical modelling of cluster systems in the Galactic potential shows that destruction processes and timescales strongly depend on cluster masses. While destruction by dynamical friction is more efficient for high mass clusters, tidal shocking and evaporation preferentially destroy low mass clusters.
Bright YSCs are observed in large numbers in interacting galaxies and merger remnants and a burning question with far-reaching implications is if these YSCs are young GCs. Masses derived for these YSCs are much higher than those of open clusters in the Milky Way. Effective radii of YSCs typically are a few pc similar to GC radii. The mass function (MF) of YSC systems in comparison with the MFs of molecular clouds or molecular cloud cores tells us about star and cluster formation processes. The MF of YSCs, as first pointed out by Meurer (1995), may differ in shape from their LF since M/L varies rapidly at young ages and the age spread within a YSC system is not much smaller than its age. The LF of open clusters in the Milky Way, the MFs of molecular clouds and molecular cloud cores, the observed LFs of YSCs, e.g. in NGC 4038/39, NGC 7252, NGC 3256, and of giant HII regions all are power laws with slopes in the range $`\alpha 1.5...1.8`$ (cf. Solomon et al. 1987, Lada et al. 1991, Kennicutt 1989, and reviews by Harris & Pudritz 1994, Elmegreen & Efremov 1997). Yet the LF of old GCs is Gaussian with typical parameters $`\mathrm{M}_\mathrm{V}7.3\mathrm{mag},\sigma (\mathrm{M}_\mathrm{V})1.3\mathrm{mag}`$, their MF is log-normal with typically $`\mathrm{Log}(\mathrm{M}/\mathrm{M}_{})5.5,\sigma 0.5`$ (e.g. Ashman et al. 1995).
Hence the question as to the MF of YSCs has profound implications. If the MF of YSCs were a power law like their LF and if YSC systems are to evolve into something similar to old GC systems, dynamical destruction processes would have to transform the power law MF into a log-normal MF over a Hubble time. If, on the other hand, the MF of YSCs were log-normal (and their LF distorted to a power law by age spread effects), then the star/cluster formation process would have to transform the power law MF of the molecular clouds into the log-normal MF of YSCs. Or, else, might already the MF of molecular clouds/cloud cores in violently star forming mergers (where to my knowledge it has not yet been observed) be different from what it is in quietly star forming ‘normal galaxies’ (where it is observed)?
## 2. Evolutionary synthesis of SSPs
Star clusters are simple stellar populations (SSPs), formed in one short burst of star formation with one metallicity Z. With evolutionary synthesis models, the time evolution of SSPs of various metallicities $`\mathrm{Z}_\mathrm{i}`$ is studied in terms of luminosities $`\mathrm{L}_\lambda `$, colours, spectrum, absorption features, stellar mass loss, and hence M/L by many groups (e.g. Bruzual & Charlot 1993, Worthey 1994, Bressan et al. 1994, F.-v.A. & Burkert 1995, Leitherer et al. 1999, Kurth et al. 1999). Basic parameters of this approach are the IMF and the set of stellar evolutionary tracks used (e.g. from the Padova or Geneva groups). All models agree that the changes in luminosities and colours are rapid in early evolutionary stages and become slower with increasing age. The colour evolution depends significantly on metallicity, already in very early stages, and the fading also depends on metallicity, in particular during the first Gyr (cf. F.-v.A. & Burkert 1995, Kurth et al. 1999).
## 3. YSC observations vs. SSP models
### 3.1. Procedure
If the metallicity of YSCs is known, individual ages can be obtained on the basis of their observed UBVI colours. Unfortunately, metallicity information from spectroscopy is only available yet for a handful of YSCs in NGC 7252 (Schweizer & Seitzer 1993, 1998) and NGC 1275 (Brodie et al. 1998). In all other cases we have to go back to some educated guess. In gas rich galaxy mergers, YSCs form out of the gas from the progenitor (spiral) galaxies, the abundance of which gives a lower limit to YSC abundances. E.g., for the YSCs in NGC 7252, which is a merger of two luminous Sc type spirals, we had predicted $`\mathrm{Z}_{\mathrm{YSC}}>\mathrm{Z}_{\mathrm{ISM}}^{\mathrm{Sc}}{}_{}{}^{1}/_{2}^{}\mathrm{Z}_{}`$ from the ISM abundance evolution in our 1-zone spiral models (F.-v.A. & Gerhard 1994). Spectroscopy of the brightest YSCs yielded $`\mathrm{Z}_{\mathrm{YSC}}\mathrm{Z}_{}`$ with some tentative indication of self-enrichment during the burst from the comparison of Mgb and Fe lines (cf. F.-v.A. & Burkert 1995). Distances are known for all YSC systems. Hence absolute luminosities $`\mathrm{L}_\mathrm{V}`$ can be combined with $`\mathrm{M}/\mathrm{L}_\mathrm{V}`$ values from SSP models of appropriate age and metallicity to derive the masses of YSCs.
### 3.2. Sources of uncertainty
Sources of uncertainty for these mass estimates come both from observations and models. Observational errors on luminosities and colours might be inhomogeneous and probably not independent of each other, dust extinction is not known for individual YSCs, the dust distribution in some systems is inhomogeneous, the metallicity is, at most, known for a small subsample of YSCs and might show an intrinsic scatter, and, finally, the completeness limit need not be homogeneous over the region of YSC observations. Intrinsic model uncertainties are estimated to be $`<0.1`$ mag in (optical) luminosities and colours. Differences between models from various authors are mainly due to differences in the stellar input physics. While serious colour discrepancies at very young ages $`10`$ Myr are seen, e.g. comparing models from Bruzual & Charlot with those of Leitherer et al. , probably due to the inclusion/non-inclusion of emission lines, fairly good agreement is reached among models from various groups for the same metallicity and IMF at all ages $`>60100`$ Myr. E.g., $`\mathrm{\Delta }(\mathrm{V}\mathrm{I})_{12\mathrm{Gyr}}<0.1`$ mag and $`\mathrm{\Delta }(\mathrm{M}/\mathrm{L}_\mathrm{V})_{12\mathrm{Gyr}}<10\%`$.
$`\mathrm{M}/\mathrm{L}_\lambda `$, however, depends on wavelength $`\lambda `$, metallicity Z, and IMF, i.e. on its slope and lower mass limit. In Tab.1, I briefly sketch out these dependencies as obtained from our models using Padova stellar evolutionary tracks for stars in the mass range $`0.160\mathrm{M}_{}`$ (cf. Kurth et al. 1999) and two different IMFs (Salpeter vs. Scalo 1986).
In general and as reported by others before, model $`\mathrm{M}/\mathrm{L}_\mathrm{V}`$ are about twice as large as are $`\mathrm{M}/\mathrm{L}_\mathrm{V}`$ values derived from observations of old GCs, even when a Scalo IMF is assumed. For a Salpeter IMF the discrepancy is higher by another factor of two. Observational $`\mathrm{M}/\mathrm{L}`$ values are obtained by measuring the central velocity dispersion $`\sigma _0`$, central surface brightness $`\mathrm{I}_0`$, and half light radius $`\mathrm{r}_{1/2}`$, and assuming isotropic orbits, no radial gradients in $`\mathrm{M}/\mathrm{L}`$, and no DM halos around GCs. Then
$`\mathrm{M}/\mathrm{L}\frac{\sigma _0^2}{\mathrm{I}_0\mathrm{r}_{1/2}}`$
and typical values quoted for Galactic GCs are $`\mathrm{M}/\mathrm{L}_\mathrm{V}2`$. There are two effects that might invalidate the assumptions going into this derivation. First of all, mass segregation in GCs (cf. Meylan this conf.) will result in an $`\mathrm{M}/\mathrm{L}`$ increasing with radius, as e.g. observed by Côté et al. 1995 for NGC 3201. Second, low mass stars with their high $`\mathrm{M}/\mathrm{L}`$ values are preferentially lost by evaporation (e.g. Gerhard this conf.). While both processes are undoubtedly at work in GCs, attempts to examine quantitatively the validity of the assumptions involved seem to still give ambiguous results. Leonard et al. 1992 use proper motion data and radial velocity measurements for stars in the GC M13 and find that $`50\%`$ of the mass of M13 is in low mass stars and brown dwarfs. Including this unseen mass leads to an increased $`\mathrm{M}/\mathrm{L}4`$, a value well compatible with the models. On the other hand, observations of tidal tails on some GCs in the Milky Way potential are interpreted to indicate that there is not much DM around those GCs, constraining their mass-to-light ratios to $`\mathrm{M}/\mathrm{L}<2.5`$ (Moore 1996).
## 4. YSCs in the Antennae: a 1<sup>st</sup> example
The Antennae galaxies (= NGC 4038/39) is an interacting pair of two gas rich spirals, probably Sc, of comparable mass. In the ongoing starburst triggered by the interaction a population of bright YSCs is formed, numerous enough to allow for the first time to reasonably define a LF. The LF from WFPC1 observations is a power law with slope $`\alpha 1.8`$ (Whitmore & Schweizer 1995 (WS95)). WFPC2 reobservations are going to be presented by Miller (this conf.).
Assuming a homogeneous metallicity $`\mathrm{Z}1/2\mathrm{Z}_{}`$ lack of individual cluster spectroscopy and in analogy to the YSC system in NGC 7252, I analysed the WFPC1 data of WS95 with evolutionary synthesis models for SSPs. In a first step, the average dereddened $`(\mathrm{V}\mathrm{I})`$ colour of the 550 YSCs is used to derive a mean age of the YSC population of $`(2\pm 2)10^8`$ yr, consistent with Barnes’ (1988) dynamical model for the interaction between NGC 4038 and 4039 and consistent with Kurth’s (1996) global starburst age. SSP models describe both the fading and the reddening of the YSC population as it ages. Assuming a mean age for the YSC population, both the LF and the colour distribution of the YSCs are simply shifted towards fainter magnitudes and redder colours, respectively, without changing shape.
### 4.1. Evolution of the YSC LF
In a second step, individual ages are derived for all the YSCs on the basis of their individual $`(\mathrm{V}\mathrm{I})`$ and $`(\mathrm{U}\mathrm{V})`$ colours. Interestingly, the resulting age distribution not only shows a peak at very young ages $`0410^8`$ yr for the YSCs, but also a contribution from $`12`$ Gyr old GCs from the parent galaxies (Fig. 1a). A small number of apparent interlopers are probably due to inhomogenities in the dust distribution. It is improbable that many of the old GCs we identify are highly reddened YSCs, since they would have to be exeedingly bright intrinsically. The fact that age estimates from $`(\mathrm{U}\mathrm{V})`$ colours, as far as available, agree with age estimates from $`(\mathrm{V}\mathrm{I})`$ supports our metallicity assumption and makes us hope that the average $`\mathrm{E}_{\mathrm{B}\mathrm{V}}`$ is correct for most of the clusters. It is clear, however, that the individual YSC extinctions are a major source of uncertainty in this analysis. With individual YSC ages and observed luminosities, SSP models can be used to predict the time evolution of cluster luminosities. Neglecting any kind of dynamical cluster destruction effects, and only using the young star clusters identified, we obtain the surprising result that by an age of 12 Gyr, when age differences among individual YSCs will be negligible, their LF will have evolved from the presently observed power law into a fairly normal Gaussian GC LF with parameters $`\mathrm{M}_\mathrm{V}=6.9`$ mag and $`\sigma (\mathrm{M}_\mathrm{V})=1.3`$ mag (Fig. 1b). The turn-over then is $`>1`$ mag brighter than the completeness limit which evolves to $`\mathrm{M}_\mathrm{V}=5.7`$ mag, and fainter by $`0.4`$ mag than for typical GC systems. This is a consequence of the enhanced metallicity of the YSC population with respect to that of old GC systems, and in agreement both with observations of old GC systems with a range of metallicities (Ashman et al. 1995) and with our SSP model predictions. With this surprising result we confirm and quantify Meurer’s conjecture that over a Hubble time, age spread effects can transform an observed power law LF of YSCs into the Gaussian LF of old GCs (cf. F.-v.A. 1998 for details).
The bright end of the LF defined by the old GCs from the parent spirals is well described by a Gaussian LF with parameters $`\mathrm{M}_\mathrm{V}=7.3`$ mag and $`\sigma (\mathrm{M}_\mathrm{V})=1.2\mathrm{mag},`$ normalised to the total number of GCs in the Milky Way and Andromeda galaxies together (Fig. 2a). Hence, if the two interacting spirals NGC 4038 and 4039 had a similar number of GCs as those galaxies, and if the bulk of the YSC population really are young GCs, Zepf & Ashman’s (1993) requirement, that the number of secondary GCs formed in mergers should be comparable to the number of primary GCs in the two progenitor spirals, would be fulfilled. This requirement is necessary if the higher specific GC frequency in ellipticals – as compared to spirals – is to be compatible with a spiral-spiral merger origin for those ellipticals.
### 4.2. MF of YSCs
To investigate whether the LFs and MFs of old GC systems are determined by the cluster formation process or whether they are the result of secular dynamical destruction processes, we derive the MF of the very young star cluster system in the Antennae. We restrict ourselves to YSCs brighter than the completeness limit, use individual YSC ages and luminosities, and combine them with model M/L for the respective YSC ages to determine their individual masses. Our tentative conclusion is that for all the 393 YSCs brighter than the completeness limit and with $`(\mathrm{V}\mathrm{I})`$ colours available, the MF is compatible with a Gaussian MF with parameters $`\mathrm{Log}(\mathrm{M}_{\mathrm{YSC}}/\mathrm{M}_{})=5.6,\sigma 0.46`$ $`\mathrm{Mass}(\mathrm{YSC})410^5\mathrm{M}_{}`$ (Fig. 2b). If we include YSCs fainter than the completeness limit, neither the shape nor the parameters of the MF are changed (cf. F.-v.A. 1999 for details).
The uncertainty in the YSC metallicity leads to age uncertainties, which, in turn, lead to uncertainties in model M/L values of the order of 10 % in these early stages. Inhomogenities in the dust distribution do not seem to be very important for the brightest YSCs for which ages from $`(\mathrm{U}\mathrm{V})`$ agree with ages from $`(\mathrm{V}\mathrm{I})`$ colours. We do not know, however, how important they are for those YSCs that are not detected in U. Hence, on the basis of $`(\mathrm{V}\mathrm{I})`$ alone, we cannot quantify in how far an inhomogeneous internal reddening might affect our results. For this reason and since our analysis is based on WFPC1 data, we caution that our conclusions concerning the evolution of the YSC LF and their MF can only be preliminary. As soon as the reduced WFPC2 data presented by Miller will become available to us, we shall repeat our analysis and supplement it with simulations to estimate the effects of differential observational uncertainties.
### 4.3. Implications
If, however, our preliminary result became confirmed, this would imply that the log-normal MF of old GCs is produced by the cluster formation process rather than by secular dynamical evolution of the cluster system. In this context. it seems very interesting to obtain observational information about the molecular cloud mass spectrum in massive interacting galaxies. Jog & Solomon (1992) conjecture that the strongly enhanced ambient pressure in massive gas rich mergers might affect the molecular cloud structure. In Ultraluminous Infrared Galaxies – which all are mergers with strong starbursts – the fraction of gas at very high densities of $`\mathrm{n}10^4\mathrm{and}10^5\mathrm{cm}^3`$, as traced by HCN and CS lines, with respect to gas at $`\mathrm{n}500\mathrm{cm}^3`$, as traced by CO, is indeed observed to be higher by 1 – 2 orders of magnitude (e.g. Solomon et al. 1992).
### 4.4. Discussion
Even with deeper WFPC2 data, however, the MF of the YSCs in the Antennae does not yet seem to be unambiguously settled. Zhang & Fall (1999) use reddening-free colour indices Q<sub>1</sub>, Q<sub>2</sub>, solar metallicity Bruzual & Charlot models, consider incompleteness and stellar contamination, and find power law MFs with slopes $`\alpha 2`$ for YSCs in the two age intervals 2.5 – 6.3 Myr and 25 – 160 Myr where Q<sub>1</sub> and Q<sub>2</sub> give unambiguous results. In an independent analysis of the same data, Miller (this conf.) finds a significant flattening in the observed LF at $`\mathrm{M}_\mathrm{V}10.7`$ mag. Intriguingly, this value is exactly the turnover luminosity implied by our MF at the mean YSC age in the absence of an age spread. Miller’s method of reconstruction, however, leads to a power law YSC MF. Clearly, more work both from the modelling and observational sides is needed before we really understand the YSC MF.
### 4.5. Other systems
For the YSCs in NGC 7252 and NGC 3921, Miller & Fall (1997) and Miller (this conf.) prefer power law MFs. Distances to those systems, however, are larger by factors 3 and 4, respectively, pushing the completeness limit to higher luminosities, even in WFPC2 data. Moreover, the YSCs in these dynamically old merger remnants have higher mean ages (650 – 750 Myr for NGC 7252, 250 – 750 Myr for NGC 3921, as compared to 200 Myr for the Antennae), and hence, are intrinsically fainter already by 0.8 – 1 mag, on average. If we assume that the MFs for the YSCs in NGC 7252 and NGC 3921 were identical to the one we derived for the YSCs in the Antennae, we obtain turn-over luminosities $`\mathrm{M}_\mathrm{V}=10\mathrm{}9.5`$ mag for NGC 7252 and $`\mathrm{M}_\mathrm{V}=9.5\mathrm{}9`$ mag for NGC 3921 at their respective mean YSC ages. In both cases, these turn-over luminosities are close to the completeness limits. The difficulty to disentangle the old GC and the YSC populations also increases with increasing YSC age. We therefore doubt that it is possible to track the MF beyond a possible turnover in NGC 7252 and 3921. In any case, the YSCs in NGC 7252 and 3921 have survived $`10`$ crossing times and from this fact alone are young GCs rather that open ones, as also indicated by their combination of small effective radii and high luminosities (cf. Miller et al. 1997, Schweizer et al. 1996).
### 4.6. Dynamical Evolution
Stellar mass loss is included in our models for the evolution of M/L. Over 12 Gyr, the masses of YSCs will decrease by $`10\mathrm{and}15\%`$ for a Salpeter and Scalo-IMF, respectively. About half of the entire stellar mass loss occurs during the $`1^{\mathrm{st}}`$ Gyr. External dynamical effects from an interaction of the clusters with the potential of the interacting galaxy system seem extremely difficult to model. Comparison of YSC LFs and MFs in an age sequence of interacting galaxies, mergers, merger remnants, and dynamically young ellipticals will allow to “see these processes at work”.
In the Milky Way potential, Vesperini (1998) has shown that an assumed initially log-normal GC MF is conserved in shape and parameters during self-similar evolution over a Hubble time despite the destruction of $`50\%`$ of the cluster population. If, on the other hand, he starts with a power law MF, severe fine-tuning is required for his model parameters to secularly transform it into the observed log-normal MF of old GC systems.
It is clearly important to analyse more YSC systems in order to see if (and in how far) their MFs are universal or might depend on environment. Old GC systems have their turn-over around $`\mathrm{M}_\mathrm{V}7.2`$ mag. With 10m telescopes they are accessible to more than Virgo cluster distances. MOS – e.g. with FORS on the VLT – in combination with HST imaging will allow to determine cluster abundances and hence to more precisely age-date them, and may even provide kinematic information for independent mass estimates.
An open question seems to me if the (globular) cluster formation process in the high metallicity environment of interacting spirals today is the same or not as it was in the Early Universe when the radiation field was stronger and the metallicity lower.
#### Acknowledgments.
I wish to thank the organisers and in particular Ariane Lançon for a stimulating workshop in a very pleasant and warm atmosphere in cold and rainy Strasbourg.
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# Punctured polygons and polyominoes on the square lattice
## 1 Introduction
A self-avoiding polygon (SAP) can be defined as a walk on a lattice which returns to the origin and has no other self-intersections. Alternatively we can define a SAP as a connected sub-graph (of a lattice) whose vertices are of degree 2. The history and significance of this problem is nicely discussed in . Generally SAP are considered distinct up to a translation, so if there are $`p_n`$ SAP of length $`n`$ there are $`2np_n`$ returns to the origin (the factor of two arising from the two possible directions in which the loop can be traveled). Staircase polygons are a special case of SAP. A staircase polygon can be viewed as the intersection of two directed walks starting at the origin, taking steps either to the right or down and terminating once the walks meet. We define a punctured SAP as a SAP with one or more holes, with the perimeter of each hole itself being a SAP. In other words a punctured SAP is a SAP enclosing one or more SAP. To avoid any possible confusion in our definition of punctured polygons, note that the punctures are disjoint — there are no degree four vertices in the punctured polygons we are considering. A similar definition can be given for punctured staircase polygons. We have shown an example of each case in figure 1. The two principal questions one can ask are “how many polygons, distinct up to a translation, are there of perimeter $`2n`$ (including the perimeter of the holes) with $`k`$ punctures?” and “how many polygons, distinct up to a translation, are there of area $`n`$ with $`k`$ punctures?”
A polyomino is defined as the union of connected cells, where a cell is a unit square with 4 edges (and its interior). Two cells are said to be joined if they share a common edge, and are said to be connected if there exists a sequence of cells joining the two cells such that successive pairs of cells are joined. A $`k`$-punctured polyomino is a polyomino with $`k`$ holes. Unlike the situation for punctured polygons, degree 4 vertices are permitted, but two holes meeting only at a vertex do count as two holes, rather than one, as they are not joined. Punctured polygons are a proper subset of punctured polyominoes. The difference between them is that punctured polygons do not include those polyominoes where there are 2 cells touching each other only at a vertex. An example is shown in figure 2.
Reverting to $`k`$-punctured polygons, for $`k=0`$ we have the simpler — but still unsolved — problem of the number of polygons of perimeter $`2n`$ or area $`n.`$ Both these problems have been extensively studied for several decades. The most recent result for polygon perimeters appears to be where polygons of perimeter up to 90 steps are given. In that paper our analysis of the polygon perimeter generating function led us to conclude that
$$P^{(0)}(x)=\underset{n}{}p_{2n}^{(0)}x^nA^{(0)}(x)+B^{(0)}(x)(1\mu ^2x)^{2\alpha },$$
where $`p_{2n}^{(0)}`$ is the number of SAP of perimeter $`2n`$, distinct up to translations. The analysis in yielded a very accurate estimate for the connective constant $`\mu =2.63815853034(10)`$ and confirmed the theoretical prediction $`\alpha =1/2`$ . Furthermore, we obtained estimates for the critical amplitudes $`A^{(0)}(x_c)0.036`$ and $`B^{(0)}(x_c)0.234913,`$ where $`x_c=1/\mu ^2`$. We also concluded that there was no evidence for a non-analytic correction-to-scaling exponent, so that we expect the asymptotic form of the coefficients to behave as:
$$p_{2n}^{(0)}\mu ^{2n}n^{\frac{5}{2}}[a_1+a_2/n+a_3/n^2+a_4/n^3+\mathrm{}].$$
The connective constant $`\mu `$ is of course the same as that for self-avoiding walks on the same lattice .
For polygon areas the most recent published work appears to be in which the first 20 terms of the area generating function were given and analysed. In that work it was found that
$$A^{(0)}(y)=\underset{n}{}a_n^{(0)}y^nC^{(0)}(y)+D^{(0)}(y)\mathrm{log}(1\kappa y),$$
where $`a_n^{(0)}`$ is the number of SAP of area $`n`$, $`\kappa 3.97087,`$ and the amplitudes $`C^{(0)}`$ and $`D^{(0)}`$ were not estimated. Recently this series has been extended to 26 terms , but in the present work we have devised a new and exponentially faster algorithm, as a result of which we have extended the series to 42 terms, and we present an analysis of this longer series. The connective constant $`\kappa `$ is found to be slightly smaller than that for the related problem of polyominoes .
In the following, we refer to the boundary of a polygon and its interior as a disc and so we will be discussing punctured discs. An unpunctured disc is a SAP.
For punctured discs, the basic problem is, analogously, the calculation of the generating functions
$$P^{(k)}(x)=\underset{n}{}p_{2n}^{(k)}x^nB^{(k)}(x)+C^{(k)}(x)(1(\mu ^{(k)})^2x)^{2\alpha _k},$$
(1)
and
$$A^{(k)}(y)=\underset{n}{}a_n^{(k)}y^nD^{(k)}(y)+E^{(k)}(y)(1\kappa ^{(k)}y)^{\beta _k},$$
(2)
where the superscript $`k`$ refers to the number of holes, or punctures. From the generating function, one wishes to deduce the asymptotic behaviour, believed to be as shown on the r.h.s. of the above equations. The major problem to be investigated is how the behaviour of $`P^{(k)}(x)`$ and $`A^{(k)}(y)`$ changes as $`k`$ is increased. Previous work has been confined to the study of punctured SAP by area. There it was proved that $`\kappa ^{(k)}=\kappa ^{(0)}=\kappa ,`$ and that if the exponent exists, $`\beta _k=\beta _0+k.`$ These results apply more generally to punctured surfaces, but in this work we are confining ourselves to two dimensions. As far as we are aware, there has been no previous work on the problem of the perimeter generating function of punctured discs.
The problem is interesting for several reasons. The effect of a change in geometry is a much studied topic in lattice statistics, and our study of the change in perimeter exponent with punctures seems to be entirely new. It has only been possible by the algorithms we have designed and implemented, which are exponentially faster than pre-existing algorithms. A number of related models have been studied previously, such as $`c`$-animals and the behaviour of prime knots in polygons . $`c`$-animals are lattice animals with exactly $`c`$ cycles. In it was proved that if the number of such animals $`a_n(c)\lambda _c^nn^{\theta _c}`$ as $`n\mathrm{}`$, $`c`$ fixed, then $`\theta _c=\theta _0+c`$ provided that $`\theta _0`$ exists. It had been previously proved that $`\lambda _c=\lambda _0.`$ The change in connective constant of $`c`$-animals as the number of cycles per vertex changes from zero to non-zero is discussed — among other results — in . Similarly, in a numerical study of knotted polymers , it was conjectured that the exponent $`\alpha `$ depends on the number of prime knots $`n_p`$ that arise in the knot decomposition of a given SAP via the relation $`\alpha (n_p)=\alpha (0)+n_p.`$
Furthermore, there is considerable pedagogical connection between some of these previous studies and our work here, and also between our work and the study of branched polymers. To sketch this connection, we first remark that the number of polyominoes, also called lattice animals, is just the number of (strongly embedded) site animals. This connection is readily seen by placing a site at the centre of every cell of the polyomino, and joining those sites corresponding to joined cells by a bond. In this way, every lattice animal is mapped to a distinct site animal and vice versa. All other models we are considering map similarly to a different subset of strongly embedded site animals. For example, SAP map to site animals whose only cycles are 4-cycles (which may be isolated or joined). Punctured polygons map to site animals with larger cycles, as well as 4-cycles, whereas punctured polyominoes correspond to site animals with more complex restrictions. Thus this study also complements the earlier studies of $`c`$-animals. In those studies, the variation of exponent with the number of cycles is considered, whereas in this study we are varying the types of allowed cycles.
In order to study these and related systems, when an exact solution can’t be found one has to resort to numerical methods. For many problems the method of series expansions is by far the most powerful method of approximation. For other problems Monte Carlo methods are superior. For the analysis of $`P^{(k)}(x)`$ and $`A^{(k)}(y)`$, series analysis is undoubtedly the most appropriate choice. This method consists of calculating the first few coefficients in the expansion of the generating function. Given such a series, using the numerical technique known as differential approximants , highly accurate estimates can frequently be obtained for the critical point and exponents, as well as the location and critical exponents of possible non-physical singularities. Other numerical methods are discussed in , and those used in the present study are described more fully below.
In the next section we will describe the finite lattice method for enumerating punctured polygons. In Section 3 we prove the invariance of $`\mu ^{(k)}`$ as $`k`$ changes, and give an heuristic argument for the exponent shift with $`k.`$ The results of the analysis of the series are presented in Section 4. Results analogous to those known for punctured polygons by area are proved for punctured polyominoes in the Appendix.
## 2 Enumeration of punctured polygons
### 2.1 Enumeration of punctured self-avoiding polygons
The method used to enumerate punctured self-avoiding polygons on the square lattice is a generalisation of the method devised by Enting for the enumeration of ordinary SAP. In the following we first describe the original method in some detail and show how simple it is to generalize the method to the enumeration of punctured polygons. The first terms in the series for the generating function can be calculated using transfer matrix techniques to count the number of polygons in rectangles $`W+1`$ edges wide and $`L+1`$ edges long. The transfer matrix technique involves drawing a line through the rectangle intersecting a set of $`W+2`$ edges. For each configuration of occupied or empty edges along the intersection we maintain a (perimeter) generating function for loops to the left of the line cutting the intersection in that particular pattern. Polygons in a given rectangle are enumerated by moving the intersection so as to add one vertex at a time, as shown in figure 3. Since the loops are non-intersecting, each configuration can be represented by an ordered set of edge states $`\{n_i\}`$, where
$$n_i=\{\begin{array}{cc}\hfill 0& \text{empty edge},\hfill \\ \hfill 1& \text{lower part of loop closed to the left},\hfill \\ \hfill 2& \text{upper part of loop closed to the left}.\hfill \end{array}$$
Configurations are read from the bottom to the top. So the configuration along the intersection of the polygon in figure 3 is $`\{0112122\}`$. In passing it is worth noting that there are some major restrictions on the possible configurations. Firstly, since all loop-ends are connected to the left of the intersection, every lower loop end must have a corresponding upper end, and it is therefore clear that the total number of ‘1’s is equal to the total number of ‘2’s. Secondly, as we look through the configuration starting from the bottom the number of ‘1’s is never smaller than the number of ‘2’s.
In Table 1 we have listed the possible local ‘input’ states and the ‘output’ states which arise as the kink in the intersection is propagated by one step. Some of these update rules are illustrated further in Fig 4. The first panel represents the input states ‘10’ and ‘01’ and the possible output states are also ‘01’ and ‘10’. The second panel represents the input state ‘11’ as part of the configuration $`\{01122\}`$. In this case we connect the two loop ends, but in doing so we see that the upper part of the second loop before the move becomes the lower part of the one remaining loop after the move. That is the configuration $`\{01122\}`$ becomes $`\{00012\}`$. This relabelling of the other loop-end when connecting two ‘1’s (or two ‘2’s) is denoted by over-lining in Table 1. In general there could be more loops nested in between the two ‘1’s and the corresponding ‘2’ at the other end of the loop. Say for instance we had the configuration $`\{11121222\}`$ and connected the first two ‘1’s then the new configuration of unconnected loop ends would be $`\{00121212\}`$ (drawing a little figure makes this quite clear). The general rule for the relabelling is as follows: When connecting two ‘1’s (‘2’s) we work upward (downward) in the configuration, counting the number of ‘1’s and ‘2’s we pass until the number of ‘2’s (‘1’s) exceeds the number of ‘1’s (‘2’s). This ‘2’ (‘1’) is the other end of the inner loop and it should now be changed to a ‘1’ (‘2’), thus becoming the lower (upper) end of the outer loop (again drawing a few pictures should make this clearer). The weights corresponding to these configuration transformations are simply calculated by counting the number of steps which have been added to the polygon. Note that the input state ‘12’ is special because connecting the two ends results in a closed loop, so this is only allowed if there are no other loops cut by the intersection and the result is a valid polygon, which is then accumulated in the total count for that particular length. Failure to observe this restriction would result in graphs with disconnected components, either one polygon over another or a polygon within another (this latter case is of course of interest when we wish to enumerate punctured polygons). This is illustrated in figure 5 where we show the possible ways a pair of loops can be placed relative to one another (first panels), how the loops can be connected to produce a valid configuration (middle panels) and the ways of connecting loops that lead to invalid graphs containing disconnected components (last panels). In this figure the invalid SAP in the last panel on the bottom could result in a valid punctured SAP. Note also that from the input state ‘00’ we can produce the output state ‘12’ only if there are other loops crossing the intersection (otherwise we would produce disconnected polygons sitting side by side). We refer the interested reader to for further details regarding the encoding and relabelling of configurations.
Due to the obvious symmetry of the lattice one need only consider rectangles with $`LW`$. In the original approach valid polygons were required to span the enclosing rectangle in the lengthwise direction. So it is clear that polygons with projection on the $`y`$-axis $`<W`$, that is polygons which are narrower than the width of the rectangle, are counted many times. It is however easy to obtain the polygons of width exactly $`W`$ and length exactly $`L`$ from this enumeration . Any polygon spanning such a rectangle has a perimeter of length at least $`2(W+L)`$. By adding the contributions from all rectangles of width $`WW_{\mathrm{max}}`$ (where the choice of $`W_{\mathrm{max}}`$ depends on available computational resources) and length $`WL2W_{\mathrm{max}}W+1`$, with contributions from rectangles with $`L>W`$ counted twice, the number of polygons per vertex of an infinite lattice is obtained correctly up to perimeter $`4W_{\mathrm{max}}+2`$.
With the original algorithm the number of configurations required as $`W_{\mathrm{max}}`$ increased grew asymptotically as $`3^{W_{\mathrm{max}}}`$ . In a recent improvement of the algorithm valid polygons were required to span the rectangle in both directions. In other words we directly enumerate polygons of width exactly $`W`$ and length $`L`$. For each configuration of partially completed polygons we keep track of the current minimum number of steps $`N_{\mathrm{cur}}`$ that have been inserted to the left of the intersection and we calculate the minimum number of additional steps $`N_{\mathrm{add}}`$ required to produce a valid polygon that spans a rectangle of size at least $`W\times W`$. If the sum $`N_{\mathrm{cur}}+N_{\mathrm{add}}>4W_{\mathrm{max}}+2`$ the partial generating function for that configuration was discarded because it would make no contribution to the polygon count up to the perimeter lengths we were trying to obtain. Numerical evidence indicated that the computational complexity was reduced significantly. While the number of configurations still grew exponentially as $`\lambda ^{W_{\mathrm{max}}}`$ the value of $`\lambda `$ was reduced from $`\lambda =3`$ to $`\lambda 2`$ with the improved algorithm. Furthermore, for any $`W`$ we know that contributions will start at $`4W`$ since the smallest polygons have to span a $`W\times W`$ rectangle, so for each configuration we need only retain $`4(W_{\mathrm{max}}W)+2`$ terms of the generating functions while in the original algorithm contributions started at $`2W+2`$ because the polygons were required to span only the length-wise direction.
The generalization to enumeration of punctured polygons is obtained by noting that a closed loop is formed whenever we connect a 1-edge to a 2-edge immediately above. If these two edges were the only loop ends in the intersection we would have formed a valid polygon. In other cases we need to ensure that the resulting polygon is a valid punctured polygon. That is, we must ensure that the separate SAP just formed will be completely enclosed within a larger polygon. So we wish to avoid forming separable polygons, as shown in the upper right panel of figure 5, and ‘holes within holes’. As it turns out the rule for a valid ‘12’ closure is simply that we can connect the two loop ends provided there is an odd number of loop ends below the loop being closed. To see this consider that as we go through a configuration we note that each time we pass a loop end we go from the outside of the polygon to the inside and visa versa. So in a configuration all lattice cells between the first and second loop ends will lie inside the finished polygon, lattice cells between the second and third loop end will lie outside the polygon, and so on. Thus we see that by closing a loop which has an odd number of loop ends below it we are closing off a part of the lattice which will lie inside the finished polygon. In particular we see that we avoid the situation shown in the upper right panel of figure 5 with graphs containing disconnected pieces one over another. Likewise we avoid creating graphs with disconnected pieces sitting side by side.
We also avoid forming holes within holes because closing a loop around the hole would be prohibited since there would be an even number of loop ends below the ‘1’-edge (except of course when forming a completed punctured polygon). Let us look at the possible edge-configurations around a puncture in some detail. First look at the configuration {$`\mathrm{}1122\mathrm{}`$} (where the $`\mathrm{}`$ are any edge configurations with an even number of edges making the total configuration of edges a valid intersection). The outer ‘12’-edges can’t be connected (the number of edges below the ‘1’ is even) so we have to connect two other edges on either side of the hole diminishing the number of edges by two and possibly changing the edge labels on either side of the hole (in which case we end up with one of the subsequent cases). In the top panel of figure 6 we show how the simplest interesting case {11112222} leads to only valid punctured polygons. Secondly, in the case {$`\mathrm{}1121\mathrm{}`$}, we can connect the two ‘1’-edges forming a partial loop enclosing the hole. In doing so the matching ‘2’-edge of the second ‘1’-edge is changed to a ‘1’-edge because it now is the new lower edge of the larger loop formed by connecting the two original ‘1’-edges. The important thing to notice is that the number of edges below this new ‘1’-edge is even so we cannot connect it to a ‘2’-edge immediately above and we therefore do not form a hole within a hole. In the middle panel of figure 6 we show the simplest interesting case {1111212222} and demonstrate that it leads to valid punctured polygons. Thirdly, the configuration $`\{\mathrm{}2121\mathrm{}`$} is not very interesting since connecting the ‘21’-edges in front of the hole does not result in a loop partially enclosing the hole. Finally, the configuration {$`\mathrm{}2122\mathrm{}`$} is obviously merely the mirror image of the second case. Note that this covers all cases since more complicated configurations merely would correspond to more convoluted loop structures.
In this work we use the generalisation of the algorithm of to count the number of punctured SAPs with up to 3 holes. Obviously, the smallest hole we can make in a SAP has perimeter 4 so the number of punctured SAPs is obtained correctly up to perimeter $`4W_{\mathrm{max}}+2+4k`$.
The algorithm used for the enumeration of punctured SAPs by area is a simple variation of the algorithm described above. The encoding of configurations along the intersection and the transformation of these configurations as the intersection is moved remain the same. The only change is that the weights are different. In order to count the enclosed area we proceed as follows: A unit of area may be added as the kink is moved to a new lattice cell (in figure 3 this is the cell in which the dotted lines meet). Whether or not a unit of area is added is determined by whether or not this lattice cell is inside or outside the polygon. But we already know from the arguments given above that a lattice cell is inside the polygon if the number of loop ends below the cell is odd (note that any loop end along the vertical edge cutting the horizontal part of the kink is included in the count). So in the case shown in figure 3 the lattice cell to which the kink is moved lies outside the polygon (there are two loop ends below the kink) and no unit of area is added. Note that a unit of area may be added for any given output configuration. In this case the area generating function is obtained correctly to $`2W_{\mathrm{max}}+3k`$. The factor $`3k`$ arise since it takes 3 lattice cells to completely surround the simplest puncture, which is just a single cell.
As far as we are aware, series for punctured polygons have not been derived previously. In such circumstances it is even more important than usual to undertake careful tests to ensure that the series are correct. To this end, a second algorithm was implemented, by one of the authors, to independently evaluate the series coefficients. The complete agreement we obtained between the two data sets reassures us as to the correctness of our results.
### 2.2 Enumeration of punctured staircase polygons
The enumeration of punctured staircase polygons is much simpler. In fact as we shall demonstrate it is a problem for which the computational complexity grows only as a polynomial in the number of terms. As stated in the Introduction we can think of staircase polygons as consisting of two non-intersecting directed walks on the square lattice. Punctured staircase polygons naturally will have more than two walks, in fact up to $`2k+2`$ walkers can be present in any given column. Due to the restrictions on staircase polygons it follows that a punctured staircase polygon is formed by requiring that the walks be directed and that any two walkers starting at one point join each other later without intersecting other walks. This can be encoded in a transfer matrix calculation as follows: Again we count the number of punctured polygons in rectangles as before and draw a line through the lattice intersecting the $`W+2`$ edges. In this case we need only specify whether or not each edge is part of the polygon or not, so each configuration can be represented by an ordered set of edge states $`\{n_i\}`$, where
$$n_i=\{\begin{array}{cc}\hfill 0& \text{empty edge},\hfill \\ \hfill 1& \text{part of loop closed to the left}.\hfill \end{array}$$
This uniquely specifies the configuration because the first occupied edge is connected to the last occupied edge and any edges in between are paired, e.g., the second and third (fourth and fifth and so on) edges form a loop to the left and has to be connected to each other later on. The rules for updating the configurations are as follows: From the local input state ‘00’ we can always get the output ‘00’, and the output ‘11’ provided there is an odd number of edges below the kink and also that the edge directly below the kink is empty. The rules for the output ‘11’ ensure that the new pair of walkers lie within the enclosing staircase polygon but not inside an internal staircase polygon (thus preventing holes within holes), and that the lower walk does not intersect other walks. From the local input state ‘01’ and ‘10’ we can always produce the output ‘01’, and the output ‘10’ provided either that the edge directly below the kink is empty or there is an even number of edges below the kink. These rules ensure that the walkers do not intersect other walks except when we close a valid loop. Finally from the local input state ‘11’ we always get the output state ‘00’. The weights associated with these updates are obtained in the same way as for the SAP enumeration, whether the enumeration is done by perimeter or area.
As in the previous case we calculate the number of punctured staircase polygons spanning the rectangles. Adding the contributions from all rectangles of width $`WW_{\mathrm{max}}`$ and length $`WL2W_{\mathrm{max}}W+1`$ the number of punctured staircase polygons is obtained correctly up to perimeter $`4W_{\mathrm{max}}+2+4k`$. Note that for fixed $`k`$ the maximum number of configurations $`N_C`$ grows as a polynomial in $`W_{\mathrm{max}}`$
$$N_C=\underset{j=0}{\overset{k}{}}\left(\begin{array}{cc}W_{\mathrm{max}}+1& \\ 2j+2& \end{array}\right).$$
So in this case the algorithm is of polynomial complexity.
## 3 Expected behaviour
As mentioned in the introduction, the problem of enumerating SAP by area has been extensively studied . It has been shown that $`\kappa ^{(k)}=\kappa ^{(0)},`$ and that if the exponents exist, $`\beta _k=\beta _0+k`$ . That is to say, for $`k`$ finite, the connective constant for $`k`$-punctured discs, by area, is the same as that for unpunctured discs, while the critical exponent increases by 1 for each puncture. We repeat that the punctures are disjoint.
As far as we are aware, there has been no previous work on the problem of the perimeter generating function of punctured SAP. We first give a proof that a 1-punctured SAP on the square lattice has the same connective constant as an unpunctured SAP, and indicate how that proof can be generalised to $`k`$-punctured SAP.
Before stating and proving the relevant theorem, we require certain preliminary results. For an (unpunctured) SAP of perimeter $`2n`$ one knows, eqn. (7.101), that
$$\mathrm{exp}[b\sqrt{n}]\mu ^{2n}/4np_{2n}^{(0)}\mu ^{2n},$$
(3)
where $`b`$ is a constant. Further, it is obvious that a polygon of perimeter $`2n`$ has maximum area $`n^2/4,`$ which occurs when the shape is an $`n/2\times n/2`$ square. Thus the area of a polygon of perimeter $`2n,`$ denoted $`𝒜_{2n},`$ satisfies $`𝒜_{2n}n^2/4.`$
Now consider 1-punctured polygons of total perimeter $`2n,`$ with inner perimeter $`2m`$ and outer perimeter $`2(nm)`$. It is possible for the inner polygon to be Hamiltonian, in which case its perimeter is greater than that of the surrounding polygon. Indeed, a square polygon of side $`2n+1,`$ and hence of perimeter $`8n+4`$ can contain an internal polygon of perimeter as large as $`4n^2.`$ Hence the semi-perimeter of the inner polygon, $`m,`$ can range from a minimum value of 2 to a maximum value of $`n2\sqrt{2n}+2<n\sqrt{n}`$ for $`n>1.`$
With these preliminaries, we can now state and prove the following theorem:
Theorem: $`lim_n\mathrm{}\frac{1}{2n}\mathrm{log}p_{2n}^{(1)}=\mathrm{log}\mu ,`$ where $`\mu `$ is the same constant as appears in the corresponding limit for unpunctured SAP.
Proof: 1-punctured polygons of total perimeter $`2n`$ are constructed by placing polygons $`P`$ of perimeter $`2m`$ inside polygons $`Q`$ of perimeter $`2n2m.`$ Let $`w(P,Q)`$ denote the number of ways of placing polygon $`P`$ inside polygon $`Q,`$ (which of course depends on both $`P`$ and $`Q.`$) Then
$$p_{2n}^{(1)}=\underset{m=2}{\overset{n\sqrt{n}}{}}\underset{P}{}\underset{Q}{}w(P,Q).$$
We bound this summand by the product of three factors. The first two factors are the number of polygons of perimeter $`2m`$ and $`2n2m`$ respectively. The third is the number of ways the smaller polygon can be placed inside the surrounding polygon and is clearly less than or equal to the area of the surrounding polygon. Explicitly,
$$p_{2n}^{(1)}\underset{m=2}{\overset{n\sqrt{n}}{}}(nm)^2\mu ^{2m}\mu ^{2(nm)}/4\mu ^{2n}/4\underset{m=2}{\overset{n}{}}(nm)^2n^3\mu ^{2n}.$$
To obtain a lower bound, consider a $`3\times 3`$ square polygon with a unit square hole at its centre. This unique realisation of $`p_{16}^{(1)}`$ can be uniquely concatenated with each unpunctured polygon by joining them at a solitary specified edge, and then deleting that edge. For each unpunctured polygon we take the set of left-most vertical edges and choose the bottom edge from this set, and we choose the right-most, top-most vertical edge of $`p_{16}^{(1)}`$ (a similar operation is shown for polyominoes in figure A1). The concatenation operation then gives, for each of the $`p_{2n14}`$ unpunctured polygons, a unique member of the set of 1-punctured polygons in which the puncture is a single cell. Thus we have $`p_{2n14}^{(0)}p_{2n}^{(1)}.`$ From the above equations we thus obtain
$$\mathrm{exp}[b\sqrt{n7}]\mu ^{2n14}/4(n7)p_{2n14}^{(0)}p_{2n}^{(1)}n^3\mu ^{2n}$$
for $`n>8.`$ The theorem then follows immediately on taking logarithms, dividing through by $`2n`$ and taking the limit as $`n\mathrm{}.`$
This proof can clearly be extended to two-punctured polygons, then to three-punctured polygons etc., by concatenating unpunctured polygons with minimal two-punctured, three-punctured etc. polygons. In the appendix we prove that $`k`$-punctured polyominoes have the same growth constant as unpunctured SAP by area.
We have been unable to prove a result analogous to the result for the exponent of punctured SAP by area, but give an argument that depends on certain assumptions that are generally accepted, though not proved. In the case of staircase polygons however our assumptions have been proved, and so our result will be rigorously true.
The key results we need are those obtained in to the effect that the mean area of SAP of perimeter $`2n`$ is proportional to $`n^{1.5}.`$ This is true both for SAP and for staircase polygons, and in the latter case it has been proved. More precisely, we need the following result. There exists constants $`D_1`$ and $`D_2`$ such that the mean area $`\overline{A}_{2n}`$ of polygons of perimeter $`2n`$ satisfies
$$D_1n^{\frac{3}{2}}\overline{A}_{2n}=\frac{_QA_Q}{p_{2n}}D_2n^{\frac{3}{2}},$$
where the sum is taken over all $`p_{2n}`$ polygons $`Q`$ of perimeter $`2n.`$
Further, there exists constants $`C_1`$ and $`C_2`$ such that the number of polygons of perimeter $`2n`$ satisfies
$$C_1\mu ^{2n}n^{\frac{5}{2}}p_{2n}^{(0)}C_2\mu ^{2n}n^{\frac{5}{2}}.$$
For staircase polygons, $`\mu =2`$ and the exponent is $`\frac{3}{2}`$ instead of $`\frac{5}{2}`$ in the above equation.
As above, let $`Q`$ denote a polygon of perimeter $`2n2m,`$ with area $`A_Q.`$ The number of ways of placing a given polygon $`P`$ of perimeter $`2m`$ inside $`Q`$ is clearly less than $`A_Q.`$
Thus the number of ways of placing $`P`$ inside any polygon of perimeter $`2n2m`$ is less than
$$\underset{Q}{}A_Q=p_{2n2m}^{(0)}\overline{A}_{2n2m},$$
(4)
where the sum is over all polygons $`Q`$ of perimeter $`2n2m.`$
Hence the number of ways of placing all of the $`p_{2m}^{(0)}`$ polygons of perimeter $`2m`$ inside any polygon of perimeter $`2n2m`$ is less than $`p_{2m}^{(0)}p_{2n2m}^{(0)}\overline{A}_{2n2m},`$ and so
$$p_{2n}^{(1)}<\underset{m=2}{\overset{n\sqrt{n}}{}}p_{2m}^{(0)}p_{2n2m}^{(0)}\overline{A}_{2n2m}<D_2C_2^2\mu ^{2n}2^5S_{2n},$$
(5)
where
$$S_{2n}=\underset{m=2}{\overset{n\sqrt{n}}{}}m^{\frac{5}{2}}/(nm).$$
The last sum may be evaluated in a variety of ways. Using Maple we find that
$$S_{2n}=0.341../n+O(n^{\frac{3}{2}}).$$
Thus we obtain the bound
$$p_{2n}^{(1)}<const.\mu ^{2n}/n.$$
(6)
For staircase polygons the analogous calculation is slightly simpler, and we find
$$S_{2n}1+\zeta (3/2)+O(n^{\frac{1}{2}}),$$
where $`\zeta (z)`$ is the Riemann zeta function.
To obtain a lower bound, we restrict the inner polygon to be of perimeter 4, that is, a unit square. The number of polygons punctured by a unit square clearly provides a lower bound to the number of 1-punctured polygons. A unit square can be placed anywhere in a polygon $`Q`$ of perimeter $`2n4`$ except on a boundary site. There are $`2n4`$ boundary sites. The mean area of $`Q`$ is greater than $`D_1(n2)^{\frac{3}{2}},`$ so we obtain the bound
$$p_{2n}^{(1)}>p_{2n4}(D_1(n2)^{\frac{3}{2}}2n+4)>C_1\frac{\mu ^{2n4}}{(2n4)^{\frac{5}{2}}}(D_1(n2)^{\frac{3}{2}}2n)>const.\mu ^{2n}/n.$$
(7)
Combining the two bounds gives the result
$$E_1\mu ^{2n}/np_{2n}^{(1)}E_2\mu ^{2n}/n,$$
where $`E_1<E_2`$ are constants. Accepting the usual asymptotic form that is expected for such models, we conclude that
$$p_{2n}^{(1)}const.\mu ^{2n}/n.$$
(8)
(For staircase polygons the analogous result is $`p_{2n}^{(1)}const.4^n.`$) Since the number of unpunctured polygons grows like $`\mu ^nn^{5/2},`$ we see that the exponent is predicted to increase by $`\frac{3}{2}`$ as the result of a single puncture, while in the case of the area generating function, the exponent is found to increase only by 1. We show in the next section that this prediction is borne out by our numerical calculations.
For punctured staircase polygons, a similar conclusion holds. That is, the connective constant is unaltered at $`\mu _{stair}=2,`$ but the exponent increases by $`1.5`$ over its unpunctured counterpart when enumerating 1-punctured staircase polygons by perimeter. In fact, we have been able to calculate the generating function for staircase polygons with a single puncture of perimeter 4, and also with a single puncture of perimeter 6. The generating functions for these special cases are given below, and are precisely in accordance with the more general results given above.
We also note that Cardy recently considered the problem of the number of punctured SAP with $`k`$ concentric, mutually self-avoiding SAPs surrounding a fixed point of the dual lattice. Thus for $`k=1`$ this corresponds to 1-punctured SAPs surrounding a fixed point. In that case Cardy finds for the number of such configurations $`b_{2n}^{(1)}`$ that $`b_{2n}^{(1)}=\mu ^{2n}\frac{\mathrm{ln}n}{64\pi ^2n}.`$ The above calculation may be repeated for Cardy’s problem. All that is required is to add a factor $`\overline{A}_m`$ to the summand in eqn. 5, as the surrounded point can be anywhere inside the inner polygon. Thus we must multiply by the mean area of that polygon. The sum defining $`S_{2n}`$ then becomes $`S_{2n}=_{m=2}^{n\sqrt{n}}\frac{1}{m(nm)}\mathrm{log}n/n,`$ in agreement with Cardy’s result. (Our lower bound for this problem is too weak, but with more effort could be improved.)
## 4 Analysis of the series
All the series we have investigated are characterised by coefficients that grow exponentially, with sub-dominant term given by a critical exponent. The generic generating function behaviour is $`G(z)=_ng_nz^nA(z)(1\sigma z)^\xi ,`$ and hence the coefficients of the generating function $`g_n=[z^n]G(z)A(1/\sigma )/\mathrm{\Gamma }(\xi )\sigma ^nn^{\xi 1}.`$ Generally speaking the existence of the growth constant $`\sigma `$ has been proved, but except for exactly solvable models, such as staircase polygons, the existence of the critical exponent $`\xi `$ has only been conjectured, though its existence has never been doubted. The radius of convergence of the generating function is usually given by the critical point, which is at $`z=1/\sigma .`$
We principally used two methods to analyse all the series studied in this paper. Firstly, to obtain the singularity structure of the generating function we used the numerical method of differential approximants . In particular, we used this method to estimate the growth constant $`\sigma `$ and the critical exponent $`\xi .`$ We were invariably able to conjecture an exact value for $`\xi ,`$ which was always an integer or half-integer for all the problems we investigated. Imposing this conjectured exponent permitted a refinement of the estimate of the growth constant — providing so-called biased estimates.
Once the exact value of the exponent was conjectured, and the growth constant accurately estimated, we turned our attention to the “fine structure” of the asymptotic form of the coefficients, by fitting the coefficients to the assumed form $`g_n=[z^n]G(z)\sigma ^nn^{\xi 1}_{i0}c_i/n^{f(i)}.`$ If there is no non-analytic correction term, then $`f(i)=i,`$ while a square-root correction term means $`f(i)=i/2.`$ For all the series studied, only these two situations were encountered.
In all cases, our procedure is to assume a particular form for $`f(i),`$ and see how it fits the data. With the very long series we now have at our disposal, it is usually easy to see if the wrong assumption has been made — the sequence of amplitude estimates $`c_i`$ either diverges to infinity or converges to zero. Once the correct assumption is made, convergence is usually rapid and obvious. A detailed demonstration of the method can be found in .
As an example of the sort of results we obtained, we show sequences of estimates of the coefficients of the perimeter generating function of unpunctured SAP in Table 2. In that case we conjectured that $`f(i)=i.`$ Because of the large amount of tabular data generated by the method, we have not given this level of detail for the many series investigated here. We show only the results for two series. For the others, we just give our assessment of the apparent convergence of the sequences $`c_i,`$ and the estimated value of the limits. As the equations involved are linear, the method is easy to implement, and interested readers can readily generate the relevant data themselves. Some subtleties nevertheless exist. For example, for punctured staircase polygons, the perimeter generating function has two singularities on the circle of convergence, and so both must be taken into account. We discuss this in more detail in the relevant section below.
As for the first stage of the analysis, the method of differential approximants, we proceeded as follows: Estimates of the critical point and critical exponent were obtained by averaging values obtained from first order $`[L/N;M]`$ and second order $`[L/N;M;K]`$ inhomogeneous differential approximants. For each order $`L`$ of the inhomogeneous polynomial we averaged over those approximants to the series which used at least the first 80% - 90% of the terms of the series, and used approximants such that the difference between $`N`$, $`M`$, and $`K`$ didn’t exceed 2. These are therefore “diagonal” approximants. Some approximants were excluded from the averages because the estimates were obviously spurious. The error quoted for these estimates reflects the spread (basically one standard deviation) among the approximants. Note that these error bounds should not be viewed as a measure of the true error as they cannot include possible systematic sources of error. However systematic error can also be taken into account in favourable situations, as for example in the case of SAP enumerated by perimeter . Again, in the interests of space, we present only our results, and not the intermediate detail from which our estimates were made. An example in full detail for one of the series investigated in this study can be found in . We turn now to the analysis of all the series.
### 4.1 Staircase polygons
For (unpunctured) staircase polygons, the multi-variable width, height and area generating function is known . As usual, we denote
$$(a)_n=\underset{i=0}{\overset{n1}{}}(1aq^i).$$
Further, denoting the first two $`q`$-Bessel functions as:
$$J_0(x,y,q)=\underset{n0}{}\frac{(1)^nx^nq^{\left(\genfrac{}{}{0pt}{}{n+1}{2}\right)}}{(q)_n(yq)_n},$$
and
$$J_1(x,y,q)=\underset{n1}{}\frac{(1)^{n1}x^nq^{\left(\genfrac{}{}{0pt}{}{n+1}{2}\right)}}{(q)_{n1}(yq)_{n1}(1yq^n)},$$
the perimeter and area generating function is simply
$$P(x,y,q)=y\frac{J_1(x,y,q)}{J_0(x,y,q)},$$
where $`x`$ $`(y)`$ is the variable conjugate to the horizontal (vertical) semi-perimeter, while $`q`$ counts the area. No analogous results are known for punctured staircase polygons, though the calculation for SAP by area can be carried over mutatis mutandis to prove that the area generating function for a punctured staircase polygon with $`k`$ holes has the same radius of convergence as the area generating function for unpunctured staircase polygons. Further, the critical exponent increases by 1 for each puncture, and unlike the case for SAP, we can not only prove the existence of a critical exponent for unpunctured staircase polygons, but we know its value. Hence we know the leading term in the asymptotic expansion of the generating function by area for $`k`$-punctured staircase polygons.
For the expected behaviour of the perimeter generating function of punctured staircase polygons, the arguments of the preceding section apply directly. The radius of convergence, and hence the connective constant remains unchanged, and the argument given in the preceding section suggests that the critical exponent should increase by 1.5 for each puncture. The results of our analysis, presented below, bear this out.
#### 4.1.1 Area generating function
For unpunctured staircase polygons, the area generating function is given by
$$A(q)=\underset{n1}{}a_n^{(0)}q^n=\frac{J_1(1,1,q)}{J_0(1,1,q)}.$$
By inspection, this has poles at the zeros of $`J_0(1,1,q).`$ The nearest zero is at $`1/q=\eta =2.30913859330,`$ and there is a simple pole at that point. The next zero is well separated (at $`1/q=\lambda =1.4435..`$) and so the asymptotic form of the generating function is
$$A(q)D/(1\eta q)+E/(1\lambda q)+\mathrm{},$$
and hence
$$a_n^{(0)}=[q^n]A(q)\eta ^n(c_0+O((\lambda /\eta ))^n).$$
Our analysis bears this out, and we estimate $`c_0=0.12881579.`$
For 1-punctured discs, our analysis, based on more than 100 series coefficients, convincingly suggests the following asymptotic form:
$$A^{(1)}(q)D^{(1)}/(1\eta q)^2+E^{(1)}/(1\eta q)^{1.5}+F^{(1)}/(1\eta q)+\mathrm{},$$
and hence
$$a_n^{(1)}=[q^n]A^{(1)}(q)\eta ^nn\underset{i0}{}c_i/n^{i/2}.$$
The sequences of amplitude estimates, assuming this asymptotic form, are shown in Table 3. The apparent convergence of the amplitude estimates is, as explained above, our source of evidence for this asymptotic form.
For 2-punctured discs, a similar analysis, based on some 86 series coefficients, convincingly suggests the following asymptotic form:
$$A^{(2)}(q)D^{(2)}/(1\eta q)^3+E^{(2)}/(1\eta q)^{2.5}+\mathrm{},$$
and hence
$$a_n^{(2)}=[q^n]A^{(2)}(q)\eta ^nn^2\underset{i0}{}a_i/n^{i/2}.$$
For 3-punctured discs our analysis was based on an 89 term series. We found that the above pattern persists, so that the generating function has the following asymptotic form:
$$A^{(3)}(q)D^{(3)}/(1\eta q)^4+E^{(3)}/(1\eta q)^{3.5}+\mathrm{},$$
and hence
$$a_n^{(3)}=[q^n]A^{(3)}(q)\eta ^nn^3\underset{i0}{}a_i/n^{i/2}.$$
Estimates of the various amplitudes defined above are shown in Table 4.
#### 4.1.2 Perimeter generating function
For unpunctured staircase polygons, the perimeter generating function (ignoring the distinction between height and width) is given by
$$P(x)=\frac{12x\sqrt{14x}}{2},$$
which is, apart from suppression of the first term, the generating function for Catalan numbers. Hence
$$p_{2n}^{(0)}=[x^n]P(x)=\frac{1}{n}\left(\genfrac{}{}{0pt}{}{2n2}{n1}\right)4^n/n^{\frac{3}{2}}\underset{i0}{}c_i/n^i.$$
The values of $`c_i`$ follow immediately from the exact solution. They are given in Table 4.
For 1-punctured discs, our analysis, again based on more than 100 series coefficients, was more equivocal than that of the 1-punctured area generating function. The method of differential approximants clearly identified a singularity at the known critical point, $`x_c=1/4,`$ but almost all approximants had a double root, implying a confluent singularity. The leading exponent was estimated to be $`1,`$ implying a pole in the generating function, (and hence immediately lending support to our conjectured change in the critical exponent of 3/2 as a result of puncturing), but we were unable, from this method, to identify the confluent exponent. Further, a second singularity was identified, of the form const.$`(1+x/x_c)^{6.5}.`$ At this point we wish to remark on the close similarity between a recently solved model of polygons, the three-choice polygon model and 1-punctured staircase polygons. 1-punctured staircase polygons can be thought of as being constructed from two three-choice polygons, with common edges deleted. This geometric similarity is borne out by the fact that the two models have identical singularity distributions, with even the exponents being the same at all non-physical singularities.
For 1-punctured discs then, the perimeter generating function is expected to be of the following asymptotic form:
$$P^{(1)}(x)=\underset{n}{}p_{2n}^{(1)}x^nB^{(1)}(x)+C^{(1)}(x)(14x)^1+D^{(1)}(x)(1+4x)^{6.5},$$
and hence
$$p_n^{(1)}=[x^n]P^{(1)}(x)4^n\underset{i0}{}c_i/n^{f(i)}+(4)^nn^{7.5}\underset{i0}{}d_i/n^i,$$
where $`B^{(1)}(x),`$ $`C^{(1)}(x)`$ and $`D^{(1)}(x)`$ are assumed regular in the disc $`|4x|1.`$
Assuming $`f(i)=i,`$ which implies only analytic correction-to-scaling terms, gave unsatisfactory results. Notably, we observed that the estimates of the amplitude $`c_1`$ in the above asymptotic form were steadily increasing, suggesting that the assumed form did not properly account for the correction-to-scaling terms. With $`f(i)=i/2`$ the amplitude estimates were much more stable. This then implies that the generating function in fact behaves as
$$P^{(1)}(x)B^{(1)}(x)+C^{(1)}(x)(14x)^1+E^{(1)}(x)(14x)^{\frac{1}{2}}+D^{(1)}(x)(1+4x)^{6.5}.$$
(9)
The amplitude estimates, assuming this asymptotic form, are shown in Table 4. The apparent convergence of the amplitude estimates is, as explained above, our source of evidence for this asymptotic form. We did not tabulate the non-physical amplitudes, as they are of little interest to our investigation. However they need to be included to stabilise estimates of the physical amplitudes. We mention in passing that $`d_00.14.`$
We also generated series for staircase polygons with two and three punctures, with perimeter 150 and 134 steps respectively. Our differential approximant analysis lent support to our expectation that the critical exponent increases by $`1.5`$ per puncture. This turned out to be true for the non-physical singularity at $`x=\frac{1}{4}`$ as well. A similar analysis to that for the one-punctured polygons strongly supported the analogous asymptotic forms,
$$P^{(k)}(x)=\underset{n}{}p_{2n}^{(k)}x^nB^{(k)}(x)+C^{(k)}(x)(14x)^{0.51.5k}+D^{(1)}(x)(1+4x)^{81.5k}$$
(10)
for $`k>0.`$ Hence the asymptotic form of the coefficients is conjectured to be
$$p_{2n}^{(k)}4^nn^{\frac{3(k1)}{2}}\underset{i0}{}c_i^{(k)}/n^{i/2}+(4)^nn^{1.5k9}\underset{i0}{}d_i^{(k)}/n^i,$$
where the amplitude estimates $`c_i^{(k)}`$ are given in Table 4.
We have also determined the exact generating function for staircase polygons punctured by a single hole of perimeter 4, and also the generating function for staircase polygons punctured by a single hole of perimeter 6. We obtained these generating functions by generating the coefficients using the algorithm discussed, and then searching for an underlying differential equation. As a result we find the perimeter generating functions $`P_4^s(x)`$ and $`P_6^s(x)`$ for 1-punctured staircase polygons with a hole of perimeter 4 and 6 respectively.
$$P_4^s(x)=\frac{2x^416x^3+20x^28x+1}{2(14x)}\frac{16x+10x^24x^3}{2\sqrt{14x}}$$
(11)
and
$`P_6^s(x)`$ $`=`$ $`{\displaystyle \frac{126x+228x^2906x^3+1709x^41378x^5+322x^6}{2(14x)^{\frac{5}{2}}}}`$ (12)
$`{\displaystyle \frac{32x^6404x^5+815x^4586x^3+182x^224x+1}{2(14x)^2}}.`$
Note that both these exact solutions display the confluent square root correction that we have found in our numerical investigations in the more general case. In the next subsection we analyse the analogous generating function for SAP.
### 4.2 Self-avoiding polygons
For unpunctured SAP, the perimeter generating function was recently extended to 90 step polygons, and the asymptotics clearly identified. The polygon generating function is defined to be
$$P^{(0)}(x)=\underset{n}{}p_{2n}^{(0)}x^nA^{(0)}(x)+B^{(0)}(x)(1\mu ^2x)^{\frac{3}{2}},$$
(13)
where the functions $`A^{(0)}`$ and $`B^{(0)}`$ are believed to be regular in the vicinity of $`x_c=1/\mu ^2.`$ We estimated $`\mu =2.63815853034(10)`$. From this equation follows the asymptotic form of the coefficients,
$$p_{2n}^{(0)}\mu ^{2n}n^{\frac{5}{2}}[c_1+c_2/n+c_3/n^2+c_4/n^3+\mathrm{}].$$
(14)
We show in Table 2 the sequence of estimates of $`c_i,`$ and in Table 4 our estimates of the amplitudes, being the limits of the sequences $`\{c_i\}.`$
The area generating function was first studied in , where the first 20 terms were given, and the asymptotic form estimated to be
$$A^{(0)}(q)=\underset{n}{}a_nq^nD(q)+E(q)\mathrm{log}(1\kappa q),$$
(15)
where $`\kappa 3.97087,`$ and the logarithm in the above equation was understood to include the possibility of some power of a logarithm other than unity. (Though our analysis below implies that this is not the case.) In the present work we extend the series to 42 terms.
#### 4.2.1 Area generating function
Using our greatly extended 42 term series, our analysis of the unpunctured disc area generating function was carried out by standard methods. We used differential approximants and found unbiased critical point and critical exponent estimates. The unbiased exponent estimate had absolute value less than $`10^5,`$ totally supporting our view that it is exactly zero. Assuming this, a biased estimate of the critical point is possible, and in this way we estimate $`\kappa =3.97094397(9).`$
We then proceeded to seek the asymptotic form of the coefficients by writing
$$a_n=[q^n]A^{(0)}(q)\kappa ^n/n\underset{i0}{}c_i/n^{f(i)}.$$
Our numerical results were well converged, demonstrating very convincingly that $`f(i)=i.`$ This is the asymptotic form consistent with a pure logarithmic singularity, not raised to any power. Estimates of the amplitudes are given in Table 4.
For 1-punctured discs, our analysis, based on the series known to $`q^{45}`$, convincingly suggests the following asymptotic form:
$$A^{(1)}(q)D^{(1)}(q)+E^{(1)}(q)/(1\kappa q),$$
and hence
$$a_n^{(1)}=[q^n]A^{(1)}(q)\kappa ^n\underset{i0}{}c_i/n^i.$$
Attempts to fit to alternative forms, corresponding to a confluent logarithm or a non-analytic correction-to-scaling term were unsuccessful, adding to our confidence that the above form is correct. Estimates of the amplitudes are given in Table 4.
For twice-punctured discs, our analysis, also based on the series known to $`q^{48}`$ convincingly suggests the following asymptotic form:
$$A^{(2)}(q)D^{(2)}(q)+E^{(2)}(q)/(1\kappa q)^2,$$
and hence
$$a_n^{(2)}=[q^n]A^{(2)}(q)\kappa ^nn\underset{i0}{}c_i/n^i.$$
The conjectured asymptotic form for $`k`$ punctured polygons, by area, is thus
$$A^{(k)}(q)D^{(k)}(q)+E^{(k)}(q)/(1\kappa q)^k.$$
#### 4.2.2 Perimeter generating function
The results for unpunctured polygons are fully discussed at the beginning of subsection 4.2. As we found with staircase polygons, the generating function for punctured discs by perimeter was a more challenging numerical analysis problem than either its unpunctured counterpart, or its area counterpart.
We found that the method of differential approximants was not particularly satisfactory. Given that we needed some 100 terms to successfully analyse the (presumably simpler) problem of punctured staircase polygons by perimeter, it is not surprising that for punctured SAP, for which we have 33 non-zero coefficients (corresponding to perimeters up to 84 steps), the method was not satisfactory. However, it did indicate the presence of a confluent singularity. As we found a confluent square-root singularity for staircase polygons, it is hardly surprising that a confluent singularity is detected for the punctured SAP generating function. In fact an exponent shift of around 1.7 was seen, compared to the expected value 1.5. We attribute this to the “short” series, coupled with the well known deleterious effect of confluent terms in such an analysis. Nevertheless, subsequent analysis of the asymptotic form of the coefficients, assuming an exponent shift of 1.5, together with a square-root confluent term, as found for punctured staircase polygons, gave satisfactory results.
We denote the generating function for $`k`$-punctured SAP, by perimeter, as
$$P^{(k)}(x)=\underset{n}{}p_{2n}^{(k)}x^nB^{(k)}(x)+C^{(k)}(x)(1\mu ^2x)^{1.51.5k},$$
(16)
where the exponent is conjectured. For 1-punctured SAP, the vanishing of the exponent implies a logarithmic singularity. We fitted the coefficients to the asymptotic form appropriate to $`\mathrm{log}(1\mu ^2x),`$ so that the asymptotic form of the coefficients just involves decreasing integer powers of $`n.`$ We then found that the estimates of the leading amplitude were monotonically increasing, which implies that the asymptotic form is wrong — too weak. Including a confluent square root singularity, as was found for punctured staircase polygons, stabilised the estimates. Accordingly, we conjecture that the asymptotic form is dominated by a logarithmic singularity, with a sub-dominant square root singularity, so that
$$p_{2n}^{(1)}=[x^n]P^{(1)}(x)\mu ^{2n}/n\underset{i0}{}c_i/n^{\frac{i}{2}}$$
(17)
Estimates of the amplitudes $`c_i`$ are given in Table 4.
For twice punctured discs, a similar analysis suggested that the asymptotic form of the generating function is
$$P^{(2)}(x)=\underset{n}{}p_{2n}^{(2)}x^nB^{(2)}(x)+C^{(2)}(x)(1\mu ^2x)^{1.5},$$
(18)
again with evidence of a square root confluent term. As for 1-punctured SAP, we give estimates of the amplitudes $`c_i`$ defined by
$$p_{2n}^{(2)}\mu ^{2n}n^{\frac{1}{2}}\underset{i0}{}c_i/n^{\frac{i}{2}}$$
in Table 4.
For thrice punctured discs, a similar analysis suggested that the asymptotic form of the generating function is
$$P^{(3)}(x)=\underset{n}{}p_{2n}^{(3)}x^nB^{(3)}(x)+C^{(3)}(x)(1\mu ^2x)^3,$$
(19)
again with evidence of a square root confluent term. As for 1-punctured SAP, we give estimates of the amplitudes $`c_i`$ defined by
$$p_{2n}^{(3)}\mu ^{2n}n^2\underset{i0}{}c_i/n^{\frac{i}{2}}$$
in Table 4.
The conjectured asymptotic form for $`k`$ punctured polygons, by perimeter, is then
$$P^{(k)}(x)=\underset{n}{}p_{2n}^{(k)}x^nB^{(k)}(x)+C^{(k)}(x)(1\mu ^2x)^{1.51.5k}$$
(20)
where for $`k>0`$ we find strong evidence for a square root correction-to-scaling term.
### 4.3 Polyominoes
The problem of polyominoes has a long and interesting history, and has been well discussed in the popular scientific literature .
The enumeration of square lattice polyominoes to 24 steps was given in 1981, extended to 25 steps in 1995 and currently stands at 28 steps .
We have analysed the latest series by the method of differential approximants, and find the generating function behaves as
$$𝒫(y)=\underset{n}{}a_ny^nG(y)+H(y)\mathrm{log}(1\tau y),$$
(21)
where $`a_n`$ is the number of polyominoes of area $`n`$. In the estimate $`\tau =4.06265(5)`$ was given. The extra terms now available allow us to make the refined estimate $`\tau =4.062591(9).`$ Analysis of the asymptotic form of the coefficients is totally consistent with a simple logarithm in the generating function. Thus
$$a_n=[y^n]𝒫(y)\tau ^n\underset{i0}{}c_i/n^{i+1}.$$
The amplitude estimates are given in Table 4. The leading amplitude is in complete agreement with, but 3 orders of magnitude more accurate than that given in , while the order of the leading term — $`O(1/n)`$ — was predicted by physical arguments back in 1981 . In that work, the logarithmic singularity in the generating function of strongly embedded site-animals was obtained. As discussed in the introduction, these are isomorphic to polyominoes.
As well as extending the polyomino series, Oliveira e Silva has enumerated $`k`$-punctured polyominoes for $`k6.`$ Clearly, $`k=0`$ polyominoes are just SAP, and as we have seen in the previous section, these grow as $`\kappa ^n`$ where $`\kappa =3.9709..<\tau .`$ The arguments in can be modified and applied to show that a $`k`$-punctured polyomino has the same growth constant $`\lambda `$ as its unpunctured counterpart, and this is done in the Appendix.
Thus the situation is that, for any finite number of punctures,
$$a_n^{(k)}=[y^n]𝒫^{(k)}(y)\kappa ^n,$$
but that
$$a_n=[y^n]𝒫(y)=\underset{k0}{}a_n^{(k)}\tau ^n.$$
We have analysed the series for $`k`$-punctured polyominoes, $`k=0,1,2,`$ and find the asymptotic form of the coefficients to be
$$a_n^{(k)}=[y^n]𝒫^{(k)}(y)\kappa ^nn^{k1}\underset{i0}{}c_i^{(k)}/n^i,$$
corresponding to a generating function for $`k`$-punctured polyominoes having a $`k^{th}`$ order pole, viz:
$$𝒫^{(k)}(y)=\underset{n}{}a_n^{(k)}y^nH^{(k)}(y)(1\kappa y)^k,$$
(22)
where $`k=0`$ is to be interpreted as a logarithm. Thus just as for punctured polygons, it is found that the exponent of $`k`$-punctured polyominoes increases by 1 for each puncture. This is also proved in the Appendix.
The coefficients $`c_i^{(k)}`$ are given in Table 4.
## 5 Conclusion
We have investigated the effect of punctures on SAP and staircase polygons enumerated both by area and perimeter.
In order to do this, we have developed a new algorithm, exponentially faster than direct counting, whereby we have radically extended a number of series. This extension was necessary in order to probe some rather subtle numerical behaviour.
We found that, in every case, a finite number of punctures does not change the exponential growth factor associated with the unpunctured counterpart of the punctured object being enumerated.
This latter conclusion was also proved for polyominoes. The effect of punctures was also investigated numerically for polyominoes.
Writing $`b_n^{(k)}`$ for the $`n^{th}`$ coefficient in the generating function for some $`k`$punctured object, so that
$$b_n^{(k)}(\omega ^{(k)})^nn^{\theta (k)},$$
we found $`\omega ^{(k)}=\omega ^{(0)}`$ in all cases. We found, further, that $`\theta (k)=\theta (0)+k`$ if enumerating any of the objects we have considered by area. This can be proved, though subject, in some cases, to the existence of the exponent in question. Subject to certain unproved assumptions we also showed that $`\theta (k)=\theta (0)+3k/2`$ if enumerating by perimeter.
We have, for the first time, obtained good numerical estimates of the sub-dominant terms for a range of problems, thus identifying both the nature of the generating function and any correction-to-scaling terms.
We have also obtained an exact solution for the generating functions of staircase polygons, enumerated by perimeter, punctured by a single hole of perimeter 4 and of perimeter 6. These exact solutions provide additional support for our numerically based conjectures of the correction-to-scaling exponent in the general case.
A more accurate estimate of the growth constant for SAP enumerated by area has been given, complementing our earlier work on the perimeter growth constants .
## E-mail or WWW retrieval of series
The series for the various generating functions studied in this paper can be obtained via e-mail by sending a request to I.Jensen@ms.unimelb.edu.au or via the world wide web on the URL http://www.ms.unimelb.edu.au/~iwan/ by following the instructions.
## Acknowledgements
We have derived great benefit from discussions of aspects of this problem with Mireille Bousquet-Mélou, John Cardy, Aleks Owczarek, Buks van Rensburg, and Stuart Whittington. We are particularly grateful to Mireille Bousquet-Mélou for her careful reading of the manuscript, to Buks van Rensburg for clarifying aspects of his earlier work, and to John Cardy for querying some earlier, incorrect results. IJ and AJG gratefully acknowledge financial support from the Australian Research Council. LHW would like to thank The University of Melbourne for the Melbourne Research Scholarship.
## Appendix. Growth constants and exponents for $`k`$-punctured polyominoes
In this appendix we show that the growth constants are the same for all $`k`$-punctured polyominoes. Our method of proof is based on that of van Rensburg and Whittington , making the necessary changes for the polyomino problem, and discussing in detail certain special cases. That is to say, if $`\kappa `$ denotes the connective constant for SAP enumerated by area, then this is the growth constant for $`k`$-punctured polyominoes for any finite $`k`$. Further, if the usual asymptotic form for the number of $`k`$-punctured polyominoes is assumed, $`s_n^{(k)}C_kn^{\varphi _k}\kappa ^n`$, and $`\varphi _0`$ exists, then $`\varphi _k=\varphi _0k.`$
### A.1 Operations and mappings on punctured polyominoes
Let the set of all $`k`$-punctured polyominoes with $`n`$ cells be denoted by $`\mathrm{\Phi }_n^{(k)}`$, and the set of all polyominoes of $`n`$ cells be $`\mathrm{\Phi }_n`$. Then, for each $`n`$, $`\mathrm{\Phi }_n=_k\mathrm{\Phi }_n^{(k)}`$. Throughout, let $`s_n^{(k)}`$ denote the cardinality of $`\mathrm{\Phi }_n^{(k)}`$ and let $`s_n`$ denote the cardinality of $`\mathrm{\Phi }_n`$. Hence $`s_n=_ks_n^{(k)}`$.
Following we now define operations on punctured polyominoes, including concatenation, drilling and surgery, which are needed in subsequent proofs. Concatenation allows us to change the size of polyominoes, while drilling and surgery are concerned with changing the number of holes in a polyomino.
#### A.1.1 Concatenation
The concatenation mapping defined here is similar to that in . Consider the bounding rectangle $`R(P)`$ of any polyomino $`P\mathrm{\Phi }_n^{(k)}`$. Define the top (bottom) edge of $`P`$ to be the top (bottom) edge along the east (west) side of $`R(P)`$.
Now, the concatenation of two polyominoes $`P\mathrm{\Phi }_n^{(h)}`$ and $`Q\mathrm{\Phi }_m^{(k)}`$ is defined by joining $`P`$ and $`Q`$ while superimposing the top edge of $`P`$ and the bottom edge of $`Q`$. The result is an $`(h+k)`$-punctured polyomino with $`m+n`$ cells, see figure A1. Hence we have a map
$$T:\mathrm{\Phi }_n^{(h)}\times \mathrm{\Phi }_m^{(k)}\mathrm{\Phi }_{m+n}^{(h+k)}$$
(A1)
###### Lemma A1
For all non-negative values of $`h`$ and $`k`$, $`s_n^{(h)}s_m^{(k)}s_{m+n}^{(h+k)}`$ and in particular, $`s_{n+1}^{(h)}s_n^{(h)}`$
Proof: For the mapping $`T`$ defined above, every pair of polyominoes in the domain can be concatenated to form a larger polyomino in the codomain. Conversely, every such polyomino in the codomain can be uniquely broken up into the original ones. However, there are some polyominoes in the codomain which cannot be formed by concatenating two smaller polyominoes. An example is a $`2\times 2`$ polyomino. Hence we get the first part of lemma as
$$|\mathrm{\Phi }_n^{(h)}\times \mathrm{\Phi }_m^{(k)}||\mathrm{\Phi }_{m+n}^{(h+k)}|.$$
Putting $`k=0`$ and $`m=1`$ and noting that $`|\mathrm{\Phi }_1^{(0)}|=1`$, we get the second part. $`\mathrm{}`$
#### A.1.2 Drilling
Simply creating a hole inside an unpunctured polyomino by removing some interior cells does not allow us to drill certain classes of polyominoes. For instance, a polyomino composed of a single linear sequence of cells cannot be ‘drilled’ since removing any cell will either disconnect the polyomino or shorten the sequence.
The following definition of drilling differs somewhat from that in , though the underlying idea is the same.
First, we drill one hole. Let $`P_0\mathrm{\Phi }_n^{(0)}`$ be the polyomino that we are to puncture. Cover $`P_0`$ by a grid system, with each grid square of size $`b\times b`$ cells for any $`b5`$ (the minimum size will be justified later). Say $`b=5.`$ Pick any grid square $`G`$ that covers at least one cell of $`P_0`$ and drill a hole there, as detailed below.
Step 1 Remove all cells within the selected grid square $`G`$. If, after this step, we are left with a $`1`$-punctured polyomino, we are done. If not, go to step 2.
Step 2 Check the corners of $`G`$: If there are 2 disconnected components touching each other only at the corner, connect them by adding a cell at the appropriate corner of $`G`$.
Step 3 Put a $`1`$-punctured polyomino with 8 cells at the centre of $`G`$.
Step 4 Reconnect the disconnected components outside $`G`$ to the $`1`$-punctured polyomino by adding linear sequences of cells (non-unique). These steps are illustrated in figure A2.
Note that the minimum value of $`b`$ is 5 because if we have anything less than that, we might create extra holes unexpectedly as the $`1`$-punctured polyomino in Step 3 must touch the boundary of the grid square. An example in figure A3 will illustrate this. In this example the original polyomino has no holes and the drilled polyomino has two holes, whereas with a grid of size 5 (or more), the number of holes in the drilled polyomino is only one.
After this operation, the maximum number of cells that could be removed is $`b^2=25`$ (finishing at step 1). The maximum number of cells that could be added is $`b^22`$ (this occurs when there was only 1 cell in the square before drilling and we end up with a 1-punctured polyomino with $`b^21`$ cells after the operation). So, depending on each instance of the drilling operation, we obtain a resulting polyomino $`P_1\mathrm{\Phi }_j^{(1)}`$ where $`j`$ could be anything from $`nb^2`$ to $`n+(b^22)`$. The drilling operation thus defines a map
$$D:\mathrm{\Phi }_n^{(0)}\underset{j=nb^2}{\overset{n+(b^22)}{}}\mathrm{\Phi }_j^{(1)}.$$
(A2)
###### Theorem A1
There exists a real constant $`C`$ such that for all $`b5`$,
$$s_n^{(0)}Cs_{n+(b^22)}^{(1)}.$$
(A3)
Proof: Consider the intermediate (possibly disconnected) polyomino after Step 1, call it $`P^i`$. There are many possible initial polyominoes which gives the same $`P^i`$, i.e., all polyominoes with the same configuration outside $`G`$. They all map to “almost” the same resulting polyomino (the uncertainty implied in “almost” comes from the non-uniqueness of reconnecting in Step 4 which will be discussed later). Therefore, the mapping from domain to codomain is $`M^{}`$-to-one where $`M^{}`$ is bounded by the number of ways that at most $`b^2`$ cells can be connected within $`G`$. Let the bound be $`M`$.
On the other hand, Step 4 of the drilling process is not unique. There might be more than one way to connect those disconnected components to the punctured polyomino. But since $`G`$ is finite, the number of ways of reconnection is bounded above. Let that bound be $`C^{}`$.
Together with the drilling mapping $`D`$ defined above, we can write,
$$s_n^{(0)}C^{}M\underset{j=nb^2}{\overset{n+(b^22)}{}}s_j^{(1)}.$$
By the increasing monotonicity of $`s_n^{(h)}`$ over $`n`$ (proven in lemma (A1)),
$$s_n^{(0)}(2b^21)C^{}Ms_{n+(b^22)}^{(1)}.$$
Set $`C=(2b^21)C^{}M`$ and the result follows. $`\mathrm{}`$
Now consider the drilling of $`h`$ holes. Denote by $`c`$ the ceiling of $`c`$, the smallest integer greater than or equal to $`c`$. Now, we could choose $`h`$ drilling locations from at least $`n/b^2`$ grid squares. To see this, consider first a polyomino with $`n`$ cells, where $`nb^2`$. The situation where we have the least number of drilling sites is when $`n`$ cells fall exactly within 1 grid square. In this way, we only have 1 possible grid square where we could drill holes. Similarly, if we have $`b^2n+w`$ cells, where $`n1`$ and $`0w<b^2`$, the minimal number of grid squares is when $`b^2n`$ cells fall into exactly $`n`$ grid squares and the other $`w`$ cells falls into another single grid. Then we have $`n+1`$ possible grid squares to drill holes. Therefore we have at least $`n/b^2`$ grid squares where we could drill holes.
Letting $`H`$ be a set of grid squares, let $`\mathrm{\Phi }_n^{[H]}`$ denote the set of all polyominoes with $`n`$ cells and $`|H|`$ holes where there is a hole in each grid square of $`H`$. Let $`s_n^{[H]}`$ denote its cardinality. One property of this set is that, for any $`n_1,n_2Z^+`$ such that $`n_1n_2`$,
$$s_{n_1}^{[H]}s_{n_2}^{[H]}.$$
###### Theorem A2
There exists a real constant $`K`$ such that for all $`b5`$,
$$\left(\genfrac{}{}{0pt}{}{n/b^2}{h}\right)s_n^{(0)}Ks_{n+(b^22)h}^{(h)},bn/b^2.$$
(A4)
Proof: Place the polyomino in the grid system. Let $`A^{}`$ be the set of all grid squares where we could drill holes. From previous results, we know there are at least $`n/b^2`$ elements in $`A^{}`$. Truncate the set $`A^{}`$ with only the first $`n/b^2`$ elements and call this set of drilling sites $`A`$.
Pick a subset $`HA`$ such that $`|H|=h`$ $`(n/b^2)`$ and drill holes one by one in each grid square in $`H`$, leading to a series of mappings:
$$M:\mathrm{\Phi }_n^{(0)}\mathrm{\Phi }_{m_1}^{(1)}\mathrm{\Phi }_{m_2}^{(2)}\mathrm{}\mathrm{\Phi }_{m_h}^{[H]}$$
where $`m_i`$’s are appropriate constants depending on each instance of operation and $`m_in+(b^22)h`$, $`i`$. So, from this mapping $`M`$, we have
$$s_n^{(0)}C^hs_{m_h}^{[H]}C^hs_{n+(b^22)h}^{[H]}.$$
(A5)
Hence
$`{\displaystyle \underset{HA}{}}s_n^{(0)}`$ $``$ $`{\displaystyle \underset{HA}{}}C^hs_{n+(b^22)h}^{[H]}`$
$`\left({\displaystyle \genfrac{}{}{0pt}{}{n/b^2}{h}}\right)s_n^{(0)}`$ $``$ $`C^h{\displaystyle \underset{HA}{}}s_{n+(b^22)h}^{[H]}`$
$`\left({\displaystyle \genfrac{}{}{0pt}{}{n/b^2}{h}}\right)s_n^{(0)}`$ $``$ $`Ks_{n+(b^22)h}^{(h)}.`$
The last line arises since $`_{HA}\mathrm{\Phi }_n^{[H]}\mathrm{\Phi }_n^{(|H|)}\mathrm{\Phi }_n^{(h)}`$ and put $`K=C^h`$. $`\mathrm{}`$
#### A.1.3 Surgery
The surgery operation removes a linear sequence of cells inside the polyomino and concatenates the sequence to the external boundary of the polyomino. Our objective is to join two holes thereby reducing the number of holes by one.
We divide our domain into three classes:
(1) the set of polyominoes which have at least one hole with size one cell; denote this set by $`\dot{\mathrm{\Phi }}_n^{(h)}`$ with cardinality $`\dot{s}_n^{(h)},`$ see figure A4(a);
(2) the set of polyominoes which are not in (1), and have three holes touching each other at corners around a single cell; denote this set by $`\ddot{\mathrm{\Phi }}_n^{(h)}`$ with cardinality $`\ddot{s}_n^{(h)}`$, see figure A4(b);
(3) the set of polyominoes which are not in (1) nor (2) (the general case), denote this set by $`\stackrel{~}{\mathrm{\Phi }}_n^{(h)}=\mathrm{\Phi }_n^{(h)}\backslash (\dot{\mathrm{\Phi }}_n^{(h)}\ddot{\mathrm{\Phi }}_n^{(h)})`$ with cardinality $`\stackrel{~}{s}_n^{(h)}`$.
##### The general case:
First, let’s look at polyominoes in $`\stackrel{~}{\mathrm{\Phi }}_n^{(h)}`$. Consider a polyomino $`\alpha _n^{(h)}\stackrel{~}{\mathrm{\Phi }}_n^{(h)}`$. Let $`L`$ be the set of all loop-free sequences<sup>1</sup><sup>1</sup>1A loop-free sequence is a sequence of cells where each successive pair of cells are joined and no cell appears more than once in the whole sequence of cells in $`\alpha _n^{(h)}`$, one end of which must touch the boundary of one hole and the other end must touch the boundary of another hole (or the exterior boundary). Define the length of the sequence to be the number of cells in the sequence. Let the set $`Z`$ be the set of sequences in $`L`$ that has minimum length.
Now pick one sequence $`zZ`$ and cut it out. This step will join two holes together (or join one hole and the exterior together). Next, using the definition in section A.1.1, concatenate $`\alpha _n^{(h)}czDz`$ ($`D`$ is a $`3\times 3`$ polyomino block), and the resulting polyomino has one less hole and nine more cells than the original one. So we have a mapping
$$S_3:\stackrel{~}{\mathrm{\Phi }}_n^{(h)}\mathrm{\Phi }_{n+9}^{(h1)}$$
This mapping is at most $`n`$-to-one because for any polyomino in the codomain, we could find at most $`n`$ locations to connect $`z`$ back. As a result, we find that,
$$\stackrel{~}{s}_n^{(h)}ns_{n+9}^{(h1)}.$$
(A6)
##### Special cases:
There are some subclasses of case (1) where we cannot apply the general surgery operation. One example is shown in figure A5 . In this example, the minimum length of the loop-free sequence connecting two boundaries is two, but when we try to remove any such minimum sequence, we will disconnect the polyomino.
To deal with case (1) polyominoes, we simply fill such a hole with a single cell. This satisfies our objective of reducing the number of holes by one. So we find a mapping
$$S_1:\dot{\mathrm{\Phi }}_n^{(h)}\mathrm{\Phi }_{n+1}^{(h1)}.$$
The mapping $`S_1`$ is $`n^{}`$ to 1 where $`n^{}`$ is less than $`n`$. To see this, consider a polyomino in the codomain. To map back to the domain, we can choose any interior cell to remove and the number of choices is obviously less than the total number of cells in the polyomino, so $`n^{}n`$. Therefore,
$$\dot{s}_n^{(h)}ns_n^{(h1)}.$$
(A7)
Finally, consider polyominoes in $`\ddot{\mathrm{\Phi }}_n^{(h)}`$. The problem with these polyominoes is that when we try to remove the cell between the 3 holes, we will inevitably join 3 holes together instead of joining 2. One way to get around this is to deliberately choose another sequence (which is also a cell in this class) in the set $`Z`$. In particular, we choose to remove the cell $`z^{}Z`$ such that $`z^{}`$ is the top cell (in lexicographic ordering) of $`Z`$ which is not the problem cell. Similarly to the general case, we find a mapping
$$S_2:\ddot{\mathrm{\Phi }}_n^{(h)}\mathrm{\Phi }_{n+9}^{(h1)}$$
and hence
$$\ddot{s}_n^{(h)}ns_{n+9}^{(h1)}.$$
(A8)
Adding (A6), (A7), (A8), we have
$$s_n^{(h)}=\dot{s}_n^{(h)}+\ddot{s}_n^{(h)}+\stackrel{~}{s}_n^{(h)}ns_{n+9}^{(h1)}+ns_{n+9}^{(h1)}+ns_n^{(h1)}$$
and hence:
###### Theorem A3
The number of polyominoes of area $`n`$ with $`h`$ holes is bounded above by $`3n`$ times the number of polyominoes of area (n+9) with one less hole. That is,
$$s_n^{(h)}3ns_{n+9}^{(h1)}$$
(A9)
### A.2 Growth constants (by area)
The following theorem proves the existence and equality of all the growth constants for $`k`$-punctured polyominoes for all finite $`k`$.
###### Theorem A4
There exists a constant $`\beta _0`$ such that for all $`h0`$
$$\underset{n\mathrm{}}{lim}n^1\mathrm{log}(s_n^{(h)})=\mathrm{log}(\beta _0).$$
(A10)
Proof: We use induction. First, $`\beta _0`$ exists . Assume $`\beta _h=lim_n\mathrm{}n^1\mathrm{log}(s_n^{(h)})`$ exists. From the results of the concatenation and surgery operations, we have
$$s_{nm}^{(h)}s_m^{(1)}s_n^{(h+1)}3ns_{n+9}^{(h)}.$$
Choose some value $`m`$ such that $`0<s_m^{(1)}<\mathrm{}`$, for example $`m=8`$. Take the logarithm, divide by $`n`$ and take the limit $`n\mathrm{}`$. This gives
$$\mathrm{log}(\beta _h)\mathrm{log}(\beta _{h+1})\mathrm{log}(\beta _h)$$
and hence $`\beta _h=\beta _{h+1}`$. Iterating from $`\beta _0`$ gives $`\beta _0=\beta _h`$ for all finite $`h`$. $`\mathrm{}`$
### A.3 Relationships between critical exponents
The following theorem establishes the relationships between critical exponents should they exist.
###### Theorem A5
Assume for all $`h`$ $`s_n^{(h)}C_hn^{\varphi _h}\beta _0^n`$ where $`C_h`$ is a $`h`$-dependent constant. Then
$$\varphi _h=\varphi _0h.$$
(A11)
Proof: From the results of the drilling operation, we have
$$C^h\left(\genfrac{}{}{0pt}{}{n/b^2}{h}\right)s_{n(b^21)}^{(0)}s_n^{(h)}.$$
(A12)
Since
$`\left({\displaystyle \genfrac{}{}{0pt}{}{n/b^2}{h}}\right)`$ $``$ $`{\displaystyle \frac{1}{h!}}n/b^2^h,`$
substituting the assumed asymptotic form $`s_n^{(h)}C_hn^{\varphi _h}\beta _h^n`$ into (A12), dividing by $`\beta _0^n`$, taking logarithms, dividing by $`\mathrm{log}(n)`$, letting $`n\mathrm{}`$ and using the above result, we get
$$h\varphi _0\varphi _h.$$
(A13)
Next, from the result of the surgery operation, we have
$`s_n^{(h)}3ns_{n+9}^{(h1)}`$ $``$ $`(3n)3(n+9)s_{n+18}^{(h2)}\mathrm{}`$
$``$ $`(3n)3(n+9)\mathrm{}3(n+9h)s_{n+9(h1)}^{(0)}[3(n+9h)]^hs_{n+9h}^{(0)}.`$
Again by substituting the assumed asymptotic form, dividing by $`\beta _0^{(n)}`$, taking logarithms, dividing by $`\mathrm{log}(n)`$ and letting $`n\mathrm{}`$, we get
$$\varphi _hh\varphi _0.$$
(A14)
Hence $`\varphi _h=\varphi _0h.`$ $`\mathrm{}`$
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# Noncanonical quantum optics
## I Introduction
The standard quantization of a harmonic oscillator is based on quantization of $`p`$ and $`q`$ but $`\omega `$ is a parameter. To have, say, two different frequencies one has to consider two independent oscillators. On the other hand, it is evident that there can exist oscillators which are in a quantum superposition of different frequencies. The example is an oscillator wave packet associated with distribution of center-of-mass momenta. It is known that the superposition of momenta gets translated into a superposition od Doppler shifts and therefore also of frequencies. We stress here the word “quantum” since the superpositions we have in mind are not those we know from classical oscillations.
This trivial observation raises the question of the role of superpositions of frequencies for a description of a single harmonic oscillator. The motivation behind the problem is associated with the question of field quantization: Is it possible that a quantum field consists of oscillators whose frequencies are indefinite? If so, maybe to quantize the field it is sufficient to use only one oscillator which exists in a quantum superposition of all the possible frequencies allowed by the boundary conditions of a given problem?
The idea is very simple. It is known that a “one-particle” state vector can be regarded as a representation of an ensemble of particles in a given pure state. On the other hand, the classical electromagnetic field can be regarded as an ensemble of oscillators. The standard idea of quantization, going back to 1925 , is to treat the field as an ensemble of quantum oscillators. But the ensemble itself is, in a sense, a classical one since for each frequency we need a separate oscillator. This is analogous to a classical ensemble of particles forming a classical wave on a lake surface. For each point on the surface we need a separate particle because a classical particle can ocupy only a single point in space. A quantum wave is of course different and we are all accustomed to the idea of a single-particle wave. In this case the properties of the entire ensemble are somehow encoded in properties of a single element of the ensemble.
For some reasons, probably partly historical, it seems that the idea of a single-particle state vector representation of the ensemble of oscillators has never been considered. The historical reason may be the fact that the very concept of field quantization occured already in 1925. At that stage quantum mechanics existed still in a matrix form and the Schrödinger paper “Quantisierung als Eigenwertproblem”, where the Schrödinger equation occured for the first time and the role of eigenvalues was explained, was not yet published. Actually, as explained by Jammer , Heisenberg quantized the harmonic oscillator without having heard of matrices so no wonder he could not treat the parameter $`\omega `$ entering $`E=p^2/2m+m\omega ^2q^2/2`$ as an eigenvalue of some operator.
Heisenberg’s quantization leads to the well known algebra of canonical commutation relations (CCR) $`[a_\omega ,a_\omega ^{}]=\mathrm{𝟏}`$, where $`\mathrm{𝟏}`$ is the identity operator and $`\omega `$ a classical parameter. As it turns out the replacemwnt of $`\omega `$ by an operator $`\widehat{\omega }`$ leads to non-canonical commutation relations (non-CCR) $`[a_\omega ,a_\omega ^{}]=\mathrm{𝟏}_\omega `$, where $`\omega `$ is an eigenvalue of $`\widehat{\omega }`$ and $`_\omega \mathrm{𝟏}_\omega =\mathrm{𝟏}`$. $`\mathrm{𝟏}_\omega `$, similarly to $`\mathrm{𝟏}`$ commutes with all creation and annihilation operators and the remaining commutators of non-CCR are the same as those of CCR. This subtle difference of the right-hand sides of CCR and non-CCR immediately explains why a non-CCR-quantized electromagnetic field will have vacuum with finite energy. Since the non-CCR and CCR algebras are so similar to each other it is not surprizing that the resulting theories are also very similar.
The main consequence of the non-CCR modification of CCR is different normalization of $`n`$-photon states. For example,
$`{\displaystyle \underset{\omega }{}}0|a_\omega a_\omega ^{}|0=1`$ (1)
and therefore $`0|a_\omega a_\omega ^{}|0<1`$. This should be contrasted with the CCR result $`0|a_\omega a_\omega ^{}|0=1`$, and the resulting divergence
$`{\displaystyle \underset{\omega }{}}0|a_\omega a_\omega ^{}|0=\mathrm{}`$ (2)
which is the source of infinite vacuum energy.
To end these introductory remarks one should mention that several approaches towards an alternative description of the electromagnetic field at a fundamental level were already proposed (e.g. Janes’ neoclassical theory, stochastic electrodynamics ). But the main idea of all such alternatives was to treat the field in classical terms and to associate the observed discreteness of emission/absorbtion phenomena with the quantum nature of atoms and not with the field itself.
The approach we will discuss in this paper does not belong to this tradition, is much more radical and, so to say, goes in the opposite direction. We will not try to make the field more classical. What we will try to do is to make it even more quantum by replacing classical parameters with eigenvalues.
## II Harmonic oscillator in superposition of frequencies
We know that frequency is typically associated with an eigenvalue of some Hamiltonian or, which is basically the same, with boundary conditions. A natural way of incorporating different frequencies into a single harmonic oscillator is by means of the frequency operator
$`\mathrm{\Omega }={\displaystyle \underset{\omega _k,j_k}{}}\omega _k|\omega _k,j_k\omega _k,j_k|`$ (3)
where all $`\omega _k0`$. For simplicity we have limited the discussion to the discrete spectrum but it is useful to include from the outset the possibility of degeneracies, represented here by the additional discrete quantum numbers $`j_k`$. The corresponding Hamiltonian is defined by
$`H`$ $`=`$ $`\mathrm{}\mathrm{\Omega }{\displaystyle \frac{1}{2}}\left(a^{}a+aa^{}\right)`$ (4)
where $`a=_{n=0}^{\mathrm{}}\sqrt{n+1}|nn+1|`$. The eigenstates of $`H`$ are $`|\omega _k,j_k,n`$ and satisfy the required formula
$`H|\omega _k,j_k,n=\mathrm{}\omega _k\left(n+{\displaystyle \frac{1}{2}}\right)|\omega _k,j_k,n`$ (5)
justifying our choice of $`H`$. The standard case of the oscillator whose frequency is just $`\omega `$ coresponds either to $`\mathrm{\Omega }=\omega \mathrm{𝟏}`$ or to the subspace spanned by $`|\omega _k,j_k,n`$ with fixed $`\omega _k=\omega `$. Introducing the operators
$`a_{\omega _k,j_k}=|\omega _k,j_k\omega _k,j_k|a`$ (6)
we find that
$`H`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\omega _k,j_k}{}}\mathrm{}\omega _k\left(a_{\omega _k,j_k}^{}a_{\omega _k,j_k}+a_{\omega _k,j_k}a_{\omega _k,j_k}^{}\right).`$ (7)
The algebra of the oscillator is “noncanonical”:
$`[a_{\omega _k,j_k},a_{\omega _k,j_k}^{}]`$ $`=`$ $`|\omega _k,j_k\omega _k,j_k|\mathrm{𝟏}=1_{\omega _k,j_k}`$ (8)
$`[a_{\omega _k,j_k},a_{\omega _l,j_l}^{}]`$ $`=`$ $`0\mathrm{for}(\omega _k,j_k)(\omega _l,j_l)`$ (9)
$`[a_{\omega _k,j_k},a_{\omega _l,j_l}]`$ $`=`$ $`0`$ (10)
$`[a_{\omega _k,j_k}^{},a_{\omega _l,j_l}^{}]`$ $`=`$ $`0`$ (11)
The dynamics in the Schrödinger picture is given by
$`i\mathrm{}_t|\mathrm{\Psi }`$ $`=`$ $`H|\mathrm{\Psi }=\mathrm{}\mathrm{\Omega }\left(a^{}a+{\displaystyle \frac{1}{2}}\mathrm{𝟏}\right)|\mathrm{\Psi }.`$ (12)
In the Heisenberg picture we obtain the important formula
$`a_{\omega _k,j_k}(t)`$ $`=`$ $`e^{iHt/\mathrm{}}a_{\omega _k,j_k}e^{iHt/\mathrm{}}`$ (13)
$`=`$ $`|\omega _k,j_k\omega _k,j_k|e^{i\omega _kt}a=e^{i\omega _kt}a_{\omega _k,j_k}.`$ (14)
Taking a general state
$`|\psi ={\displaystyle \underset{\omega _k,j_k,n}{}}\psi (\omega _k,j_k,n)|\omega _k,j_k,n`$ (15)
we find that the average energy of the oscillator is
$`H=\psi |H|\psi ={\displaystyle \underset{\omega _k,j_k,n}{}}|\psi (\omega _k,j_k,n)|^2\mathrm{}\omega _k\left(n+{\displaystyle \frac{1}{2}}\right).`$ (16)
The average clearly looks as an average energy of an ensemble of different and independent oscillators. The ground state of the ensemble, i.e. the one with $`\psi (\omega _k,j_k,n>0)=0`$ has energy
$`H={\displaystyle \frac{1}{2}}{\displaystyle \underset{\omega _k,j_k}{}}|\psi (\omega _k,j_k,0)|^2\mathrm{}\omega _k`$ (17)
which is finite if
$`{\displaystyle \underset{\omega _k,j_k}{}}\psi (\omega _k,j_k,0)|\omega _k,j_k`$ (18)
belongs to the domain of $`\mathrm{\Omega }`$. The result is not surprising but still quite remarkable if one thinks of the problem of field quantization.
The very idea of quantizing the electromagnetic field, as put forward by Born, Heisenberg, Jordan and Dirac , is based on the observation that the mode decomposition of the electromagnetic energy is analogous to the energy of an ensemble of independent harmonic oscillators. In 1925, after the work of Heisenberg, it was clear what to do: One had to replace each classical oscillator by a quantum one. But since each oscillator had a definite frequency, to have an infinite number of different frequencies one needed an infinite number of oscillators. The price one payed for this assumption was the infinite energy of the electromagnetic vacuum.
The infinity is regarded as an “easy” one since one can get rid of it by redefining the Hamiltonian and removing the infinite term. The result looks correct and many properties typical of a quantum harmonic oscillator are indeed observed in electromagnetic field. However, subtraction of infinite terms is in mathematics as forbidden as division by zero so to avoid evident absurdities one is forced to invent various ad hoc regularizations whose only justification is that otherwise the theory would not work. In larger perspective (say, in cosmology) it is not at all clear that an infinite (or arbitrarily cut off at the Planck scale) energy of the vacuum does not lead to contradictions with observational data . Finally, Dirac himself had never been fully satisfied by the theory he created. As Weinberg put it, Dirac’s “demand for a completely finite theory is similar to a host of other aesthetic judgements that theoretical physicists always need to make” .
The oscillator that can exist in superpositions of different frequencies is a natural candidate as a starting point for Dirac-type field quantization. Symbolically, if the Heisenberg quantization is $`p^2+\omega ^2q^2\widehat{p}^2+\omega ^2\widehat{q}^2`$, where $`\omega `$ is a parameter, the new scheme is $`p^2+\omega ^2q^2\widehat{p}^2+\widehat{\omega }^2\widehat{q}^2`$, where $`\widehat{\omega }`$ is an operator. Its spectrum can be related to boundary conditions imposed on the fields.
We do not need to remove the ground state energy since in the Hilbert space of physical states the correction is finite. The question we have to understand is whether one can obtain the well known quantum properties of the radiation field by this type of quantization.
## III “First quantization” — One-oscillator field operators
The new quantization will be performed in two steps. In this section we describe the first step, a kind of first quantization. In next sections we shall perform an analogue of second quantization which will lead to the final framework.
The energy and momentum operators of the field are defined in analogy to $`H`$ from the previous section
$`H`$ $`=`$ $`{\displaystyle \underset{s,\stackrel{}{\kappa }}{}}\mathrm{}\omega _\stackrel{}{\kappa }|s,\stackrel{}{\kappa }s,\stackrel{}{\kappa }|{\displaystyle \frac{1}{2}}\left(a^{}a+aa^{}\right)`$ (19)
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{s,\stackrel{}{\kappa }}{}}\mathrm{}\omega _\stackrel{}{\kappa }\left(a_{s,\stackrel{}{\kappa }}^{}a_{s,\stackrel{}{\kappa }}+a_{s,\stackrel{}{\kappa }}a_{s,\stackrel{}{\kappa }}^{}\right)`$ (20)
$`\stackrel{}{P}`$ $`=`$ $`{\displaystyle \underset{s,\stackrel{}{\kappa }}{}}\mathrm{}\stackrel{}{\kappa }|s,\stackrel{}{\kappa }s,\stackrel{}{\kappa }|{\displaystyle \frac{1}{2}}\left(a^{}a+aa^{}\right)`$ (21)
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{s,\stackrel{}{\kappa }}{}}\mathrm{}\stackrel{}{\kappa }\left(a_{s,\stackrel{}{\kappa }}^{}a_{s,\stackrel{}{\kappa }}+a_{s,\stackrel{}{\kappa }}a_{s,\stackrel{}{\kappa }}^{}\right)`$ (22)
where $`s=\pm 1`$ corresponds to circular polarizations. Denote $`P=(H/c,\stackrel{}{P})`$ and $`Px=Ht\stackrel{}{P}\stackrel{}{x}`$. We employ the standard Dirac-type definitions for mode quantization in volume $`V`$
$`\widehat{\stackrel{}{A}}(t,\stackrel{}{x})`$ $`=`$ $`{\displaystyle \underset{s,\stackrel{}{\kappa }}{}}\sqrt{{\displaystyle \frac{\mathrm{}}{2\omega _\stackrel{}{\kappa }V}}}\left(a_{s,\stackrel{}{\kappa }}e^{i\omega _\stackrel{}{\kappa }t}\stackrel{}{e}_{s,\stackrel{}{\kappa }}e^{i\stackrel{}{\kappa }\stackrel{}{x}}+a_{s,\stackrel{}{\kappa }}^{}e^{i\omega _\stackrel{}{\kappa }t}\stackrel{}{e}_{s,\stackrel{}{\kappa }}^{}e^{i\stackrel{}{\kappa }\stackrel{}{x}}\right)`$ (23)
$`=`$ $`e^{iPx/\mathrm{}}\widehat{\stackrel{}{A}}e^{iPx/\mathrm{}}`$ (24)
$`\widehat{\stackrel{}{E}}(t,\stackrel{}{x})`$ $`=`$ $`i{\displaystyle \underset{s,\stackrel{}{\kappa }}{}}\sqrt{{\displaystyle \frac{\mathrm{}\omega _\stackrel{}{\kappa }}{2V}}}\left(a_{s,\stackrel{}{\kappa }}e^{i\omega _\stackrel{}{\kappa }t}e^{i\stackrel{}{\kappa }\stackrel{}{x}}\stackrel{}{e}_{s,\stackrel{}{\kappa }}a_{s,\stackrel{}{\kappa }}^{}e^{i\omega _\stackrel{}{\kappa }t}e^{i\stackrel{}{\kappa }\stackrel{}{x}}\stackrel{}{e}_{s,\stackrel{}{\kappa }}^{}\right)`$ (25)
$`=`$ $`e^{iPx/\mathrm{}}\widehat{\stackrel{}{E}}e^{iPx/\mathrm{}}`$ (26)
$`\widehat{\stackrel{}{B}}(t,\stackrel{}{x})`$ $`=`$ $`i{\displaystyle \underset{s,\stackrel{}{\kappa }}{}}\sqrt{{\displaystyle \frac{\mathrm{}\omega _\stackrel{}{\kappa }}{2V}}}\stackrel{}{n}_\kappa \times \left(a_{s,\stackrel{}{\kappa }}e^{i\omega _\stackrel{}{\kappa }t}e^{i\stackrel{}{\kappa }\stackrel{}{x}}\stackrel{}{e}_{s,\stackrel{}{\kappa }}a_{s,\stackrel{}{\kappa }}^{}e^{i\omega _\stackrel{}{\kappa }t}e^{i\stackrel{}{\kappa }\stackrel{}{x}}\stackrel{}{e}_{s,\stackrel{}{\kappa }}^{}\right)`$ (28)
$`=`$ $`e^{iPx/\mathrm{}}\widehat{\stackrel{}{B}}e^{iPx/\mathrm{}},`$ (29)
where
$`a_{s,\stackrel{}{\kappa }}`$ $`=`$ $`|s,\stackrel{}{\kappa }s,\stackrel{}{\kappa }|a`$ (30)
$`a_{s,\stackrel{}{\kappa }}^{}`$ $`=`$ $`|s,\stackrel{}{\kappa }s,\stackrel{}{\kappa }|a^{}.`$ (31)
For later purposes we introduce the notation
$`[a_{s,\stackrel{}{\kappa }},a_{s,\stackrel{}{\kappa }}^{}]=1_{s,\stackrel{}{\kappa }}=|s,\stackrel{}{\kappa }s,\stackrel{}{\kappa }|\mathrm{𝟏}.`$ (32)
Now take a state (say, in the Heisenberg picture)
$`|\mathrm{\Psi }`$ $`=`$ $`{\displaystyle \underset{s,\stackrel{}{\kappa },n}{}}\mathrm{\Psi }_{s,\stackrel{}{\kappa },n}|s,\stackrel{}{\kappa },n`$ (33)
$`=`$ $`{\displaystyle \underset{s,\stackrel{}{\kappa }}{}}\mathrm{\Phi }_{s,\stackrel{}{\kappa }}|s,\stackrel{}{\kappa }|\alpha _{s,\stackrel{}{\kappa }}`$ (34)
where $`|\alpha _{s,\stackrel{}{\kappa }}`$ form a family of one-oscillator coherent states:
$`a|\alpha _{s,\stackrel{}{\kappa }}=\alpha _{s,\stackrel{}{\kappa }}|\alpha _{s,\stackrel{}{\kappa }}`$ (35)
The averages of the field operators are
$`\mathrm{\Psi }|\widehat{\stackrel{}{A}}(t,\stackrel{}{x})|\mathrm{\Psi }`$ $`=`$ $`{\displaystyle \underset{s,\stackrel{}{\kappa }}{}}|\mathrm{\Phi }_{s,\stackrel{}{\kappa }}|^2\sqrt{{\displaystyle \frac{\mathrm{}}{2\omega _\stackrel{}{\kappa }V}}}\left(\alpha _{s,\stackrel{}{\kappa }}e^{i\kappa x}\stackrel{}{e}_{s,\stackrel{}{\kappa }}+\alpha _{s,\stackrel{}{\kappa }}^{}e^{i\kappa x}\stackrel{}{e}_{s,\stackrel{}{\kappa }}^{}\right)`$ (36)
$`\mathrm{\Psi }|\widehat{\stackrel{}{E}}(t,\stackrel{}{x})|\mathrm{\Psi }`$ $`=`$ $`{\displaystyle \underset{s,\stackrel{}{\kappa }}{}}|\mathrm{\Phi }_{s,\stackrel{}{\kappa }}|^2\sqrt{{\displaystyle \frac{\mathrm{}\omega _\stackrel{}{\kappa }}{2V}}}\left(\alpha _{s,\stackrel{}{\kappa }}e^{i\kappa x}\stackrel{}{e}_{s,\stackrel{}{\kappa }}\alpha _{s,\stackrel{}{\kappa }}^{}e^{i\kappa x}\stackrel{}{e}_{s,\stackrel{}{\kappa }}^{}\right)`$ (37)
$`\mathrm{\Psi }|\widehat{\stackrel{}{B}}(t,\stackrel{}{x})|\mathrm{\Psi }`$ $`=`$ $`i{\displaystyle \underset{s,\stackrel{}{\kappa }}{}}|\mathrm{\Phi }_{s,\stackrel{}{\kappa }}|^2\sqrt{{\displaystyle \frac{\mathrm{}\omega _\stackrel{}{\kappa }}{2V}}}\left(\alpha _{s,\stackrel{}{\kappa }}e^{i\kappa x}\stackrel{}{n}_\kappa \times \stackrel{}{e}_{s,\stackrel{}{\kappa }}\alpha _{s,\stackrel{}{\kappa }}^{}e^{i\kappa x}\stackrel{}{n}_\kappa \times \stackrel{}{e}_{s,\stackrel{}{\kappa }}^{}\right)`$ (38)
These are just the classical fields. More precisely, the fields look like averages of monochromatic coherent states with probabilities $`|\mathrm{\Phi }_{s,\stackrel{}{\kappa }}|^2`$. The energy-momentum operators satisfy also the standard relations
$`H`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _V}d^3x\left(\widehat{\stackrel{}{E}}(t,\stackrel{}{x})\widehat{\stackrel{}{E}}(t,\stackrel{}{x})+\widehat{\stackrel{}{B}}(t,\stackrel{}{x})\widehat{\stackrel{}{B}}(t,\stackrel{}{x})\right),`$ (39)
$`\stackrel{}{P}`$ $`=`$ $`{\displaystyle _V}d^3x\widehat{\stackrel{}{E}}(t,\stackrel{}{x})\times \widehat{\stackrel{}{B}}(t,\stackrel{}{x}).`$ (40)
To end this section let us note that
$`\mathrm{\Psi }|H|\mathrm{\Psi }`$ $`=`$ $`{\displaystyle \underset{s,\stackrel{}{\kappa }}{}}\mathrm{}\omega _\stackrel{}{\kappa }|\mathrm{\Phi }_{s,\stackrel{}{\kappa }}|^2\left(|\alpha _{s,\stackrel{}{\kappa }}|^2+{\displaystyle \frac{1}{2}}\right)`$ (41)
$`\mathrm{\Psi }|\stackrel{}{P}|\mathrm{\Psi }`$ $`=`$ $`{\displaystyle \underset{s,\stackrel{}{\kappa }}{}}\mathrm{}\stackrel{}{\kappa }|\mathrm{\Phi }_{s,\stackrel{}{\kappa }}|^2\left(|\alpha _{s,\stackrel{}{\kappa }}|^2+{\displaystyle \frac{1}{2}}\right).`$ (42)
The contribution from the vacuum fluctuations is nonzero but finite. One can phrase the latter property also as follows. The noncanonical algebra of creation-annihilation operators satisfies the resolution of identity
$`{\displaystyle \underset{s,\stackrel{}{\kappa }}{}}[a_{s,\stackrel{}{\kappa }},a_{s,\stackrel{}{\kappa }}^{}]=\mathrm{𝟏}`$ (43)
wheras the canonical algebra would impliy
$`{\displaystyle \underset{s,\stackrel{}{\kappa }}{}}[a_{s,\stackrel{}{\kappa }},a_{s,\stackrel{}{\kappa }}^{}]=\mathrm{}\mathrm{𝟏}.`$ (44)
## IV “Second quantization”
The Hilbert space of states of the field we have constructed is spanned by vectors $`|s,\stackrel{}{\kappa },n`$. Still there is no doubt that both in reality (and the standard formalism) there exist multiparticle entangled states such as those spanned by tensor products of the form
$`|+,\stackrel{}{\kappa }_1,1|,\stackrel{}{\kappa }_2,1,`$ (45)
and the similar. It seems that there is no reason to limit our discussion to a single Hilbert space of a single oscillator. What we have done so far was a quantization of the electromagnetic field at the level of a “one-particle” Hilbert space. Similarly to quantization of other physical systems the next step is to consider many particles. What is not obvious (physically) is whether the oscillators should be considered as noninteracting. This physical freedom leads to two natural candidates for a free-field Hamiltonian.
The noninteracting extension is essentially clear. Having the one-particle energy-momentum operators $`P_a`$ (i.e. generators of 4-translations in the 1-particle Hilbert space) we define in the standard way their extensions to the Fock-type space
$`𝒫_a`$ $`=`$ $`P_a`$ (46)
$`\left(P_a\mathrm{𝟏}+\mathrm{𝟏}P_a\right)`$ (47)
$`\left(P_a\mathrm{𝟏}\mathrm{𝟏}+\mathrm{𝟏}P_a\mathrm{𝟏}+\mathrm{𝟏}\mathrm{𝟏}P_a\right)`$ (48)
$`\mathrm{}.`$ (49)
The $`x`$-dependence of fields is introduced similarly to the one-particle level
$`\stackrel{}{}(t,\stackrel{}{x})`$ $`=`$ $`e^{i𝒫x/\mathrm{}}\stackrel{}{}e^{i𝒫x/\mathrm{}}`$ (50)
but the field itself has yet to be defined. Assume
$`\stackrel{}{}`$ $`=`$ $`c_1\stackrel{}{F}`$ (51)
$`c_2\left(\stackrel{}{F}\mathrm{𝟏}+\mathrm{𝟏}\stackrel{}{F}\right)`$ (52)
$`c_3\left(\stackrel{}{F}\mathrm{𝟏}\mathrm{𝟏}+\mathrm{𝟏}\stackrel{}{F}\mathrm{𝟏}+\mathrm{𝟏}\mathrm{𝟏}\stackrel{}{F}\right)`$ (53)
$`\mathrm{}`$ (54)
where $`c_k`$ are constants discussed below, and $`\stackrel{}{F}`$ is $`\widehat{\stackrel{}{A}}`$, $`\widehat{\stackrel{}{E}}`$, or $`\widehat{\stackrel{}{B}}`$. The multi-oscillator annihilation operator associated with such fields must be therefore of the form
$`𝒂_{s,\stackrel{}{\kappa }}`$ $`=`$ $`c_1a_{s,\stackrel{}{\kappa }}`$ (55)
$`c_2\left(a_{s,\stackrel{}{\kappa }}\mathrm{𝟏}+\mathrm{𝟏}a_{s,\stackrel{}{\kappa }}\right)`$ (56)
$`c_3\left(a_{s,\stackrel{}{\kappa }}\mathrm{𝟏}\mathrm{𝟏}+\mathrm{𝟏}a_{s,\stackrel{}{\kappa }}\mathrm{𝟏}+\mathrm{𝟏}\mathrm{𝟏}a_{s,\stackrel{}{\kappa }}\right)`$ (57)
$`\mathrm{}.`$ (58)
Having two 1-particle operators, say $`X`$ and $`Y`$, one can easily establish a relation between the 1-particle commutator $`[X,Y]`$ and the commutator of the extensions $`𝒳`$, $`𝒴`$:
$`[𝒳,𝒴]`$ $`=`$ $`c_1^2[X,Y]`$ (59)
$`c_2^2\left([X,Y]\mathrm{𝟏}+\mathrm{𝟏}[X,Y]\right)`$ (60)
$`c_3^2\left([X,Y]\mathrm{𝟏}\mathrm{𝟏}+\mathrm{𝟏}[X,Y]\mathrm{𝟏}+\mathrm{𝟏}\mathrm{𝟏}[X,Y]\right)`$ (61)
$`\mathrm{}.`$ (62)
The annihilation operators so defined satisfy therefore the algebra
$`[𝒂_{s,\stackrel{}{\kappa }},𝒂_{s^{},\stackrel{}{\kappa }^{}}^{}]`$ $`=`$ $`0\mathrm{for}(s,\stackrel{}{\kappa })(s^{},\stackrel{}{\kappa }^{}),`$ (63)
$`[𝒂_{s,\stackrel{}{\kappa }},𝒂_{s,\stackrel{}{\kappa }}^{}]`$ $`=`$ $`\mathrm{𝟏}_{s,\stackrel{}{\kappa }},`$ (64)
$`[𝒂_{s,\stackrel{}{\kappa }},𝒂_{s^{},\stackrel{}{\kappa }^{}}]`$ $`=`$ $`0`$ (65)
$`[𝒂_{s,\stackrel{}{\kappa }}^{},𝒂_{s^{},\stackrel{}{\kappa }^{}}^{}]`$ $`=`$ $`0`$ (66)
where the operator $`\mathrm{𝟏}_{s,\stackrel{}{\kappa }}`$ is defined by
$`\mathrm{𝟏}_{s,\stackrel{}{\kappa }}`$ $`=`$ $`c_1^21_{s,\stackrel{}{\kappa }}`$ (67)
$`c_2^2\left(1_{s,\stackrel{}{\kappa }}\mathrm{𝟏}+\mathrm{𝟏}1_{s,\stackrel{}{\kappa }}\right)`$ (68)
$`c_3^2\left(1_{s,\stackrel{}{\kappa }}\mathrm{𝟏}\mathrm{𝟏}+\mathrm{𝟏}1_{s,\stackrel{}{\kappa }}\mathrm{𝟏}+\mathrm{𝟏}\mathrm{𝟏}1_{s,\stackrel{}{\kappa }}\right)`$ (69)
$`\mathrm{},`$ (70)
and $`1_{s,\stackrel{}{\kappa }}`$ is a single-oscillator operator (32).
An important property of the 1-oscillator description was the resolution of identity (43). The requirement that the same be valid at the multi oscillator level leads to $`c_n=1/\sqrt{n}`$. In such a case one finds that
$`\mathrm{𝟏}_{s,\stackrel{}{\kappa }}^2\mathrm{𝟏}_{s,\stackrel{}{\kappa }}`$ (71)
but nevertheless
$`{\displaystyle \underset{s,\stackrel{}{\kappa }}{}}\mathrm{𝟏}_{s,\stackrel{}{\kappa }}=\mathrm{𝟏},`$ (72)
that is $`\mathrm{𝟏}_{s,\stackrel{}{\kappa }}`$ are the so-called positive operator valued (POV) measures . Below we shall give still another justification of this particular choice of $`c_n`$.
We can finally write
$`\stackrel{}{𝒜}(t,\stackrel{}{x})`$ $`=`$ $`{\displaystyle \underset{s,\stackrel{}{\kappa }}{}}\sqrt{{\displaystyle \frac{\mathrm{}}{2\omega _\stackrel{}{\kappa }V}}}\left(𝒂_{s,\stackrel{}{\kappa }}e^{i\omega _\stackrel{}{\kappa }t}\stackrel{}{e}_{s,\stackrel{}{\kappa }}e^{i\stackrel{}{\kappa }\stackrel{}{x}}+𝒂_{s,\stackrel{}{\kappa }}^{}e^{i\omega _\stackrel{}{\kappa }t}\stackrel{}{e}_{s,\stackrel{}{\kappa }}^{}e^{i\stackrel{}{\kappa }\stackrel{}{x}}\right)`$ (73)
$`=`$ $`e^{i𝒫x/\mathrm{}}\stackrel{}{𝒜}e^{i𝒫x/\mathrm{}}`$ (74)
$`\stackrel{}{}(t,\stackrel{}{x})`$ $`=`$ $`i{\displaystyle \underset{s,\stackrel{}{\kappa }}{}}\sqrt{{\displaystyle \frac{\mathrm{}\omega _\stackrel{}{\kappa }}{2V}}}\left(𝒂_{s,\stackrel{}{\kappa }}e^{i\omega _\stackrel{}{\kappa }t}e^{i\stackrel{}{\kappa }\stackrel{}{x}}\stackrel{}{e}_{s,\stackrel{}{\kappa }}𝒂_{s,\stackrel{}{\kappa }}^{}e^{i\omega _\stackrel{}{\kappa }t}e^{i\stackrel{}{\kappa }\stackrel{}{x}}\stackrel{}{e}_{s,\stackrel{}{\kappa }}^{}\right)`$ (75)
$`=`$ $`e^{i𝒫x/\mathrm{}}\stackrel{}{}e^{i𝒫x/\mathrm{}}`$ (76)
$`\stackrel{}{}(t,\stackrel{}{x})`$ $`=`$ $`i{\displaystyle \underset{s,\stackrel{}{\kappa }}{}}\sqrt{{\displaystyle \frac{\mathrm{}\omega _\stackrel{}{\kappa }}{2V}}}\stackrel{}{n}_\kappa \times \left(𝒂_{s,\stackrel{}{\kappa }}e^{i\omega _\stackrel{}{\kappa }t}e^{i\stackrel{}{\kappa }\stackrel{}{x}}\stackrel{}{e}_{s,\stackrel{}{\kappa }}𝒂_{s,\stackrel{}{\kappa }}^{}e^{i\omega _\stackrel{}{\kappa }t}e^{i\stackrel{}{\kappa }\stackrel{}{x}}\stackrel{}{e}_{s,\stackrel{}{\kappa }}^{}\right)`$ (78)
$`=`$ $`e^{i𝒫x/\mathrm{}}\stackrel{}{}e^{i𝒫x/\mathrm{}}.`$ (79)
These operators form a basis of the modified version of nonrelativistic quantum optics.
Let us return for the moment to the case of a general $`c_n`$. A straightforward calculation shows that
$`𝐇`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _V}d^3x\left(\stackrel{}{}(t,\stackrel{}{x})\stackrel{}{}(t,\stackrel{}{x})+\stackrel{}{}(t,\stackrel{}{x})\stackrel{}{}(t,\stackrel{}{x})\right)={\displaystyle \frac{1}{2}}{\displaystyle \underset{s,\stackrel{}{\kappa }}{}}\mathrm{}\omega _\stackrel{}{\kappa }\left(𝒂_{s,\stackrel{}{\kappa }}^{}𝒂_{s,\stackrel{}{\kappa }}+𝒂_{s,\stackrel{}{\kappa }}𝒂_{s,\stackrel{}{\kappa }}^{}\right)`$ (80)
$`=`$ $`{\displaystyle \underset{s,\stackrel{}{\kappa }}{}}\mathrm{}\omega _\stackrel{}{\kappa }[c_1^2\frac{1}{2}\{a_{s,\stackrel{}{\kappa }},a_{s,\stackrel{}{\kappa }}^{}\}`$ (81)
$`c_2^2\left(\frac{1}{2}\{a_{s,\stackrel{}{\kappa }},a_{s,\stackrel{}{\kappa }}^{}\}\mathrm{𝟏}+\mathrm{𝟏}\frac{1}{2}\{a_{s,\stackrel{}{\kappa }},a_{s,\stackrel{}{\kappa }}^{}\}+a_{s,\stackrel{}{\kappa }}^{}a_{s,\stackrel{}{\kappa }}+a_{s,\stackrel{}{\kappa }}a_{s,\stackrel{}{\kappa }}^{}\right)`$ (82)
$`c_3^2(\frac{1}{2}\{a_{s,\stackrel{}{\kappa }},a_{s,\stackrel{}{\kappa }}^{}\}\mathrm{𝟏}\mathrm{𝟏}+\mathrm{𝟏}\frac{1}{2}\{a_{s,\stackrel{}{\kappa }},a_{s,\stackrel{}{\kappa }}^{}\}\mathrm{𝟏}+\mathrm{𝟏}\mathrm{𝟏}\frac{1}{2}\{a_{s,\stackrel{}{\kappa }},a_{s,\stackrel{}{\kappa }}^{}\}`$ (83)
$`+a_{s,\stackrel{}{\kappa }}a_{s,\stackrel{}{\kappa }}^{}\mathrm{𝟏}+a_{s,\stackrel{}{\kappa }}^{}a_{s,\stackrel{}{\kappa }}\mathrm{𝟏}+\mathrm{𝟏}a_{s,\stackrel{}{\kappa }}a_{s,\stackrel{}{\kappa }}^{}`$ (84)
$`+\mathrm{𝟏}a_{s,\stackrel{}{\kappa }}^{}a_{s,\stackrel{}{\kappa }}+a_{s,\stackrel{}{\kappa }}^{}\mathrm{𝟏}a_{s,\stackrel{}{\kappa }}+a_{s,\stackrel{}{\kappa }}\mathrm{𝟏}a_{s,\stackrel{}{\kappa }}^{})`$ (85)
$`\mathrm{}]`$ (86)
where $`\{,\}`$ denotes the anti-commutator. Comparing this with the generator of time translations
$`=c𝒫_0`$ $`=`$ $`{\displaystyle \underset{s,\stackrel{}{\kappa }}{}}\mathrm{}\omega _\stackrel{}{\kappa }[\frac{1}{2}\{a_{s,\stackrel{}{\kappa }},a_{s,\stackrel{}{\kappa }}^{}\}`$ (87)
$`\left(\frac{1}{2}\{a_{s,\stackrel{}{\kappa }},a_{s,\stackrel{}{\kappa }}^{}\}\mathrm{𝟏}+\mathrm{𝟏}\frac{1}{2}\{a_{s,\stackrel{}{\kappa }},a_{s,\stackrel{}{\kappa }}^{}\}\right)`$ (88)
$`\left(\frac{1}{2}\{a_{s,\stackrel{}{\kappa }},a_{s,\stackrel{}{\kappa }}^{}\}\mathrm{𝟏}\mathrm{𝟏}+\mathrm{𝟏}\frac{1}{2}\{a_{s,\stackrel{}{\kappa }},a_{s,\stackrel{}{\kappa }}^{}\}\mathrm{𝟏}+\mathrm{𝟏}\mathrm{𝟏}\frac{1}{2}\{a_{s,\stackrel{}{\kappa }},a_{s,\stackrel{}{\kappa }}^{}\}\right)`$ (89)
$`\mathrm{}]`$ (90)
we can see that there is a relation between $``$ and $`𝐇`$ but the latter contains terms describing interactions between the oscillators. The contribution from these interactions vanishes on vacuum states. Below, when we introduce the notion of a generalized coherent state, we will be able to relate averages of $``$ and $`𝐇`$. In a similar way one can introduce the “Pointing operator”
$`\stackrel{}{𝐏}`$ $`=`$ $`{\displaystyle _V}d^3x\stackrel{}{}(t,\stackrel{}{x})\times \stackrel{}{}(t,\stackrel{}{x})={\displaystyle \frac{1}{2}}{\displaystyle \underset{s,\stackrel{}{\kappa }}{}}\mathrm{}\stackrel{}{\kappa }\left(𝒂_{s,\stackrel{}{\kappa }}^{}𝒂_{s,\stackrel{}{\kappa }}+𝒂_{s,\stackrel{}{\kappa }}𝒂_{s,\stackrel{}{\kappa }}^{}\right).`$ (91)
Its relation to the generator of 3-translations $`\stackrel{}{𝒫}`$ is similar to this between $``$ and $`𝐇`$.
In the above construction the only element which is beyond a simple transition to many oscillators is the choice of $`c_n`$. For different choices of these constants we obtain different representations of non-CCR and therefore also different quantization schemes. Several different ways of reasoning lead to $`c_n=1/\sqrt{n}`$ as we shall also see in the next sections.
## V Some particular states
We assume that all the multi-oscillator states are symmetric with respect to permutations of the oscillators.
### A Generalized coherent states
For general $`c_n`$ an eigenstate of $`𝒂_{s,\stackrel{}{\kappa }}`$ corresponding to the eigenvalue $`\alpha _{s,\stackrel{}{\kappa }}`$ is of the form
$`|𝜶_{s,\stackrel{}{\kappa }}`$ $`=`$ $`f_1(s,\stackrel{}{\kappa })|s,\stackrel{}{\kappa },\alpha _{s,\stackrel{}{\kappa }}/c_1`$ (92)
$`f_2(s,\stackrel{}{\kappa })|s,\stackrel{}{\kappa },\alpha _{s,\stackrel{}{\kappa }}/(2c_2)|s,\stackrel{}{\kappa },\alpha _{s,\stackrel{}{\kappa }}/(2c_2)`$ (93)
$`f_3(s,\stackrel{}{\kappa })|s,\stackrel{}{\kappa },\alpha _{s,\stackrel{}{\kappa }}/(3c_3)|s,\stackrel{}{\kappa },\alpha _{s,\stackrel{}{\kappa }}/(3c_3)|s,\stackrel{}{\kappa },\alpha _{s,\stackrel{}{\kappa }}/(3c_3)`$ (94)
$`\mathrm{}`$ (95)
where
$`|s,\stackrel{}{\kappa },\alpha _{s,\stackrel{}{\kappa }}=|s,\stackrel{}{\kappa }|\alpha _{s,\stackrel{}{\kappa }},`$ (96)
$`_{n=1}^{\mathrm{}}|f_n(s,\stackrel{}{\kappa })|^2=1`$, and $`a|\alpha _{s,\stackrel{}{\kappa }}=\alpha _{s,\stackrel{}{\kappa }}|\alpha _{s,\stackrel{}{\kappa }}`$. What is interesting not all $`f_n`$ have to be nonvanishing.
The average “energies” of the field in the above eigenstate are
$`𝜶_{s,\stackrel{}{\kappa }}||𝜶_{s,\stackrel{}{\kappa }}`$ $`=`$ $`\mathrm{}\omega _\stackrel{}{\kappa }|\alpha _{s,\stackrel{}{\kappa }}|^2{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{nc_n^2}}|f_n(s,\stackrel{}{\kappa })|^2+{\displaystyle \frac{1}{2}}\mathrm{}\omega _\stackrel{}{\kappa }{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}n|f_n(s,\stackrel{}{\kappa })|^2`$ (97)
and
$`𝜶_{s,\stackrel{}{\kappa }}|𝐇|𝜶_{s,\stackrel{}{\kappa }}`$ $`=`$ $`\mathrm{}\omega _\stackrel{}{\kappa }|\alpha _{s,\stackrel{}{\kappa }}|^2+{\displaystyle \frac{1}{2}}\mathrm{}\omega _\stackrel{}{\kappa }{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}nc_n^2|f_n(s,\stackrel{}{\kappa })|^2.`$ (98)
The two averages will differ only by the value of the vacuum contribution if $`c_n=1/\sqrt{n}`$ which leads us back to the above mentioned choice of $`c_n`$. With this choice and taking the general combination of coherent states
$`|𝚿={\displaystyle \underset{s,\stackrel{}{\kappa }}{}}\mathrm{\Phi }_{s,\stackrel{}{\kappa }}|𝜶_{s,\stackrel{}{\kappa }}`$ (99)
we find
$`𝚿||𝚿`$ $`=`$ $`{\displaystyle \underset{s,\stackrel{}{\kappa }}{}}\mathrm{}\omega _\stackrel{}{\kappa }|\mathrm{\Phi }_{s,\stackrel{}{\kappa }}|^2|\alpha _{s,\stackrel{}{\kappa }}|^2+{\displaystyle \frac{1}{2}}{\displaystyle \underset{s,\stackrel{}{\kappa }}{}}\mathrm{}\omega _\stackrel{}{\kappa }|\mathrm{\Phi }_{s,\stackrel{}{\kappa }}|^2{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}n|f_k(s,\stackrel{}{\kappa })|^2`$ (100)
and
$`𝚿|𝐇|𝚿`$ $`=`$ $`{\displaystyle \underset{s,\stackrel{}{\kappa }}{}}\mathrm{}\omega _\stackrel{}{\kappa }|\mathrm{\Phi }_{s,\stackrel{}{\kappa }}|^2|\alpha _{s,\stackrel{}{\kappa }}|^2+{\displaystyle \frac{1}{2}}{\displaystyle \underset{s,\stackrel{}{\kappa }}{}}\mathrm{}\omega _\stackrel{}{\kappa }|\mathrm{\Phi }_{s,\stackrel{}{\kappa }}|^2.`$ (101)
One may wonder, then, what is the more natural choice of the free-field Hamiltonian: $`𝐇`$ describing interacting oscillators, or $``$ describing the noninteracting ones? The coherent-state average of $`𝐇`$ does not depend on the average number of the oscillators and naturally includes the process of energy exchange between different oscillators. Moreover, both $`𝐇`$ and $`\stackrel{}{𝐏}`$ are defined in the standard way in terms of the multi-oscillator non-CCR algebra. With this choice of free dynamics we find
$`e^{i𝐇t/\mathrm{}}𝒂_{s,\stackrel{}{\kappa }}e^{i𝐇t/\mathrm{}}=e^{i\omega _\stackrel{}{\kappa }t\mathrm{𝟏}_{s,\stackrel{}{\kappa }}}𝒂_{s,\stackrel{}{\kappa }}`$ (102)
as opposed to the standard formula
$`e^{it/\mathrm{}}𝒂_{s,\stackrel{}{\kappa }}e^{it/\mathrm{}}=e^{i\omega _\stackrel{}{\kappa }t}𝒂_{s,\stackrel{}{\kappa }}.`$ (103)
The latter choice is simpler because it leads to the standard form of the interaction-picture Hamiltonian and therefore will be the basis of our non-canonical quantum optics. The version based on $`𝐇`$ is a subject of ongoing study.
### B Vacuum
Similarly to the one-oscillator case the traditional notion of a vacuum state is replaced in our formalism by a vacuum subspace consisting of all the vectors annihilated by all $`𝒂_{s,\stackrel{}{\kappa }}`$. Their general form is
$`|\mathrm{𝟎}`$ $`=`$ $`{\displaystyle \underset{s,\stackrel{}{\kappa }}{}}O_{s,\stackrel{}{\kappa },0}^{(1)}|s,\stackrel{}{\kappa },0`$ (104)
$`{\displaystyle \underset{s_j,\stackrel{}{\kappa }_j}{}}O_{s_1,s_2,\stackrel{}{\kappa }_1,\stackrel{}{\kappa }_2,0,0}^{(2)}|s_1,\stackrel{}{\kappa }_1,0|s_2,\stackrel{}{\kappa }_2,0`$ (105)
$`{\displaystyle \underset{s_j,\stackrel{}{\kappa }_{\lambda _j}}{}}O_{s_1,s_2,s_3\stackrel{}{\kappa }_{\lambda _1},\stackrel{}{\kappa }_{\lambda _2},\stackrel{}{\kappa }_{\lambda _3},0,0,0}^{(3)}|s_1,\stackrel{}{\kappa }_1,0|s_2,\stackrel{}{\kappa }_2,0|s_3,\stackrel{}{\kappa }_3,0`$ (106)
$`\mathrm{}`$ (107)
It seems that there is no reason for introducing the standard unique “vacuum state” understood as the cyclic vector of the GNS construction.
As an example of a vacuum state consider
$`|\mathrm{𝟎}`$ $`=`$ $`\sqrt{p_1}|O`$ (108)
$`\sqrt{p_2}|O|O`$ (109)
$`\sqrt{p_3}|O|O|O`$ (110)
$`\mathrm{}`$ (111)
The average energy of the free-field vacuum state is therefore
$`\overline{}=\mathrm{𝟎}||\mathrm{𝟎}={\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}np_nO|H|O=\overline{n}\overline{H}`$ (112)
where $`\overline{n}`$ and $`\overline{H}`$ are, respectively, the average number of oscillators and the average energy of a single oscillator. For the sake of completeness let us note that
$`\overline{𝐇}=\mathrm{𝟎}|𝐇|\mathrm{𝟎}={\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}c_n^2np_nO|H|O=\overline{H}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}c_n^2np_n`$ (113)
For $`c_n=1`$ we obtain $`\overline{𝐇}=\overline{}`$; for $`c_n=1/\sqrt{n}`$ $`\overline{𝐇}=\overline{H}`$ the latter being independent of the number of oscillators. In both cases no problem with infinite vacuum energy is found. Obviously, one can contemplate also other vacua, say, in entangled or mixed states.
### C Multi-photon states
Assume $`|\mathrm{𝟎}`$ is a vacuum state. The non-canonical algebra (non-CCR) differes from the canonical one in the commutator
$`[𝒂_\lambda ,𝒂_\lambda ^{}]=\mathrm{𝟏}_\lambda `$ (114)
where for any $`\lambda `$, $`\lambda ^{}`$
$`[\mathrm{𝟏}_\lambda ^{},𝒂_\lambda ]=[\mathrm{𝟏}_\lambda ^{},𝒂_\lambda ^{}]=0`$ (115)
where $`\lambda `$ stands for $`(s,\stackrel{}{\kappa })`$. A normalized state describing a collection of photons is defined in analogy to the standard formalism by
$`{\displaystyle \frac{1}{\sqrt{n_1!n_2!\mathrm{}n_N!\mathrm{𝟎}|\mathrm{𝟏}_{\lambda _N}^{n_N}\mathrm{}\mathrm{𝟏}_{\lambda _2}^{n_2}\mathrm{𝟏}_{\lambda _1}^{n_1}|\mathrm{𝟎}}}}(𝒂_{\lambda _N}^{})^{n_N}\mathrm{}(𝒂_{\lambda _2}^{})^{n_2}(𝒂_{\lambda _1}^{})^{n_1}|\mathrm{𝟎}=:|𝒏_{\lambda _1},𝒏_{\lambda _2},\mathrm{},𝒏_{\lambda _N}`$ (116)
States corresponding to the same $`\lambda `$ but different $`n`$’s, or to the same $`n`$ but different $`\lambda `$’s, are orthogonal. As a consequence a non-canonical vacuum average of any product of non-canonical creation and annihilation operators vanishes if and only if an analogous expression formulated in terms of the canonical objects does. This property is a consequence of three facts which hold true in both formalisms: (a) annihilation operators annihilate vacuum states, (b) creation operators are obtained by a Hermitian conjugate of the annihilation operators, and (c) the RHS of a commutator of creation and annihilation operators commutes with all creation and annihilation operators.
## VI Perturbation theory
It is essential that, similarly to the one-oscillator formalism, the free Hamiltonian (defined simply as a generator of time translations) generates the standard dynamics of annihilation operators:
$`e^{it/\mathrm{}}𝒂_{s,\stackrel{}{\kappa }}e^{it/\mathrm{}}=e^{i\omega _\stackrel{}{\kappa }t}𝒂_{s,\stackrel{}{\kappa }}.`$ (117)
Accordingly, the form of the interaction-picture Hamiltonian will by the same as in the standard theory. This would not be quite the same if we have chosen $`𝐇`$ in the role of the free Hamiltonian (an option which, nevertheless, should be investigated). In what follows we start with $`H=H_0+V`$, where
$`H_0`$ $`=`$ $`H_A+`$ (118)
$`V`$ $`=`$ $`{\displaystyle \frac{e}{m}}\stackrel{}{𝒜}(\stackrel{}{x})\stackrel{}{p}`$ (119)
$`=`$ $`{\displaystyle \frac{e}{m}}{\displaystyle \underset{s,\stackrel{}{\kappa }}{}}\sqrt{{\displaystyle \frac{\mathrm{}}{2\omega _\stackrel{}{\kappa }V}}}\left(𝒂_{s,\stackrel{}{\kappa }}e^{i\stackrel{}{\kappa }\stackrel{}{x}}\stackrel{}{e}_{s,\stackrel{}{\kappa }}\stackrel{}{p}+𝒂_{s,\stackrel{}{\kappa }}^{}e^{i\stackrel{}{\kappa }\stackrel{}{x}}\stackrel{}{e}_{s,\stackrel{}{\kappa }}^{}\stackrel{}{p}\right).`$ (120)
In the interaction picture we get
$`V(t)`$ $`=`$ $`{\displaystyle \frac{e}{m}}\stackrel{}{𝒜}(t,\stackrel{}{x})\stackrel{}{p}(t)`$ (121)
$`=`$ $`{\displaystyle \frac{e}{m}}{\displaystyle \underset{s,\stackrel{}{\kappa }}{}}\sqrt{{\displaystyle \frac{\mathrm{}}{2\omega _\stackrel{}{\kappa }V}}}\left(𝒂_{s,\stackrel{}{\kappa }}e^{i(\omega _\stackrel{}{\kappa }t\stackrel{}{\kappa }\stackrel{}{x})}\stackrel{}{e}_{s,\stackrel{}{\kappa }}\stackrel{}{p}(t)+𝒂_{s,\stackrel{}{\kappa }}^{}e^{i(\omega _\stackrel{}{\kappa }t\stackrel{}{\kappa }\stackrel{}{x})}\stackrel{}{e}_{s,\stackrel{}{\kappa }}^{}\stackrel{}{p}(t)\right),`$ (122)
and
$`\stackrel{}{p}(t)=e^{iH_At}\stackrel{}{p}e^{iH_At}.`$ (123)
Since we are purposefully neglecting the “$`\stackrel{}{A}^2`$” term in the Hamiltonian, one should restrict the analysis to the dipole approximation and therefore it is justified to set $`\stackrel{}{x}=0`$:
$`V(t)`$ $`=`$ $`{\displaystyle \underset{\lambda }{}}\left(𝒂_\lambda \widehat{g}_\lambda (t)+𝒂_\lambda ^{}\widehat{g}_\lambda ^{}(t)\right).`$ (124)
The operators
$`\widehat{g}_{s,\stackrel{}{\kappa }}(t)={\displaystyle \frac{e}{m}}\sqrt{{\displaystyle \frac{\mathrm{}}{2\omega _\stackrel{}{\kappa }V}}}e^{i\omega _\stackrel{}{\kappa }t}\stackrel{}{e}_{s,\stackrel{}{\kappa }}\stackrel{}{p}(t)=\widehat{g}_\lambda (t)`$ (125)
are identical to those from the standard formalism and act only on atomic degrees of freedom (i.e. commute with $`𝒂_{s,\stackrel{}{\kappa }}`$).
### A Spontaneous decay of an excited state
The first problem we shall treat in the non-canonical way is a lifetime of an excited atomic state. The problem, as we shall see, is of particular importance for the physical interpretation of the non-canonical formalism.
Assume that at $`t=0`$ the atom-field system is described by the state $`|\mathrm{\Psi }(0)=|\mathrm{𝟎},A`$. The amplitude that the atom remains in the excited state is
$`\mathrm{𝟎},A|\mathrm{\Psi }(t)`$ $`=`$ $`1`$ (126)
$`+{\displaystyle \frac{1}{(i\mathrm{})^2}}{\displaystyle _0^t}𝑑t_2{\displaystyle _0^{t_2}}𝑑t_1{\displaystyle \underset{\lambda _1\lambda _2}{}}\mathrm{𝟎}|𝒂_{\lambda _2}𝒂_{\lambda _1}^{}|\mathrm{𝟎}A|\widehat{g}_{\lambda _2}(t_2)\widehat{g}_{\lambda _1}^{}(t_1)|A`$ (127)
$`+{\displaystyle \frac{1}{(i\mathrm{})^4}}{\displaystyle _0^t}𝑑t_4{\displaystyle _0^{t_4}}𝑑t_3{\displaystyle _0^{t_3}}𝑑t_2{\displaystyle _0^{t_2}}𝑑t_1{\displaystyle \underset{\lambda _1\lambda _2\lambda _3\lambda _4}{}}\mathrm{𝟎}|𝒂_{\lambda _4}𝒂_{\lambda _3}^{}𝒂_{\lambda _2}𝒂_{\lambda _1}^{}|\mathrm{𝟎}A|\widehat{g}_{\lambda _4}(t_4)\widehat{g}_{\lambda _3}^{}(t_3)\widehat{g}_{\lambda _2}(t_2)\widehat{g}_{\lambda _1}^{}(t_1)|A`$ (128)
$`+{\displaystyle \frac{1}{(i\mathrm{})^4}}{\displaystyle _0^t}𝑑t_4{\displaystyle _0^{t_4}}𝑑t_3{\displaystyle _0^{t_3}}𝑑t_2{\displaystyle _0^{t_2}}𝑑t_1{\displaystyle \underset{\lambda _1\lambda _2\lambda _3\lambda _4}{}}\mathrm{𝟎}|𝒂_{\lambda _4}𝒂_{\lambda _3}𝒂_{\lambda _2}^{}𝒂_{\lambda _1}^{}|\mathrm{𝟎}A|\widehat{g}_{\lambda _4}(t_4)\widehat{g}_{\lambda _3}(t_3)\widehat{g}_{\lambda _2}^{}(t_2)\widehat{g}_{\lambda _1}^{}(t_1)|A`$ (129)
$`+\mathrm{}`$ (130)
$`=`$ $`1`$ (131)
$`+{\displaystyle \frac{1}{(i\mathrm{})^2}}{\displaystyle _0^t}𝑑t_2{\displaystyle _0^{t_2}}𝑑t_1{\displaystyle \underset{\lambda _1\lambda _2}{}}\widehat{0}|\widehat{a}_{\lambda _2}\widehat{a}_{\lambda _1}^{}|\widehat{0}A|\widehat{g}_{\lambda _2}(t_2)\widehat{g}_{\lambda _1}^{}(t_1)|AX_{\lambda _2\lambda _1}^{01}`$ (132)
$`+{\displaystyle \frac{1}{(i\mathrm{})^4}}{\displaystyle _0^t}𝑑t_4{\displaystyle _0^{t_4}}𝑑t_3{\displaystyle _0^{t_3}}𝑑t_2{\displaystyle _0^{t_2}}𝑑t_1{\displaystyle \underset{\lambda _1\lambda _2\lambda _3\lambda _4}{}}\widehat{0}|\widehat{a}_{\lambda _4}\widehat{a}_{\lambda _3}^{}\widehat{a}_{\lambda _2}\widehat{a}_{\lambda _1}^{}|\widehat{0}A|\widehat{g}_{\lambda _4}(t_4)\widehat{g}_{\lambda _3}^{}(t_3)\widehat{g}_{\lambda _2}(t_2)\widehat{g}_{\lambda _1}^{}(t_1)|A`$ (133)
$`\times X_{\lambda _4\lambda _3\lambda _2\lambda _1}^{0101}`$ (134)
$`+{\displaystyle \frac{1}{(i\mathrm{})^4}}{\displaystyle _0^t}𝑑t_4{\displaystyle _0^{t_4}}𝑑t_3{\displaystyle _0^{t_3}}𝑑t_2{\displaystyle _0^{t_2}}𝑑t_1{\displaystyle \underset{\lambda _1\lambda _2\lambda _3\lambda _4}{}}\widehat{0}|\widehat{a}_{\lambda _4}\widehat{a}_{\lambda _3}\widehat{a}_{\lambda _2}^{}\widehat{a}_{\lambda _1}^{}|\widehat{0}A|\widehat{g}_{\lambda _4}(t_4)\widehat{g}_{\lambda _3}(t_3)\widehat{g}_{\lambda _2}^{}(t_2)\widehat{g}_{\lambda _1}^{}(t_1)|A`$ (135)
$`\times X_{\lambda _4\lambda _3\lambda _2\lambda _1}^{0011}`$ (136)
$`+\mathrm{}`$ (137)
In the above perturbative expansion we have explicitly shown all the nonvanishing terms up to the fifth order of perturbation theory. Here $`\widehat{0}|\widehat{a}_{\lambda _4}\widehat{a}_{\lambda _3}\widehat{a}_{\lambda _2}^{}\widehat{a}_{\lambda _1}^{}|\widehat{0}`$ etc. are the expressions one would have obtained in ordinary canonical formalism and
$`X_{\lambda _4\lambda _3\lambda _2\lambda _1}^{0011}={\displaystyle \frac{\mathrm{𝟎}|𝒂_{\lambda _4}𝒂_{\lambda _3}𝒂_{\lambda _2}^{}𝒂_{\lambda _1}^{}|\mathrm{𝟎}}{\widehat{0}|\widehat{a}_{\lambda _4}\widehat{a}_{\lambda _3}\widehat{a}_{\lambda _2}^{}\widehat{a}_{\lambda _1}^{}|\widehat{0}}}`$ (138)
etc. Such expressions are well defined since whenever their denominators vanish the whole term it containing vanishes as well. This fact is of crucial importance and shows that the perturbative expansions in both canonical and noncanonical frameworks contain terms of exactly the same type but differing by the numerical factors $`X_{\mathrm{}}^{\mathrm{}}`$.
Let us note that in the above calculation we have not used the explicit realization of the non-canonical algebra but only the algebra itself. At such a general level both the canonincal and non-canonical theories can be regarded as particular cases of a more general theory characterized by the algebra
$`[a_\lambda ,a_\lambda ^{}^{}]`$ $`=`$ $`\delta _{\lambda \lambda ^{}}I_\lambda ,`$ (139)
$`[a_\lambda ,a_\lambda ^{}]`$ $`=`$ $`0,`$ (140)
$`[a_\lambda ^{},a_\lambda ^{}^{}]`$ $`=`$ $`0,`$ (141)
$`[a_\lambda ^{},I_\lambda ^{}]`$ $`=`$ $`0,`$ (142)
$`[a_\lambda ,I_\lambda ^{}]`$ $`=`$ $`0.`$ (143)
The canonical choice, based on oscillators with classical parameter $`\lambda `$, is $`I_\lambda =\mathrm{𝟏}`$; the choice based on oscillators with eigenvalue $`\lambda `$ is $`I_\lambda =\mathrm{𝟏}_\lambda `$.
To proceed further and get more information as to the physical meaning of the non-canonical dynamics we have to make the analysis less general. First of all let us stick to the particular choice of $`\mathrm{𝟏}_\lambda `$ in terms of POV measures we have introduced earlier and assume that $`_\lambda \mathrm{𝟏}_\lambda =\mathrm{𝟏}`$. Second, let us take the vacuum state in the form (111). Under such assumptions we can explicitly compute the factors $`X_{\mathrm{}}^{\mathrm{}}`$:
$`X_{\lambda \lambda }^{01}`$ $`=`$ $`{\displaystyle \frac{\mathrm{𝟎}|𝒂_\lambda 𝒂_\lambda ^{}|\mathrm{𝟎}}{\widehat{0}|\widehat{a}_\lambda \widehat{a}_\lambda ^{}|\widehat{0}}}=\mathrm{𝟎}|\mathrm{𝟏}_\lambda |\mathrm{𝟎}=|O_\lambda |^2`$ (144)
$`X_{\lambda \lambda \lambda ^{}\lambda ^{}}^{0101}`$ $`=`$ $`{\displaystyle \frac{\mathrm{𝟎}|𝒂_\lambda 𝒂_\lambda ^{}𝒂_\lambda ^{}𝒂_\lambda ^{}^{}|\mathrm{𝟎}}{\widehat{0}|\widehat{a}_\lambda \widehat{a}_\lambda ^{}\widehat{a}_\lambda ^{}\widehat{a}_\lambda ^{}^{}|\widehat{0}}}=\mathrm{𝟎}|\mathrm{𝟏}_\lambda \mathrm{𝟏}_\lambda ^{}|\mathrm{𝟎}={\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}p_n\left(1{\displaystyle \frac{1}{n}}\right)|O_\lambda |^2|O_\lambda ^{}|^2=\left(11/n\right)|O_\lambda |^2|O_\lambda ^{}|^2`$ (145)
$`X_{\lambda \lambda \lambda \lambda }^{0101}`$ $`=`$ $`{\displaystyle \frac{\mathrm{𝟎}|𝒂_\lambda 𝒂_\lambda ^{}𝒂_\lambda 𝒂_\lambda ^{}|\mathrm{𝟎}}{\widehat{0}|\widehat{a}_\lambda \widehat{a}_\lambda ^{}\widehat{a}_\lambda \widehat{a}_\lambda ^{}|\widehat{0}}}=\mathrm{𝟎}|\mathrm{𝟏}_\lambda ^2|\mathrm{𝟎}=\left(11/n\right)|O_\lambda |^4+1/n|O_\lambda |^2`$ (146)
$`X_{\lambda \lambda \lambda \lambda }^{0011}`$ $`=`$ $`{\displaystyle \frac{\mathrm{𝟎}|𝒂_\lambda 𝒂_\lambda 𝒂_\lambda ^{}𝒂_\lambda ^{}|\mathrm{𝟎}}{\widehat{0}|\widehat{a}_\lambda \widehat{a}_\lambda \widehat{a}_\lambda ^{}\widehat{a}_\lambda ^{}|\widehat{0}}}=\mathrm{𝟎}|\mathrm{𝟏}_\lambda ^2|\mathrm{𝟎}=\left(11/n\right)|O_\lambda |^4+1/n|O_\lambda |^2`$ (147)
$`X_{\lambda \lambda ^{}\lambda \lambda ^{}}^{0011}`$ $`=`$ $`{\displaystyle \frac{\mathrm{𝟎}|𝒂_\lambda 𝒂_\lambda ^{}𝒂_\lambda ^{}𝒂_\lambda ^{}^{}|\mathrm{𝟎}}{\widehat{0}|\widehat{a}_\lambda \widehat{a}_\lambda ^{}\widehat{a}_\lambda ^{}\widehat{a}_\lambda ^{}^{}|\widehat{0}}}=\mathrm{𝟎}|\mathrm{𝟏}_\lambda \mathrm{𝟏}_\lambda ^{}|\mathrm{𝟎}=\left(11/n\right)|O_\lambda |^2|O_\lambda ^{}|^2.`$ (148)
What is interesting (and very characteristic) all these factors are smaller than 1 (this follows trivially from $`_\lambda |O_\lambda |^2=1`$). An analysis of higher order terms shows that this is a generic property of the non-canonical perturbation theory. The $`n`$ occuring in the average $`1/n`$ is the number-of-oscillators operator. For realistic vacua one may expect the average number of oscillators to be large and therefore $`1/n0`$. Taking a more general vacuum state we arrive at $`1/n_\lambda `$ instead of $`1/n`$, which means that the influence of the vacuum may vary from frequency to frequency (i.e. from point to point in space).
Now, if we ignore the corrections coming from $`1/n`$ we can see that the non-canonical perturbative expansion of the amplitude is the same we would have obtained by using the standard theory but with $`\widehat{g}_\lambda `$ regularized by $`\widehat{g}_\lambda O_\lambda \widehat{g}_\lambda `$. As a consequence there exists a natural cut-off in the theory which follows only from the fact that the vacuum wave function is square-integrable and therefore $`O_{s,\stackrel{}{\kappa }}0`$ for $`|\stackrel{}{\kappa }|\mathrm{}`$. It is quite remarkable that the same mechanism that eliminated the infinite vacuum energy plays a similar role in the other parts of the theory. As we shall see shortly the actual role of the vacuum can be analyzed only in a fully relativistic setting since then the charge and mass renormalization come into play.
However, for the sake of concreteness and to make some rough estimates of the effects involved let us take the trivial example where the vacuum amplitudes are constant, say, $`O_{s,\stackrel{}{\kappa }}=C`$ for all $`\omega _\stackrel{}{\kappa }<\omega _{\mathrm{max}}`$ and zero otherwise. The dynamics of the amplitude is then (up to $`1/n0`$) equivalent to the standard one with the cut-off at $`\omega _{\mathrm{max}}`$ and the coupling constant $`e/m`$ replaced by $`Ce/m`$. This implies that $`e`$ and $`m`$ have to be treated as bare parameters and the experimantal value is $`e_{\mathrm{ex}}/m_{\mathrm{ex}}=Ce/m`$. With this observation in mind we can discuss non-canonically other perturbative effects which are widely believed to be a consequence of the standard canonical quantization. Below we find it useful to make the bare parameter $`e/m`$ explicit in calculations and for this reason will use the notation $`\widehat{g}_\lambda =(e/m)\widehat{f}_\lambda `$.
### B Spontaneous emission of $`N`$ identical photons in $`N`$-th order perturbation theory
Assume the atomic transition is $`|A|B`$. Up to the $`N`$-th order perturbative correction we get
$`𝑵_\lambda ,B|\mathrm{\Psi }(t)`$ $`=`$ $`\left({\displaystyle \frac{e}{m}}\right)^N{\displaystyle \frac{1}{(i\mathrm{})^N}}{\displaystyle _0^t}𝑑t_N{\displaystyle _0^{t_N}}\mathrm{}{\displaystyle _0^{t_2}}𝑑t_1𝑵_\lambda |𝒂_\lambda ^N|\mathrm{𝟎}B|\widehat{f}_\lambda ^{}(t_N)\mathrm{}\widehat{f}_\lambda ^{}(t_1)|A`$ (149)
$`=`$ $`\left({\displaystyle \frac{e}{m}}\right)^N{\displaystyle \frac{𝑵_\lambda |𝒂_\lambda ^N|\mathrm{𝟎}}{\widehat{N}_\lambda |\widehat{a}_\lambda ^N|\widehat{0}}}\times (\mathrm{relevant}\mathrm{canonical}\mathrm{formula})`$ (150)
$`=`$ $`\left({\displaystyle \frac{e}{m}}\right)^N\sqrt{\mathrm{𝟎}|\mathrm{𝟏}_\lambda ^N|\mathrm{𝟎}}\times (\mathrm{relevant}\mathrm{canonical}\mathrm{formula})`$ (151)
where the “hatted” expressions are those from the canonical theory. As we can see the task is reduced to computing $`\mathrm{𝟎}|\mathrm{𝟏}_\lambda ^N|\mathrm{𝟎}`$. The general formula, valid for any $`N`$, is somewhat complicated and not very illuminating. The cases $`N=1`$ and $`N=2`$ we have already met. Making the simplifying choice of a very “flat” distribution of the vacuum modes (i.e. $`O_{s,\stackrel{}{\kappa }}=C`$ below some threshold) we get
$`\mathrm{𝟏}_\lambda ,B|\mathrm{\Psi }(t)`$ $`=`$ $`C\widehat{1}_\lambda ,B|\widehat{\mathrm{\Psi }}(t)`$ (153)
$`\mathrm{𝟐}_\lambda ,B|\mathrm{\Psi }(t)`$ $`=`$ $`\sqrt{C^4+1/n(C^2C^4)}\widehat{2}_\lambda ,B|\widehat{\mathrm{\Psi }}(t).`$ (154)
Identifying $`Ce/m`$ with $`e_{\mathrm{ex}}/m_{\mathrm{ex}}`$ we obtain the standard quantum-optics results provided
$`\sqrt{C^4+1/n(C^2C^4)}C^2.`$ (155)
For $`N=3`$
$`\mathrm{𝟎}|\mathrm{𝟏}_\lambda ^3|\mathrm{𝟎}`$ $`=`$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n^3}}\left(nC^2+3n(n1)C^4+(n^33n^2+2n)C^6\right)p_n`$ (156)
$`=`$ $`C^6+3C^4\left(1C^2\right)1/n+C^2\left(13C^2+2C^4\right)1/n^2.`$ (157)
As before the result becomes the standard one if the approximation
$`\sqrt{C^6+3C^4\left(1C^2\right)1/n+C^2\left(13C^2+2C^4\right)1/n^2}C^3`$ (158)
is justified. Let us note that in the standard cananonical quantum optics one considers vacuum consisting of an infinite number of oscillators and such subtleties are trivially ignored.
### C Spontaneous emission of two different photons
By the same argument as before the proportionality factor we need to estimate in second-order perturbation theory is
$`{\displaystyle \frac{\mathrm{𝟏}_\lambda ,\mathrm{𝟏}_\lambda ^{}|𝒂_\lambda ^{}𝒂_\lambda ^{}^{}|\mathrm{𝟎}}{\widehat{1}_\lambda ,\widehat{1}_\lambda ^{}|\widehat{a}_\lambda ^{}\widehat{a}_\lambda ^{}^{}|\widehat{0}}}`$ $`=`$ $`\sqrt{\mathrm{𝟎}|\mathrm{𝟏}_\lambda \mathrm{𝟏}_\lambda ^{}|\mathrm{𝟎}}=\sqrt{11/n}C^2C^2`$ (159)
so the result agrees with the canonical one. It is also in perfect agreement with the explicit calculations given in .
### D Stimulated emission
The last example we will discuss is the first order calculation of the transition amplitude $`|𝑵_\lambda ,A|(𝑵\mathbf{+}\mathrm{𝟏})_\lambda ,B`$. The appropriate proportionality coeeficient is
$`\sqrt{{\displaystyle \frac{\mathrm{𝟎}|\mathrm{𝟏}_\lambda ^{N+1}|\mathrm{𝟎}}{\mathrm{𝟎}|\mathrm{𝟏}_\lambda ^N|\mathrm{𝟎}}}}C`$ (160)
provided the values of $`1/n`$, $`1/n^2`$, etc. are sufficiently small. This is the correct result since $`C`$ gets absorbed into the renormalized coupling constant.
## VII Blackbody radiation
The final test of the new formalism we want to perform is the problem of blackbody radiation. Planck’s famous formula
$`\varrho (\omega )={\displaystyle \frac{\mathrm{}}{\pi ^2c^3}}{\displaystyle \frac{\omega ^3}{e^{\beta \mathrm{}\omega }1}}={\displaystyle \frac{\mathrm{}}{\pi ^2c^3}}\omega ^3\overline{N}_\omega ,`$ (161)
where $`\overline{N}_\omega `$ is the average number of excitations of an oscillator in inverse temperature $`\beta `$, is one of the first great sucesses of quantum radiation theory and marks the beginning of quantum mechanics. Contemporary measurements of $`\varrho (\omega )`$ performed by means of COBE (Cosmic Background Explorer) are in a very good agreement with the Planck law. The data have been carefully analyzed in the context of nonextensive statistics in search of possible deviations from extensivity. The result that comes out systematically is $`|q1|<10^4`$ where $`q`$ is the Tsallis parameter. The case $`q=1`$ corresponds to the exact Planck formula. If there are any corrections whatever, they must be quite small.
The derivation of the formula given by Einstein is based on the properties of spontaneous and stimulated emissions. As we have seen above there may occur differences with respect to the standard formalism but under reasonable assumptions they may be expected to be small.
Below we follow another standard route which consists basically of two steps. First, one counts the number of different wave vectors $`\stackrel{}{k}`$ such that $`c|\stackrel{}{k}|[\omega ,\omega +\mathrm{\Delta }\omega ]`$. Second, one associates with each such a vector an oscillator and counts the average number of its excitations assuming the Boltzmann-Gibbs probability distribution at temperature $`T`$ and chemical potential $`\mu =0`$. The latter assumption is justified by the fact that the number of excitations of the electromagnetic field is not conserved in atom-light interactions.
In the new model the situation is slightly different since there exists an additional conserved quantum number: The number of oscillators. As we have seen in previous calculations the Hamiltonian is block-diagonal with respect to $``$ but changes the number of excitations in each $`N`$-oscillator subspace of the direct sum. The state vectors at the multi-oscillator level are symmetric with respect to permutations of the oscillators and therefore the oscillators themselves have to be regarded as bosons whose number is conserved and their chemical potential is $`\mu 0`$. However, their excitations should be regarded as bosons with vanishing chemical potential.
The eigenvalues of $``$
$`E_{m,n}=m\mathrm{}\omega \left(n+{\displaystyle \frac{1}{2}}\right).`$ (162)
corresponding to the oscillator whose frequency is $`\omega `$ are parametrized by two natural numbers: $`m`$ (the number of oscillators) and $`n`$ (the number of excitations). Assuming the standard Boltzmann-Gibbs statistics we obtain the probabilities
$`p_{m,n}=Z^1e^{\beta [m\mathrm{}\omega (n+\frac{1}{2})m\mu ]}`$ (163)
where
$`Z`$ $`=`$ $`{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}e^{\beta m(\mu +\mathrm{}\omega /2)}{\displaystyle \frac{e^{\beta m\mathrm{}\omega }}{1e^{\beta m\mathrm{}\omega }}}.`$ (164)
The Lambert series
$`{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}a_m{\displaystyle \frac{x^m}{1x^m}}`$ (165)
is convergent for any $`x`$ if $`_{m=1}^{\mathrm{}}a_m`$ is convergent. Otherwise (165) converges for exactly those $`x`$ for which the power series $`_{m=1}^{\mathrm{}}a_mx^m`$ does. In (164) $`a_m=e^{\beta m(\mu +\mathrm{}\omega /2)}`$ and $`_{m=1}^{\mathrm{}}a_m<\mathrm{}`$ if $`\mu +\mathrm{}\omega /2<0`$. If $`\mu +\mathrm{}\omega /20`$ we still have convergence of (164) as long as $`_{m=1}^{\mathrm{}}e^{\beta m[\frac{1}{2}\mathrm{}\omega \mu ]}<\mathrm{}`$. The upper limit imposed on $`\mu `$ by the finiteness of $`Z`$ is therefore $`\mu <\frac{1}{2}\mathrm{}\omega `$. In what follows we assume that $`\mu `$ is $`\omega `$-independent and therefore $`\mu 0`$.
The appropriate average number of excitations is
$`\overline{n}_\omega `$ $`=`$ $`Z^1{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}mne^{\beta [m\mathrm{}\omega (n+\frac{1}{2})m\mu ]}`$ (166)
and the Planck formula is replaced by
$`\varrho _{\mathrm{new}}(\omega )`$ $`=`$ $`{\displaystyle \frac{\mathrm{}}{\pi ^2c^3}}\omega ^3\overline{n}_\omega .`$ (167)
It is easy to show that $`\varrho _{\mathrm{new}}(\omega )`$ tends to the Planck distribution with $`\mu \mathrm{}`$. To see this consider a more general series
$`Z^1{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}mnq_me^{\beta m\mathrm{}\omega (n+\frac{1}{2})}`$ (168)
where $`Z`$ is the normalization factor and $`_{m=1}^{\mathrm{}}q_m<\mathrm{}`$. If $`q_1=1`$ and $`q_m=0`$ for $`m>1`$ then (168) is just the exact Planckian formula. Factoring out $`e^{\beta |\mu |}`$ in both the numerator and the denominator of $`\overline{n}_\omega `$ we obtain $`q_1=1`$ and $`q_m=e^{\beta |\mu |(m1)}`$ for $`m>1`$. For $`|\mu |\mathrm{}`$ all $`q_m`$, for $`m>1`$, vanish and the limiting distribution is Planckian.
This proves that an experimental agreement with the ordinary Planck’s $`\varrho (\omega )`$ cannot rule out our modification but can, at most, set a lower bound on an admissible value of $`|\mu |`$. However, assuming that $`\mu `$ has some finite and fixed value it should be in principle measurable. The plots show that the modifications become visible around $`\mu 3k_BT`$. Assuming that the chemical potential is temperature independent, say $`\mu =k_BT_0`$, we obtain a kind of critical temperature $`T_{\mathrm{critical}}T_0/3`$ above which the ratio $`\mu /(k_BT)`$ is small enough to make the modifications of the distribution observable. For $`T<T_{\mathrm{critical}}`$ the distribution should be given by the Planck law; for $`T>T_{\mathrm{critical}}`$ the distribution should approach the $`\mu =0`$ distribution, i.e. this would be a Planck-type curve but with the maximum lowered and shifted towards higher energies.
Fig. 1 shows the plots of $`\varrho _{\mathrm{new}}(\omega )`$ for $`\mu =0`$ (lower dotted), $`\mu =0.8k_BT`$ (upper dotted), and $`\mu =10k_BT`$ (solid). The thick dashed curve is the Planck distribution. The curve obtained for $`\mu =10k_BT`$ is indistinguishable from the Planck distribution. The plot does not change if one takes $`\mu <10k_BT`$ and differences are not visible even if one plots the distributions in logarithmic scales (not shown here). This is a numerical proof that the distribution we have obtained on the basis of the modified quantization tends very quickly to the Planck one as $`\mu \mathrm{}`$. It is instructive to compare the modification we have predicted with those arising from nonextensive statistics. The two thin dashed lines represent Tsallis distributions resulting from nonextensive formalism for $`q=0.95`$ (lower) and $`q=1.05`$ (upper). The modifications we have derived are therefore qualitatively different from those resulting from Tsallis statistics.
## VIII Conclusions
“A theory that is as spectacularly successful as quantum electrodynamics has to be more or less correct, although we may not be formulating it in just the right way” . The above quotation from Weinberg could serve as a motto opening our paper. The main idea we have tried to advocate was that the standard canonical quantization procedure is, in certain sense, too classical to be good.
The reasons for such a choice of quantization could be both historical and sociological and may be rooted in the fact that the idea of quantizing the field was formulated before the real development of modern quantum mechanics. In oscillations of a simple pendulum it may be justified to treat $`\omega `$ as an external parameter defining the system (via, say, the length of the pendulum). But oscillations of the electromagnetic field do not seem to have such a “mechanical” origin and it is more natural to think of the spectrum of frequencies as eigenvalues of some Hamiltonian. That is exactly what happens with other quantum wave equations.
We have defined the quantum electromagnetic field as an oscillator that can exist in a superposition of different frequencies (or, rather, wave vectors). This should not be confused with the classical superpositions of frequencies created by, say, a guitar string. The superpositions we have in mind dissappear at the classical level.
Once one accepts this viewpoint it becomes clear how to quantize the field at the level of a single oscillator. We do not need many oscillators to perform the field quantization. But there is no reason to believe that all the possible fields can be described by the same single oscillator. And even more: We know that the structure of the one-oscillator Hilbert space is not rich enough to describe multi-particle entangled states and there is no doubt that such states are physical. The next step, performed already after the quantization, is to consider fields consisting of 1, 2, 3 and more oscillators, and even existing in superpositions of different numbers of them. The resulting structure is in many respects analogous to the Fock space so that the procedure can be (although somewhat misleadingly) referred to as “second quantization”. What is essential we do not need the vacuum state understood as the unique cyclic vector of the GNS construction.
On the other hand, there exist vacuum states. These are all the states describing ground states of the oscillators. They correspond to concrete finite average values of energy. A general vacuum state is therefore a superposition of different eigenstates of a free Hamiltonian and is not, in itself, an eigenstate of the Hamiltonian.
One technical assumption we have made is the resolution of identity property of the non-CCR algebra. This assumption is clearly satisfied at the one-oscillator level. The remaining assumptions are standard. The system is described by laws of ordinary quantum mechanics so that to compute concrete problems we can use standard methods. Perturbation theory leads to structures we know from the standard Feynmann diagrams. The blackbody radiation is calculated by means of the standard Boltzmann-Gibbs statistics.
There are still several unexplored possibilities. To give an example, we have seen in explicit calculations that differences between the canonical and non-canonical formalisms consist of the factors occuring in perturbative expansions of amplitudes. These factors explicitly depend on the choice of the RHS of the non-CCR algebra. As such, they point into possible experiments testing directly the algebra of canonical commutation relations. The meaning of such tests is, at the present stage, obscured by the lack of proper understanding of the role of renormalisation. Actually, one should not expect here precise results since even the nonrelativistic canonical quantum electrodynamics is non-renormalizable. The simple illustration we have used, namely the one with the cut-off and flat vacuum, shows that the new theory is not that far from the canonical one as one might expect.
Let us close these remarks with another quotation: “Present quantum electrodynamics contains many very important ‘elements of truth’, but also some clear ‘elements of nonsense’. Because of the divergences and ambiguities, there is general agreement that a rather deep modification of the theory is needed, but in some forty years of theoretical work, nobody has seen how to disentengle the truth from the nonsense. In such a situation, one needs more experimental evidence, but during that same forty years we have found no clues from the laboratory as to what specific features of QED might be modified. Even worse, in the absence of any alternative theory whose predictions differ from those of QED in known ways, we have no criterion telling us which experiments would be relevant ones to try.
It seems useful, then, to examine the various disturbing features of QED, which give rise to mathematical or conceptual difficulties, to ask whether present empirical evidence demands their presence, and to explore the consequences of the modified (although perhaps rather crude and incomplete) theories in which these features are removed. Any difference between the predictions of QED and some alternative theory, corresponds to an experiment which might distinguish between them; if it appears untried but feasible, then we have the opportunity to subject QED to a new test in which we know just what to look for, and which we would be very unlikely to think of without the alternative theory. For this purpose, the alternative theory need not be worked out as completely as QED; it is sufficient if we know in what way their predictions will differ in the area of interest. Nor does the alternative theory need to be free of defects in all other respects; for if experiment should show that it contains just a single ‘element of truth’ that is not in QED, then the alternative theory will have served its purpose; we would have the long-missing clue showing in what way QED must be modified, and electrodynamics (and, I suspect, much more of theoretical physics along with it) could get moving again “ .
###### Acknowledgements.
This work was done partly during my stays in Antwerp University, UIA, and Arnold Sommerfeld Institute in Clausthal. I gratefully acknowledge a support from the Alexander von Humboldt Foundation and the Polish-Flemish grant No. 007. I’m indebted to Prof. Iwo Białynicki-Birula, Robert Alicki, Jan Naudts, Maciej Kuna, and Wolfgang Luecke for critical comments, and Paweł Syty for a stimulating discussion on small $`\omega `$’s. I express my gratitude to prof. Heinz-Dietrich Doebner for various help.
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# Radio and X-ray Signatures of Merging Neutron Stars
## 1 Introduction
The gravitational wave-induced merger of binary neutron stars has evoked considerable interest in recent years due to their importance as a source of gravitational radiation (Thorne 1987 and references therein) and potentially also gamma-ray bursts (Blinnikov et al 1984; Paczynski 1986; Eichler et al 1989; Paczynski 1991). The goal of identifying electromagnetic signatures of the merger events is an important one, even if such manifestations are not gamma-ray bursts. Given the considerable information processing required to infer the presence of a gravitational wave burst (Cutler et al 1993), the presence of another signature will be invaluable.
In this paper we examine the magnetospheric interactions in merging neutron star binary systems and describe their pre- and post-merger signatures. In particular, we consider systems containing one low field $`B_r10^{911}G`$, rapidly spinning ($`P1100ms`$) recycled pulsar and one high field ($`B_m10^{1215}G`$), slowly spinning ($`P101000s`$) non-recycled pulsar, as expected on both empirical and evolutionary grounds. We shall examine how energy is extracted from the spin and orbital motion of the pulsar and in what manner it is radiated. Aspects of this calculation have been considered before by Vietri (1996), who considered magnetospherically induced gamma-ray bursts, and Lipunov & Panchenko (1996), who considered the the far-field dipolar and quadrupolar configurations of a dipole merging with a superconducting sphere. Our default estimates will be for systems in which the high field pulsar has a field $`10^{15}`$ G (sometimes called a magnetar), which has the potential for the strongest signal. Recent work suggests that such pulsars may constitute $`10\%`$ of the young pulsar population (Kulkarni & Thompson 1998).
In section 2 we will describe the magnetospheric interactions that remove energy from the orbit and which allow it to emerge in electromagnetic form. This section will draw heavily on concepts from pulsar electrodynamics and also the field of satellite-magnetosphere electrodynamics, such as in the Io-Jupiter system. One of the primary results is that much of the energy is released as a pair plasma into the magnetosphere. Section 3 describes the physical state and evolution of this plasma, drawing on concepts developed to describe Soft Gamma Repeaters and section 4 reviews the state of the observations appropriate to this phenomenon.
## 2 The Extraction of Spin and Orbital Energy
High magnetic field pulsars spin down rapidly, so that we consider the high field pulsar to be essentially non-rotating. A corollary to this is that the light cylinder radius for the magnetar magnetosphere (or ‘magnetarsphere’!) is $`c/\mathrm{\Omega }5\times 10^{11}\mathrm{cm}(P/100\mathrm{s})`$, so that all of our subsequent discussion concerns processes occurring deep within the closed region of this magnetosphere. This will hold true right up to the point of merger as tidal interactions cannot enforce synchronisation in a coalescing neutron star binary (Bildsten & Cutler 1992).
The extraction of energy from the pulsar spin and orbital motion is driven by how the strongly conducting neutron star interacts with the external magnetic field of the magnetar. As a model problem we consider perfectly conducting sphere moving through an externally imposed uniform magnetic field $`𝐁_0`$ with velocity $`𝐯`$ and rotating with angular velocity $`𝛀`$. Motion of a conducting sphere through magnetic field is possible only if the resistivity of the sphere (neutron star) is nonzero. But the neutron star crust is virtually a perfect conductor: the magnetic diffusion times for neutron stars are very long - in fact comparable to the age of the universe. We argue that the required resistivity is due to the dissipation of the induced magnetospheric currents far from the neutron star surface. This is analogous to the case of isolated pulsars, where currents are dissipated in the pulsar wind-ISM shocks more than $`10^9`$ light cylinder radii away.
The electrodynamics of the low-field pulsar interaction with the magnetar magnetic field is considered in appendix A. Qualitatively, this interaction has several important ingredients. The conducting neutron star excludes the external field from its interior. The induced magnetic field has a dipole structure with the magnetic dipole directed against the external field. The resultant magnetic field is
$$𝐁_{tot}=𝐁_0+\frac{R^3}{2r^3}𝐁_0\frac{3R^3(𝐁_\mathrm{𝟎}𝐫)𝐫}{2r^5}$$
(1)
At the surface the total magnetic field has only a tangential component, inside the star the magnetic field is zero. The orbital motion of the neutron star with respect to the external field will induce surface charges with a dipole structure and surface charge density
$$\sigma _{orb}=\frac{1}{4\pi cR_0}\left(𝐁_0[𝐑_0\times 𝐯]\right)$$
(2)
where $`R_0`$ is the radius of the neutron star. These surface charges will also produce electric fields which will have a component along the total magnetic field. If the neutron star orbits in a vacuum, this electric field will accelerate charges from the surface or surrounding region to relativistic energies. If the internal magnetic field of the neutron star is exactly zero, the resulting structure of the electric field will be of the ”outer gap type” (Chen, Ho & Ruderman 1986) - a region in the magnetosphere with $`E_{}`$ which does not intersect the surface of the neutron star. However, there is likely to be some small component of radial magnetic field at the surface. This could result either from whatever intrinsic field the recycled neutron star possesses or from a second-order induced field resulting from rotation of the star (see below and appendix A). In this case the electric field will draw charges from the surface.
Similarly to orbital motion, the rotation of the neutron star will produce a surface charge density
$$\sigma _{rot}=\frac{3R_0\mathrm{\Omega }B_0\mathrm{sin}\psi \mathrm{sin}\theta }{8\pi c}$$
(3)
where $`\mathrm{sin}\psi =\mathrm{cos}\mathrm{\Omega }t\mathrm{cos}\theta \mathrm{sin}\alpha \mathrm{sin}\theta \mathrm{cos}\alpha `$ and $`𝛀=\mathrm{\Omega }(\mathrm{sin}\alpha ,0,\mathrm{cos}\alpha )`$, i.e. $`\psi `$ is the polar angle in the frame aligned with $`𝛀`$ and rotating with the neutron star. This charge density is stationary in the frame of the neutron star while in the laboratory frame it yields an additional surface current $`𝐣=\sigma _{rot}𝛀\times 𝐫`$. The magnetic field due to this current is of order $`(R\mathrm{\Omega }/c)^2`$ smaller than the external field $`𝐁_0`$, but has a radial component at the surface.
Similarly to the case of the aligned rotator studied by Goldreich & Julian (1969), the strong electric field produced by surface charges will accelerate charges in an attempt to short out the component of the electric field that lies parallel to the magnetic field. The typical densities of the primary beam will be $`n_{GJ}\mathrm{\Omega }B_0/2\pi ec`$ for acceleration by $`\sigma _{rot}`$ and $`vB_0/ecR`$ for acceleration by $`\sigma _{orb}`$. After being accelerated to sufficient energies ($`\gamma 10^6`$) the initial primaries produce curvature photons and a dense population of secondary electron-positron pairs that will screen the induced electric field. This mechanism of energy extraction is essentially the same as in the classical pulsar case with a couple of small but important differences. The first is that, unlike the case of the pulsar, the near field energy density is dominated by the plasma, rather than Poynting flux (see appendix A). Secondly, the field configuration defined by (1) contains no closed magnetosphere. In the traditional pulsar case, the ‘working surface’ of the energy extraction is limited to the polar cap, a fraction $`(r\mathrm{\Omega }/c)^2`$ of the the total surface area, which is linked to the open field lines. In the case under discussion here, the polar cap effectively encompasses the entire star.
The energy extracted by accelerating primary particles is limited by the maximum energy that primary particles can reach:
$$L4\pi R^2n_{GJ}\gamma _{max}m_ec^33.1\times 10^{36}\mathrm{ergs}\mathrm{s}^1.$$
(4)
However, the energy extraction from the orbital motion is likely to be significantly more efficient than implied by (4). Once the pair production cascade has loaded the external field lines with plasma, the spiraling neutron star emits Alfvén waves along the external magnetic field (Drell, Foley & Ruderman 1965; Barnett & Olbert 1986; Wright & Southwood 1987), in much same way as Io interacting with Jupiter or various artificial satellites in the earth’s magnetosphere. In this case, the pair production front acts as a surface of finite resistivity, allowing the neutron star to ‘cross field lines’<sup>1</sup><sup>1</sup>1Even if the resistivity were considerably lower, a similar level of energy extraction would occur via the screw-instability of strongly wound magnetic field configurations (Low 1986; Aly 1991; Volwer, van Oss & Kuijper 1993; Gruzinov 1999), given only the assumption of sufficient ambient plasma to justify the force-free approximation. We assume that these waves are dissipated in the magnetar magnetosphere by non-linear damping mechanisms similar to those invoked by Soft Gamma Repeater models (e.g. Thompson & Duncan 1995). Thus we shall assume that the bulk of the energy extracted from the orbit is deposited into the magnetospheric pair plasma, and is of order (Drell et al 1965)
$`L_{orb}`$ $`4\pi R^2B_m^2\left({\displaystyle \frac{R}{a}}\right)^6{\displaystyle \frac{v^2}{c}}`$
$``$ $`7.4\times 10^{45}\mathrm{ergs}\mathrm{s}^1\left({\displaystyle \frac{B_m}{10^{15}\mathrm{G}}}\right)^2\left({\displaystyle \frac{a}{10^7\mathrm{cm}}}\right)^7.`$ (5)
Additional sources of energy are the Poynting losses due to the motion of the induced dipole (Lipunov & Panchenko 1996) and the time varying component of the induced magnetic fields (see appendix A), though the corresponding luminosities are much smaller than that given by Eq. (5). Poynting fluxes will be in a form of low frequency electromagnetic waves. Unlike the equivalent situation for pulsars, where the density of the secondary plasma is low, these low frequency electromagnetic waves may not be able to propagate through the dense secondary plasma present in the ”magnetarsphere” - they will convert their energy into plasma (Asseo et al. 1978). Thus, most of the energy lost by the neutron star will be converted into plasma in the near zone and later radiated - this is in contrast to normal pulsars where most of the losses within the light cylinder are due to the Poynting flux.
Our situation also differs somewhat from that considered by Vietri (1996), who addressed the problem of the merger of two high-field pulsars. The consequently large radii of field line curvature implied screening was ineffective and allowed efficient acceleration of particles and high energy emission. The fundamental difference in our case is that plasma screening occurs close to the low field neutron star, where the radius of curvature is smaller. The result is efficient generation of pair plasma and effective screening of parallel electric fields. The pair plasma will then mediate the dominant energy extraction by Alfvén wave emission.
### 2.1 Radio Emission
In normal pulsars, the acceleration of particles by electric fields at the surface yields coherent radiation observed as radio emission. Thus, we might hope for similar signals in this instance. The lack of a complete theory of pulsar radio emission forces us to adopt a simple parameterisation based on what we know about pulsars. We expect the radio emission to be associated with the primary beam only, whose luminosity is given by Eq. (4). We shall adopt an efficiency of $`ϵ0.1`$ for the conversion of primary beam energy to radio luminosity, based on the radio efficiencies of pulsars (see Taylor, Manchester & Lyne 1993), assuming the pulsar beam luminosity is $`10^3`$ of the spin-down luminosity (Kennel & Coroniti 1984). Then an optimistic estimate for the radio flux at 400 MHz (chosen because it is at this frequency that the pulsar fluxes are best estimated) is
$$F_\nu 2.1\mathrm{mJy}\frac{ϵ}{0.1}\left(\frac{D}{100Mpc}\right)^2B_{15}^{2/3}a_7^{5/2}.$$
(6)
This is within the range of the larger radio telescopes operating today, although somewhat less than the sensitivities of current radio transient searches.
There are several complications that may preclude generation of radio emission. If the neutron star is moving through a pre-existing plasma generated by the previous orbital cycles the electric gaps may be quenched, there will be no need to accelerate further particles and the beam luminosity may drop to zero. In addition, the formation of positronium in the magnetic fields exceeding $`4\times 10^{12}G`$ (Usov & Melrose 1996, Arons 1998) may also quench the radio emission.
Furthermore, the generated radio emission may be absorbed in the magnetarsphere. We expect that nonresonant Thomson scattering of the low frequency ($`\nu <<\nu _B`$) radio emission will not be important due to the the strong suppression ($`\sigma =\sigma _T(\nu /\nu _B)^2`$) of the scattering cross-section by the magnetic field at low frequencies. Resonant cyclotron absorption may be important in the outer regions of the magnetarsphere where the cyclotron frequency becomes comparable to the radio wave frequency ($`a10^{10}cm`$). Nevertheless, such absorption does not occur in the pulsar case, so we may reasonably expect the radio emission to escape the magnetarsphere.
Thus, the first electromagnetic signature we anticipate from a realistic merging neutron star binary is a coherent radio burst, emitted $``$ seconds before the gravitational wave burst. However, the effects of interstellar dispersion can cause delays of hours, possibly allowing for radio follow-up observations at low frequencies (Palmer 1993; Lipunova, Panchenko & Lipunov 1997).
## 3 Evolution of the Magnetospheric Pair Plasma
Most of the energy liberated by the strong electric fields of section 2 is not radiated, but is rather released into the magnetosphere of the slowly-rotating magnetar in the form of Alfvén waves and a dense pair plasma. The energy release (5) is a significant fraction of the local magnetic energy density. In such a case, a wind, driven either by hydromagnetic or plasma pressure, is likely to result (Paczynski 1986, 1990, Melia & Fatuzzo 1995; Katz 1996) while some will remain trapped, in a fashion similar to that of the Soft Gamma Repeater picture of a magnetically confined pair plasma (Thompson & Duncan 1995). We envisage that the plasma released into regions of decreasing field strength powers the wind while plasma released into regions of increasing field strength will be trapped. Figure 1 shows a schematic version of our scenario.
Let us consider first the case of the wind. A release of energy at the rate given by equation (5) results in a compactness parameter $`\eta =L/ac10^7B_{15}^2a_7`$. Thus, this is the same situation envisaged in cosmological models for gamma-ray bursts (Paczynski 1986; Goodman 1986), wherein the release of a large quantity of pure energy within a small volume leads to a relativistically expanding fireball. The energy release during the neutron star inspiral will drive a relativistically expanding wind of pairs and photons. Thermal and statistical equilibrium between photons and pair plasma is maintained during the expansion by pair production and comptonization (Cavallo & Rees 1978) until the comoving temperature drops to $`T3\times 10^8`$ K and pair production can no longer maintain the necessary electron scattering optical depth. At this point the radiation escapes, with an approximately thermal spectrum. However, the relativistic boost increases the observed temperature by a factor $`\gamma `$, the original Lorentz factor of the fireball and reduces the observed burst time by the same factor. Thus, the observed energetic and temporal properties of the wind emission may be approximately described by thermal emission at the appropriate initial temperatures and timescales, despite the fact that the true photosphere is on scales much larger than the original volume. Hence, we shall estimate the observed flux in this case as
$$F_{wind}\frac{\alpha L}{4\pi a^2}3\times 10^{30}\alpha \mathrm{ergs}\mathrm{cm}^2\mathrm{s}^1B_{15}^2a_7^9,$$
(7)
($`\alpha `$ is the fraction of the energy release lost in the wind) yielding effective temperatures just before merger $`1.5\mathrm{MeV}B_{15}^{1/2}`$.
The case of the trapped plasma is somewhat more subtle. This plasma is very optically thick (Svensson 1987; Thompson & Duncan 1995) and the inspiral time $`0.4\mathrm{s}a_7^4`$ is short. Hence, very little of the total energy contained in the magnetosphere is radiated in this time. The emission that does occur is dominated by the region just above the surface of the magnetar, where the strong magnetic field decreases the electron scattering cross-section and thereby promotes a larger photon flux. At late times, the plasma temperature is high enough that ablation of material from the surface of the star is likely to provide an Eddington limit
$$F_{edd}4.14\times 10^{26}\mathrm{ergs}\mathrm{cm}^2\mathrm{s}^1B_{15}a_7^3.$$
(8)
In the case of Eddington limited cooling, we see that the emission actually gets softer as the inspiral proceeds (the opposite of the contribution (7) from the relativistic wind) because the plasma temperature increases and acts to negate the magnetic suppression of the electron scattering cross-section.
Thus, our strongest prediction is the presence of an X-ray precursor to the neutron star merger. This precursor should be dominated by approximately thermal emission from the wind component, which progressively hardens as the binary approaches merger. For magnetars ($`B_{15}1`$), this can approach energies $``$ 1.5 MeV while binaries containing a normal pulsar ($`B_{15}10^3`$) will only get as hard as $`50`$ keV. In some cases these objects may be accompanied by softer ($``$ 1 keV) components associated with the cooling of the trapped plasma. The observed flux will be dominated by the harder wind component, with flux levels $`3\times 10^9\mathrm{ergs}\mathrm{cm}^2\mathrm{s}^1B_{15}^2a_7^7`$ for a source at 100 Mpc (the distance scale for which we expect a few neutron star mergers per year).
### 3.1 The Ultimate Fate of the Pair Plasma
Much of the energy released by the processes in section 2 is retained in the magnetospheric plasma for timescales longer than the inspiral time, i.e. the merger event will occur surrounded by a significant magnetospheric plasma. This energy totals about $`E_{\mathrm{plasma}}2\times 10^{47}\mathrm{ergs}B_{15}^2`$. Vietri (1996) proposed that this energy, most of which is released on the last few orbits, could power a Gamma-Ray Burst as it escapes it’s magnetic confinement. Our estimate of the energy release is somewhat smaller than his (strictly speaking he calculated the maximum energy a magnetosphere could contain, rather than the energy release itself). Recent determinations of the distances (Metzger et al 1997; Kulkarni et al 1998; Djorgovski et al 1998) to burst events suggest that much larger energy releases are required to explain many gamma-ray bursts (modulo beaming considerations).
A more intriguing possibility occurs if it is the magnetar which is disrupted to form a rapidly rotating disk around the compact merger remnant (as might be expected from the non-recycled and thus presumably lighter object). If the magnetic field footpoints remain tied to the disrupted material the magnetosphere is forced to co-rotate with the disk and the corotation radius must move rapidly inwards, converting closed field lines to open ones and ejecting the magnetospheric plasma. This would allow the plasma to tap the much larger reservoir of disk rotational energy to power the gamma-ray burst (as in many other gamma ray burst models. See Hartmann 1996 for a review) and could also account for significant collimation of the outflow. Furthermore, the approximate equipartition between plasma energy density and magnetic field energy is appropriate for the formation of an episodic jet (Ouyed & Pudritz 1997), which may contribute to burst temporal variability. Finally, the magnetospheric origin of the pair plasma would also avoid the baryonic loading problem encountered by mechanisms which propose to generate the pairs by neutrino annihilation close to the merger product (Janka & Ruffert 1996; Ruffert & Janka 1998).
## 4 Discussion
### 4.1 Observations
Our results predict the appearance of X-ray and Radio transients as precursors to gravitational wave bursts and possibly also Gamma-Ray Bursts.
Motivated by the search for X-ray counterparts to GRB, Gotthelf, Hamilton & Helfand (1996) and Greiner (1999) have searched the Einstein & ROSAT databases, respectively, for brief X-ray transients. The energy ranges searched are somewhat softer ($`0.14`$ keV) than we would predict for the peak of the energy distribution. Nevertheless, both searches found classes of transients of possible astronomical origin in sufficient abundance to encompass any reasonable estimate of the event rate (Phinney 1991; Narayan, Piran & Shemi 1991; Lipunov et al 1995, Arzoumanian, Cordes & Wasserman 1999). Searches for untriggered bursts in the BATSE catalogue (Kommers et al 1997) proved even more interesting, revealing a class of bursts restricted to the 25-50 keV channel. Most of these events are consistent with being intensity fluctuations in Cygnus X-1, but the remaining 10% are consistent with the expectations of our model.
If some neutron star mergers do yield GRB, then our results may provide an explanation for the subset of GRB discovered to show X-ray precursors (Murakami et al 1992; Castro-Tirado et al 1993; in ’T Zand et al 1999). These events show the characteristic soft-to-hard spectral evolution we anticipate, although the spectrum in well-studied cases such as GB900126 (Murakami et al 1991), appears somewhat softer ($`1.5`$ keV) than we would expect. The flux ($`2\times 10^9\mathrm{ergs}\mathrm{cm}^2\mathrm{s}^1`$), however, is appropriate, suggesting perhaps that an analysis more sophisticated than the black-body assumption is required.
Perhaps the best candidate for our model is the unusual transient GB900129 observed by Ginga (Strohmayer et al 1995), which yielded a thermal bremstrahlung temperature $`20`$ keV and duration 5-10 s. Strohmayer et al note the similarity to the SGR spectral characteristics, which agrees with the magnetospheric origins in our model as well. Figure 2 shows the comparison of the energetics of various observed transients with our models.
Radio transients associated with GRB are a subject of growing interest and several searches (e.g. FLIRT and STARE) are ongoing. However, radio searches for brief transients are particularly bedevilled by terrestrial interference. Most limits lie in the 10-100 kJy range at 76 MHz (FLIRT; Balsano 1999) and 611 MHz (STARE; Katz et al 1998). These don’t particularly constrain our model, which anticipates signals $``$ mJy-Jy. One radio transient uncovered by FLIRT does deserve mention. FLIRT (Balsano 1999) located a radio transient apparantly associated with GRB 980329. The transient is unique in the database and occurred within 50 s of the burst. The transient showed evidence for dispersion, with a DM$`66\mathrm{p}\mathrm{c}\mathrm{cm}^3`$. All these argue that the association is real. However, the signal was very narrow band, indicative of terrestrial interference. If one does choose to interpret this ambiguous transient as a real association, the dispersion measure would rule out a truly cosmological burst. Furthermore, the $``$ kJy flux would suggest distances $``$1 Mpc based on our luminosity estimates. All these would argue against the suggestion that the event occurred at very high redshift (Fruchter 1999) and the red optical transient would most likely arise from extinction (Reichart et al 1999).
### 4.2 Binaries with Black Holes
We have not yet discussed the signatures of mergers associated with binaries in which one of the components is a black hole rather than a neutron star, although such binaries may outnumber the double neutron star binaries (e.g. Bethe & Brown 1998). If the black hole is formed from a strongly magnetized object, then only open field lines remain. Thus, the inspiral of a neutron star through this magnetosphere will still generate the relativistic wind and it’s associated X-ray signature<sup>2</sup><sup>2</sup>2This precursor may also be cut off before the merger if the open field lines are restricted to only a small polar cap, i.e. the neutron star may eventually orbit in a field free zone., but there will be no trapped magnetospheric plasma. We would then expect to see the same X-ray and radio precursor to the event, but no soft X-ray component to the precursor or any post merger signature associated with the magnetospheric pair plasma. For binaries in which the low field component is a black hole, the effective resistivity of the event horizon (Thorne, Price & Macdonald 1986) is considerably larger than that for a neutron star crust. Consequently the distortion of the magnetic field due to the orbital motion is much smaller and the energy extraction in observable energy is similarly reduced. Furthermore, one cannot extract charged particles from the event horizon, although ‘outer gap’ acceleration is still feasible. We expect such mergers to be (electromagnetically) much quieter.
## 5 Conclusion
If neutron star mergers are not associated with GRB, then any additional electromagnetic signatures will be invaluable when the search begins in earnest for the gravitational wave signal. Li & Paczynski (1999) have suggested one such signature; namely a post-merger mini-supernova powered by radioactive decay of disrupted neutron star material. We have demonstrated the possibility of additional precursor signals in the radio and X-ray regimes, driven by the magnetospheric interactions of the neutron star and their magnetic fields. Our results differ somewhat from those of Vietri (1996) who considered a related model. We ascribe this to the much more localized interaction in our scenario, the result of a more realistic choice of parameters, and to our more complete description of the electrodynamics of the accelerated plasma.
To conclude we re-iterate the properties of what we would consider a prime candidate for an electromagnetic counterpart to a neutron star merger. Estimates of the merger rate suggest that the events typically observed would be at distances $`100`$ Mpc, suggesting X-ray fluxes $`3\times 10^9\mathrm{ergs}\mathrm{cm}^2\mathrm{s}^1`$ with effective temperatures progressing upwards through the 10-100 keV range preceding the gamma-ray event on timescales of order seconds or less. Associated radio fluxes could be as much as $`5`$ Jy at this distance, although the ability of the radio waves to propagate in the late-time plasma shroud is rather uncertain. The coincidence of the radio signal could be influenced by dispersion in both the host galaxy and ours. Dispersion in the inter-Galactic medium will be of the order of $`1\mathrm{c}\mathrm{m}^2\mathrm{pc}(D/100\mathrm{M}\mathrm{p}\mathrm{c})`$ for an ionized IGM mass fraction $`10^2`$ of the critical density and thus is unlikely to contribute significantly for any detectable events. There are also several possible signatures of the merger event itself, depending on how the orbital and binding energies of the binary and components is disbursed between remnant and ejecta.
We thank Steve Thorsett, Vicky Kaspii & Jackie Hewitt for information regarding the FLIRT and STARE radio transient programs and Vladimir Lipunov and Andrei Gruzinov for discussions.
## Appendix A Electromagnetic interaction of a neutron star with external magnetic field
Our problem concerns the electrodynamics of a low field neutron star orbiting in the field of a high field neutron star. Although the intrinsic field of the neutron star will be important, let us begin with the model problem of a spinning, conducting sphere in a uniform, externally imposed field.
Consider an unmagnetized conducting sphere in an external homogeneous magnetic field $`𝐁_\mathrm{𝟎}`$. The external field will induce surface currents
$$𝐠=\frac{c}{4\pi }B_0\mathrm{sin}\theta 𝐞_\varphi $$
(9)
flowing in the azimuthal direction $`𝐞_\varphi `$ about the $`𝐁_\mathrm{𝟎}`$ axis. These currents will induce magnetic field with a dipole structure $`\mu =𝐁_\mathrm{𝟎}R_0^3/2`$ so that the total field is
$$𝐁_{tot}=𝐁_\mathrm{𝟎}+\frac{R_{0}^{}{}_{}{}^{3}𝐁_\mathrm{𝟎}}{2r^3}\frac{3R_{0}^{}{}_{}{}^{3}(𝐁_\mathrm{𝟎}𝐫)𝐫}{2r^5}$$
(10)
At the surface the radial component of the total magnetic field is zero. Consider now the same sphere, moving with velocity $`𝐯`$ and rotating with angular velocity $`\mathrm{\Omega }`$. Up to the relativistic correction of the order $`v^2/c^2`$ the magnetic field is the same in the star’s rest frame (which is moving and rotating with respect to the laboratory frame). In the star’s rest frame the electric field is a sum of electric fields due to the uniform motion and rotation. The electric field due to the rectilinear motion is uniform, given by $`𝐄_{orb}=\frac{1}{c}\left[𝐯\times 𝐁_0\right]`$, while the electric field due to the rotation is $`𝐄_{rot}=\frac{1}{c}\left[\left[𝛀\times 𝐫\right]\times 𝐁_{tot}\right]`$. $`𝐄_{rot}`$ has radial and meridional components:
$`E_{rot,r}={\displaystyle \frac{\left(2r^3+R_{0}^{}{}_{}{}^{3}\right)}{2cr^4}}\left(r^2(𝐁_{\mathrm{𝐭𝐨𝐭}}𝛀)(𝐁_{\mathrm{𝐭𝐨𝐭}}𝐫)(𝛀𝐫)\right)`$
$$E_{rot,\psi }=\frac{\left(r^3R_{0}^{}{}_{}{}^{3}\right)}{cr^5}(𝐁_{\mathrm{𝐭𝐨𝐭}}𝐫)\left(\left(r^2𝛀\right)+(𝛀𝐫)𝐫\right)$$
(11)
where $`\mathrm{sin}\psi =\left(\mathrm{cos}\varphi \mathrm{cos}\theta \mathrm{sin}\alpha \mathrm{cos}\alpha \mathrm{sin}\theta \right)`$ is the polar angle in the frame alinged with $`𝛀`$ and we assumed that $`𝐁`$ is antiparallel to $`𝐳`$ and $`𝛀=\mathrm{\Omega }\{\mathrm{sin}\alpha ,0,\mathrm{cos}\alpha \}`$. There is also a non-inertial, spatially distributed charge density associated with this electric field (the Goldreich-Julian density)
$$\rho =\frac{1}{4\pi e}\mathrm{div}𝐄=\frac{𝛀𝐁_{\mathrm{𝐭𝐨𝐭}}}{2\pi ec}.$$
(12)
The ample supply of charges inside the star will short out the total electric field inside the star by producing surface charge density $`\sigma =\frac{1}{4\pi }E_r`$:
$`\sigma _{orb}`$ $`=`$ $`{\displaystyle \frac{1}{4\pi cR_0}}\left(𝐁_\mathrm{𝟎}[𝐑_0\times 𝐯]\right)`$
$`\sigma _{rot}`$ $`=`$ $`{\displaystyle \frac{3B_{tot}\mathrm{\Omega }R_0\mathrm{sin}(\psi )\mathrm{sin}\theta }{8\pi c}}`$ (13)
The surface charge distribution $`\sigma _{orb}`$ has dipole structure with $`𝐛=\frac{R_0^2}{8\pi c}\left[𝐯\times 𝐁_\mathrm{𝟎}\right]`$, while $`\sigma _{rot}`$ has monopole (total charge $`Q_{rot}=B_0\mathrm{\Omega }R_{0}^{}{}_{}{}^{3}\mathrm{cos}\alpha /c`$) and quadrupole contributions. Both types of surface charges will produce electric fields outside of the star with nonzero component along the magnetic field line of the order
$`E_{,orb}`$ $``$ $`{\displaystyle \frac{R_{0}^{}{}_{}{}^{3}}{cr^4}}\left(𝐁_0\left[𝐫\times 𝐯\right]\right)\mathrm{cos}\theta `$
$`E_{,rot}`$ $``$ $`{\displaystyle \frac{B_0\mathrm{\Omega }R_{0}^{}{}_{}{}^{3}\mathrm{cos}\alpha \mathrm{cos}\theta }{cr^2}}`$ (14)
These electric fields will accelerate the primary charges to relativistic energies.
The surface charge distribution $`\sigma _{orb}`$ is stationary in the moving but non-rotating frame, while surface charges $`\sigma _{rot}`$ are stationary in the neutron star frame (which is moving and rotating with respect to the lab frame). An observer in the neutron star frame will detect three types of electric currents: due to rotation of $`\sigma _{orb}`$ and inertial currents due to electric fields of the surface charges $`𝐣_{in}=\left[(𝐄_{\sigma _{orb}}+𝐄_{\sigma _{rot}})\times 𝛀\right]/(4\pi )`$. Inertial current due to $`𝐄_{\sigma _{rot}}`$ will generate magnetic field of the order $`B_{rot}B_0(\mathrm{\Omega }R_0/c)^2`$ with a component perpendicular to the surface of the millisecond pulsar. Equivalently, in the laboratory frame the charges $`\sigma _{rot}`$ rotating with the star will produce surface currents along the $`\psi `$ direction $`𝐠_{rot}=\sigma _{rot}\left[𝛀\times 𝐫\right]`$ that will generate magnetic field $`B_{rot}`$.
The Poynting losses due to rotating dipole will be proportional to the time varying component of the induced magnetic field $`B_0(\mathrm{\Omega }R_0/c)^2\mathrm{sin}\alpha \mathrm{cos}(\mathrm{\Omega }t)`$, so that the resulting Poynting flux would be $`PB^2(\mathrm{\Omega }R_0/c)^4\mathrm{sin}^2\alpha `$. It is suppressed by a small factor $`(\mathrm{\Omega }R_0/c)^41`$ if compared with the rotating dipole with the strength equal to the external magnetic field.
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# Neutrino Proper Time?
## 1 Introduction
The quantum phase of a free electron on a given path $`\mathrm{\Gamma }`$ is the product of its mass and the proper time along $`\mathrm{\Gamma },`$ divided by the constant $`\mathrm{}.`$ Expressed in a general reference frame the product of mass and proper time, i.e. the action, is the integral of the space-time scalar product of energy-momentum and displacement. The classical path is a path that has a quantum phase that is stationary against small fluctuations $`\delta \mathrm{\Gamma }`$ in the path.
If the path for a neutrino is light-like and in the direction of its light-like energy-momentum, then its quantum phase doesn’t change, because the scalar product of its energy-momentum and displacement is then zero. This is the problem; the electron quantum phase on the electron’s classical path is nonzero and proportional to its proper time while the neutrino quantum phase is completely different: the neutrino clock is frozen - it doesn’t change as the neutrino moves.
Hence, in Sec. 2 and 3, we consider neutrino paths with varying proper time, i.e. clocks that actually run. The neutrino proper time increases along the paths just as the proper time increases as an electron moves on its classical path. The neutrino quantum phase is then defined to be the proper time times the electron mass, just as it is for an electron. This definition makes the quantum phase of electrons and electron neutrinos agree. Electron, muon, and tau neutrinos have phases that disagree because the proper time is multiplied by the electron, muon, or tau mass, respectively.
In Sec. 4 we consider the possibility that muon neutrinos travel along space-like paths. We show that there is a loss of 7% of the muon neutrinos from pion decays and a 10% loss of muon neutrinos from muon decays, giving an average loss of 9%. These 9% carry their forward momentum backwards into space and are not detected.
Evidence for such an effect may have been found by the Super-Kamiokande collaboration (SK). SK finds a depletion of atmospheric muon neutrinos compared with the expected flux. Determining the expected flux is in itself a complex problem. At SK both the dependence on distance from source to detector and the overall normalization are measured. The observed overall normalization is 15.8% with an estimated uncertainty of 25%. Therefore the 9% depletion derived in Sec. 4 is about half of the central value reported, 15.8%. Thus the depletion of muon neutrinos at the source could be partly caused by some muon neutrinos moving in directions that differ from their momenta.
No electron neutrino depletion is reported by Super-Kamiokande. In Sec. 5 we show that no depletion implies that the electron neutrinos have time-like or light-like paths. The arrival of electron neutrinos only a matter of hours ahead of light from SN1987A is generally taken to mean that the neutrinos traveled to Earth from the supernova at light-speed or at near-light speed. Thus the light-like option for electron neutrinos explains both the SK observations and the SN1987A reports.
The predicted motions can be tested with current accelerator and detector technology. Arranging for a neutrino detector to be placed upstream from a pion or muon beam should detect the reverse motion muon neutrinos produced in the pion or muon decays, if such reverse motion muon neutrinos exist.
## 2 Neutrino Paths with Proper Time
Before dealing with the neutrino, we discuss the path of an electron in space-time and how electron paths are related to quantum phase. Let $`p^\mu `$ be energy-momentum, $`p^\mu `$ = $`\{p^k,E\},`$ and let $`x^\mu `$ be position, $`x^\mu `$ = $`\{x^k,t\},`$ where $`\mu `$ $`\{1,2,3,4\}`$ and $`k`$ $`\{1,2,3\}.`$ Consider a path $`\mathrm{\Gamma }`$ = $`x^\mu (\tau ),`$ where $`\tau `$ is the proper time along $`\mathrm{\Gamma },`$ $`(d\tau )^2`$ = $`(dt)^2`$ $`dx^kdx^k.`$ The quantum phase $`\varphi `$ along $`\mathrm{\Gamma }`$ is just $`S_\mathrm{\Gamma }/\mathrm{}.`$ The classical path is a straight line that has a quantum phase of
$$\varphi =\frac{S}{\mathrm{}}=\frac{Et+p^kx^k}{\mathrm{}},(\mathrm{electron})$$
(1)
where $`S`$ is the action along the classical path and repeated indices are summed.
The mass $`m,`$ $`m>`$ 0, determines the square of the energy-momentum,
$$E^2p^kp^k=m^2(\mathrm{electron})$$
(2)
The action $`S`$ is invariant under space-time transformations of the 4-vectors $`p^\mu `$ and $`x^\mu `$. So we can transform to a reference frame in which the momentum $`p^k`$ vanishes. We assume that $`x^k`$ is proportional to $`p^k`$ so that $`x^k`$ also vanishes. By (2), $`E`$ = $`m.`$ By (1), $`t`$ = $`S/m,`$ and, since this is the rest frame, the proper time $`\tau `$ is $`S/m.`$ Hence the square of the magnitude of $`x^\mu `$ in that frame as in all frames is
$$\tau ^2=t^2x_{}^{k}{}_{}{}^{2}=\frac{S^2}{m^2}(\mathrm{electron})$$
(3)
By (1) and (3) the action on the classical path in the rest frame is $`S`$ = $`Et`$ = $`m\tau ,`$ the electron mass times the proper time. In a general inertial frame we get
$$S=Etp^kx^k=m\tau (\mathrm{electron})$$
(4)
Assume that the electron and the electron neutrino have the same form for the quantum phase along the classical path,
$$\varphi =\frac{S}{\mathrm{}}=\frac{pt+p^kx^k}{\mathrm{}},(\mathrm{neutrino})$$
(5)
where the energy of the neutrino is the magnitude $`p`$ of its momentum, so that the energy-momentum is a null 4-vector.
In this paper the electron and the electron neutrino share the same quantum phase, $`m\tau .`$ We get
$$S=ptp^kx^k=m\tau ,(\mathrm{neutrino})$$
(6)
where $`m`$ is the electron mass. Each flavor has its own mass; for the muon or tau neutrino $`m`$ is the muon or tau mass.
The alternative classical paths possible for a neutrino have sub- or super-luminal speeds or have the speed of light but in the direction opposite to the momentum. For space-like paths proper distance is more appropriate than proper time; I continue to say ‘proper time’ even when I mean proper distance.
We define $`\tau `$ to be the proper time as usual for time-like paths. For space-like paths it is actually the proper distance. We have
$$t^2x_{}^{k}{}_{}{}^{2}=\eta \tau ^2,(\mathrm{neutrino})$$
(7)
where $`\eta `$ = $`+1`$ indicates a time-like path, $`\eta `$ = $`1`$ indicates a space-like path, and $`\eta `$ = $`0`$ indicates a light-like path.
By (7), a light-like path, $`\eta `$ = 0, implies that $`x^k`$ = $`\pm tn^k,`$ if $`x^k`$ and $`p^k`$ are both in the direction of the unit 3-vector $`n^k.`$ For the solution $`x^k`$ = $`+tn^k,`$ and by (6), the proper time $`\tau `$ is stuck at zero. In this paper the underlying presumption is that the action for an electron and its neutrino are one and the same function of proper time. Hence we dismiss the possibility that a neutrino path can have the solution $`x^k`$ = $`+tn^k,`$ because no electron path has zero proper time.
It may be that the neutrino follows a space-like path, i.e. $`\eta `$ = $`1`$ in (7). In this case the electron and its neutrino do not travel the same kind of path; the electron’s classical path is time-like, $`\eta `$ = $`+1.`$ This is true for the classical paths. But in quantum mechanics we can consider paths $`\mathrm{\Gamma }`$ that are space-like even for electrons. Thus having the electron and its neutrino share the same internal clock implies extending the formulas to space-like neutrino paths.
## 3 Neutrino Paths in Two Frames
Relativity protects scalar products and 4-vector magnitudes, but directions in space are not preserved under Lorentz transformations. If neutrino momentum and displacement can be in different directions, then only in special frames do momentum and displacement align, dubbed ‘Aligned Momentum’ (AM) frames. ‘Aligned’ is used in the following ways: the momentum is ‘in the same direction’ as the displacement or the momentum is ‘in the opposite direction’ of the displacement. A boost in the direction of the momentum and displacement takes one AM frame to another AM frame.
In an AM frame we introduce new notation $`X`$ and $`T`$ defined by
$$x^k=X\frac{\tau }{2}n^kt=T\frac{\tau }{2}.(\mathrm{AM})$$
(8)
where $`n^k`$ is a unit vector in the direction of the momentum, $`p^k`$ = $`pn^k.`$ By (6) and (7) we have
$$X=\eta s\frac{1}{s}T=\eta s+\frac{1}{s}(\mathrm{AM})$$
(9)
where $`s`$ stands for the ratio of momentum to mass, $`s`$ = $`p/m`$ (= $`\mathrm{sinh}w,`$ where $`\mathrm{tanh}w`$ is the speed of a particle of mass $`m`$ and momentum $`p.`$) In Fig. 2 the path $`(X,T)`$ is plotted for the three values of $`\eta .`$
The AM frame and a general frame, the ‘LAB’ frame, differ by a space-time transformation, i.e. a rotation followed by a boost. Choose the positive $`x`$-axis as the direction of the boost. Let $`\alpha `$ be the angle between the direction $`X,`$ i.e. the direction $`n^k`$ of the momentum in the AM frame, and the positive $`x`$-axis. See Fig. 3a. We get
$$X^{}=\gamma [X\mathrm{cos}\alpha +\beta T]Y^{}=X\mathrm{sin}\alpha T^{}=\gamma [T+\beta X\mathrm{cos}\alpha ](\mathrm{LAB})$$
(10)
where a prime indicates a value in the LAB frame, $`\beta `$ is the boost velocity, and $`\gamma `$ = $`(1\beta ^2)^{1/2}.`$ The neutrino energy-momentum in the LAB frame is
$$p_x^{}=\gamma [p\mathrm{cos}\alpha +\beta p]p_y^{}=p\mathrm{sin}\alpha p^{}=\gamma [p+\beta p\mathrm{cos}\alpha ](\mathrm{LAB})$$
(11)
## 4 Muon Neutrinos
Cosmic rays continually bombard the Earth’s atmosphere creating muon and electron neutrinos as well as other particles. A muon neutrino is produced together with a muon in charged pion decay and together with an electron and an electron neutrino when a muon decays. Other neutrino-producing processes are said to occur much less often.
In the Super-Kamiokande (SK) experiment atmospheric neutrinos are detected, muon neutrinos from pion and muon decays and electron neutrinos from the muon decays. We consider muon neutrinos first. We consider only what happens to muon neutrinos just after they are created at the source. This discussion has no bearing on the depletion of muon neutrinos in flight from the source in the atmosphere to the SK detector.
In the rest frame of the decaying particle, all directions in space are equivalent. Hence the rest frame, i.e. the Center of Mass (COM) frame, is convenient for finding average momenta. See Fig. 4.
In the decay of a pion or muon, the muon and the muon neutrino can share a common momentum and be aligned only in one special reference frame. To make the muon and muon neutrino even more alike in this frame assume that this special frame is also an Aligned Momentum (AM) frame so that the muon neutrino and its displacement are aligned. The muon momentum and the muon neutrino momentum are equal in magnitude and aligned and the muon and muon neutrino displacements are also aligned. See Fig. 1.
For muon neutrinos originating in pion decay, the pion rest frame is both the Center of Mass frame in which neutrinos are emitted in all directions equally and the AM frame in which the neutrino momentum and muon momentum are equal in magnitude. Let the SK rest frame be the LAB frame. Since the pion rest frame and the LAB frame typically differ by a large boost, the direction of the boost from the AM to the LAB frame, the $`x`$-direction in Fig. 3a, is along the line from the pion to SK. And, in the pion rest frame, the angle $`\alpha `$ is the angle between the emitted neutrino and the $`x`$-direction toward SK.
For muon neutrinos detected at SK with momenta of 400 MeV/$`c`$ or more, the factor $`\gamma `$ is more than 10. By (10) and (11) and as shown in Fig. 5 and Fig. 6, for large factors $`\gamma >>`$ 1, the displacement $`X_{}^{}{}_{}{}^{k}`$ is mainly either in the same direction as or in the direction opposite to the momentum $`p_{}^{}{}_{}{}^{k}`$ in the LAB frame.
From Fig. 6, we see that reverse motion, i.e. $`x_{}^{}{}_{}{}^{k}`$ opposite to momentum $`p_{}^{}{}_{}{}^{k},`$ occurs for space-like paths, $`\eta `$ = $`1.`$ By (9), (10), and (11) we find that $`X^{}`$ is opposite to $`p_{}^{}{}_{}{}^{k}`$, i.e. $`\mathrm{cos}\theta `$ = $`1,`$ for directions with $`\alpha <\alpha _0`$ where
$$\mathrm{cos}\alpha _0=\frac{T}{X}=\frac{1s^2}{1+s^2}.(\eta =1)(\gamma >>1)$$
(12)
Therefore muon neutrinos emitted in the direction of SK out to an angle of $`\alpha _0`$ actually move back toward outer space in the SK frame of reference. These muon neutrinos cannot be detected at SK.
Assume that muon neutrinos travel along space-like paths, i.e. choose $`\eta `$ = $`1`$ in (7). Since all directions $`\alpha `$ in the pion frame are equally likely, we calculate the expected normalization $`N`$ to be the fraction of the unit sphere for which the path of the muon neutrinos is toward SK. We get
$$N=\frac{\mathrm{toward}\mathrm{SK}}{\mathrm{all}}=\frac{_{\alpha _0}^\pi 2\pi \mathrm{sin}\alpha d\alpha }{4\pi }=\frac{1}{1+s^2}.(\eta =1)$$
(13)
In the rest frame of a decaying pion, by conservation of energy and momentum and the known rest masses of the pion and muon, we find that the ratio of muon neutrino momentum to muon mass is $`s`$ = $`p/m_\mu `$ = $`(m_\pi ^2m_\mu ^2)/(2m_\mu m_\pi )`$ = 0.28. By (13) the normalization $`N`$ for muon neutrinos from pion decay is $`N`$ = 0.93.
Atmospheric muon neutrinos are also produced in muon decay. A muon decays to an electron, an electron neutrino, and a muon neutrino. The electron mass is only 1/210th of the muon mass, so let us neglect the electron mass. Then we can treat the three decay particles equally, since each decay particle has a null energy-momentum 4-vector. The energy of each decay particle is the magnitude of its momentum. In the rest frame of the muon, the three momenta cancel. Thus the three momenta form a triangle with a constant perimeter equal to the muon mass, see Fig. 4b. An equilateral triangle is the symmetric solution with no one decay particle being given any more or less momentum than any other decay particle. Thus an estimate of the average muon neutrino momentum is about one-third of a muon mass, $`p_\nu `$ = $`m_\mu /3.`$
But the muon rest frame is also an Aligned Momentum frame for the muon neutrino. This follows because, as we already assumed, the muon neutrino is certain to have its momentum and displacement aligned in the frame in which the muon and the muon neutrino have aligned momenta with equal magnitudes as in Fig. 1b. Furthermore the transformation from the AM frame with equal momenta, i.e. Fig. 1b, to the muon rest frame, Fig. 4b, is a boost that is itself aligned with the momenta and displacements. Hence the muon rest frame is an Aligned Momentum frame for the muon neutrino.
Since the muon rest frame is an AM frame and since all directions in the rest frame are equivalent, we can use (13) to calculate the normalization. The estimate above gives $`s`$ = $`p/m_\mu `$ = $`1/3.`$ By (13) we get $`N`$ = $`0.90`$.
The overall normalizations estimated for muon neutrinos from pion and muon decay are collected in the following table.
| $`\nu _\mu `$ Source | $`s`$ = $`p/m_\mu `$ | $`\nu _\mu `$ Normalization |
| --- | --- | --- |
| $`\pi ^{}`$ $`\mu ^{}+\overline{\nu }_\mu `$ | $`0.28`$ | 0.93 |
| $`\mu ^{}`$ $`e^{}+\overline{\nu }_e+\nu _\mu `$ | $`0.33`$ | 0.90 |
Together the pion and muon sources of muon neutrinos give an average normalization of
$$\overline{N}=\frac{0.93+0.90}{2}=0.91,(\mathrm{neutrino})(\eta =1)$$
(14)
which is an 9% depletion. The best fit value of the source depletion at Super-Kamiokande is an overall normalization of 15.8% with an estimated 25% uncertainty. Thus the observed Super-Kamiokande depletion can be partly explained by muon neutrinos that move along space-like paths carrying light-like momenta. If future observations bring the 15.8% depletion down closer to 9%, then some other muon depletion mechanism would not be needed.
An accelerator based experiment testing the model could be devised. By the discussion of pion and muon decays above, arranging for a detector to be placed upstream of a pion or muon beam should detect upstream muon neutrinos. These muon neutrinos carry light-like momenta backwards, opposite to the direction of the neutrino momentum. Thus experiments should look for neutrino-induced events upstream that deposit downstream momenta.
## 5 Electron and Tau Neutrinos
Atmospheric electron neutrinos are produced in muon decays. By the discussion in Sec. 4 the electron neutrino is emitted in all directions equally likely in the muon rest frame, the COM frame. And the average electron neutrino has a momentum of $`m_\mu /3`$ the same as the muon neutrino. For an electron neutrino the momentum to mass ratio is very large, $`s`$ = $`p/m_e`$ = $`m_\mu /(3m_e)`$ = 70. By (13), if the electron neutrinos move on space-like paths, then 99.98% of atmospheric electron neutrinos move in the direction opposite to their momenta in the SK LAB frame. These would not be seen at SK.
But no electron depletion is reported by SK. Hence electron neutrinos cannot move on space-like paths, $`\eta `$ $`1.`$ From Fig. 5 we see that time-like and light-like paths are still open to the electron neutrino since virtually all neutrinos with $`\eta `$ = $`+1`$ or $`\eta `$ = 0 would move in the direction of their momenta and arrive at the SK detector.
The detection of neutrinos just hours before photons from the supernova SN1987A is strong evidence that neutrinos traveled on light-like or at least near light-like paths in the direction from SN1987A toward Earth. Furthermore the momentum absorbed in the detectors pointed away from SN1987A. The displacement and energy-momentum appear to be null 4-vectors pointing in the same direction. By Fig. 5 these electron neutrinos could have $`\eta `$ = $`0`$ if the factor $`\gamma `$ for the boost from the COM frame to the LAB frame is large enough. We assume that it is and conclude that electron neutrinos move along light-like paths.
Tau neutrinos have not been detected, as far as I know. By the discussion of muon decay, a tau neutrino would have an average momentum of $`p`$ = $`m_\tau /3`$ = $`592`$ MeV in the tau rest frame. For taus created essentially at rest in the LAB frame, the tau neutrinos would follow one of the paths in Fig. 2 but with $`s`$ = $`0.33.`$ Hence, for a time-like path, a space-like path, or a light-like path, all the tau neutrinos would travel in the direction opposite to the momentum carried by the tau neutrino. Experiments verifying this behavior would find tau neutrino induced events with momenta pointing towards the tau source, not away from the source.
## Appendix A Problems
1. Consider muon decay in the case when the muon neutrino has the average muon neutrino momentum of $`p`$ = $`m_\mu /3`$ in the muon rest frame, as found in Sec. 4. Find the relative velocity of the AM frame with equal momenta, Fig. 1b, with respect to the muon rest frame, Fig. 4b.
2. Show that there is no AM frame with equal momenta for neutrino scattering, i.e. no frame such as the ones illustrated in Fig. 1. The momentum before the interaction cannot be aligned with the momentum after the interaction and keep the same magnitude.
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# 1 Introduction
## 1 Introduction
We showed in part I of this paper that perturbative renormalization is a special case of a general mathematical procedure of extraction of finite values based on the Riemann-Hilbert problem. More specifically we associated to any given renormalizable quantum field theory an (infinite dimensional) complex Lie group $`G`$. We then showed that passing from the unrenormalized theory to the renormalized one was exactly the replacement of the loop $`d\gamma (d)G`$ of elements of $`G`$ obtained from dimensional regularization (for $`dD=`$ dimension of space-time) by the value $`\gamma _+(D)`$ of its Birkhoff decomposition, $`\gamma (d)=\gamma _{}(d)^1\gamma _+(d)`$.
The original loop $`d\gamma (d)`$ not only depends upon the parameters of the theory but also on the additional “unit of mass” $`\mu `$ required by dimensional analysis. We shall show in this paper that the mathematical concepts developped in part I provide very powerful tools to lift the usual concepts of the $`\beta `$-function and renormalization group from the space of coupling constants of the theory to the complex Lie group $`G`$.
We first observe, taking $`\phi _6^3`$ as an illustrative example to fix ideas and notations, that even though the loop $`\gamma (d)`$ does depend on the additional parameter $`\mu `$,
$$\mu \gamma _\mu (d),$$
$`(1)`$
the negative part $`\gamma _\mu ^{}`$ in the Birkhoff decomposition,
$$\gamma _\mu (d)=\gamma _\mu ^{}(d)^1\gamma _{\mu ^+}(d)$$
$`(2)`$
is actually independent of $`\mu `$,
$$\frac{}{\mu }\gamma _\mu ^{}(d)=0.$$
$`(3)`$
This is a restatement of a well known fact and follows immediately from dimensional analysis. Moreover, by construction, the Lie group $`G`$ turns out to be graded, with grading,
$$\theta _t\mathrm{Aut}G,tR,$$
$`(4)`$
inherited from the grading of the Hopf algebra $``$ of Feynman graphs given by the loop number,
$$L(\mathrm{\Gamma })=\text{loop number of}\mathrm{\Gamma }$$
$`(5)`$
for any 1PI graph $`\mathrm{\Gamma }`$.
The straightforward equality,
$$\gamma _{e^t\mu }(d)=\theta _{t\epsilon }(\gamma _\mu (d))tR,\epsilon =Dd$$
$`(6)`$
shows that the loops $`\gamma _\mu `$ associated to the unrenormalized theory satisfy the striking property that the negative part of their Birkhoff decomposition is unaltered by the operation,
$$\gamma (\epsilon )\theta _{t\epsilon }(\gamma (\epsilon )).$$
$`(7)`$
In other words, if we replace $`\gamma (\epsilon )`$ by $`\theta _{t\epsilon }(\gamma (\epsilon ))`$ we do not change the negative part of its Birkhoff decomposition. We settled now for the variable,
$$\epsilon =DdC\backslash \{0\}.$$
$`(8)`$
Our first result (section 2) is a complete characterization of the loops $`\gamma (\epsilon )G`$ fulfilling the above striking invariance. This characterization only involves the negative part $`\gamma _{}(\epsilon )`$ of their Birkhoff decomposition which by hypothesis fulfills,
$$\gamma _{}(\epsilon )\theta _{t\epsilon }(\gamma _{}(\epsilon )^1)\text{is convergent for}\epsilon 0.$$
$`(9)`$
It is easy to see that the limit of (9) for $`\epsilon 0`$ defines a one parameter subgroup,
$$F_tG,tR$$
$`(10)`$
and that the generator $`\beta =\left(\frac{}{t}F_t\right)_{t=0}`$ of this one parameter group is related to the residue of $`\gamma `$
$$\underset{\epsilon =0}{Res}\gamma =\left(\frac{}{u}\gamma _{}\left(\frac{1}{u}\right)\right)_{u=0}$$
$`(11)`$
by the simple equation,
$$\beta =Y\mathrm{Res}\gamma ,$$
$`(12)`$
where $`Y=\left(\frac{}{t}\theta _t\right)_{t=0}`$ is the grading.
This is straightforward but our result is the following formula (14) which gives $`\gamma _{}(\epsilon )`$ in closed form as a function of $`\beta `$. We shall for convenience introduce an additional generator in the Lie algebra of $`G`$ (i.e. primitive elements of $`^{}`$) such that,
$$[Z_0,X]=Y(X)X\text{Lie}G.$$
$`(13)`$
The scattering formula for $`\gamma _{}(\epsilon )`$ is then,
$$\gamma _{}(\epsilon )=\underset{t\mathrm{}}{lim}e^{t\left(\frac{\beta }{\epsilon }+Z_0\right)}e^{tZ_0}.$$
$`(14)`$
Both factors in the right hand side belong to the semi-direct product,
$$\stackrel{~}{G}=G\underset{\theta }{>}R$$
$`(15)`$
of the group $`G`$ by the grading, but of course the ratio (14) belongs to the group $`G`$.
This shows (section 3) that the higher pole structure of the divergences is uniquely determined by the residue and gives a strong form of the t’Hooft relations, which will come as an immediate corollary.
In section 4 we show, specializing to the massless case, that the formula for the bare coupling constant,
$$g_0=gZ_1Z_3^{3/2}$$
$`(16)`$
where both $`gZ_1=g+\delta g`$ and the field strength renormalization constant $`Z_3`$ are thought of as power series (in $`g`$) of elements of the Hopf algebra $``$, does define a Hopf algebra homomorphism,
$$_{CM}\stackrel{g_0}{},$$
$`(17)`$
from the Hopf algebra $`_{CM}`$ of coordinates on the group of formal diffeomorphisms of $`C`$ such that,
$$\phi (0)=0,\phi ^{}(0)=\mathrm{id}$$
$`(18)`$
to the Hopf algebra $``$ of the massless theory. We had already constructed in a Hopf algebra homomorphism from $`_{\mathrm{CM}}`$ to the Hopf algebra of rooted trees, but the physical significance of this construction was unclear.
The homomorphism (17) is quite different in that for instance the transposed group homomorphism,
$$G\stackrel{\rho }{}\mathrm{Diff}(C)$$
$`(19)`$
lands in the subgroup of odd diffeomorphisms,
$$\phi (z)=\phi (z)z.$$
$`(20)`$
Moreover its physical significance will be transparent. We shall show in particular that the image by $`\rho `$ of $`\beta =Y\mathrm{Res}\gamma `$ is the usual $`\beta `$-function of the coupling constant $`g`$.
We discovered the homomorphism (17) by lengthy concrete computations. We have chosen to include them in an appendix besides our conceptual proof given in section 4. The main reason for this choice is that the explicit computation allows to validate the concrete ways of handling the coproduct, coassociativity, symmetry factors$`\mathrm{}`$ that underly the theory.
As a corollary of the construction of $`\rho `$ we get an action by (formal) diffeomorphisms of the group $`G`$ on the space $`X`$ of (dimensionless) coupling constants of the theory. We can then in particular formulate the Birkhoff decomposition directly in the group,
$$\mathrm{Diff}(X)$$
$`(21)`$
of formal diffeomorphisms of the space of coupling constants.
The unrenormalized theory delivers a loop
$$\delta (\epsilon )\mathrm{Diff}(X),\epsilon 0$$
$`(22)`$
whose value at $`\epsilon `$ is simply the unrenormalized effective coupling constant.
The Birkhoff decomposition,
$$\delta (\epsilon )=\delta _+(\epsilon )\delta _{}(\epsilon )^1$$
$`(23)`$
of this loop gives directly then,
$$\delta _{}(\epsilon )=\text{ bare coupling constant}$$
$`(24)`$
and,
$$\delta _+(D)=\text{renormalized effective coupling constant.}$$
$`(25)`$
This result is now, in its statement, no longer depending upon our group $`G`$ or the Hopf algebra $``$. But of course the proof makes heavy use of the above ingredients.
Now the Birkhoff decomposition of a loop,
$$\delta (\epsilon )\mathrm{Diff}(X),$$
$`(26)`$
admits a beautiful geometric interpretation. If we let $`X`$ be a complex manifold and pass from formal diffeomorphisms to actual ones, the data (26) is the initial data to perform, by the clutching operation, the construction of a complex bundle,
$$P=(S^+\times X)_\delta (S^{}\times X)$$
$`(27)`$
over the sphere $`S=P_1(C)=S^+S^{}`$, and with fiber $`X`$,
$$XP\stackrel{\pi }{}S.$$
$`(28)`$
The meaning of the Birkhoff decomposition (23),
$$\delta (\epsilon )=\delta _+(\epsilon )\delta _{}(\epsilon )^1$$
is then exactly captured by an isomorphism of the bundle $`P`$ with the trivial bundle,
$$S\times X.$$
$`(29)`$
## 2 Asymptotic scaling in graded complex Lie groups
We shall first prove the formula (14) of the introduction in the general context of graded Hopf algebras and then apply it to the Birkhoff decomposition of the loop associated in part I to the unrenormalized theory.
We let $``$ be a connected commutative graded Hopf algebra (connected means that $`^{(0)}=C`$) and let $`\theta _t`$, $`tR`$ be the one parameter group of automorphisms of $``$ associated with the grading so that for $`x`$ of degree $`n`$,
$$\theta _t(x)=e^{tn}xtR.$$
$`(1)`$
By construction $`\theta _t`$ is a Hopf algebra automorphism,
$$\theta _t\mathrm{Aut}().$$
$`(2)`$
We also let $`Y=\left(\frac{}{t}\theta _t\right)_{t=0}`$ be the generator which is a derivation of $``$.
We let $`G`$ be the group of characters of $``$,
$$\phi :C$$
$`(3)`$
i.e. of homomorphisms from the algebra $``$ to $`C`$. The product in $`G`$ is given by,
$$(\phi _1\phi _2)(x)=\phi _1\phi _2,\mathrm{\Delta }x$$
$`(4)`$
where $`\mathrm{\Delta }`$ is the coproduct in $``$. The augmentation $`\overline{e}`$ of $``$ is the unit of $`G`$ and the inverse of $`\phi G`$ is given by,
$$\phi ^1,x=\phi ,Sx$$
$`(5)`$
where $`S`$ is the antipode in $``$.
We let $`L`$ be the Lie algebra of derivations,
$$\delta :C,$$
$`(6)`$
i.e. of linear maps on $``$ such that
$$\delta (xy)=\delta (x)\overline{e}(y)+\overline{e}(x)\delta (y)x,y.$$
$`(7)`$
Even if $``$ is of finite type so that $`^{(n)}`$ is finite dimensional for any $`nN`$ there are more elements in $`L`$ than in the Lie algebra $`P`$ of primitive elements,
$$\mathrm{\Delta }Z=Z1+1Z$$
$`(8)`$
in the graded dual Hopf algebra $`_{\mathrm{gr}}^{}`$ of $``$. But one passes from $`P`$ to $`L`$ by completion relative to the $`I`$-adic topology, $`I`$ being the augmentation ideal of $`_{\mathrm{gr}}^{}`$.
The linear dual $`^{}`$ is in general an algebra (with product given by (4)) but not a Hopf algebra since the coproduct is not necessarily well defined. It is however well defined for characters $`\phi `$ or derivations $`\delta `$ which satisfy respectively $`\mathrm{\Delta }\phi =\phi \phi `$, and $`\mathrm{\Delta }\delta =\delta 1+1\delta `$.
For $`\delta L`$ the expression,
$$\phi =\mathrm{exp}\delta $$
$`(9)`$
makes sense in the algebra $`^{}`$ since when evaluated on $`x`$ one has $`x,\delta ^n=0`$ for $`n`$ large enough (since $`x,\delta ^n=\mathrm{\Delta }^{(n1)}x,\delta \mathrm{}\delta `$ vanishes for $`n>\mathrm{deg}x`$). Moreover $`\phi `$ is a group-like element of $`^{}`$, i.e. a character of $``$. Thus $`\phi G`$.
The one parameter group $`\theta _t\mathrm{Aut}()`$ acts by automorphisms on the group $`G`$,
$$\theta _t(\phi ),x=\phi ,\theta _t(x)x$$
$`(10)`$
and the derivation $`Y`$ of $``$ acts on $`L`$ by
$$Y(\delta ),x=\delta ,Y(x),$$
$`(11)`$
and defines a derivation of the Lie algebra $`L`$ where we recall that the Lie bracket in $`L`$ is given by,
$$[\delta _1,\delta _2],x=\delta _1\delta _2\delta _2\delta _1,\mathrm{\Delta }xx.$$
$`(12)`$
Let us now consider a map,
$$\epsilon C\backslash \{0\}\phi _\epsilon G$$
$`(13)`$
such that for any $`x`$, $`\overline{e}(x)=0`$ one has,
$$\epsilon \phi _\epsilon ,x\text{is a polynomial in}\frac{1}{\epsilon }\text{without constant term.}$$
$`(14)`$
Thus $`\epsilon \phi _\epsilon `$ extends to a map from $`P_1(C)\backslash \{0\}`$ to $`G`$, such that
$$\phi _{\mathrm{}}=1.$$
$`(15)`$
For such a map we define its residue as the derivative at $`\mathrm{}`$, i.e. as,
$$\mathrm{Res}\phi =\underset{\epsilon \mathrm{}}{lim}\epsilon (\phi _\epsilon 1).$$
$`(16)`$
By construction $`\mathrm{Res}\phi L`$ is a derivation $`C`$. When evaluated on $`x`$, $`\mathrm{Res}\phi `$ is just the residue at $`\epsilon =0`$ of the function $`\epsilon \phi _\epsilon ,x`$.
We shall now assume that for any $`tR`$ the following limit exists for any $`x`$,
$$\underset{\epsilon 0}{lim}\phi _\epsilon ^1\theta _{t\epsilon }(\phi _\epsilon ),x.$$
$`(17)`$
Using (10), (4) and (5) we have,
$$\phi _\epsilon ^1\theta _{t\epsilon }(\phi _\epsilon ),x=\phi _\epsilon \phi _\epsilon ,(S\theta _{t\epsilon })\mathrm{\Delta }x,$$
$`(18)`$
so that with $`\mathrm{\Delta }x=x_{(1)}x_{(2)}`$ we get a sum of terms $`\phi _\epsilon ,Sx_{(1)}\phi _\epsilon ,\theta _{t\epsilon }(x_{(2)})`$ $`=P_1\left(\frac{1}{\epsilon }\right)e^{kt\epsilon }P_2\left(\frac{1}{\epsilon }\right)`$. Thus (17) just means that the sum of these terms is holomorphic at $`\epsilon =0`$. It is clear that the value at $`\epsilon =0`$ is then a polynomial in $`t`$,
$$F_t,x=\underset{\epsilon 0}{lim}\phi _\epsilon ^1\theta _{t\epsilon }(\phi _\epsilon ),x.$$
$`(19)`$
Let us check that $`tF_tG`$ is a one parameter group,
$$F_{t_1+t_2}=F_{t_1}F_{t_2}t_iR.$$
$`(20)`$
The group $`G`$ is a topological group for the topology of simple convergence, i.e.,
$$\phi _n\phi \text{iff}\phi _n,x\phi ,xx.$$
$`(21)`$
Moreover, using (10) one checks that
$$\theta _{t_1\epsilon }(\phi _\epsilon ^1\theta _{t_2\epsilon }(\phi _\epsilon ))F_{t_2}\text{when}\epsilon 0.$$
$`(22)`$
We then have $`F_{t_1+t_2}`$ $`=`$ $`\underset{\epsilon 0}{lim}\phi _\epsilon ^1\theta _{(t_1+t_2)\epsilon }(\phi _\epsilon )`$ $`=`$ $`\underset{\epsilon 0}{lim}\phi _\epsilon ^1\theta _{t_1\epsilon }(\phi _\epsilon )\theta _{t_1\epsilon }(\phi _\epsilon ^1\theta _{t_2\epsilon }(\phi _\epsilon ))`$ $`=`$ $`F_{t_1}F_{t_2}`$.
This proves (20) and we let,
$$\beta =\left(\frac{}{t}F_t\right)_{t=0}$$
$`(23)`$
which defines an element of $`L`$ such that,
$$F_t=\mathrm{exp}(t\beta )tR.$$
$`(24)`$
As above, we view $`^{}`$ as an algebra on which $`Y`$ acts as a derivation by (11). Let us prove,
Lemma 1. Let $`\epsilon \phi _\epsilon G`$ satisfy $`(17)`$ with $`\phi _\epsilon =1+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\frac{d_n}{\epsilon ^n}`$, $`d_n^{}`$. One then has
$$Y(d_1)=\beta Yd_{n+1}=d_n\beta n1.$$
Proof. Let $`x`$ and let us show that
$$\beta ,x=\underset{\epsilon 0}{lim}\epsilon \phi _\epsilon \phi _\epsilon ,(SY)\mathrm{\Delta }(x).$$
$`(25)`$
Using (18) we know by hypothesis that,
$$\phi _\epsilon \phi _\epsilon ,(S\theta _{t\epsilon })\mathrm{\Delta }(x)F_t,x$$
$`(26)`$
where the convergence holds in the space of holomorphic functions of $`t`$ in say $`|t|1`$ so that the derivatives of both sides at $`t=0`$ are also convergent thus yielding (25).
Now the function $`\epsilon \epsilon \phi _\epsilon \phi _\epsilon ,(SY)\mathrm{\Delta }x`$ is holomorphic for $`\epsilon C\backslash \{0\}`$ and also at $`\epsilon =\mathrm{}P_1(C)`$ since $`\phi _{\mathrm{}}=1`$. Moreover by (25) it is also holomorphic at $`\epsilon =0`$ and is thus a constant, which gives,
$$\phi _\epsilon \phi _\epsilon ,(SY)\mathrm{\Delta }(x)=\frac{1}{\epsilon }\beta ,x.$$
$`(27)`$
Using the product in $`^{}`$ this means that
$$\phi _\epsilon ^1Y(\phi _\epsilon )=\frac{1}{\epsilon }\beta ,$$
$`(28)`$
and multiplying by $`\phi _\epsilon `$ on the left, that,
$$Y(\phi _\epsilon )=\frac{1}{\epsilon }\phi _\epsilon \beta .$$
$`(29)`$
One has $`Y(\phi _\epsilon )={\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\frac{Y(d_n)}{\epsilon ^n}`$ and $`\frac{1}{\epsilon }\phi _\epsilon \beta =\frac{1}{\epsilon }\beta +{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\frac{1}{\epsilon ^{n+1}}d_n\beta `$. Thus (29) gives the lemma.
In particular we get $`Y(d_1)=\beta `$ and since $`d_1`$ is the residue, $`\mathrm{Res}\phi `$, this gives,
$$\beta =Y(\mathrm{Res}\phi ),$$
$`(30)`$
which shows that $`\beta `$ is uniquely determined by the residue of $`\phi _\epsilon `$.
We shall now write a formula for $`\phi _\epsilon `$ in terms of $`\beta `$. This is made possible by Lemma 1 which shows that $`\beta `$ uniquely determines $`\phi _\epsilon `$. What is not transparent from Lemma 1 is that for $`\beta L`$ the elements $`\phi _\epsilon ^{}`$ are group-like, so that $`\phi _\epsilon G`$. In order to obtain a nice formula we take the semi direct product of $`G`$ by $`R`$ acting on $`G`$ by the grading $`\theta _t`$,
$$\stackrel{~}{G}=G\underset{\theta }{>}R,$$
$`(31)`$
and similarly we let $`\stackrel{~}{L}`$ be the Lie algebra
$$\stackrel{~}{L}=LCZ_0$$
$`(32)`$
where the Lie bracket is given by
$$[Z_0,\alpha ]=Y(\alpha )\alpha L$$
$`(33)`$
and extends the Lie bracket of $`L`$.
We view $`\stackrel{~}{L}`$ as the Lie algebra of $`\stackrel{~}{G}`$ in a way which will become clear in the proof of the following,
Theorem 2. Let $`\epsilon \phi _\epsilon G`$ satisfy $`(17)`$ as above. Then with $`\beta =Y(\mathrm{Res}\phi )`$ one has,
$$\phi _\epsilon =\underset{t\mathrm{}}{lim}e^{tZ_0}e^{t\left(\frac{\beta }{\epsilon }+Z_0\right)}.$$
The limit holds in the topology of simple convergence in $`G`$. Both terms $`e^{tZ_0}`$ and $`e^{t\left(\frac{\beta }{\epsilon }+Z_0\right)}`$ belong to $`\stackrel{~}{G}`$ but their product belongs to $`G`$.
Proof. We endow $`^{}`$ with the topology of simple convergence on $``$ and let $`\theta _t`$ act by automorphisms of the topological algebra $`^{}`$ by (10). Let us first show, with,
$$\phi _\epsilon =1+\underset{n=1}{\overset{\mathrm{}}{}}\frac{d_n}{\epsilon ^n},d_n^{},$$
$`(34)`$
that the following holds,
$$d_n=_{s_1s_2\mathrm{}s_n0}\theta _{s_1}(\beta )\theta _{s_2}(\beta )\mathrm{}\theta _{s_n}(\beta )\mathrm{\Pi }𝑑s_i.$$
$`(35)`$
For $`n=1`$, this just means that,
$$d_1=_0^{\mathrm{}}\theta _s(\beta )𝑑s,$$
$`(36)`$
which follows from (30) and the equality
$$Y^1(x)=_0^{\mathrm{}}\theta _s(x)𝑑sx,\overline{e}(x)=0.$$
$`(37)`$
We see from (37) that for $`\alpha ,\alpha ^{}^{}`$ such that
$$Y(\alpha )=\alpha ^{},\alpha ,1=\alpha ^{},1=0$$
$`(38)`$
one has,
$$\alpha =_0^{\mathrm{}}\theta _s(\alpha ^{})𝑑s.$$
$`(39)`$
Combining this equality with Lemma 1 and the fact that $`\theta _s\mathrm{Aut}^{}`$ is an automorphism, gives an inductive proof of (35). The meaning of this formula should be clear, we pair both sides with $`x`$, and let,
$$\mathrm{\Delta }^{(n1)}x=x_{(1)}x_{(2)}\mathrm{}x_{(n)}.$$
$`(40)`$
Then the right hand side of (35) is just,
$$_{s_1\mathrm{}s_n0}\beta \mathrm{}\beta ,\theta _{s_1}(x_{(1)})\theta _{s_2}(x_{(2)})\mathrm{}\theta _{s_n}(x_{(n)})\mathrm{\Pi }𝑑s_i$$
$`(41)`$
and the convergence of the multiple integral is exponential since,
$$\beta ,\theta _s(x_{(i)})=O(e^s)\text{for}s+\mathrm{}.$$
$`(42)`$
We see moreover that if $`x`$ is homogeneous of degree $`\mathrm{deg}(x)`$, and if $`n>\mathrm{deg}(x)`$, at least one of the $`x_{(i)}`$ has degree 0 so that $`\beta ,\theta _s(x_{(i)})=0`$ and (41) gives 0. This shows that the pairing of $`\phi _\epsilon `$ with $`x`$ only involves finitely many non zero terms in the formula,
$$\phi _\epsilon ,x=\overline{e}(x)+\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{\epsilon ^n}d_n,x.$$
$`(43)`$
With all convergence problems out of the way we can now proceed to prove the formula of Theorem 2 without care for convergence.
Let us first recall the expansional formula ,
$$e^{(A+B)}=\underset{n=0}{\overset{\mathrm{}}{}}_{{\scriptscriptstyle u_j}=1,u_j0}e^{u_0A}Be^{u_1A}\mathrm{}Be^{u_nA}\mathrm{\Pi }𝑑u_j$$
$`(44)`$
(cf. for the exact range of validity of (44)).
We apply this with $`A=tZ_0`$, $`B=t\beta `$, $`t>0`$ and get,
$$e^{t(\beta +Z_0)}=\underset{n=0}{\overset{\mathrm{}}{}}_{{\scriptscriptstyle v_j}=t,v_j0}e^{v_0Z_0}\beta e^{v_1Z_0}\beta \mathrm{}\beta e^{v_nZ_0}\mathrm{\Pi }𝑑v_j.$$
$`(45)`$
Thus, with $`s_1=tv_0`$, $`s_1s_2=v_1,\mathrm{},s_{n1}s_n=v_{n1}`$, $`s_n=v_n`$ and replacing $`\beta `$ by $`\frac{1}{\epsilon }\beta `$, we obtain,
$$e^{t(\beta /\epsilon +Z_0)}=\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{\epsilon ^n}_{ts_1s_2\mathrm{}s_n0}e^{tZ_0}\theta _{s_1}(\beta )\mathrm{}\theta _{s_n}(\beta )\mathrm{\Pi }𝑑s_i.$$
$`(46)`$
Multiplying by $`e^{tZ_0}`$ on the left and using (41) thus gives,
$$\phi _\epsilon =\underset{t\mathrm{}}{lim}e^{tZ_0}e^{t(\beta /\epsilon +Z_0)}.$$
$`(47)`$
It is obvious conversely that this formula defines a family $`\epsilon \phi _\epsilon `$ of group-like elements of $`^{}`$ associated to any preassigned element $`\beta L`$.
Corollary 3. For any $`\beta L`$ there exists a (unique) map $`\epsilon \phi _\epsilon G`$ satisfying $`(17)`$ and $`(34)`$.
## 3 The renormalization group flow
Let us now apply the above results to the group $`G`$ associated in part I to the Hopf algebra $``$ of 1PI Feynman graphs of a quantum field theory. We choose $`\phi _6^3`$ for simplicity. As explained in part I the group $`G`$ is a semi-direct product,
$$G=G_0>G_c$$
$`(1)`$
of an abelian group $`G_0`$ by the group $`G_c`$ associated to the Hopf subalgebra $`_c`$ constructed on 1PI graphs with two or three external legs and fixed external structure. Passing from $`G_c`$ to $`G`$ is a trivial step and we shall thus concentrate on the group $`G_c`$. The unrenormalized theory delivers, using dimensional regularization with the unit of mass $`\mu `$, a loop,
$$\epsilon \gamma _\mu (\epsilon )G_c,$$
$`(2)`$
and we first need to see the exact $`\mu `$ dependence of this loop. We consider the grading of $`_c`$ and $`G_c`$ given by the loop number of a graph,
$$L(\mathrm{\Gamma })=IV+1$$
$`(3)`$
where $`I`$ is the number of internal lines and $`V`$ the number of vertices.
One has,
$$\gamma _{e^t\mu }(\epsilon )=\theta _{t\epsilon }(\gamma _\mu (\epsilon ))tR.$$
$`(4)`$
Let us check this using the formulas of section 3 of part I. For $`N=2`$ external legs the dimension $`B`$ of $`\sigma ,U_\mathrm{\Gamma }`$ is equal to 0 by (12) of loc.cit. Thus the $`\mu `$ dependence is given by
$$\mu ^{\frac{\epsilon }{2}V_3}$$
$`(5)`$
where $`V_3`$ is the number of 3-point vertices of $`\mathrm{\Gamma }`$. One checks that $`\frac{1}{2}V_3=L`$ as required. Similarly if $`N=3`$ the dimension $`B`$ of $`\sigma ,U_\mathrm{\Gamma }`$ is equal to $`\left(1\frac{3}{2}\right)d+3`$, $`d=6\epsilon `$ by (12) of loc.cit. so that the $`\mu `$-dependence is,
$$\mu ^{\frac{\epsilon }{2}V_3}\mu ^{\epsilon /2}.$$
$`(6)`$
But this time, $`V_3=2L+1`$ and we get
$$\mu ^{\epsilon L}$$
$`(7)`$
as required.
We now reformulate a well known result, the fact that counterterms, once appropriately normalized, are independent of $`m^2`$ and $`\mu ^2`$,
Lemma 4. Let $`\gamma _\mu =(\gamma _\mu ^{})^1(\gamma _{\mu ^+})`$ be the Birkhoff decomposition of $`\gamma _\mu `$. Then $`\gamma _\mu ^{}`$ is independent of $`\mu `$.
As in part I we perform the Birkhoff decomposition with respect to a small circle $`C`$ with center $`D=6`$ and radius $`<1`$.
The proof of the lemma follows immediately from . Indeed the dependence in $`m^2`$ has in the minimal subtraction scheme the same origin as the dependence in $`p^2`$ and we have chosen the external structure of graphs (eq. (41) of part I) so that no $`m^2`$ dependence is left <sup>2</sup><sup>2</sup>2This can be easily achieved by maintaining non-vanishing fixed external momenta. $`\gamma _\mu ^{}`$ is independent on such external structures by construction .. But then, since $`\mu ^2`$ is a dimensionful parameter, it cannot be involved any longer.
Corollary 5. Let $`\phi _\epsilon =(\gamma _\mu ^{})^1(\epsilon )`$, then for any $`tR`$ the following limit exists in $`G_c`$,
$$\underset{\epsilon 0}{lim}\phi _\epsilon ^1\theta _{t\epsilon }(\phi _\epsilon ).$$
In other words $`\epsilon \phi _\epsilon G_c`$ fulfills the condition (17) of section 2.
Proof. The product $`\phi _\epsilon ^1\gamma _\mu (\epsilon )`$ is holomorphic at $`\epsilon =0`$ for any value of $`\mu `$. Thus by (4), for any $`tR`$, both $`\phi _\epsilon ^1\gamma _\mu (\epsilon )`$ and $`\phi _\epsilon ^1\theta _{\epsilon t}(\gamma _\mu (\epsilon ))`$ are holomorphic at $`\epsilon =0`$. The same holds for $`\theta _{\epsilon t}(\phi _\epsilon ^1)\gamma _\mu (\epsilon )`$ and hence for the ratio
$$\phi _\epsilon ^1\gamma _\mu (\epsilon )(\theta _{\epsilon t}(\phi _\epsilon ^1)\gamma _\mu (\epsilon ))^1=\phi _\epsilon ^1\theta _{\epsilon t}(\phi _\epsilon ).$$
We let $`\gamma _{}(\epsilon )=\phi _\epsilon ^1`$ and translate the results of section 2.
Corollary 6. Let $`F_t=\underset{\epsilon 0}{lim}\gamma _{}(\epsilon )\theta _{t\epsilon }(\gamma _{}(\epsilon )^1)`$. Then $`F_t`$ is a one parameter subgroup of $`G_c`$ and $`F_t=\mathrm{exp}(t\beta )`$ where $`\beta =Y\mathrm{Res}\phi _\epsilon `$ is the grading operator $`Y`$ applied to the residue of the loop $`\gamma (\epsilon )`$.
In general, given a loop $`\epsilon \gamma (\epsilon )G`$ it is natural to define its residue at $`\epsilon =0`$ by first performing the Birkhoff decomposition on a small circle $`C`$ around $`\epsilon =0`$ and then taking,
$$\mathrm{Res}_{\epsilon =0}\gamma =\frac{}{u}(\phi _{1/u})_{u=0}$$
$`(8)`$
where $`\phi _\epsilon =\gamma _{}(\epsilon )^1`$ and $`\gamma (\epsilon )=\gamma _{}(\epsilon )^1\gamma _+(\epsilon )`$ is the Birkhoff decomposition.
As shown in section 2, the residue or equivalently $`\beta =Y\mathrm{Res}`$ uniquely determines $`\phi _\epsilon =\gamma _{}(\epsilon )^1`$ and we thus get, from Theorem 2,
Corollary 7. The negative part $`\gamma _{}(\epsilon )`$ of the Birkhoff decomposition of $`\gamma _\mu (\epsilon )`$ is independent of $`\mu `$ and given by,
$$\gamma _{}(\epsilon )=\underset{t\mathrm{}}{lim}e^{t\left(\frac{\beta }{\epsilon }+Z_0\right)}e^{tZ_0}.$$
As above we adjoined the primitive element $`Z_0`$ to implement the grading $`Y`$ (cf. section 2). Our choice of the letter $`\beta `$ is of course not innocent and we shall see in section 5 the relation with the $`\beta `$-function.
## 4 The action of $`G_c`$ on the coupling constants
We shall show in this section that the formula for the bare coupling constant $`g_0`$ in terms of 1PI graphs, i.e. the generating function,
$$g_0=(gZ_1)(Z_3)^{3/2}$$
$`(1)`$
where we consider the right hand side as a formal power series with values in $`_c`$ given explicitly by, (with $`\mathrm{}=L(\mathrm{\Gamma })`$ the loop number of the graphs),
$$g_0=\left(x+\underset{\text{}}{}x^{2l+1}\frac{\mathrm{\Gamma }}{S(\mathrm{\Gamma })}\right)\left(1\underset{\text{}}{}x^{2l}\frac{\mathrm{\Gamma }}{S(\mathrm{\Gamma })}\right)^{3/2}$$
$`(2)`$
does define a Hopf algebra homomorphism,
$$\mathrm{\Phi }:_{\mathrm{CM}}$$
$`(3)`$
from the Hopf algebra $`_{\mathrm{CM}}`$ of coordinates on the group of formal diffeomorphisms of $`C`$ with
$$\phi (0)=0,\phi ^{}(0)=1,$$
$`(4)`$
to the Hopf algebra $`_c`$ of 1PI graphs.
This result is only valid if we perform on $`_c`$ the simplification that pertains to the massless case, $`m=0`$, but because of the $`m`$-independence of the counterterms all the corollaries will be valid in general. The desired simplification comes because in the case $`m=0`$ there is no need to indicate by a cross left on an internal line the removal of a self energy subgraph. Indeed and with the notations of part I we can first of all ignore all the $`\text{}_{(0)}`$ since $`m^2=0`$, moreover the yield a $`k^2`$ term which exactly cancels out with the additional propagator when we remove the subgraph and replace it by . This shows that we can simply ignore all these crosses and write coproducts in the simplest possible way. To get familiar with this coproduct and with the meaning of the Hopf algebra morphism (3) we urge the reader to begin by the concrete computation done in the appendix, which checks its validity up to order six in the coupling constant.
Let us now be more explicit on the meaning of formula (2). We first expand $`g_0`$ as a power series in $`x`$ and get a series of the form,
$$g_0=x+\underset{2}{\overset{\mathrm{}}{}}\alpha _nx^n$$
$`(5)`$
where the even coefficients $`\alpha _{2n}`$ are zero and the coefficients $`\alpha _{2n+1}`$ are finite linear combinations of products of graphs, so that,
$$\alpha _{2n+1}n1.$$
$`(6)`$
We let $`_{\mathrm{CM}}`$ be the Hopf algebra of the group of formal diffeomorphisms such that (4) holds. We take the generators $`a_n`$ of $`_{\mathrm{CM}}`$ given by the equality
$$\phi (x)=x+\underset{n2}{}a_n(\phi )x^n.$$
$`(7)`$
and define the coproduct in $`_{\mathrm{CM}}`$ by the equality
$$\mathrm{\Delta }a_n,\phi _1\phi _2=a_n(\phi _2\phi _1)$$
$`(8)`$
We then define uniquely the algebra homomorphism
$$\mathrm{\Phi }:_{\mathrm{CM}}$$
by the condition,
$$\mathrm{\Phi }(a_n)=\alpha _n.$$
$`(9)`$
By construction $`\mathrm{\Phi }`$ is a morphism of algebras. We shall show that it is comultiplicative, i.e.
$$(\mathrm{\Phi }\mathrm{\Phi })\mathrm{\Delta }x=\mathrm{\Delta }\mathrm{\Phi }(x)x_{\mathrm{CM}}$$
$`(10)`$
and comes from a group morphism,
$$\rho :G_cG_2.$$
$`(11)`$
where $`G_2`$ is the group of characters of $`_{\mathrm{CM}}`$ which is by construction the opposite of the group of formal diffeomorphisms. In fact we shall first describe the corresponding Lie algebra morphism, $`\rho `$.
Let us first recall from part I that a 1PI graph $`\mathrm{\Gamma }`$ defines a primitive element $`(\mathrm{\Gamma })`$ of $`_{\mathrm{gr}}^{}`$ which only pairs nontrivially with the monomial $`\mathrm{\Gamma }`$ of $``$ and satisfies $`(\mathrm{\Gamma }),\mathrm{\Gamma }=1`$. We take the following natural basis $`\underset{¯}{\mathrm{\Gamma }}=S(\mathrm{\Gamma })(\mathrm{\Gamma })`$ for the Lie algebra of primitive elements of $`_{\mathrm{gr}}^{}`$, labelled by 1PI graphs with two or three external legs. By part I, theorem 2, their Lie bracket is given by,
$$[\underset{¯}{\mathrm{\Gamma }},\underset{¯}{\mathrm{\Gamma }}^{}]=\underset{v^{}}{}\underset{¯}{\mathrm{\Gamma }^{}^v^{}\mathrm{\Gamma }}\underset{v}{}\underset{¯}{\mathrm{\Gamma }^v\mathrm{\Gamma }^{}},$$
$`(12)`$
where $`\mathrm{\Gamma }^{}_v^{}\mathrm{\Gamma }`$ is the graph obtained by grafting $`\mathrm{\Gamma }`$ at the vertex $`v^{}`$ of $`\mathrm{\Gamma }^{}`$. (Our basis differs from the one used in loc. cit. by an overall - sign, but the present choice will be more convenient). In our context of the simplified Hopf algebra the places where a given graph $`\mathrm{\Gamma }`$ can be inserted in another graph $`\mathrm{\Gamma }^{}`$ are no longer always labelled by vertices of $`\mathrm{\Gamma }^{}`$. They are when $`\mathrm{\Gamma }`$ is a vertex graph but when $`\mathrm{\Gamma }`$ is a self energy graph such places are just labelled by the internal lines of $`\mathrm{\Gamma }^{}`$, as we could discard the use of external structures and two-point vertices for self energy graphs.
We also let $`Z_n^{}`$ be the natural basis of the Lie algebra of primitive elements of $`_{\mathrm{CM}}^{}`$ which corresponds to the vector fields $`x^{n+1}\frac{}{x}`$. More precisely, $`Z_n^{}`$ is given as the linear form on $`_{\mathrm{CM}}`$ which only pairs with the monomial $`a_{n+1}`$,
$$Z_n^{},a_{n+1}=1,$$
$`(13)`$
and the Lie bracket is given by,
$$[Z_n^{},Z_m^{}]=(mn)Z_{n+m}^{}.$$
$`(14)`$
We then first prove,
Lemma 8. Let $`\rho _\mathrm{\Gamma }=\frac{3}{2}`$ for $`2`$-point graphs and $`\rho _\mathrm{\Gamma }=1`$ for $`3`$-point graphs. The equality $`\rho (\underset{¯}{\mathrm{\Gamma }})=\rho _\mathrm{\Gamma }Z_2\mathrm{}^{}`$, where $`\mathrm{}=L(\mathrm{\Gamma })`$ is the loop number, defines a Lie algebra homomorphism.
Proof. We just need to show that $`\rho `$ preserves the Lie bracket. Let us first assume that $`\mathrm{\Gamma }_1,\mathrm{\Gamma }_2`$ are vertex graphs and let $`V_i`$ be the vertex number of $`\mathrm{\Gamma }_i`$. One has,
$$V=2L+1$$
$`(15)`$
for any vertex graph $`\mathrm{\Gamma }`$. Thus the Lie bracket $`\rho [\underset{¯}{\mathrm{\Gamma }}_1,\underset{¯}{\mathrm{\Gamma }}_2]`$ provides $`V_2V_1=2(L_2L_1)`$ vertex graph contributions all equal to $`Z_{2(L_1+L_2)}^{}`$ so that
$$\rho ([\underset{¯}{\mathrm{\Gamma }}_1,\underset{¯}{\mathrm{\Gamma }}_2])=2(L_2L_1)Z_{2(L_1+L_2)}^{}$$
$`(16)`$
which is exactly $`[\rho (\underset{¯}{\mathrm{\Gamma }}_1),\rho (\underset{¯}{\mathrm{\Gamma }}_2)]`$ by (14).
Let then $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_2`$ be 2-point graphs. For any such graph one has,
$$I=3L1$$
$`(17)`$
where $`I`$ is the number of internal lines of $`\mathrm{\Gamma }`$. Thus $`\rho ([\underset{¯}{\mathrm{\Gamma }}_1,\underset{¯}{\mathrm{\Gamma }}_2])`$ gives $`I_2I_1=3(L_2L_1)`$ 2-point graph contributions, each equal to $`\frac{3}{2}Z_{2(L_1+L_2)}^{}`$. Thus,
$$\rho ([\underset{¯}{\mathrm{\Gamma }}_1,\underset{¯}{\mathrm{\Gamma }}_2])=\frac{3}{2}3(L_2L_1)Z_{2(L_1+L_2)}^{}$$
$`(18)`$
but the right hand side is $`\rho _{\mathrm{\Gamma }_1}\rho _{\mathrm{\Gamma }_2}\mathrm{\hspace{0.17em}2}(L_2L_1)Z_{2(L_1+L_2)}^{}`$ so that,
$$\rho ([\underset{¯}{\mathrm{\Gamma }}_1,\underset{¯}{\mathrm{\Gamma }}_2])=[\rho (\underset{¯}{\mathrm{\Gamma }}_1),\rho (\underset{¯}{\mathrm{\Gamma }}_2)],$$
$`(19)`$
as required.
Finally if say $`\mathrm{\Gamma }_1`$ is a 3-point graph and $`\mathrm{\Gamma }_2`$ a 2-point graph, we get from $`[\underset{¯}{\mathrm{\Gamma }}_1,\underset{¯}{\mathrm{\Gamma }}_2]`$ a set of $`V_2`$ 2-point graphs minus $`I_1`$ 3-point graphs which gives,
$$\left(\frac{3}{2}V_2I_1\right)Z_{2(L_1+L_2)}^{}.$$
$`(20)`$
One has $`V_2=2L_2`$, $`I_1=3L_1`$ so that $`\frac{3}{2}V_2I_1=3(L_2L_1)=\rho _{\mathrm{\Gamma }_1}\rho _{\mathrm{\Gamma }_2}\mathrm{\hspace{0.17em}2}(L_2L_1)`$ which gives (19) as required.
We now have the Lie algebra morphism $`\rho `$ and the algebra morphism $`\mathrm{\Phi }`$. To $`\rho `$ corresponds a morphism of groups,
$$\rho :G_cG_2$$
$`(21)`$
and we just need to check that the algebra morphism $`\mathrm{\Phi }`$ is the transposed of $`\rho `$ on the coordinate algebras,
$$\mathrm{\Phi }(a)=a\rho a_{\mathrm{CM}}.$$
$`(22)`$
To prove (22) it is enough to show that $`\mathrm{\Phi }`$ is equivariant with respect to the action of the Lie algebra $`L`$ of primitive elements of $`_{\mathrm{gr}}^{}`$. More precisely, given a primitive element
$$Z_{\mathrm{gr}}^{},\mathrm{\Delta }Z=Z1+1Z,$$
$`(23)`$
we let $`_Z`$ be the derivation of the algebra $``$ given by,
$$_Z(y)=Z\mathrm{id},\mathrm{\Delta }yy.$$
$`(24)`$
What we need to check is the following,
Lemma 9. For any $`a_{\mathrm{CM}}`$, $`ZL`$ one has $`_Z\mathrm{\Phi }(a)=\mathrm{\Phi }(_{(\rho Z)}(a))`$.
Proof. It is enough to check the equality when $`Z`$ is of the form $`\underset{¯}{\mathrm{\Gamma }}=S(\mathrm{\Gamma })(\mathrm{\Gamma })`$ with the above notations. Thus we let $`\mathrm{\Gamma }`$ be a 1PI graph and let $`_\mathrm{\Gamma }`$ be the corresponding derivation of $``$ given by (24) with $`Z=S(\mathrm{\Gamma })(\mathrm{\Gamma })`$. Now by definition of the primitive element $`(\mathrm{\Gamma })`$ one has (cf. (48) section 2 of part I),
$$(\mathrm{\Gamma })\mathrm{id},\mathrm{\Delta }\mathrm{\Gamma }^{}=n(\mathrm{\Gamma },\mathrm{\Gamma }^{\prime \prime };\mathrm{\Gamma }^{})\mathrm{\Gamma }^{\prime \prime }$$
$`(25)`$
where the integer $`n(\mathrm{\Gamma },\mathrm{\Gamma }^{\prime \prime };\mathrm{\Gamma }^{})`$ is the number of subgraphs of $`\mathrm{\Gamma }^{}`$ which are isomorphic to $`\mathrm{\Gamma }`$ while $`\mathrm{\Gamma }^{}/\mathrm{\Gamma }\mathrm{\Gamma }^{\prime \prime }`$. By Theorem 2 of part I we have,
$$S(\mathrm{\Gamma })S(\mathrm{\Gamma }^{\prime \prime })n(\mathrm{\Gamma },\mathrm{\Gamma }^{\prime \prime };\mathrm{\Gamma }^{})=i(\mathrm{\Gamma },\mathrm{\Gamma }^{\prime \prime };\mathrm{\Gamma }^{})S(\mathrm{\Gamma }^{})$$
$`(26)`$
where $`i(\mathrm{\Gamma },\mathrm{\Gamma }^{\prime \prime };\mathrm{\Gamma }^{})`$ is the number of times $`\mathrm{\Gamma }^{}`$ appears in $`\mathrm{\Gamma }^{\prime \prime }\mathrm{\Gamma }`$. We thus get,
$$_\mathrm{\Gamma }\frac{\mathrm{\Gamma }^{}}{S(\mathrm{\Gamma }^{})}=i(\mathrm{\Gamma },\mathrm{\Gamma }^{\prime \prime };\mathrm{\Gamma }^{})\frac{\mathrm{\Gamma }^{\prime \prime }}{S(\mathrm{\Gamma }^{\prime \prime })}$$
$`(27)`$
which shows that $`_\mathrm{\Gamma }`$ admits a very simple definition in the generators $`\frac{\mathrm{\Gamma }^{}}{S(\mathrm{\Gamma }^{})}`$ of $``$.
The derivation $`_{(\rho Z)}`$ of $`_{\mathrm{CM}}`$ is also very easy to compute. One has by construction, (Lemma 8),
$$\rho (Z)=\rho _\mathrm{\Gamma }Z_2\mathrm{}^{}\mathrm{}=L(\mathrm{\Gamma })$$
$`(28)`$
and the derivation $`d_k`$ of $`_{\mathrm{CM}}`$ associated to the primitive element $`Z_k^{}`$ of $`_{\mathrm{CM}}^{}`$ is simply given, in the basis $`a_n_{\mathrm{CM}}`$ by,
$$d_k(a_n)=(nk)a_{nk}.$$
$`(29)`$
We thus get,
$$_{(\rho Z)}=\rho _\mathrm{\Gamma }d_2\mathrm{},\mathrm{}=L(\mathrm{\Gamma }).$$
$`(30)`$
Now by construction both $`_Z\mathrm{\Phi }`$ and $`\mathrm{\Phi }_{(\rho Z)}`$ are derivations from the algebra $`_{\mathrm{CM}}`$ to $``$ viewed as a bimodule over $`_{\mathrm{CM}}`$, i.e. satisfy,
$$\delta (ab)=\delta (a)\mathrm{\Phi }(b)+\mathrm{\Phi }(a)\delta (b).$$
$`(31)`$
Thus, to prove the lemma we just need to check the equality,
$$_\mathrm{\Gamma }\mathrm{\Phi }(a_n)=\rho _\mathrm{\Gamma }\mathrm{\Phi }(d_2\mathrm{}(a_n))\mathrm{}=L(\mathrm{\Gamma })$$
$`(32)`$
or equivalently using the generating function
$$g_0=x+\mathrm{\Phi }(a_n)x^n,$$
$`(33)`$
that
$$_\mathrm{\Gamma }g_0=\rho _\mathrm{\Gamma }x^{2\mathrm{}+1}\frac{}{x}g_0.$$
$`(34)`$
Now by construction of $`\mathrm{\Phi }`$ we have,
$$g_0=(xZ_1)(Z_3)^{3/2},Z_3=1\delta Z$$
$`(35)`$
where,
$$Z_1=1+\underset{\text{}}{}x^{2l}\frac{\mathrm{\Gamma }}{S(\mathrm{\Gamma })},\delta Z=\underset{\text{}}{}x^{2l}\frac{\mathrm{\Gamma }}{S(\mathrm{\Gamma })}$$
$`(36)`$
Thus, since both $`_\mathrm{\Gamma }`$ and $`\frac{}{x}`$ are derivations we can eliminate the denominators in (34) and rewrite the desired equality, after multiplying both sides by $`(1\delta Z)^{5/2}`$ as,
$$\begin{array}{cc}\left(x\frac{}{\mathrm{\Gamma }}Z_1\right)(1\delta Z)+\frac{3}{2}\left(\frac{}{\mathrm{\Gamma }}\delta Z\right)(xZ_1)\text{ }\hfill & (37)\hfill \\ =\rho _\mathrm{\Gamma }\left(x^{2\mathrm{}+1}\left(\frac{}{x}(xZ_1)\right)(1\delta Z)+\frac{3}{2}(xZ_1)x^{2\mathrm{}+1}\frac{}{x}\delta Z\right).\text{ }\hfill & \end{array}$$
Both sides of this formula are bilinear expressions in the 1PI graphs. We first need to compute $`\frac{}{\mathrm{\Gamma }}Z_1`$ and $`\frac{}{\mathrm{\Gamma }}\delta Z`$. One has
$$\frac{}{\mathrm{\Gamma }}Z_1=\underset{\text{}}{}x^{2l+2l^{}}c(\mathrm{\Gamma },\mathrm{\Gamma }^{})\frac{\mathrm{\Gamma }^{}}{S(\mathrm{\Gamma }^{})}$$
$`(38)`$
and,
$$\frac{}{\mathrm{\Gamma }}\delta Z=\underset{\text{}}{}x^{2l+2l^{}}c(\mathrm{\Gamma },\mathrm{\Gamma }^{})\frac{\mathrm{\Gamma }^{}}{S(\mathrm{\Gamma }^{})},$$
$`(39)`$
where $`\mathrm{}=L(\mathrm{\Gamma })`$, $`\mathrm{}^{}=L(\mathrm{\Gamma }^{})`$ are the loop numbers and the integral coefficient $`c(\mathrm{\Gamma },\mathrm{\Gamma }^{})`$ is given by,
$$c(\mathrm{\Gamma },\mathrm{\Gamma }^{})=V^{}\text{if}\rho _\mathrm{\Gamma }=1\text{and}c(\mathrm{\Gamma },\mathrm{\Gamma }^{})=I^{}\text{if}\rho _\mathrm{\Gamma }=3/2$$
$`(40)`$
(where $`V^{}`$ and $`I^{}`$ are respectively the number of vertices and of internal lines of $`\mathrm{\Gamma }^{}`$).
To prove (38) and (39) we use (27) and we get in both cases expressions like (38), (39) with
$$c(\mathrm{\Gamma },\mathrm{\Gamma }^{})=\underset{\mathrm{\Gamma }^{\prime \prime }}{}i(\mathrm{\Gamma },\mathrm{\Gamma }^{};\mathrm{\Gamma }^{\prime \prime }).$$
$`(41)`$
But this is exactly the number of ways we can insert $`\mathrm{\Gamma }`$ inside $`\mathrm{\Gamma }^{}`$ and is thus the same as (40). Let now $`\mathrm{\Gamma }_1`$ be a 3-point graph and $`\mathrm{\Gamma }_2`$ a 2-point graph. The coefficient of the bilinear term,
$$\frac{\mathrm{\Gamma }_1}{S(\mathrm{\Gamma }_1)}\frac{\mathrm{\Gamma }_2}{S(\mathrm{\Gamma }_2)},$$
$`(42)`$
in the left hand side of (37) is given by
$$\left(\frac{3}{2}c(\mathrm{\Gamma },\mathrm{\Gamma }_2)c(\mathrm{\Gamma },\mathrm{\Gamma }_1)\right)x^{2\mathrm{}+2\mathrm{}_1+2\mathrm{}_2+1}.$$
$`(43)`$
Its coefficient in the right hand side of (37) is coming from the terms,
$$\begin{array}{cc}x^{2\mathrm{}+1}\frac{}{x}\left(x^{2\mathrm{}_1+1}\frac{\mathrm{\Gamma }_1}{S(\mathrm{\Gamma }_1)}\right)\left(\frac{\mathrm{\Gamma }_2}{S(\mathrm{\Gamma }_2)}\right)x^{2\mathrm{}_2}+\text{ }\hfill & (44)\hfill \\ \frac{3}{2}x^{2\mathrm{}_1+1}\frac{\mathrm{\Gamma }_1}{S(\mathrm{\Gamma }_1)}x^{2\mathrm{}+1}\frac{}{x}\left(x^{2\mathrm{}_2}\frac{\mathrm{\Gamma }_2}{S(\mathrm{\Gamma }_2)}\right)\text{ }\hfill & \end{array}$$
which gives,
$$(3\mathrm{}_22\mathrm{}_11)x^{1+2\mathrm{}+2\mathrm{}_1+2\mathrm{}_2}.$$
$`(45)`$
We thus only need to check the equality,
$$\frac{3}{2}c(\mathrm{\Gamma },\mathrm{\Gamma }_2)c(\mathrm{\Gamma },\mathrm{\Gamma }_1)=(3\mathrm{}_22\mathrm{}_11)\rho _\mathrm{\Gamma }.$$
$`(46)`$
Note that in general, for any graph with $`N`$ external legs we have,
$$V=2(L1)+N,I=3(L1)+N.$$
$`(47)`$
Let us first take for $`\mathrm{\Gamma }`$ a 3-point graph so that $`\rho _\mathrm{\Gamma }=1`$. Then the left hand side of (46) gives $`\frac{3}{2}V_2V_1=\frac{3}{2}(2\mathrm{}_2)(2(\mathrm{}_11)+3)=3\mathrm{}_22\mathrm{}_11`$.
Let then $`\mathrm{\Gamma }`$ be a 2-point graph, i.e. $`\rho _\mathrm{\Gamma }=\frac{3}{2}`$. Then the left hand side of (46) gives $`\frac{3}{2}I_2I_1=\frac{3}{2}(3\mathrm{}_21)3\mathrm{}_1=\rho _\mathrm{\Gamma }(3\mathrm{}_22\mathrm{}_11)`$ which gives the desired equality.
Finally we also need to check the scalar terms and the terms linear in $`\mathrm{\Gamma }_1`$ or in $`\mathrm{\Gamma }_2`$. The only scalar terms in the left hand side of (37) are coming from $`x\frac{}{\mathrm{\Gamma }}Z_1+\frac{3}{2}x\frac{}{\mathrm{\Gamma }}\delta Z`$ and this gives,
$$x^{2\mathrm{}+1}\rho _\mathrm{\Gamma }.$$
$`(48)`$
The only scalar term in the right hand side of (37) comes from $`x^{2\mathrm{}+1}`$ thus they fulfill (37).
The terms linear in $`\mathrm{\Gamma }_1`$ in the left hand side of (37) come only from $`x\frac{}{\mathrm{\Gamma }}Z_1`$ if $`\mathrm{\Gamma }`$ is a 3-point graph and the coefficient of $`\mathrm{\Gamma }_1/S(\mathrm{\Gamma }_1)`$ is thus,
$$c(\mathrm{\Gamma },\mathrm{\Gamma }_1)x^{1+2\mathrm{}_1+2\mathrm{}}.$$
$`(49)`$
In the right hand side of (37) we just get
$$(2\mathrm{}_1+1)x^{1+2\mathrm{}_1+2\mathrm{}}.$$
$`(50)`$
We thus need to check that $`c(\mathrm{\Gamma },\mathrm{\Gamma }_1)=2\mathrm{}_1+1`$ which follows from (40) and (47) since $`V_1=2\mathrm{}_1+1`$.
Similarly, if $`\mathrm{\Gamma }`$ is a 2-point graph the left side of (37) only contributes by $`x\frac{}{\mathrm{\Gamma }}Z_1+\frac{3}{2}x^{2\mathrm{}+1}Z_1`$, so that the coefficient of $`\mathrm{\Gamma }_1/S(\mathrm{\Gamma }_1)`$ is
$$\left(c(\mathrm{\Gamma },\mathrm{\Gamma }_1)+\frac{3}{2}\right)x^{1+2\mathrm{}+2\mathrm{}_1}.$$
$`(51)`$
In the right hand side of (37) we get just as above
$$(2\mathrm{}_1+1)x^{1+2\mathrm{}+2\mathrm{}_1}$$
$`(52)`$
multiplied by $`\rho _\mathrm{\Gamma }=3/2`$.
Now here, since $`\mathrm{\Gamma }`$ is a 2-point graph, we have $`c(\mathrm{\Gamma },\mathrm{\Gamma }_1)=I_1=3(\mathrm{}_11)+3=3\mathrm{}_1`$ so that
$$c(\mathrm{\Gamma },\mathrm{\Gamma }_1)+\frac{3}{2}=\frac{3}{2}(2\mathrm{}_1+1)=\rho _\mathrm{\Gamma }(2\mathrm{}_1+1)$$
as required.
The check for terms linear in $`\mathrm{\Gamma }_2`$ is similar.
We can now state the main result of this section,
Theorem 10. The map $`\mathrm{\Phi }=_{\mathrm{CM}}`$ given by the effective coupling is a Hopf algebra homomorphism. The transposed Lie group morphism is $`\rho :G_cG_2`$.
The proof follows from Lemma 9 which shows that the map from $`G_c`$ to $`G_2`$ given by the transpose of the algebra morphism $`\mathrm{\Phi }`$ is the Lie group morphism $`\rho `$. By construction the morphism $`\mathrm{\Phi }`$ is compatible with the grading $`\mathrm{\Theta }`$ of $``$ and $`\alpha `$ of $`_{CM}`$ given by $`\mathrm{deg}(a_n)=n1`$ (cf. ), one has indeed,
$$\mathrm{\Phi }\alpha _t=\mathrm{\Theta }_{2t}\mathrm{\Phi },tR.$$
$`(51)`$
Finally we remark that our proof of Theorem 10 is similar to the proof of the equality
$$F_{\varphi _1\varphi _2}=F_{\varphi _2}F_{\varphi _1}$$
$`(52)`$
for the Butcher series used in the numerical integration of differential equations, but that the presence of the $`Z_3`$ factor makes it much more involved in our case.
## 5 The $`\beta `$-function and the Birkhoff decomposition of the unrenormalized effective coupling in the diffeomorphism group
Let us first recall our notations from part I concerning the effective action. We work in the Euclidean signature of space time and in order to minimize the number of minus signs we write the functional integrals in the form,
$$𝒩e^{S(\phi )}P(\phi )[D\phi ]$$
$`(1)`$
so that the Euclidean action is<sup>3</sup><sup>3</sup>3We know of course that the usual sign convention is better to display the positivity of the action functional.,
$$S(\phi )=\frac{1}{2}(_\mu \phi )^2\frac{1}{2}m^2\phi ^2+\frac{g}{6}\phi ^3.$$
$`(2)`$
The effective action, which when used at tree level in (1) gives the same answer as the full computation using (2), is then given in dimension $`d=6\epsilon `$ by,
$$\begin{array}{cc}S_{\mathrm{eff}}(\phi )=S(\phi )+\underset{1\mathrm{P}\mathrm{I}}{}\frac{(\mu ^{\epsilon /2}g)^{n2}}{n!}\frac{\mathrm{\Gamma }}{S(\mathrm{\Gamma })}(p_1,\mathrm{},p_n)\text{ }\hfill & (3)\hfill \\ \times \phi (p_1)\mathrm{}\phi (p_n)\mathrm{\Pi }dp_i\text{ }\hfill & \end{array}$$
where, as in part I, we do not consider tree graphs as 1PI and the integral is performed on the hyperplane $`p_i=0`$. To be more precise one should view the right hand side of (3) as a formal power series with values in the Hopf algebra $``$. The theory provides us with a loop
$$\gamma _\mu (\epsilon )=\gamma _{}(\epsilon )^1\gamma _{\mu _+}(\epsilon )$$
$`(4)`$
of characters of $``$. When we evaluate $`\gamma _\mu (\epsilon )`$ (resp. $`\gamma _{}`$, $`\gamma _{\mu ^+}`$) on the right hand side of (3) we get respectively the unrenormalized effective action, the bare action and the renormalized effective action (in the MS scheme).
Our notation is hiding the $`g`$-dependence of $`\gamma _\mu (\epsilon )`$, but this dependence is entirely governed by the grading. Indeed with $`t=\mathrm{log}(g)`$ one has, with obvious notations,
$$\gamma _{\mu ,g}(\epsilon )=\mathrm{\Theta }_{2t}(\gamma _{\mu ,1}(\epsilon ))$$
$`(5)`$
Since $`\mathrm{\Theta }_t`$ is an automorphism the same equality holds for both $`\gamma _{\mu _+}`$ and $`\gamma _{}`$.
As in section 4 we restrict ourselves to the massless case and let $`\gamma _\mu (\epsilon )=\gamma _{}(\epsilon )^1\gamma _{\mu _+}(\epsilon )`$ be the Birkhoff decomposition of $`\gamma _\mu (\epsilon )=\gamma _{\mu ,1}(\epsilon )`$.
Lemma 11. Let $`\rho :G_cG_2`$ be given by Theorem 10. Then $`\rho (\gamma _\mu (\epsilon ))(g)`$ is the unrenormalized effective coupling constant, $`\rho (\gamma _{\mu _+}(0))(g)`$ is the renormalized effective coupling constant and $`\rho (\gamma _{}(\epsilon ))(g)`$ is the bare coupling constant $`g_0`$.
This follows from (3) and Theorem 10. It is now straightforward to translate the results of the previous sections in terms of diffeomorphisms. The only subtle point to remember is that the group $`G_2`$ is the opposite of the group of diffeomorphisms so that if we view $`\rho `$ as a map to diffeomorphisms it is an antihomomorphism,
$$\rho (\gamma _1\gamma _2)=\rho (\gamma _2)\rho (\gamma _1).$$
$`(6)`$
Theorem 12. The renormalization group flow is the image $`\rho (F_t)`$ by $`\rho :G_c\mathrm{Diff}`$ of the one parameter group $`F_tG_c`$.
Proof. The bare coupling constant $`g_0`$ governs the bare action,
$$S_{\mathrm{bare}}(\phi _0)=\frac{1}{2}(_\mu \phi _0)^2+\mu ^{\epsilon /2}\frac{g_0}{6}\phi _0^3$$
$`(7)`$
in terms of the bare field $`\phi _0`$.
Now when we replace $`\mu `$ by
$$\mu ^{}=e^t\mu $$
$`(8)`$
we can keep the bare action, and hence the physical theory, unchanged provided we replace the renormalized coupling constant $`g`$ by $`g^{}`$ where
$$g_0(\epsilon ,g^{})=e^{\epsilon \frac{t}{2}}g_0(\epsilon ,g).$$
$`(9)`$
By construction we have,
$$g^{}=\psi _\epsilon ^1(e^{\epsilon \frac{t}{2}}\psi _\epsilon (g))$$
$`(10)`$
where $`\psi _\epsilon `$ is the formal diffeomorphism given by
$$\psi _\epsilon =\rho (\gamma _{}(\epsilon )).$$
$`(11)`$
Now the behaviour for $`\epsilon 0`$ of $`g^{}`$ given by (10) is the same as for,
$$\psi _\epsilon ^1\alpha _{\epsilon t/2}(\psi _\epsilon )$$
$`(12)`$
where $`\alpha _s`$ is the grading of $`\mathrm{Diff}`$ given as above by
$$\alpha _s(\psi )(x)=e^s\psi (e^sx).$$
$`(13)`$
Thus, since the map $`\rho `$ preserves the grading,
$$\rho (\theta _t(\gamma ))=\alpha _{t/2}\rho (\gamma )$$
$`(14)`$
(by (51) of section 4), we see by Corollary 6 of section 3 that,
$$g^{}\rho (F_t)g\text{when}\epsilon 0.$$
$`(15)`$
As a corollary we get of course,
Corollary 13. The image by $`\rho `$ of $`\beta L`$ is the $`\beta `$-function of the theory.
In fact all the results of section 3 now translate to the group $`G_2`$. We get the formula for the bare coupling constant in terms of the $`\beta `$-function, namely,
$$\psi _\epsilon =\underset{t\mathrm{}}{lim}e^{tZ_0}e^{t\left(\frac{\beta }{\epsilon }+Z_0\right)}$$
$`(16)`$
where $`Z_0=x\frac{}{x}`$ is the generator of scaling. But we can also express the main result of part I independently of the group $`G`$ or of its Hopf algebra $``$. Indeed the group homomorphism $`\rho :GG_2`$ maps the Birkhoff decomposition of $`\gamma _\mu (\epsilon )`$ to the Birkhoff decomposition of $`\rho (\gamma _\mu (\epsilon ))`$. But we saw above that $`\rho (\gamma _\mu (\epsilon ))`$ is just the unrenormalized effective coupling constant. We can thus state,
Theorem 14. Let the unrenormalized effective coupling constant $`g_{\mathrm{eff}}(\epsilon )`$ viewed as a formal power series in $`g`$ be considered as a loop of formal diffeomorphisms and let $`g_{\mathrm{eff}}(\epsilon )=(g_{\mathrm{eff}_{}})^1(\epsilon )g_{\mathrm{eff}_+}(\epsilon )`$ be its Birkhoff decomposition in the group of formal diffeomorphisms. Then the loop $`g_{\mathrm{eff}_{}}(\epsilon )`$ is the bare coupling constant and $`g_{\mathrm{eff}_+}(0)`$ is the renormalized effective coupling.
Note that $`G_2`$ is naturally isomorphic to the opposite group of Diff so we used the opposite order in the Birkhoff decomposition.
This result is very striking since it no longer involves the Hopf algebra $``$ or the group $`G`$ but only the idea of thinking of the effective coupling constant as a formal diffeomorphism. The proof is immediate, by combining Lemma 11, Theorem 10 of section 4 with Theorem 4 of part I.
Now in the same way as the Riemann-Hilbert problem and the Birkhoff decomposition for the group $`G=\mathrm{GL}(n,C)`$ are intimately related to the classification of holomorphic $`n`$-dimensional vector bundles on $`P_1(C)=C_+C_{}`$, the Birkhoff decomposition for the group $`G_2=\mathrm{Diff}^0`$ is related to the classification of one dimensional complex (non linear) bundles,
$$P=(C_+\times X)_{g_{\mathrm{eff}}}(C_{}\times X).$$
$`(17)`$
Here $`X`$ stands for a formal one dimensional fiber and $`C_\pm `$ are, as in part I, the components of the complement in $`P_1(C)`$ of a small circle around $`D`$. The total space $`P`$ should be thought of as a 2-dimensional complex manifold which blends together the $`\epsilon =Dd`$ and the coupling constant of the theory.
## 6 Conclusions
We showed in this paper that the group $`G`$ of characters of the Hopf algebra $``$ of Feynman graphs plays a key role in the geometric understanding of the basic ideas of renormalization including the renormalization group and the $`\beta `$-function.
We showed in particular that the group $`G`$ acts naturally on the complex space $`X`$ of dimensionless coupling constants of the theory. Thus, elements of $`G`$ are a refined form of diffeomorphisms of $`X`$ and as such should be called diffeographisms.
The action of these diffeographisms on the space of coupling constants allowed us first of all to read off directly the bare coupling constant and the renormalized one from the Riemann–Hilbert decomposition of the unrenormalized effective coupling constant viewed as a loop of formal diffeomorphisms. This showed that renormalization is intimately related with the theory of non-linear complex bundles on the Riemann sphere of the dimensional regularization parameter $`\epsilon `$. It also allowed us to lift both the renormalization group and the $`\beta `$-function as the asymptotic scaling in the group of diffeographisms. This used the full power of the Riemann–Hilbert decomposition together with the invariance of $`\gamma _{}(\epsilon )`$ under a change of unit of mass. This gave us a completely streamlined proof of the existence of the renormalization group and more importantly a closed formula of scattering nature, delivering the full higher pole structure of minimal subtracted counterterms in terms of the residue.
In the light of the predominant role of the residue in NCG we expect this type of formula to help us to decipher the message on space-time geometry burried in the need for renormalization.
Moreover, thanks to the previous results no longer depend upon dimensional regularization but can be formulated in any regularization or renormalization scheme. Also, we could discard a detailed discussion of anomalous dimensions, since it is an easy corollary of the knowledge of the $`\beta `$-function.
For reasons of simplicity our analysis was limited to the case of one coupling constant. The generalization to a higher dimensional space $`X`$ of coupling constants is expected to involve the same ingredients as those which appear in higher dimensional diffeomorphism groups and Gelfand-Fuchs cohomology . We left aside the detailed study of the Lie algebra of diffeographisms and its many similarities with the Lie algebra of formal vector fields. This, together with the interplay between Hopf algebras, rational homotopy theory, BRST cohomology, rooted trees and shuffle identities will be topics of future joint work.
## 7 Appendix: up to three loops
We now want to check the Hopf algebra homomorphism $`_{CM}`$ up to three loops as an example. We regard $`g_0`$ as a series in a variable $`x`$ (which can be thought of as a physical coupling) up to order $`x^6`$, making use of $`g_0=xZ_1Z_3^{3/2}`$ and the expression of the $`Z`$-factors in terms of 1PI Feynman graph. The challenge is then to confirm that the coordinates $`\delta _n`$ on $`G_2`$, implicitly defined by
$$\mathrm{log}\left[g_0(x)^{}\right]^{(n)}$$
commute with the Hopf algebra homomorphism: calculating the coproduct $`\mathrm{\Delta }_{CM}`$ of $`\delta _n`$ and expressing the result in Feynman graphs must equal the application of the coproduct $`\mathrm{\Delta }`$ applied to $`\delta _n`$ expressed in Feynman graphs.
By (2) of section 4 we write $`g_0=xZ_1Z_3^{3/2}`$,
$$Z_1=1+\underset{k=1}{\overset{\mathrm{}}{}}z_{1,2k}x^{2k},$$
$$Z_3=1\underset{k=1}{\overset{\mathrm{}}{}}z_{3,2k}x^{2k},$$
and
$$Z_g=Z_1Z_3^{3/2},z_{i,2k}_c,i=1,3,$$
as formal series in $`x^2`$. Using
$$\mathrm{log}\left(\frac{}{x}xZ_g\right)=\underset{k=1}{\overset{\mathrm{}}{}}\frac{\delta _{2k}}{(2k)!}x^{2k},$$
which defines $`\delta _{2k}`$ as the previous generators $`a_n(\varphi )`$ of coordinates of $`G_2`$, we find
$`{\displaystyle \frac{1}{2!}}\delta _2\stackrel{~}{\delta _2}`$ $`=`$ $`3z_{1,2}+{\displaystyle \frac{9}{2}}z_{3,2},`$ (1)
$`{\displaystyle \frac{1}{4!}}\delta _4\stackrel{~}{\delta _4}`$ $`=`$ $`5[z_{1,4}+{\displaystyle \frac{3}{2}}z_{3,4}]{\displaystyle \frac{9}{2}}z_{1,2}^26z_{1,2}z_{3,2}{\displaystyle \frac{3}{4}}z_{3,2}^2,`$ (2)
$`{\displaystyle \frac{1}{6!}}\delta _6\stackrel{~}{\delta _6}`$ $`=`$ $`9z_{1,2}^3+18z_{1,2}^2z_{3,2}5[3z_{1,2}z_{1,4}+{\displaystyle \frac{3}{2}}z_{3,2}z_{3,4}]`$
$`+12[z_{1,2}z_{3,2}^2z_{1,2}z_{3,4}z_{1,4}z_{3,2}]+7[z_{1,6}+{\displaystyle \frac{1}{2}}z_{3,2}^3+{\displaystyle \frac{3}{2}}z_{3,6}].`$
The algebra homomorphism $`_c`$ of section(4) is effected by expressing the $`z_{i,2k}`$ in Feynman graphs, with 1PI graphs with three external legs contributing to $`Z_1`$, and 1PI graphs with two external legs, self-energies, contributing to $`Z_3`$.
Explicitly, we have
$`z_{1,2}`$ $`=`$ $`\text{},`$
$`z_{3,2}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\text{},`$
$`z_{1,4}`$ $`=`$ $`\text{}+\text{}+\text{}+{\displaystyle \frac{1}{2}}\left[\text{}+\text{}+\text{}\right]+{\displaystyle \frac{1}{2}}\text{},`$
$`z_{3,4}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[\text{}+\text{}\right].`$
The symmetry factor
$$2=S\left(\text{}\right)$$
is most obvious if we redraw
$$\text{}.$$
Further, we have
$`z_{1,6}`$ $`=`$ $`\text{}+\text{}+\text{}+\text{}+\text{}+\text{}+\text{}+\text{}+\text{}`$
$`+{\displaystyle \frac{1}{2}}[\text{}+\text{}+\text{}+\text{}+\text{}+\text{}`$
$`+\text{}+\text{}+\text{}]+\text{}+\text{}+\text{}`$
$`+{\displaystyle \frac{1}{2}}[\text{}+\text{}+\text{}+\text{}+\text{}+\text{}`$
$`+\text{}+\text{}+\text{}]`$
$`+{\displaystyle \frac{1}{4}}\left[\text{}+\text{}+\text{}+\text{}+\text{}+\text{}\right]`$
$`+{\displaystyle \frac{1}{2}}\left[\text{}+\text{}+\text{}\right]+\text{}`$
$`+{\displaystyle \frac{1}{2}}\left[\text{}+\text{}+\text{}+\text{}+\text{}+\text{}\right]`$
$`+{\displaystyle \frac{1}{2}}\left[\text{}+\text{}+\text{}\right]`$
$`+{\displaystyle \frac{1}{2}}\left[\text{}+\text{}+\text{}+\text{}+\text{}+\text{}\right]+\text{primitive terms},`$
and
$`z_{3,6}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\text{}+{\displaystyle \frac{1}{8}}\text{}+{\displaystyle \frac{1}{4}}\text{}+{\displaystyle \frac{1}{2}}\text{}+{\displaystyle \frac{1}{2}}\text{}`$
$`+{\displaystyle \frac{1}{2}}\text{}+{\displaystyle \frac{1}{4}}\text{}+{\displaystyle \frac{1}{2}}\text{}+{\displaystyle \frac{1}{2}}\text{}+\text{}.`$
Here, primitive terms refers to 1PI three-loop vertex graphs without subdivergences. They fulfil all desired identities below trivially, and are thus not explicitly given. On the level of diffeomorphisms, we have the coproducts
$`\mathrm{\Delta }_{CM}[\delta _4]`$ $`=`$ $`\delta _41+1\delta _4+4\delta _2\delta _2,`$ (4)
$`\mathrm{\Delta }_{CM}[\delta _6]`$ $`=`$ $`\delta _61+1\delta _6+20\delta _2\delta _4+6\delta _4\delta _2`$ (5)
$`+28\delta _2^2\delta _2,`$
where we skip odd gradings.
We have to check that the coproduct $`\mathrm{\Delta }`$ of Feynman graphs reproduces these results.
Applying $`\mathrm{\Delta }`$ to the rhs of (2) gives, using the expressions for $`z_{i,k}`$ in terms of Feynman graphs,
$`\mathrm{\Delta }\left(\stackrel{~}{\delta _4}\right)`$ $`=`$ $`6\text{}\text{}+{\displaystyle \frac{9}{2}}\left[\text{}\text{}+\text{}\text{}\right]`$
$`+{\displaystyle \frac{27}{8}}\text{}\text{}+\stackrel{~}{\delta _4}1+1\stackrel{~}{\delta _4}.`$
This has to be compared with $`\stackrel{~}{\delta _4}1+1\stackrel{~}{\delta _4}+\frac{2!2!}{4!}4\stackrel{~}{\delta _2}\stackrel{~}{\delta _2}`$, which matches perfectly, as
$`\stackrel{~}{\delta _2}\stackrel{~}{\delta _2}`$ $`=`$ $`9\text{}\text{}+{\displaystyle \frac{27}{4}}\left[\text{}\text{}+\text{}\text{}\right]`$
$`+{\displaystyle \frac{81}{16}}\text{}\text{}.`$
After this warming up, let us do the check at order $`g^6`$, which will be much more demanding, as the coproduct will be noncocommutative now.
We need that $`\mathrm{\Delta }_{CM}(\delta _6)`$ in (5) is equivalent to
$`\mathrm{\Delta }_{CM}(\stackrel{~}{\delta _6})`$ $`=`$ $`\stackrel{~}{\delta _6}1+1\stackrel{~}{\delta _6}+20{\displaystyle \frac{2!4!}{6!}}\stackrel{~}{\delta _2}\stackrel{~}{\delta _4}+6{\displaystyle \frac{2!4!}{6!}}\stackrel{~}{\delta _4}\stackrel{~}{\delta _2}`$ (6)
$`+28{\displaystyle \frac{2!2!2!}{6!}}\stackrel{~}{\delta _2}^2\stackrel{~}{\delta _2}.`$
Applying the Hopf algebra homomorphism to Feynman graphs on both sides of the tensor product delivers
$`\mathrm{\Delta }\left(\stackrel{~}{\delta _6}\right)`$ $`=`$ $`\stackrel{~}{\delta _6}1+1\stackrel{~}{\delta _6}`$ (7)
$`+{\displaystyle \frac{28}{90}}\left[3\text{}+{\displaystyle \frac{9}{4}}\text{}\right]^2\left[3\text{}+{\displaystyle \frac{9}{4}}\text{}\right]`$
$`+{\displaystyle \frac{6}{15}}[5[\text{}+\text{}+\text{}+{\displaystyle \frac{1}{2}}[\text{}+\text{}+\text{}+\text{}]`$
$`+{\displaystyle \frac{3}{2}}[{\displaystyle \frac{1}{2}}\text{}+{\displaystyle \frac{1}{2}}\text{}]]`$
$`3[{\displaystyle \frac{3}{2}}\text{}\text{}+\text{}\text{}+{\displaystyle \frac{1}{16}}\text{}\text{}]]`$
$`\left[3\text{}+{\displaystyle \frac{9}{4}}\text{}\right]`$
$`+{\displaystyle \frac{20}{15}}\left[3\text{}+{\displaystyle \frac{9}{4}}\text{}\right]`$
$`5[\text{}+\text{}+\text{}+{\displaystyle \frac{1}{2}}[\text{}+\text{}+\text{}+\text{}]`$
$`+{\displaystyle \frac{3}{2}}[{\displaystyle \frac{1}{2}}\text{}+{\displaystyle \frac{1}{2}}\text{}]]`$
$`3\left[{\displaystyle \frac{3}{2}}\text{}\text{}+\text{}\text{}+{\displaystyle \frac{1}{16}}\text{}\text{}\right].`$
Multiplying this out, we find the following result
$`\mathrm{\Delta }\left(\stackrel{~}{\delta _6}\right)=\stackrel{~}{\delta _6}`$ $``$ $`1`$
$`+1`$ $``$ $`\stackrel{~}{\delta _6}`$
$`+10\text{}`$ $``$
$`+3\text{}`$ $``$
$`+{\displaystyle \frac{9}{4}}\text{}`$ $``$
$`+{\displaystyle \frac{15}{2}}\text{}`$ $``$
$`+{\displaystyle \frac{9}{4}}\left[\text{}+\text{}+\text{}\right]`$ $``$
$`+{\displaystyle \frac{15}{2}}\text{}`$ $``$ $`\left[\text{}+\text{}+\text{}\right]`$
$`+{\displaystyle \frac{9}{2}}\left[\text{}+\text{}+\text{}\right]`$ $``$
$`+15\text{}`$ $``$ $`\left[\text{}+\text{}+\text{}\right]`$
$`+{\displaystyle \frac{9}{2}}\text{}`$ $``$
$`+15\text{}`$ $``$
$`+{\displaystyle \frac{9}{2}}\text{}`$ $``$
$`+15\text{}`$ $``$
$`+{\displaystyle \frac{27}{8}}\text{}`$ $``$
$`+{\displaystyle \frac{45}{4}}\text{}`$ $``$
$`+{\displaystyle \frac{27}{8}}\text{}`$ $``$
$`+{\displaystyle \frac{45}{4}}\text{}`$ $``$
$`+3\left[\text{}+\text{}+\text{}\right]`$ $``$
$`+10\text{}`$ $``$ $`\left[\text{}+\text{}+\text{}\right]`$
$`+6\left[\text{}+\text{}+\text{}\right]`$ $``$
$`+20\text{}`$ $``$ $`\left[\text{}+\text{}+\text{}\right]`$
$`9\text{}`$ $``$
$`{\displaystyle \frac{9}{2}}\text{}\text{}`$ $``$
$`{\displaystyle \frac{3}{4}}\text{}`$ $``$
$`{\displaystyle \frac{9}{16}}\text{}`$ $``$
$`{\displaystyle \frac{27}{2}}\text{}`$ $``$
$`12\text{}`$ $``$
$`18\text{}`$ $``$
$`+{\displaystyle \frac{27}{4}}\text{}\text{}`$ $``$
$`+\text{}\text{}`$ $``$
$`+9\text{}\text{}`$ $``$
$`+{\displaystyle \frac{9}{4}}\text{}\text{}`$ $``$
$`+3\text{}\text{}`$ $``$ $`\text{}.`$ (8)
Now we have to compare with $`\mathrm{\Delta }(\stackrel{~}{\delta _6})`$, so we first apply the homomorphism to graphs and use the coproduct $`\mathrm{\Delta }`$ on them. For this, we need
$`\mathrm{\Delta }\left[z_{1,2}\right]`$ $`=`$ $`\text{}1+1\text{},`$ (9)
$`\mathrm{\Delta }\left[z_{3,2}\right]`$ $`=`$ $`\text{}1+1\text{},`$ (10)
$`\mathrm{\Delta }\left[z_{1,4}\right]`$ $`=`$ $`z_{1,4}1+1z_{1,4}+3\text{}\text{}+{\displaystyle \frac{3}{2}}\text{}\text{},`$ (11)
$`\mathrm{\Delta }\left[z_{3,2}\right]`$ $`=`$ $`z_{3,4}1+1z_{3,4}+{\displaystyle \frac{1}{2}}\text{}\text{}+\text{}\text{},`$ (12)
$`\mathrm{\Delta }\left[z_{1,6}\right]`$ $`=`$ $`z_{1,6}1+1z_{1,6}+3\text{}\left[\text{}+\text{}+\text{}\right]`$
$`+3\left[\text{}+\text{}+\text{}\right]\text{}`$
$`+2\text{}\left[\text{}+\text{}+\text{}\right]`$
$`+3\text{}\text{}\text{}+\text{}\left[\text{}+\text{}+\text{}\right]`$
$`+{\displaystyle \frac{3}{2}}\text{}\text{}\text{}`$
$`+{\displaystyle \frac{3}{2}}\text{}\left[\text{}+\text{}+\text{}\right]+{\displaystyle \frac{3}{2}}\text{}\text{}`$
$`+{\displaystyle \frac{1}{2}}\text{}\left[\text{}+\text{}+\text{}\right]+{\displaystyle \frac{3}{2}}\text{}\text{}`$
$`+{\displaystyle \frac{3}{2}}\text{}\left[\text{}+\text{}+\text{}\right]`$
$`+\text{}\left[\text{}+\text{}+\text{}\right]`$
$`+{\displaystyle \frac{9}{2}}\text{}\text{}\text{}+{\displaystyle \frac{3}{2}}\text{}\left[\text{}+\text{}+\text{}\right]`$
$`+{\displaystyle \frac{3}{2}}\left[\text{}+\text{}+\text{}\right]\text{}+{\displaystyle \frac{3}{2}}\text{}\text{}`$
$`+{\displaystyle \frac{5}{2}}\text{}\text{}+3\text{}\text{},`$
$`\mathrm{\Delta }\left[z_{3,6}\right]`$ $`=`$ $`z_{3,6}1+1z_{3,6}+{\displaystyle \frac{1}{2}}\text{}\text{}`$
$`+{\displaystyle \frac{1}{2}}\text{}\text{}+{\displaystyle \frac{1}{4}}\text{}\text{}`$
$`+{\displaystyle \frac{1}{8}}\text{}\text{}\text{}`$
$`+{\displaystyle \frac{1}{2}}\text{}\text{}+{\displaystyle \frac{1}{4}}\text{}\text{}\text{}+\text{}\text{}`$
$`+{\displaystyle \frac{1}{2}}\text{}\text{}+\text{}\text{}+{\displaystyle \frac{1}{2}}\text{}\text{}\text{}`$
$`+\text{}\text{}+{\displaystyle \frac{1}{4}}\text{}\text{}`$
$`+{\displaystyle \frac{1}{2}}\text{}\text{}+\text{}\text{}`$
$`+{\displaystyle \frac{1}{2}}\left[\text{}+\text{}\right]\text{}+\text{}\text{}+\text{}\text{}`$
$`+\left[\text{}+\text{}+\text{}\right]\text{}.`$
It is now only a matter of using the rhs of (7) for $`\stackrel{~}{\delta _6}`$ to confirm that we reproduce the result (8). For example, for the contribution to $`\text{}\text{}`$ in $`\mathrm{\Delta }(\stackrel{~}{\delta _6})`$ we find
$$\frac{5\times 3}{2}\text{}\text{}+7\times \frac{5}{2}\text{}\text{}=10\text{}\text{},$$
as desired. Similarly, one checks all of the 32 tensorproducts of (8).
## Acknowledgements
Both authors thank the IHES for generous support during this collaboration. D.K. is grateful to the DFG for a Heisenberg Fellowship.
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# Cosmic crystallography: the euclidean isometries
## 1 Introduction
Various methods have been proposed to investigate the shape of the universe, and cosmic crystallography (CC) is one of them . As CC ponders, if the universe is multiply connected (MC) then multiple images of a same cosmic object may be seen in the sky. The separations between these images are correlated by the geometry and the topology of the spacetime; so if one selects a catalog of $`n`$ observed images of various cosmic objects and performs a histogram of the $`n(n1)/2`$ separations between these $`n`$ images, then the existing correlations must somehow show up.
It was recently examined in what respects the histogram for a multiply connected observed universe should differ from that of a simply connected (SC) one, with same geometry and radius. It was found that each isometry of the MC universe individually imprints either a small localized deformity on the histogram of the SC universe , or a sharp spike if the isometry is a Clifford translation .
Since each observed universe model has a unique pair separation normalized histogram, a strategy to unveil the cosmic topology seems straightforward: one should compare the histogram obtained from observational astronomy with histograms obtained from computer simulated universe models with prescribed topologies.
Both types of histograms (observational and simulated) are infected with statistical noises, and methods to reduce these noises are desirable. A suggestion was made, to replace the histogram related to the SC component of the simulated model by an exact continuous probability density function. For each geometry with constant curvature the corresponding function was then derived ; however, appropriate functions were still lacking, to replace the histograms related to each isometry component .
In the present study we derive some of these functions, namely the normalized pair separation probability densities for the euclidean isometries. Following the prescriptions of ref. , if we now merge these new functions with the functions already given in then we obtain noiseless normalized probability densities more suitable for comparing with the normalized observational histograms. For completeness we also examine the euclidean isometries with fixed point, as well as the orientation reversing isometries.
## 2 The euclidean isometries
The isometries of the euclidean space $`E^3`$ are the (pure) translations $`t`$, (pure) rotations $`\omega `$, screw motions $`t\omega `$, (pure) reflections $`ϵ`$, and glide reflections $`ϵt`$. While $`t`$, $`\omega `$, and $`t\omega `$ are orientation preserving, $`ϵ`$ and $`ϵt`$ reverse the orientation of $`E^3`$; and while $`t`$, $`t\omega `$, and $`ϵt`$ act freely on $`E^3`$, $`\omega `$ and $`ϵ`$ have fixed points. We study all five isometries for the sake of completeness, although cosmic crystallography is presently interested on the orientation preserving isometries without fixed points, namely $`t`$ and $`t\omega `$ only.
Our general approach to investigate an isometry $`g`$ in cosmic crystallography is first consider a solid ball $``$ with radius $`a`$, then apply $`g`$ to $``$ thus producing a new solid ball $`_g`$, next consider the set of points $`P`$ $``$ $``$ whose corresponding transformed points $`P_g`$ are also in $``$, finally randomly select a pair ($`P`$, $`P_g`$) and ask for the probability $`𝒫_g^{}(l)dl`$ that their separation lies between the values $`l`$ and $`l+dl`$. The normalization condition
$`{\displaystyle _0^{\mathrm{}}}𝒫_g^{}(l)𝑑l=1`$ (2.1)
must be obeyed.
It is clear that the balls $``$ and $`_g`$ need intersect, so the separation $`m`$ between their centers $`C`$ and $`C_g`$ has to satisfy
$`m<2a.`$ (2.2)
The intersection $`_g`$ is a rotationally symmetric solid lens whose diameter, thickness, and volume are (see figure 1)
$`D`$ $`=`$ $`\sqrt{4a^2m^2},`$
$`T`$ $`=`$ $`2am,`$ (2.3)
$`V_g^{}`$ $`=`$ $`{\displaystyle \frac{\pi }{12}}(2am)^2(4a+m).`$
Figure 1 The solid balls $``$ and $`_g`$ both with radius $`a`$ intersect in a solid lens with diameter $`D`$ and thickness $`T`$. $`\mathrm{}`$
## 3 Translations, screw mo-tions, and rotations
Translations
If the euclidean space is subjected to a nonzero translation $`t`$, the probability density that a point $`P`$ be displaced a distance $`l`$ is clearly
$`𝒫_t(l)=\delta (lt),`$ (3.1)
the Dirac delta; the density satisfies the normalization condition (2.1), and does not depend on the ball $``$.
Basics on screw motions of a solid ball
In $`E^3`$, imagine a straight line $``$ placed at a distance $`b`$ from the center $`C`$ of the solid ball $``$ with radius $`a`$; $``$ and $``$ may intersect ($`b<a`$), be tangent ($`b=a`$) or be disjoint ($`b>a`$); see figure 2.
Figure 2 Relative positions of a solid ball $``$ with radius $`a`$ and a straight line $``$ at a distance $`b`$ from the center of $``$. $`\mathrm{}`$
Now consider a screw motion $`g`$ of the ball, with nonzero translation $`t`$ parallel to the line $``$ and nonzero rotation $`\omega `$ around the line; for our purposes the senses of $`t`$ and $`\omega `$ are irrelevant, so for definiteness and simplicity we assume $`t>0`$ and $`0<\omega \pi `$.
The separation $`m`$ between the centers $`C`$ and $`C_g`$ of the intersecting balls $``$ and $`_g`$ is
$`m=\sqrt{t^2+4b^2\mathrm{sin}^2\omega /2},`$ (3.2)
and the condition $`m<2a`$ implies the constraint
$`t^2+4b^2\mathrm{sin}^2\omega /2<4a^2`$ (3.3)
between the four independent parameters $`a`$, $`b`$, $`t`$, and $`\omega `$; see figure 3.
Figure 3 A solid ball (not displayed) with center $`C`$ rotates $`\omega `$ around the line $``$ (the $`z`$ axis) and translates $`t`$ parallel to the axis, eventually reaching the position centered at $`C_g`$ ; both centers $`C`$ and $`C_g`$ are at a distance $`b`$ from the axis, and their separation $`m`$ must be smaller than the sum $`2a`$ of the two radii in order that the balls intersect. $`\mathrm{}`$
We next introduce the auxiliary variable $`r`$ and a function $`Q_g^{}(r)`$ which help simplify our study. Call $`r`$ the separation between the point $`P`$ (also $`P_g`$) and the axis $``$ (see figure 4); it is related to $`l`$ (the separation from $`P`$ to $`P_g`$) through
$`l^2=t^2+4r^2\mathrm{sin}^2\omega /2;`$ (3.4)
for each isometry ($`t,\omega `$) this is a bijective relation between $`l`$ and $`r`$, so the probability density $`𝒫_g^{}(l)`$ of finding a $`g`$-pair with mutual separation $`l`$ parallels the akin probability density $`Q_g^{}(r)`$ of finding a $`g`$-pair whose members (both in $``$) are at a distance $`r`$ from the axis of the motion:
$`𝒫_g^{}(l)dl=Q_g^{}(r)dr.`$ (3.5)
Figure 4 Under the screw motion $`g`$ with axis $``$ , translation $`t`$ , and rotation $`\omega `$ , a point $`P`$ separated $`r`$ from $``$ moves to the position $`P_g`$ at the distance $`l`$ given by eq. $`(\text{3.4})`$. $`\mathrm{}`$
Clearly the normalization condition
$`{\displaystyle _0^{\mathrm{}}}Q_g^{}(r)𝑑r=1`$ (3.6)
is also to be satisfied. Once we obtain $`Q_g^{}(r)`$, and since from (3.4) we have
$`r={\displaystyle \frac{\sqrt{l^2t^2}}{2\mathrm{sin}\omega /2}},`$ (3.7)
then we will finally compute
$`𝒫_g^{}(l)={\displaystyle \frac{dr}{dl}}𝒬_g^{}[r(l)].`$ (3.8)
To obtain $`Q_g^{}(r)`$ we first consider a sufficiently long cylinder $`𝒞_r`$ with radius $`r`$ and axis $``$. A short reflection gives that the probability density $`Q_g^{}(r)`$ is linearly proportional to the area $`S_g^{}(r)`$ of the intersection $`_g𝒞_r`$; the coefficient of proportionality is the inverse of the volume $`V_g^{}`$ of the intersection $`_g`$ , eq.(2):
$`Q_g^{}(r)={\displaystyle \frac{area(_g𝒞_r)}{vol(_g)}}={\displaystyle \frac{S_g^{}(r)}{V_g^{}}}.`$ (3.9)
To find the area $`S_g^{}(r)`$ we first assume the line $``$ (the axis of the screw motion) along the cartesian $`z`$-axis; the cylinder $`𝒞_r`$ with axis $``$ and radius $`a`$ then has equation
$`x^2+y^2=r^2.`$ (3.10)
Still without loss of generality we take the center $`C`$ of the solid ball $``$ with radius $`a`$ at the cartesian position ($`b`$, 0, 0); the points of $``$ then satisfy
$`(xb)^2+y^2+z^2a^2.`$ (3.11)
Finally the center $`C_g`$ of the solid ball $`_g`$ is at the cartesian position ($`b`$ cos$`\omega `$, $`b`$ sin$`\omega `$, $`t`$), so the points of $`_g`$ satisfy
$`(xb\mathrm{cos}\omega )^2`$ $`+`$ $`(yb\mathrm{sin}\omega )^2`$ (3.12)
$`+`$ $`(zt)^2a^2.`$
$`S_g^{}(r)`$ is then the area of the surface whose points ($`x,y,z`$) satisfy (3.10), (3.11), and (3.12) simultaneously. To visualize this surface $`_g𝒞_r`$ we first consider the surface $`𝒞_r`$ , then the similar surface $`_g𝒞_r`$, and finally the combined intersection $`(𝒞_r)(_g𝒞_r)`$.
The form of the intersection $`𝒞_r`$ depends on the relative values of $`a,b`$, and $`r`$ (see figure 5):
(i) if $`0<r<ab`$ then all generatrices of $`𝒞_r`$ penetrate into $``$: the intersection is a topological ring (an annulus, a compact cylinder, a disc with a hole in it);
(ii) if $`|ab|<r<a+b`$ then only a part of the generatrices of $`𝒞_r`$ goes through $``$: the intersection is a topological disc;
(iii) if $`a+b<r`$, or if $`0<r<ba`$, then there is no intersection.
Figure 5 (A) The solid sphere $``$ and the sufficiently long cylindrical surface $`𝒞_r`$ , in perspective; in the other drawings the line of sight is the vertical;
(i) the intersection $`𝒞_r`$ is a topological ring when $`0<r<ab`$ (loosely saying, $`a`$ is too large);
(ii) $``$ and $`𝒞_r`$ intersect in a topological disc when $`|ab|<r<a+b`$ (then $`a,b`$, and $`r`$ may form a triangle);
(iii) $``$ and $`𝒞_r`$ do not intersect when $`a+b<r`$ nor when $`0<r<ba`$ (loosely saying, either $`r`$ or $`b`$ is too large). $`\mathrm{}`$
Since the intersections $`𝒞_r`$ are drawn on the cylinder $`𝒞_r`$ itself, we use the cylindrical coordinates $`\varphi `$ (the azimuthal angle) and $`z`$ (altitude) to visualize them. A ring-like intersection $`𝒞_r`$ is then bounded by the two curves (see figure 6)
$`z=\pm \sqrt{a^2b^2r^2+2br\mathrm{cos}\varphi },`$ (3.13)
where $`\varphi (\pi ,\pi ]`$; the positive curve oscillates between the extreme values
$`z_{max}`$ $`=`$ $`\sqrt{a^2(br)^2},`$
$`z_{min}`$ $`=`$ $`\sqrt{a^2(b+r)^2}.`$ (3.14)
Figure 6 The ring-like intersection $`𝒞_r`$ , eq.(3.13), when $`a=4`$ , $`b=19/5`$ , and $`r=1/10`$ . The boundaries at $`\varphi =\pi `$ and $`\varphi =\pi `$ are identified. The upper and lower boundaries are not sinusoidal. $`\mathrm{}`$
In the disc-like intersection $`𝒞_r`$ (see figure 7) the surface is again read from (3.13), but now $`\varphi [\varphi _{max},\varphi _{max}]`$ with
$`\varphi _{max}=\mathrm{cos}^1{\displaystyle \frac{b^2+r^2a^2}{2br}};`$ (3.15)
while $`z_{max}`$ is still given in (3), $`z_{min}=0`$ now.
Figure 7 The disc-like intersection $`𝒞_r`$ , eq.(3.13), when $`a=1`$ and $`b=r=1/\sqrt{2}`$ ; it is an oval (not an ellipse) centered at $`\varphi =z=0`$ , with $`\varphi _{max}=\pi /2`$ and $`z_{max}=1`$ . $`\mathrm{}`$
A disc-like intersection $`_g𝒞_r`$ is a $`g`$-transported copy of the corresponding intersection $`𝒞_r`$, parallelly dragged on the cylinder: while $`𝒞_r`$ is centered at ($`\varphi =0,z=0`$), the center of $`_g𝒞_r`$ is at ($`\varphi =\omega ,z=t`$).
Similar statements are true for the ring-like intersections $`_g𝒞_r`$. By inspection it is then easy to visualize the form of the combined intersection $`_g𝒞_r`$ in both disc-like and ring-like cases. Since $`\varphi `$ has a cyclic character, the nonempty intersection of two discs on $`𝒞_r`$ can be either 1 or 2 discs, while two rings intersect in either 1 or 2 discs, or 1 ring (see figures 8 and 9).
Figure 8 Instance of disconnected intersection $`_g𝒞_r`$ (dashed areas) when $`𝒞_r`$ (and $`_g𝒞_r`$) are ring-like; here $`a=4.505,b=2,r=2.5,t=5.3`$ , and $`\omega =\pi `$ . $`\mathrm{}`$
Figure 9 Instance of disconnected intersection $`_g𝒞_r`$ (dashed areas) when $`𝒞_r`$ (and $`_g𝒞_r`$) are disc-like; here $`a=3.5,b=r=2,t=0.7`$ , and $`\omega =14\pi /15`$ . $`\mathrm{}`$
We are now ready to evaluate the area $`S_g^{}(r)`$ of the combined intersection $`_g𝒞_r`$; given the fixed values of $`a,b,t`$, and $`\omega `$, then for each value of $`r`$ we need to integrate the differential
$`dS_g^{}(r)=\mathrm{\Delta }z.rd\varphi ,`$ (3.16)
where $`\mathrm{\Delta }z(r,\varphi )`$ is the varying height of the intersection, and where the limits of integration in $`\varphi `$ may vary with $`r`$ (see eq.(3.15)).
As is evident from figures 8 and 9, the combined intersection is always bounded by one of the two curves
$`z_1(\varphi )=\sqrt{a^2b^2r^2+2br\mathrm{cos}\varphi },`$ (3.17)
$`z_3(\varphi )=t+\sqrt{a^2b^2r^2+2br\mathrm{cos}(\varphi \omega )},`$ (3.18)
from the upper side, and one of the two curves
$`z_2(\varphi )=z_1(\varphi ),`$ (3.19)
$`z_4(\varphi )=t\sqrt{a^2b^2r^2+2br\mathrm{cos}(\varphi \omega )},`$ (3.20)
from the lower side. Further, we only have nonempty combined intersection when $`z_1(\varphi )>z_4(\varphi )`$. We then obtain
$`S_g^{}(r)`$ $`=`$ $`r{\displaystyle _{\varphi _{max}}^{\varphi _{max}}}\mathrm{\Theta }(z_1z_4)`$
$`\times `$ $`\left[min(z_1,z_3)max(z_2,z_4)\right]d\varphi ,`$ (3.21)
with $`\varphi _{max}`$ given by (3.15) for disc-like intersections, and $`\varphi _{max}=\pi `$ for ring-like intersections; $`\mathrm{\Theta }`$ is the step function with values 0 and 1.
Elliptic integrals
As is seen from eqs. (3.17)-(3), the area $`S_g^{}(r)`$ involves terms of the form
$`A(r)={\displaystyle _0^{\alpha (r)}}\sqrt{a^2b^2r^2+2br\mathrm{cos}\varphi }𝑑\varphi ,`$ (3.22)
with $`\mathrm{cos}(\varphi \omega )`$ sometimes replacing $`\mathrm{cos}\varphi `$ in the integrand. To evaluate these terms we define
$$\begin{array}{ccc}f^2(r)& =& a^2(br)^2,\\ & & \\ k^2(r)& =& \frac{4br}{f^2(r)},\end{array}$$
(3.23)
after noting that the condition $`f^2(r)>0`$ is satisfied by all nonempty intersections $`𝒞_r`$. Then introduce the (incomplete) elliptic integral of $`2^{nd}`$ kind
$`E(\gamma ,k)={\displaystyle _0^\gamma }{\displaystyle \frac{\sqrt{1k^2x^2}}{\sqrt{1x^2}}}𝑑x,`$ (3.24)
and finally have the integral (3.22) expressed as
$`A(r)=2f(r)E(\mathrm{sin}{\displaystyle \frac{1}{2}}\alpha (r),k(r)),`$ (3.25)
where $`f=\sqrt{f^2}`$ and $`k=\sqrt{k^2}`$ . When $`\gamma =1`$ in (3.24), or equivalently $`\alpha =\pi `$ in (3.22), we have the complete elliptic integral of $`2^{nd}`$ kind
$`E(k)=E(1,k);`$ (3.26)
the functions $`E(k)`$ only appear in the ring-like intersections.
Which of $`z_1`$ or $`z_3`$ is minimum in (3), and which of $`z_2`$ or $`z_4`$ is maximum and whether $`z_1>z_4`$or$`z_1<z_4`$ , usually depends on the angle $`\varphi `$; clearly also the values of the parameters $`a,b,t,`$ and $`\omega `$ are relevant, well as the radius $`r`$.
A sample screw motion with $`b0`$
We consider the case with $`a=t=1/\sqrt{2},b=1/\sqrt{8},\omega =\pi `$; the intersection $`𝒞_r`$ is then ring-like for $`0<r<1/\sqrt{8}`$ (equivalently $`1/\sqrt{2}<l<1`$ ), is disc-like for $`1/\sqrt{8}<r<3/\sqrt{8}`$ (or $`1<l<\sqrt{5}`$), and is empty for larger values of $`r`$ and $`l`$. The combined intersection $`_g𝒞_r`$ has a ring-like regime for $`0<r<1/\sqrt{8}`$ $`(1/\sqrt{2}<l<1)`$ , then a two separate discs regime for $`1/\sqrt{8}<r<1/2`$ $`(1<l<\sqrt{3/2})`$, and is empty for larger $`r`$ and $`l`$.
We first consider the ring-like regime: preliminarly find, from (3.23),
$$\begin{array}{ccc}f^2(r)& =& \frac{3}{8}+\frac{r}{\sqrt{2}}r^2,\\ & & \\ k^2(r)& =& \frac{\sqrt{2}r}{f^2},\end{array}$$
(3.27)
and the area (3) of the intersection $`_g𝒞_r`$ ,
$`S_g^{}(r)=r\left(8fE(k)\pi \sqrt{2}\right);`$ (3.28)
now since
$`r={\displaystyle \frac{1}{2}}\sqrt{l^21/2},dr/dl=l/(4r),`$ (3.29)
the probability density (3.8) is, for $`1/\sqrt{2}<l<1`$ ,
$`𝒫_g^{}(l)={\displaystyle \frac{l}{4V_g^{}}}\left(8fE(k)\pi \sqrt{2}\right),`$ (3.30)
where the volume of $`_g`$ is
$`V_g^{}=\pi (4\sqrt{2}5)/12.`$ (3.31)
Clearly $`f(r)`$ and $`k(r)`$ in (3.30) need consider the dependence $`r(l)`$ given in (3.29).
To describe the two-discs regime we further define
$`g(r)`$ $`=`$ $`\sqrt{1/(2r)^21},`$
$`g_\pm (r)`$ $`=`$ $`\sqrt{{\displaystyle \frac{1\pm g}{2}}},`$ (3.32)
then obtain for the area (3)
$`S_g^{}(r)`$ $`=`$ $`8rf\left(E(g_+,k)E(g_{},k)\right)`$ (3.33)
$`+`$ $`\sqrt{2}r(2\mathrm{cos}^1g\pi ),`$
and finally, for $`1<l<\sqrt{3/2}`$ , the probability density
$`𝒫_g^{}(l)`$ $`=`$ $`{\displaystyle \frac{l}{4V_g^{}}}[8f(E(g_+,k)E(g_{},k))`$ (3.34)
$`+`$ $`\sqrt{2}(2\mathrm{cos}^1g\pi )].`$
In (3.34) we clearly need replace $`r`$ by its value $`r(l)`$ in (3.29) in the functions $`f,g,g_+,g_{}`$, and $`k`$. In figure 10 we reproduce the function $`𝒫_g^{}(l)`$ for this sample screw motion, for the whole interval $`1/\sqrt{2}<l<\sqrt{3/2}`$.
Figure 10 The probability density $`𝒫_g^{}(l)`$ for the screw motion of the solid ball $``$ with $`a=2b=t=1/\sqrt{2}`$ , and $`\omega =\pi `$. For $`l[1/\sqrt{2},1]`$ we use the ring-like equation (3.30), while for $`l[1,\sqrt{3/2}]`$ the two-discs equation (3.34) is used. The integrated area is $`1`$. $`\mathrm{}`$
The special cases with $`b=0`$
In these cases the axis $``$ of the screw motion contains the centers $`C`$ and $`C_g`$ of the solid balls $``$ and $`_g`$; the intersections $`𝒞_r`$ and $`_g𝒞_r`$ are ring-like and have constant width, as well as the combined intersections $`_g𝒞_r`$. From (3.2) we have $`m=t`$, then (2) gives
$`V_g^{}={\displaystyle \frac{\pi }{12}}(2at)^2(4a+t)[b=0]`$ (3.35)
provided $`t<2a`$; this sine qua non is assumed whenever $`b=0`$. The equations (3.17)-(3.20) simplify to
$`z_1=z_2`$ $`=`$ $`\sqrt{a^2r^2},`$
$`z_3`$ $`=`$ $`t+z_1,z_4=tz_1,`$ (3.36)
provided $`r<a`$; since we assumed $`t>0`$, we have $`min(z_1,z_3)=`$ $`\sqrt{a^2r^2}`$ and $`max(z_2,z_4)=t\sqrt{a^2r^2}`$, while the condition $`z_1>z_4`$ implies
$`r<r_{max}=\sqrt{a^2t^2/4}.`$ (3.37)
The area of the combined intersection is then
$`S_g^{}(r)=2\pi r(2\sqrt{a^2r^2}t)[b=0]`$ (3.38)
whenever positive, otherwise it is zero. Finally the probability density is, when $`b=0`$,
$`𝒫_g^{}(l)={\displaystyle \frac{dr}{dl}}{\displaystyle \frac{S_g^{}(r)}{V_g^{}}}`$
$`={\displaystyle \frac{6l(2\sqrt{a^2r^2}t)}{(2at)^2(4a+t)}}\mathrm{csc}^2\omega /2`$ (3.39)
if positive, otherwise it is zero; in (3) one clearly has to substitute $`r`$ for its expression eq.(3.7). As a matter of fact we have $`𝒫_g^{}(l)0`$ when $`b=0`$ only for $`l_{min}<l<l_{max}`$, where
$`l_{min}=t,`$
$`l_{max}=\sqrt{t^2\mathrm{cos}^2\omega /2+4a^2\mathrm{sin}^2\omega /2},`$ (3.40)
and we have in these extreme limits
$`𝒫_g^{}(l_{min})=6t(2at)^1(4a+t)^1\mathrm{csc}^2\omega /2,`$
$`𝒫_g^{}(l_{max})=0;`$ (3.41)
a sample graph of $`𝒫_g^{}(l)`$ for screw motions with $`b=0`$ is given in figure 11.
Figure 11 The probability density $`𝒫_g^{}(l)`$ for the screw motion of the solid ball $``$ with $`b=0,a=1/\sqrt{2},t=1`$ , and $`\omega =\pi `$ , eq.(3). The integrated area is $`1`$. Note that $`l_{min}=1`$ , $`l_{max}=\sqrt{2}`$ , and $`𝒫_g^{}(1)=6(3+\sqrt{2})/73.8`$ . $`\mathrm{}`$
Rotations
A rotation is the limit of a screw motion when the translation tends to zero, so it is formally described by simply setting $`t=0`$ in the appropriate preceding equations. In particular, eq. (3.4) becomes $`l=2r\mathrm{sin}\omega /2`$. Rotations of the solid sphere $``$ are classified into two categories:
(i) if $`a<b`$ then $``$ is exempt of fixed points;
(ii) if $`b<a`$ then the axis $``$ traverses $``$, which has fixed points.
In category $`a<b`$ the intersection condition (3.3) clearly must be satisfied (see figure 12),
$`b\mathrm{sin}\omega /2<a.`$ (3.42)
Figure 12 The solid lens $`_g`$ when the solid ball $``$ rotates $`\omega `$ around the $`z`$-axis and $`b\mathrm{sin}\omega /2<a<b`$ ; the lens lays between the radial positions $`r_{min}`$ and $`r_{max}`$ given by eq. (3). $`\mathrm{}`$
The lens $`_g`$ lays between the radial positions
$`r_{min}=r_mD/2,r_{max}=r_m+D/2,`$
$`r_m=b\mathrm{cos}\omega /2,`$ (3.43)
where $`r_m`$ is the radial coordinate of the center of the lens and $`D`$ is its diameter (2) with
$`m=2b\mathrm{sin}\omega /2.`$ (3.44)
All intersections $`𝒞_r`$ (and $`_g𝒞_r`$) are topological discs. Since $`z_1>z_4`$ and $`z_4=z_3`$ in rotations, the area (3) has the simpler expression
$`S_g^{}(r)=2r{\displaystyle _{\varphi _{max}}^{\varphi _{max}}}min(z_1,z_3)𝑑\varphi .`$ (3.45)
When $`b\mathrm{sin}\omega /2<a<b`$ and $`r_mD/2<r<r_m+D/2`$ we obtain
$`S_g^{}(r)=8rf[E(1/k,k)E(\mathrm{sin}\omega /4,k)],`$ (3.46)
with $`f(r)`$ and $`k(r)`$ given in (3.23). The probability density (3.8) is easily written if we take in succession $`l=2r\mathrm{sin}\omega /2`$ , then $`Q_g^{}(r)=S_g^{}(r)/V_g^{}`$ with (3.46), finally $`V_g^{}`$ in (2) with $`m=2b\mathrm{sin}\omega /2`$ ; in figure 13 a sample graph of $`𝒫_g^{}(l)`$ for rotation with $`a<b`$ is given.
Figure 13 The probability density $`𝒫_g^{}(l)`$ for a rotation $`\omega `$ of the solid ball $``$ with radius $`a<b`$ . Here $`a=1,b=2,`$and$`\omega =\pi /6`$ . The integrated area is 1. $`\mathrm{}`$
In category $`b<a`$ the intersections $`𝒞_r`$ are either ring-like or disc-like, and the combined intersections $`_g𝒞_r`$ can be of three types according to the relative values of $`a,b,\omega ,`$ and $`r`$ (see figure 14):
(i) a ring, when $`0<r<ab`$ ; then
$`S_g^{}(r)=8rf[\mathrm{\hspace{0.17em}2}E(k)`$
$`E(\mathrm{sin}\omega /4,k)E(\mathrm{cos}\omega /4,k)];`$ (3.47)
(ii) a pair of discs, when $`ab<r<D/2r_m`$ ; then
$`S_g^{}(r)=8rf[\mathrm{\hspace{0.17em}2}E(1/k,k)`$
$`E(\mathrm{sin}\omega /4,k)E(\mathrm{cos}\omega /4,k)];`$ (3.48)
(iii) a disc, when $`D/2r_m<r<D/2+r_m`$ ; then
$`S_g^{}(r)=8rf\left[E(1/k,k)E(\mathrm{sin}\omega /4,k)\right].`$ (3.49)
Figure 14 The solid ball $``$ with center $`C`$ and radius $`a>b`$ rotates $`\omega `$ around the $`z`$-axis and gives the new ball $`_g`$ with center $`C_g`$ . The solid lens $`_g`$ is intersected by the cylindrical surface $`𝒞_r`$ with radius $`r`$ and axis along the $`z`$-axis. The intersection $`_g𝒞_r`$ is a ring if $`r`$ is small $`(0r<ab)`$, is a pair of discs for $`ab<r<D/2r_m`$, and is a single disc when $`r`$ is larger $`(D/2r_m<r<D/2+r_m)`$. $`\mathrm{}`$
In figures 15 and 16 sample graphs of the probability $`𝒫_g^{}(l)`$ for rotation with $`b<a`$ are given, for the entire interval $`0<r<D/2+r_m`$ .
Figure 15 The probability density $`𝒫_g^{}(l)`$ for a (pure) rotation of the solid ball $``$ with radius $`a>b`$ . Here $`a=2,b=1.3`$ , and $`\omega =\pi /3`$. The integrated area is 1. $`\mathrm{}`$
Figure 16 The probability density $`𝒫_g^{}(l)`$ for a (pure) rotation of the solid ball $``$ with radius $`a>b`$ . Here $`a=1,b=5/7`$ , and $`\omega =2\mathrm{sec}^15`$. The region where each of the three expressions (3)-(3.49) for $`S_g^{}(r)`$ is used is indicated. The integrated area is 1. $`\mathrm{}`$
Rotations with $`b=0`$
In the category $`b<a`$ the special cases $`b=0`$ deserve a few words. The solid spheres $``$ and $`_g`$ coincide, and we easily find that
$`l=2r\mathrm{sin}\omega /2,S_g^{}(r)=4\pi r\sqrt{a^2r^2},`$
$`V_g^{}=4\pi a^3/3.`$ (3.50)
The probability density is then
$`𝒫_g^{}(l)={\displaystyle \frac{3}{8}}{\displaystyle \frac{l\sqrt{4a^2\mathrm{sin}^2\omega /2l^2}}{a^3\mathrm{sin}^3\omega /2}}`$ (3.51)
whenever $`0l2a\mathrm{sin}\omega /2`$ ; a sample graph is given in figure 17.
Figure 17 The probability density $`𝒫_g^{}(l)`$ for a (pure) rotation with $`b=0`$. Here $`a=1`$ , and $`\omega =\pi /6`$ . The integrated area is 1. $`\mathrm{}`$
## 4 Reflections and glide reflections
In $`E^3`$, let a solid ball $``$ with radius $`a`$ have its center $`C`$ at the cartesian position $`(b,0,0)`$. Next consider a glide reflection $`g`$ on the plane $`x=0`$ (the plane $`𝒳_0`$) with nonzero translation $`t`$ in the direction $`+z`$. The center $`C_g`$ of the glide reflected ball $`_g`$ is at the cartesian position $`(b,0,t)`$, and the separation $`m`$ between $`C`$ and $`C_g`$ is
$`m=\sqrt{4b^2+t^2}.`$ (4.1)
The condition of intersection $`_g\varphi `$ implies the constraint
$`4b^2+t^2<4a^2`$ (4.2)
between the three independent parameters $`a,b,`$ and $`t`$; a fortiori $`b<a`$, so both balls intersect the plane $`𝒳_0`$; see figure 18.
Figure 18 The solid ball $``$ with center $`C`$ and radius $`a`$ is first reflected on the plane $`x=0`$ a distance $`b<a`$ apart, then translated $`t`$ upwards. $`\mathrm{}`$
Now randomly choose a point $`P=(x,y,z)`$ of $``$, such that $`P_g=(x,y,z+t)`$ is also in $``$; clearly $`P_g`$. $`l`$ being the separation from $`P`$ to $`P_g`$, we ask for the probability density $`𝒫_g^{}(l)`$ as described in sec. 2; we readily find
$`𝒫_g^{}(l)`$ $`=`$ $`2{\displaystyle \frac{area(_g𝒳_x)}{vol(_g)}}`$ (4.3)
$`=`$ $`2{\displaystyle \frac{S_g^{}(x)}{V_g^{}}},`$
where $`𝒳_x`$ is one of the two planes
$`x=\pm {\displaystyle \frac{1}{2}}\sqrt{l^2t^2};`$ (4.4)
the multiplying factor 2 in (4.3) accounts for these two possibilities, and the volume $`V_g^{}`$ is given in (2) .
The intersections $`_g𝒳_x`$ are of two types, depending on whether or not the plane $`𝒳_x`$ intersects the equator $``$ of the lens $`_g`$ (see figure 19):
Figure 19 The solid lens $`_g`$ and the two types of intersections $`_g𝒳_x`$ when the isometry is a glide reflection. In the equator $``$ of the lens, $`M`$ is the point placed farthest from the reflector plane $`x=0`$ . In (A), $`M`$ is below the center $`C_g`$ of the glide reflected ball $`_g`$ , then $`x_{max}=x_M=\mathrm{\Delta }`$ given in (4); only bicircular sections (i) are possible. In (B), $`M`$ is above $`C_g`$ , then $`x_{max}=ab>x_M`$ and two types of vertical cross-sections can occur: (i) bicircular, and (ii) circular. $`\mathrm{}`$
(i) a bicircular disc, when $`𝒳_x\varphi `$ ; the disc is enclosed by two circles with unequal radii
$`r_\pm (x)=\sqrt{a^2(bx)^2}.`$ (4.5)
It has area
$`S_g^{}(x)=r_{}^2\mathrm{cos}^1{\displaystyle \frac{h_{}}{r_{}}}+r_+^2\mathrm{cos}^1{\displaystyle \frac{h_+}{r_+}}{\displaystyle \frac{t\delta }{2}},`$ (4.6)
whenever $`|x|<\mathrm{\Delta }`$ , where
$`h_\pm (x)=t/2\pm 2bx/t,`$ (4.7)
$`\mathrm{\Delta }={\displaystyle \frac{t}{2}}\sqrt{{\displaystyle \frac{a^2}{b^2+t^2/4}}1},`$
$`\delta (x)=2\sqrt{(1+4b^2/t^2)(\mathrm{\Delta }^2x^2)}.`$ (4.8)
(ii) a circular disc with radius
$`\rho =\sqrt{a^2+(b+|x|)^2},`$ (4.9)
which appears only when
$`q=4b(ab)t^2>0,\mathrm{\Delta }<|x|<ab;`$ (4.10)
the area of the disc is clearly $`\pi \rho ^2`$ .
Collecting together preceding terms we finally obtain
$`𝒫_g^{}(l)={\displaystyle \frac{l}{\sqrt{l^2t^2}}}{\displaystyle \frac{1}{V_g^{}}}[\mathrm{\Theta }(\mathrm{\Delta }|x|)S_g^{}(x)`$
$`+\mathrm{\Theta }(q)\mathrm{\Theta }(|x|\mathrm{\Delta })\mathrm{\Theta }(ab|x|)\pi \rho ^2],`$ (4.11)
valid for $`t<ll_{max}`$ where
$`l_{max}=\sqrt{t^2+4(ab)^2}\mathrm{\Theta }(q)`$
$`+{\displaystyle \frac{at}{\sqrt{b^2+t^2/4}}}\mathrm{\Theta }(q).`$ (4.12)
A graph of $`𝒫_g^{}(l)`$ is given in figure 20.
Figure 20 The probability density $`𝒫_g^{}(l)`$ for a solid ball under a glide reflection, eq.(4). Here $`a=2,b=1`$ , and $`t=0.25`$ . The regions where the sections $`_g𝒳_x`$ are bicircular or circular are displayed. The function diverges when $`lt`$ , nevertheless the integrated area from $`t`$ to $`l_{max}`$ is finite. $`\mathrm{}`$
Reflections
(Pure) reflections are glide reflections whose translation is $`t=0`$ .
Figure 21 The solid lens $`_g`$ when the isometry $`g`$ is a (pure) reflection on the plane $`x=0`$ . The sections $`_g𝒳_x`$ are circles with variable radius $`\rho `$ . $`\mathrm{}`$
As is evident from figure 21, the probability density $`𝒫_g^{}(l)`$ is proportional to the area of a disc with radius
$`\rho =\sqrt{a^2(b+l/2)^2},`$ (4.13)
and is given by
$`𝒫_g^{}(l)={\displaystyle \frac{3}{2}}{\displaystyle \frac{a^2(b+l/2)^2}{(ab)^2(2a+b)}},`$ (4.14)
for $`0l2(ab)`$. A graph is presented in figure 22.
Figure 22 The probability density $`𝒫_g^{}(l)`$ for a solid ball when the isometry $`g`$ is a (pure) reflection, eq.(4.14). Here $`a=2`$ and $`b=1`$ . The integrated area is 1. $`\mathrm{}`$
## 5 Discussions
We initially aimed to write out one single expression for the probability density $`𝒫_g^{}(l)`$ for screw motions of solid balls, valid for whatever values of the four parameters $`a,b,t`$ , and $`\omega `$. However, we soon found that such expression would demand a quite large number of step functions to account for all sort of possibilities. Since in practice the isometries are dealt with one at each time, we found more appropriate to present a simple method to have the exact $`𝒫_g^{}(l)`$ for each individual screw motion with fixed values of the four parameters. Nevertheless, for those isometric motions of solid balls described by three or less free parameters the exact expression for $`𝒫_g^{}(l)`$ for any euclidean isometry is short enough and was displayed.
As promised in the Introduction, we exhibited the analytic counterpart $`𝒫_g^{}(l)`$ of the computer simulations of pair separations histograms of the euclidean isometries in cosmic crystallography thus far obtained in the literature. The graph of $`𝒫_g^{}(l)`$ in figure 10 corresponds to the isometries $`b`$ and $`c`$ in the Fagundes and Gausmann study, or equivalently the isometries $`\beta `$ and $`\delta `$ in Gomero . The discontinuity in $`l0.7`$ is not observed in the two upper figures 1 (Universe E4) of due to the strong statistical noise present in these histograms; nevertheless it is clearly seen in the position $`s=l^2=0.5`$ in the mean histogram 5b of as well as in the position $`l0.7`$ in the mean histograms 1a, 1b, 2b, 4a, and 5 of . Similarly, figure 11 corresponds to the isometry $`a`$ of and $`\alpha `$ in and ; the discontinuity in $`l=1`$ has their counterparts again in the histogram 5b of and in the histograms 1a, 1b, 2a, 4a, and 5 of .
In contrast with the screw motions (figures 10 and 11), the pure rotations (figures 13, 15, and 16, all with $`t=0`$) do not show discontinuity of $`𝒫_g^{}(l)`$ . Oppositely to figure 13, where $`a<b`$ , figures 15 and 16 correspond to $`a>b`$ , so the ball $``$ now has fixed points and the graph of $`𝒫_g^{}(l)`$ effectively starts from $`l=0`$. The strange-looking graph in figure 15 was confirmed in a computer simulation; the irregular behavior near $`l=0.7`$ corresponds to the narrow $`r`$-interval where the pair-of-discs combined intersection eq.(3) occurs.
In figure 16 we have chosen values for $`a,b`$, and $`\omega `$ such that the three types of combined intersection (ring, pair of discs, and one disc) have equal range in the $`l`$ scale. En passant, the pair-of-discs $`l`$range is now wide, and does not originate a bump as did in figure 15.
Figure 17 corresponds to rotation of the solid ball $``$ around a diameter $`𝒟`$; from (3.51) we find that defining $`l_{max}=2a\mathrm{sin}\omega /2`$ then the graph of $`l_{max}𝒫_g^{}(l)`$ against $`l/l_{max}`$ does not depend on $`l_{max}`$. The points of $``$ along the diameter $`𝒟`$ are fixed under the isometry, so the graph again effectively starts from the origin.
Figure 20 corresponds to a glide reflection whose sections $`x=const`$ in the intersection $`_g`$ are either bicircular discs (for small $`|x|`$) or circular (for larger $`|x|`$). The minimum displacement $`l`$ occurs for the points of $``$ in the intersection with the reflector plane $`x=0`$, giving $`l_{min}=t`$, the translation. Since for all points near the reflector plane we have $`lt(1+2|x|^2)`$, then these points are displaced almost the same value $`lt`$; as a consequence, $`𝒫_g^{}(l)`$ diverges in the vicinity of $`l=t`$. Nevertheless the integrated area is finite, with value 1.
Between $`l=t=0.25`$ and $`l0.5`$ we have bicircular sections (4.6), while for $`0.5<l<l_{max}2.0`$ the sections $`_g𝒳_x`$ are circles.
The transition from a glide reflection to a (pure) reflection is worth describing: if in figure 20 we continuously displace the vertical line $`l=t=0.25`$ towards $`l=0`$, then the region of divergence of $`𝒫_g^{}(l)`$ shrinks continuously and disappears when $`l=0`$, eventually giving the graph of figure 22.
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# Declarative Representation of Revision Strategies
## Introduction
Formal models of belief revision differ in what they consider as representations of the epistemic states of an agent. In the AGM approach (?) epistemic states are identified with logical theories, that is, sets of formulas closed under classical inference. Other approaches like those discussed in (?) consider finite sets of formulas, sometimes called belief bases, as epistemic states. They investigate how to revise such belief bases. Only a rather small fraction of work in belief revision has studied an obvious alternative: the revision of epistemic states expressed as nonmonotonic theories (????).
This is somewhat surprising since close relationships between properties of nonmonotonic inference relations and postulates for belief revision have been established (??). Indeed, one of the reasons why nonmonotonic logics were invented is their ability to handle conflicts and inconsistencies, one of the major issues in belief revision. If this is the case, shouldn’t it be possible to use the power of nonmonotonic inference to simplify revision? In fact, what we have in mind is a complete trivialization of the revision problem. We want to be able to revise nonmonotonic theories simply by adding new information, and we want to leave everything else to the nonmonotonic inference relation.
An early approach in this spirit was the author’s paper (?) where an extension of Poole-systems (?), the so-called preferred subtheory approach, was used. New information, possibly equipped with information about the reliability level of this information, was simply added to the available information. The nonmonotonic inference relation determined the acceptable beliefs.
We are not satisfied with this approach any longer for several reasons. Existing theories of belief revision, including the one presented in the earlier paper, have difficulties to model the way real agents revise their beliefs. One of the reasons for this is that they do not represent information which is commonly used by agents for this purpose. For instance, new information always comes together with certain meta-information (formulas don’t fly into the agent’s mind): Where does the information come from? Was it an observation? Did you read it in the newspaper? Did someone tell you, and if so, who? Did the person who gave you the information have a motive to lie? and so on. In most cases we reason with and about this meta-information when revising our beliefs. We strongly believe that realistic models of revision should provide the necessary means to represent this kind of information.
The meta-information is used to determine the entrenchment of pieces of information. The less entrenched the information is, the more willing we are to give it up. Again, entrechment relations are not just there, they result from reasoning processes. To model this kind of reasoning, entrenchment should be expressible in the logical language. Once we have the possibility to express entrenchment (or plausibility, or preference) in the language, it will also become possible to represent revision strategies declaratively. This in turn makes it possible to revise the revision strategies themselves.
Here is a real life example that can be used to illustrate what we have in mind. Assume Peter tells you that your girl-friend Anne went out for dinner with another man yesterday. Peter even knows his name: it was John, a highly attractive person known for having numerous affairs. You are concerned and talk to Anne about this. She tells you she was at home yesterday evening waiting for you to call. Peter insists that he saw Anne with that man. You are not sure what to believe. Luckily, you find out that Anne has a twin sister Mary. Mary indeed went out with her new boy-friend John. This explains why Peter got mixed up. You now believe Anne and happily continue your relationship.
What this example nicely illustrates is the way we reason about the reliability of information. There is no given fixed entrenchment ordering to start with. In the example there is also, at least in the beginning, no reason to trust Peter more than the girl-friend, or vice versa. And obviously, it is not the new information that is accepted in each situation. It is the additional context information which is relevant here: it gives us an explanation for Peter’s mistake and decreases the reliability of Peter’s observation enough to break the tie.
To be able to formalize examples of this kind we propose in this paper an approach to belief revision where
* nonmonotonic belief bases represent epistemic states and nonmonotonic inference is used to completely trivialize revision,
* it is possible to express and reason about meta-information, including the reliability of formulas,
* revision strategies can be represented declaratively, that is, logical formulas express how conflicts among different pieces of information are resolved.
The outline of the paper is as follows. In the next section we introduce the nonmonotonic formalism we use here to represent epistemic states. In the following section we show how to use this formalism for representing revision strategies. We then discuss the AGM postulates for revision and show that almost all of them are not valid in our approach (which does not bother us). In the following section we briefly deal with contraction. We then discuss forgetting in the context of our approach. Finally, we discuss related work and conclude.
## Representing reliability relations
In this section we introduce the formalism used in this paper. As mentioned in the introduction, one of the distinguishing features of our approach is that we want to be able to reason about the reliability of the available information in the logical language. In the AGM approach (??) entrenchment relations are used to represent how strongly an agent sticks to his beliefs: the more entrenched a formula, the less willing to give it up the agent is. Entrenchment relations have several properties which are based on the logical strength of the formulas. For instance, logically weaker formulas are not less entrenched than logically stronger ones. The intuition is that if a weaker formula has to be given up, the stronger formula has to be given up anyway.
In our approach we do not require such properties. We may even have equivalent formulas $`p`$ and $`p^{}`$ with different reliability. This may happen when, for instance, $`p`$ and $`p^{}`$ come from different sources $`s`$ and $`s^{}`$ with different reliability. Note that although the less reliable information does not add to the accepted beliefs as long as the more reliable equivalent information is in force, the situation may change when new information about the reliability of $`s`$ is obtained. Should $`s`$ turn out to be highly unreliable later (of course, beliefs about the reliability of sources may be revised as any other beliefs) then it becomes important to have $`p^{}`$ with, say, somewhat lower reliability available.
All we require, therefore, is the existence of a strict partial order $``$ between formulas which tells us how to resolve potential conflicts. To avoid misunderstandings we will not call $``$ an entrenchment relation. Instead, we speak of reliability, or simply priority among formulas.
Since we want to represent $``$ in the logical language we need to be able to refer to formulas. Instead of using a quoting mechanism for this purpose, we will use named formulas, that is pairs consisting of a formula and a name for the formula. Technically, names are just ground terms that can be used everywhere in the language.
We will present our formalism in two steps: we first introduce an extension of Poole systems which allows us to express preference information in the language, together with an appropriate definition of extensions. It turns out that due to the potential self-referentiality of preference information not all theories expressed in this formalism possess extensions, that is, acceptable sets of beliefs. In a second step, we therefore introduce a new notion of prioritized inference defined as the least fixed point of a monotone operator. Epistemic states, then, are identified with preferential default theories under this least fixed point semantics.
Our basic formalism extends the well-known Poole systems (?). Recall that Poole systems consist of a consistent set of (first order) formulas $`F`$, the facts, and a possibly inconsistent set of formulas $`D`$, the defaults. A set of formulas $`E`$ is an extension of a Poole-system $`(F,D)`$ iff $`E=Th(FD^{})`$ where $`D^{}`$ is a maximal $`F`$-consistent subset of $`D`$.
Our formalism differs from this approach in the following respects:
1. In the context of belief revision it seems inappropriate to consider some information as absolutely certain and unrevisable. We therefore do not use $`F`$. Instead, we have a single set $`T`$ containing all the information.<sup>1</sup><sup>1</sup>1One of the reviewers of this paper points out that using $`F`$ may have representational advantages since it eliminates the need to use preferences to indicate the most reliable information, see our examples in the rest of the paper. We might therefore reintroduce $`F`$ in future versions of this paper for purely practical reasons. This does not seem to pose any technical problems.
2. We represent preference and other meta-information in the language. We therefore introduce names for formulas and a special symbol $`<`$. $`d<d^{}`$ intuitively says that in case of a conflict $`d^{}`$ should be given up rather than $`d`$ since the latter is more reliable. We require that $`<`$ represents a strict partial order.<sup>2</sup><sup>2</sup>2We assume that the properties of $`<`$, like those of equality, are part of the underlying logic and need not be represented through explicit axioms in our default theories.
3. We introduce a new notion of extension which takes the preference information into account adequately.
To avoid confusion we want to emphasize that $`<`$ belongs to the logical language, whereas $``$ is a meta level symbol. For the following definitions it is essential to clearly separate between these levels.
For simplicity, we only consider finite default theories in this paper. A generalization to the infinite case would have to reduce partial orderings to well-orderings rather than total orders.
###### Definition 1
A named formula is a structure of the form $`d:p`$, where $`p`$ is a first order formula and $`d`$ a ground term representing the name of the formula.
We use the functions $`name`$ and $`form`$ to extract the name respectively formula of a named formula, that is $`name(d:p)=d`$ and $`form(d:p)=p`$. We will also apply both functions to sets of named formulas with the obvious meaning.
###### Definition 2
A preference default theory $`T`$ is a finite set of named formulas such that
* $`form(T)`$ is a set of first order formulas whose logical language contains a reserved symbol $`<`$ representing a strict total order, and
* $`d_1:pT`$, $`d_2:qT`$ and $`pq`$ implies $`d_1d_2`$.
The last item in the definition guarantees that different formulas have different names.
###### Definition 3
Let $`T`$ be a preference default theory, $``$ a total order on $`T`$. The extension of $`T`$ generated by $``$, denoted $`E_T^{}`$, is the set $`E_T^{}=Th(_{i=0}^{|T|}E_i)`$ where
* $`E_0=\mathrm{}`$, and for $`0<i|T|`$
* $`E_i=E_{i1}\{form(d_i)\}`$ if this set is consistent, $`E_{i1}`$ otherwise.
Here $`d_i`$ is the $`i`$-th element of $`T`$ according to the total order $``$.
The set $`_{i=0}^{|T|}E_i`$ is called the extension base of $`E_T^{}`$.
We say $`E`$ is an extension of $`T`$ if there is some total order $``$ such that $`E=E_T^{}`$. Obviously, all maximal consistent subsets of $`form(T)`$ are extension bases. We now consider the general case of partial orders.
###### Definition 4
Let $`T`$ be a preference default theory, $``$ a strict partial order on $`T`$. The set of extensions of $`T`$ generated by $``$ is
$$Ext_T^{}=\{E_T^{^{}}\text{ }|\text{ }^{}\text{ is a total order extending }\}.$$
We next define two notions of compatibility:
###### Definition 5
Let $`T`$ be a preference default theory, $``$ a strict partial ordering of $`T`$, $`S`$ a set of formulas. We say $``$ is compatible with $`S`$ iff
$$S\{d<d^{}\text{ }|\text{ }d:pd^{}:q\}\{\neg (d<d^{})\text{ }|\text{ }d:pd^{}:q\}$$
is consistent.
An extension $`E`$ of $`T`$ is compatible with $`S`$ iff there is a strict partial ordering $``$ of $`T`$ compatible with $`S`$ such that $`EExt_T^{}`$.
The set of extensions of $`T`$ compatible with $`S`$ is denoted $`Ext_T^S`$.
###### Definition 6
Let $`T`$ be a preference default theory. A set of formulas $`E`$ is called a preferred extension of $`T`$ iff $`EExt_T^E`$.
Intuitively, $`E`$ is a preferred extension if it is the deductive closure of a maximal consistent subset of $`T`$ which can be generated through a total preference ordering compatible with the preference information in $`E`$ itself. The preference information in $`E`$ certainly does not have to be total.
Here is a simple example illustrating preference default theories:
> $`d_1(x):bird(x)\mathrm{𝑓𝑙𝑖𝑒𝑠}(x)\text{ }|\text{ }x\text{ ground object term}`$
> $`d_2:x.penguin(x)\neg \mathrm{𝑓𝑙𝑖𝑒𝑠}(x)`$
> $`d_3:bird(tweety)penguin(tweety)`$
> $`d_4:x.d_3<d_1(x)`$
> $`d_5:x.d_2<d_1(x)`$
As is common in Poole systems, rules with exceptions, that is, formulas whose instances can be defeated without defeating the formula as a whole (here $`d_1`$), are represented as schemata used as abbreviations for all of their ground instances. As above we will make the intended instances explicit in all examples. To make sure that the different ground instances can be distinguished by name we have to parameterize the names also. We assume that terms used as names can be distinguished from other terms which we call object terms.<sup>3</sup><sup>3</sup>3A more elaborate formalization would be based on sorted logic with sorts for names and other types of objects from the beginning. We do not pursue this here since we want to keep things as simple as possible. In our case, $`d_1(tweety)`$ is a proper rule name, $`d_1(d_1)`$ is not. Since we only consider finite theories we must also assume that the set of object terms is finite.
In our example we obtain 3 extensions $`E_1`$, $`E_2`$ and $`E_3`$. In $`E_1`$ the instance of $`d_1(x)`$ with $`x=tweety`$ is rejected, in $`E_2`$ $`d_2`$ is rejected, and $`E_3`$ rejects $`d_3`$. All extensions contain $`d_4`$ and $`d_5`$. It is not difficult to see that only $`E_1`$ can be constructed using a total ordering of $`T`$ which is compatible with this information. $`E_1`$ is thus the single preferred extension of this preference default theory.
Preference default theories under extension semantics are very flexible and expressive. The reason we are not yet fully satisfied with them is that they can express unsatisfiable preference information: there are theories which do not possess any preferred extensions. The simplest example is as follows:
> $`d_1:d_2<d_1`$
> $`d_2:d_1<d_2`$
Accepting the first of the two contradictory formulas requires to give preference to the second, and vice versa. No preferred extension exists for this theory.
This means that preference default theories together with the standard notion of nonmonotonic inference where a formula is considered derivable whenever it is contained in all (preferred) extensions do not seem fully adequate for representing epistemic states of rational agents.
We will therefore introduce another, somewhat less standard notion of nonmonotonic consequence.<sup>4</sup><sup>4</sup>4An alternative way of handling this problem would be to introduce some kind of “stratification” into our theories. Stratification, a term taken form the area of logic programming, would prohibit formulas from speaking, directly or indirectly, about their own priority. Unfortunately, it turns out that only highly restrictive forms of stratification guarantee existence of extensions. For this reason we do not pursue this approach here. This approach shares some intuition with the fixed point formulation of well-founded semantics for logic programs with negation due to Baral and Subrahmanian (?). In particular, it is based on the least fixed point of a monotone operator.
Let us first explain the underlying idea. Starting with the empty set, we iteratively compute the intersection of those extensions which are compatible with the information obtained so far. Since the set of formulas computed in each step may contain new preference information the number of extensions may be reduced, and their intersection thus may grow. We continue like this until no further change happens, that is, until a fixed point is reached.
###### Definition 7
Let $`T`$ be a preference default theory, $`S`$ a set of formulas. We define an operator $`C_T`$ as follows:
$$C_T(S)=Ext_T^S$$
###### Proposition 1
The operator $`C_T`$ is monotone.
Proof: $`SS^{}`$ implies that an ordering $``$ is compatible with $`S`$ whenever it is compatible with $`S^{}`$. We thus have $`Ext_T^S^{}Ext_T^S`$ and therefore $`Ext_T^SExt_T^S^{}`$. $`\mathrm{}`$
Monotone operators, according to the well-known Knaster-Tarski theorem (?), possess a least fixed point. This fixed point can be computed by iterating the operator on the empty set. We, therefore, can define the accepted conclusions of a preference default theory as follows:
###### Definition 8
Let $`T`$ be a preference default theory. A formula $`p`$ is an accepted conclusion of $`T`$ iff $`p\mathrm{𝑙𝑓𝑝}(C_T)`$, where $`\mathrm{𝑙𝑓𝑝}(C_T)`$ is the least fixed point of the operator $`C_T`$.
We call extensions which are compatible with $`\mathrm{𝑙𝑓𝑝}(C_T)`$ accepted extensions.
Several illustrative examples will be given in the next section. Here we just show how the theory without preferred extension is handled in this approach. We have $`T=\{d_1:(d_2<d_1),d_2:(d_1<d_2)\}`$. We first compute $`C_T(\mathrm{})`$. Since no preference information is available in the empty set we obtain $`Th(\{d_2<d_1\})Th(\{d_1<d_2\})`$ which is equivalent to $`Th(\{d_2<d_1d_1<d_2\})`$. This set is already the least fixed point.
###### Proposition 2
Let $`T`$ be a preference default theory, $`p`$ an accepted conclusion of $`T`$. Then $`p`$ is contained in all preferred extensions of $`T`$.
Proof: If $`T`$ has no preferred extension the proposition is trivially true. So assume $`T`$ possesses preferred extension(s). A simple induction shows that each preferred extension is among the extensions compatible with the formulas computed in each step of the iteration of $`C_T`$. Therefore each preferred extension is also an accepted extension. $`\mathrm{}`$
###### Proposition 3
Let $`T`$ be a preference default theory. The set of accepted conclusions of $`T`$ is consistent.
Proof: We show by induction that, for arbitrary $`n`$, the set of formulas obtained after $`n`$ applications of $`C_T`$ is consistent. For $`n=0`$ this is trivial. Assume the set of formulas $`S`$ obtained after $`n1`$ iterations is consistent. Since $`S`$ is consistent and $`<`$ formalizes a strict partial ordering there must be at least one strict partial ordering $``$ compatible with $`S`$, so the set of all extensions compatible with $`S`$ is nonempty. Since each extension is by definition consistent the intersection of an arbitrary nonempty set of extensions must also be consistent. $`\mathrm{}`$
Since preference default theories under accepted conclusion semantics always lead to consistent beliefs, we will in the next section identify epistemic states with preference default theories and belief sets with their accepted conclusions.
## Revising epistemic states
### The revision operator
Given an agent’s epistemic state is identified with a preference default theory as introduced in the last section, it is natural to identify the set of beliefs accepted by the agent with the accepted conclusions of this theory. We therefore define belief sets as follows:
###### Definition 9
Let $`T`$ be an epistemic state. $`Bel(T)`$, the belief set induced by $`T`$, is the set of accepted conclusions of $`T`$.
It is a basic assumption of our approach that belief sets cannot be revised directly. Revision of belief sets is always indirect, through the revision of the epistemic state inducing the belief set. Note that since two different epistemic states may induce the same belief set, the revision function which takes an epistemic state and a formula and produces a new epistemic state does not induce a corresponding function on belief sets.
Given an epistemic state $`T`$, revising it with new information simply means generating a new name for it and adding the corresponding named formula.
###### Definition 10
Let $`T`$ be an epistemic state, $`p`$ a formula. The revision of $`T`$ with $`p`$, denoted $`Tp`$, is the epistemic state $`(T\{n:p\})`$ where $`n`$ is a new name not appearing in $`T`$.
Notation: in the rest of the paper we assume that names are of the form $`d_j`$ where $`j`$ is a numbering of the formulas. If $`T`$ has $`j`$ elements and a new formula is added, then its new name is $`d_{j+1}`$.
### Representing revision strategies
In this subsection we show how revision strategies used by an agent can be represented in our approach. We first discuss an example where the strategy is based on the type of the available information. We distinguish between strict rules, observations and defaults. Strict rules have highest priority because they represent well-established or terminological information. Observations can be wrong, but they are considered more reliable than default information. Consider the following epistemic state $`T`$:
> $`d_1:penguin(tweety)`$
> $`d_2:x.penguin(x)bird(x)`$
> $`d_3:x.penguin(x)\neg \mathrm{𝑓𝑙𝑖𝑒𝑠}(x)`$
> $`d_4(x):bird(x)\mathrm{𝑓𝑙𝑖𝑒𝑠}(x)\text{ }|\text{ }x\text{ ground object term}`$
> $`d_5:observation(d_1)`$
> $`d_6:rule(d_2)`$
> $`d_7:rule(d_3)`$
> $`d_8:x.\mathrm{𝑑𝑒𝑓𝑎𝑢𝑙𝑡}(d_4(x))`$
> $`d_9:n,n^{}.rule(n)observation(n^{})n<n^{}`$
> $`d_{10}:n,n^{}.observation(n)default(n^{})n<n^{}`$
$`T`$ has 4 extensions. The corresponding extension bases are obtained from $`T`$ by leaving out $`d_1`$, $`d_2`$, $`d_3`$, or $`d_4(tweety)`$, respectively. All extensions, and thus $`C_T(\mathrm{})`$, contain information stating that $`d_4(tweety)`$ has lower preference than the other three formulas. Therefore, the only extension compatible with $`C_T(\mathrm{})`$ is the one generated by leaving out $`d_4(tweety)`$. This set is also the least fixpoint of $`C_T`$. $`Bel(T)`$ thus does not contain $`\mathrm{𝑓𝑙𝑖𝑒𝑠}(tweety)`$.
The next example formalizes the revision strategy of an agent who prefers newer information over older information and information from a more reliable source over information from a less reliable source. In case of a conflict between the two criteria the latter one wins.
Assume the following specific scenario: At time 10 Peter informs you that $`p`$ holds. At time 11 John tells you this is not true. Although you normally prefer later information, you also have reason to prefer what Peter told you since you believe Peter is more reliable than John. Since you consider reliability of your sources even more important than the temporal order you believe $`p`$.
Here is the formal representation of this scenario. We use $`X<d`$ where $`X`$ is a finite set of names as an abbreviation for $`_{xX}x<d`$. Note that we have to make sure, by adding adequate preferences, that the rules representing our revision strategy cannot be used - via contraposition - to defeat our meta-knowldge about $`d_1`$ and $`d_2`$:
| $`d_1:p`$ |
| --- |
| $`d_2:\neg p`$ |
| $`d_3:time(d_1)=10`$ |
| $`d_4:time(d_2)=11`$ |
| $`d_5:source(d_1)=Peter`$ |
| $`d_6:source(d_2)=John`$ |
| $`d_7:more\text{-}rel(John,Peter)`$ |
| $`d_8(n,n^{}):more\text{-}rel(source(n),source(n^{}))n<n^{}`$ |
| | $`|`$ $`n,n^{}\{d_1,\mathrm{},d_7\}`$ |
| $`d_9(n,n^{}):time(n)<time(n^{})n^{}<n`$ |
| | $`|`$ $`n,n^{}\{d_1,\mathrm{},d_7\}`$ |
| $`d_{10}:n,n^{}.\{d_3,\mathrm{},d_7\}<d_8(n,n^{})<d_9(n,n^{})`$ |
This preference default theory has 14 extensions which are obtained by leaving out one of $`\{d_1,d_2\}`$ and one of $`\{d_3,d_4,d_5,d_6,d_7,d_8(d_1,d_2),d_9(d_1,d_2)\}`$. All extensions contain $`d_{10}`$. This means that after the next iteration of the $`C_T`$-operator we are left with 2 extensions which are obtained by leaving out one of $`\{d_1,d_2\}`$ and $`d_9(d_1,d_2)`$. Both extensions contain the formula $`d_1<d_2`$. The next and final iteration of $`C_T`$ thus eliminates the extension containing $`d_2`$. We are left with a single extension and $`p`$ is among the accepted conclusions.
We next present the example from the introduction. This time we use categories $`low`$, $`medium`$ and $`high`$<sup>5</sup><sup>5</sup>5We assume uniqueness of names for the categories. Otherwise the set $`\{d_3,d_5,d_6,d_7\}`$ would be consistent and could be used to defeat $`d_9`$ which, obviously, is unintended. to express reliability: the reliability of a formula with name $`n`$ is $`rel(n)`$. We have the following information:
| $`d_1:date(Anne,John)`$ |
| --- |
| $`d_2:\neg date(Anne,John)`$ |
| $`d_3:rel(d_1)=medium`$ |
| $`d_4:rel(d_2)=medium`$ |
| $`d_5:date(Mary,John)`$ |
| $`d_6:twins(Mary,Anne)`$ |
| $`d_7:date(Mary,John)twins(Mary,Anne)`$ |
| | $`rel(d_1)=low`$ |
| $`d_8:n,n^{}.rel(n)=highrel(n^{})=medium`$ |
| | $`n<n^{}`$ |
| $`d_9:n,n^{}.rel(n)=mediumrel(n^{})=low`$ |
| | $`n<n^{}`$ |
| $`d_{10}:rel(d_5)=rel(d_6)=rel(d_7)=rel(d_8)=`$ |
| | $`rel(d_9)=high`$ |
| $`d_{11}:rel(d_3)=rel(d_4)=medium`$ |
Although the agent initially considers $`d_1`$ and $`d_2`$ as equally reliable, the information that Anne has a twin sister Mary who is dating John decreases the reliability of $`d_1`$ to $`low`$. $`d_8`$ and $`d_9`$ say how the reliability categories are to be translated to preferences. $`d_{10}`$ and $`d_{11}`$ make sure that meta-information is preferred, and that $`d_7`$ can defeat $`d_3`$.Taking all reliability information into account the agent accepts $`\neg date(Anne,John)`$.
Note that we do not discuss here how agents form meta-beliefs of the kind required to represent the example (induction, folk psychology?). We simply assume that this information is available to the agent.
## Postulates
We now discuss the postulates for revision which are at the heart of the AGM approach (?). Since our approach uses epistemic states rather than deductively closed sets of formulas (belief sets) as substrate of revision, some of the postulates need reformulation. In particular, AGM use the expansion operator $`+`$ in some postulates. Expansion of a belief set $`K`$ with a formula $`p`$ means adding $`p`$ to the belief set and closing under deduction, that is $`K+p=Th(K\{p\})`$. Since epistemic states always induce consistent belief sets the distinction between revising and expanding an epistemic state does not seem to make much sense in our context. We therefore translate expansion in the following postulates to expansion of the induced belief set.
In the following we present the AGM postulates (K\*i) in their original form together with our corresponding reformulations (T\*i). In each case K is a belief set in the sense of AGM, $`T`$ an epistemic state as defined in this paper, $`A,B`$ are formulas:
> (K\*1) $`KA`$ is a belief set.
> (T\*1) $`Bel(TA)`$ is belief set.
Obviously satisfied.
> (K\*2) $`AKA`$
> (T\*2) $`ABel(TA)`$
Not satisfied. New information is not necessarily accepted in our approach. We see this as an advantage since otherwise belief sets would always depend on the order in which information was obtained.
> (K\*3) $`KAK+A`$
> (T\*3) $`Bel(TA)Bel(T)+A`$
Not satisfied. Assume we have $`T=\{d_1:p,d_2:\neg p\}`$, that is $`Bel(T)`$ is the set of tautologies. Let $`A=d_1<d_2`$. Now $`Bel(TA)`$ contains $`p`$ which is not contained in $`Bel(T)+A`$.
| (K\*4) if $`\neg AK`$ then $`K+AKA`$ |
| --- |
| (T\*4) if $`\neg ABel(T)`$ then |
| | $`Bel(T)+ABel(TA)`$ |
Not satisfied. It may be the case that $`\neg A`$, although not in the belief set, is contained in one of the accepted extensions. Adding $`A`$ to the epistemic state does not necessarily lead to a situation where this extension disappears.
> (K\*5) $`KA`$ iff $`\neg A`$
> (T\*5) $`Bel(TA)`$ iff $`\neg A`$
Not satisfied. Revising an epistemic state with logically inconsistent information has no effect whatsoever. The information is simply disregarded. Inconsistent belief sets are impossible in our approach, so the right to left implication does not hold.
> (K\*6) If $`AB`$ then $`KA=KB`$
> (T\*6) If $`AB`$ then $`Bel(TA)=Bel(TB)`$
Satisfied under the condition that $`A`$ and $`B`$ are given the same name, or the names of $`A`$ and $`B`$ do not yet appear in $`S`$. But note that logically equivalent information may have different impact on the belief sets when different meta-information is available. For instance, $`d_1:p`$ and $`d_2:p`$ may have different effects if different meta-information about the sources of $`d_1`$ and $`d_2`$, respectively, is available.
> (K\*7) $`K(AB)(KA)+B`$
> (T\*7) $`Bel(T(AB))Bel(TA)+B`$
Not satisfied. Here is a counterexample. Assume we have $`T=\{d_1:p,d_2:\neg p,d_3:\neg p\}`$. Now let $`A=d_1<d_2`$ and $`B=d_1<d_3`$. Clearly, revising the epistemic state with $`AB`$ leads to a single accepted extension containing $`p`$ since the two conflicting formulas are less preferred. $`p`$ is thus in the belief set induced by the revised state. On the other hand, revising the epistemic state with $`A`$ leads to two extensions, one containing $`p`$, the other $`\neg p`$. $`p`$ is thus not in the belief set induced by the new state. This does not change when we expand the belief set with $`d_1<d_3`$.
| (K\*8) If $`\neg BKA`$ then $`(KA)+BK(AB)`$ |
| --- |
| (T\*8) If $`\neg BBel(TA)`$ then |
| | $`Bel(TA)+BBel(T(AB))`$ |
Not satisfied. This is immediate from the fact that $`Bel(T(AB))`$ does not necessarily contain $`B`$, that is from the failure of (T\*2).
This analysis shows that the intuitions captured by the AGM postulates are indeed very different from those underlying our approach.
## Contraction
Contraction means making a formula underivable without assuming its negation. There may be different reasons for this, not all of them requiring extensions of our framework. For instance, the reliability of a source of a certain piece of information may be in doubt due to extra information. In that case it may happen that a belief $`pBel(T)`$ is no longer in the belief set $`Bel(Tq)`$ for appropriate $`q`$ even if $`\neg pBel(Tq)`$. Such effects are handled implicitly in our approach.
If, however, the agent may obtain information of the kind “do not believe $`p`$” rather than “believe $`\neg p`$”, then extra mechanisms seem necessary. In the context of AGM-style approaches the contraction operator $``$ can be defined through revision on the basis of the so-called Harper identity: $`KA=(K\neg A)K`$. The intuition here is that revision with $`\neg A`$ removes the formulas used to derive $`A`$, and the intersection with $`K`$ guarantees that no new information is derived from $`\neg A`$.
This intuition can, to a certain extent, be captured using Poole’s constraints (?). Constraints, basically, are formulas used in the construction of maximal consistent subsets of the premises, but not used for derivations.
To model contraction of epistemic states we must distinguish between these two types of formulas, premises and constraints. Extension bases consist of both types and also the compatibility of preference orderings is checked against premises and contraints. Extensions, however, are generated only from the premises. Constraints, as regular formulas, have names and may come with meta-information, e.g., information about their reliability.
We do not want to go into further technical detail here. Instead, we illustrate contraction using an example. We indicate constraints by choosing names of the form $`c_j`$ for them. Assume the epistemic state is as follows:
> $`d_1:peng(tweety)`$
> $`d_2(x):peng(x)\neg \mathrm{𝑓𝑙𝑖𝑒𝑠}(x)`$
The agent receives the information “do not believe $`\neg flies(tweety)`$”. The following constraint is added:
> $`c_1:\mathrm{𝑓𝑙𝑖𝑒𝑠}(tweety)`$
Note that the constraint is not necessarily preferred to the premises. Let $`inst(d_2)`$ denote the set of all ground instances of $`d_2`$. We obtain three extension bases
> $`E_1=\{peng(tweety)\}inst(d_2)`$
> $`E_2=\{peng(tweety),\mathrm{𝑓𝑙𝑖𝑒𝑠}(tweety)\}inst(d_2)\{peng(tweety)\neg \mathrm{𝑓𝑙𝑖𝑒𝑠}(tweety)\}`$
> $`E_3=\{\mathrm{𝑓𝑙𝑖𝑒𝑠}(tweety)\}inst(d_2)`$
Although $`E_2`$ and $`E_3`$ contain $`\mathrm{𝑓𝑙𝑖𝑒𝑠}(tweety)`$ this formula is not in the extensions generated from these extension bases, and for this reason not in the belief set, since it is a constraint.
Note that constraints do not necessarily prohibit formulas from being in the belief set since they may have low reliability. For example, if we revise the epistemic state obtained above with $`d_1<c_1`$ and $`x.d_2(x)<c_1`$ then the belief set contains $`\neg \mathrm{𝑓𝑙𝑖𝑒𝑠}(tweety)`$.
Although we used the Harper identity above to motivate the use of constraints for contraction, its natural reformulation
$$Bel(TA)=Bel(T\neg A)Bel(T)$$
is not valid in our approach. Assume $`T=\{d_1:p,d_2:\neg p\}`$. Obviously, $`Bel(T)`$ is the set of tautologies. Now let $`A=\neg (d_1<d_2)`$. We contract by adding the constraint $`c_1:(d_1<d_2)`$. Now the single accepted extension of the new epistemic state and thus its belief set is $`Th(p)`$, a strict superset of $`Bel(T)`$.
## Forgetting
In our approach revising a knowledge base means adding a formula to the epistemic state. Even in the case of contraction the epistemic state grows. For ideal agents this may be adequate since every piece of information, whether it contributes to the current belief set or not, may turn out to be relevant later. However, for agents with limited resources the expansion of the epistemic state cannot go on forever.
This raises the question how and when pieces of information should be forgotten. What we need is some kind of a mental garbage collection strategy. In LISP systems garbage collection is the process of identifying unaccessable memory space which is then made available again. In our context there is no clear distinction between garbage and non-garbage. As mentioned before, every information may become relevant through additional information, so it would not be reasonable to throw away information just because it is, say, not contained in any extension base. On the other hand, even those formulas contributing to the current belief set may be considered as garbage if the corresponding part of the belief set is not relevant to the agent. It appears that a satisfactory treatment of forgetting would have to take the utility of information for the agent into account. This is beyond the scope of this paper and a topic of further research.
## Related work and discussion
In this paper we proposed a framework for belief revision where preference default theories together with a corresponding nonmonotonic inference relation are used to represent epistemic states and belief sets, respectively. Our underlying formalism draws upon ideas developed in (?) and (?), the notion of accepted conclusions introduced to guarantee consistency of belief sets and its application to belief revision is new. The framework is expressive enough to represent and reason about reliability and other properties of information. It thus can be used to represent revision strategies of agents declaratively. Another advantage of the framework is that it lends itself to iteration in an obvious and natural way.
In an earlier paper (?) the author used nonmonotonic belief bases in the preferred subtheories framework to model revision. This approach, however, did not represent reliability information explicitly. Williams and Antoniou (?) investigated revision of Reiter default theories. In a similar spirit, Antoniou et al (?) discuss revision of theories expressed in Nute’s defeasible logic. Also these approaches do not reason about the reliability of information. This is also true for existing work in revising logic programs, see (?) for an example.
Forms of revision where new information is not necessarily accepted were investigated by Hansson (?). This form of revision is sometimes referred to as non-prioritized belief revision. Hansson called his version ”semi-revision”. Explicit reasoning about the available information is not modelled in Hansson’s approach.
Structured belief bases were investigated by Wassermann (?). Rather than using the structure to model meta-level and preference information, Wassermann uses structure to determine relevant parts of the belief base. The focus is thus on local revision operations and related complexity issues. Chopra and Parikh (?) propose a model where belief bases are partitioned into subbases according to syntactic criteria. Belnap’s four-valued logic is used for query answering. Again the focus is on keeping the effects of revision as local as possible. It is assumed that the local revision operators used satisfy the AGM postulates. The approach is thus very different from ours.
## Acknowledgements
The work presented in this paper was funded by DFG (Deutsche Forschungsgemeinschaft), Forschergruppe Kommunikatives Verstehen. I thank R. Booth, S. Lange, H. Sturm and F. Wolter for helpful comments. Thanks also to the anonymous reviewers of the paper.
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# 1 Introduction and Summary
## 1 Introduction and Summary
The study of brane-antibrane systems in string and M-theory has brought a new perspective onto the understanding of non-perturbative extended objects (see -, and references therein). BPS branes arise as by-products when the tachyonic mode of the open strings connecting the brane and the antibrane condenses in a vortex-like configuration. Also, one can deduce in this framework the presence of new non-BPS extended objects when the tachyon condenses, instead, in a kink-like configuration. These objects have subsequently been shown to play a key role in testing string dualities beyond the BPS level. In Type II and M-theory they are however unstable, because the open strings ending on them contain real tachyonic excitations, but the instability can be cured when the theory is projected out by a certain symmetry that removes the tachyons from the spectrum. This happens in particular for some non-BPS branes in Type I and in certain orbifold/orientifold constructions of Type II and M-theory.
An interesting and still open problem is the construction of the supergravity solutions corresponding to these non-BPS branes. In the Type II theories given that they are unstable one does not expect to find stable classical solutions of the supergravity that could be associated to these branes. Only in some cases it has been possible to construct classical solutions, like the one corresponding to a Kaluza-Klein monopole-antimonopole pair in M-theory, that is stabilized by suspending the system in an external magnetic field (see and for related work), or the solution of the, unstable, non-BPS D-instanton of Type IIA constructed in . For the stable non-BPS branes one however expects to find solutions of the supergravity equations of motion that describe these branes in the strong coupling regime.
Recently, using the boundary state formalism, Eyras and Panda found the asymptotic behavior of the solution corresponding to the non-BPS D0-brane of the Type IIB theory orbifolded by $`(1)^{F_L}I_4`$ <sup>1</sup><sup>1</sup>1Here $`F_L`$ denotes the left-moving spacetime fermion number, and $`I_4:x^ix^i;i=1,\mathrm{}4`$.. This theory is S-dual to Type IIB orientifolded by $`\mathrm{\Omega }I_4`$, $`\mathrm{\Omega }`$ being the worldsheet parity reversal operation. The twisted sector of this orientifold construction of Type IIB consists, in the compact case, on 16 O5 orientifold fixed planes together with a D5-brane on top of each plane, each D5-O5 system carrying an SO(2) vector potential. This theory contains massive non-BPS states in the perturbative spectrum arising from open strings stretched between a D5-brane and its image . These states are stable, since they are the lightest states charged under the SO(2) gauge field of the twisted sector, and correspond in the strong coupling limit to the non-BPS D0-branes of IIB/$`(1)^{F_L}I_4`$ . These non-BPS D0-branes are also charged with respect to the SO(2) vector field of the, S-dual, twisted sector, and this renders them stable. Moreover, for a critical value of the radii of the compact orbifold a pair of branes satisfies a no-force condition at least when the distance is larger than the string scale <sup>2</sup><sup>2</sup>2In it is shown that for coincident branes a vacuum configuration in which the branes attract each other seems to be more favourable even at the critical radius. As mentioned in that reference this result is however challenged by the fact that the expectation value for the tachyon fields is beyond the range of validity of the approximation. It would be very interesting to clarify this point further.. This opens the possibility of constructing solutions corresponding to a large number of parallel non-BPS D0-branes, which could correctly describe the weak coupling regime of the theory. More importantly for our discussion, it allows the possibility of constructing infinite arrays of non-BPS D0-branes, from where one can derive other non-BPS D-brane solutions via T-duality transformations.
In this paper we concentrate on this and other stable non-BPS branes, that occur in the six dimensional orbifold/orientifold constructions obtained by projecting out the Type IIB theory by $`\mathrm{\Omega }I_4`$ and duality-related operations. In particular, we focus on the Type IIB and Type IIA theories divided out by $`(1)^{F_L}I_4`$. The $`(1)^{F_L}`$ operation is identified in the strong coupling limit as the transformation reversing the orientation of open D-strings (D2-branes) in the Type IIB (Type IIA) theory, as implied by its connexion via S-duality with the $`\mathrm{\Omega }`$ symmetry of Type IIB (in Type IIA a further T-duality transformation is required). Therefore, the twisted sector of Type IIB/$`(1)^{F_L}I_4`$ can be described non-perturbatively, in the compact case, as 16 O5-NS5B systems , and that of Type IIA/$`(1)^{F_L}I_4`$ as the same number of O5-NS5A systems . An O5-plane with a NS5B-brane on top of it contains an SO(2) gauge field associated to open D-strings stretched between the NS5-brane and its image, and the non-BPS D0-branes are charged with respect to this twisted field. In the Type IIA theory, the O5-NS5A system contains a self-dual 2-form field associated to open D2-branes stretched between a NS5A-brane and its image, and therefore the twisted sector contributes with SO(2) self-dual 2-form potentials<sup>3</sup><sup>3</sup>3See and for a more detailed description of the twisted sector..
In Sen conjectured that together with the non-BPS D0-brane, coupled electrically to the SO(2) vector field of the NS5B-O5 system, there is a non-BPS D2-brane, placed on the orbifold plane, that couples magnetically to the same vector field, and arises from open D3-branes stretched between a NS5B-brane and its mirror. Also, T-duality predicts a non-BPS D1-brane in IIA/$`(1)^{F_L}I_4`$ located on the orbifold plane, which should couple, electrically and magnetically, to the self-dual 2-form potential of the twisted sector, and arise from open D2-branes stretched between a NS5A-brane and its mirror.
In this note we construct these D-brane solutions. We also show that there is a non-BPS M1-brane solution of M-theory orientifolded by $`\mathrm{\Omega }_\rho I_5`$ , where $`\mathrm{\Omega }_\rho `$ reverses the orientation of the M2-brane, to which all these branes are related by reduction and dualities. This provides a unifying picture within M-theory of the stable non-BPS branes that occur in the six dimensional orbifold/orientifold constructions related to $`\mathrm{\Omega }_\rho I_5`$.
M-theory on a 5-torus orientifolded by $`\mathrm{\Omega }_\rho I_5`$ contains a twisted sector that can be identified as 16 O5-M5 systems, together with other 16 O5 orientifold fixed planes which do not contribute any twisted states . This theory contains non-BPS M1-branes that arise from open M2-branes stretched between an M5-brane and its mirror . They couple, electrically and magnetically, to the self-dual 2-form potential living in the M5-brane, what makes them stable. In reducing to the Type IIA theory one can consider two possibilities:
1. Reduce along a worldvolume direction of the M5-O5 system. In this case one obtains Type IIA orientifolded by $`\mathrm{\Omega }I_5`$, whose twisted sector states arise from 16 D4-O4 systems. This theory contains perturbative massive non-BPS states, which can be interpreted in M-theory as non-BPS M1-branes wrapped on the eleventh direction , as well as non-perturbative non-BPS strings coming from open D2-branes stretched between a D4-brane and its mirror, which correspond to unwrapped M1-branes in M-theory . T-duality along one of the orbifolded directions gives then rise to the Type IIB theory orientifolded by $`\mathrm{\Omega }I_4`$, whose twisted sector is described by 16 D5-O5 systems. This theory contains perturbative non-BPS particle states and non-perturbative non-BPS 2-branes, connected by T-duality with the non-BPS objects of IIA.
2. Reduce M-theory/$`\mathrm{\Omega }_\rho I_5`$ along one of the orbifolded directions. In this case one obtains Type IIA projected out by $`(1)^{F_L}I_4`$, with a twisted sector consisting of 16 NS5-O5 systems. The non-BPS M1-brane gives rise to a non-BPS D1-brane in IIA that couples (electrically and magnetically) to the self-dual 2-form potential living in the worldvolume of the NS5A-O5. Now T-duality along a worldvolume direction of the NS5A-O5 maps the theory onto Type IIB divided by $`(1)^{F_L}I_4`$, with a twisted sector identified as 16 NS5-O5 systems. Non-BPS D0-branes are coupled electrically to the SO(2) vector field of the twisted sector, and non-BPS D2-branes magnetically. These branes are related to the non-BPS D1-brane by T-duality.
Consistently with the whole duality picture , the two theories that are obtained by either reducing along an M5-brane direction and then T-dualizing along a transverse direction, or viceversa, are related by S-duality. In particular one obtains IIB/$`\mathrm{\Omega }I_4`$ and IIB/$`(1)^{F_L}I_4`$ respectively. We also see that the non-BPS M1-brane of M-theory/$`\mathrm{\Omega }_\rho I_5`$ is the eleven dimensional origin of the non-BPS branes that can be defined in the Type II orbifolds/orientifolds obtained by reduction.
## 2 The non-BPS D0-brane solution of
The asymptotic behavior of the solution corresponding to the non-BPS D0-brane of Type IIB orbifolded by $`(1)^{F_L}I_4`$ has been derived in using the boundary state formalism. In this formalism one can compute the long distance behavior of the massless fields generated by the D-brane and predict in this manner the asymptotic form of the corresponding classical solution . A pair of non-BPS D0-branes satisfies a no-force condition when the orbifold is compactified to a particular critical value of the radii . When this happens it is possible to construct periodic infinite arrays of non-BPS D0-branes and compute T-dual solutions, which is what we shall be doing in the next sections. Note that both the no-force condition and the validity of the classical solution and the T-duality rules hold for distances larger than the string scale.
The asymptotic form of the solution of , corresponding to a D0-brane situated at one of the fixed points of the orbifold, reads, in string frame<sup>4</sup><sup>4</sup>4In a somewhat more general solution depending on a free parameter $`a`$ is given, derived by impossing the no-force condition of a pair of branes at the critical radii as a constraint for the background fields. Here we have chosen to work with the strictly linearized solution, though the same kind of generalization can be done for our solutions.:
$`ds_{D0}^2=(1{\displaystyle \frac{1}{3}}{\displaystyle \frac{\kappa _6T_0}{2\pi ^2\mathrm{\Omega }_4}}{\displaystyle \frac{1}{|y|^3}}+\mathrm{})dt^2+`$
$`+(1+{\displaystyle \frac{1}{3}}{\displaystyle \frac{\kappa _6T_0}{2\pi ^2\mathrm{\Omega }_4}}{\displaystyle \frac{1}{|y|^3}}+\mathrm{})(\delta _{mn}dy^mdy^n+`$
$`+\delta _{ij}dx^idx^j);m,n=1,\mathrm{}5;i,j=1,\mathrm{}4,`$
$`e^\varphi =1+{\displaystyle \frac{1}{2}}{\displaystyle \frac{\kappa _6T_0}{2\pi ^2\mathrm{\Omega }_4}}{\displaystyle \frac{1}{|y|^3}}+\mathrm{},`$
$`C_0^{(1)}={\displaystyle \frac{1}{3}}{\displaystyle \frac{\kappa _6Q_0}{\sqrt{2}\mathrm{\Omega }_4}}{\displaystyle \frac{1}{|y|^3}}+\mathrm{}.`$ (2.1)
Here we have taken $`\alpha ^{}=1`$ but otherwise the notation is that in . Namely, $`\kappa _D^2=8\pi G_D`$, $`\kappa _{Dd}^2=\kappa _D^2/V_d`$, with $`V_d`$ the volume of the $`d`$ dimensional space, $`\mathrm{\Omega }_4`$ is the area of a unit sphere surrounding the D0-brane, $`T_0`$ is the tension of the brane, $`Q_0`$ its charge<sup>5</sup><sup>5</sup>5Impossing open-closed string consistency for boundary states, $`T_0`$ and $`Q_0`$ are fixed to: $`T_0=8\pi ^{7/2}`$, $`Q_0=8\sqrt{2}\pi ^{3/2}`$, (see )., $`y^m`$, $`m=1,\mathrm{}5`$ the longitudinal directions along the NS5B-O5 worldvolume, and $`x^i`$, $`i=1,\mathrm{}4`$, the transverse, orbifolded directions. The critical value of the radii: $`R_c=1/\sqrt{2}`$, has already been substituted in the solution. $`C^{(1)}`$ is the vector potential coming from the twisted sector, under which the D0-brane is charged.
## 3 The non-BPS D1-brane of IIA/$`(1)^{F_L}I_4`$
Considering a periodic infinite array of non-BPS D0-branes along the $`y^5`$ direction we can construct via T-duality a non-BPS D1-brane solution in the Type IIA theory projected out by $`(1)^{F_L}I_4`$. This brane is situated at one of the fixed points of the orbifold with its worldsheet extended along the non-compact spacetime. We find:
$`ds_{D1}^2=(1{\displaystyle \frac{1}{2}}{\displaystyle \frac{\kappa _6T_1}{2\pi ^2\mathrm{\Omega }_3}}{\displaystyle \frac{1}{|y|^2}}+\mathrm{})(dt^2+d\sigma ^2)+`$
$`+(1+{\displaystyle \frac{1}{2}}{\displaystyle \frac{\kappa _6T_1}{2\pi ^2\mathrm{\Omega }_3}}{\displaystyle \frac{1}{|y|^2}}+\mathrm{})\left(\delta _{mn}dy^mdy^n+\delta _{ij}dx^idx^j\right);`$
$`m,n=1,\mathrm{}4;i,j=1,\mathrm{}4,`$
$`e^\varphi =1+{\displaystyle \frac{1}{2}}{\displaystyle \frac{\kappa _6T_1}{2\pi ^2\mathrm{\Omega }_3}}{\displaystyle \frac{1}{|y|^2}}+\mathrm{},`$
$`C_{0\sigma }^{(2)}={\displaystyle \frac{1}{2}}{\displaystyle \frac{\kappa _6Q_1}{\sqrt{2}\mathrm{\Omega }_3}}{\displaystyle \frac{1}{|y|^2}}+\mathrm{}.`$ (3.1)
Here $`T_1`$ is the tension of the brane, $`Q_1`$ its charge and $`\mathrm{\Omega }_3`$ the area of the unit 3-sphere surrounding the string. The twisted sector consists on a 2-form potential, under which the D1-brane is charged. T-duality implies that this field is to be interpreted as the 2-form potential living in a NS5A-O5 system.
This D1-brane solution can be interpreted as a ten dimensional 1-brane located at the origin of the four-dimensional compact space. Considering first a single compactified direction, a 1-brane sitting at the origin of the $`S^1`$ can be seen from the point of view of the covering space of the $`S^1`$ as an equally spaced array of D1-branes in the $`S^1`$ direction. If $`\stackrel{}{x}`$ denotes a vector in the full eight dimensional transverse space, we can then approximate<sup>6</sup><sup>6</sup>6The choice of this function will be clear below. $`1/|x|^6`$ by a sum $`_{nZ}1/(r^2+(x^42\pi nR_4)^2)^3`$, with $`r^2=_{m=1}^4y^my_m+_{i=1}^3x^ix_i`$, and, assuming that the size of the compact direction is smaller than the distance in the non-compact space, we can further approximate the sum by an integral. Repeating this process for the four compact directions we can finally write (see for instance for more details):
$$\frac{1}{|y|^2}\frac{1}{|x|^6}\mathrm{\Pi }_{i=1}^4(2\pi R_i)(I_1I_2I_3I_4)^1,$$
(3.2)
where $`R_i`$ are the radii of the compactified orbifold and $`I_n_0^\pi 𝑑\theta \mathrm{sin}^n\theta `$. Substituting back in the expression for the D1-brane solution in the compact orbifold we can then obtain the corresponding solution in the uncompactified case:
$`ds_{D1}^2=(1{\displaystyle \frac{1}{6}}{\displaystyle \frac{\kappa _{10}T_1}{\mathrm{\Omega }_7}}{\displaystyle \frac{1}{|x|^6}}+\mathrm{})(dt^2+d\sigma ^2)+`$
$`+(1+{\displaystyle \frac{1}{6}}{\displaystyle \frac{\kappa _{10}T_1}{\mathrm{\Omega }_7}}{\displaystyle \frac{1}{|x|^6}}+\mathrm{})\left(\delta _{mn}dy^mdy^n+\delta _{ij}dx^idx^j\right);`$
$`m,n=1,\mathrm{}4;i,j=1,\mathrm{}4,`$
$`e^\varphi =1+{\displaystyle \frac{1}{6}}{\displaystyle \frac{\kappa _{10}T_1}{\mathrm{\Omega }_7}}{\displaystyle \frac{1}{|x|^6}}+\mathrm{},`$
$`C_{0\sigma }^{(2)}={\displaystyle \frac{1}{2}}{\displaystyle \frac{\kappa _6Q_1}{\sqrt{2}\mathrm{\Omega }_3}}{\displaystyle \frac{1}{|y|^2}}+\mathrm{}.`$ (3.3)
Now $`\mathrm{\Omega }_7`$ is the area of the unit 7-sphere surrounding the string. The expression for the 2-form potential is the same as in the compactified case since it lives in the twisted sector. In it is shown that for the non-BPS D0-brane this kind of approach relating the uncompactified and the compactified solutions gives the same answer than the boundary state analysis of the uncompactified orbifold. This should be the case also for the D1-brane.
It is straightforward to check that this D1-brane solution solves the equations of motion derived from an action $`S_{\mathrm{untwisted}}+S_{\mathrm{twisted}}`$, where:
$$S_{\mathrm{untwisted}}=\frac{1}{2\kappa _{10}^2}d^{10}xe^{2\varphi }\sqrt{|\mathrm{det}g|}\left(R+4(\varphi )^2\right),$$
(3.4)
and the action corresponding to the twisted sector is proportional to the worldvolume effective action associated to a pair NS5A-O5. This reads, in string frame and to first order in $`\alpha ^{}`$:
$$\begin{array}{ccc}\hfill S_{\mathrm{twisted}}& & d^6y\sqrt{|\mathrm{det}g|}\left(1+(^{(3)})^2+\mathrm{}\right)d^6y\sqrt{|\mathrm{det}g|}=\hfill \\ & & \\ & =& d^6y\sqrt{|\mathrm{det}g|}(^{(3)})^2+\mathrm{}.\hfill \end{array}$$
(3.5)
Here $`^{(3)}`$ is the field strength of the self-dual 2-form potential of the NS5A-brane, and the self-duality condition is impossed at the level of the equations of motion. The metric is restricted to the position of the orientifold fixed plane.
## 4 The non-BPS D2-brane of IIB/$`(1)^{F_L}I_4`$
Performing now a T-duality transformation of the D1-brane solution along the $`y^4`$ direction we obtain a D2-brane solution of IIB/$`(1)^{F_L}I_4`$. This brane is located at one of the fixed points of the orbifold, with its worldvolume extending along the non-compact spacetime. Taking the non-BPS D1-brane in the compactified orbifold at the critical radii, where it is possible to construct periodic infinite arrays of strings, and applying the T-duality rules we find:
$`ds_{D2}^2=(1{\displaystyle \frac{\kappa _6T_2}{2\pi ^2\mathrm{\Omega }_2}}{\displaystyle \frac{1}{|y|}}+\mathrm{})(dt^2+d\sigma _1^2+d\sigma _2^2)+`$
$`+(1+{\displaystyle \frac{\kappa _6T_2}{2\pi ^2\mathrm{\Omega }_2}}{\displaystyle \frac{1}{|y|}}+\mathrm{})\left(\delta _{mn}dy^mdy^n+\delta _{ij}dx^idx^j\right);`$
$`m,n=1,\mathrm{}3;i,j=1,\mathrm{}4,`$
$`e^\varphi =1+{\displaystyle \frac{1}{2}}{\displaystyle \frac{\kappa _6T_2}{2\pi ^2\mathrm{\Omega }_2}}{\displaystyle \frac{1}{|y|}}+\mathrm{},`$
$`C_{0\sigma _1\sigma _2}^{(3)}={\displaystyle \frac{\kappa _6Q_2}{\sqrt{2}\mathrm{\Omega }_2}}{\displaystyle \frac{1}{|y|}}+\mathrm{}.`$ (4.1)
Here $`T_2`$ is the tension of the brane, $`Q_2`$ its charge and $`\mathrm{\Omega }_2`$ the area of the unit 2-sphere surrounding the 2-brane.
The same analysis of the previous section gives the following form for the solution in the uncompactified case:
$`ds_{D2}^2=(1{\displaystyle \frac{1}{5}}{\displaystyle \frac{\kappa _{10}T_2}{\mathrm{\Omega }_6}}{\displaystyle \frac{1}{|x|^5}}+\mathrm{})(dt^2+d\sigma _1^2+d\sigma _2^2)+`$
$`+(1+{\displaystyle \frac{1}{5}}{\displaystyle \frac{\kappa _{10}T_2}{\mathrm{\Omega }_6}}{\displaystyle \frac{1}{|x|^5}}+\mathrm{})\left(\delta _{mn}dy^mdy^n+\delta _{ij}dx^idx^j\right);`$
$`m,n=1,\mathrm{}3;i,j=1,\mathrm{}4,`$
$`e^\varphi =1+{\displaystyle \frac{1}{10}}{\displaystyle \frac{\kappa _{10}T_2}{\mathrm{\Omega }_6}}{\displaystyle \frac{1}{|x|^5}}+\mathrm{},`$
$`C_{0\sigma _1\sigma _2}^{(3)}={\displaystyle \frac{\kappa _6Q_2}{\sqrt{2}\mathrm{\Omega }_2}}{\displaystyle \frac{1}{|y|}}+\mathrm{}.`$ (4.2)
Now $`\mathrm{\Omega }_6`$ is the area of the unit 6-sphere surrounding the membrane.
This brane is electrically charged with respect to the 3-form potential of the NS5B-O5 system, or equivalently, magnetically charged with respect to its vector potential. Therefore, it solves the equations of motion derived from $`S_{\mathrm{untwisted}}+S_{\mathrm{twisted}}`$, with:
$$\begin{array}{ccc}\hfill S_{\mathrm{twisted}}& & d^6y\sqrt{|\mathrm{det}g|}\left(1+(\stackrel{~}{}^{(4)})^2+\mathrm{}\right)d^6y\sqrt{|\mathrm{det}g|}=\hfill \\ & & \\ & =& d^6y\sqrt{|\mathrm{det}g|}(\stackrel{~}{}^{(4)})^2+\mathrm{},\hfill \end{array}$$
(4.3)
where we have dualized the vector field of the NS5B-brane onto a 3-form potential with field strength $`\stackrel{~}{}^{(4)}`$, and the metric is restricted to the position of the orientifold fixed plane.
## 5 The non-BPS M1-brane of M-theory/$`\mathrm{\Omega }_\rho I_5`$
Oxidizing the D1-brane solution of the Type IIA theory on the orbifold we can obtain the expression for a stable M1-brane solution of M-theory orientifolded by $`\mathrm{\Omega }_\rho I_5`$. In the compact case we obtain:
$`d\widehat{s}_{M1}^2=(1{\displaystyle \frac{5}{6}}{\displaystyle \frac{\kappa _6T_1}{2\pi ^2\mathrm{\Omega }_3}}{\displaystyle \frac{1}{|y|^2}}+\mathrm{})(dt^2+d\sigma ^2)+`$
$`+(1+{\displaystyle \frac{1}{6}}{\displaystyle \frac{\kappa _6T_1}{2\pi ^2\mathrm{\Omega }_3}}{\displaystyle \frac{1}{|y|^2}}+\mathrm{})\left(\delta _{mn}dy^mdy^n+\delta _{ij}dx^idx^j\right)+`$
$`+(1+{\displaystyle \frac{2}{3}}{\displaystyle \frac{\kappa _6T_1}{2\pi ^2\mathrm{\Omega }_3}}{\displaystyle \frac{1}{|y|^2}}+\mathrm{})dz^2;`$
$`m,n=1,\mathrm{}4;i,j=1,\mathrm{}4,`$
$`\widehat{C}_{0\sigma }^{(2)}={\displaystyle \frac{1}{2}}{\displaystyle \frac{\kappa _6Q_1}{\sqrt{2}\mathrm{\Omega }_3}}{\displaystyle \frac{1}{|y|^2}}+\mathrm{}.`$ (5.1)
This solution has an $`SO(1,1)\times SO(4)\times SO(4)\times U(1)`$ symmetry, i.e. it corresponds to an asymmetric orbifold. This is in agreement with the results of , which show that the orbifold corresponding to M-theory on $`T^5/\mathrm{\Omega }_\rho I_5`$ has to be asymmetric so that the twisted sector cancels the gravitational anomalies involved in the construction. We see that asymptotically, i.e. in the region where the dilaton of the Type IIA theory is of order 1, the orbifold regains isotropy, also in agreement with .
Uplifting the D1-brane solution corresponding to Type IIA on the uncompactified orbifold we find an M1-brane on $`R^{1,5}\times (R^4\times S^1)/\mathrm{\Omega }_\rho I_5`$:
$`d\widehat{s}_{M1}^2=(1{\displaystyle \frac{5}{18}}{\displaystyle \frac{\kappa _{10}T_1}{\mathrm{\Omega }_7}}{\displaystyle \frac{1}{|x|^6}}+\mathrm{})(dt^2+d\sigma ^2)+`$
$`+(1+{\displaystyle \frac{1}{18}}{\displaystyle \frac{\kappa _{10}T_1}{\mathrm{\Omega }_7}}{\displaystyle \frac{1}{|x|^6}}+\mathrm{})\left(\delta _{mn}dy^mdy^n+\delta _{ij}dx^idx^j\right)+`$
$`+(1+{\displaystyle \frac{2}{9}}{\displaystyle \frac{\kappa _{10}T_1}{\mathrm{\Omega }_7}}{\displaystyle \frac{1}{|x|^6}}+\mathrm{})dz^2;`$
$`m,n=1,\mathrm{}4;i,j=1,\mathrm{}4,`$
$`\widehat{C}_{0\sigma }^{(2)}={\displaystyle \frac{1}{2}}{\displaystyle \frac{\kappa _6Q_1}{\sqrt{2}\mathrm{\Omega }_3}}{\displaystyle \frac{1}{|y|^2}}+\mathrm{}.`$ (5.2)
As we did in the previous sections we can interpret this solution as an M1-brane located at the origin of the $`z`$-circle, and find the expression for the corresponding solution in the completely uncompactified case:
$`d\widehat{s}_{M1}^2=(1{\displaystyle \frac{5}{21}}{\displaystyle \frac{\widehat{\kappa }_{11}\widehat{T}_1}{\mathrm{\Omega }_8}}{\displaystyle \frac{1}{|x|^7}}+\mathrm{})(dt^2+d\sigma ^2)+`$
$`+(1+{\displaystyle \frac{1}{21}}{\displaystyle \frac{\widehat{\kappa }_{11}\widehat{T}_1}{\mathrm{\Omega }_8}}{\displaystyle \frac{1}{|x|^7}}+\mathrm{})\left(\delta _{mn}dy^mdy^n+\delta _{ij}dx^idx^j\right)+`$ (5.3)
$`+(1+{\displaystyle \frac{4}{21}}{\displaystyle \frac{\widehat{\kappa }_{11}\widehat{T}_1}{\mathrm{\Omega }_8}}{\displaystyle \frac{1}{|x|^7}}+\mathrm{})dz^2.`$
Here $`\mathrm{\Omega }_8`$ is the area of the unit 8-sphere surrounding the string, and we have used $`\widehat{\kappa }_{11}=\kappa _{10}(2\pi R_z)^{1/2}`$, $`\widehat{T}_1=T_1\widehat{\kappa }_{11}/\kappa _{10}`$. $`\widehat{C}^{(2)}`$ remains the same since it lives in the twisted sector.
Finally, the contribution to the supergravity action from the twisted sector is proportional to the worldvolume action describing an M5-O5 system, which in quadratic approximation reads:
$$\widehat{S}_{\mathrm{twisted}}d^6\widehat{y}\sqrt{|\mathrm{det}\widehat{g}|}(\widehat{}^{(3)})^2+\mathrm{}$$
(5.4)
Here $`\widehat{}^{(3)}`$ is the field strength associated to the self-dual 2-form potential of the M5-brane worldvolume, with the self-duality condition impossed at the level of the equations of motion, and the metric is restricted to the position of the orientifold fixed plane.
### Acknowledgements
It is a pleasure to thank Laurent Houart for very interesting discussions, and Eduardo Eyras for pointing out some mistakes in a previous version of this paper.
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# Surface tension between kaon condensate and normal nuclear matter phase
## I INTRODUCTION
The possibility of different phase transitions taking place in the superdense interior of neutron stars has been the target of considerable interest during the last few decades, where pion and kaon condensation as well as quark deconfinement have been investigated. But only less than a decade ago was it realized that if the phase transition is of first order, then a geometrically structured extended region will form in the superdense interior of the neutron star, where the two phases are in equilibrium . The reason for this richness in structure is that a neutron star has two globally conserved charges, baryon number and electric charge, and two chemical potentials associated with these charges. Previous studies using the Maxwell construction could only ensure one chemical potential was common in the two phases, whereas the general phase equilibrium criteria by Gibbs ensure thermodynamical equilibrium for a system with any number of chemical potentials. The consequences are that the system is not locally charge neutral and a competition between Coulomb and surface energies are responsible for the geometrical structures. Moreover, the common pressure will vary with the proportion of the phases, and thus create an extended mixed phase region with structure in the neutron star.
For a first order deconfinement transition, studies of the detailed crystalline structure of the mixed phase region have always been hindered by the lack of a single good model describing both phases. This is in contrast to the first order transition to a kaon condensed phase described in , where both the normal nuclear matter phase and the kaon condensed phase are described by the same relativistic mean-field model, which allow us to calculate the profiles of all important quantities across the interface. From these profiles the Coulomb energy and the surface tension can be found, where, especially for the latter, only educated guesses were previously possible. A condensate consisting of negatively charged kaons is favored in neutron stars because they, contrary to other kaon types, can replace electrons as neutralizing agents .
Detailed knowledge of the structure of a possible mixed phase region at densities above saturation is important, irrespective of which first order phase transition is responsible for it, as it may have important consequences for transport and superfluid properties and rotation in the form of r-mode instabilities, non-canonical values of the braking index, and glitch phenomena in pulsars .
In the present paper we calculate the surface properties, i.e., the surface tension and the curvature coefficient, for a semi-infinite slab of normal nuclear matter in phase equilibrium with a semi-infinite slab of kaon condensed matter and show the resulting crystalline structure in the central part of the neutron star. Since we assure compliance with Gibbs phase equilibrium criteria, the two phases cannot be separately charge neutral, though overall the net charge vanishes. Thus for a infinite system we cannot explicitly take Coulomb interactions into account in the Lagrangian. However, it turns out that the typical radii of the geometrical structures are smaller than the Debye screening lengths of about 10 fm , and therefore it is a reasonable first approximation to ignore this effect in the calculation of the surface tension.
Section II contains a description of the non-uniform relativistic mean-field model used to describe both phases. In Sec. III the surface properties are described, whereas the consequences for the crystalline structure in a neutron star are illustrated in Sec. IV. Finally our results are summarized in Sec. V.
## II NON-UNIFORM RELATIVISTIC MEAN-FIELD MODEL
The relativistic mean-field model used here is described in detail in . In this model the nucleon and kaon interactions are treated on an equal footing which means both couple to the scalar meson $`\sigma `$, the vector meson $`𝝎`$, and the isovector meson $`𝝆`$, denoted $`\sigma `$, $`V_\mu `$, and $`\stackrel{}{R}_\mu `$, respectively. The Lagrangian for the nucleons is given by
$`_N=`$ $`\overline{\mathrm{\Psi }}_N(i\gamma ^\mu _\mu m_N^{}g_{\omega N}\gamma ^\mu V_\mu g_{\rho N}\gamma ^\mu \stackrel{}{\tau }_N\stackrel{}{R}_\mu )\mathrm{\Psi }_N`$ (3)
$`+{\displaystyle \frac{1}{2}}_\mu \sigma ^\mu \sigma {\displaystyle \frac{1}{2}}m_\sigma ^2\sigma ^2U(\sigma ){\displaystyle \frac{1}{4}}V_{\mu \nu }V^{\mu \nu }`$
$`+{\displaystyle \frac{1}{2}}m_\omega ^2V_\mu V^\mu {\displaystyle \frac{1}{4}}\stackrel{}{R}_{\mu \nu }\stackrel{}{R}^{\mu \nu }+{\displaystyle \frac{1}{2}}m_\rho ^2\stackrel{}{R}_\mu \stackrel{}{R}^\mu ,`$
where $`m_N^{}m_Ng_{\sigma N}\sigma `$, $`V_{\mu \nu }_\mu V_\nu _\nu V_\mu ,`$ $`\stackrel{}{R}_{\mu \nu }_\mu \stackrel{}{R}_\nu _\nu \stackrel{}{R}_\mu `$, and the scalar self-interactions are on the form $`U(\sigma )=(1/3)bm_N(g_{\sigma N}\sigma )^3+(1/4)c(g_{\sigma N}\sigma )^4`$, where $`b`$ and $`c`$ are constants. $`\mathrm{\Psi }_N`$ is the nucleon field operator and $`\stackrel{}{\tau }_N`$ is the isospin operator.
The kaon Lagrangian is given by
$$_K=𝒟_\mu ^{}K^{}𝒟^\mu Km_K^2K^{}K,$$
(4)
where $`𝒟_\mu _\mu +ig_{\omega N}V_\mu +ig_{\rho N}\stackrel{}{\tau }_K\stackrel{}{R}_\mu `$ and $`m_K^{}m_Kg_{\sigma K}\sigma `$. $`K`$ denotes the $`K^{}`$ field and $`\stackrel{}{\tau }_K`$ the kaon isospin operator.
In the mean-field approximation only the 0’th component of the vector fields and the isospin 3-component of the isovector field have finite mean values. The equation of motion for the kaon is
$$\left(𝒟_\mu 𝒟^\mu +m_K^2\right)K=0.$$
(5)
In a semi-infinite system the kaon field is only translational invariant in the transverse direction $`𝐳_{}`$ and can then be written as
$$K=\varphi (z)e^{i\left(\omega _K+𝐤_{}𝐳_{}\right)},$$
(6)
where $`𝐤_{}`$ is the transverse momentum, $`\omega _K`$ is the in-medium kaon energy, and $`\varphi (z)`$ is the kaon amplitude.
For s-wave condensation, i.e. $`𝐤=0`$, the kaon field equation is given by
$$^2\varphi +\left(\left(\omega _K+g_{\omega K}V_0+g_{\rho K}R_{0,3}\right)^2m_K^2\right)\varphi =0.$$
(7)
In the semi-infinite case $`^2=d^2/dz^2`$, where $`z`$ is the direction perpendicular to the surface.
The kaon density is
$$\rho _K=2\left(\omega _K+g_{\omega K}V_0+g_{\rho K}R_{0,3}\right)K^{}K,$$
(8)
and only for an infinite system $`\rho _K=2m_K^{}K^{}K`$ as seen from Eq. (7).
The equations of motion for the position dependent meson fields become
$`^2\sigma m_\sigma ^2\sigma =`$ $`g_{\sigma N}\left(\rho _s+bm_N(g_{\sigma N}\sigma )^2+c(g_{\sigma N}\sigma )^3\right)`$ (10)
$`2g_{\sigma K}m_K^{}K^{}K,`$
$$^2V_0m_\omega ^2V_0=g_{\omega N}(\rho _n+\rho _p)+g_{\omega K}\rho _K,$$
(11)
and
$$^2R_{0,3}m_\rho ^2R_{0,3}=g_{\rho N}(\rho _p\rho _n)+g_{\rho K}\rho _K,$$
(12)
where the neutron and proton densities, $`\rho _n`$ and $`\rho _p`$, and the scalar density $`\rho _s`$ are calculated in the Thomas-Fermi (local density) approximation.
$$\rho _i=\frac{1}{3\pi ^2}k_i^3,i=n,p$$
(13)
and
$`\rho _s={\displaystyle \frac{m_N^{}}{2\pi ^2}}`$ $`{\displaystyle \underset{i=n,p}{}}(k_i\sqrt{k_i^2+m_N^2}`$ (15)
$`m_N^2\mathrm{ln}\left({\displaystyle \frac{k_i+\sqrt{k_i^2+m_N^2}}{m_N^{}}}\right)).`$
Expressions for the local Fermi momenta $`k_i`$ are obtained from the Dirac equation for the nucleons
$$\mu _i=g_{\omega N}V_0g_{\rho N}R_{0,3}+\sqrt{k_i^2+m_N^2},$$
(16)
where the upper sign is used for neutrons. The electron chemical potential $`\mu _e=\mu _n\mu _p`$. The kaon amplitude is zero unless the condition $`\omega _K=\mu _K^{}=\mu _e`$ is fulfilled.
The model parameters $`g_{iN}/m_i`$ ($`i=\sigma ,\omega ,\rho `$), $`b`$, and $`c`$ can be algebraically determined from the five bulk properties of nuclear matter, which we take as: $`E/A=16.3`$ MeV, $`\rho _0=0.153`$ fm<sup>-3</sup>, $`a_{sym}=32.5`$ MeV, $`K=240`$ MeV, and $`m^{}/m_N=0.78`$. Because of the Laplacian terms in the meson field equations explicit values for the masses are required. We use for the $`𝝎`$ and $`𝝆`$ their rest masses $`m_\omega =782`$ MeV and $`m_\rho =768`$ MeV, while the $`\sigma `$ mass is determined from the surface properties of symmetric nuclear matter as described in , $`m_\sigma =390`$ MeV. The kaon coupling constants for the vector mesons are determined from the quark and isospin counting rule,
$$g_{\omega K}=g_{\omega N}/3\text{and}g_{\rho K}=g_{\rho N},$$
(17)
and the scalar coupling constant is fixed to the optical potential of the $`K^{}`$ at $`\rho _0`$ from
$$U_K(\rho _0)=g_{\sigma K}\sigma (\rho _0)g_{\omega K}V_0(\rho _0),$$
(18)
where we for the optical potential use $`U_K(\rho _0)=120`$ MeV.
The Gibbs conditions for an infinite system of a kaon condensed phase and a normal nuclear matter phase to be in thermodynamical equilibrium at zero temperature are
$`\mu _{i,N}`$ $`=`$ $`\mu _{i,K}i=n,e`$ (19)
$`P_N(\mu _n,\mu _e)`$ $`=`$ $`P_K(\mu _n,\mu _e),`$ (20)
where $`P_N`$ and $`P_K`$ are the pressures of the normal nuclear matter phase and the kaon condensed phase, respectively. These conditions combined with the global condition of electric charge conservation,
$$q_{total}=(1\chi )q_N(\mu _n,\mu _e)+\chi q_K(\mu _n,\mu _e),$$
(21)
where $`q`$ denotes the charge of the corresponding phase, and the field equations (without the Laplacian terms) allow us to solve for the bulk values of the fields, densities, and chemical potentials in each phase for any volume fraction of kaon phase $`\chi `$ between zero and one.
In order to calculate the surface tension between the kaon condensed phase and the normal nuclear matter phase, the profiles of the fields and densities have to be determined across the interface between semi-infinite slabs of each phase. This is done by simultaneously solving the four coupled differential equations for the $`K^{}`$ and the meson fields Eqs. (7,10-12) through a relaxation procedure, where initial guesses for the different profiles are relaxed to their equilibrium values. The boundary conditions at $`\pm \mathrm{}`$ are provided by the bulk values of the kaon amplitude and meson fields for each phase at a fixed $`\chi `$. In practice a 30 fm region with the interface placed approximately in the center is sufficient to fulfill the boundary conditions, see Fig. 1.
It is a surprisingly good approximation and a significant simplification to neglect the Laplacian terms in the meson field equations, and only keeping the term in the kaon field equation. Differentiating twice brings the kaon equation to the form $`d^2\varphi /dz^2=X/`$ $`\left(\left(\omega _K+g_{\omega K}V_0+g_{\rho K}R_{0,3}\right)^2m_K^2\right)`$. The denominator vanishes in the bulk kaon condensed phase in contrast to the meson field equations, which take the form $`d^2\sigma /dz^2Y/m_\sigma ^2`$ ($`X`$ and $`Y`$ denote two complicated functions). Thus the Laplacian term is generally most important for the kaon field equation. In Ref. the approximation was used to calculate the weak charge of droplets of kaon condensed phase to study effects on the neutrino opacity. We do both calculations and use the results from the approximative procedure as initial guess for the solution to the full set of four differential equations. This procedure ensures rapid convergence.
## III Surface Properties
In the transition zone between two phases the pressure tensor is no longer isotropic as it is for the homogeneous bulk phases. For a plane interface with the normal in the $`z`$direction, the pressure tensor can be split into a normal and a tangential component. The normal component of the pressure $`P`$ stays constant across the interface, cf. to Gibbs phase equilibrium conditions, whereas the tangential pressure component $`P_T(z)`$ changes as function of $`z`$ across the interface. The mechanical definition of the surface tension $`\sigma _G`$ of a plane interface is
$$\sigma _G=_{\mathrm{}}^{\mathrm{}}(PP_T(z))𝑑z,$$
(22)
where $`P_T(\mathrm{})=P_T(\mathrm{})=P`$. In nuclear physics this definition of the surface tension is denoted Gibbs definition .
Because we have assumed the leptonic species are homogeneously distributed throughout the system, they do not contribute to the surface tension. The pressure of the hadronic species $`P_H`$ is given by
$$P_H=\rho \frac{ϵ_H}{\rho }ϵ_H=\mu _n\rho _n+\mu _p\rho _p+\mu _e\rho _kϵ_H.$$
(23)
Therefore the surface tension of a plane interface between a kaon condensed phase and a normal nuclear matter phase can be written as
$`\sigma _G={\displaystyle _{\mathrm{}}^{\mathrm{}}}`$ $`dz`$ $`(ϵ_H(z)ϵ_{H,N}\mu _n(\rho _n(z)\rho _{n,N})`$ (25)
$`\mu _p(\rho _p(z)\rho _{p,N})\mu _e\rho _k(z)),`$
where the $`z`$ dependent quantities at minus infinity take the bulk values of the kaon condensed phase and at plus infinity take the bulk values of the normal nuclear matter phase. The energy density of the hadrons across the interface is given by
$`ϵ_H(z)={\displaystyle \frac{1}{2}}\left((\sigma )^2+m_\sigma ^2\sigma ^2\right)+{\displaystyle \frac{1}{2}}\left((V_0)^2+m_\omega ^2V_0^2\right)`$ (26)
$`+{\displaystyle \frac{1}{2}}\left((R_{0,3})^2+m_\rho ^2R_{0,3}^2\right)+2\left(\omega _K+g_{\omega K}V_0+g_{\rho K}R_{0,3}\right)^2\varphi ^2`$ (27)
$`+U(\sigma )+{\displaystyle \underset{i=n,p}{}}{\displaystyle \frac{1}{\pi ^2}}{\displaystyle _0^{k_i}}\sqrt{k^2+m_N^2}k^2𝑑k.`$ (28)
Notice, if the Laplacian terms are neglected except for the kaon field, the only consequence is that the squared gradient terms of $`V_0`$ and $`R_{0,3}`$ change sign.
The curvature coefficient is in Gibbs definition given by
$`\gamma _G`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}dz(zz_0)(ϵ_H(z)ϵ_{H,N}\mu _n(\rho _n(z)\rho _{n,N})`$ (30)
$`\mu _p(\rho _p(z)\rho _{p,N})\mu _e\rho _k(z)),`$
where the surface location, i.e., the position of the equivalent sharp surface,
$$z_0=\frac{_{\mathrm{}}^{\mathrm{}}z\rho ^{}(z)𝑑z}{_{\mathrm{}}^{\mathrm{}}\rho ^{}(z)𝑑z},$$
(31)
and the prime denotes the derivative with respect to $`z`$.
In Fig. 2 the surface tension and curvature coefficient densities (i.e., the integrands of Eqs. (25) and (30)), $`ϵ_s`$ and $`ϵ_c`$, respectively, are plotted both for the full set of differential equations and for the approximation with only a differential equation for the kaon field. The differences between the two surface tension density profiles are small. The exact calculation result in a slightly broader and more symmetric profile with a smaller peak value than the approximation. The curvature coefficient density $`ϵ_c=(zz_0)ϵ_s`$, and is therefore naturally more sensitive to the exact form, especially the width, of $`ϵ_s`$. This is also seen from the plot. The densities shown in Fig. 2 are for bulk parameters corresponding to $`\chi =0.107`$. For $`\chi `$ increasing from 0 to 1 the width of $`ϵ_s`$ increases from about 4 fm to 7 fm, while the peak value decreases.
The overall consequences for the surface tension and curvature coefficient are shown in Fig. 3, where $`\sigma _G`$ and $`\gamma _G`$ have been plotted as a function of $`\chi `$. Notice, the curvature coefficient is negative. The surface tension decreases monotonically with increasing volume fraction of the kaon phase, whereas the curvature coefficient increases. The values of $`\sigma _G`$ from the exact and approximative procedure deviate less than 0.4%. Thus the two curves lie on top of each other in this plot. This is in contrast to the curvature, where the difference is about 33% between the exact and approximative calculation. Furthermore, the absolute value of the curvature coefficient is comparable to the surface tension. This means the curvature energy is not negligible compared to the surface energy. But of course even higher order effects may mask the effect of the curvature energy as in the nuclear mass formula, where the apparent absence of a curvature term seem to be caused by higher order effects . For a spherical droplet of kaon phase with radius $`R=5`$ fm in a system where $`\chi 0`$, the surface energy is reduced about $`2\gamma _G/(\sigma _GR)=18`$ %, whereas the surface energy increases, since the curvature is a signed quantity, about 40 % for a spherical bubble of normal nuclear matter phase in the $`\chi 1`$ limit. We have here assumed, that both $`\sigma _G`$ and $`\gamma _G`$ calculated from a plane interface remain unaltered.
In Fig. 4 the exact calculation of $`\sigma _G`$ and $`\gamma _G`$ as a function of $`\chi `$ have been compared with the fits
$`\sigma _{fit}`$ $`=`$ $`\left(0.00786{\displaystyle \frac{ϵ_Kϵ_N}{\text{MeV/fm}\text{3}}}\right)^2\text{MeV/fm}^2`$ (32)
$`\gamma _{fit}`$ $`=`$ $`0.0976\left({\displaystyle \frac{ϵ_Kϵ_N}{\text{MeV/fm}\text{3}}}\right)^{3/4}\text{MeV/fm},`$ (33)
where the energy densities are in units of MeV/fm<sup>3</sup>. The fit deviate less than 5% for the curvature coefficient and less than 0.5% for the surface tension over the region where $`ϵ_Kϵ_N`$ decreases from 682 MeV to 358 MeV, i.e., about 50%. It is not exactly trivial to understand why the fits, especially for the surface tension, are so good. For the surface tension, Eq. (32) is equivalent to saying
$$(ϵ(z)ϵ_N)\frac{dϵ}{dz}P_NP(z),$$
(34)
whenever $`ϵ_s=P_NP(z)`$ is not negligible. This is indeed the case as long as $`d^2ϵ/dz^2`$ is negative but certainly not otherwise. We have not pursued the issue any further, but suspect this may be a generic property of phase equilibria described in mean-field theory in the Thomas-Fermi approximation.
## IV CRYSTALLINE STRUCTURE IN NEUTRON STARS
In Ref. the surface tension was assumed to be proportional to the difference in energy density between the two phases and the constant chosen rather arbitrarily, but in a way that ensured the sum of surface and Coulomb energies was always much smaller than the bulk energy. Furthermore, the Coulomb energy was, as usual, calculated under the assumption that the phase boundary is sharp. Figure 18 of that paper show the diameter $`D`$ and the spacing $`S`$ of the geometrical structures as a function of the radial coordinate for a star at the mass limit. As we have now shown, the surface tension is proportional to the difference in energy density squared, and it is between a factor of 2-5 smaller than what was guessed at in . Knowledge of the density profiles of the charged particles across the plane interface allow us to calculate the Coulomb energy for a soft boundary. It turns out that the Coulomb energy is reduced less than 15% for geometries with the typical small sizes encountered here. This translate into about a 5% reduction of the sum of surface and Coulomb energies. Therefore it is a good approximation to treat the phase boundary as being sharp with respect to calculating the Coulomb energy. Figure 5 is an updated version of Fig. 18 of Ref. . (Details of this type of calculation can be found in .) The differences between the two figures are minor. Overall the sizes of the geometrical structures just decrease about 30%. This is due to the fact, that to first approximation the location of the transition from one geometry to another is independent of $`\sigma _G`$. Furthermore, the sum of surface and Coulomb energies scale only as $`\sigma _G^{1/3}`$.
We have ignored the curvature term completely in these calculations, even though it is straight forward to include the curvature energy, see . However, if the curvature energy is negative, as it is for kaon droplets, the sum of finite size and Coulomb energies will not have a local minimum but will decrease to minus infinity with decreasing size of the structure. This is of course unphysical and higher order effects will prevent it from happening. These complicated higher order effects may, as previously mentioned, even to some extent mask the presence of a curvature term. Moreover, the curvature coefficient is very sensitive to the exact profile of the phase boundary as we have seen, so its real value is uncertain for the rather small structures encountered here. Therefore we choose to consider the value of the curvature coefficient as some measure for the amount of uncertainty in the surface tension calculation, but not more than that.
## V SUMMARY
Little is known about the equation of state for nuclear matter above saturation density, but it is expected that at least a phase transition to deconfined quark matter happens at a few times the saturation density. However, this phase transition may not be the only one encountered in neutron stars. A transition to a phase with a $`K^{}`$ condensate is also a possibility. If any of these phase transitions is of first order, a neutron star will have a mixed phase region in its dense interior, which is very likely to have some observable consequences.
For the first time the surface properties in the interface between normal nuclear matter and kaon condensed matter have been calculated, which makes it possible to study this mixed phase region in greater detail. Our calculations are only a first approximation, there are a number of complicating aspects, which we have ignored or only treated approximately - e.g., explicit consideration of the Coulomb field which results in screening effects; the validity of the assumption that the surface tension is the same for semi-infinite slabs as it is for small slabs, rods, and drops; and the importance of the curvature and even higher order terms compared to the surface term. Concerning the first two points, we can generally say that screening effects will reduce the Coulomb energy, likewise the surface energy will decrease due to the decrease in the surface tension when the system is squeezed. The reason for the latter is that the surface region will dominate a small system, so that only $`ϵ_s`$ and not $`dϵ_s/dz`$ is zero at the boundaries. For example a slab which is 10 fm thick and with a volume fraction of kaon condensed phase of about 0.4, the surface tension is reduced about 12% compared to the surface tension of the infinite system (30 fm thick), while the absolute value of the curvature coefficient drops about a factor of three. These two effects pull in opposite directions with regard to the size of the geometry which minimize the sum of Coulomb and surface energies, thus the overall effect on the size is expected to be only minor.
We have taken another step towards a better understanding of the mixed phase region for a first order transition involving a kaon condensate. We do, however, realize there is still much room for improvement.
## ACKNOWLEDGMENTS
We wish to thank M. Centelles and X. Viñas for providing their relaxation codes to us. M.B.C. acknowledges financial support by the Carlsberg Foundation. One of us (NKG) was supported by the Director, Office of Energy Research, Office of High Energy and Nuclear Physics, Division of Nuclear Physics, of the U.S. Department of Energy under Contract DE-AC03-76SF00098.
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# Theoretical interpretation of the experimental electronic structure of lens shaped, self-assembled InAs/GaAs quantum dots
## I Introduction: Using theory as a bridge between the structure and the electronic properties of quantum dots
Self-assembled, Stranski-Krastanow (SK) grown semiconductor quantum dots have recently received considerable attention as they exhibit a rich spectrum of phenomena including quantum-confinementhawrylak ; gaponenko ; mrs\_feb98 , exchange-splittingslandin98 , Coulomb charging/blockadedrexler94 ; medeiros95 ; fricke96 ; miller97 ; warburton97 ; schmidt97 ; warburton98 ; schmidt98 ; brunkov98 and multi-exciton transitionslandin98 ; dekel98 . Over the past few years a considerable number of high quality measurements of the electronic level stucture of these dot systems have been performed, using photoluminescence (PL)schmidt97 ; schmidt98 ; yang99 ; yang2000 ; berryman97 ; itskevich98 ; itskevich99 ; fry2000 ,photoluminescence luminescence excitation (PLE)landin98 ; dekel98 , capacitancedrexler94 ; medeiros95 ; fricke96 ; brunkov98 and far infra red (FIR) spectroscopyfricke96 ; pan98 ; sauvage99 ; sauvage97 ; pan98:2 . These measurements have been able to determine the electronic level structure to relatively high precision. In parallel with these measurements, several groups have also attempted to measure the geometry and composition of these dotsyang99 ; yang2000 ; metzger99 ; garcia97 ; rubin96 . So far, however, these measurements have failed to provide details of the shape, size, inhomogeneous strain and alloying profiles to a similar level of accuracy to that in which the electronic structure has been determined. As a result, the size of the dots were often used as adjustable parameters in models that fit experimental spectra. For example, using a single-band effective mass model, Dekel et. aldekel98 defined an “effective shape” (cuboid) and “effective dimension” that reproduced the measured excitonic transitions. Similar “parabolic dot” models have been assumed by Hawrylak et. alhawrylak .
The accuracy of single-band and multi-band effective mass methods was recently examined in a series of papers wood96 ; fu\_pressure ; wang\_cdse ; pryor98:2 ; wang2000 . In these works, the shape, size and composition of nanostructures were arbitrarily fixed, and the electronic structure was evaluated by successively improving the basis set, starting from single-band methods (effective-mass), going to six and eight band methods (k.p), and finally, using a converged, multi-band approach (plane-wave pseudopotentials). It was found that conventional effective-mass and k.p methods can sometimes significantly misrepresent the fully converged results even when the shape, size and composition was given. The observed discrepancies were both quantitative (such as band gap values, level spacings, Coulomb energies) and qualitative (absence of polarization anisotropy in square based pyramidal dotswang2000 , missing energy levelswang\_cdse ). As a result of these limitations these methods may not offer a reliable bridge between the electronic structure and atomic structure.
In this paper, we offer a bridge between recent measurements of the electronic structure and measurements of the atomic structure of the dots using accurate theoretical modeling. Modeling can determine if the calculated electronic structure resulting from an assumed shape, size, strain and alloying profiles agrees with the measured electronic structure or not. A theory that can perform such a “bridging function” must be accurate and reliable. The pseudopotential approach to this problem qualifies, in that any discrepancy between the predicted and measured electronic properties can be attributed to incorrectly assumed shape, size or alloying profile. We have studied a range of shapes, sizes and alloy profiles and find that a lens-shaped InAs dot with an inhomogeneous Ga alloying profile is in closest agreement with current measurements. In the following sections we attempt to provide a consistent theoretical interpretation of numerous spectroscopic properties of InAs/GaAs dots.
## II Outline of the Method of calculation
We aim to calculate the energy associated with various electronic excitations in InAs/GaAs quantum dots. These energies can be expressed as total energy differences and require four stages of calculation:
(i) Assume the shape, size and composition and compute the equilibrium displacements: We first construct a supercell containing both the quantum dot and surrounding GaAs barrier material. The shape, size and composition profile are taken as input and subsequently refined. Sufficient GaAs barrier is used, so that when periodic boundary conditions are applied to the system, the electronic and strain interactions between dots in neighboring cells is negligible. The atomic positions within the supercell are then relaxed by minimizing the strain energy described by an atomistic force fieldkeating ; pryor98 including bond bending, bond stretching and bond bending-bond stretching interactions (see section III.1). An atomic force field is similar to continuum elasticity approachespryor98 in that both methods are based on the elastic constants, $`\{C_{ij}\}`$, of the underlying bulk materials. However, atomistic approaches are superior to continuum methods in two ways, (a) they can contain anharmonic effects, and (b) they capture the correct point group symmetry, e.g. the point group symmetry of a square based, zinc blende pyramidal dot is $`C_{2v}`$, since the and \[1$`\overline{1}`$0\] directions are inequivalent while continuum methodspryor98 , find $`C_{4v}`$. More details of the atomistic relaxation are given in section III.1.
(ii) Setup and solve the pseudopotential single-particle equation: A single-particle Schrödinger equation is set up at the relaxed atomic positions, $`\{𝐑_{n\alpha }\}`$
$$\widehat{H}\psi _i(𝐫)=\{\frac{\beta }{2}^2+\underset{n\alpha }{}\widehat{v}_\alpha (𝐫𝐑_{n\alpha })\}\psi _i(𝐫)=ϵ_i\psi _i(𝐫).$$
(1)
The potential for the system is written as a sum of strain-dependent, screened atomic pseudopotentials, $`v_\alpha `$, that are fit to bulk properties extracted from experiment and first-principles calculations (see section III.2). The Schrödinger equation is solved by expanding $`\psi `$ in a linear combination of bulk states, $`\varphi _{nk}`$, from bands, $`n`$, and k-points, $`k`$,
$$\psi _i(𝐫,ϵ)=\underset{n,k}{}c_{n,k}^{(i)}\varphi _{nk}(𝐫,ϵ),$$
(2)
taken at a few strain values. The solution of Eqs. (1) and (2) provides the level structure and dipole transition matrix elements. More details on the solution of the Schrödinger equation are given in section III.3.
(iii) Calculate the screened, inter-particle many-body interactions: The calculated single particle wavefunctions are used to compute the electron-electron, electron-hole and hole-hole direct , $`J_{ee},J_{eh},J_{hh}`$, and exchange $`K_{ee},K_{eh},K_{hh}`$ Coulomb energies (see section III.4).
(iv) Calculate excitation energies as differences in total, many particle energies: For example, the difference between the total energy $`E_{11}[h_0^1e_0^1]`$ of a dot with a hole in level $`h_0`$ and an electron in level $`e_0`$ and the total energy $`E_{00}[h_0^0e_0^0]`$ of the unexcited dot is
$`E_{11}[h_0^1e_0^1]E_{00}[h_0^0e_0^0]`$ $`=`$ $`\left(ϵ_{e_0}ϵ_{h_0}\right)`$ (3)
$``$ $`J_{e_0h_0}+2K_{e_0h_0}\delta _{S,0},`$
where (in the absence of spin-orbit coupling) $`\delta _{S,0}=1`$ for triplet states, and 0 for singlet states. Analagous expressions exist for electron-addition experiments (see section III.4).
The main approximations involved in our method are: (a) the fit of the pseudopotential to the experimental data of bulk materials is never perfect (see Table 1) and (b) we neglect self-consistent iterations in that we assume that the screened pseudopotential drawn from a bulk calculation is appropriate for the dot. Our numerical convergence parameters are (i) the size of the GaAs barrier separating periodic images of the dots, and (ii) the number of bulk wavefunctions used in the LCBB expansion of the wavefunctions. To examine the effects of these approximations and convergences on the ultimate level of accuracy that can be obtained with our methodology we have first applied these methods to an InGaAs/GaAs quantum well(see section III.5), where experimental measurements of the shape, size, composition and transition energies are more established (see section III.5). We next describe the details of our method.
## III Details of the Method of Calculations
### III.1 Calculation of equilibrium atomic positions for a given shape
To calculate the relaxed atomic positions within the supercell, we use a generalization (G-VFF) of the original valence force field (VFF)keating model. Our implementation of the VFF includes bond stretching, bond angle bending and bond-length/bond-angle interaction terms in the VFF Hamiltonian. This enables us to accurately reproduce the $`C_{11}`$, $`C_{12}`$ and $`C_{44}`$ elastic constants in a zincblende bulk material. We have also included higher order bond stretching terms, which lead to the correct dependence of the Young’s modulus with pressure. The G-VFF total energy can be expressed as:
$`E_{VFF}`$ $`=`$ $`{\displaystyle \underset{i}{}}{\displaystyle \underset{j}{\overset{nn_i}{}}}{\displaystyle \frac{3}{8}}[\alpha _{ij}^{(1)}\mathrm{\Delta }d_{ij}^2+\alpha _{ij}^{(2)}\mathrm{\Delta }d_{ij}^3]`$ (4)
$`+`$ $`{\displaystyle \underset{i}{}}{\displaystyle \underset{k>j}{\overset{nn_i}{}}}{\displaystyle \frac{3\beta _{jik}}{8d_{ij}^0d_{ik}^0}}[(𝐑_j𝐑_i)(𝐑_k𝐑_i)`$
$``$ $`cos\theta _{jik}^0d_{ij}^0d_{ik}^0]^2`$
$`+`$ $`{\displaystyle \underset{i}{}}{\displaystyle \underset{k>j}{\overset{nn_i}{}}}{\displaystyle \frac{3\sigma _{ijk}}{d_{ik}^0}}\mathrm{\Delta }d_{ij}[(𝐑_j𝐑_i)(𝐑_k𝐑_i)`$
$``$ $`cos\theta _{jik}^0d_{ij}^0d_{ik}^0],`$
where $`\mathrm{\Delta }d_{ij}^2=\left[[(R_iR_j)^2d_{ij}^{0}{}_{}{}^{2}]/d_{ij}^0\right]^2`$. Here $`𝐑_i`$ is the coordinate of atom i and $`d_{ij}^0`$ is the ideal (unrelaxed) bond distance between atom types of $`i`$ and $`j`$. Also, $`\theta _{jik}^0`$ is the ideal (unrelaxed) angle of the bond angle $`jik`$. The $`^{nn_i}`$ denotes summation over the nearest neighbors of atom $`i`$. The bond stretching, bond angle bending, and bond-length/bond-angle interaction coefficients $`\alpha _{ij}^{(1)}(\alpha )`$, $`\beta _{jik}`$, $`\sigma _{jik}`$ are related to the elastic constants in a pure zincblende structure in the following way,
$`C_{11}+2C_{12}`$ $`=`$ $`{\displaystyle \frac{\sqrt{3}}{4d_0}}(3\alpha +\beta 6\sigma )`$
$`C_{11}C_{12}`$ $`=`$ $`{\displaystyle \frac{\sqrt{3}}{d_0}}\beta `$
$`C_{44}`$ $`=`$ $`{\displaystyle \frac{\sqrt{3}}{d_0}}{\displaystyle \frac{[(\alpha +\beta )(\alpha \beta \sigma ^2)2\sigma ^3+2\alpha \beta \sigma ]}{(\alpha +\beta +2\sigma )^2}}.`$ (5)
The second-order bond stretching coefficient $`\alpha ^{(2)}`$ is related to the pressure derivative of the Young’s modulus by $`\frac{dB}{dP}`$, where $`B=(C_{11}+2C_{12})/3`$ is the Young’s modulus. Note that in the standardkeating VFF which we have used previouslyjkim98 ; williamson99 ; williamson98 the last terms of Eq.(4) are missing, so $`\sigma =0`$ in Eq.(III.1). Thus there were only two free parameters ($`\alpha ,\beta `$) and therefore three elastic constants could not, in general, be fit exactly. The G-VFF parameters and the resulting elastic constants are shown in Table 2 for GaAs and InAs crystals. For an InGaAs alloy system, the bond angle and bond-length/bond-angle interaction parameters $`\beta `$, $`\sigma `$ for the mixed cation Ga-As-In bond-angle are taken as the algebraic average of the In-As-In and Ga-As-Ga values. The ideal bond angle $`\theta _{jik}^0`$ is 109 for the pure zincblende crystal. However, to satisfy Vegas’s law for the alloy volume, we find that it is necessary to use $`\theta _{GaAsIn}^0=110.5^{}`$ for the cation mixed bond angle.
As a simple test of this G-VFF for alloy systems, we compared the relaxed atomic positions from G-VFF with pseudopotential LDA results for a (100) (GaAs)<sub>1</sub>/(InAs)<sub>1</sub> superlattice where the $`c/a`$ ratio is fixed to 1, but we allow energy minimizing changes in the overall lattice constant ($`a_{eq}`$) and the atomic internal degrees of freedom ($`u_{eq}`$). We find $`a_{eq}^{LDA}=5.8612`$ Å and $`u_{eq}^{LDA}=0.2305`$, while the G-VFF results are $`a_{eq}^{GVFF}=5.8611`$ Å and $`u_{eq}^{GVFF}=0.2305`$. In comparison the original VFF yields $`a_{eq}^{VFF}=5.8476`$ Å and $`u_{eq}^{VFF}=0.2303`$.
### III.2 The Empirical Pseudopotential Hamiltonian
We set up the single-particle Hamiltonian as
$$\widehat{H}=\frac{\beta }{2}^2+\underset{n\alpha }{}\widehat{v}_\alpha (𝐫𝐑_{n\alpha }),$$
(6)
where $`𝐑_{n\alpha }`$ is the G-VFF relaxed position of the n<sup>th</sup> atom of type $`\alpha `$. Here $`\widehat{v}_\alpha (𝐫)`$ is a screened empirical pseudopotential for atomic type $`\alpha `$. It contains a local part and a nonlocal, spin-orbit interaction part.
The local potential part is designed to include dependence on the local hydrostatic strain Tr$`(ϵ)`$:
$$v_\alpha ^{loc}(r;ϵ)=v_\alpha ^{eq}(r;0)[1+\gamma _\alpha \mathrm{Tr}(ϵ)],$$
(7)
where the $`\gamma _\alpha `$ is a fitting parameter. The zero strain potential $`v_\alpha ^{eq}(r;0)`$ is expressed in reciprocal space q as
$$v(q)=a_0(q^2a_1)/[a_2e^{a_3q^2}1].$$
(8)
The local hydrostatic strain $`Tr(ϵ)`$ for a given atom at $`𝐑`$ is defined as $`\mathrm{\Omega }_R/\mathrm{\Omega }_01`$, where $`\mathrm{\Omega }_R`$ is the volume of the tetrahedron formed by the four atoms bonded to the atom at $`𝐑`$. $`\mathrm{\Omega }_0`$ is the volume of that tetrahedron in the unstrained condition. The need for explicit dependence of the atomic pseudopotential on strain in Eq.(7) results from the following: While the description in Eq.(6) of the total pseudopotential as a superposition of atomic potentials situated at specific sites, $`\{R_{n\alpha }\}`$, does capture the correct local symmetries in the system, the absence of a self-consistent treatment of the Schrödinger equation deprives the potential from changing in response to strain. In the absence of a strain-dependent term, the volume dependence of the energy of the bulk valence band maximum is incorrect. While self-consistent descriptions show that the volume deformation potential $`a_v=dE_v/d\mathrm{ln}\mathrm{\Omega }`$ of the valence band maximum is negative for GaAs, GaSb, InAs, InSb and for all II-VI this qualitative behavior can not be obtained by a non-self-consistent calculation that lacks a strain dependent pseudopotential.
The nonlocal part of the potential describes the spin-orbit interaction,
$`H_{so}`$ $`=`$ $`{\displaystyle \underset{n\alpha }{}}\widehat{V}_\alpha ^{so}(R_{n\alpha })`$
$``$ $`{\displaystyle \underset{n\alpha }{}}{\displaystyle \underset{l}{}}V_{l,\alpha }^{so}(rR_{n\alpha })|l_{R_{n\alpha }}𝐋𝐒l|_{R_{n\alpha }},`$
where $`|l_{R_{n\alpha }}`$ is a projector of angular momentum $`l`$ centered at $`R_{n\alpha }`$, $`𝐋`$ is the spatial angular momentum operator, $`𝐒`$ is the Dirac spin operator, and $`V_{l,\alpha }^{so}(r)`$ is a potential describing the spin-orbit interaction.
In Eq(6), the kinetic energy of the electrons has been scaled by a factor of $`\beta `$. The origin of this term is as follows: In an accurate description of the crystal band structure, such as the GW methodhedin99 , a general, spatially non-local potential, $`V(r,r^{})`$, is needed to describe the self-energy term. In the absence of such a term the occupied band width of an inhomogeneous electron gas is too large compared to the exact many-body result. To a first approximation, however, the leading effects of this non-local potential, $`V(r,r^{})`$, can be represented by scaling the kinetic energy. This can be seen by Fourier transforming $`V(r,r^{})`$ in reciprocal space, $`q`$, then making a Taylor expansion of $`q`$ about zero. We find that the introduction of such a kinetic energy scaling, $`\beta `$ permits a simultaneous fit of both the effective masses and energy gaps. In this study, we fit $`\beta =1.23`$ for both GaAs and InAs.
The pseudopotential parameters in Eqs(7) and (8) were fitted to the bulk band structures, experimental deformation potentials and effective masses and first-principles calculations of the valence band offsets of of GaAs and InAs. The alloy bowing parameter for the GaInAs band gap (0.6 eV) is also fitted. The pseudopotential parameters are given in Table 3 and their fitted properties are given in Table 1bornstein . We see that unlike the LDA, here we accurately reproduce the bulk band gaps and the bulk effective masses. One significant difference in our parameter set, to that used in conventional k.p studies, is our choice of a negative magnitude for the valence band deformation potential, $`a_v`$, which we have obtained from LAPW calculationsfranceschetti94 .
The present InAs and GaAs pseudopotentials have been systematically improved relative to our previous InAs and GaAs potentialsjkim98 ; williamson99 ; williamson98:2 ; wang99 ; wang99:2 , although the functional form has remained the same. Firstly, the pseudopotentials for InAs and GaAs used in Ref.jkim98 ; wang99 did not include the spin-orbit interaction. In Refs.williamson98:2 ; williamson99 ; wang99:2 we used potentials that included the spin-orbit interaction, but were not able to simultaneously, accurately fit the electron effective and the zone center band gap, due to the lack of the above $`\beta `$ parameter. The potential used here is identical to that used in Refs.wang99 ; wang2000 .
### III.3 Calculating the single particle eigenstates
One could use a straight forward expansion of the single particle wavefunctions in a plane wave basis set, as we have previously done in Refs.williamson98 ; williamson99 ; jkim98 . However, as was shown in Refs.wang99 ; wang2000 ; wang97:2 , a more economical representation is to use the Linear Combination of Bulk Bands (LCBB) methodwang99 ; wang2000 ; wang97:2 . Within the LCBB the eigenstates of the pseudopotential Hamiltonian are expanded in a basis of bulk Bloch orbitals
$$\psi _i(𝐫)=\underset{s}{}\underset{n,k}{}c_{s,n,k}^{(i)}u_{s,n,k}(𝐫)e^{i𝐤.𝐫},$$
(10)
where $`u_{s,n,k}(𝐫)`$ is the cell periodic part of the bulk Bloch wavefunction for structure, $`s`$, at the $`n^{th}`$ band and the k<sup>th</sup> k-point, $`k`$. These states form a physically more intuitive basis than traditional plane waves therefore the number of bands and k-points can be significantly reduced to keep only the physically important bands and k-points (around the $`\mathrm{\Gamma }`$ point in this case). This method was recently generalized to strained semiconductor heterostructure systemswang99 and to include to spin-orbit interactionwang99:2 . In this paper use an LCBB basis derived from four structures, $`s`$. These structures are (i) unstrained, bulk InAs at zero pressure, (ii) unstrained, bulk GaAs at zero pressure, (iii) bulk InAs subjected to the strain value in the center of the InAs dot, and (iv) bulk InAs subjected to the strain value at the tip of the InAs dot. By interpolating the strain profile between these four structures, the basis is able to accurately describe all the strain in the system. The wavevectors, $`\{k\}`$, used here include all allowed values within $`16\pi /L`$ of the zone center, where $`L`$ is the supercell size. For calculations of electron states, the band indices, $`n`$, include only the band around the $`\mathrm{\Gamma }_{1c}`$ point. For the hole states we also include the three bands around the $`\mathrm{\Gamma }_{15v}`$ point. This basis set produces single particle energies that are converged with respect to basis size, to within 1 meV.
### III.4 Constructing the energies of different electronic configurations
Using screened Hartree Fock theory, the energy associated with loading $`N`$ electrons into a quantum dot can be expressedfranceschetti2000 as
$$E_N=\underset{i}{}(ϵ_i+\mathrm{\Sigma }_i^{pol})n_i+\underset{i<j}{}(J_{ij}^{ee}K_{ij}^{ee})n_in_j,$$
(11)
where $`ϵ_i`$ are the single-particle energies and $`\mathrm{\Sigma }_i^{pol}`$ are the polarization self-energies of the $`i^{th}`$ electron state , $`J_{ij}^{ee}`$ and $`K_{ij}^{ee}`$ are the direct and exchange Coulomb integrals between the $`i^{th}`$ and $`j^{th}`$ electronic states and $`n_i`$ are the occupation numbers ($`_in_i=N`$). As shown in Ref.franceschetti2000 , for free standing, colloidal quantum dots the dielectric constant inside the dot is dramatically different to that outside (vacuum) and hence the polarization self-energy, $`\mathrm{\Sigma }_i^{pol}`$, is very significant ($``$1 eV). For self assembled InAs dots embedded in GaAs, the dielectric constants of InAs and GaAs are similar ($`ϵ_{\mathrm{}}=12.3,10.6`$) and we calculate this term as $``$1 meV, and hence we choose to neglect it here. The direct and exchange Coulomb energies, are defined as
$`J_{ij}`$ $`=`$ $`{\displaystyle \frac{|\psi _i(𝐫_1)|^2|\psi _j(𝐫_2)|^2}{\overline{ϵ}(𝐫_1𝐫_2)|𝐫_1𝐫_2|}𝑑𝐫_1𝑑𝐫_2}`$
$`K_{ij}`$ $`=`$ $`{\displaystyle \frac{\psi _i^{}(𝐫_1)\psi _i(𝐫_2)\psi _j^{}(𝐫_2)\psi _j(𝐫_1)}{\overline{ϵ}(𝐫_1𝐫_2)|𝐫_1𝐫_2|}𝑑𝐫_1𝑑𝐫_2},`$ (12)
where $`\overline{ϵ}`$ is a phenomenological, screened dielectric functionwilliamson98:2 containing a Thomas Fermi electronic component and an ionic component from Ref.haken . Our exchange automatically includes both short and long range components.
Denoting electron levels as $`e_0,e_1,e_2`$…, hole levels as $`h_0,h_1,h_2`$… and the number of electrons and holes as $`N`$ and $`M`$, the total energy, $`E_{MN}`$, is
$`E_{MN}`$ $`=`$ $`{\displaystyle \underset{i}{}}ϵ_{h_i}m_i+{\displaystyle \underset{i<j}{}}(J_{ij}^{hh}K_{ij}^{hh})m_im_j`$ (13)
$`+`$ $`{\displaystyle \underset{i}{}}ϵ_{e_i}n_i+{\displaystyle \underset{i<j}{}}(J_{ij}^{ee}K_{ij}^{ee})n_in_j`$
$``$ $`{\displaystyle \underset{ij}{}}(J_{ij}^{eh}K_{ij}^{eh})n_im_j,`$
where $`n_i`$ and $`m_i`$ are the electron and hole occupation numbers respectively and $`_in_i=N`$ and $`_im_i=M`$. Using Eq.(13), in the strong confinement regime where kinetic energy effects dominate over the effects of exchange and correlation, an exciton involving electrons excited from hole state $`i`$ to electron state $`j`$ can be expressed as
$$E_{ij}^{exciton}=\left(ϵ_{e_j}ϵ_{h_i}\right)J_{ji}^{eh}+K_{ji}^{eh}\delta _{S,0}.$$
(14)
To study charged dots, if one assumes the electron states are occupied in order of increasing energy (Aufbau principle), the total energy of a dot charged with $`N`$ electrons, $`E_{0N}`$, is
$`E_{00}[e_0^0]`$ $`=`$ $`0`$
$`E_{01}[e_0^1]`$ $`=`$ $`ϵ_{e_0}`$
$`E_{02}[e_0^2]`$ $`=`$ $`2ϵ_{e_0}+J_{e_0,e_0}`$
$`E_{03}[e_0^2e_1^1]`$ $`=`$ $`\left(2ϵ_{e_0}+ϵ_{e_1}\right)+\left[J_{e_0,e_0}+2J_{e_0,e_1}\right]K_{e_0,e_1}`$
$`E_{04}[e_0^2e_1^2]`$ $`=`$ $`\left(2ϵ_{e_0}+2ϵ_{e_1}\right)+\left[J_{e_0,e_0}+J_{e_1,e_1}+4J_{e_0,e_1}\right]`$ (15)
$``$ $`2K_{e_0,e_1}.`$
As indicated in section II, our wavefunctions, $`\{\psi _i\}`$, are not iterated to self-consistency. This affects the magnitude of the direct and exchange Coulomb integrals. We have previously examined the accuracy of this perturbative treatment for colloidal InAs dots by comparing the non-self-consistent Coulomb energy with that obtained self consistentlyfranceschetti97 . The differences were negligible.
### III.5 Quantum well tests
To test the above methods, we first calculated the energy levels in a quantum well, and compared the results with experiment. In Fig. 1(a), we compare the calculated electron-heavy hole transition energies for a 96 Å In<sub>0.24</sub>Ga<sub>0.76</sub>As quantum well inside a GaAs matrix. The peaks in the experimental spectra occurgershoni89 at 1.275, 1.395 and 1.538 eV. Our calculated transitions occur at 1.290, 1.404 and 1.545 eV respectively. Figure 1(b) compares the band gap of a In<sub>0.22</sub>Ga<sub>0.78</sub>As quantum well as a function of its thickness. The measured band gapslaymarie95 for quantum wells with thicknesses of 6 and 18 ML are 1.458 and 1.351 eV. Our calculated values are 1.466 and 1.366 eV.
## IV Physical quantities to compare with experiment
The quantities we use to characterize the electronic structure are illustrated in Fig. 2 which shows a schematic layout of the electron and hole single-particle energy levels in a quantum dot. Assuming that all levels are spatially nondegenerate (thus having only spin degeneracy), we mark the electron levels as $`e_0,e_1,e_2`$…. and the hole levels as $`h_0,h_1,h_2`$. The level $`e_0`$ is sometimes called “s-like”, whereas $`e_1`$ and $`e_2`$ are called “p-like”, and $`e_3`$ and $`e_4`$ are called “d-like”. Since the GaAs environment of the InAs dots is largely unstrained, it is convenient to set as a reference energy the VBM of GaAs as $`E=0`$, and the CBM of GaAs as E=1520 meV. All energy levels can be referenced with respect to these band edges.
For the electron levels, the quantities that we consider are:
(i) The number of dot-confined electron states, $`N_e`$.
(ii) The spacing $`\delta _{sp}=ϵ_{e_1}ϵ_{e_0}`$ between “s-like” and “p-like” electron states.
(iii) The splitting $`\delta _{pp}=ϵ_{e_2}ϵ_{e_1}`$ between the “p-like” electron states.
(iv) The spacing $`\delta _{pd}=ϵ_{e_3}ϵ_{e_2}`$ between “p-like” and “d-like” electron states.
(v) The “binding energy” of the first electron level, $`e_0`$, with respect to the GaAs conduction band minimum, $`\mathrm{\Delta }E(e)=E_{GaAs,CBM}ϵ_{e_0}`$.
(vi) The position of the bottom of the band for the 2D InAs “wetting layer” (WL) with respect to the GaAs CBM, $`\mathrm{\Delta }E_{WL}^{(e)}=E_{GaAs,CBM}E_{WL}^{(e)}`$.
(vii) Inter-electron direct $`J_{e_i,e_j}^{ee}`$ and exchange $`K_{e_i,e_j}^{ee}`$ Coulomb energies.
For the hole levels we consider are:
(i) The number, $`N_h`$, of dot-confined hole states.
(ii) The intra-band spacings of the hole levels, $`\delta _{ij}^{(h)}=ϵ_{h_j}ϵ_{h_i}`$.
(iii) The “binding energy” of the first hole level, $`h_0`$, with respect to the GaAs valence band maximum, $`\mathrm{\Delta }E(h)=E_{GaAs,VBM}+ϵ_{h_0}`$.
(iv) The position of the top of the band for the 2D InAs “wetting layer” (WL) with respect to the GaAs VBM, $`\mathrm{\Delta }E_{WL}^{(h)}=E_{GaAs,VBM}+E_{WL}^{(h)}`$.
(v) Inter-hole direct $`J_{h_i,h_j}^{hh}`$ and exchange $`K_{h_i,h_j}^{hh}`$ Coulomb energies.
Finally, for the recombination of electrons and holes, we consider:
(i) The excitonic energies, $`E_{ij}^{exciton}`$, as defined in Eq.(14). By subtracting calculated values for the single particle energies $`ϵ_{e_j}`$ and $`ϵ_{h_i}`$ from measured optical excitation energies one can estimate the electron-hole direct Coulomb energies $`J_{h_ie_j}`$.
(ii) The ratio of absorption intensities for light polarized along and \[1$`\overline{1}`$0\] directions, defined as
$$\lambda =\frac{P_{[110]}}{P_{[1\overline{1}0]}}=\frac{<\psi _{e_0}|r_{[110]}|\psi _{h_0>^2}}{<\psi _{e_0}|r_{[1\overline{1}0]}|\psi _{h_0}>^2}.$$
(16)
This ratio can deviate from unity due to three reasons; (a) The dots has different dimensions in the and \[1$`\overline{1}`$0\] directions. We refer to this as the the “geometric factor”. (b) The atomistic zincblende symmetry makes the two directions symmetry inequivalent even if the lengths along the two directions are equal. We refer to this as the “atomic symmetry factor”. One manifestation of this affect is that if the strain is calculated atomistically, it is different in the two directions even in the absence of a geometric factorpryor98 . (c) A piezoelectric field that breaks the symmetry. Previous studiesyang2000 have shown that this effect is negligible in InAs/GaAs dots so we will neglect it here. k.p calculations neglect the “atomic symmetry” factor (except for the small effect of strain asymmetry), but retain the “geometric factor”. Pseudopotential calculations retain both effects. For example, in a square based pyramid (where by definition the “geometric factor” does not contribute), k.p produces $`\lambda =1`$, while pseudopotential theory gives $`\lambda =1.2`$ (see Table 5). This shows that there is not a simple mapping from dot shape to polarization anisotropy, $`\lambda `$.
(iii) Excitonic dipole: As the center of the electron and hole wavefunctions do not exactly coincide with each other, it is possible that an exciton will exhibit a detectable dipole moment,
$$d_{h_i,e_j}=\psi _{h_i}|\widehat{r}|\psi _{h_i}\psi _{e_j}|\widehat{r}|\psi _{e_j}.$$
(17)
The quantities defined above characterize the electronic structure. Next, in section V, we will next provide all of these quantities from our calculations, and then in section VI we will extract measured values of these quantities from the available experiments.
## V Theoretical Results
The electronic structure of a series of GaInAs/GaAs self-assembled quantum dots was calculated using the methodology described in Section II. We have chosen to focus on the well established “lens shaped” dot geometry from Refs.drexler94 ; medeiros95 ; fricke96 ; miller97 ; warburton97 ; schmidt97 ; warburton98 ; schmidt98 . The shape of this dot is shown in Fig. 3. The profile is obtained by selecting the section of a pure InAs sphere that yields a circular base with diameter 252Å and a height of 35 Å. The main experimental uncertainty about this dot is the composition profile. It is not known if the dots are pure InAs or if Ga has diffused into the dots. For comparison, we also calculate the electronic structure of a square based InAs pyramid with a base of 113Å and a height of 56Å. This is not believed to be a realistic geometry, however, it has been used as a benchmark for many previous theoretical calculationswang2000 ; jkim98 ; williamson99 ; grundman95 ; jaros96 and we include it here for comparison purposes. In the following sections these two geometries will be referred to as the “lens” and the “pyramid”. The results of our calculations are shown in Table 5 and Fig. 4.
### V.1 Confined electron states
Figure 4 shows the calculated square of the envelope function for the electron states in the pyramidal and lens shaped InAs/GaAs quantum dots. For the lens shaped dot, the electron states can be approximately interpreted as eigenstates of the $`\widehat{L_z}`$ operatorhawrylak . Here we plot only the first 6 bound states corresponding to $`l_z=0,\pm 1`$ and $`\pm 2`$. The first state $`e_0`$, has $`l_z=0`$ and is commonly described as $`s`$-like as it has no nodes. The $`e_1`$ and $`e_2`$ states have $`l_z=\pm 1`$, and are $`p`$-like with nodal planes (110) and ($`\overline{1}`$10). The $`e_3,e_4`$ and $`e_5`$ states have $`l_z=\pm 2`$ and 0 respectively and are commonly described as $`d_{x^2y^2}`$, $`d_{xy}`$ and $`2s`$ respectively. Due to the underlying zincblende atomistic structure, the $`C_{\mathrm{}}`$ symmetry is reduced to $`C_{2v}`$. Hence, the $`e_0`$ to $`e_5`$ states correspond to the $`a_1,b_1,b_2,a_1,a_2`$ and $`a_1`$ irreducible representations of the $`C_{2v}`$ group, rather than eigenstates of $`\widehat{L_z}`$. This allows states $`e_0,e_3`$ and $`e_5`$ to couple. This coupling is evident, for example, in the larger charge density along compared to \[1$`\overline{1}`$0\] in the $`e_3`$ state, due to its coupling with $`e_1`$. The observable effect of this $`C_{2v}`$ symmetry is to split the $`e_1`$ and $`e_2`$ $`p`$-states, $`\delta _{pp}`$, and the $`e_3`$ and $`e_4`$ $`d`$-states, $`\delta _{dd}`$. The alignment of the $`e_1`$ and $`e_2`$ $`p`$-states states along the and \[1$`\overline{1}`$0\] directions also results from the underlying zincblende lattice structure. Note, this analysis neglects the effects of the spin-orbit interaction which reduces the $`C_{2v}`$ group to a double group with the same single representation for all the states. In our calculations the spin-orbit interaction is included, but is produces no significant effects for the electron states.
The electron states in the pyramidal dot also belong to the $`C_{2v}`$ group and show a one-to-one correspondence with those in the lens shaped dot. However, there are only 5 bound states in the pyramidal dot due to its smaller size. Here we define an electron state as bound if its energy is below that of the unstrained, bulk GaAs conduction band edge.
The calculated values of the $`s`$-$`p`$ and $`p`$-$`d`$ energy spacings, $`\delta _{sp}`$, and, $`\delta _{pd}`$, for the lens and pyramidal shaped dots are 65 and 68 meV and 108 and 64 meV respectively. The splitting of the two $`p`$ states, $`\delta _{pp}=e_2e_1`$ are 2 and 26 meV respectively. The calculated values of the electron binding energy, $`\mathrm{\Delta }E(e)`$, are 271 and 171 meV respectively. The electron-electron direct Coulomb energies, $`J_{e_0e_0}^{ee}`$, $`J_{e_1e_1}^{ee}`$ and $`J_{e_0e_1}^{ee}`$ in the lens and pyramidal dots are calculated as 32, 25 and 25 meV and 40, 35 and 36 respectively. On applying a magnetic field in the growth direction, we calculate an increase in the splitting of the two $`p`$ states ($`e_2e_1`$) in the lens shaped dot from 2 to 20 meV. Details of this magnetic field calculation will be given in a future publicationshumway2000 . Finally, the energy of the electron wetting layer level, $`\mathrm{\Delta }E_{WL}^{(e)}`$, with thicknesses of 1 and 2 ML is 15 and 24 meV below the CBM of unstrained bulk GaAs.
### V.2 Confined hole states
Figure 4 shows calculated wavefunctions squared for the hole states in pyramidal and lens shaped InAs/GaAs quantum dots. Unlike the electron states, the hole states cannot be approximated by the solutions of a single band Hamiltonian. Instead there is a strong mixing between the original bulk Bloch states with $`\mathrm{\Gamma }_{8v}`$ and $`\mathrm{\Gamma }_{7v}`$ symmetry. The larger effective mass for holes results in a reduced quantum confinement of the hole states and consequently many more bound hole states. Only the 6 bound hole states with the highest energy are shown in Figure 4.
The calculated values of the $`h_0`$-$`h_1`$ , $`h_1`$ -$`h_2`$ and $`h_2`$-$`h_3`$ hole level spacings for the pyramidal and lens shaped dots are 8,7 and 6 meV and 15, 20 and 1 meV respectively. The calculated hole binding energies, $`\mathrm{\Delta }E(e)`$, are 194 and 198 meV. We calculate the highest energy hole level in pure InAs wetting layers, $`\mathrm{\Delta }E_{WL}^{(h)}`$, with thicknesses of 1 and 2 ML to reside 30 and 50 meV above the VBM of unstrained bulk GaAs. The hole-hole Coulomb energies, $`J_{h_0h_0}^{hh}`$, are 25 and 31 meV.
### V.3 Electron-hole excitonic recombination
Figure 5 shows our calculated single exciton absorption spectrum for a pure InAs, lens shaped dot with a base of 252 Å and a height of 35 Å. The energies of each of the absorption peaks are calculated from Eq.(14). The ratios of the dipole matrix elements for light polarized along and \[1$`\overline{1}`$0\] are calculated from Eq.(16). Figure 5, illustrates that, for a lens shaped dot, both the conventional $`e_ih_i`$ transitions and additional, $`e_1h_2,e_2h_1,e_3h_4`$ and $`e_4h_3`$ transitions are strongly allowed. The ratio of the polarization anisotropies, $`\lambda `$, are shown in Table 4. As a result of the circular symmetry of the lens shaped dot, we calculate a polarization ratio of $`\lambda =1.03`$ for the $`e_0h_0`$ transition. This value is in contrast to that calculated value for a pyramidal dot of $`\lambda =1.2`$wang99:2 . For the higher angular momentum transitions we find larger deviations from unity. The magnitude of the ratios, follows the polarization of the wavefunctions shown in Fig. 4. For example we find ratios greater and then less than unity for the $`e_1h_1`$ and $`e_2h_2`$ transitions, as reflected by the elongations of the $`e_1,h_1`$ and $`e_2,h_2`$ wavefunctions along the and \[1$`\overline{1}`$0\] directions.
We calculate ground state electron-hole direct Coulomb energies, $`J_{e_0h_0}^{eh}`$, of 37 and 31 meV in the lens shaped and pyramidal dots. The calculated ground state electron-hole exchange energies, $`K_{e_0h_0}^{eh}`$ are an order of magnitude smaller, with values of 3 and 0.2 meV. These yield excitonic band gaps of 1.03 and 1.12 respectively. The calculated polarization anisotropy ratios \[Eq.(16)\] for light polarized along and \[1$`\overline{1}`$0\] directions are $`\lambda =1.03`$ and 1.20 for the lens and pyramidal shapes respectively. The calculated excitonic dipoles \[Eq.(17)\] are -3.1 and 0.16Å respectively. A positive dipole is defined as the center of the hole wavefunction being located above the center of the electron wavefunction.
## VI Analysis of pertinent experimental measurements
### VI.1 The intra-band $`s`$-$`p`$ and $`p`$-$`d`$ electron energy spacings
Measurements of the spacing between the $`e_0`$ and $`e_1`$ like electron levels ($`s`$-like and $`p`$-like) are based on infra red absorption. For the lens shaped dots, Fricke et. al.fricke96 load electrons into the dots by growing a sample consisting of a n-type doped layer, a tunneling barrier, a layer of InAs/GaAs lens shaped dots, a GaAs spacer and a GaAs/AlAs short period superlattice (SPS). By applying a voltage between the n-doped layer at the bottom of the sample and a Cr contact grown on top of the SPS, electrons are attracted from the n-doped layer into the InAs dots. Infra-red photons were used to excite electrons from the occupied $`e_0`$ level into the $`e_1`$ level. Neglecting the small exchange energy, the energy differences for the $`e_1e_0`$ excitations when 1 and 2 electrons are present in the dot are
$`E_{01}[e_1^1]E_{01}[e_0^1]`$ $`=`$ $`\left(ϵ_{e_1}ϵ_{e_0}\right)`$
$`E_{02}[e_0^1e_1^1]E_{02}[e_0^2]`$ $`=`$ $`\left(ϵ_{e_1}ϵ_{e_0}\right)+\left[J_{e_1,e_0}^{ee}J_{e_0,e_0}^{ee}\right],`$ (18)
The first of these energy differences yields a direct measurement of the $`sp`$ energy spacing, $`\delta _{sp}`$, of 49.1 meV. The second energy difference was measured at 50.1 meV. Drexler et. al.drexler94 also used infra red transmission spectroscopy to measure the energy spacing, $`\delta _{sp}=41`$ meV. Pan et. al.pan98 ; pan98:2 have also performed infra-red absorption measurements on truncated pyramidal dots with a base of 180 Å and height of $``$60 Å. In these experiments, no gate voltage is applied, and therefore the excitations take place from the ground state of the samples, $`E_{00}`$. They observe multiple infra-red absorption peaks between 89 and 103 meV. These could be associated either with the $`s`$-$`p`$ spacing of the electron levels or spacings of the hole states (see below).
Itskevich et. al.itskevich99 perform high power PL measurements of pyramidal dots with a base of 150 Å and a height of 30 Å . This high power excitation is able to simultaneously load multiple excitons into the dots. Due to state filling, these multiple excitons will occupy ground state ($`e_0,h_0`$) and higher ($`e_1,e_2,h_0,h_1,h_2`$) single particle levels. Therefore the PL measurements are able to simultaneously measure recombination between electrons occupying the $`e_0,e_1,e_2`$ and $`e_3`$ levels with holes in the $`h_0,h_1,h_2`$ and $`h_3`$ levels. In general, to describe the total energy differences associated with decay from $`N`$ to $`N1`$ excitons occupying a dot requires a treatment that includes the exchange and correlation between multiple occupational configurations of the $`N`$ and $`N1`$ excitonic states. Such a configurational interaction (CI) approach has previously been considered for model parabolic dotshawrylak and will be discussed for realistic dots in a future publicationwilliamson2000:2 . For the purposes of this discussion, we limit ourselves to discussing the energy differences associated with the lowest energy configurations on the $`N`$ exciton state, i.e. those predicted by the Aufbau principle. Within this approximation, the peaks in the high power, PL spectra can be interpreted as corresponding to \[see Eq.(III.4) where exchange is neglected\]
Peak 1: $`E_{11}[h_0^1e_0^1]E_{00}=\left(ϵ_{e_0}ϵ_{h_0}\right)J_{e_0,h_0}^{eh}`$
Peak 2: $`E_{33}[h_0^2h_1^1e_0^2e_1^1]E_{22}[h_0^2e_0^2]=\left(ϵ_{e_1}ϵ_{h_1}\right)J_{e_1,h_1}^{eh}`$
$`+2\left[J_{e_1,h_0}^{eh}J_{e_0,h_1}^{eh}+J_{e_1,e_0}^{ee}+J_{h_1,h_0}^{hh}\right]`$
Peak 3: $`E_{77}[h_0^2h_1^2h_2^2h_3^1e_0^2e_1^2e_2^2e_3^1]E_{66}[h_0^2h_1^2h_2^2e_0^2e_1^2e_2^2]`$ (19)
$`=\left(ϵ_{e_3}ϵ_{h_3}\right)J_{e_3,h_3}^{eh}`$
$`+2{\displaystyle \underset{i=0}{\overset{2}{}}}\left[J_{e_3,h_i}^{eh}J_{e_i,h_3}^{eh}+J_{e_3,e_i}^{ee}+J_{h_3,h_i}^{hh}\right].`$
Note, peak 3 is not assigned to a recombination from $`e_2`$ to $`h_2`$ as this is almost degenerate with peak 2. Itskevich et. al. assume that (i) the Coulomb integrals in the square brackets on the right hand side of Eq.(VI.1) cancel, (ii) that $`J_{e_0,h_0}^{eh}=J_{e_1,h_1}^{eh}=J_{e_3,h_3}^{eh}`$, and (iii) that the hole spacings, $`\delta _{h_0h_1},\delta _{h_1h_2}`$, are small compared to the electron level spacings. With these assumptions the spacings between peaks 1 and 2 and peaks 2 and 3 can be assigned to the $`s`$-$`p`$ and $`p`$-$`d`$ energy spacings. They find spacings $`\delta _{sp}`$ and $`\delta _{pd}`$ of 75 and 80 meV respectively. Our calculations suggest that assumptions (i), (ii) and (iii) probably introduce errors of $`\pm `$10, $`\pm `$5 and +10 meV respectively. The neglect of exchange interactions in the above discussion also introduces an error of $`\pm `$5 meV.
### VI.2 The intra-band electron $`p`$ level splitting
For the lens shaped dots, the Capacitance Voltage spectroscopy of Fricke et. al.fricke96 can be used to estimate the splitting of the $`p`$ states, $`\delta _{pp}`$, by loading two electrons into the dot and exciting them using FIR spectroscopy. The relevant energy differences are
$`E_{02}[e_0^1e_1^1]E_{02}[e_0^2]`$ $`=`$ $`\left(ϵ_{e_1}ϵ_{e_0}\right)+\left[J_{e_1,e_0}^{ee}J_{e_0,e_0}^{ee}\right]`$
$`E_{02}[e_0^1e_2^1]E_{02}[e_0^2]`$ $`=`$ $`\left(ϵ_{e_2}ϵ_{e_0}\right)+\left[J_{e_2,e_0}^{ee}J_{e_0,e_0}^{ee}\right].`$ (20)
By assuming $`J_{e_1,e_0}^{ee}=J_{e_2,e_0}^{ee}`$, the difference in the two above expressions yields the energy spacing $`ϵ_{e_2}ϵ_{e_1}`$. They find a value of $``$2 meV. To measure the effect of a magnetic field on the splitting of the $`p`$ states Fricke et. al.fricke96 use infra-red absorption to measure the above energy differences in an applied magnetic field. At a field of 15 Tesla they measure an energy spacing of 19 meV. More theoretically, Dekel et. aldekel98 have demonstrated that one must assume a splitting of the $`p`$ states to explain the number of multi-exciton levels observed in their single dot, micro PL measurements.
### VI.3 The intra-band hole energy spacings
There are currently no measurements available for the energy spacings between the hole states in the lens shaped dots. Itskevich et. al.itskevich99 have performed high power PL measurements of dots estimated to be square based, truncated pyramids with a base of 150Å and a height of 30Å , under hydrostatic pressure to estimate hole level spacings. At hydrostatic pressures above 55 kbar, they measure the PL associated with transitions from the $`X_{1c}`$-state in the GaAs matrix to the $`h_0`$ and $`h_1`$ levels in the dots. They estimate a hole level spacing of $``$15 meV. However, the nature of the initial and final hole states is unclear.
Sauvage et. al.sauvage99 use polarized photoinduced intraband absorption spectroscopy to measure the energy spacing between the lowest hole state, $`h_{000}`$, and the hole state with a single node in the growth direction, $`h_{001}`$. This corresponds to \[see Eq.(III.4)\]
$`E_{11}[h_{000}^1e_0^1]E_{11}[h_{001}^1e_0^1]`$ $`=`$ $`\left(ϵ_{h_{001}}ϵ_{h_{000}}\right)`$ (21)
$`+`$ $`\left[J_{e_0,h_{000}}^{eh}J_{e_0,h_{001}}^{eh}\right].`$
By assuming that $`J_{e_0,h_{000}}^{eh}=J_{e_0,h_{001}}^{eh}`$ Sauvage et. al. estimate the $`h_{001}h_{000}`$ spacing to be $``$120 meV. Note $`h_{001}`$ is almost certainly higher in energy than states with nodes perpendicular to the growth direction ($`h_{010}`$ and $`h_{100}`$) due to the smaller dimension of the dot in the growth direction. Consequently, the energy difference $`h_{001}h_{000}`$ is not the spacing of the first two hole states $`h_0`$-$`h_1`$.
Tang et. al.tang98 measure activation energies for excitations from $`h_0`$ and $`h_1`$ to the hole wetting layer of 48 and 30 meV respectively, implying an $`h_0`$-$`h_1`$ spacing of $``$18 meV.
### VI.4 The electron and hole binding energies, $`\mathrm{\Delta }E(e)`$ and $`\mathrm{\Delta }E(h)`$
There have been no direct measurements of the electron or hole binding energy for lens shaped InAs dots. However, it has been measured in other dots by several groups using a range of techniques. Berryman et. al.berryman97 placed pyramidal InAs dots estimated to have a base of 100 Å and height of 15 Å in a p-n junction and measured the temperature dependence of the ac conductance as a function of frequency. These measurements predict a hole binding, $`\mathrm{\Delta }E(h)`$, energy of $``$240 meV. When subtracted from the bulk GaAs band gap, this yields a value for the electron binding energy, $`\mathrm{\Delta }E(e)`$, of $``$60 meV. The authors obtain similar results from temperature dependent Hall measurements of thermal hole trapping. Itskevich et. al.itskevich98 measured the pressure at which PL measurements could detect a $`\mathrm{\Gamma }`$-X crossing in pyramidal InAs/GaAs quantum dot samples. By extracting these PL measurements back to zero pressure they were able to extrapolate a value for the electron binding energy, $`\mathrm{\Delta }E(e)`$, of $``$50 meV. Itskevich et. al.itskevich99 also used high pressure PL to measure the energy difference between the $`X_{1c}`$ level in bulk GaAs and the $`h_0`$ level in the quantum dots. By extrapolating this value back to zero pressure, they predict a value for the hole binding energy, $`\mathrm{\Delta }E(h)`$, of $``$250 meV. Brunkov et. al.brunkov98 performed capacitance-voltage spectroscopy measurements, which when fitted to a capacitance model for their dot geometry predict an electron binding energy of 80 meV for dots with bases of 250 Å. Tang et. al.tang98 measured the spacing of the electron wetting layer to both the GaAs CBM and the $`e_0`$ level and hence deduced a value for the electron binding energy, $`\mathrm{\Delta }E(e)`$, of $``$80 meV, for dots with an estimated base of 130 to 170 Å.
### VI.5 The position of the electron and hole wetting layer level
The presence of a distinct wetting layer signal in the PL spectra of a sample of self assembled quantum dots is the hallmark of a high quality sample. In samples where the wetting layer has “dissolved” due to the growth conditions, it is likely that the geometry and composition of the quantum dots will also have dramatically altered from their uncapped state.
In the lens shaped InAs dots Schmidt et. al.schmidt97 observe PL emission from the ground state of the wetting layer at 1.34 eV. Photovoltage measurementsschmidt97 on the same samples show a strong peak corresponding to absorption into the ground state of the wetting layer at 1.35 eV. There are currently no measurements available for the position of the individual electron and hole wetting layers in the lens shaped InAs/GaAs quantum dot samples. In Ref.sauvage97 Sauvage et. al. grow lens shaped InAs dots with an estimated base of 150 Å and a height of 30 Å on a substrate that is n-doped with silicon. This n-doping loads electrons into the $`e_0`$ state in the dot, which they excite into the wetting layer using infra-red excitation. In these samples they estimate the wetting layer to be 150 meV above the $`e_0`$ level. In Ref. sauvage99 Sauvage et. al. load electrons into the $`e_0`$ state of similar dots using an optical interband pump. Subsequent infra-red absorption places the wetting layer 190 meV above the $`e_0`$ state. Tang et. al.tang98 measure thermal transfer of holes to the wetting layer, and obtain a spacing between the $`h_0`$ level and the hole wetting layer, $`\mathrm{\Delta }E_{WL}^{(e)}`$ of $``$48 meV. They also measure thermal transfer from an excited state, possibly involving $`h_1`$, which places the hole wetting layer $``$30 meV below the $`h_1`$ level.
### VI.6 The number of confined electron and hole states
The actual number of confined electron and hole states in a self assembled InAs/GaAs quantum dot depends on the size and composition of the dot. Early single band, effective mass calculationsgrundman95 for pure InAs pyramidal dots with a base of 120 Å and height 60 Å predicted only a single bound electron state and several bound hole states. Consequently, many experiments were then interpreted in this light. More accurate multi-band $`k.p`$stier99 ; jiang97 ; jaros96 and pseudopotentialjkim98 ; williamson99 calculations have predicted 5 or more bound electron states in the same dots.
The high power PL experiments of Itskevich et. al.itskevich99 show the gradual disappearance of 5 peaks as a function of external pressure. This is interpreted as direct evidence for 5 confined electron levels in their samples. The single dot, multi-exciton measurements of Dekel et. al.dekel98 require the assumption of at least 3 bound electron states to explain their experimental spectra. Similarly, the capacitance-voltage spectroscopy of Fricke et. al.fricke96 shows two peaks corresponding to the capacitance of $`s`$-like states and the nearly degenerate $`p`$-like states, providing evidence for at least 3 bound electron states.
### VI.7 Electron and hole Coulomb and Exchange Interactions
By loading multiple electrons and holes into quantum dots it is possible to measure the Coulomb and exchange interactions between these additional electrons and holes. The magnitude of these interactions is a function of the shape of the electronic wavefunctions (see section II) and provides an additional quantity to test the accuracy of theoretical models.
To study electron-electron interactions, Fricke et. al.fricke96 use the same experimental setup discussed in section VI.1. From Eq.(III.4) we see the energy differences corresponding to the peaks in the Capacitance Voltage (CV) spectra associated with loading one and two electrons into the $`e_0`$ level in the dots is
$`E_{01}[e_0^1]E_{00}`$ $`=`$ $`ϵ_{e_0}`$
$`E_{02}[e_0^2]E_{01}[e_0^1]`$ $`=`$ $`ϵ_{e_0}+J_{e_0,e_0}^{ee}.`$ (22)
The electron-electron Coulomb interaction, $`J_{e_0,e_0}^{ee}`$, can therefore be directly measured as the splitting of these two CV peaks. They find a value of $`J_{e_0,e_0}^{ee}`$23 meV. From Eq.(VI.1) we see that
$$J_{e_0,e_1}^{ee}=J_{e_0,e_0}^{ee}+(50.149.1)=24\text{meV}.$$
(23)
Finally by fitting 4 equidistant bell curves to the CV spectra corresponding to loading 3,4,5 and 6 electrons into the dots, an approximate value for the charging energy between the $`p`$ states, $`J_{e_1,e_1}^{ee}`$, of $``$18 meV is obtained. From Eq.(III.4) we that the spacing $`E_{04}E_{03}=J_{e_1,e_1}^{ee}+2J_{e_0,e_1}^{ee}`$, while $`E_{06}E_{05}=3J_{e_1,e_1}^{ee}+2J_{e_0,e_1}^{ee}`$ and hence the approximation of equidistant peaks will introduce some error into this estimate for $`J_{e_1,e_1}^{ee}`$.
### VI.8 Electron-hole excitonic recombination
To our knowledge there have so far been no measurements of the polarization anisotropy in the lens shaped dots discussed here. The polarization anisotropy for $`e_ih_i`$ excitonic recombination in InAs/GaAs was measured by Yang et. al.yang99 ; yang2000 , who find a ratio of $`\lambda _{e_0,h_0}=1.2`$ and $`\lambda _{e_2,h_2}=1.3`$ for InAs dots whose geometry is measured to be formed by four {136} faceted planes with bases ranging from 150 to 250 Å and a base to height ratio of 4:1. Yang et. al.yang99 ; yang2000 have performed k.p calculations for this dot geometry, which include the “geometric factor” but not the “atomic symmetry factor” discussed in section IV. They find $`\lambda _{e_0,h_0}=1.8`$ and $`\lambda _{e_2,h_2}=3.5`$. The authors suggest the k.p simulations of the measured polarization ratio can be used to deduce the geometric shape anisotropy. However, we demonstrate here that when the “atomic symmetry factor” is included an anisotropy of $`\lambda _{e_0,h_0}=1.3`$ is obtained even for a square based pyramid. Thus k.p simulations lacking the “atomic symmetry” factor can not be used to reliably deduce the geometric shape anisotropy.
To our knowledge there have so far been no measurements of the excitonic dipole in lens shaped InAs dots. Fry et. al.fry2000 used photocurrent spectroscopy within an applied electric field to measure the excitonic dipole moment, $`d_{h_i,e_j}`$, \[Eq.(17)\]. They find the center of the hole wavefunctions to be located $``$4 Å above the center of the electron wavefunction (positive dipole). Fry et. al.fry2000 also perform single-band, effective mass calculations, in an attempt to isolate the origin of this dipole. They predict that in the absence of alloying the dipole is -3 Å, i.e. the opposite sign , but a linear composition profile with Ga<sub>0.5</sub>In<sub>0.5</sub>As at the base and pure InAs at the top of a truncated pyramid with a base of 155 Å, and height 55 Å, reproduces the correct dipole of 4 Å. They suggest that this alloying profile explains the observed dipole. We have repeated these calculations and confirm that, within a single-band model, such a composition profile causes both electrons and holes to move up in the dot compared to their positions in a pure InAs dot. The heavier effective mass of the holes, results in less kinetic energy associated with confinement at the top of the dot and hence the holes move up more than electrons on the introduction of Ga, producing the correct dipole. However, when we repeat these calculations in the more sophisticated, multi-band LCBB basis used here, we find significant heavy hole-light hole mixing in the $`h_0`$ state, which acts to reduce the above effect and produce a smaller dipole of $``$1 Å, in contradiction with experiment (+4 Å). We therefore conclude that in a more complete calculation the shape and alloy profile postulated by Fry et. al.fry2000 does not produce the observed excitonic dipole.
## VII Comparison of experiment and theory
In Table 5 we show the results of our calculations for pure InAs, lens shaped quantum dots embedded within GaAs \[column (a)\]. Table 5 also shows the experimentally measured splittings of the electron levels, the electron-electron and electron-hole Coulomb energies, the magnetic field dependence and the excitonic band gap measured in Refs.fricke96 and warburton98 . The agreement between the measured energy level spacings, Coulomb energies and magnetic field response with our theoretical lens shaped model is generally good. Both the model and experiment find (i) a large spacing, $`\delta _{sp}`$, ($``$50-60 meV) between the $`s`$-like $`e_0`$ state and the $`p`$-like $`e_1`$ state, (ii) a small spacing, $`\delta _{pp}`$, ($``$3 meV) between the two $`p`$-like $`e_1`$ and $`e_2`$ states and (iii) a large spacing ($``$55 meV) between the $`p`$-like $`e_2`$ state and the $`d`$-like $`e_3`$ state.
These electron level spacings are similar to those found for pyramidal quantum dotsjkim98 (see Table 5). However, in the pyramidal dot, the spacings of the two $`p`$-like and $`d`$-like states, $`\delta _{pp},\delta _{dd}`$, is larger (26 and 23 meV) as a result of the lower C<sub>2v</sub> symmetry of a zincblende pyramid. Both the model and experiment also find similar values for the Coulomb energies, $`J(e_0e_0)`$ and $`J(e_0h_0)`$ ($``$25 meV).
The calculated hole binding energy of $`\mathrm{\Delta }E(h)=193`$ meV is in good agreement with those of Berryman et. al.berryman97 ($``$240 meV) and Itskevich et. al.itskevich99 ($``$250 meV). Our calculated electron binding energies, $`\mathrm{\Delta }E(h)`$, are considerable larger (271 meV) than those of Berryman et. al.berryman97 ($``$60 meV) and Tang et. al.tang98 ($``$80 meV). We attribute this difference to the larger size of our dots. The assumption of a pure InAs dot also affects the comparison. The agreement improves when we include Ga in our dots (see section VII.2).
The calculated electron-electron and electron-hole Coulomb energies are in reasonable agreement with those extracted from Refs.fricke96 and warburton98 . For the integrals $`J_{e_0e_0}^{ee},J_{e_0e_1}^{ee},J_{e_1e_1}^{ee}`$ and $`J_{e_0h_0}^{eh}`$ we calculate values of 31, 25, 25 and 37 respectively, compared to measured values of 23, 24, 18 and 33.3 meV
The calculated polarization anisotropies, $`\lambda `$, for the $`e_0h_0`$ recombination in lens and pyramidal shaped, pure InAs dots are $`\lambda =1.03`$ and 1.2 respectively. A future measurement of this anisotropy ratio for lens shaped dots would provide an important piece of evidence for determining the detailed geometry of the dots.
In the lens shaped dot we find a difference in the average positions of the $`h_0`$ and $`e_0`$ states, $`d_{h_i,e_j}`$, of around 1 Å. This is smaller than the value we calculate for a pyramidal quantum dot, where we find the hole approximately 3.1 Å higher than the electron.
In summary, the assumed lens shaped geometry, with a pure InAs composition produces a good agreement with measured level splitting, Coulomb energies and magnetic field dependence. A closer inspection of the remaining differences reveals that the calculations systematically overestimate the splittings between the single particle electron levels ($`\delta _{sp}`$: 65 vs. 50 meV, $`\delta _{pd}`$:68 vs. 48 meV) and underestimates the excitonic band gap (1032 vs. 1098 meV).
### VII.1 Pure InAs dots: The effects of shape and size
Focusing on the lens shape only, we examine the effect of changing the height and base of the assumed geometry. Calculations were performed on similar lens shaped, pure InAs dots where (i) the base of the dot was increased from 252 to 275Å, while keeping the height fixed at 35Å, \[column (b)\] and (ii) the height of the dot was decreased from 35 to 25Å, while keeping the base fixed at 252Å, \[column (c)\]. It shows that decreasing the height of the dot increases the quantum confinement and hence increases the splittings of the electron and hole levels ($`\delta _{sp}`$: from 65 to 69 meV and $`\delta _{h_0,h_1}`$: from 8 to 16 meV). Decreasing the height of the dot also acts to increase the excitonic band gap from 1032 to 1131 meV by pushing up the energy of the electron levels and pushing down the hole levels. Conversely, increasing the base of the dot decreases both the splittings of the single particle levels ($`\delta _{sp}`$: from 66 to 61 meV) and the band gap (1032 to 1016 eV). These small changes in the geometry of the lens shaped dot have only a small effect on electronic properties that depend on the shape of the wavefunctions. The electron-electron and electron-hole Coulomb energies remain relatively unchanged, the magnetic field induced splitting remain at 20 meV, the polarization anisotropy, $`\lambda `$, remains close to 1.0 and the excitonic dipole, $`d_{h_i,e_j}`$, remains negligible. In summary, reducing either the height or the base of the dot increases quantum confinement effects and hence increases energy spacings and band gaps, while not significantly effecting the shape wavefunctions.
### VII.2 Interdiffused In(Ga)As/GaAs lens shaped dots
We next investigate the effect of changing the composition of the quantum dots, while keeping the geometry fixed. There have recently been several experimentsfry2000 ; metzger99 ; garcia97 suggesting that a significant amount of Ga diffuses into the nominally pure InAs quantum dots during the growth process. We investigate two possible mechanisms for this Ga in-diffusion; (i) Ga diffuses into the dots during the growth process from all directions producing a dot with a uniform Ga composition Ga<sub>x</sub>In<sub>1-x</sub>As, and (ii) Ga diffuses up from the substrate, as suggested in Ref.fry2000 . To investigate the effects of these two methods of Ga in-diffusion on the electronic structure of the dots, we compare pure InAs dots embedded in GaAs with Ga<sub>x</sub>In<sub>1-x</sub>As, random alloy dots embedded in GaAs, where the Ga composition, $`x`$, (i) is fixed at 0.15, \[column (d)\] and (ii) varies linearly from 0.3 at the base to 0 at the top of the dot, \[column (e)\].
The electronic structure of these dots is compared in Table 5. It shows that increasing the amount of Ga in the dots acts to decrease the electron level spacings ($`\delta _{sp}`$: from 65 to 58 and 64 meV for $`x=0.15`$ and $`x=0.3`$ to $`x=0`$ respectively). It also acts to increase the excitonic band gap from 1032 to 1080 and 1125 meV respectively. The electron binding energy, $`\mathrm{\Delta }E(e)`$, is decreased by the in diffusion of Ga (from 271 to 209 and 192 meV), while the hole binding energy, $`\mathrm{\Delta }E(h)`$, is relatively unaffected. This significant decrease in the electron binding energy considerably improves the agreement with experiments on other dot geometriesberryman97 ; tang98 .
As with changing the size of the dots, we find that Ga in-diffusion has only a small effects on properties that depend on the shape of the wavefunctions. The calculated electron-electron and electron-hole Coulomb energies are almost unchanged, while the average separation of the electron and hole, $`d_{h_i,e_j}`$, increases from 0.16 to to 0.5 and 1.2 Å and the polarization ratio, $`\lambda `$, and magnetic field response are also unchanged.
Table 5 shows that the dominant contribution to the increase in the excitonic band gap and reduction in electron binding energy, results mostly from an increase in the energy of the electron levels as the Ga composition is increased. This can be understood by considering the electronic properties of the bulk Ga<sub>x</sub>In<sub>1-x</sub>As random alloy. The unstrained valence band offset between GaAs and InAs is $``$ 50 meVwei98 , while the conduction band offset in $``$ 1100 meV and hence changing the Ga composition, $`x`$, has a large effect on the energy of the electron states and only a small effect on the hole states. In summary, the effect of Ga in-diffusion is to reduce the spacing of the electron levels while significantly increasing their energy and hence increasing the band gap. We find that only the average Ga composition in the dots is important to their electronic properties. Whether this Ga is uniformly or linearly distributed throughout the dots has a negligible effect. Note, in Ref.fry2000 it is suggested that a linear composition profile is required to produce an excitonic dipole moment in agreement with that measured by stark experiments. For the lens shaped geometry discussed here there have so far been no such measurements of the dipole, but our calculations suggest that it should be small ($``$1Å).
## VIII Discussion
The effects of changing the geometry of the lens shaped, pure InAs dots on the single particle energy levels can be qualitatively understood from single band, effective mass arguments. These predict that decreasing any dimension of the dot, increases the quantum confinement and hence the energy level spacings and the single particle band gap will increase. Note that as the dominant quantum confinement in these systems arises from the vertical confinement of the electron and hole wavefunctions, changing the height has a stronger effect of the energy levels than changing the base. In this case decreasing the height by 10Å has a much stronger effect on the energy spacings and on the band gap than increasing the base by 23Å.
As increasing(decreasing) the dimensions of the dot acts to decrease(increase) both the level spacings and the gap, it is clear that changing the dot geometry alone will not significantly improve the agreement with experiment as this requires a simultaneous decrease in the energy level splittings and increase in the band gap. However, Ga in-diffusion into the dots acts to increase the band gap of the dot while decreasing the energy level spacings. Table 5 shows that adopting a geometry with a base of 275 Å and a height of 35 Åand a uniform Ga composition of Ga<sub>0.15</sub>In<sub>0.85</sub>As produces the best fit to the measurements in Refs.fricke96 and warburton98 .
In conclusion, our results strongly suggest that to obtain very accurate agreement between theoretical models and experimental measurements for lens shaped quantum dots, one needs to adopt a model of the quantum dot that includes some Ga in-diffusion within the quantum dot. When Ga in-diffusion is included, we obtain an excellent agreement between state of the art multi-band pseudopotential calculations and experiments for a wide range of electronic properties. We are able to predict most observable properties to an accuracy of $`\pm 5`$ meV, which is sufficient to make predictions of both the geometry and composition of the dot samples.
Acknowledgments We thank J. Shumway and A. Franceschetti for many useful discussions and their comments on the manuscript. This work was supported DOE – Basic Energy Sciences, Division of Materials Science under contract No. DE-AC36-99GO10337.
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# Multivariate hypergeometric functions as tau functions of Toda lattice and Kadomtsev-Petviashvili equation
## 1 Introduction
Hypergeometric functions play an important role both in physics and in mathematics . Many special functions and polynomials (such as q-Askey-Wilson polynomials, q-Jacobi polynomials, q-Racah polynomials, q-Hahn polynomials, expressions for Clebsch-Gordan coefficients) are just certain hypergeometric functions evaluated in special values of parameters. In physics hypergeometric functions and their q-deformed counterparts sometimes play the role of wave functions and correlation functions for quantum integrable systems. In the present paper we shall construct hypergeometric functions as tau-functions of the Kadomtsev-Petviashvili (KP) hierarchy of equations. It is interesting that the KP equation, which originally serves in plasma physics, now plays a very important role in physics (for modern applications see review ) and in mathematics. This peculiarity of Kadomtsev-Petviashvili equation appeared in the paper , where L-A pair of the KP equation was presented, and mainly in the paper of V.E.Zakharov and A.B.Shabat in 1974 where this equation was integrated by the dressing method. Actually it was the paper where so-called hierarchy of higher KP equations appeared. Another very important equation is the two-dimensional Toda lattice integrated first in . In the present paper we use these equations to construct hypergeometric functions which depend on many variables, these variables are KP and Toda lattice higher times. Here we shall use the general approach to integrable hierarchies of Kyoto school . Especially a set of papers about Toda lattice is important for us.
We briefly outline the connection of what we do with Zakharov-Shabat dressing method and with the nonlocal $`\overline{}`$ problem ; we mention the related system of orthogonal polynomials. The lack of space does not allow us to develop these topics.
We devote this paper to Vladimir Evgen’evich Zakharov on his 60 birthday.
## 2 Notations
There are several well-known different multivariate generalizations of hypergeometric series of one variable . Let $`\left|q\right|<1`$ and let $`𝐱_{\left(N\right)}=(x_1,\mathrm{},x_N)`$ be indeterminates. The multiple basic hypergeometric series is
$`{}_{p}{}^{}\mathrm{\Phi }_{s}^{}\left({\displaystyle \genfrac{}{}{0pt}{}{a_1,\mathrm{},a_p}{b_1,\mathrm{},b_s}}|q,𝐱_{\left(N\right)}\right)={\displaystyle \underset{l\left(𝐧\right)N}{}}{\displaystyle \frac{(q^{a_1};q)_𝐧\mathrm{}(q^{a_p};q)_𝐧}{(q^{b_1};q)_𝐧\mathrm{}(q^{b_s};q)_𝐧}}{\displaystyle \frac{q^{n\left(𝐧\right)}}{H_𝐧\left(q\right)}}s_𝐧\left(𝐱_{\left(N\right)}\right).`$ (1)
The sum is over all different partitions $`𝐧=(n_1,n_2,\mathrm{},n_r)`$, where $`n_1n_2\mathrm{}n_r`$, $`r\left|𝐧\right|`$, $`\left|𝐧\right|=n_1+\mathrm{}+n_r`$ and whose length $`l\left(𝐧\right)=rN`$. Schur polynomial $`s_𝐧\left(𝐱_{\left(N\right)}\right)`$, with $`Nl\left(𝐧\right)`$, is a symmetric function of variables $`𝐱_{\left(N\right)}`$ and defined as follows :
$$s_𝐧\left(𝐱_{\left(N\right)}\right)=\frac{a_{𝐧+\delta }}{a_\delta },a_𝐧=det\left(x_i^{n_j}\right)_{1i,jN},\delta =(N1,N2,\mathrm{},1,0).$$
(2)
Schur function $`s_𝐧\left(𝐱_{\left(N\right)}\right)=0`$ for $`N<l\left(𝐧\right)`$. In the theory of the KP hierarchy it is convenient to define Schur functions in terms of KP higher times $`𝐭=(t_1,t_2,\mathrm{})`$ as
$$s_𝐧\left(𝐭\right)=det\left(p_{n_ii+j}\left(𝐭\right)\right)_{1i,jr},\underset{m=0}{\overset{+\mathrm{}}{}}p_m\left(𝐭\right)z^m=\mathrm{exp}\left(\underset{i=1}{\overset{+\mathrm{}}{}}t_iz^i\right)=e^{\xi (𝐭,z)}.$$
(3)
The functions $`s_𝐧\left(𝐭\right)`$ are related to $`s_𝐧\left(𝐱_{\left(N\right)}\right)`$ via the Miwa’s change of variables:
$$t_m=\underset{i=1}{\overset{N}{}}\frac{x_i^m}{m}.$$
(4)
Each coefficient $`(q^a;q)_𝐧`$ in (1) is expressed in terms of the so-called q-deformed Pochhammer symbols $`(q^a;q)_{n_i}`$:
$`(q^a;q)_𝐧=(q^a;q)_{n_1}(q^{a1};q)_{n_2}\mathrm{}(q^{ar+1};q)_{n_r}`$ (5)
$`(q^a;q)_{n_i}=\left(1q^a\right)\left(1q^{a+1}\right)\mathrm{}\left(1q^{a+n_i1}\right),(q^a;q)_0=1`$ (6)
$`q`$-deformed ’hook polynomials’ $`H_𝐧\left(q\right)`$ are
$`H_𝐧\left(q\right)={\displaystyle \underset{(i,j)𝐧}{}}\left(1q^{h_{ij}}\right),h_{ij}=\left(n_i+n_j^{}ij+1\right),`$ (7)
where $`𝐧^{}=\left(n_{}^{}{}_{1}{}^{}+n_{}^{}{}_{2}{}^{}+\mathrm{}+n_{}^{}{}_{r^{}}{}^{}\right)`$ is the conjugated partition and $`q^{n\left(𝐧\right)}=q^{_{i=1}^r\left(i1\right)n_i}`$.
The formula
$${}_{p}{}^{}\mathrm{\Phi }_{s}^{}\left(\genfrac{}{}{0pt}{}{a_1,\mathrm{},a_p}{b_1,\mathrm{},b_s}|q,𝐱_{\left(N\right)},𝐲_{\left(N\right)}\right)=$$
$$\underset{l\left(𝐧\right)N}{}\frac{(q^{a_1};q)_𝐧\mathrm{}(q^{a_p};q)_𝐧}{(q^{b_1};q)_𝐧\mathrm{}(q^{b_s};q)_𝐧}\frac{q^{n\left(𝐧\right)}}{H_𝐧\left(q\right)}\frac{s_𝐧\left(𝐱_{\left(N\right)}\right)s_𝐧\left(𝐲_{\left(N\right)}\right)}{s_𝐧(1,q,q^2,\mathrm{},q^{N1})}$$
(8)
defines the multiple basic hypergeometric function of two sets of variables .
Another generalization of hypergeometric series is so-called hypergeometric function of matrix argument $`𝐗`$ with indices $`𝐚`$ and $`𝐛`$ :
$`{}_{p}{}^{}F_{s}^{}\left({\displaystyle \genfrac{}{}{0pt}{}{a_1,\mathrm{},a_p}{b_1,\mathrm{},b_s}}|𝐗\right)={\displaystyle \underset{𝐧}{}}{\displaystyle \frac{\left(a_1\right)_𝐧\mathrm{}\left(a_p\right)_𝐧}{\left(b_1\right)_𝐧\mathrm{}\left(b_s\right)_𝐧}}{\displaystyle \frac{Z_𝐧\left(𝐗\right)}{\left|𝐧\right|!}}.`$ (9)
Here $`𝐗`$ is $`N\times N`$ matrix, and $`Z_𝐧\left(𝐗\right)`$ is zonal spherical polynomial for the symmetric space $`GL(N,C)/U\left(N\right)`$, see ,. Let us note that in the limit $`q1`$ series (1) coincides with (9), see .
Now let us review some facts from the KP theory . We have fermionic fields $`\psi \left(z\right)=_{kZ}\psi _kz^k`$ and $`\psi ^{}\left(z\right)=_{kZ}\psi _k^{}z^{k1}dz`$, where fermionic operators satisfy the canonical anti-commutation relations:
$$[\psi _m,\psi _n]_+=[\psi _m^{},\psi _n^{}]_+=0;[\psi _m,\psi _n^{}]_+=\delta _{mn}.$$
(10)
Let us introduce left and right vacuums by the properties:
$`\psi _m|0=0\left(m<0\right),\psi _m^{}|0=0\left(m0\right),`$ (11)
$`0|\psi _m=0\left(m0\right),0|\psi _m^{}=0\left(m<0\right).`$ (12)
Throughout the text the subscript $``$ does not denote the complex conjugation. The vacuum expectation value is defined by relations:
$`0\left|1\right|0=1,0\left|\psi _m\psi _m^{}\right|0=1m<0,0\left|\psi _m^{}\psi _m\right|0=1m0,`$ (13)
$$0\left|\psi _m\psi _n\right|0=0\left|\psi _m^{}\psi _n^{}\right|0=0,0\left|\psi _m\psi _n^{}\right|0=0mn.$$
(14)
Let us notice that if a function $`h`$ has no poles and zeroes for integer value of argument, then relations (10)-(14) are invariant under the following transformation
$$\psi _n\frac{1}{h\left(n\right)}\psi _n,\psi _n^{}h\left(n\right)\psi _n^{}.$$
(15)
Let us denote $`\widehat{gl}\left(\mathrm{}\right)=\mathrm{𝐿𝑖𝑛}\{1,:\psi _i\psi _j^{}:|i,jZ\}`$, with usual normal ordering $`:\psi _i\psi _j^{}:=\psi _i\psi _j^{}0\left|\psi _i\psi _j^{}\right|0`$. We define the operator $`g`$ which is an element of the group $`\widehat{GL}\left(\mathrm{}\right)`$ corresponding to the infinite dimensional Lie algebra $`\widehat{gl}\left(\mathrm{}\right)`$. The tau-function of the KP equation and the tau-function of the two-dimensional Toda lattice (TL) sometimes are defined as
$$\tau _{KP}(M,𝐭)=M\left|e^{H\left(𝐭\right)}g\right|M,MZ,$$
(16)
$$\tau _{TL}(M,𝐭,𝐭^{})=M\left|e^{H\left(𝐭\right)}ge^{H^{}\left(𝐭^{}\right)}\right|M,MZ,$$
(17)
where $`𝐭=(t_1,t_2,\mathrm{})`$ and $`𝐭^{}=(t_1^{},t_2^{},\mathrm{})`$ are called higher Toda lattice times (the first set $`𝐭`$ is in the same time the set of higher KP times). $`H\left(𝐭\right)`$ and $`H^{}\left(𝐭^{}\right)`$ belong to the following $`\widehat{gl}\left(\mathrm{}\right)`$ Cartan subalgebras:
$$H\left(𝐭\right)=\underset{n=1}{\overset{+\mathrm{}}{}}t_nH_n,H^{}\left(𝐭^{}\right)=\underset{n=1}{\overset{+\mathrm{}}{}}t_n^{}H_n,H_n=\frac{1}{2\pi i}:z^n\psi \left(z\right)\psi ^{}\left(z\right):.$$
(18)
According to the integer $`M`$ in (17) plays the role of discrete Toda lattice variable and defines the following charged vacuums
$`M|=0|\mathrm{\Psi }_M^{},|M=\mathrm{\Psi }_M|0,`$ (19)
$`\mathrm{\Psi }_M=\psi _{M1}\mathrm{}\psi _1\psi _0M>0,\mathrm{\Psi }_M=\psi _M^{}\mathrm{}\psi _2^{}\psi _1^{}M<0,`$
$`\mathrm{\Psi }_M^{}=\psi _0^{}\psi _1^{}\mathrm{}\psi _{M1}^{}M>0,\mathrm{\Psi }_M^{}=\psi _1\psi _2\mathrm{}\psi _MM<0.`$ (20)
Lemma 1 For $`j_1<\mathrm{}<j_r<0i_s<\mathrm{}<i_1`$, $`sr0`$ the following formula is valid:
$$sr\left|e^{H\left(𝐭\right)}\psi _{j_1}^{}\mathrm{}\psi _{j_r}^{}\psi _{i_s}\mathrm{}\psi _{i_1}\right|0=\left(1\right)^{j_1+\mathrm{}+j_r+\left(rs\right)\left(rs+1\right)/2}s_𝐧\left(𝐭\right),$$
(21)
where the partition $`𝐧=(n_1,\mathrm{},n_{sr},n_{sr+1},\mathrm{},n_{sr+j_1})`$ is defined by the following pair of partitions:
$`(n_1,\mathrm{},n_{sr})=(i_1\left(sr\right)+1,i_2\left(sr\right)+2,\mathrm{},i_{sr}),`$ (22)
$`(n_{sr+1},\mathrm{},n_{sr+j_1})=(i_{sr+1},\mathrm{},i_s|j_11,\mathrm{},j_r1).`$ (23)
Here $`\left(\mathrm{}|\mathrm{}\right)`$ is another notation for a partition .
So-called vertex operators $`V\left(z\right)`$ and $`V^{}\left(z\right)`$ are defined by:
$$V\left(z\right)=z^Me^{\xi (𝐭,z)}e^{\xi (\stackrel{~}{},z^1)},V^{}\left(z\right)=z^Me^{\xi (𝐭,z)}e^{\xi (\stackrel{~}{},z^1)},$$
(24)
where $`\stackrel{~}{}=(\frac{}{t_1},\frac{1}{2}\frac{}{t_2},\mathrm{})`$ and the function $`\xi (𝐭,z)`$ is the same as in (3). The Baker-Akhiezer function and the conjugated one are
$$w(M,𝐭;z)=\frac{V\left(z\right)\tau }{\tau }w^{}(M,𝐭;z)=\frac{V^{}\left(z\right)\tau }{\tau }.$$
(25)
The function $`u=2_{t_1}^2\mathrm{log}\tau `$ solves the celebrated KP equation
$$u_{t_3}=\frac{1}{4}u_{t_1t_1t_1}+\frac{3}{2}uu_{t_1}+\frac{3}{4}^{t_1}u_{t_2t_2}𝑑t_1^{}.$$
(26)
It is well-known that the Schur functions $`s_𝐧\left(𝐭\right)`$ are tau-functions for some rational solutions of the KP hierarchy. It is also known that not any linear combination of Schur functions turns out to be a KP tau-function. In order to find these combinations one should solve bilinear difference equation (see ), a version of discrete Hirota equation. Below we shall present KP tau-functions which are infinite hypergeometric series of Schur polynomials (1),(9). We shall use the fermionic representations of tau-functions .
## 3 Hypergeometric tau-functions
### 3.1 Additional symmetries and tau-function
Let $`r`$ be a function of one variable. We shall assume that $`r\left(n\right)`$ is finite for $`nZ`$. Let $`D=z\frac{d}{dz}`$ acts on the space of functions $`\left\{z^n;nZ\right\}`$ . Then $`r\left(D\right)z^n=r\left(n\right)z^n`$. All functions of operator $`D`$ which we consider below are given via their eigenvalues on this basis.
Let us consider an abelian subalgebra in $`\widehat{gl}\left(\mathrm{}\right)`$ formed by the following set of fermionic operators
$$A_k=\frac{1}{2\pi i}\psi ^{}\left(z\right)\left(\frac{1}{z}r\left(D\right)\right)^k\psi \left(z\right),k=1,2,\mathrm{},$$
(27)
where the operator $`r\left(D\right)`$ acts on all functions of $`z`$ from the right hand side. For the collection of independent variables $`\beta =(\beta _1,\beta _2,\mathrm{})`$ we denote
$$A\left(\beta \right)=\underset{n=1}{\overset{\mathrm{}}{}}\beta _nA_n.$$
(28)
For the partition $`𝐧=(n_1,\mathrm{},n_k)`$ and a function of one variable $`r`$, let us introduce the following notation
$$r_𝐧\left(M\right)=\underset{i=1}{\overset{k}{}}r\left(1i+M\right)r\left(2i+M\right)\mathrm{}r\left(n_ii+M\right).$$
(29)
We set $`r_\mathrm{𝟎}\left(M\right)=1`$. Using the notation from (21) we have
Lemma 2 The following formula holds
$$0\left|\psi _{i_1}^{}\mathrm{}\psi _{i_s}^{}\psi _{j_s}\mathrm{}\psi _{j_1}e^{A\left(\beta \right)}\right|0=\left(1\right)^{j_1+\mathrm{}+j_s}r_𝐧\left(0\right)s_𝐧\left(\beta \right).$$
(30)
Let us consider the following tau-function of the KP hierarchy
$$\tau _r(M,𝐭,\beta ):=M\left|e^{H\left(𝐭\right)}e^{A\left(\beta \right)}\right|M.$$
(31)
Proposition 1 We have the following expansion:
$$\tau _r(M,𝐭,\beta )=\underset{n=0}{\overset{+\mathrm{}}{}}\underset{\left|𝐧\right|=n}{}r_𝐧\left(M\right)s_𝐧\left(𝐭\right)s_𝐧\left(\beta \right).$$
(32)
We shall not consider the problem of convergence of this series. The variables $`M,𝐭`$ play the role of KP higher times, $`\beta `$ is a collection of group times for a commuting subalgebra of additional symmetries of KP (see and Remark 7 in ). From different point of view (32) is a tau-function of two-dimensional Toda lattice with two sets of continuous variables $`𝐭`$, $`\beta `$ and one discrete variable $`M`$. Formula (32) is symmetric with respect to $`𝐭\beta `$. This ’duality’ supplies us with the string equation which characterizes a tau-function of hypergeometric type (see below). In the similar expansions to (32) were considered, without specifying the coefficients and in a different context.
Let us introduce
$$\stackrel{~}{A}_k=\frac{1}{2\pi i}\psi ^{}\left(z\right)\left(\stackrel{~}{r}\left(D\right)z\right)^k\psi \left(z\right),\left(k=1,2,\mathrm{}\right),\stackrel{~}{A}\left(\beta \right)=\underset{n=1}{\overset{\mathrm{}}{}}\stackrel{~}{\beta }_n\stackrel{~}{A}_n.$$
(33)
Then we have the following generalization of Proposition 1:
Proposition 2
$$M\left|e^{\stackrel{~}{A}\left(\stackrel{~}{\beta }\right)}e^{A\left(\beta \right)}\right|M=\underset{𝐧}{}\left(\stackrel{~}{r}r\right)_𝐧\left(M\right)s_𝐧\left(\stackrel{~}{\beta }\right)s_𝐧\left(\beta \right).$$
(34)
For an interpretation of this expansion see Remark 2 below. Let us mark that the function $`r\left(M\right)`$ is connected with the Sato $`b`$-function .
Let us consider the following difference equations on functions $`\stackrel{~}{h}\left(D\right),h\left(D\right)`$:
$$\stackrel{~}{h}\left(D\right)\stackrel{~}{r}\left(D\right)\stackrel{~}{h}^1\left(D1\right)=1,h^1\left(D1\right)r\left(D\right)h\left(D\right)=1.$$
(35)
The similar equation appeared in the paper of Graev (see formula (6) in ) and was used for generating of different hypergeometric series. As in , in terms of the operator $`r\left(D\right)`$, it is possible to construct the differential equations (and q-difference equations) for hypergeometric functions. We shall present some examples below in (50) and (54).
Our fermionic representation is equivalent to that used in Toda lattice theory if $`r\left(n\right)0,nZ`$. We define a Hamiltonian $`H_0\left(h\right)\widehat{gl}\left(\mathrm{}\right)`$ and a set of $`C_n,nZ`$:
$$H_0\left(h\right):=\frac{1}{2\pi i}:\psi ^{}\left(z\right)\mathrm{log}\left(h\left(D\right)\right)\psi \left(z\right):$$
(36)
$$C_n=\frac{1}{h\left(n1\right)}\mathrm{}\frac{1}{h\left(1\right)}\frac{1}{h\left(0\right)}\frac{1}{\stackrel{~}{h}\left(n1\right)}\mathrm{}\frac{1}{\stackrel{~}{h}\left(1\right)}\frac{1}{\stackrel{~}{h}\left(0\right)},n>0,$$
(37)
$$C_n=h\left(n\right)\mathrm{}h\left(2\right)h\left(1\right)\stackrel{~}{h}\left(n\right)\mathrm{}\stackrel{~}{h}\left(2\right)\stackrel{~}{h}\left(1\right),n<0.$$
(38)
Proposition 3 If function $`r`$ has no zeroes at integer values of argument then
$`\tau (n,\stackrel{~}{\beta },\beta ):=n\left|e^{\stackrel{~}{A}\left(\stackrel{~}{\beta }\right)}e^{A\left(\beta \right)}\right|n=n\left|e^{H\left(\stackrel{~}{\beta }\right)}ge^{H^{}\left(\beta \right)}\right|nC_n^1,`$ (39)
$`g=e^{H_0\left(\stackrel{~}{h}h\right)}=e^{H_0\left(\stackrel{~}{h}\right)+H_0\left(h\right)},C_n=n\left|g\right|n.`$ (40)
For $`\stackrel{~}{r}=1`$ we can put $`\stackrel{~}{\beta }=𝐭`$. Then the following equations hold
$`_{t_1}_{\beta _1}\varphi _n=r\left(n\right)e^{\varphi _{n1}\varphi _n}r\left(n+1\right)e^{\varphi _n\varphi _{n+1}},e^{\varphi _n}={\displaystyle \frac{\tau (n+1,𝐭,\beta )}{\tau (n,𝐭,\beta )}},`$ (41)
$`\tau \left(n\right)_{\beta _1}_{t_1}\tau \left(n\right)_{t_1}\tau \left(n\right)_{\beta _1}\tau \left(n\right)=r\left(n\right)\tau \left(n1\right)\tau \left(n+1\right).`$ (42)
If the function $`r`$ has no integer zeroes, then after the change $`\phi _n=\varphi _n\mathrm{log}h\left(n\right)`$ we obtain Toda lattice equation in standard form :
$$_{t_1}_{\beta _1}\phi _n=e^{\phi _{n+1}\phi _n}e^{\phi _n\phi _{n1}}.$$
(43)
The main point of the paper is based on the observation that if $`r\left(D\right)`$ is a rational function of $`D`$ then $`\tau _r`$ is a hypergeometric series. If $`r\left(D\right)`$ is a rational function of $`q^D`$ we obtain q-deformed hypergeometric series. We shall see both cases in the following examples. (In a separate paper the case of rational expressions of elliptic theta-functions will be considered).
### 3.2 Examples of the tau-functions
Now let us consider various $`r\left(D\right)`$.
Example 1 Let $`r\left(M\right)=M`$ and $`\beta =(\beta _1,0,0,\mathrm{})`$. For $`M=\pm 1`$ we get
$$\tau (1,𝐭,\beta _1)=e^{\xi (𝐭,\beta _1)},\tau (1,𝐭,\beta _1)=e^{\xi (𝐭,\beta _1)}.$$
(44)
Thus $`\beta _1`$ plays the role of a spectral parameter for the vacuum Baker-Akhiezer function. This is in accordance to the meaning of $`\beta _1`$ as a group time for the Galilean transformation .
Example 2 Let all parameters $`b_k`$ be nonintegers.
$${}_{p}{}^{}r_{s}^{}\left(D\right)=\frac{\left(D+a_1\right)\left(D+a_2\right)\mathrm{}\left(D+a_p\right)}{\left(D+b_1\right)\left(D+b_2\right)\mathrm{}\left(D+b_s\right)}.$$
(45)
If all $`a_k`$ are also nonintegers the relevant $`h\left(D\right)`$ is:
$${}_{p}{}^{}h_{s}^{}\left(D\right)=\frac{\mathrm{\Gamma }\left(D+b_1+1\right)\mathrm{\Gamma }\left(D+b_2+1\right)\mathrm{}\mathrm{\Gamma }\left(D+b_s+1\right)}{\mathrm{\Gamma }\left(D+a_1+1\right)\mathrm{\Gamma }\left(D+a_2+1\right)\mathrm{}\mathrm{\Gamma }\left(D+a_p+1\right)}.$$
(46)
For the correlator (32) we have:
$${}_{p}{}^{}\tau _{s}^{}(M,𝐭,\beta )=\underset{n=0}{\overset{+\mathrm{}}{}}\underset{\left|𝐧\right|=n}{}s_𝐧\left(𝐭\right)s_𝐧\left(\beta \right)\frac{\left(a_1+M\right)_𝐧\mathrm{}\left(a_p+M\right)_𝐧}{\left(b_1+M\right)_𝐧\mathrm{}\left(b_s+M\right)_𝐧}.$$
(47)
One can get the hypergeometric function related to Schur functions by putting $`\beta _1=1`$ and $`\beta _i=0`$ for $`i=(2,3,\mathrm{})`$ in (47):
$`{}_{p}{}^{}F_{s}^{}\left({\displaystyle \genfrac{}{}{0pt}{}{a_1+M,\mathrm{},a_p+M}{b_1+M,\mathrm{},b_s+M}}|t_1,t_2,\mathrm{}\right)=`$
$`{\displaystyle \underset{n=0}{\overset{+\mathrm{}}{}}}{\displaystyle \underset{\left|𝐧\right|=n}{}}{\displaystyle \frac{\left(a_1+M\right)_𝐧\mathrm{}\left(a_p+M\right)_𝐧}{\left(b_1+M\right)_𝐧\mathrm{}\left(b_s+M\right)_𝐧}}{\displaystyle \frac{s_𝐧\left(𝐭\right)}{H_𝐧}}`$ (48)
with hook polynomial $`H_𝐧=s_𝐧\left(\beta \right)=_{(i,j)𝐧}h_{ij}`$, $`h_{ij}=\left(n_i+n_j^{}ij+1\right)`$ .
We obtain the ordinary hypergeometric function of one variable of type
$${}_{p1}{}^{}F_{s}^{}\left(\pm t_1\beta _1\right)=\tau (\pm 1,𝐭,\beta ),$$
(49)
if we take $`a_1=0`$, $`𝐭=(t_1,0,0,\mathrm{})`$, $`\beta =(\beta _1,0,0,\mathrm{})`$ .
Now we consider the following change of variables $`t_m=_{i=1}^N\frac{x_i^m}{m}`$. In this case the formula (3.2) turns out to be the hypergeometric function of matrix argument, see (9). Taking $`N=1`$ we obtain the ordinary hypergeometric function of one variable, $`x_1`$. Formula (67) will explain the connection between this hypergeometric function and the function (49). The ordinary hypergeometric function satisfies well-known hypergeometric equation
$$\left(_{x_1}{}_{p}{}^{}r_{s}^{}\left(D_1\right)\right){}_{p}{}^{}F_{s}^{}\left(x_1\right)=0,D_1:=x_1_{x_1}.$$
(50)
Example 3 The q-generalization of the Example 2:
$${}_{p}{}^{}r_{s}^{\left(q\right)}\left(D\right)=\frac{_{i=1}^p\left(1q^{a_i+D}\right)}{_{i=1}^s\left(1q^{b_i+D}\right)}.$$
(51)
For the variables $`\beta _i=\frac{1}{i\left(1q^i\right)}`$ ($`i=1,2,\mathrm{}`$) and $`t_m=_{j=1}^N\frac{x_j^m}{m}`$ we obtain the formula (1) (see )
$`{}_{p}{}^{}\tau _{s}^{\left(q\right)}\left({\displaystyle \genfrac{}{}{0pt}{}{a_1+M,\mathrm{},a_p+M}{b_1+M,\mathrm{},b_s+M}}|q,𝐱_{\left(N\right)}\right)=`$
$`{\displaystyle \underset{n=0}{\overset{+\mathrm{}}{}}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{\left|𝐧\right|=n}{l\left(𝐧\right)N}}{}}{\displaystyle \frac{(q^{a_1+M};q)_𝐧\mathrm{}(q^{a_p+M};q)_𝐧}{(q^{b_1+M};q)_𝐧\mathrm{}(q^{b_s+M};q)_𝐧}}{\displaystyle \frac{q^{n\left(𝐧\right)}}{H_𝐧\left(q\right)}}s_𝐧\left(𝐱_{\left(N\right)}\right).`$ (52)
When $`N=1`$ we have the ordinary q-deformed hypergeometric function:
$`{}_{p}{}^{}\mathrm{\Phi }_{s}^{}\left({\displaystyle \genfrac{}{}{0pt}{}{a_1+M,\mathrm{},a_p+M}{b_1+M,\mathrm{},b_s+M}}|q,x_1\right)={\displaystyle \underset{n=0}{\overset{+\mathrm{}}{}}}{\displaystyle \frac{(q^{a_1+M};q)_n\mathrm{}(q^{a_p+M};q)_n}{(q^{b_1+M};q)_n\mathrm{}(q^{b_s+M};q)_n}}{\displaystyle \frac{x_1^n}{(q;q)_n}}`$ (53)
which satisfies the following q-difference equation
$$\left(\frac{1}{x_1}\left(1q^{D_1}\right){}_{p}{}^{}r_{s}^{\left(q\right)}\left(D_1\right)\right){}_{p}{}^{}\mathrm{\Phi }_{s}^{}\left(x_1\right)=0,D_1:=x_1_{x_1}.$$
(54)
There are various applications for series (53), for instance see , and . Bosonic representation of (53) was found in . Let us note that operator $`q^D`$ which acts on fermions $`\psi \left(z\right)`$ was used in in different context.
Example 4 q-Askey-Wilson polynomials
Let operator $`{}_{4}{}^{}r_{3}^{\left(q\right)}\left(D\right)`$ be
$${}_{4}{}^{}r_{3}^{\left(q\right)}\left(D\right)=\frac{\left(1q^{n+D}\right)\left(1abcdq^{n1+D}\right)\left(1ae^{i\eta }q^D\right)\left(1ae^{i\eta }q^D\right)}{\left(1abq^D\right)\left(1acq^D\right)\left(1adq^D\right)}.$$
(55)
By choosing $`\beta _i=\frac{1}{i\left(1q^i\right)}`$ and $`𝐭=(q,\frac{q^2}{2},\frac{q^3}{3},\mathrm{})`$ we get
$`{}_{4}{}^{}\tau _{3}^{\left(q\right)}(M,𝐭,\beta )={}_{4}{}^{}\phi _{3}^{}\left({\displaystyle \genfrac{}{}{0pt}{}{q^{Mn},q^{M+n1}abcd,aq^Me^{i\eta },aq^Me^{i\eta }}{q^Mab,q^Mac,q^Mad}}|q,q\right)=`$
$`{\displaystyle \underset{m=0}{\overset{+\mathrm{}}{}}}{\displaystyle \frac{(q^{Mn};q)_m(q^{M+n1}abcd;q)_m(aq^Me^{i\eta };q)_m(aqe^{i\eta };q)_m}{(abq^M;q)_m(acq^M;q)_m(adq^M;q)_m}}{\displaystyle \frac{q^m}{(q;q)_m}}.`$ (56)
For $`M=0`$ we obtain q-Askey-Wilson polynomials:
$`p_n(\mathrm{cos}\eta ;a,b,c,d|q)=aq^n(ab;q)_n(ac;q)_n(ad;q)_n{}_{4}{}^{}\tau _{3}^{\left(q\right)}(𝐭,\beta ,0).`$ (57)
Example 5 Clebsch-Gordan coefficients $`C_q(𝐥,𝐣)`$ see .
Let us take
$${}_{3}{}^{}r_{2}^{\left(q\right)}\left(D\right)=\frac{\left(1q^{jl_1+D}\right)\left(1q^{l_1+j+1+D}\right)\left(1q^{l+m+D}\right)}{\left(1q^{l_2l+j+1+D}\right)\left(1q^{ll_2+j+D}\right)}.$$
(58)
For the variables $`𝐭=(q,\frac{q^2}{2},\frac{q^3}{3},\mathrm{})`$ and $`\beta _i=\frac{1}{i\left(1q^i\right)}`$
$${}_{3}{}^{}\tau _{2}^{\left(q\right)}(0,𝐭,\beta )={}_{3}{}^{}\mathrm{\Phi }_{2}^{}\left(\genfrac{}{}{0pt}{}{jl_1,l_1+j+1,l+m}{l_2l+j+1,ll_2+j}|q,q\right).$$
(59)
Thus we have the following fermionic representation: $`C_q(𝐥,𝐣)=`$
$$\frac{\left(1\right)^{l_1j}q^B\mathrm{\Delta }\left(𝐥\right)\left[l+l_2j\right]!\left([𝐥,𝐣]\left[2l+1\right]\right)^{\frac{1}{2}}}{\left[l_1l_2+l\right]!\left[l+l_2l_1\right]!\left[l_2l+j\right]!\left[l_1j\right]!\left[l_2+k\right]!\left[lm\right]!}{}_{3}{}^{}\tau _{2}^{\left(q\right)}(0,𝐭,\beta ),$$
(60)
$$\left[a\right]:=q^{\left(1a\right)/2}\frac{1q^a}{1q},\left[n\right]!:=\left[1\right]\left[2\right]\mathrm{}\left[n\right],m=j+k,$$
$$\mathrm{\Delta }\left(𝐥\right)\mathrm{\Delta }(l_1,l_2,l):=\left(\frac{\left[l_1+l_2l\right]!\left[l_1l_2+l\right]!\left[ll_1+l_2\right]!}{\left[l_1+l_2+l+1\right]!}\right)^{\frac{1}{2}},$$
$$[𝐥,𝐣]=\left[l_1+j\right]!\left[l_1j\right]!\left[l_2+k\right]!\left[l_2k\right]!\left[l+m\right]!\left[lm\right]!,$$
$$B=\frac{1}{4}\left(l_2\left(l_2+1\right)l_1\left(l_1+1\right)l\left(l+1\right)+2j\left(m+1\right)\right).$$
### 3.3 Different representations
Let us rewrite hypergeometric series in different way representing all Pochhammer coefficients $`(q^a;q)_𝐧`$ and $`\left(a\right)_𝐧`$ through Schur functions. This gives us the opportunity to interchange the role of Pochhammer coefficients and Schur functions in (8),(3.2), and to present different fermionic representations of the hypergeometric functions. We have the following relations (see ):
$$\underset{(i,j)𝐧}{}\left(1q^{a+ji}\right)=\frac{s_𝐧\left(𝐭(a,q)\right)}{s_𝐧\left(𝐭(+\mathrm{},q)\right)},\underset{(i,j)𝐧}{}\left(a+ji\right)=\frac{s_𝐧\left(𝐭\left(a\right)\right)}{s_𝐧(𝐭(+\mathrm{})},$$
(61)
where parameters $`t_m(a,q)`$ and $`t_m\left(a\right)`$ are chosen via generalized Miwa transform with multiplicity $`a`$
$`t_m(a,q)={\displaystyle \frac{1\left(q^a\right)^m}{m\left(1q^m\right)}},t_m\left(a\right)={\displaystyle \frac{a}{m}},m=1,2,\mathrm{},`$ (62)
$`s_𝐧\left(𝐭(+\mathrm{},q)\right)=\underset{a+\mathrm{}}{lim}s_𝐧\left(𝐭(a,q)\right)={\displaystyle \frac{q^{n\left(𝐧\right)}}{H_𝐧\left(q\right)}},`$ (63)
$`s_𝐧\left(𝐭\left(+\mathrm{}\right)\right)=\underset{a+\mathrm{}}{lim}s_𝐧({\displaystyle \frac{t_1\left(a\right)}{a}},{\displaystyle \frac{t_2\left(a\right)}{a^2}},\mathrm{})=\underset{a+\mathrm{}}{lim}{\displaystyle \frac{1}{a^{\left|𝐧\right|}}}s_𝐧\left(𝐭\left(a\right)\right)={\displaystyle \frac{1}{H_𝐧}}.`$ (64)
Now we rewrite the series (52) and (3.2) only in terms of Schur functions:
$`{}_{p}{}^{}\mathrm{\Phi }_{s}^{}\left({\displaystyle \genfrac{}{}{0pt}{}{a_1+M,\mathrm{},a_p+M}{b_1+M,\mathrm{},b_s+M}}|q,𝐱_{\left(N\right)}\right)=\tau _r(M,𝐭(+\mathrm{},q),𝐭)`$
$`={\displaystyle \underset{n=0}{\overset{+\mathrm{}}{}}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{\left|𝐧\right|=n}{l\left(𝐧\right)N}}{}}{\displaystyle \frac{_{k=1}^ps_𝐧\left(𝐭(a_k+M,q)\right)}{_{k=1}^ss_𝐧\left(𝐭(b_k+M,q)\right)}}\left(s_𝐧\left(𝐭(+\mathrm{},q)\right)\right)^{sp+1}s_𝐧\left(𝐱_{\left(N\right)}\right),`$ (65)
$`{}_{p}{}^{}F_{s}^{}\left({\displaystyle \genfrac{}{}{0pt}{}{a_1+M,\mathrm{},a_p+M}{b_1+M,\mathrm{},b_s+M}}|𝐱_{\left(N\right)}\right)=\tau _r(M,𝐭\left(+\mathrm{}\right),𝐭)=`$
$`{\displaystyle \underset{n=0}{\overset{+\mathrm{}}{}}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{\left|𝐧\right|=n}{l\left(𝐧\right)N}}{}}{\displaystyle \frac{_{k=1}^ps_𝐧\left(𝐭\left(a_k+M\right)\right)}{_{k=1}^ss_𝐧\left(𝐭\left(b_k+M\right)\right)}}\left(s_𝐧\left(𝐭\left(+\mathrm{}\right)\right)\right)^{sp+1}s_𝐧\left(𝐱_{\left(N\right)}\right).`$ (66)
We obtain different fermionic representations of hypergeometric functions (3.3), (65) and they are parametrized by a complex number $`b`$:
Proposition 4. For $`bC`$ and for $`r={}_{p}{}^{}r_{s}^{}`$ (see (45)) we have
$$\tau _r(M,𝐭\left(+\mathrm{}\right),𝐭)=\tau _{r_b}(M,𝐭\left(b+M\right),𝐭),r_b=\frac{r}{b+D}.$$
(67)
For $`r={}_{p}{}^{}r_{s}^{\left(q\right)}`$ (see (51)) we have
$$\tau _r(M,𝐭(+\mathrm{},q),𝐭)=\tau _{r_b}(M,𝐭(b+M,q),𝐭),r_b=\frac{r}{1q^{b+D}}.$$
(68)
Remark 1. There are two ways to restrict the sum (32) to the sum over partitions of length $`l\left(𝐧\right)N`$. First, if we use Miwa’s change (4), then $`s_𝐧\left(𝐱_{\left(N\right)}\right)=0`$, for $`𝐧`$ with length $`l\left(𝐧\right)>N`$. The second way is to restrict the Pochhammer coefficients: if we put $`a_i=N`$ for one $`i`$ from (51) equal to $`N`$, then the coefficient (5) vanishes for $`l\left(𝐧\right)>N`$. Since we expressed Pochhammer coefficients in terms of Schur functions in (61) both ways have the same explanation. Indeed
$$t_m(N,q)=\frac{1}{m}\frac{1\left(q^N\right)^m}{1q^m}=\frac{1}{m}\left(1+\left(q\right)^m+\left(q^2\right)^m+\mathrm{}+\left(q^{N1}\right)^m\right).$$
(69)
Therefore we obtain for Miwa’s change: $`x_1=1,x_2=q,\mathrm{},x_N=q^{N1}`$ and
$$s_𝐧\left(𝐭(N,q)\right)=s_𝐧(1,q,\mathrm{},q^{N1})=0,l\left(𝐧\right)>N.$$
(70)
The same we have for the sum over partitions $`𝐧`$ such that $`l\left(𝐧^{}\right)<K`$. Again the first way has to be realized through the following Miwa’s change of variables:
$$t_m=\underset{i=1}{\overset{K}{}}\frac{x_i^m}{m},s_𝐧\left(𝐭\right)=s_𝐧^{}\left(𝐱_{\left(K\right)}\right).$$
(71)
The second way is to make one of the parameters, for example $`a_j`$ from (51) equal to $`\left(K\right)`$. In this case
$$s_𝐧\left(𝐭(K,q)\right)=s_𝐧^{}(\frac{1}{q},\frac{1}{q^2},\mathrm{},\frac{1}{q^K})=0,l\left(𝐧^{}\right)>K.$$
(72)
### 3.4 Zakharov-Shabat dressing and string equations
Now we shall get hypergeometric functions from another point of view - the Zakharov-Shabat factorization problem. Let us introduce infinite matrices to describe KP and Toda lattice flows and symmetries . We denote the Zakharov-Shabat dressing matrices by $`K`$ and $`\overline{K}`$. $`K`$ is a lower triangular matrix with unit main diagonal, and $`\overline{K}`$ is an upper triangular matrix. They solve Gauss factorization problem for infinite matrices:
$$\overline{K}=KU(M,𝐭,\beta ),U(M,𝐭,\beta )=U^+\left(𝐭\right)U^{}(M,\beta ).$$
(73)
Here $`U^\pm `$ belong to the different abelian multiparametric subgroups in $`GL\left(\mathrm{}\right)`$ with the infinite sets of group times $`𝐭`$ and $`\beta `$,
$$U^+\left(𝐭\right)=\mathrm{exp}\left(\xi (𝐭,\mathrm{\Lambda })\right),U^{}(M,\beta )=\mathrm{exp}\left(\xi (\beta ,\mathrm{\Lambda }^1r\left(\mathrm{\Delta }+M\right))\right),$$
(74)
where the matrices $`\left(\mathrm{\Lambda }\right)_{jk}=\delta _{j,k1}`$, $`\left(\mathrm{\Delta }\right)_{jk}=j\delta _{j,k}`$, for $`\xi `$ see (3). The function $`r`$ is the same as in (27). Then following Zakharov-Shabat arguments we find that the variables $`\mathrm{log}\left(\overline{K}_{ii}\right)=\varphi _{i+M}`$ solve Toda lattice equation in the form (41). At the same time (73) describes a set of KP equations parametrized by integer $`M`$.
The tau-function can be obtained as follows. By taking the projection $`UU_{}`$ for nonpositive values of matrix indices we obtain a determinant representation of the tau-function (32):
$$\tau _r(M,𝐭,\beta )=\frac{detU_{}(M,𝐭,\beta )}{det\left(U_{}^+\left(𝐭\right)\right)det\left(U_{}^{}(M,\beta )\right)}=detU_{}(M,𝐭,\beta ),$$
(75)
since both determinants in the denominator are equal to one. Formula (75) is also a Segal-Wilson formula for $`GL\left(\mathrm{}\right)`$ 2-cocycle . Choosing the function $`r`$ as in Section 3.2 we obtain hypergeometric functions listed in the Introduction.
Remark 2. Therefore the hypergeometric functions which were considered above have the meaning of $`GL\left(\mathrm{}\right)`$ two-cocycle on the two multiparametrical group elements $`U^+\left(𝐭\right)`$ and $`U^{}(M,\beta )`$. Both elements $`U^+\left(𝐭\right)`$ and $`U^{}(M,\beta )`$ can be considered as elements of group of pseudodifferential operators on the circle. The corresponding Lie algebras consist of the multiplication operators $`\left\{z^n;nN_0\right\}`$ and of the pseudodifferential operators $`\left\{\left(\frac{1}{z}r\left(z\frac{d}{dz}+M\right)\right)^n;nN_0\right\}`$. Two sets of group times $`𝐭`$ and $`\beta `$ play the role of indeterminates of the hypergeometric functions (47). Formulas (32) and (34) mean the expansion of $`GL\left(\mathrm{}\right)`$ group 2-cocycle in terms of corresponding Lie algebra 2-cocycle
$$\omega (z,\frac{1}{z}r\left(D+M\right))=r\left(M\right),\omega (\stackrel{~}{r}\left(D+M\right)z,\frac{1}{z}r\left(D+M\right))=\stackrel{~}{r}\left(M\right)r\left(M\right).$$
(76)
Now following let us briefly describe additional symmetries mentioned in Section 3.1 and the string equations. Below we use the notations from the papers for $`L,\overline{L},\widehat{M};[\widehat{M},L]=1`$. The Toda lattice additional symmetries and higher flows are given by the following Lax equations
$$_{\beta _n}L=[L,\left(\left(r\left(\widehat{M}\right)L^1\right)^n\right)_{}],_{\beta _n}\overline{L}=[\overline{L},\left(\left(r\left(\widehat{M}\right)L^1\right)^n\right)_+],$$
(77)
$$_{t_n^{}}L=[L,\left(\overline{L}^n\right)_{}],_{t_n^{}}\overline{L}=[\overline{L},\left(\overline{L}^n\right)_+].$$
(78)
Let us impose the condition that the group times $`\beta _m`$ of additional symmetries can be identified with the Toda lattice times $`t_m^{}`$. Then one can obtain a set of string equations , which characterizes the TL hypergeometric solutions:
$$h\left(\widehat{M}1\right)L^m=\overline{L}^mh\left(\widehat{M}1\right),m=2,3,\mathrm{},$$
(79)
where functions $`h`$ and $`r`$ are connected by (35). The simplest string equation is $`L=\overline{L}r\left(\widehat{M}\right)`$. When $`r\left(M\right)=M+a`$ the string equation describes $`c=1`$ string .
At last let us mention the $`\overline{}`$ problem for the KP-1 equation
$$\frac{w(M,𝐭,𝐭^{},z,\overline{z})}{\overline{z}}=w(M,𝐭,𝐭^{},\overline{z},z)T(z,\overline{z}),$$
(80)
where for $`z\mathrm{}`$ and for $`z0`$ the asymptotics of Baker function $`w`$ are
$$w=z^Me^{\xi (𝐭,z)}\left(1+O\left(z^1\right)\right),w=z^Me^{\xi (𝐭^{},z^1)}e^{\phi _M(𝐭,𝐭^{})}\left(1+O\left(z\right)\right)$$
(81)
with some function $`\phi _M`$. Instead of writing down the explicit form of the function $`T`$, which gives the hypergeometric solution of the KP-1 equation, we only give a set of constraints on $`T`$ (which are equivalent to the conditions mentioned above: $`𝐭^{}`$ may be identified with times of additional symmetries $`\beta `$ ):
$$\left(\frac{1}{z}r\left(D\right)\right)^mT(z,\overline{z})=\left(\frac{1}{\overline{z}}r\left(\overline{D}\right)\right)^mT(z,\overline{z}),m=1,2,\mathrm{}.$$
(82)
For the KP-2 equation we have nonlocal Riemann problem, see
$$w_+\left(z\right)w_{}\left(z\right)=_Rw_{}\left(z^{}\right)R(z^{},z)𝑑z^{}$$
(83)
with the similar constraints on the kernel $`R`$:
$$\left(\frac{1}{z}r\left(z_z\right)\right)^mR(z^{},z)=\left(\frac{1}{z^{}}r\left(z^{}_z^{}\right)\right)^mR(z^{},z),m=1,2,\mathrm{}.$$
(84)
The statement is the following. If a hypergeometrical solution belongs to the class described by (80) (or by (83)), then the function $`T`$ (respectively $`R`$) solves equation (82) (respectively (84)).
We present the infinitesimal version of the Zakharov-Shabat dressing
$$\frac{1}{2\pi i}\left(\left(z^{}\frac{1}{z}r\left(D\right)\right)^1w\left(z\right)\right)_x^1w^{}\left(z\right)dz=\underset{n=1}{}z_{}^{}{}_{}{}^{n1}A_n$$
(85)
as a generating function for the Zakharov-Shabat zero curvature equations for the additional symmetries $`[_{t_n}B_n,_{\beta _m}A_m]=0`$, compare with .
From soliton theory and bosonization formula one can obtain various relations for the tau-function (32). Let us consider an example. We introduce
$$\mathrm{\Omega }_r\left(\beta \right):=\frac{1}{2\pi i}\underset{ϵ0}{lim}V^{}\left(z+ϵ\right)\underset{m=1}{\overset{\mathrm{}}{}}\beta _m\left(\frac{1}{z}r(z,D)\right)^mV\left(z\right)dz,$$
(86)
where $`V\left(z\right)`$ and $`V^{}\left(z\right)`$ are defined by (24). Now we have the following shift argument formula
$$e^{\mathrm{\Omega }_r\left(\gamma \right)}\tau _r(M,𝐭,\beta )=\tau _r(M,𝐭,\beta +\gamma ).$$
(87)
### 3.5 Orthogonal polynomials
It is known that the hypergeometric functions (9) appear in the group representation theory and are connected with the so-called matrix integrals . On the other hand the set of examples reveals a connection between the matrix integrals and the soliton theory. To establish this connection it is useful to consider the related systems of the orthogonal polynomials. Let us briefly describe how to write down these polynomials. Let $`M_+`$ be the largest integer zero of the function $`r`$. Then the function
$$f\left(zz^{}\right)=\underset{n=0}{\overset{+\mathrm{}}{}}\left(zz^{}\right)^{n+M_+}h\left(n+M_+\right)$$
(88)
is the eigenfunction of the operator $`\frac{1}{z}r\left(D\right)`$ with the eigenvalue $`z^{}`$. We use this function as weight function for a system of orthogonal polynomials $`p_n^\pm ,n=0,1,2,\mathrm{}`$, related to the hypergeometric solution of KP:
$$_\gamma p_n^{}(z,𝐭,\beta )e^{\xi (z,𝐭)+\xi (z^{},\beta )}f\left(zz^{}\right)p_m^+(z^{},𝐭,\beta )𝑑z𝑑z^{}=e^{\varphi _{M_++n}(𝐭,\beta )}\delta _{n,m}.$$
(89)
### 3.6 Further generalization
Formula (39) is related to ’Gauss decomposition’ of operators inside vacuums $`0\left|\mathrm{}\right|0`$ into diagonal operator $`e^{H_0\left(h\right)}`$, upper triangular operator $`e^{H\left(𝐭\right)}`$ and lower triangular operator $`e^{H^{}\left(𝐭^{}\right)}`$ (the last two have the Toeplitz form). Now let us consider more general two-dimensional Toda lattice tau-function
$`\tau =M\left|e^{H\left(𝐭\right)}ge^{A\left(\beta \right)}\right|M,g=e^{\stackrel{~}{A}_1\left(\stackrel{~}{\gamma }_1\right)}\mathrm{}e^{\stackrel{~}{A}_k\left(\stackrel{~}{\gamma }_k\right)}e^{A_l\left(\gamma _l\right)}\mathrm{}e^{A_1\left(\gamma _1\right)},`$ (90)
where operators $`A_i\left(\gamma _i\right)=_{j=1}^+\mathrm{}\gamma _{ij}A_{ij}`$ and $`\stackrel{~}{A_i}\left(\stackrel{~}{\gamma }_i\right)=_{j=1}^+\mathrm{}\stackrel{~}{\gamma }_{ij}\stackrel{~}{A}_{ij}`$ are defined like $`A\left(\beta \right)`$ of (27) and $`\stackrel{~}{A}\left(\stackrel{~}{\beta }\right)`$ of (33) respectively and correspond to operators $`r^i\left(D\right)`$ and $`\stackrel{~}{r}^i\left(D\right)`$. Collections of variables $`\stackrel{~}{\gamma }=\left\{\stackrel{~}{\gamma }_{ij}\right\},\gamma =\left\{\gamma _{ij}\right\}`$ play the role of coordinates for some wide enough class of Clifford group elements $`g`$. Let us calculate this tau-function. We introduce a set consisting of $`m+1`$ partitions:
$$\left(𝐧_\mathrm{𝟏},\mathrm{},𝐧_𝐦,𝐧_{𝐦+\mathrm{𝟏}}=𝐧\right),0𝐧_\mathrm{𝟏}𝐧_\mathrm{𝟐}\mathrm{}𝐧_𝐦𝐧_{𝐦+\mathrm{𝟏}}=𝐧.$$
(91)
Also we define a set $`\mathrm{\Theta }_𝐧^m=(𝐧_\mathrm{𝟏},𝐧_\mathrm{𝟐}/𝐧_\mathrm{𝟏},\mathrm{},𝐧/𝐧_𝐦)`$. We take
$$s_{\mathrm{\Theta }_𝐧^m}\left(\mu \right)=s_{𝐧_\mathrm{𝟏}}\left(\mu _1\right)s_{𝐧_\mathrm{𝟐}/𝐧_\mathrm{𝟏}}\left(\mu _2\right)\mathrm{}s_{𝐧/𝐧_𝐦}\left(\mu _{m+1}\right),\mu _i=\left\{\mu _{ij}\right\}.$$
(92)
Here $`s_{𝐧_{𝐢+\mathrm{𝟏}}/𝐧_𝐢}\left(\mu _{i+1}\right)`$ is a skew Schur function :
$$s_{𝐧_{𝐢+\mathrm{𝟏}}/𝐧_𝐢}\left(\mu _{i+1}\right)=det\left(p_{n_\rho ^{i+1}n_\nu ^i\rho +\nu }\left(\mu _{i+1}\right)\right)_{1\rho ,\nu r},𝐧_{𝐢+\mathrm{𝟏}}=(n_1^{\left(i+1\right)},\mathrm{},n_r^{\left(i+1\right)}).$$
(93)
A skew analogy of (29) is
$`r_{\mathrm{\Theta }_𝐧^m}\left(M\right)=r_{𝐧_\mathrm{𝟏}}\left(M\right)r_{𝐧_\mathrm{𝟐}/𝐧_\mathrm{𝟏}}^1\left(M\right)\mathrm{}r_{𝐧/𝐧_𝐦}^m\left(M\right),`$ (94)
$`r_{𝐧_{𝐢+\mathrm{𝟏}}/𝐧_𝐢}\left(M\right)={\displaystyle \underset{j=1}{\overset{r}{}}}r\left(n_j^{\left(i\right)}j+1+M\right)\mathrm{}r\left(n_j^{\left(i+1\right)}j+M\right).`$ (95)
If the function $`r^i\left(m\right)`$ has no poles and zeroes at integer points, then the relation $`r_{\theta _i}^i\left(M\right)=\frac{r_{𝐧_{𝐢+\mathrm{𝟏}}}^i\left(M\right)}{r_{𝐧_𝐢}^i\left(M\right)}`$ ($`i=1,\mathrm{},m`$) is correct. To calculate the tau function we need the following Lemma
Lemma 3 Let partitions
$$𝐧=(i_1,\mathrm{},i_s|j_11,\mathrm{},j_s1),\stackrel{~}{𝐧}=(\stackrel{~}{i}_1,\mathrm{},\stackrel{~}{i}_r|\stackrel{~}{j}_11,\mathrm{},\stackrel{~}{j}_r1)$$
(96)
satisfy the relation $`𝐧\stackrel{~}{𝐧}`$. Then the next formula is valid:
$`0\left|\psi _{\stackrel{~}{i}_1}^{}\mathrm{}\psi _{\stackrel{~}{i}_r}^{}\psi _{\stackrel{~}{j}_r}\mathrm{}\psi _{\stackrel{~}{j}_1}e^{A^i\left(\gamma _i\right)}\psi _{j_1}^{}\mathrm{}\psi _{j_s}^{}\psi _{i_s}\mathrm{}\psi _{i_1}\right|0=`$
$`=\left(1\right)^{\stackrel{~}{j}_1+\mathrm{}+\stackrel{~}{j}_r+j_1+\mathrm{}+j_s}s_\theta \left(\gamma _i\right)r_\theta \left(0\right),\theta =𝐧/\stackrel{~}{𝐧}.`$ (97)
The proof is achieved by direct calculation (see for help).
Now we obtain the following generalization of Proposition 1:
Proposition 5
$$\tau (M,𝐭,\beta ;\gamma ,\stackrel{~}{\gamma })=\underset{𝐧}{}\underset{\mathrm{\Theta }_𝐧^k}{}\underset{\mathrm{\Theta }_𝐧^l}{}\stackrel{~}{r}_{\mathrm{\Theta }_𝐧^k}\left(M\right)r_{\mathrm{\Theta }_𝐧^l}\left(M\right)s_{\mathrm{\Theta }_𝐧^k}(𝐭,\stackrel{~}{\gamma })s_{\mathrm{\Theta }_𝐧^l}(\beta ,\gamma ).$$
(98)
With the help of this series one can obtain different hypergeometric series.
Example 6. Let us consider the tau function given by the correlator
$`\tau (M,\stackrel{~}{\beta },\beta ,\gamma _1)=M\left|e^{\stackrel{~}{A}\left(\stackrel{~}{\beta }\right)}e^{A_1\left(\gamma _1\right)}e^{A\left(\beta \right)}\right|M,`$ (99)
$`\stackrel{~}{r}^1\left(D\right)={\displaystyle \frac{\stackrel{~}{a}_1+D}{\stackrel{~}{b}_1+D}},r\left(D\right)={\displaystyle \frac{a_1+D}{b_1+D}},r^1\left(D\right)=1.`$ (100)
We put $`\stackrel{~}{\beta }=(x,\frac{x^2}{2},\frac{x^3}{3},\mathrm{})`$, $`\beta =(y_1,0,0,\mathrm{})`$, $`\gamma _1=(y_2,0,\mathrm{})`$. Thus we have
$$\tau (M,x,y_1,y_2)=\underset{n_1=0}{\overset{+\mathrm{}}{}}\underset{n_2=0}{\overset{+\mathrm{}}{}}\frac{\left(\stackrel{~}{a}_1+M\right)_{n_1+n_2}\left(a_1+M\right)_{n_1}}{\left(\stackrel{~}{b}_1+M\right)_{n_1+n_2}\left(b_1+M\right)_{n_1}}\frac{y_1^{n_1}y_2^{n_2}}{n_1!n_2!}x^{n_1+n_2}.$$
(101)
## Conclusion
We get multivariate hypergeometric functions as certain tau-functions of the KP hierarchy and also as the ratios of Toda lattice tau-functions considered in , evaluated at certain values of higher Toda lattice times. It means that multivariate hypergeometric functions solve a set of continuous and discrete bilinear Hirota equations . We shall write down these equations explicitly in the different paper. Hypergeometric solution of the KP equation is $`u=2_{t_1}^2\mathrm{log}\tau `$, where $`\tau `$ is a hypergeometric function. To investigate the properties of this new class of solutions is a separate interesting problem which we leave for future investigation. It is quite unexpected that we get q-deformed version of these hypergeometric functions as tau-functions not of a q-deformed KP hierarchy but of the usual KP hierarchy. It is now an interesting problem to establish links between these results, group-theoretic approach to the q-special functions and matrix integrals. We expect to work out connections with matrix models of Kontsevich type and two-matrix models related to 2D Toda lattice . We can present links between our construction in the present paper and so-called generalized Miwa’s change of variables in the three dimensional three-wave systems , and with the multicomponent KP hierarchy , it will be published separately. We hope to generalize our results to the KP hierarchies of $`B_{\mathrm{}}`$, $`C_{\mathrm{}}`$ and $`D_{\mathrm{}}`$ types to get different hypergeometric series.
## Acknowledgements
One of the authors (A.O.) is pleased to thank T.Shiota and especially Vl.Dragovich for the helpful discussions, and also L.A.Dickey for the interest. D.S. would like to thank S.Senchenko for helpful discussions.
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# A Theory of Gamma-Ray Bursts
## 1 Introduction
The discovery of afterglows to gamma-ray bursts has greatly increased the possibility of studying their physics. Since these afterglows have thus far only been seen for long gamma-ray bursts (duration $`\text{ }>2`$ s), we shall concentrate on the mechanism for this subclass. The shorter bursts (duration $`\text{ }<2`$ s) may have a different origin; specifically, it has been suggested that they are the result of compact-object mergers and therefore offer the intriguing possibility of associated outbursts of gravity waves. (Traditionally, binary neutron stars have been considered in this category (Eichler et al. 1989, Janka et al. 1999). More recently, Bethe & Brown (1998) have shown that low-mass black-hole, neutron-star binaries, which have a ten times greater formation rate and are stronger gravity-wave emitters, may be the more promising source of this kind.)
An important recent clue to the origin of long bursts is the probable association of some of them with ultra-bright type Ibc supernovae (Galama et al. 1998, Bloom et al. 1999, Galama et al. 2000). The very large explosion energy<sup>1</sup><sup>1</sup>1Höflich et al. (1999) have proposed that the explosion energy was not much larger than usual, but that the explosion was very asymmetric; this model also provides a reasonable fit to the light curve of SN 1998bw. implied by fitting the light curve of SN 1998bw, which was associated with GRB 980425, indicates that a black hole was formed in this event (Iwamoto et al. 1998). This provides two good pieces of astrophysical information: it implicates black holes in the origin of gamma-ray bursts, and it demonstrates that a massive star can explode as a supernova even if its core collapses into a black hole.
In this paper, we start from the viewpoint that the gamma-ray burst is powered by electromagnetic energy extraction from a spinning black hole, the so-called Blandford-Znajek (1977) mechanism. This was worked out in detail by Lee, Wijers, & Brown (1999), and further details and comments were discussed by Lee, Brown, & Wijers (2000), who built on work by Thorne et al. (1986) and Li (2000). They have shown that with the circuitry in a 3$`+`$1 dimensional description using the Boyer-Lindquist metric, one can have a simple pictorial model for the BZ mechanism.
The simple circuitry which involves steady state current flow is, however, inadequate for describing dissipation of the black hole rotational energy into the accretion disk formed from the original helium envelope. In this case the more rapidly rotating black hole tries to spin up the inner accretion disk through the closed field lines coupling the black hole and disk. Electric and magnetic fields vary wildly with time. Using the work of Blandford & Spruit (2000) we show that this dissipation occurs in an oscillatory fashion, giving a fine structure to the GRB, and that the total dissipation should furnish an energy comparable to that of the GRB to the accretion disk. We use this energy to drive the hypernova explosion.
Not any black-hole system will be suitable for making GRB: the black hole must spin rapidly enough and be embedded in a strong magnetic field. Moreover, the formation rate must be high enough to get the right rate of GRB even after accounting for substantial collimation of GRB outflows. We explore a variety of models, and give arguments why some will have sufficient energy and extraction efficiency to power a GRB and a hypernova. We argue that the systems known as black-hole transients are the relics of GRBs, and discuss the recent evidence from high space velocities and chemical abundance anomalies that these objects are relics of hypernovae and GRBs; we especially highlight the case of Nova Scorpii 1994 (GRO J1655$``$40).
The plan of this paper is as follows. We first show that it is reasonable to expect similar energy depositions into the GRB outflow and the accretion disk (Sect. 2) and discuss the amount of available energy to be extracted (Sect. 3). Then we show the agreement of those results with the detailed numerical simulations by MacFadyen & Woosley, and use those simulations to firm up our numbers (Sect. 4). We continue by presenting a simple derivation of the energetics that approximates the full results well (Sect. 5). Finally, we discuss some previously suggested progenitors (Sect. 6) and present our preferred progenitors: soft X-ray transients (Sect. 7).
## 2 Simple Circuitry
Although our numbers are based on the detailed review of Lee, Wijers, & Brown (1999), which confirms the original Blandford-Znajek (1977) paper, we illustrate our arguments with the pictorial treatment of Thorne et al. (1986) in “The Membrane Paradigm”. Considering the time as universal in the Boyer-Lindquist metric, essential electromagnetic and statistical mechanics relations apply in their 3$`+`$1 dimensional manifold. We summarize their picture in our Fig. 1.
The surface of the black hole can be considered as a conductor with surface resistance $`R_{\mathrm{BH}}=4\pi /c=377`$ ohms. A circuit that rotates rigidly with the black hole can be drawn from the loading region, the low-field region up the axis of rotation of the black hole in which the power to run the GRB is delivered, down a magnetic field line, then from the North pole of the black hole along the (stretched) horizon to its equator. From the equator we continue the circuit through part of the disk and then connect it upwards with the loading region. We can also draw circuits starting from the loading region which pass along only the black hole or go through only the disk, but adding these would not change the results of our schematic model.
Using Faraday’s law, the voltage $`V`$ can be found by integrating the vector product of charge velocity, $`\stackrel{}{v}`$, and magnetic field, $`\stackrel{}{B}`$, along the circuit:
$`V={\displaystyle [\stackrel{}{v}\times \stackrel{}{B}]𝑑\stackrel{}{l}},`$ (1)
($`d\stackrel{}{l}`$ is the line element along the circuit). Because this law involves $`\stackrel{}{v}\times \stackrel{}{B}`$ the integrals along the field lines make no contribution. We do get a contribution $`V`$ from the integral from North pole to equator along the black hole surface. Further contributions to $`V`$ will come from cutting the field lines from the disk. We assume the field to be weak enough in the loading region to be neglected.
The GRB power, $`E_{\mathrm{GRB}}`$, will be
$`\dot{E}_{\mathrm{GRB}}=I_{\mathrm{BH}+\mathrm{D}}^2R_\mathrm{L}`$ (2)
where $`R_\mathrm{L}`$ is the resistance of the loading region, and the current is given by
$`I_{\mathrm{BH}+\mathrm{D}}^2=\left({\displaystyle \frac{V_\mathrm{D}+V_{\mathrm{BH}}}{R_\mathrm{D}+R_{\mathrm{BH}}+R_\mathrm{L}}}\right)^2.`$ (3)
(The index BH refers to the black hole, L to the load region, and D to the disk.)
The load resistance has been estimated in various ways and for various assumptions by Lovelace, MacAuslan, & Burns (1979) and by MacDonald & Thorne (1982), and by Phinney (1983). All estimates agree that to within a factor of order unity $`R_\mathrm{L}`$ is equal to $`R_{\mathrm{BH}}`$.
In a similar fashion, some power will be deposited into the disk
$`\dot{E}_{disk}=I_{BH+D}^2R_D`$ (4)
but this equilibrium contribution will be small because of the low disk resistance $`R_D`$.
Blandford & Spruit (2000) have shown that important dissipation into the disk comes through magnetic field lines coupling the disk to the black hole rotation. As shown in Fig. 2 these lines, anchored in the inner disk, thread the black hole.
The more rapidly rotating black hole will provide torques, along its rotation axis, which spin up the inner accretion disk, in which the closed magnetic field lines are anchored. With increasing centrifugal force the material in the inner disk will move outwards, cutting down the accretion. Angular momentum is then advected outwards, so that the matter can drift back inwards. It then delivers more matter to the black hole and is flung outwards again. The situation is like that of a ball in a roulette wheel (R.D. Blandford, private communication). First of all it is flung outwards and then drifts slowly inwards. When it hits the hub it is again thrown outwards. The viscous inflow time for the fluctuations is easily estimated to be
$`\tau _d\mathrm{\Omega }_{disk}^1\left({\displaystyle \frac{r}{H}}\right)^2\alpha _{vis}^1`$ (5)
where $`H`$ is the height of the disk at radius $`r`$, $`\mathrm{\Omega }_{disk}`$ its angular velocity, and $`\alpha _{vis}`$ is the usual $`\alpha `$-parameterization of the viscosity. We choose $`\alpha _{vis}0.1`$, $`r/H10`$ for a thin disk and then arrive at $`\tau _d0.1`$ s. We therefore expect variability on all time scales between the Kepler time (sub-millisecond) and the viscous time, which may explain the very erratic light curves of many GRBs.
We suggest that the GRB can be powered by $`\dot{E}_{\mathrm{GRB}}`$ and a Type Ibc supernova explosion by $`\dot{E}_{\mathrm{SN}}`$ where $`\dot{E}_{\mathrm{SN}}`$ is the power delivered through dissipation into the disk. To the extent that the number of closed field lines coupling disk and black hole is equal to the number of open field lines threading the latter, the two energies will be equal. In the spectacular case of GRB 980326 (Bloom et al. 1999), the GRB lasts about $`5`$ s, which we take to be the time that the central engine operates. We shall show that up to $`10^{53}`$ erg is available to be delivered into the GRB and into the accretion disk, the latter helping to power the supernova (SN) explosion. This is more energy than needed and we suggest that injection of energy into the disk shuts off the central engine by blowing up the disk and thus removing the magnetic field needed for the energy extraction from the black hole. If the magnetic field is high enough the energy will be delivered in a short time, and the quick removal of the disk will leave the black hole still spinning quite rapidly.
## 3 Energetics of GRBs
The maximum energy that can be extracted from the BZ mechanism (Lee, Wijers, & Brown 1999) is
$`(E_{BZ})_{max}0.09M_{\mathrm{BH}}c^2.`$ (6)
This is 31% of the black hole rotational energy, the remainder going toward increasing the entropy of the black hole. This maximum energy is obtained if the extraction efficiency is
$`ϵ_\mathrm{\Omega }={\displaystyle \frac{\mathrm{\Omega }_{disk}}{\mathrm{\Omega }_H}}=0.5.`$ (7)
In Appendix A we give numerical estimates for this ratio for various $`\omega =\mathrm{\Omega }_{disk}/\mathrm{\Omega }_K`$ and various radii in the region of parameter space we consider. As explained in Section 2 we expect the material in the inner disk to swing in and out around the marginally stable radius, $`r_{ms}`$. It can be seen from the Table 2 and Appendix A that the relevant values of $`ϵ_\mathrm{\Omega }`$ are close to that of eq. (7).
For a $`7M_{}`$ black hole, such as that found in Nova Sco 1994 (GRO J1655$``$40),
$`E_{max}1.1\times 10^{54}\mathrm{erg}.`$ (8)
We estimate below that the energy available in a typical case will be an order of magnitude less than this. Without collimation, the estimated gamma-ray energy in GRB 990123 is about $`4.5\times 10^{54}`$ erg (Andersen et al. 1999). The BZ scenario entails substantial beaming, so this energy should be multiplied by $`d\mathrm{\Omega }/4\pi `$, which may be a small factor (perhaps $`10^2`$).
The BZ power can be delivered at a maximum rate of
$`P_{\mathrm{BZ}}=6.7\times 10^{50}\left({\displaystyle \frac{B}{10^{15}\mathrm{G}}}\right)^2\left({\displaystyle \frac{M_{\mathrm{BH}}}{M_{}}}\right)^2\mathrm{erg}\mathrm{s}^1,`$ (9)
(Lee et al. 1999) so that high magnetic fields are necessary for rapid delivery.
The above concerns the maximum energy output into the jet and the disk. The real energy available in black-hole spin in any given case, and the efficiency with which it can be extracted, depend on the rotation frequency of the newly formed black hole and the disk or torus around it. The state of the accretion disk around the newly formed black hole, and the angular momentum of the black hole, are somewhat uncertain. However, the conditions should be bracketed between a purely Keplerian, thin disk (if neutrino cooling is efficient) and a thick, non-cooling hypercritical advection-dominated accretion disk (HADAF), of which we have a model (Brown, Lee & Bethe 2000). Let us examine the result for the Keplerian case. In terms of
$`\stackrel{~}{a}{\displaystyle \frac{Jc}{M^2G}},`$ (10)
where $`J`$ is the angular momentum of the black hole, we find the rotational energy of a black hole to be
$`E_{rot}=f(\stackrel{~}{a})Mc^2,`$ (11)
where
$`f(\stackrel{~}{a})=1\sqrt{{\displaystyle \frac{1}{2}}(1+\sqrt{1\stackrel{~}{a}^2})}.`$ (12)
For a maximally rotating black hole one has $`\stackrel{~}{a}=1`$<sup>2</sup><sup>2</sup>2 As an aside, we note a nice mnemonic: if we define a velocity $`v`$ from the black-hole angular momentum by $`J=MR_{\mathrm{Sch}}v`$, so that $`v`$ carries the quasi-interpretation of a rotation velocity at the horizon, then $`\stackrel{~}{a}=2v/c`$. A maximal Kerr hole, which has $`R_{\mathrm{event}}=R_{\mathrm{Sch}}/2`$, thus has $`v=c`$. For $`\stackrel{~}{a}\text{ }<0.5`$, the rotation energy is well approximated by the easy-to-remember expression $`E_{\mathrm{rot}}=\frac{1}{2}Mv^2`$..
We begin with a neutron star in the middle of a Keplerian accretion disk, and let it accrete enough matter to send it into a black hole. In matter free regions the last stable orbit of a particle around a black hole in Schwarzschild geometry is
$`r_{\mathrm{lso}}=3R_{\mathrm{Sch}}=6{\displaystyle \frac{GM}{c^2}}.`$ (13)
This is the marginally stable orbit $`r_{\mathrm{ms}}`$. However, under conditions of hypercritical accretion, the pressure and energy profiles are changed and it is better to use (Abramowicz et al. 1988)
$`r_{\mathrm{lso}}\text{ }>2R_{\mathrm{Sch}}.`$ (14)
With the equal sign we have the marginally bound orbit $`r_{\mathrm{mb}}`$. With high rates of accretion we expect this to be a good approximation to $`r_{\mathrm{lso}}`$. The accretion disk can be taken to extend down to the last stable orbit (refer to Appendix B for the details).
We take the angular velocity to be Keplerian, so that the disk velocity $`v`$ at radius $`2R_{\mathrm{Sch}}`$ is given by
$`v^2={\displaystyle \frac{GM}{2R_{\mathrm{Sch}}}}={\displaystyle \frac{c^2}{4}},`$ (15)
or $`v=c/2`$. The specific angular momentum, $`l`$, is then
$`l2R_{\mathrm{Sch}}v=2{\displaystyle \frac{GM}{c}}`$ (16)
which in Kerr geometry indicates $`\stackrel{~}{a}1`$. Had we taken one of the slowest-rotating disk flows that are possible, the advection-dominated or HADAF case (Narayan and Yi 1994, Brown, Lee & Bethe 2000), which has $`\mathrm{\Omega }^2=2\mathrm{\Omega }_K^2/7`$, we would have arrived at $`\stackrel{~}{a}0.54`$, so the Kerr parameter will always be high.
Further accretion will add angular momentum to the black hole at a rate determined by the angular velocity of the inner disk. The material accreting into the black hole is released by the disk at $`r_{\mathrm{lso}}`$, where the angular momentum delivered to the black hole is determined. This angular momentum is, however, delivered into the black hole at the event horizon $`R_{\mathrm{Sch}}`$, with velocity at least double that at which it is released by the disk, since the lever arm at the event horizon is only half of that at $`R_{\mathrm{Sch}}`$, and angular momentum is conserved. With more rapid rotation involving movement towards a Kerr geometry where the event horizon and last stable orbit coincide at
$`r_{\mathrm{lso}}=R_{\mathrm{event}}={\displaystyle \frac{GM}{c^2}}.`$ (17)
Although we must switch over to a Kerr geometry for quantitative results, we see that $`\stackrel{~}{a}`$ will not be far from its maximum value of unity. Again, for the lower angular-momentum case of a HADAF, the expected black-hole spin is not much less.
## 4 Comparison with Numerical Calculation
Our schematic model has the advantage over numerical calculations that one can see analytically how the scenario changes with change in parameters or assumptions. However, our model is useful only if it reproduces faithfully the results of more complete calculations which involve other effects and much more detail than we include. We here make comparison with Fig.19 of MacFadyen & Woosley (1999). Accretion rates, etc., can be read off from their figure which we reproduce as our Fig.3. MacFadyen & Woosley prefer $`\stackrel{~}{a}_{initial}=0.5`$ (We have removed their curve for $`\stackrel{~}{a}_{initial}=0`$). This is a reasonable value if the black hole forms from a contracting proto-neutron star near breakup. MacFadyen & Woosley find that $`\stackrel{~}{a}_{initial}=0.5`$ is more consistent with the angular momentum assumed for the mantle than $`\stackrel{~}{a}_{initial}=0`$. (They take the initial black hole to have mass $`2M_{}`$; we choose the Brown & Bethe (1994) mass of $`1.5M_{}`$.) We confirm this in the next section.
After 5 seconds (the duration of GRB 980326) the MacFadyen & Woosley black hole mass is $`3.2M_{}`$ and their Kerr parameter $`\stackrel{~}{a}0.8`$, which gives $`f(\stackrel{~}{a})`$ of our eq.(12) of 0.11. With these parameters we find $`E=2\times 10^{53}`$ erg, available for the GRB and SN explosion.
One can imagine that continuation of the MacFadyen & Woosley curve for $`M_{BH}(M_{})`$ would ultimately give something like our $`7M_{}`$, but the final black hole mass may not be relevant for our considerations. This is because more than enough energy is available to power the supernova in the first 5 seconds; as the disk is disrupted, the magnetic fields supported by it will also disappear, which turns off the Blandford-Znajek mechanism.
Power is delivered at the rate given by eq.(9). Taking a black hole mass relevant here, $`3.2M_{}`$, we require a field strength of $`5.8\times 10^{15}`$ G in order for our estimated energy ($`4\times 10^{52}`$ erg) to be delivered in 5 s (the duration of GRB 980326). For such a relatively short burst, we see that the required field is quite large, but it is still not excessive if we bear in mind that magnetic fields of $`10^{15}`$ G have already been observed in magnetars (Kouveliotou 1998, 1999). Since in our scenario we have many more progenitors than there are GRBs, we suggest that the necessary fields are obtained only in a fraction of all potential progenitors.
Thus we have an extremely simple scenario for powering a GRB and the concomitant SN explosion in the black hole transients, which we will discuss in Section.7.2. After the first second the newly evolved black hole has $`10^{53}`$ erg of rotational energy available to power these. The time scale for delivery of this energy depends (inversely quadratically) on the magnitude of the magnetic field in the neighborhood of the black hole, essentially that on the inner accretion disk. The developing supernova explosion disrupts the accretion disk; this removes the magnetic fields anchored in the disk, and self-limits the energy the B-Z mechanism can deliver.
## 5 An Even More Schematic Model
Here we calculate the energy available in a rotating black hole just after its birth (before accretion adds more). Our model is to take a $`1.5M_{}`$ neutron star which co-rotates with the inner edge of the accretion disk in which it is embedded. The neutron star then collapses to a black hole, conserving its angular momentum. Since the accretion disk is neutrino cooled, but perhaps not fully thin, its angular velocity will be somewhere between the HADAF value and the Keplerian value. We parameterize it as $`\mathrm{\Omega }=\omega \mathrm{\Omega }_\mathrm{K}`$, where $`\omega =1`$ for Keplerian and $`\omega =\sqrt{2/7}0.53`$ for the HADAF.
The moment of inertia, $`I`$, of a neutron star is well fitted for many different equations of state with the simple expression
$`I={\displaystyle \frac{0.21MR^2}{12GM/Rc^2}}`$ (18)
(Lattimer & Prakash 2000). With $`J=\omega I\mathrm{\Omega }_\mathrm{K}`$ and a neutron star of $`1.5M_{}`$, with a radius of 10 km, we find
$`\stackrel{~}{a}^2=\left({\displaystyle \frac{Jc}{GM^2}}\right)^2=0.64\omega ^2.`$ (19)
We choose $`\omega 1.0`$ to roughly reproduce the MacFadyen & Woosley value of $`\stackrel{~}{a}`$, see our Fig. 3. We do not really believe the disk to be so efficiently neutrino cooled that its angular velocity is Keplerian; i.e. $`\omega =1`$, but it may be not far from it. Our $`\omega `$ should be more properly viewed as a fudge factor which allows us to match the more complete MacFadyen & Woosley calculation. MacFadyen & Woosley find that, while the accretion disk onto the black hole is forming, an additional solar mass of material is added to it “as the dense stellar core collapses through the inner boundary at all polar angles”. We shall add this to our $`1.5M_{}`$ and take the black hole mass to be $`2.5M_{}`$. We neglect the increase in spin of the black hole by the newly accreted matter; this is already included in the MacFadyen & Woosley results. For $`\stackrel{~}{a}^2=0.64`$ we find $`f(\stackrel{~}{a}^2)=0.11`$, so that the black hole rotation energy becomes
$`E_{\mathrm{BZ}}=1.5\times 10^{53}\mathrm{erg}`$ (20)
in rough agreement with the estimates of MacFadyen & Woosley in the last section.
## 6 Previous Models
### 6.1 Collapsar
We have not discussed the Collapsar model of Woosley (1993), and MacFadyen & Woosley (1999). In this model the center of a rotating Wolf-Rayet star evolves into a black hole, the outer part being held out by centrifugal force. The latter evolves into an accretion disk and then by hypercritical accretion spins the black hole up. MacFadyen & Woosley point out that “If the helium core is braked by a magnetic field prior to the supernova explosion to the extent described by Spruit & Phinney (1998) then our model will not work for single stars.” Spruit & Phinney argue that magnetic fields maintained by differential rotation between the core and envelope of the star will keep the whole star in a state of approximately uniform rotation until 10 years before its collapse. As noted in the last section, with the extremely high magnetic fields we need the viscosity would be expected to be exceptionally high, making the Spruit & Phinney scenario probable. Livio & Pringle (1998) have commented that one finds evidence in novae that the coupling between layers of the star by magnetic fields may be greatly suppressed relative to what Spruit & Phinney assumed. However, we note that even with this suppressed coupling, they find pulsar periods from core collapse supernovae no shorter than 0.1 s. Independent evidence for the fact that stellar cores mostly rotate no faster than this comes from the study of supernova remnants: Bhattacharya (1990, 1991) concludes that the absence of bright, pulsar-powered plerions in most SNRs indicates that typically pulsar spin periods at birth are no shorter than 0.03–0.05 s. Translated to our black holes, such spin periods would imply $`\stackrel{~}{a}\text{ }<0.01`$, quite insufficient to power a GRB. As a cautionary note, we might add that without magnetic coupling the cores of evolved stars can spin quite rapidly (Heger et al. 2000). This rapid initial spin may be reconciled with Bhattacharya’s limit if r-mode instabilities cause very rapid spindown in the first few years of the life of a neutron star (e.g., Heger, Langer, & Woosley 2000, Lindblom & Owen 1999).
### 6.2 Coalescing Low-Mass Black Holes and Helium Stars
Fryer & Woosley (1998) suggested the scenario of a black hole spiraling into a helium star. This is an efficient way to spin up the black hole.
Bethe & Brown (1998) evolved low-mass black holes with helium star companion, as well as binaries of compact objects. In a total available range of binary separation $`0.04<a_{13}<4`$, low-mass black-hole, neutron-star binaries were formed when $`0.5<a_{13}<1.4`$ where $`a_{13}`$ is the initial binary separation in units of $`10^{13}`$ cm. The low-mass black hole coalesces with the helium star in the range $`0.04<a_{13}<0.5`$. Binaries were distributed logarithmically in $`a`$. Thus, coalescences are more common than low-mass black-hole, neutron-star binaries by a factor of $`\mathrm{ln}(0.5/0.04)/\mathrm{ln}(1.9/0.5)=1.9`$
In Bethe & Brown (1998), the He-star, compact-object binary was disrupted $`50\%`$ of the time by the He-star explosion. This does not apply to the coalescence. Thus, the rate of low-mass black-hole, He-star mergers is 3.8 times the formation rate of low-mass black-hole, neutron-star binaries, or
$`R=3.8\times 10^4\mathrm{yr}^1`$ (21)
in the Galaxy. The estimated empirical rate of GRBs, with a factor of 100 for beaming, is $`10^5`$ yr<sup>-1</sup> in the Galaxy (Appendix C of Brown et al. 1999). Thus, the number of progenitors is more than adequate.
In Bethe & Brown (1998) the typical black hole mass was $`2.4M_{}`$, somewhat more massive than their maximum assumed neutron star mass of $`1.5M_{}`$. As it enters the helium star companion an accretion disk is soon set up and the accretion scenario will follow that described above, with rotating black holes of various masses formed. Brown, Lee, & Bethe (2000) find that the black hole will be spun up quickly. We have not pursued this scenario beyond the point that it was developed by Fryer & Woosley (1998).
## 7 Soft X-ray Transients as Relics of Hypernovae and GRB
### 7.1 Our Model: Angular Momentum
We favor a model of hypernovae similar to MacFadyen & Woosley (1999) in that it involves a failed supernova as a centerpiece. But, in distinction to MacFadyen & Woosley, our initial system is a binary, consisting of a massive star A (which will later become the failed SN) and a lighter companion B, which serves to provide ample angular momentum.
Failed supernovae require a ZAMS mass of $`2035M_{}`$, according to the calculations of Woosley & Weaver (1995) as interpreted by Brown, Lee, & Bethe (1999). The limits 20 and $`35M_{}`$ are not accurately known, but it is a fairly narrow range, so we shall in many of our calculations assume a “typical” ZAMS mass of $`25M_{}`$. The heavy star A must not be in a close binary because then its hydrogen envelope would be removed early in its evolution and therefore the star would lose mass by wind at a very early stage and become a low-mass compact object (Brown, Weingartner, & Wijers 1996). Instead, we assume a wide binary, with a separation, $`a`$ in the range
$`a=5001000R_{},`$ (22)
so star A evolves essentially as a single star through its first few burning stages. It is essential that most of the He core burning is completed before its hydrogen envelope is removed (Wellstein & Langer 1999; Heger & Wellstein 2000). We assume the initial distance $`a`$ between the two stars to be in this range. When star A fills its Roche lobe, the companion, star B, will spiral inwards.
The initiation and early development of the common envelope has been best treated by Rasio & Livio (1996). This is the only phase that can at present be modeled in a realistic way. They find a short viscous time in the envelope, but emphasize that numerical viscosity may play an important role in their results. However, we believe the viscosity to be large. Torkelsson et al. (1996) showed the Shakura-Sunyaev (1973) viscosity parameter, $`\alpha _{SS}`$, to range from 0.001 to 0.7, with the higher values following from the presence of vertical magnetic fields. Since in our Blandford-Znajek model extremely high magnetic fields $`10^{15}`$ G are needed in the He envelope to deliver the energy rapidly, we believe $`\alpha _{SS}`$ to be not much less than unity. Given such high viscosities, it seems reasonable to follow the Rasio-Livio extrapolation, based on a short viscous transport time, to later times. The most significant new result of Rasio & Livio “is that, during the dynamical phase of common envelope evolution, a corotating region of gas is established near the central binary. The corotating region has the shape of an oblate spheroid encasing the binary (i.e., the corotating gas is concentrated in the orbital plane).”
A helium core, which we deal with, is not included in their calculations, because they do not resolve the inner part of the star numerically. However, since the physics of the spiral-in does not really change as it proceeds past the end of their calculations, it seems most likely that during further spiral-in, the spin-up of material inside the orbit of the companion will continue to be significant.
Star B will stop spiraling in when it has ejected the H envelope of A. Since we assume that all stars A have about the same mass, and that $`a_i`$ is very large, we expect
$`{\displaystyle \frac{M_B}{a_f}}\mathrm{const}.`$ (23)
From section 7.2 we conclude that $`a_f`$ is a few $`R_{}`$ for $`M_B=(0.41)M_{}`$. Now the He cores of stars of ZAMS mass $`M=2035M_{}`$ have a radius about equal to $`R_{}`$. Therefore small $`M_B`$ stars will spiral into the He core of A. There they cannot be stopped but will coalesce with star A. However, they will have transmitted their angular momentum to star A.
Star B of larger mass will stop at larger $`a_fR_{}`$. It is then not clear whether they will transfer all of their angular momentum to star A. In any case, they must generally wait until they evolve off the main sequence into the subgiant or possibly even the giant stage before they can fill their Roche Lobes and later accrete onto the black hole resulting from star A.
The Kepler velocity of star B at $`a_f`$ is
$`V_K^2=G{\displaystyle \frac{M_{a_f}}{a_f}}.`$ (24)
We estimate the final mass of A, after removal of its hydrogen envelope, to be about $`10M_{}`$; then
$`V_K1.2\times 10^8a_{f,11}^{1/2}\mathrm{cm}\mathrm{s}^1,`$ (25)
where $`a_{f,11}`$ is $`a_f`$ in units of $`10^{11}`$ cm. The specific angular momentum of B is then
$`j(B)=a_fV_K=1.2\times 10^{19}a_{f,11}^{1/2}\mathrm{cm}^2\mathrm{s}^1.`$ (26)
If B and A share their angular momentum, the specific angular momentum is reduced by a factor $`M_B/(M_{A,f}+M_B)`$ which we estimate to be $`0.1`$. Since $`a_f`$ should be $`\text{ }>3R_{}`$ (See Table 1), the specific angular momentum of A should be
$`j(A)\text{ }>10^{18}\mathrm{cm}^2\mathrm{s}^1.`$ (27)
Star B has now done its job and can be disregarded.
### 7.2 Supernova and collapse
Star A now goes through its normal evolution, ending up as a supernova. But since we have chosen its mass to be between 20 and $`35M_{}`$, the SN shock cannot penetrate the heavy envelope but is stopped at some radius
$`R_{\mathrm{SN}}10^{10}\mathrm{cm},`$ (28)
well inside the outer edge of the He envelope. We estimate $`R_{\mathrm{SN}}`$ by scaling from SN 1987A: in that supernova, with progenitor mass $`18M_{}`$, most of the He envelope was returned to the galaxy. The separation between compact object and ejecta was estimated to occur at $`R5\times 10^8`$ cm (Woosley 1988, Bethe 1990) at mass point 1.5 $`M_{}`$ (gravitational). Woosley and Weaver (1995) find remnant masses of $`2M_{}`$, although with large fluctuations, for ZAMS masses in the range 20–35$`M_{}`$, which go into high-mass black holes. From table 3 of Brown, Weingartner, and Wijers (1996) we see that fallback between $`R=3.5`$ and $`4.5\times 10^8`$ cm is 0.03$`M_{}`$. Using this we can extrapolate to $`R=10^{10}`$ cm as the distance within which matter has to begin falling in immediately in our heavier stars, to make up a compact object of 2 $`M_{}`$. Unlike in 1987A the shock energy in the more massive star does not suffice to eject the envelope beyond this point, and the remaining outer envelope will also eventually fall back.
At $`R_{\mathrm{SN}}`$, the specific angular momentum of Kepler motion around a central star of mass $`10M_{}`$ is, cf. eq.(26)
$`j_K(10M_{})=1.2\times 10^{19}R_{f,11}^{1/2}\mathrm{cm}^2\mathrm{s}^1=4\times 10^{18}\mathrm{cm}^2\mathrm{s}^1.`$ (29)
In reality, at this time the central object has a mass $`M1.5M_{}`$ (being a neutron star) and since $`j_KV_KM^{1/2}`$
$`j_K(1.5M_{})=1.5\times 10^{18}\mathrm{cm}^2\mathrm{s}^1.`$ (30)
The angular momentum inherent in star A, eq.(27), is therefore greater than the Kepler angular momentum. This would not be the case had our initial object been a single star, a collapsar. (The collapsar may work none the less, but our binary model is more certain to work.)
The supernova material is supported by pressure inside the cavity, probably mostly due to electromagnetic radiation. The cavity inside $`R_{\mathrm{SN}}`$ is rather free of matter. After a while, the pressure in the cavity will reduce. This may happen by opening toward the poles, in which case the outflowing pressure will drive out the matter near the poles and create the vacuum required for the gamma ray burst. Reduction of pressure will also happen by neutrino emission. As the pressure gets reduced, the SN material will fall in toward the neutron star in the center. But because the angular momentum of the SN material is large (eq.27) the material must move more or less in Kepler orbits; i.e., it must spiral in. This is an essential point in the theory.
If $`j(A)`$ is less than $`j_K`$ at $`R_{\mathrm{SN}}`$, the initial motion will have a substantial radial component in addition to the tangential one. But as the Kepler one decreases, cf. eq.29, there will come a point of $`r`$ at which $`j_K=j(A)`$. At this point an accretion disk will form, consisting of SN material spiraling in toward the neutron star. The primary motion is circular, but viscosity will provide a radial component inward
$`v_r\alpha v_K`$ (31)
where $`\alpha `$ is the viscosity parameter. It has been argued by Brandenburg et al. (1996) that $`\alpha 0.1`$ in the presence of equipartition magnetic fields perpendicular to the disk, and it may be even larger with the high magnetic fields required for GRBs. Narayan & Yi (1994) have given analytical solutions for such accretion disks. The material will arrive at the neutron star essentially tangentially, and therefore its high angular momentum will spin up the neutron star substantially. Accretion will soon make the neutron star collapse into a black hole. The angular momentum will be conserved, so the angular velocity is increased since the black hole has smaller radius than the neutron star. Thus the black hole is born with considerable spin.
A large fraction of the material of the failed supernova will accrete onto the black hole, giving it a mass of order $`7M_{}`$. All this material adds to the angular momentum of the black hole since all of it has the Kepler velocity at the black hole radius. Our estimates show that the black hole would be close to an extreme Kerr hole (Section 5), were it to accrete all of this material. It may, however, be so energetic that it drives off part of the envelope in the explosion before it can all accrete (see Section 5).
### 7.3 Soft X-ray Transients with Black-Hole Primaries
Nine binaries have been observed which are black-hole X-ray transients. All contain a high-mass black hole, of mass $`7M_{}`$. In seven cases the lower-mass companion (star B) has a mass $`\text{ }<M_{}`$. The two stars are close together, their distance being of order $`5R_{}`$. Star B fills its Roche Lobe, so it spills over some material onto the black hole. The accretion disk near the black hole emits soft X rays. Two of the companions are subgiants, filling their Roche lobes at a few times larger separations from the black hole.
In fact, however, the accretion onto the central object is not constant, so there is usually no X-ray emission. Instead, the material forms an accretion disk around the black hole, and only when enough material has been assembled, it falls onto the black hole to give observable X rays. Hence, the X-ray source is transient. Recent observation of a large space velocity of Cygnus X-1 (Nelemans et al. 1999) suggests that it has evolved similarly to the transient sources, with the difference that the companion to the black hole is an $`18M_{}`$ O star. The latter pours enough matter onto the accretion disk so that Cyg X-1 shines continuously. We plan to describe the evolution of Cyg X-1 in a future paper (Brown et al. 2000).
Table 1 is an abbreviated list of data on transient sources. A more complete table is given in Brown et al. (1999b). Two of the steady X-ray sources, in the LMC, have been omitted, because we believe the LMC to be somewhat special because of its low metallicity; also masses, etc., of these two are not as well measured. Of the others, 6 are main-sequence K stars, one is main-sequence M, and the other two have masses greater than the Sun. The masses given are geometric means of the maximum and minimum masses given by the observers. The distance $`a`$ between the black hole and the optical (visible) star is greater for the heavier stars than for the K- and M stars (except the more evolved one of them) as was expected in Section 7.1 for the spiraling in of star B. The table also gives the radius of the Roche Lobe and the specific orbital angular momentum of star B.
Five K stars have almost identical distance $`a5R_{}`$, and also Roche Lobe sizes, $`1.0R_{}`$. These Roche Lobes can be filled by K stars on the main sequence. The same is true for the M star. Together, K and M stars cover the mass range from 0.3 to 1$`M_{}`$. The two heavier stars have Roche Lobes of 3 and $`5R_{}`$ which cannot possibly be filled by main-sequence stars of mass $`2M_{}`$. We must therefore assume that these stars are subgiants, in the Herzsprung gap. These stars spend only about 1% of their life as subgiants, so we must expect that there are many “silent” binaries in which the $`2M_{}`$ companion has not yet evolved off the main sequence and sits well within its Roche lobe, roughly 100 times more. The time as subgiants is even shorter for more massive stars; this explains their absence among the transient sources.
Therefore we expect a large number of “silent partners”: stars of more than $`1M_{}`$, still on their main sequence, which are far from filling their Roche Lobe and therefore do not transfer mass to their black hole partners. In fact, we do not see any reason why the companion of the black hole could not have any mass, up to the ZAMS mass of the progenitor of the black hole; it must only evolve following the formation of the black hole. It then crosses the Herzsprung gap in such a short time, less than the thermal time scale, that star A cannot accept the mass from the companion, so that common envelope evolution must ensue. If we include these ‘silent partners’ in the birth rate, assuming a flat mass ratio distribution, we enhance the total birth rate of black-hole binaries by a factor 25 over the calculations by Brown, Lee, & Bethe (1999).
On the lower mass end of the companions, there is only one M star. This is explained in terms of the model of Section 7.1 by the fact that stars of low mass will generally spiral into the He core of star A, and will coalesce with A, see below eq.(23), so no relic is left. (Since the core is left spinning rapidly, these complete merger cases could also be suitable GRB progenitors.) As the outcome of the spiral-in depends also on other factors, such as the initial orbital separation and the primary mass, one may still have an occasional survival of an M star binary (note that the one M star companion is M0, very nearly in the K star range).
The appearance of the black hole transient X-ray binaries is much like our expectation of the relic of the binary which has made a hypernova: a black hole of substantial mass, and an ordinary star, possibly somewhat evolved, of smaller mass. We expect that star B would stop at a distance $`a_f`$ from star A which is greater if the mass of B is greater (see Section 7.1). This is just what we see in the black-hole binaries: the more massive companion stars ($`2M_{}`$) are further from the black hole than the K stars. We also note that the estimated birth rate of these binaries is high enough for them to be the progenitors of GRB, even if only in a modest fraction of them the conditions for GRB powering are achieved.
### 7.4 Nova Scorpii 1994 (GRO J1655-40)
Nova Sco 1994 is a black hole transient X-ray source. It consists of a black hole of $`7M_{}`$ and a subgiant of about $`2M_{}`$. Their separation is $`17R_{}`$. Israelian et al. (1999) have analyzed the spectrum of the subgiant and have found that the $`\alpha `$-particle nuclei O, Mg, Si and S have abundances 6 to 10 times the solar value. This indicates that the subgiant has been enriched by the ejecta from a supernova explosion; specifically, that some of the ejecta of the supernova which preceded the present Nova Sco (a long time ago) were intercepted by star B, the present subgiant. Israelian et al. (1999) estimate an age since accretion started from the assumption that enrichment has only affected the outer layers of the star. We here reconsider this: the time that passed since the explosion of the progenitor of the black hole is roughly the main-sequence lifetime of the present subgiant companion, which given its mass of $``$2$`M_{}`$ will be about 1 Gyr. This is so much longer than any plausible mixing time in the companion that the captured supernova ejecta must by now be uniformly mixed into the bulk of the companion. This rather increases the amount of ejecta that we require the companion to have captured. (Note that the accretion rate in this binary is rather less than expected from a subgiant donor, though the orbital period leaves no doubt that the donor is more extended than a main-sequence star (Regős, Tout, and Wickramasinghe 1998). It is conceivable that the high metal abundance has resulted in a highly non-standard evolution of this star, in which case one might have to reconsider its age.)
The presence of large amounts of S is particularly significant. Nomoto et al. (2000) have calculated the composition of a hypernova from an $`11M_{}`$ CO core, see Fig. 4. This shows substantial abundance of S in the ejecta. Ordinary supernovae produce little of this element, as shown by the results of Nomoto et al. (2000) in Fig. 4. The large amount of S, as well as O, Mg and Si we consider the strongest argument for considering Nova Sco 1994 as a relic of a hypernova, and for our model, generally.
Fig. 4 also shows that <sup>56</sup>Ni and <sup>52</sup>Fe are confined to the inner part of the hypernova, and if the cut between black hole and ejecta is at about $`5M_{}`$, there will be no Fe-type elements in the ejecta, as observed in Nova Scorpii 1994. By contrast hypernova 1998bw shows a large amount of Ni, indicating that in this case the cut was at a lower included mass.
The massive star A in Nova Sco will have gone through a hypernova explosion when the F-star B was still on the main sequence, its radius about $`1.5R_{}`$. Since the explosion caused an expansion of the orbit, the orbital separation $`a`$ was smaller at the time of the supernova than it is now, roughly by a factor
$`a_{\mathrm{then}}=a_{\mathrm{now}}/(1+\mathrm{\Delta }M/M_{\mathrm{now}}).`$ (32)
($`\mathrm{\Delta }M`$ is the mass lost in the explosion; see, e.g., Verbunt, Wijers, and Burm 1990). With $`\mathrm{\Delta }M0.8M_{\mathrm{now}}`$, as required by the high space velocity, this means $`a_{\mathrm{then}}=10R_{}`$. Therefore the fraction of solid angle subtended by the companion at the time of explosion was
$`{\displaystyle \frac{\mathrm{\Omega }}{4\pi }}={\displaystyle \frac{\pi (1.5R_{})^2}{4\pi (10R_{})^2}}6\times 10^3.`$ (33)
Assuming the ejecta of the hypernova to have been at least $`5M_{}`$ (Nelemans et al. 1999), the amount deposited on star B was
$`M_D\text{ }>0.03M_{}.`$ (34)
The solar abundance of oxygen is about 0.01 by mass, so with the abundance in the F star being 10 times solar, and oxygen uniformly mixed, we expect $`0.1\times 2.5=0.25M_{}`$ of oxygen to have been deposited on the companion, much more than the total mass it could have captured from a spherically symmetry supernova. \[Si/O\] is 0.09 by mass in the Sun, and \[S/O\] is 0.05, so since the over-abundances of all three elements are similar we expect those ratios to hold here, giving about 0.02$`M_{}`$ of captured Si and 0.01$`M_{}`$ of captured S. We therefore need a layer of stellar ejecta to have been captured which has twice as much Si as S, at the same time as having about 10 times more O. From fig. 4, we see that this occurs nowhere in a normal supernova, but does happen in the hypernova model of Nomoto et al. (2000) at mass cuts of 6$`M_{}`$ or more. This agrees very nicely with the notion that a hypernova took place in this system, and that the inner 7$`M_{}`$ or so went into a black hole.
What remains is to explain how the companion acquired ten times more mass than the spherical supernova model allows, and once again we believe that the answer is given in recent hypernova calculations (MacFadyen and Woosley 1999, Wheeler et al. 2000): hypernovae are powered by jet flows, which means they are very asymmetric, with mass outflow along the poles being much faster and more energetic than along the equator. The disk provides a source for easily captured material in two ways: First, it concentrates mass in the equatorial plane, which will later be ejected mostly in that plane. Second, the velocity acquired by the ejecta is of the order of the propagation speed of the shock through it. This propagation speed is proportional to $`\sqrt{P_2/\rho _1}`$, where $`P_2`$ is the pressure behind the shock and $`\rho _1`$ the density ahead of it. The driving pressure will be similar in all directions (or larger, due to the jet injection, in the polar regions), whereas the disk density is much higher than the polar density. Hence, the equatorial ejecta will be considerably slower than even normal supernova ejecta, greatly increasing the possibility of their capture by the companion. Other significant effects of the disk/jet geometry are (1) that the companion is shielded from ablation of its outer layers by fast ejecta, which is thought to occur in spherical supernovae with companion stars (Marietta, Burrows & Fryxell 2000) and (2) that there is no iron enrichment of the companion, because the iron —originating closest to the center— is either all captured by the black hole or ejected mainly in the jet, thus not getting near the companion (Wheeler et al. 2000; note that indeed no overabundance of Fe is seen in the companion of GRO J1655$``$40).
For the companion to capture the required 0.2–0.3$`M_{}`$ of ejecta it is sufficient that the ejecta be slow enough to become gravitationally bound to it. However, the material may not stay on: when the companion has so much mass added on a dynamical time scale it will be pushed out of thermal equilibrium, and respond by expanding, as do main-sequence stars that accrete mass more gradually on a time scale faster than their thermal time scale (e.g., Kippenhahn & Meyer-Hofmeister 1977). During this expansion, which happens on a time scale much longer than the explosion, the star may expand beyond its Roche lobe and transfer some of its mass to the newly formed black hole. However, because the dense ejecta mix into the envelope on a time scale between dynamical and thermal, i.e., faster than the expansion time, this back transfer will not result in the bulk of the ejecta being fed back, though probably the material lost is still richer in heavy elements than the companion is now. Since the outer layers of the star are not very dense, and the mass transfer is not unstable because the black hole is much more massive than the companion, the total amount of mass transferred back is probably not dramatic. However, the expansion does imply that the pre-explosion mass of the companion was somewhat higher than its present mass, and that the amount of ejecta that needs to be captured in order to explain the abundances observed today is also somewhat higher than the present mass of heavy elements in the companion.
A further piece of evidence that may link Nova Sco 1994 to our GRB/hypernova scenario are the indications that the black hole in this binary is spinning rapidly. Zhang, Cui, & Chen (1997) argue from the strength of the ultra-soft X–ray component that the black hole is spinning near the maximum rate for a Kerr black hole. However, studies by Sobczak et al. (1999) show that it must be spinning with less than 70% maximum. Gruzinov (1999) finds the inferred black hole spin to be about 60% of maximal from the 300 Hz QPO. Our estimates of the last section indicate that enough rotational energy will be left in the black hole so that it will still be rapidly spinning.
We have already mentioned the unusually high space velocity of $`150\pm 19`$ km s<sup>-1</sup>. Its origin was first discussed by Brandt et al. (1995), who concluded that significant mass must have been loss in the formation of the black hole in order to explain this high space velocity: it is not likely to acquire a substantial velocity in its own original frame of reference, partly because of the large mass of the black hole. But the mass lost in the supernova explosion is ejected from a moving object and thus carries net momentum. Therefore, momentum conservation demands that the center of mass of the binary acquire a velocity; this is the Blaauw–Boersma kick (Blaauw 1961, Boersma 1961). Note that the F-star companion mass is the largest among the black-hole transient sources, so the center of mass is furthest from the black hole and one would expect the greatest kick. Nelemans et al. (1999) estimate the mass loss in this kick to be $`510M_{}`$.
In view of the above, we consider it well established that Nova Sco 1994 is the relic of a hypernova. We believe it highly likely that the other black-hole transient X-ray sources are also hypernova remnants. We believe it likely that the hypernova explosion was accompanied by a GRB if, as in GRB 980326, the energy was delivered in a few seconds. It is not clear what will happen if the magnetic fields are so low that the power is delivered only over a much longer time. There could then still be intense power input for a few seconds due to neutrino annihilation deposition near the black hole (Janka et al. 1999), but that may not be enough for the jet to pierce through the He star and cause a proper GRB (MacFadyen and Woosley 1999). At this point, we recall that the GRB associated with SN1998bw was very sub-luminous, $`10^5`$ times lower than most other GRB. While it has been suggested that this is due to us seeing the jet sideways, it is in our view more likely that the event was more or less spherical (Kulkarni et al. 1998) and we see a truly lower-power event. A good candidate would be the original suggestion by Colgate (1968, 1974) of supernova shock break-out producing some gamma rays. Indications are that the expansion in SN1998bw was mildly relativistic (Kulkarni et al. 1998) or just sub-relativistic (Waxman and Loeb 1999). In either case, what we may have witnessed is a natural intermediate event in our scenario: we posit that there is a continuum of events varying from normal supernovae, delivering 1 foe more or less spherically in ten seconds, to extreme hypernovae/GRB that deliver 100 foes in a highly directed beam. In the middle, there will be cases where the beam cannot pierce through the star, but the total energy delivered is well above a supernova, with as net result a hypernova accompanied by a very weak GRB.
### 7.5 Numbers
Nearly all observed black hole transient X-ray sources are within 5 kpc of the Sun. Extrapolating to the entire Galaxy, a total of 8,800 black-hole transients with main-sequence K companions has been suggested (Brown, Lee, & Bethe 1999).
The lifetime of a K star in a black hole transient X-ray source is estimated to be $`10^{10}`$ yr (Van Paradijs 1996) but we shall employ $`10^9`$ yr for the average of the K-stars and the more massive stars, chiefly those in the “silent partners”. In this case the birth rate of the observed transient sources would be
$`\lambda _K=10^4/10^9=10^5\mathrm{per}\mathrm{galaxy}\mathrm{yr}^1.`$ (35)
We see no reason why low-mass companions should be preferred, so we assume that the formation rate of binaries should be independent of the ratio
$`q=M_{B,i}/M_{A,i}.`$ (36)
In other discussions of binaries, e.g., in Portegies Zwart & Yungelson (1998), it has often been assumed that the distribution is uniform in $`q`$. This is plausible but there is no proof. Since all primary masses $`M_A`$ are in a narrow interval, 20 to $`35M_{}`$, this means that $`M_B`$ is uniformly distributed between zero and some average $`M_A`$, let us say $`25M_{}`$. Then the total rate of creation of binaries of our type is
$`\lambda ={\displaystyle \frac{25}{0.7}}\lambda _K=3\times 10^4\mathrm{galaxy}^1\mathrm{yr}^1.`$ (37)
This is close to the rate of mergers of low mass black holes with neutron stars which Bethe & Brown (1998) have estimated to be
$`\lambda _m2\times 10^4\mathrm{galaxy}^1\mathrm{yr}^1.`$ (38)
These mergers have been associated speculatively with short GRBs, while formation of our binaries is supposed to lead to “long” GRBs (Fryer, Woosley, & Hartmann 1999). We conclude that the two types of GRB should be equally frequent, which is not inconsistent with observations. In absolute number both of our estimates eqs. (37) and (38) are substantially larger than the observed rate of $`10^7`$ galaxy<sup>-1</sup> yr<sup>-1</sup> (Wijers et al. 1998); this is natural, since substantial beaming is expected in GRBs produced by the Blandford-Znajek mechanism. Although we feel our mechanism to be fairly general, it may be that the magnetic field required to deliver the BZ energy within a suitable time occurs in only a fraction of the He cores.
## 8 Discussion and Conclusion
Our work here has been based on the Blandford-Znajek mechanism of extracting rotational energies of black holes spun up by accreting matter from a helium star. We present it using the simple circuitry of “The Membrane Paradigm” (Thorne et al. 1986). Energy delivered into the loading region up the rotational axis of the black hole is used to power a GRB. The energy delivered into the accretion disk powers a SN Ib explosion.
We also discussed black-hole transient sources, high-mass black holes with low-mass companions, as possible relics for both GRBs and Type Ib supernova explosions, since there are indications that they underwent mass loss in a supernova explosion. In Nova Sco 1994 there is evidence from the atmosphere of the companion star that a very powerful supernova explosion (‘hypernova’) occurred.
We estimate the progenitors of transient sources to be formed at a rate of 300 GEM (Galactic Events per Megayear). Since this is much greater than the observed rate of GRBs, there must be strong collimation and possible selection of high magnetic fields in order to explain the discrepancy.
We believe that there are strong reasons that a GRB must be associated with a black hole, at least those of duration several seconds or more discussed here. Firstly, neutrinos can deliver energy from a stellar collapse for at most a few seconds, and sufficient power for at most a second or two. Our quantitative estimates show that the rotating black hole can easily supply the energy as it is braked, provided the ambient magnetic field is sufficiently strong. The black hole also solves the baryon pollution problem: we need the ejecta that give rise to the GRB to be accelerated to a Lorentz factor of 100 or more, whereas the natural scale for any particle near a black hole is less than its mass. Consequently, we have a distillation problem of taking all the energy released and putting it into a small fraction of the total mass. The use of a Poynting flux from a black hole in a magnetic field (Blandford & Znajek 1977) does not require the presence of much mass, and uses the rotation energy of the black hole, so it provides naturally clean power.
Of course, nature is extremely inventive, and we do not claim that all GRBs will fit into the framework outlined here. We would not expect to see all of the highly beamed jets following from the BZ mechanism head on, the jets may encounter some remaining hydrogen envelope in some cases, jets from lower magnetic fields than we have considered here may be much weaker and delivered over longer times, etc., so we speculate that a continuum of phenomena may exist between normal supernovae and extreme hypernovae/GRBs. This is why we call our effort “A Theory of Gamma Ray Bursts” and hope that it will be a preliminary attempt towards systematizing the main features of the energetic bursts.
We would like to thank Stan Woosley for much useful information. Several conversations with Roger Blandford made it possible for us to greatly improve our paper, as did valuable comments from Norbert Langer. This work is partially supported by the U.S. Department of Energy Grant No. DE-FG02-88ER40388. HKL is supported also in part by KOSEF Grant No. 1999-2-112-003-5 and by the BK21 program of the Korean Ministry of Education.
## Appendix A Estimates of $`ϵ_\mathrm{\Omega }=\mathrm{\Omega }_{disk}/\mathrm{\Omega }_H`$
We collect here useful formulas needed to calculate $`ϵ_\mathrm{\Omega }=\mathrm{\Omega }_{disk}/\mathrm{\Omega }_H`$. First of all
$`\mathrm{\Omega }_H`$ $`=`$ $`{\displaystyle \frac{\stackrel{~}{a}}{1+\sqrt{1\stackrel{~}{a}^2}}}\left({\displaystyle \frac{c^3}{2MG}}\right)`$ (A1)
$`=`$ $`{\displaystyle \frac{\sqrt{2}\stackrel{~}{a}}{1+\sqrt{1\stackrel{~}{a}^2}}}\left({\displaystyle \frac{R}{R_{Sch}}}\right)^{3/2}\mathrm{\Omega }_K`$
$`\mathrm{\Omega }_{disk}`$ $`=`$ $`\omega \mathrm{\Omega }_K\left[1+\stackrel{~}{a}{\displaystyle \frac{GM}{c^2}}\sqrt{{\displaystyle \frac{GM}{c^2R^3}}}\right]^1`$ (A2)
$`=`$ $`\omega \mathrm{\Omega }_K\left[1+\stackrel{~}{a}\left({\displaystyle \frac{R_{Sch}}{2R}}\right)^{3/2}\right]^1`$
where $`\mathrm{\Omega }_K\sqrt{GM/R^3}`$ and $`\omega `$ is dimensionless parameter ($`0<\omega <1`$). Thus
$`{\displaystyle \frac{\mathrm{\Omega }_{disk}}{\mathrm{\Omega }_H}}`$ $`=`$ $`\omega {\displaystyle \frac{1+\sqrt{1\stackrel{~}{a}^2}}{\sqrt{2}\stackrel{~}{a}}}\left({\displaystyle \frac{R_{Sch}}{R}}\right)^{3/2}\left[1+\stackrel{~}{a}\left({\displaystyle \frac{R_{Sch}}{2R}}\right)^{3/2}\right]^1.`$ (A3)
The numerical estimates are summarized in Table 2 for various $`\omega `$ and radii.
## Appendix B Spin-up of Black Holes by Accretion
The specific angular momentum and energy of test particles in Keplerian circular motion, with rest mass $`\delta m`$, are
$`\stackrel{~}{E}`$ $``$ $`{\displaystyle \frac{E}{\delta m}}=c^2\left[{\displaystyle \frac{r^2R_{\mathrm{Sch}}r+a\sqrt{R_{\mathrm{Sch}}r/2}}{r(r^2\frac{3}{2}R_{\mathrm{Sch}}r+a\sqrt{2R_{\mathrm{Sch}}r})^{1/2}}}\right]`$
$`\stackrel{~}{l}`$ $``$ $`{\displaystyle \frac{l}{\delta m}}=c\sqrt{{\displaystyle \frac{R_{\mathrm{Sch}}r}{2}}}\left[{\displaystyle \frac{(r^2a\sqrt{2R_{\mathrm{Sch}}r}+a^2)}{r(r^2\frac{3}{2}R_{\mathrm{Sch}}r+a\sqrt{2R_{\mathrm{Sch}}r})^{1/2}}}\right]`$ (B1)
where $`R_{\mathrm{Sch}}=2GM/c^2`$ and BH spin $`a=J/Mc=\stackrel{~}{a}(GM/c^2)`$. The accretion of $`\delta m`$ changes the BH’s total mass and angular momentum by $`\mathrm{\Delta }M=\stackrel{~}{E}\delta m`$ and $`\mathrm{\Delta }J=\stackrel{~}{l}\delta m`$. The radii of marginally bound ($`r_{mb}`$) and stable ($`r_{ms})`$ orbits are given as
$`r_{mb}`$ $`=`$ $`R_{\mathrm{Sch}}a+\sqrt{R_{\mathrm{Sch}}(R_{\mathrm{Sch}}2a)}`$
$`r_{ms}`$ $`=`$ $`{\displaystyle \frac{R_{\mathrm{Sch}}}{2}}\left(3+Z_2\left[(3Z_1)(3+Z_1+2Z_2)\right]^{1/2}\right)`$
$`Z_1`$ $`=`$ $`1+\left(1{\displaystyle \frac{4a^2}{R_{\mathrm{Sch}}^2}}\right)^{1/3}\left[\left(1+{\displaystyle \frac{2a}{R_{\mathrm{Sch}}}}\right)^{1/3}+\left(1{\displaystyle \frac{2a}{R_{\mathrm{Sch}}}}\right)^{1/3}\right]`$
$`Z_2`$ $`=`$ $`\left(3{\displaystyle \frac{4a^2}{R_{\mathrm{Sch}}^2}}+Z_1^2\right)^{1/2}.`$ (B2)
The numerical values of the specific angular momentum and energy of test particles are summarized in Table 3 and Fig.5. In Fig.6, we test how much mass we need in order to spin up the non-rotating black hole up to given $`\stackrel{~}{a}`$. Note that the last stable orbit is almost Keplerian even with the accretion disk, and we assume 100% efficiency of angular momentum transfer from the last stable Keplerian orbit to BH. In order to spin-up the BH up to $`\stackrel{~}{a}=0.9`$, we need $`68\%`$ ($`52\%`$) of original non-rotating BH mass in case of $`r_{lso}=r_{ms}`$ ($`r_{mb}`$). For a very rapidly rotating BH with $`\stackrel{~}{a}=0.99`$, we need $`122\%`$ and 82%, respectively. For $`r_{lso}=r_{ms}`$, there is an upper limit, $`\stackrel{~}{a}=0.998`$, which can be obtained by accretion (Thorne 1974). In the limit where $`r_{lso}=r_{mb}`$, however, spin-up beyond this limit is possible because the photons can be captured inside thick accretion disk, finally into BH (Abramowicz et al. 1988).
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# Contents
## 1 Introduction
### 1.1 Preamble
The item of the low energy behavior of a strong interaction attracts more and more interest along with the further experimental data accumulation. In the perturbative quantum chromodynamics (pQCD) this behavior is spoiled by unphysical singularities lying at the three flavour region $`(f=3)`$ and associated with the scale parameter $`\mathrm{\Lambda }_{f=3}350`$ MeV. In the “small energy” and “small momentum transfer” regions $`(\sqrt{s},Q\sqrt{Q^2}3\mathrm{\Lambda })`$ these singularities complicate theoretical interpretation of data. On the other hand, their existence contradicts some general statements of the local QFT.
Meanwhile, this issue has a rather elegant solution. As it has been shown (see, also fresh review ), by combining three elements:
1. Usual Feynman perturbation theory for effective coupling(s) and observables,
2. Renormalizability, i.e., renormalization–group (RG) invariance, and
3. General principles of local QFT — like causality, unitarity, Poincaré invariance and spectrality — in the form of spectral representations of Källen–Lehmann and Jost–Lehmann–Dyson type
it turns out to be possible to formulate an Invariant Analytic Approach (IAA) for the pQCD invariant coupling and observables in which the central theoretical object is a spectral density. Being calculated by usual, RG–improved, perturbation theory it defines and relates $`Q^2`$–analytic, RG-invariant expressions in the both Euclidean and Minkowskian channels.
The IAA obeys several remarkable properties:
— It enables one to obtain modified perturbation expressions for observables, free of unphysical singularities, poles and cuts, with behavior correlated in both space-like and time-like domains.
— In particular, the IAA results in modified ghost-free expressions for invariant QCD coupling $`\alpha _{\mathrm{an}}(Q^2;f)`$ and $`\stackrel{~}{\alpha }(s;f)`$ which obey reduced higher–loops and renormalization–scheme sensitivity . See, Fig.1.
— Then, it yields changing the structure of perturbation expansion for observables: instead of common power series, as a result of its integral transformation, there appear non-power asymptotic series à la Erdélyi over the sets of specific functions $`𝒜_k(Q^2;f)`$ and $`𝔄_k(s;f)`$. These functions are defined via integral transformations of related powers $`\alpha _s^k(Q^2;f)`$ in terms of relevant spectral densities. At small and moderate argument values, they diminish with the $`k`$ growth much quicker than the corresponding powers $`\alpha _{\mathrm{an}}^k(Q^2;f)`$ and $`\stackrel{~}{\alpha }^k(s;f)`$, (and even oscillate in the region $`\sqrt{s},Q\mathrm{\Lambda }`$) thus improving essentially the convergence of perturbation expansion for observables.
We review all these IAA features, important for our further developments, in the second part of Section 1.
The first purpose of this work is to elucidate relation between the Radyushkin–Krasnikov–Pivovarov “pipization” trick and the Solovtsov–Milton construction of effective $`s`$–channel QCD coupling within the IAA scheme.
In the course of this analysis — see Section 2 — we reveal a spectacular “distorting mirror” correlation between analyticized and pipizated invariant QCD coupling in space-like $`\alpha _{\mathrm{an}}(Q^2;f)`$ and time-like $`\stackrel{~}{\alpha }(s;f)`$ regions as well as between corresponding expansion functions $`𝒜_k(Q^2;f)`$ and $`𝔄_k(s;f)`$. See Fig.2.
Then, in Section 2.3, we consider an issue of transition across the heavy quark thresholds, for constructing a “global” picture valid in the whole physical region $`M_\tau \sqrt{s},QM_Z`$.
It should be noted, that all precedent papers Refs. dealt only with the massless quarks at fixed flavour number, $`f`$, case. This can be justified, to some extent, for analysis inside a narrow interval of relevant energy $`\sqrt{s}`$ or momentum transfer $`Q`$ values. Meanwhile, the ultimate goal of all the pQCD business is a correlation of QCD effective coupling values extracted from different experiments.
To construct the global invariant analytic couplings, one needs recipe of relating expressions with different $`f`$ values. For this goal, we use the same guideline as previously, at the “fixed $`f`$ case”, that is we start with adequately defined “global” spectral functions — see expression (14). Here, an essential point is the matching condition that relates $`\mathrm{\Lambda }_f`$ parameters for different fixed flavor number $`f`$ values. We use standard $`\overline{\mathrm{MS}}`$ prescription ascending to early 80s. This results in a smooth global Euclidean $`\alpha _{\mathrm{an}}(Q)`$, $`𝒜_k(Q^2)`$ and spline–continuous Minkowskian $`\stackrel{~}{\alpha }(s)`$, $`𝔄_k(s)`$ expressions.
In the concluding Section, using examples of inclusive $`\tau `$ decay, $`e^+e^{}`$ hadrons annihilation and sum rules, we shortly comment the possible implication of our new global scheme on perturbative analysis of QCD processes.
The main results of this work are summed concisely in the Subsection 3.3.
### 1.2 The $`s`$–channel: early attempts
As it is well known, the notion of invariant (or effective) coupling originally was introduced in the RG treatment of renormalizable QFT. In the RG formalism, invariant coupling function $`\overline{\alpha }`$ was defined as a product of propagator and vertex amplitudes initially related with a product of real finite Dyson’s renormalization constants. This construction is valid only in the space-like domain <sup>1</sup><sup>1</sup>1 Physically, in QED it can be considered as a Fourier transform of spatial electron charge distribution first discussed by Dirac . and can be directly used for analysis of corresponding observables. However, the RG formalism does not provide us with analogous object in time-like region.
It is worth noting that sporadic attempts to define the effective coupling $`\alpha (s)`$ in the Minkowskian, time-like, domain were made in late 70s. Omitting an early simple–minded trick with “mirror reflection” of singular function
$$\alpha _s(Q^2;f)\alpha (s;f)|\alpha _s(s;f)|,$$
we mention here the practically simultaneous results of Radyushkin and Krasnikov and Pivovarov . In both the papers, the integral transformation $`\stackrel{~}{\alpha }(s;f)=𝐑[\overline{\alpha }_s(Q^2;f)]`$ reverse $`𝐑=𝐃^1`$ to “dipole representation” for the Adler function
$$D(Q^2)=\frac{Q^2}{\pi }_0^{\mathrm{}}\frac{ds}{(s+Q^2)^2}R(s)𝐃\left\{R(s)\right\}$$
(1)
in terms of an observable $`R(s)`$ in the time-like region, has been used.
In , as a starting point for observables in the Euclidean, i.e., space-like domain $`Q^2>0`$, the perturbation series
$$D_{\mathrm{pt}}(Q^2)=1+\underset{k1}{}d_k\overline{\alpha }_{s}^{}{}_{}{}^{k}(Q^2;f)$$
(2)
has been assumed. It contains powers of usual, RG summed, invariant coupling $`\overline{\alpha }_s(Q^2;f)`$ that obeys unphysical singularities in the infrared (IR) region around $`Q^2\mathrm{\Lambda }_3^2`$.
By using the reverse transformation
$$R(s)=\frac{i}{2\pi }_{si\epsilon }^{s+i\epsilon }\frac{dz}{z}D_{\mathrm{pt}}(z)𝐑\left[D_{\mathrm{pt}}(Q^2)\right]$$
(3)
these authors arrived at the “$`𝐑`$–transformed” expansion that, in our notation, reads
$$R_\pi (s)=1+\underset{k1}{}d_k𝔄_k(s;f);𝔄_k(s;f)=𝐑\left[\overline{\alpha }_{s}^{}{}_{}{}^{k}(Q^2;f)\right].$$
(4)
For example
$$𝐑\left[\frac{1}{l}\right]=\frac{1}{2}\frac{1}{\pi }\mathrm{arctan}\frac{L}{\pi };\text{with}l=\mathrm{ln}\frac{Q^2}{\mathrm{\Lambda }^2};L=\mathrm{ln}\frac{s}{\mathrm{\Lambda }^2},$$
$$𝐑\left[\frac{\mathrm{ln}l}{l^2}\right]=\frac{\mathrm{ln}\left[\sqrt{L^2+\pi ^2}\right]+1L𝐑\left[1/l\right]}{L^2+\pi ^2};𝐑\left[\frac{1}{l^2}\right]=\frac{1}{L^2+\pi ^2}.$$
This yields<sup>2</sup><sup>2</sup>2This expression we give in the form equivalent to that one used in . In papers it was given in an another form, non-adequate at $`L0`$. See, also .
$$\stackrel{~}{\alpha }^{(1)}(s;f)=\frac{1}{\beta _0}\left[\frac{1}{2}\frac{1}{\pi }\mathrm{arctan}\frac{L}{\pi }\right];\beta _0=\frac{332f}{12\pi }.$$
(5)
At the two–loop iterative case with
$$\beta _{[f]}\overline{\alpha }_{s}^{}{}_{}{}^{(2)}(Q^2;f)=\frac{1}{l}b_f\frac{\mathrm{ln}l}{l^2},\beta _{[f]}\beta _0;\beta _1=\frac{10238f}{12\pi },b_f=\frac{\beta _1}{\beta _0^2},$$
by combining $`𝐑\left[1/l\right]b_f𝐑[\mathrm{ln}l/l^2]`$ one obtains explicit expression for the “iterative” two-loop effective $`s`$–channel coupling $`\stackrel{~}{\alpha }^{(2)}(s;f)=𝔄_1^{(2)}(s;f),`$
$$\stackrel{~}{\alpha }_{iter}^{(2)}(s;f)=\left(1\frac{b_fL}{L^2+\pi ^2}\right)\stackrel{~}{\alpha }^{(1)}(s;f)+\frac{b_f}{\beta _{[f]}}\frac{\mathrm{ln}\left[\sqrt{L^2+\pi ^2}\right]+1}{L^2+\pi ^2}.$$
Obtained $`\stackrel{~}{\alpha }^{(1)}`$ and $`\stackrel{~}{\alpha }_{\mathrm{iter}}^{(2)}`$ are monotonous functions with finite IR limit, free of $`\mathrm{\Lambda }`$–singularity which is “screened” by resummed “$`\pi ^2`$–terms”. Non-singular expressions for higher functions $`𝔄_k`$ could be constructed in the same way.
The positive feature of this construction was an automatic summation of the so–called “$`\pi ^2`$ – terms” that “screen” unphysical singularities and observed property
$$\left(𝐑\left[\overline{\alpha }_{s}^{}{}_{}{}^{k+1}\right]\right)^{1/(k+1)}<\left(𝐑\left[\overline{\alpha }_{s}^{}{}_{}{}^{k}\right]\right)^{1/k}$$
that improves the convergence of perturbation series.
However, there was one essential drawback. The dipole transformation (1), that is supposed to be reverse to $`𝐑`$, being applied to (4) does not return us to the input (2)
$$𝐃\left\{R_\pi (s)\right\}=𝐃\left\{𝐑\left[D_{\mathrm{pt}}\right]\right\}D_{\mathrm{pt}}(Q^2)𝐃𝐑𝐈,$$
as far as the unphysical singularities of $`\overline{\alpha }_s(Q^2;f)`$ and of its powers are incompatible with analytic properties in the complex $`Q^2`$ plane of the integral in the r.h.s. of (1).
Resolution of this issue came 15 years later with the IAA. The “missing link” is the analyticization transformation.
### 1.3 Analyticization in the $`Q^2`$–channel
Operation
$$F(Q^2)F_{\mathrm{an}}(Q^2)=𝐀F(Q^2)$$
has been introduced in terms of the Källén–Lehmann representation and correlates with analytic properties of the Adler function contained in eq.(1).
Generally, this transformation is defined for a function $`F`$ that should be analytic in the $`Q^2`$ plane with a cut along the negative part of the real axis. In our case, this function could be either invariant coupling $`\overline{\alpha }_s`$ itself <sup>3</sup><sup>3</sup>3As it has been explained in detail in the first papers on the IAA, the QCD invariant coupling, according to general properties of local QFT, should satisfy the Källén–Lehmann spectral representation. For the original analysis of this issue see Ref. . or its power, or some series in its powers.
Operation $`𝐀`$ consists of two elements:
– use the Källen–Lehmann representation
$$F_{\mathrm{an}}(Q^2)=\frac{1}{\pi }\underset{0}{\overset{\mathrm{}}{}}\frac{d\sigma }{\sigma +Q^2}\rho _{\mathrm{pt}}(\sigma )\text{with}$$
– the spectral density defined via straightforward continuation of $`F`$ on the cut
$$\rho _{\mathrm{pt}}(\sigma )=\mathrm{}F(\sigma ).$$
A couple of comments are in order.
* Operation $`𝐀`$, being applied to the usual coupling<sup>4</sup><sup>4</sup>4For the time being, we consider the massless case with a fixed number $`f`$ of effective quark flavors in the $`\overline{\mathrm{MS}}`$ scheme. For the transition between the regions with different $`f`$ values, see Section 2.3. $`F=\overline{\alpha }_s(Q^2;f)`$, results in the analyticized coupling
$$\alpha _{\mathrm{an}}(Q^2;f)=\frac{1}{\pi }\underset{0}{\overset{\mathrm{}}{}}\frac{d\sigma }{\sigma +Q^2}\rho (\sigma ;f);\rho (\sigma ;f)=\mathrm{}\overline{\alpha }_s(\sigma ;f)$$
(6)
which is a smooth monotonic function free of unphysical singularities, with a finite value at the origin
$$\alpha _{\mathrm{an}}(0;f)=1/\beta _{[3]}1.4$$
which is remarkably independent (see, e.g., ) of higher loop contributions.
Here, $`\rho `$ is defined as an imaginary part of the usual, RG invariant, effective coupling $`\overline{\alpha }_s`$ continued on the physical cut.
* Operation $`𝐀`$, applied to power perturbation series (2) for an observable $`D_{\mathrm{pt}}(Q^2)`$, produces a non-power series
$$D_{\mathrm{an}}(Q^2;f)=1+\underset{k1}{}d_k𝒜_k(Q^2;f);\alpha _{\mathrm{an}}(Q^2;f)=𝒜_1(Q^2;f)$$
(7)
with
$$𝒜_k(x;f)=\frac{1}{\pi }\underset{0}{\overset{\mathrm{}}{}}\frac{d\sigma }{\sigma +x}\rho _k(\sigma ;f);\rho _k(\sigma ;f)=\mathrm{}\left[\overline{\alpha }_{s}^{}{}_{}{}^{k}(\sigma ;f)\right].$$
(8)
For example,
$$𝐀\left[\frac{1}{l}\right]=\frac{1}{l}+\frac{1}{1e^l},𝐀\left[\frac{1}{l^2}\right]=\frac{1}{l^2}+\frac{e^l}{(1e^l)^2},\mathrm{}$$
that is
$$𝐀\alpha _s^{[1]}(Q^2;f)=\alpha _{\mathrm{an}}^{[1]}(Q^2;f)=\frac{1}{\beta _0}\left[\frac{1}{\mathrm{ln}(Q^2/\mathrm{\Lambda }^2}+\frac{\mathrm{\Lambda }^2}{\mathrm{\Lambda }^2Q^2}\right]$$
(9)
and so on.
Here, in the invariant analytic coupling $`\alpha _{\mathrm{an}}`$ , the $`\mathrm{\Lambda }`$–pole is compensated by power term containing the non-perturbative $`Q^2/\mathrm{\Lambda }^2=(Q^2/\mu ^2)\mathrm{exp}(1/\beta _0\alpha _\mu )`$ structure.
Properties of the invariant analytic functions $`𝒜_k`$, free of ghost troubles, and non-power expansion (7) have been discussed in papers . They are quite similar to those for $`𝔄_k`$ and expansion (4) — see below.
### 1.4 Summary of the IAA
Here, we repeat in brief basic definitions of the Invariant Analytic Approach.
First, one has to transform the usual singular invariant coupling
$$\overline{\alpha }_s(Q^2;f)𝐀\overline{\alpha }_s(Q^2;f)=\alpha _{\mathrm{an}}(Q^2;f)$$
into the analyticized one, free of ghost singularities in the space-like region.
Second, with the help of the operation $`𝐑`$, one defines invariant coupling $`\stackrel{~}{\alpha }(s;f)`$ in the time-like domain
$$\alpha _{\mathrm{an}}(Q^2;f)\stackrel{~}{\alpha }(s;f)=𝐑\left[\alpha _{\mathrm{an}}\right]=\underset{s}{\overset{\mathrm{}}{}}\frac{d\sigma }{\sigma }\rho (\sigma ;f)$$
(10)
with spectral density $`\rho `$ defined in (6).
Here, we have a possibility of reconstructing the Euclidean, $`Q^2`$–channel, invariant coupling $`\alpha _{\mathrm{an}}(Q^2;f)`$ from the Minkowskian, $`s`$–channel, one $`\stackrel{~}{\alpha }(s;f)`$ by the dipole transformation
$$\alpha _{\mathrm{an}}(Q^2;f)=\frac{Q^2}{\pi }_0^{\mathrm{}}\frac{ds}{(s+Q^2)^2}\stackrel{~}{\alpha }(s;f)𝐃\left\{\stackrel{~}{\alpha }(s;f)\right\}.$$
(11)
For instance, substituting $`\stackrel{~}{\alpha }^{(1)}(s;f)`$ into the integrand, one obtains after integration by parts
$$\frac{Q^2}{\pi \beta _0}_0^{\mathrm{}}\frac{d\sigma }{(\sigma +Q^2)^2}\left(\frac{1}{2}\frac{1}{\pi }\mathrm{arctan}\frac{\mathrm{ln}(\sigma /\mathrm{\Lambda }^2)}{\pi }\right)=$$
$$\frac{Q^2}{\pi \beta _0}_0^{\mathrm{}}\frac{d\sigma }{(\sigma +Q^2)}\frac{1}{\mathrm{ln}^2(\sigma /\mathrm{\Lambda }^2)+\pi ^2}=\alpha _{\mathrm{an}}^{(1)}(Q^2,f)$$
precisely in the form (9). This simple calculation elucidates the connection between the ghost–free expressions both in the $`s`$– and $`Q^2`$–channels. They are connected also by the reverse transformation $`\stackrel{~}{\alpha }^{(1)}(s;f)=𝐑\left[\alpha _{\mathrm{an}}^{}{}_{}{}^{(1)}(Q^2;f)\right].`$
On the Fig. 1 we give a concise summary of the IAA results for invariant analytic couplings $`\alpha _{\mathrm{an}}(Q^2,3)`$ and $`\stackrel{~}{\alpha }(s,3)`$ calculated for one– , two– and three–loop cases in both the Euclidean and Minkowskian domains.
Here, dash–dotted curves represent one-loop IAA approximations (5) and (9). Solid IAA curves are based on exact two-loop solutions of RG equations<sup>5</sup><sup>5</sup>5 As it has been recently established, the exact solution to two-loop RG differential equation for the invariant coupling can be expressed in terms of a special function $`W`$, the Lambert function, defined by relation $`W(z)e^{W(z)}=z,`$ with an infinite number of branches $`W_n(z)`$. For some detail of analyticized solutions expressed in terms of Lambert function — see Refs. . and approximate three–loop solutions in the $`\overline{\mathrm{MS}}`$ scheme. Their remarkable coincidence (within the 1–2 per cent limit) demonstrates reduced sensitivity of the IAA with respect to the higher–loops effects in the whole Euclidean and Minkowskian regions from IR till UV limits.
For comparison, by dotted line we also give usual $`\overline{\alpha }_s(Q^2)`$ two-loop effective QCD coupling with a pole at $`Q^2=\mathrm{\Lambda }^2.`$
As it has been shown in , relations parallel to eqs.(10) and (11) are valid for powers of the pQCD invariant coupling. This can be resumed in the form of a self-consistent scheme.
## 2 Self-consistent scheme for observables
### 2.1 Relations between Euclidean and Minkowskian
First, one has to transform usual power perturbation series (2) of the $`Q^2`$ domain
$$𝐈.D_{\mathrm{pt}}(Q^2)D_{\mathrm{an}}(Q^2)=𝐀D_{\mathrm{pt}}(Q^2)$$
into the non-power one (7).
Second, with the help of the operation $`𝐑`$, one introduces
$$\mathrm{𝐈𝐈}.D_{\mathrm{an}}(Q^2)R_\pi (s)=𝐑\left[D_{\mathrm{an}}(Q^2)\right]$$
the s–channel non-power expansion $`R_\pi (s)`$ (4) with
$$𝔄_k(s)=\underset{s}{\overset{\mathrm{}}{}}\frac{d\sigma }{\sigma }\rho _k(\sigma );\rho _k(\sigma )=\mathrm{}\left[\alpha _s^k(\sigma )\right].$$
(12)
The third element is the closure of the scheme that is provided by the operation (11)
$$\mathrm{𝐈𝐈𝐈}.R_\pi (s)D_{\mathrm{an}}(Q^2)=𝐃\left\{R_\pi (s)\right\}$$
reverse to II.
In the other words, to enjoy self-consistency $`𝐑𝐃=𝐃𝐑=\mathrm{𝟏},`$ one should abandon completely the usual effective coupling $`\alpha _s(Q^2)`$ and power series $`D_{\mathrm{pt}}`$, eq.(2), applying operations $`𝐑`$ and $`𝐃=𝐑^1`$ only to IAA invariant couplings $`\alpha _{\mathrm{an}}`$ , $`\stackrel{~}{\alpha }`$ and to non-power expansions $`D_{\mathrm{an}}`$ and $`R_\pi `$.
### 2.2 Expansion of observables over non-power sets $`\left\{𝒜\right\}`$ and $`\left\{𝔄\right\}`$
To realize the effect of transition from expansion over the “traditional” power set
$$\left\{\overline{\alpha }_{s}^{}{}_{}{}^{k}(Q^2,f)\right\}=\overline{\alpha }_s(Q^2),\overline{\alpha }_{s}^{}{}_{}{}^{2},\mathrm{}\overline{\alpha }_{s}^{}{}_{}{}^{k}\mathrm{}$$
to expansions over non–power sets in the space-like and time-like domains
$$\left\{𝒜_k(Q^2,f)\right\}=\alpha _{\mathrm{an}}(Q^2,f),𝒜_2(Q^2,f),𝒜_3\mathrm{};\left\{𝔄_k(s,f)\right\}=\stackrel{~}{\alpha }(s,f),𝔄_2(s,f),𝔄_3\mathrm{},$$
it is instructive to learn properties of the latters.
In a sense, both non-power sets are similar:
— They consist of functions that are free of unphysical singularities.
— First functions, the new effective couplings, $`𝒜_1=\alpha _{\mathrm{an}}`$ and $`𝔄_1=\stackrel{~}{\alpha }`$ are monotonically decreasing. In the IR limit, they are finite and equal $`\alpha _{\mathrm{an}}(0,3)=\stackrel{~}{\alpha }(0,3)1.4`$ with the same infinite derivatives. Both have the same leading term $`1/\mathrm{ln}x`$ in the UV limit.
— All other functions (“effective coupling powers”) of both the sets start from the zero IR values $`𝒜_{k2}(0,f)=𝔄_{k2}(0,f)=0`$ and obey the UV behavior $`1/(\mathrm{ln}x)^k`$ corresponding to $`\overline{\alpha }_{s}^{}{}_{}{}^{k}(x)`$. They are no longer monotonous. The second functions $`𝒜_2`$ and $`𝔄_2`$ are positive with maximum around $`s,Q^2\mathrm{\Lambda }^2`$. Higher functions $`𝒜_{k3}`$ and $`𝔄_{k3}`$ oscillate in the region of low argument values and obey $`k2`$ zeroes.
Remarkably enough, the mechanism of liberation of unphysical singularities is quite different. While in the space-like domain it involves non-perturbative, power in $`Q^2`$, structures, in the time-like region it is based only upon resummation of the “$`\pi ^2`$ terms”. Figuratively, (non-perturbative !) analyticization in the $`Q^2`$–channel can be treated as a quantitatively distorted reflection (under $`Q^2s=Q^2`$) of (perturbative) “pipization” in the $`s`$–channel. This effect of “distorting mirror” first discussed in is illustrated on figures 1 and 2.
Summarize the main results essential for data analysis. Instead of power perturbative series in the space-like
$$D_{\mathrm{pt}}(Q^2)=1+d_{\mathrm{pt}}(Q^2);d_{\mathrm{pt}}(Q^2)=\underset{k1}{}d_k\overline{\alpha }_{s}^{}{}_{}{}^{k}(Q^2;f)(\text{2}a)$$
and time-like regions
$$R_{\mathrm{pt}}(s)=1+r_{\mathrm{pt}}(s);r_{\mathrm{pt}}(s)=\underset{k1}{}r_k\stackrel{~}{\alpha }^k(s;f);(r_{1,2}=d_{1,2},r_3=d_3d_1\frac{\pi ^2\beta _{[f]}^2}{3}),$$
one has to use asymptotic expansions (7) and (4)
$$d_{\mathrm{an}}(Q^2)=\underset{k1}{}d_k𝒜_k(Q^2,f);r_\pi (s)=\underset{k1}{}d_k𝔄_k(s,f)$$
with the same coefficients $`d_k`$ over non-power sets of functions $`\left\{𝒜\right\}`$ and $`\left\{𝔄\right\}`$.
### 2.3 Global formulation
To apply the new scheme for analysis of QCD processes, one has to formulate it “globally”, in the whole experimental domain, i.e., for regions with different values of a number $`f`$ of active quarks. For this goal, we revise the issue of the threshold crossing.
Threshold matching In a real calculation, the procedure of the threshold matching is in use. One of the simplest is the matching condition in the massless $`\overline{\mathrm{MS}}`$ scheme
$$\overline{\alpha }_s(Q^2=M_f^2;f1)=\overline{\alpha }_s(Q^2=M_f^2;f)$$
(13)
related to the mass squared $`M_f^2`$ of the f-th quark.
This condition allows one to define a “global” function $`\overline{\alpha }_s(Q^2)`$ consisting of the smooth parts
$$\overline{\alpha }_s(Q^2)=\overline{\alpha }_s(Q^2;f)\text{at}M_{f1}^2Q^2M_f^2$$
s and continuous in the whole space-like interval of positive $`Q^2`$ values with discontinuity of derivatives at the matching points. We call such a functions as the spline–continuous ones.
At the first sight, any massless matching, yielding the spline–type function, violates the analyticity in the $`Q^2`$ variable, thus disturbing the relation between the $`s`$– and $`Q^2`$–channels<sup>6</sup><sup>6</sup>6Any massless scheme is an approximation that can be controlled by the related mass–dependent scheme . Using such a scheme, one can devise a smooth transition across the heavy quark threshold. Nevertheless, from the practical point of view, it is sufficient (besides the case of data lying in close vicinity of the threshold) to use the spline–type matching (13) and forget about the smooth threshold crossing..
However, in the IAA, the original power perturbation series (2) with its unphysical singularities and possible threshold non-analyticity has no direct relation to data, being a sort of a “raw material” for defining spectral density. Meanwhile, the discontinuous density is not dangerous. Indeed, expression of the form
$$\rho _k(\sigma )=\rho _k(\sigma ;3)+\underset{f4}{}\theta (\sigma M_f^2)\left\{\rho _k(\sigma ;f)\rho _k(\sigma ;f1)\right\}$$
(14)
with $`\rho _k(\sigma ;f)=\mathrm{}\overline{\alpha }_{s}^{}{}_{}{}^{k}(\sigma ,f)`$ defines, according to (8) and (12), the smooth global
$$𝒜_k(Q^2)=\frac{1}{\pi }\underset{0}{\overset{\mathrm{}}{}}\frac{d\sigma }{\sigma +x}\rho _k(\sigma )$$
(15)
and spline–continuous global
$$𝔄_k(s)=\underset{s}{\overset{\mathrm{}}{}}\frac{d\sigma }{\sigma }\rho _k(\sigma )$$
(16)
functions <sup>7</sup><sup>7</sup>7Here, by eqs.(15),(16) and (14) we introduced new “global” effective invariant couplings and higher expansion functions different from the previous ones with fixed $`f`$ value..
We see that in this construction the role of the input perturbative invariant coupling $`\overline{\alpha }_s(Q^2)`$ is twofold. It provides us not only with spectral density (14) but with matching conditions (13) relating $`\mathrm{\Lambda }_f`$ with $`\mathrm{\Lambda }_{f+1}`$ as well.
Note that the matching condition (13) is tightly related to the renormalization procedure. Just for this profound reason we keep it untouched (compare with Ref. ).
The $`s`$-channel: shift constants As a practical result, we now observe that the “global” $`s`$–channel coupling $`\stackrel{~}{\alpha }(s)`$ and other functions $`𝔄_k(s)`$, generally, differs of effective coupling with fixed flavor number $`f`$ value $`\stackrel{~}{\alpha }(s;f)`$ and $`𝔄_k(s;f)`$ by a constants. For example, at $`M_5^2sM_6^2`$
$$\stackrel{~}{\alpha }(s)=\underset{s}{\overset{\mathrm{}}{}}\frac{d\sigma }{\sigma }\rho (\sigma )=\underset{s}{\overset{M_6^2}{}}\frac{d\sigma }{\sigma }\rho (\sigma ;5)+\underset{M_6^2}{\overset{\mathrm{}}{}}\frac{d\sigma }{\sigma }\rho (\sigma ;6)=\stackrel{~}{\alpha }(s;5)+c(5).$$
Generally,
$$\stackrel{~}{\alpha }(s)=\stackrel{~}{\alpha }(s;f)+c(f)\text{at}M_f^2sM_{f+1}^2$$
(17)
with shift constants $`c(f)`$ that can be calculated in terms of integrals over $`\rho (\sigma ;f+n)n1`$ with additional reservation $`c(6)=0`$ related to the asymptotic freedom condition.
More specifically,
$$c(f1)=\stackrel{~}{\alpha }(M_f^2;f)\stackrel{~}{\alpha }(M_f^2;f1)+c(f),c(6)=0.$$
These $`c(f)`$ reflect the $`\stackrel{~}{\alpha }(s)`$ continuity at the matching points $`M_f^2`$.
Analogous shift constants
$$𝔄_k(s)=𝔄_k(s;f)+𝔠_k(f)\text{at}M_f^2sM_{f+1}^2$$
(18)
are responsible for continuity of higher expansion functions. Meanwhile, $`𝔠_2(f)`$ relates to discontinuities of the “main” spectral function (14).
The one-loop estimate with $`\beta _{[f]}\rho (\sigma ;f)=\left\{\mathrm{ln}^2(\sigma /\mathrm{\Lambda }_f^2)+\pi ^2\right\}^1`$,
$$c(f1)c(f)=\frac{1}{\pi \beta _{[f]}}\mathrm{arctan}\frac{\pi }{\mathrm{ln}\frac{M_f^2}{\mathrm{\Lambda }_f^2}}\frac{1}{\pi \beta _{[f1]}}\mathrm{arctan}\frac{\pi }{\mathrm{ln}\frac{M_f^2}{\mathrm{\Lambda }_{f1}^2}}\frac{17f}{54}\overline{\alpha }_{s}^{}{}_{}{}^{3}(M_f^2)$$
(19)
and traditional values of the scale parameter $`\mathrm{\Lambda }_3,\mathrm{\Lambda }_4350250`$ MeV reveals that these constants
$$c(5)3.10^4;c(4)3.10^3;c(3)0.01,𝔠_2(f)3\alpha (M_f^2)c(f)$$
are essential at a few per cent level for $`\stackrel{~}{\alpha }`$ and at ca 10% level for the $`𝔄_2`$.
This means that the quantitative analysis of some $`s`$–channel events like, e.g., $`e^+e^{}`$ annihilation , $`\tau `$–lepton decay and charmonium width at the $`f=3`$ region should be influenced by these constants.
Global Euclidean functions On the other hand, in the Euclidean, instead of the spline-type function $`\overline{\alpha }_s`$, we have now continuous, analytic in the whole $`Q^2>0`$ domain, invariant coupling defined, along with (15), via the spectral integral
$$\alpha _{\mathrm{an}}(Q^2)=\frac{1}{\pi }\underset{0}{\overset{\mathrm{}}{}}\frac{d\sigma }{\sigma +Q^2}\rho (\sigma )$$
(20)
with the discontinuous density $`\rho (\sigma )`$ (14).
Unhappily, here, unlike for the time-like region, there is no possibility of enjoying any more explicit expression for $`\alpha _{\mathrm{an}}(Q^2)`$ even in the one-loop case. Moreover, the Euclidean functions $`\alpha _{\mathrm{an}}`$ and $`𝒜_k`$, being considered in a particular $`f`$–flavour region $`M_f^2Q^2M_{f+1}^2`$, do depend on all $`\mathrm{\Lambda }_3,\mathrm{},\mathrm{\Lambda }_6`$ values simultaneously.
Nevertheless, the real difference from the $`f=3`$ case, numerically, is not big at small $`Q^2`$ and in the “few GeV region”, for practical reasons, it could be of importance .
This situation is illustrated by Fig. 2. Here, by thick solid curves with maxima around $`\sqrt{s},Q\mathrm{\Lambda }`$, we draw expansion functions $`𝒜_2`$ and $`𝔄_2`$ in a few GeV region. Thin solid lines zeroes around $`\mathrm{\Lambda }`$ and negative values below, represent $`𝒜_3`$ and $`𝔄_3`$. For comparison, we give also second and third powers of relevant analytic couplings $`\alpha _{\mathrm{an}}`$ and $`\stackrel{~}{\alpha }`$ .
All these functions correspond to exact two–loop solutions expressed in terms of Lambert function <sup>8</sup><sup>8</sup>8Details of these calculations will be published elsewhere. The assistance of D.S. Kurashev and B.A. Magradze in calculation curves with Lambert functions is gratefully acknowledged..
## 3 Correlation of experiments
Another quantitative effect stems from the non-power structure of the IAA perturbative expansion. It is also emphasized at the few GeV region.
### 3.1 The $`s`$–channel
To illustrate the qualitative difference between our global scheme and other practice of data analysis, we first consider the $`f=3`$ region.
Inclusive $`\tau `$ decay. The IAA scheme with fixed $`f=3`$ was used in Ref. for analysis of the inclusive $`\tau `$–decay. Here, the observed quantity, the $`\tau `$ lepton time of half–decay, depends on the integral of the $`s`$–channel matrix element over the region $`0<s<M_\tau ^2`$. As a result of the 2–loop IAA analysis of the experimental input $`R_\tau =3.633`$ , the value $`\stackrel{~}{\alpha }^{(2)}(M_\tau ^2)=0.378`$ has been obtained that has to be compared with related result of usual analysis $`\overline{\alpha }_{s}^{}{}_{}{}^{(3)}(M_\tau ^2)=0.337`$. This shift $`\mathrm{\Delta }\alpha 0.04`$ resulted in a rather big change in the extracted $`\mathrm{\Lambda }`$ value. Meanwhile, some part of this shift can be “absorbed” by the shift constant $`c(3).`$
The process of Inclusive $`e^+e^{}`$ hadron annihilation provides us with an important piece of information on the QCD parameters. In the usual treatment, (see, e.g., Refs.) the basic relation looks like
$$\frac{R(s)}{R_0}=1+r(s);r(s)=\frac{\overline{\alpha }_s(s)}{\pi }+r_2\overline{\alpha }_{s}^{}{}_{}{}^{2}(s)+r_3\overline{\alpha }_{s}^{}{}_{}{}^{3}(s).$$
(21)
Here, the numerical coefficients $`r_1=1/\pi =0.318,r_2=0.142,r_3=0.413`$ (related to the $`f=5`$ case) are not diminishing. However, a rather big negative $`r_3`$ value comes mainly from the $`r_1\pi ^2\beta _{[5]}^2/3`$ contribution equal to $`0.456`$. Instead of (21), with due account of (4), we now have
$$r(s)=1+\frac{\stackrel{~}{\alpha }(s)}{\pi }+d_2𝔄_2(s)+d_3𝔄_3(s);$$
(22)
with reasonably decreasing coefficients $`d_1=0.318;d_2=0.142;d_3=0.043,`$ the mentioned $`\pi ^2`$ term of $`r_3`$ being “swallowed” by $`\stackrel{~}{\alpha }(s)`$<sup>9</sup><sup>9</sup>9This term contributes about $`8.10^4`$ into the $`r(M_Z^2)`$ and, correspondingly, 0.0025 into the extracted $`\overline{\alpha }_s(M_Z^2)`$ value. This means, that the main part of the “traditional three-loop term” $`r_3\overline{\alpha }_{s}^{}{}_{}{}^{3}`$ in the r.h.s. of (21) being of the one–loop origin is essential for the modern quantitative analysis of the data. In particular, it should be taken into the account even in the so-called NLLA which is a common approximation for the analysis of events at $`\sqrt{s}=M_Z.`$ .
Now, the main difference between (22) and (21) is due to the term $`d_2𝔄_2`$ standing in the place of $`d_2\stackrel{~}{\alpha }^2`$. The difference can be estimated by adding into (21) the structure $`r_4\alpha ^4`$ with $`r_4=d_2\beta _{[5]}^2\pi ^20.62.`$ This effect could be essential in the region of $`\stackrel{~}{\alpha }(s)0.200.25`$.
### 3.2 The $`Q^2`$–channel
The $`Q^2`$–channel: Bjorken and GLS sum rules. In the paper , the IAA has been applied to the Bjorken sum rules. Here, one has to deal with the $`Q^2`$–channel at small transfer momentum squared $`Q^210\mathrm{GeV}^2`$. Due to some controversy of experimental data, we give here only a part of the results of . For instance, using data of the SMC Collaboration for $`Q_0^2=10\mathrm{GeV}^2`$ the authors obtained $`\alpha _{\mathrm{an}}^{(3)}(Q_0^2)=0.301`$ instead of $`\alpha _{\mathrm{pt}}^{(3)}(Q_0^2)=0.275`$.
In the Euclidean channel, instead of power expansion like (2), we typically have
$$d(Q^2)=\frac{\alpha _{\mathrm{an}}(Q^2)}{\pi }+d_2𝒜_2(Q^2)+d_3𝒜_3(Q^2).$$
(23)
Here, the modification is related to non-perturbative power structures behaving like $`\mathrm{\Lambda }^2/Q^2`$ at $`Q^2\mathrm{\Lambda }^2`$. As it has been estimated above, these corrections could be essential in a few GeV region.
The same remark could be made with respect to analysis of the Gross–Llywellin-Smith sum rules of .
Some comments are in order:
— We see that, generally, the extracted values of $`\alpha _{\mathrm{an}}`$ and of $`\stackrel{~}{\alpha }`$ are both slightly greater in a few GeV region than the relevant values of $`\overline{\alpha }_s`$ for the same experimental input. This corresponds to the above-mentioned non-power character of new asymptotic expansions with a suppressed higher-loop contribution.
— At the same time, for equal values of $`\alpha _{\mathrm{an}}(x_{})=\stackrel{~}{\alpha }(x_{})=\overline{\alpha }_s(x_{})`$, the analytic scale parameter $`\mathrm{\Lambda }_{\mathrm{an}}`$ values extracted from $`\alpha _{\mathrm{an}}`$ and $`\stackrel{~}{\alpha }`$ are a bit greater than that $`\mathrm{\Lambda }_{\overline{\mathrm{MS}}}`$ taken from $`\overline{\alpha }_s`$. This feature is related to a “smoother” behavior of both the regular functions $`\alpha _{\mathrm{an}}`$ and $`\stackrel{~}{\alpha }`$ as compared to the singular $`\overline{\alpha }_s`$.
### 3.3 Conclusion
To summarize, we repeat once more our main points.
1. We have formulated the self-consistent scheme for analyzing data both in the space-like and time-like regions.
The fundamental equation connecting these regions is the dipole spectral relation (1) between renormalization–group invariant non-power expansions $`D_{\mathrm{an}}(Q^2)`$ and $`R_\pi (s)`$.
Just this equation, equivalent to the Källen–Lehmann representation, is responsible for non-perturbative terms in the $`Q^2`$–channel involved into $`\alpha _{\mathrm{an}}(Q^2)`$ and non-power expansion functions $`\left\{𝒜_k(Q^2)\right\}`$. These terms, non-analytic in the coupling constant $`\alpha `$, are a counterpart to the perfectly perturbative $`\pi ^2`$–terms effectively summed in the $`s`$–channel expressions $`\stackrel{~}{\alpha }(s)`$ and $`\left\{𝔄_k(s)\right\}`$.
2. As a by-product, we ascertain a new qualitative feature of the IAA, relating to its non-perturbativity in the $`Q^2`$–domain. It can be considered as a minimal nonperturbativity or minimal non-analyticity<sup>10</sup><sup>10</sup>10Compatible with the RG invariance and the $`Q^2`$ analyticity — compare with . in $`\alpha `$ as far as it corresponds to perturbativity in the $`s`$–channel.
Physically, it implies that minimal non-perturbativity cannot be referred to any mechanism producing effect in the $`s`$–channel.
3. The next result relates to the correlation between regions with different values of the effective flavor number $`f`$. Dealing with the massless $`\overline{\mathrm{MS}}`$ renormalization scheme, we argue that the usual perturbative QCD expansion provides our scheme only with step–discontinuous spectral density (14) depending simultaneously on different scale parameters $`\mathrm{\Lambda }_f;f=3,\mathrm{},6`$ connected by usual matching relations.
This step–discontinuous spectral density yields, on the one hand, smooth analytic coupling $`\alpha _{\mathrm{an}}(Q^2)`$ and higher functions $`\left\{𝒜_k(Q^2)\right\}`$ in the space-like region— eq.(15).
On the other hand, it produces the spline–continuous invariant coupling $`\stackrel{~}{\alpha }(s)`$ and functions $`\left\{𝔄_k(s)\right\}`$ in the time-like region — eq.(16).
As a result, the global expansion functions $`\left\{𝒜_k(Q^2)\right\}`$ and $`\left\{𝔄_k(s)\right\}`$ differ both from the that ones $`\left\{𝒜_k(Q^2;f)\right\}`$ and $`\left\{𝔄_k(s;f)\right\}`$ with a fixed value of a flavour number.
4. Thus, our global IAA scheme uses common invariant coupling $`\overline{\alpha }_s(Q^2,f)`$ and matching relations, only as an input. Practical calculation for an observable now involves expansions over the sets $`\left\{𝒜_k(Q^2)\right\}`$ and $`\left\{𝔄_k(s)\right\}`$, that is non-power series with usual numerical coefficients $`d_k`$ obtained by calculation of the relevant Feynman diagrams.
This means that, generally, one should check the accuracy of the bulk of extractions of the QCD parameters from diverse “low energy” experimental data. Our preliminary estimate shows that such a revision could influence the rate of their correlation.
5. Last but not least. As it has been mentioned in our recent publications , the IAA obeys an immunity with respect to higher loop and renormalization scheme effects.
Now, we got an additional insight into this item related to observables and can state that the perturbation series for an observable in the IAA have better convergence properties (than in usual RG–summed perturbation theory) in both the $`s`$– and $`Q^2`$– channels.
Acknowledgements
The author is indebted to D.Yu. Bardin, N.V. Krasnikov, B.A. Magradze, S.V. Mikhailov, A.V. Radyushkin, I.L. Solovtsov and O.P. Solovtsova for useful discussion and comments. This work was partially supported by grants of the Russian Foundation for Basic Research (RFBR projects Nos 99-01-00091 and 00-15-96691), by INTAS grant No 96-0842 and by INTAS-CERN grant No 2000-377.
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# On Mostow rigidity for variable negative curvature
## 1 Introduction
According to the Mostow rigidity theorem, the isometry type of a complete finite volume locally symmetric negatively curved Riemannian $`n`$-manifold with $`n3`$ is uniquely determined by its fundamental group. This is no longer true for finite volume manifolds of variable sectional curvature. In fact, there even exist negatively curved manifolds which are homeomorphic but not diffeomorphic to finite volume hyperbolic manifolds \[FJO98\]. Yet it turns out that there are essentially finitely many possibilities for the geometry and topology of such manifolds provided the sectional curvature is pinched between two negative fixed constants.
For reals $`ab0`$, let $`_{a,b,\pi ,n}`$ be the class of complete finite volume Riemannian manifolds of dimension $`n3`$ with sectional curvatures in $`[a,b]`$ and fundamental groups isomorphic to $`\pi `$. Note that $`n`$ can be read off the fundamental group $`\pi `$, namely, $`n=\mathrm{max}(\mathrm{cd}(\pi ),\mathrm{cd}(N)+1)`$ where “cd” is the cohomological dimension and $`N`$ is a maximal nilpotent subgroup of $`\pi `$. Also, $`\mathrm{cd}(\pi )=n`$ iff each manifold in $`_{a,b,\pi ,n}`$ is closed. Here is our main result.
###### Theorem 1.1.
The class $`_{a,b,\pi ,n}`$ falls into finitely many diffeomorphism types. Furthermore, for any sequence of manifolds in $`_{a,b,\pi ,n}`$, there exists a subsequence $`(M_k,g_k)`$, a smooth manifold $`M`$, and diffeomorphisms $`f_k:MM_k`$ such that the pullback metrics $`f_k^{}g_k`$ converge in $`C^{1,\alpha }`$ topology uniformly on compact subsets to a complete finite volume $`C^{1,\alpha }`$ Riemannian metric on $`M`$.
Although this result does not appear in the literature, it is not new except in dimension $`4`$ where only homeomorphism finiteness has been known. However, the proof we present is very different from the existing argument that runs as follows.
M. Gromov and W. Thurston \[Gro82\] used straightening and bounded cohomology to deduce that volume is bounded by the simplicial volume:
$$vol(M)C(a/b,n)M<\mathrm{}$$
for any $`M_{a,b,\pi ,n}`$. Since all the manifolds in $`_{a,b,\pi ,n}`$ belong to the same proper homotopy type, they must have equal simplicial volumes. Thus, volume is uniformly bounded from above on $`_{a,b,\pi ,n}`$. Furthermore, a universal lower volume bound comes from the Heintze-Margulis theorem \[BGS85\]. For closed manifolds of dimension $`n4`$ and sectional curvatures within $`[1,0)`$, the diameter can be bounded in terms of volume \[Gro78\]. Hence the conclusion of 1.1 follows from the Cheeger-Gromov compactness theorem. Similarly, the conclusion of the theorem 1.1 for finite volume manifolds of dimension $`5`$ can be deduced from the work of K. Fukaya \[Fuk84\]. The dimension restriction comes from treating the ends by the weak h-cobordism theorem. In fact, Fukaya proves a similar statement for $`n=4`$ where diffeomorphism finiteness is replaced by homotopy finiteness. (The topological $`4`$-dimensional weak $`h`$-cobordism theorem, unavailable at the time of \[Fuk84\], can be also applied here because we deal with virtually nilpotent fundamental groups \[Gui92, FT95\]. Yet this only gives homeomorphism rather than diffeomorphism finiteness.)
By contrast, main ideas in our approach come from Kleinian groups and geometric group theory. Essentially, given a degenerating sequence of manifolds $`M_k`$, one can use rescaling in the universal covers to produce a nontrivial action of $`\pi `$ on an $`𝐑`$-tree with virtually nilpotent arc stabilizers (cf. \[Bes88, Pau88, Pau91\]). Then results of Rips, Bestvina and Feighn \[BF95\] imply that $`\pi `$ splits over a virtually nilpotent subgroup. We prove that this does not happen if $`_{a,b,\pi ,n}`$ is nonempty. Then the methods of the Cheeger-Gromov compactness theorem imply that $`M_k`$ subconverges in pointed $`C^{1,\alpha }`$ topology to a complete $`C^{1,\alpha }`$ Riemannian manifold $`M`$. We then prove that $`\pi _1(M)`$ contains a subgroup isomorphic to $`\pi `$ which implies that $`M`$ has finite volume and, in fact, $`\pi _1(M)\pi `$. Now the convergence $`M_kM`$ is analogous to strong convergence of Kleinian groups. Studying this convergence yields the theorem 1.1. Instead of using the h-cobordism theorem to deal with ends, we find a direct geometric argument that works in all dimensions. Similarly to \[Gro82\], our methods provide a uniform upper bound on the volume of manifolds in $`_{a,b,\pi ,n}`$.
Note that the real Schwarz lemma of Besson-Courtois-Gallot \[BCG98\] gives yet another way to get a uniform upper bound on volume of compact manifolds in $`_{a,b,\pi ,n}`$, and hence another proof of 1.1 in this case. It is still an open question whether their method extends to finite volume manifolds.
Technically, our proof is much easier for closed manifolds; this case was previously treated in \[Bel99\]. One of the key facts needed for the theorem 1.1 is the following lemma which is of independent interest.
###### Lemma 1.2.
If the class $`_{a,b,\pi ,n}`$ is nonempty, then $`\pi `$ does not split over a virtually nilpotent subgroup. Furthermore, $`\mathrm{Out}(\pi )`$ is finite and $`\pi `$ is cohopfian.
That $`\pi `$ does not split over a virtually nilpotent subgroup can be also deduced (after some extra work) from a recent paper of B. Bowditch \[Bow98\] who studied the structure of the splittings of relatively hyperbolic groups over subgroups of peripheral groups. Unlike Bowditch’s work that applies in a more general situation, our argument is completely elementary.
Recall that a group $`\pi `$ is called cohopfian if it has no proper subgroups isomorphic to $`\pi `$. That $`\mathrm{Out}(\pi )`$ is finite and $`\pi `$ is cohopfian is due to G. Prasad \[Pra76\] in locally symmetric case. Here we follow the idea of F. Paulin \[Pau91\], and E. Rips and Z. Sela \[RS94\] who proved these properties for a large class of word-hyperbolic groups.
Topology of complete finite volume negatively curved manifolds seems to be encoded in the fundamental group. It is a deep recent result of F. T. Farrell and L. Jones \[FJ93b, FJ98\] that any homotopy equivalence of two manifolds in the class $`_{n,a,b,\pi }`$ with $`n5`$ is homotopic to a homeomorphism (which is a diffeomorphism away from compact subsets). Then the smoothing theory \[HM74, KS77\] implies that there exist at most finitely many nondiffeomorphic manifolds in the class $`_{a,b,\pi ,n}`$ with $`n5`$. Furthermore, if $`n6`$, there are homeomorphic negatively curved manifolds that are not diffeomorphic \[FJ89b, FJ89a, FJ93a, FJO98\]. Such examples are still unknown if $`n=4,5`$. If $`n=3`$, one expects that any two manifolds in the class $`_{a,b,\pi ,n}`$ are diffeomorphic. This is known for Haken manifolds \[Wal68\] (note that any noncompact finite volume manifold is Haken). For non-Haken manifolds this would follow from (as yet unproved) Thurston’s hyperbolization conjecture and Mostow Rigidity.
Note that 1.1 gives diffeomorphism finiteness in all dimensions. As we explained above this is mostly interesting if $`n=4`$; in fact, the interior of any compact smooth $`4`$-manifold with nonempty boundary has at least countably many smooth structures \[BE98\].
Theorem 1.1 combined with results of Gao \[Gao90\] immediately implies that the class of finite volume Einstein manifolds that lie in $`_{a,b,\pi ,n}`$ is compact in pointed $`C^{\mathrm{}}`$ topology (i.e. for any sequence of Einstein manifolds in $`_{a,b,\pi ,n}`$, there is a subsequence $`M_k`$, a manifold $`M`$, and diffeomorphisms $`f_k:MM_k`$ such that $`f_k^{}g_k`$ converge in $`C^{\mathrm{}}`$ topology uniformly on compact subsets to a complete finite volume Einstein metric on $`M`$). Furthermore, since negatively curved Einstein metrics on compact manifolds of dimension $`>2`$ are isolated in the moduli space of the Einstein metrics \[Bes88, 12.73\], we conclude that, up to homothety, there are only finitely many compact Einstein manifolds in $`_{a,b,\pi ,n}`$. It is a tantalizing open problem to decide whether any compact negatively curved Einstein manifold is locally symmetric. A related question was recently resolved in dimension four: any Einstein metric on a compact negatively curved locally symmetric $`4`$-manifold is locally symmetric \[BCG95, LeB95\]
One can also use 1.1 to deduce several pinching results. For example, given a group $`\pi `$, there exists an $`ϵ=ϵ(\pi )>0`$ such that any finite volume manifold from $`_{1ϵ,1,\pi ,n}`$ is diffeomorphic to a real hyperbolic manifold. Note that $`ϵ`$ has to depend on the topology of the manifolds (in our case on the fundamental group) as examples \[FJ89b, FJ93a, FJO98, GT87\] show. Similar results hold for almost quarter pinched Kähler and quaternionic-Kähler manifold manifolds.
The structure of the paper is as follows. In the section $`2`$ we give some background on finite volume negatively curved manifolds. Sections $`3`$ and $`5`$ contain a proof of 1.2. Convergence of finite volume negatively curved manifolds is discussed in Section $`4`$. Theorem 1.1 is proved in Section $`6`$. Applications to pinching are discussed in Section $`7`$.
The author is grateful to Christoph Böhm, Brian Bowditch, Harish Seshadri, Peter Shalen, Christopher Stark, and McKenzie Wang for helpful discussions or communications.
## 2 Preliminaries
Let $`M`$ be a finite volume complete Riemannian manifold of sectional curvature within $`[a,b]`$, $`ab<0`$ of dimension $`n>2`$. In this section we list some properties of $`M`$ which we are going to use throughout the paper without explicit references. A comprehensive account on nonpositive curvature can be found in \[BGS85\]. Also see \[HH77, Sch84, Bow93, Bow95\].
###### 2.1.
Virtually nilpotent subgroups. Write $`M`$ as $`X/\pi `$ where $`X`$ is the universal cover of $`M`$ and $`\pi \pi _1(M)`$ is the covering group. By the Cartan-Hadamard theorem $`X`$ is diffeomorphic to the Euclidean space, thus $`M`$ is aspherical. An infinite order isometry $`\gamma `$ of $`X`$ either stabilizes a bi-infinite geodesic or fix exactly one point at infinity; such a $`\gamma `$ is called hyperbolic or parabolic, respectively.
Virtually nilpotent discrete subgroups of $`Isom(X)`$ are finitely generated \[Bow93\]. Since $`X`$ is contractible, $`\pi `$ is torsion free. Any nontrivial virtually nilpotent subgroup of $`\pi `$ is either an infinite cyclic group generated by a loxodromic isometry and stabilizing a bi-infinite geodesic, or a group that consists of parabolic isometries with a common fixed point at infinity. Any virtually nilpotent subgroup of $`\pi `$ lies in a unique maximal virtually nilpotent subgroup of $`\pi `$. See \[BGS85, Bow93, Bow95\] for more information.
###### 2.2.
Thin/thick decomposition. For a positive $`ϵ`$, write $`M_{[ϵ,\mathrm{})}`$ for the $`ϵ`$-thick part of $`M`$ which is the set of points of $`M`$ with injectivity radius $`ϵ`$. Similarly, $`M_{(0,ϵ)}=MM_{[ϵ,\mathrm{})}`$ is called the $`ϵ`$-thin part of $`M`$.
According to the Margulis lemma, there exists a universal constant $`\mu _{n,a}`$ such that for each $`ϵ<\mu _{n,a}`$ the $`ϵ`$-thin part of $`M`$ is a union of finitely many connected components. Unbounded components are called cusps while bounded components are called tubes. Each tube contains a closed geodesic of length $`2ϵ`$ and is homeomorphic to the tubular neighborhood of the geodesic. Since there are only finitely many tubes, we can assume that $`M_{(0,ϵ)}`$ consists of cusps by taking $`ϵ`$ small enough.
Each cusp is a union of geodesic rays emanating from a common point at infinity. Also, given a cusp $`C`$, let $`\stackrel{~}{C}`$ be a connected component of the preimage of $`C`$ in the universal cover. The group $`\{\gamma \pi :\gamma (\stackrel{~}{C})=\stackrel{~}{C}\}`$ coincides with the stabilizer $`\mathrm{\Gamma }_z`$ in $`\mathrm{\Gamma }`$ of a point $`z`$ at infinity. The group $`\mathrm{\Gamma }_z`$ preserves horospheres centered at $`z`$ and acts on each horosphere with compact quotient since $`C=\stackrel{~}{C}/\mathrm{\Gamma }_z`$ has finite volume. Horospheres are $`C^2`$ submanifolds of $`X`$ each diffeomorphic to the Euclidean space \[HH77\]. Each horosphere centered at $`z`$ is orthogonal to geodesics asymptotic to $`z`$. Tangent vectors to such geodesics form the so-called radial vector field on $`X`$ and $`C`$; this is a $`C^1`$-vector field \[HH77\]. (Throughout the paper all geodesic are assumed to have unit speed.) In fact, radial vector field is the gradient of the so called Busemann function. Any Busemann function defines a $`C^2`$-Riemannian submersion of a cusp region bounded by a horosphere into the real line. Each cusp is diffeomorphic to the product of a real line and a closed aspherical manifold which is the quotient of a horosphere by $`\mathrm{\Gamma }_z`$. Note that the boundary $`C`$ of a cusp is generally nonsmooth. Pushing along geodesic rays asymptotic to $`x`$ defines a homeomorphism of $`C`$ and the $`\mathrm{\Gamma }_z`$-quotient of a horosphere.
###### 2.3.
Compactification. According to \[BGS85\], $`M`$ is diffeomorphic to the interior of a compact manifold with (possibly empty) boundary. If the boundary is nonempty, its connected components are quotients of horospheres by maximal parabolic subgroups of $`\pi `$. Each boundary component corresponds to a conjugacy class of maximal parabolic subgroups of $`\pi `$. The inclusion of each boundary component is $`\pi _1`$-injective.
###### 2.4.
Exponential convergence of geodesics. To simplify notations we let $`a=K^2`$, $`b=k^2`$. Let $`\gamma _1`$, $`\gamma _2`$ be two geodesics asymptotic to a point $`z=\gamma _1(+\mathrm{})=\gamma _2(+\mathrm{})`$ such that $`\gamma _1(0)`$, $`\gamma _2(0)`$ lie on the same horosphere centered at $`z`$. Then for any $`t`$, $`\gamma _1(t)`$, $`\gamma _2(t)`$ lie on the same horosphere centered at $`z`$. Denote by $`h(t)`$ the distance between $`\gamma _1(t)`$ and $`\gamma _2(t)`$ on this horosphere equipped with the induced Riemannian metric. Also denote by $`d(t)`$ the distance between $`\gamma _1(t)`$ and $`\gamma _2(t)`$ in $`X`$. It is proved in \[HH77\] that for $`t0`$
$$e^{kt}\frac{h(0)}{h(t)}e^{Kt}\mathrm{and}\frac{2}{k}\mathrm{sinh}\left(\frac{kd(t)}{2}\right)h(t)\frac{2}{K}\mathrm{sinh}\left(\frac{Kd(t)}{2}\right).$$
Therefore, we deduce that $`d(t)h(t)h(0)e^{kt}\frac{2}{K}\mathrm{sinh}\left(\frac{Kd(0)}{2}\right)e^{kt}`$ and $`d(0)e^{Kt}h(0)e^{Kt}h(t)\frac{2}{K}\mathrm{sinh}\left(\frac{Kd(t)}{2}\right)`$ for all $`t0`$. It is straightforward to check that $`\mathrm{sinh}(x)2x`$ whenever $`x[0,1]`$. Also, it follows from comparison with Euclidean triangles that $`d(t)d(0)`$ for $`t0`$. Hence, if $`d(0)2/K`$ and $`t0`$, then $`e^{kt}/2d(0)/d(t)2e^{Kt}`$.
## 3 No splitting
Throughout this section $`\pi `$ is the fundamental group of a finite volume noncompact complete Riemannian manifold $`M`$ of dimension $`n>2`$ and with sectional curvatures within $`[a,b]`$ for $`ab<0`$. By \[BGS85\] $`M`$ is the interior of a compact manifold with boundary which we denote by $`M_\pi `$.
A group $`G`$ is said to split over a subgroup $`C`$ if $`G=A_CB`$ or $`G=A_C`$ where $`ACB`$. It is well-known that $`A`$ and $`B`$ necessarily have infinite index in $`G`$. Note that $`G`$ splits over $`C`$ iff $`G`$ act without edge inversions on a simplicial tree with no proper invariant subtree, no global fixed point, and exactly one orbit of edges such that $`C`$ is a stabilizer of some edge \[Ser80\].
The purpose of this section is to prove that $`\pi `$ does not split over a virtually nilpotent group. When $`M`$ is a closed manifold, this can be shown simply by looking at the Mayer-Vietoris sequence of the splitting (see e.g. \[Bel99\]).
Noncompact case is more subtle. Main ingredients of the proof are again Mayer-Vietoris sequence, and the following splitting-theoretic lemma of B. Bowditch. More general (and harder to prove) results can be deduced from a recent paper of Bowditch \[Bow98\] who studies the structure of splittings of relatively hyperbolic groups over subgroups of peripheral groups. By contrast, our approach is elementary: we use basic manifold topology and Bass-Serre theory.
###### Lemma 3.1.
Suppose that $`\pi `$ splits over a virtually nilpotent subgroup. Then $`\pi `$ also splits over a virtually nilpotent subgroup $`C`$ as $`A_CB`$ or $`A_C`$ where $`ACB`$ in such a way that a conjugate of any maximal parabolic subgroup lies in $`A`$ or $`B`$. Furthermore, if $`C`$ is parabolic, then the splitting can be chosen so that the maximal parabolic subgroup containing $`C`$ lies in $`A`$.
###### Proof.
The proof is essentially borrowed from \[Bow98, 3.5, 5.2\] where a more general situation is considered. We specialize terminology to our case and give more details when seems appropriate.
Since $`\pi `$ splits over a virtually nilpotent subgroup, $`\pi `$ acts without edge inversions on a simplicial tree $`T`$ with no proper invariant subtree, no global fixed point, and exactly one orbit of edges. We seek to construct a $`\pi `$-action on a (perhaps another) tree with the above properties such that every maximal parabolic subgroup of $`\pi `$ fixes a vertex.
Fix an edge $`e`$ of $`T`$ and denote its stabilizer in $`\pi `$ by $`\pi _e`$. Let $`P`$ be the maximal virtually nilpotent subgroup of $`\pi `$ that contains $`\pi _e`$. First, note that any maximal parabolic subgroup $`P^{}`$ other than $`P`$ fixes a vertex of $`T`$. (Indeed, every edge stabilizer of $`T`$ is a conjugate of $`\pi _e`$, in particular, it lies in a conjugate of $`P`$, hence it must have trivial intersection with $`P^{}`$. So, unless $`P^{}`$ fixes a vertex of $`T`$, we get that $`P^{}`$ splits over the trivial group which is impossible because $`P^{}`$ is noncyclic virtually nilpotent, and hence one-ended.) In particular, if $`P`$ is not parabolic, then the $`\pi `$-action on $`T`$ satisfies the desired properties.
Now assume that $`P`$ fixes a vertex $`v`$ of $`T`$. Then the fixed-point-set of $`\pi _e`$ contains the segment joining $`v`$ and a vertex of $`e`$. In particular, $`\pi _e`$ lies in $`\pi _{\overline{e}}`$, the stabilizer of an edge $`\overline{e}`$ adjacent to $`v`$. Since $`\pi _{\overline{e}}`$ is conjugate to $`\pi _e`$, it is virtually nilpotent. Hence $`\pi _{\overline{e}}`$ lies in a maximal parabolic subgroup containing $`\pi _e`$ which is $`P`$. Thus the splitting of $`\pi `$ over $`\pi _{\overline{e}}`$ satisfies the desired properties.
It remains to consider the case when $`P`$ is parabolic and $`P`$ does not fix a vertex of $`T`$. Let $`\tau `$ be the unique $`P`$-invariant minimal subtree of $`T`$. Let $`g`$ be an element of $`\pi `$ such that $`g^1e`$ is an edge of $`\tau `$ (such a $`g`$ exists since there is only one orbit of edges). Then the stabilizer of $`g^1e`$ in $`\pi `$ is $`g^1\pi _eg`$. Note that the group $`g^1\pi _egP`$ (which is the stabilizer of $`g^1e`$ in $`P`$) is infinite because otherwise $`P`$ splits over a finite group $`g^1\pi _egP`$ and $`P`$ cannot since noncyclic nilpotent groups are one ended. Since $`\pi _eP`$, we conclude that $`g^1PgP`$ is infinite, therefore $`P=g^1Pg`$ and, hence, $`gP`$ because any maximal parabolic subgroup is equal to its normalizer. The above argument has several implications as follows.
$``$ $`e`$ must be an edge of $`\tau `$. (Indeed, using that there is only one orbit of edges we find $`g\pi `$ such that $`g^1e`$ is an edge of $`\tau `$. By above $`g`$ must belong to $`P`$ and since $`P`$ stabilizes $`\tau `$, $`e`$ is an edge of $`\tau `$.)
$``$ $`P`$ is equal to the setwise stabilizer of $`\tau `$. (Indeed, if $`h\tau =\tau `$, than $`h`$ takes the edge $`e`$ of $`\tau `$ to an edge of $`\tau `$. Hence, again $`hP`$.)
$``$ Any edge of $`T`$ lies in a unique $`\pi `$-image of $`\tau `$. (Indeed, any edge of $`T`$ is $`\pi `$-equivalent to $`e`$, so it is an edge of a tree that is a $`\pi `$-image of $`\tau `$. Uniqueness of this tree is deduced as follows. Suppose, arguing by contradiction, that two trees that are in $`\pi `$-image of $`\tau `$ share an edge. Applying an element of $`\pi `$ we can assume these two trees are $`\tau `$ and, say, $`f\tau `$ for some $`f\pi `$. This common edge is of the form $`g^1e`$ for some $`g\pi `$, so by above $`gP`$. Hence, $`e`$ is the common edge of $`\tau =g\tau `$ and $`gf\tau `$. So $`e`$ is $`gf`$-image of an edge of $`\tau `$, or equivalently, $`f^1g^1e`$ is an edge of $`\tau `$. By the same argument, $`gfP`$, so $`fP`$ which contradicts to $`\tau f\tau `$.)
Let $`Q`$ be the set of $`\pi `$-images of $`\tau `$. Construct now a graph $`S(Q)`$ with the vertex set $`V(T)Q`$ where we deem two vertices $`v,\eta `$ adjacent iff $`vV(T)`$, $`\eta Q`$, and $`v\eta `$.
First, prove that $`S(Q)`$ is a (simplicial) tree. Indeed, to show that $`S(Q)`$ is connected it suffices to find an arc that joins two arbitrary vertices $`x,yV(T)`$. The shortest arc that joins $`x`$ and $`y`$ in $`T`$ can be uniquely written in the form $`x_0x_1\mathrm{}x_n`$ where $`x_0=x`$, $`x_n=y`$ and the arc $`x_{i1}x_i`$ is contained in the tree $`\tau _iQ`$ (any edge of $`T`$ lies in a unique tree from $`Q`$). Then the arc $`x_0\tau _1x_1\mathrm{}\tau _nx_n`$ joins $`x`$ and $`y`$ in $`S(Q)`$. Second, show that $`S(Q)`$ has no circuits. Indeed, suppose $`x_1\tau _1x_2\mathrm{}\tau _nx_n`$ is a circuit in $`S(Q)`$ with $`x_1=x_n`$. Let $`\alpha _i`$ be the arc in $`\tau _i`$ that connects $`x_i`$ to $`x_{i+1}`$. Then $`\alpha _1\alpha _2\mathrm{}\alpha _n`$ is a circuit in $`T`$, a contradiction.
The group $`\pi `$ acts on $`S(Q)`$ without edge inversions. Any maximal parabolic subgroup now fixes a vertex. In particular, $`P`$ is the stabilizer of $`\tau `$. Since $`\pi `$ is not virtually nilpotent, there exists a vertex $`v`$ of $`\tau `$ whose stabilizer $`\pi _v`$ is not a subgroup of $`P`$. (If the stabilizers of two adjacent vertex groups of $`\tau `$ were subgroups of $`P`$, then, since $`T`$ has only one $`\pi `$-orbit of edges, the groups $`A`$ and $`B`$ would lie in a conjugate of $`P`$. Hence $`\pi =A,B`$ would lie in the conjugate of $`P`$.)
Now look at the edge $`d`$ that joins vertices $`v`$ and $`\tau `$. The stabilizer $`\pi _d`$ of $`d`$ is a subgroup of $`P`$, so it is not equal to $`\pi _v`$. If $`\pi _d=P`$, then $`P`$ has to stabilize $`v`$ as well, hence $`P`$ fixes a point of $`T`$ which is not the case. Thus, $`\pi _d`$ is not equal to the stabilizers of $`v`$ and $`\tau `$.
Now if any vertex of $`T`$ is $`\pi `$-equivalent to $`v`$, then $`S(Q)`$ has one orbit of edges and $`\pi \pi _v_{\pi _d}P`$ is the desired splitting. Otherwise, $`T`$ has two orbits of vertices represented by $`v`$ and $`v^{}\tau `$ and hence $`S(Q)`$ has two orbits of edges $`d`$, $`d^{}`$. This defines a splitting $`\pi \pi _v_{\pi _d}(P_{\pi _d^{}}\pi _v^{})`$ over $`\pi _d`$ as needed. ∎
###### 3.2.
Gluing aspherical cell complexes. In this section we frequently use the following standard construction. Let $`f:XY`$ and $`g:XZ`$ be (cellular) maps of cell complexes that induce $`\pi _1`$-injections on each connected component of $`X`$. Assume that $`Y`$, $`Z`$ and each connected component of $`X`$ are aspherical. Form a cell complex by gluing $`X\times [0,1]`$ to $`Y`$ and $`Z`$ where $`(x,0)`$ and $`(x,1)`$ are identified with $`f(x)`$ and $`g(x)`$, respectively. The result is an aspherical cell complex (use Mayer-Vietoris in the universal covers and then the Hurewitz theorem). When $`X`$ is connected, its fundamental group is isomorphic to $`\pi _1(Y)_{\pi _1(X)}\pi _1(Z)`$. Similarly, if $`Y=Z`$ we can form a cell complex by gluing $`X\times [0,1]`$ to $`Y`$ where $`(x,0)`$ and $`(x,1)`$ are identified with $`f(x)`$ and $`g(x)`$, respectively. The result is, again, an aspherical cell complex and, if $`X`$ is connected, its fundamental group is isomorphic to $`\pi _1(Y)_{\pi _1(X)}`$.
###### 3.3.
Topological model for the splitting: amalgamated product case. Assume that $`\pi `$ splits over a virtually nilpotent group so that 3.1 gives an isomorphism $`A_CB\pi `$ such that any maximal parabolic subgroup is conjugate to a subgroup of $`A`$ or $`B`$, and the maximal virtually nilpotent subgroup $`P`$ containing $`C`$ lies in $`A`$. The inclusions of $`A`$ and $`B`$ into $`\pi `$ define coverings $`M_AM_\pi `$, $`M_BM_\pi `$.
Note that for each connected component $`L`$ of $`M_\pi `$, the inclusion $`LM_\pi `$ lifts to $`M_A`$ or $`M_B`$. (Indeed, maps into aspherical manifolds are homotopic iff the induced $`\pi _1`$-homomorphisms are conjugate so $`LM_\pi `$ is homotopic to a map that takes $`\pi _1(B)`$ to $`A`$ or $`B`$. This map lifts to the corresponding cover, and hence so does $`LM_\pi `$ by the covering homotopy theorem.) For every boundary component of $`M_\pi `$ fix such a lift thereby defining an injection $`\pi _0(M_\pi )\pi _0(M_AM_B)`$. We refer to these components of $`M_AM_B`$ as lifted.
Consider a closed aspherical manifold $`N_P`$ with $`\pi _1(N_P)P`$ defined as follows. If $`P`$ is parabolic, $`N_P`$ is the boundary component of $`M_\pi `$ corresponding to the inclusion $`P\pi `$. If $`P`$ is loxodromic, $`N_P`$ is a circle embedded in $`M_\pi `$ that realizes a generator of $`P`$. If $`P`$ is trivial, $`N_P`$ is a point of $`M_\pi `$.
In any case the inclusion $`N_PM_\pi `$ lifts to $`M_A`$, and we fix such a lift (in case $`P`$ is parabolic, we use the same lift onto a boundary component of $`M_A`$ that has been chosen above.) We use the notation $`\stackrel{~}{N}_P`$ for the lifted manifold. The covering $`N_CN_P`$ induced by $`CP`$ lifts to a covering onto $`\stackrel{~}{N}_PM_A`$.
Also fix a lift of $`N_CN_P`$ to $`M_B`$. This lift turns out to be a homeomorphism of $`N_C`$ onto a boundary component of $`M_B`$ because $`C=AB=PB`$. This component is denoted by $`\stackrel{~}{N}_C`$.
The above maps of $`N_C`$ into $`M_A`$ and $`M_B`$ certainly become homotopic after projecting to $`M_\pi `$. Now we use this homotopy to build an aspherical cell complex $`Y`$ by gluing $`M_A`$, $`M_B`$, and $`N_C\times [0,1]`$ identifying $`N_C\times \{0\}`$ with $`\stackrel{~}{N}_P`$ via the covering $`N_C\stackrel{~}{N}_P`$, and $`N_C\times \{1\}`$ with $`\stackrel{~}{N}_C`$ via the chosen lift. Now we get a homotopy equivalence $`YM_\pi `$ which extends the coverings $`M_AM_\pi `$ and $`M_BM_\pi `$.
Doubling $`M_\pi `$ along the boundary produces a closed aspherical manifold $`DM_\pi `$. Similarly, the doubles of $`M_A`$, $`M_B`$ along the lifted boundary components are denoted by $`DM_A`$, $`DM_B`$. Finally, let $`DY`$ be the “double” of $`Y`$ along the lifted boundary components, that is, $`DY`$ is obtained by gluing $`DM_A`$, $`DM_B`$, and two copies of $`N_C\times [0,1]`$ as above. The homotopy equivalence $`YM_\pi `$ now extends to a homotopy equivalence $`DYDM_\pi `$.
It is useful to give another description of $`DY`$. Namely, let $`\overline{D}M_B`$ be the manifold obtained from $`DM_B`$ by identifying two copies of $`\stackrel{~}{N}_CDM_B`$. Then $`DY`$ is obtained by gluing $`DM_A`$, $`\overline{D}M_B`$, and $`N_C\times [0,1]`$ where $`N_C\times \{1\}`$ is identified with $`\stackrel{~}{N}_C\overline{D}M_B`$, and $`N_C\times \{0\}`$ is glued to $`\stackrel{~}{N}_PM_A`$ via the covering $`N_C\stackrel{~}{N}_P`$.
###### 3.4.
Topological model for the splitting: HNN-extension case. As above, fix an isomorphism $`A_C\pi `$ such that any maximal parabolic subgroup is conjugate to a subgroup of $`A`$ or $`B`$, and the maximal virtually nilpotent subgroup $`P`$ containing $`C`$ lies in $`A`$. (We think of $`A_C`$ as $`A,t|\varphi (C)=tCt^1`$ where $`\varphi :CA`$ is a monomorphism.)
We keep essentially the same notations as in the amalgamated product case. Thus the inclusion $`A\pi `$ defines a covering $`M_AM_\pi `$ and, we fix lifts of all boundary components of $`M_\pi `$ to $`M_A`$. Again, a lift of $`N_P`$ to $`M_A`$ is denoted by $`\stackrel{~}{N}_P`$. The covering $`N_CN_P`$ induced by the inclusion $`CP`$ of course lifts to a covering onto $`\stackrel{~}{N}_P`$ using the chosen lift of $`N_P`$.
Also the composition $`N_CN_PM_\pi `$ can be lifted to a homeomorphism onto a component $`\stackrel{~}{N}_P`$ of $`M_A`$ so that the lift induces $`\varphi :CA`$. (Indeed, let $`f`$ be the lift of $`N_CN_PM_\pi `$ to the universal covers such that $`f`$ is equivariant with respect to the inclusion $`C\pi `$. Then $`g=tf`$ is a $`\varphi `$-equivariant homeomorphism onto a boundary component of the universal cover of $`M_\pi `$. So $`g`$ descends to a covering of $`N_C`$ onto a component of $`M_A`$ which is, in fact, a homeomorphism because $`\varphi (C)`$ is maximal parabolic in $`A`$. The last statement is true since the maximal virtually nilpotent subgroup of $`\pi `$ containing $`\varphi (C)`$ is $`tPt^1`$ and $`\varphi (C)=tAt^1A=tPt^1A`$.)
Again, $`\overline{D}M_A`$ is a double of $`M_A`$ along the union of the lifted boundary components and $`\stackrel{~}{N}_C`$. As before we get a homotopy equivalence between $`DM_\pi `$ and $`DY`$ where $`DY`$ is obtained by gluing $`\overline{D}M_A`$ and $`N_C\times [0,1]`$ so that $`N_C\times \{0\}`$ is identified with a copy of $`\stackrel{~}{N}_C`$ sitting inside $`\overline{D}M_A`$, while $`N_C\times \{1\}`$ is glued to $`\stackrel{~}{N}_P`$ via the covering $`N_C\stackrel{~}{N}_P`$.
###### Theorem 3.5.
Let $`\pi `$ be the fundamental group of a finite volume complete Riemannian manifold $`M`$ of dimension $`n>2`$ and with sectional curvatures within $`[a,b]`$ for $`ab<0`$. Then $`\pi `$ does not split over a virtually nilpotent group.
###### Proof.
Assume first that the splitting of $`\pi `$ given by 3.1 is $`A_CB`$. Think of $`DY`$ as glued from $`DM_A`$, $`\overline{D}M_B`$, and $`N_C\times [0,1]`$. Since the splitting is nontrivial, both $`A`$ and $`B`$ have infinite index in $`\pi `$ so $`M_A`$, $`M_B`$ are noncompact. Hence $`DM_A`$, $`\overline{D}M_B`$ are noncompact since they are glued from noncompact spaces along compact subset. Now look at the Mayer-Vietoris sequence for homology with $`/2`$-coefficients.
$$0H_n(DY)H_{n1}(N_C)H_{n1}(DM_A)H_{n1}(\overline{D}M_B)\mathrm{}$$
The map $`H_{n1}(N_C)H_{n1}(DM_A)H_{n1}(\overline{D}M_B)`$ can be written as $`i_ai_b`$ where $`i_a:N_CDM_A`$ is the covering onto $`\stackrel{~}{N}_P`$, and $`i_b:N_C\overline{D}M_B`$ is the homeomorphism onto $`\stackrel{~}{N}_C`$. Show that $`i_b`$ is injective. This is clear if the cohomological dimension of $`C`$ is $`<n1`$. Otherwise, by the exact sequence of the pair $`(\overline{D}M_B,\stackrel{~}{N}_C)`$ it suffices to show that $`H_n(\overline{D}M_B,\stackrel{~}{N}_C)=0`$ which is true by \[Dol72, VIII.3.4\]. (As stated, \[Dol72, VIII.3.4\] only applies when the complement of $`\stackrel{~}{N}_C`$ in $`\overline{D}M_B`$ is connected. However, if $`\overline{D}M_B\stackrel{~}{N}_C`$ is nonconnected, it has two noncompact components each equal to $`M_B`$ and the result again follows since by the relative Mayer-Vietoris sequence $`H_n(\overline{D}M_B,\stackrel{~}{N}_C)`$ is a sum of two copies of $`H_n(M_B,\stackrel{~}{N}_C)=0`$.)
By exactness $`H_n(DY)=0`$ and we get a contradiction with the fact that $`DY`$ is homotopy equivalent to the closed $`n`$-manifold $`DM_\pi `$.
Similarly, if $`\pi A_C`$, then we get the following Mayer-Vietoris sequence with $`/2`$-coefficients \[Bro82\].
$$0H_n(DY)H_{n1}(N_C)H_{n1}(\overline{D}M_A)\mathrm{}$$
The map $`H_{n1}(N_C)H_{n1}(\overline{D}M_A)`$ can be written as $`i_{}f_{}`$ where $`i:N_C\overline{D}M_A`$ is the homeomorphism onto $`\stackrel{~}{N}_C`$, and $`f`$ is the covering of $`N_C`$ onto $`\stackrel{~}{N}_PDM_A`$. It remains to show that $`i_{}f_{}`$ is injective which is clear if the cohomological dimension of $`C`$ is $`<n1`$. Otherwise, look at the exact sequence of the pair $`(\overline{D}M_A,\stackrel{~}{N}_C\stackrel{~}{N}_P)`$ with $`/2`$-coefficients. Again, by \[Dol72, VIII.3.4\] $`H_n(\overline{D}M_A,\stackrel{~}{N}_C\stackrel{~}{N}_P)=0`$, hence the inclusion $`\stackrel{~}{N}_C\stackrel{~}{N}_P\overline{D}M_A`$ induces an injection on $`(n1)`$-homology. Therefore,
$$(i_{}f_{})[N_C]=[N_C][N_P]\mathrm{deg}(N_CN_P)$$
is nonzero as wanted an we get a contradiction as before. ∎
## 4 Convergence of finite volume manifolds
We refer to \[Bel98a\] or \[Bel98b\] for background. Here we only recall basic definitions and prove several new lemmas specific to the finite volume case. By an action of an abstract group $`\pi `$ on a space $`X`$ we mean a group homomorphism $`\rho :\pi \mathrm{Homeo}(X)`$. An action $`\rho `$ is called free if $`\rho (\gamma )(x)x`$ for all $`xX`$ and all $`\gamma \pi \text{id}`$. In particular, if $`\rho `$ is a free action, then $`\rho `$ is injective.
###### 4.1.
Equivariant pointed Lipschitz topology. Let $`\mathrm{\Gamma }_k`$ be a discrete subgroup of the isometry group of a complete Riemannian manifold $`X_k`$ and $`p_k`$ be a point of $`X_k`$. The class of all such triples $`\{(X_k,p_k,\mathrm{\Gamma }_k)\}`$ can be given the so-called equivariant pointed Lipschitz topology \[Fuk86\]; when $`\mathrm{\Gamma }_k`$ is trivial this reduces to the usual pointed Lipschitz topology.
If $`(X_k,p_k)`$ is a sequence of simply-connected complete Riemannian $`n`$-manifolds with $`asec(X_k)0`$, then $`(X_k,p_k)`$ subconverges in the pointed Lipschitz topology to $`(X,p)`$ where $`X`$ is a $`C^{\mathrm{}}`$–manifold with a complete $`C^{1,\alpha }`$–Riemannian metric of Alexandrov curvature $`a`$ and $`0`$. In fact, in a suitable harmonic atlas on $`X`$, the sequence $`(X_k,p_k)`$ subconverges to $`(X,p)`$ in pointed $`C^{1,\alpha }`$ topology \[Fuk86, And90\].
###### 4.2.
Pointwise convergence topology. Suppose that the sequence $`(X_k,p_k)`$ converges to $`(X,p)`$ in the pointed Lipschitz topology. This allows one to talk about the convergence of a sequence of points $`x_kX_k`$ to $`xX`$ Furthermore, a sequence of isometries $`\gamma _k\mathrm{Isom}(X_k)`$ we say that $`\gamma _k`$ converges, if for any $`xX`$ and any sequence $`x_kX_k`$ that converges to $`x`$, $`\gamma _k(x_k)`$ converges. The limiting transformation $`\gamma `$ that takes $`x`$ to the limit of $`\gamma _k(x_k)`$ is necessarily an isometry of $`X`$.
Let $`\rho _k:\pi \mathrm{Isom}(X_k)`$ be a sequence of isometric actions of a group $`\pi `$ on $`X_k`$. We say that a sequence of actions $`(X_k,p_k,\rho _k)`$ converges in the pointwise convergence topology if $`\rho _k(\gamma )`$ converges for every $`\gamma \pi `$. The map $`\rho :\mathrm{\Gamma }\text{Isom}(X)`$ that takes $`\gamma `$ to the limit of $`\rho _k(\gamma )`$ is necessarily a homomorphism.
###### Lemma 4.3.
Let $`\rho _k:\pi \mathrm{Isom}(X_k)`$ be a sequence of free isometric actions of a discrete group $`\pi `$ on Hadamard $`n`$-manifolds $`X_k`$ of sectional curvatures within $`[a,b]`$ for $`ab<0`$. Assume that the sequence $`(X_k,p_k,\rho _k(\pi ))`$ converges in the equivariant pointed Lipschitz topology to $`(X,p,\mathrm{\Gamma })`$ and $`(X_k,p_k,\rho _k)`$ converges to $`(X,p,\rho )`$ in the pointwise convergence topology. If $`\rho (\pi )`$ has finite index in $`\mathrm{\Gamma }`$, then $`\mathrm{\Gamma }=\rho (\pi )`$.
###### Proof.
Being the fundamental group of a pinched negatively curved manifold, $`\pi `$ is torsion free and any abelian subgroup of $`\pi `$ is finitely generated \[BS87, Bow93\]. Therefore, every element in $`\pi `$ is a power of a primitive element (an element of a group is called primitive it is not a proper power). Fix an arbitrary $`h\mathrm{\Gamma }`$ and find $`g\pi `$ so that $`h^m=\rho (g)`$ for some $`m1`$. By passing to the appropriate root we can assume $`g`$ is primitive. There always exists a sequence $`\rho _k(g_k)`$ that converges to $`h`$. Then both $`\rho _k(g_k^m)`$ and $`\rho _k(g)`$ converge to $`h^m`$. Hence $`g_k^m=g`$ for large $`k`$ by \[Bel98a, lem 2.6\]. Since $`g`$ is primitive, $`m=1`$ so $`h\rho (\pi )`$ as wanted. ∎
In general, if $`|sec|`$ is bounded above, the injectivity radius is at best upper-semicontinuous. However, it becomes continuous if the curvature is also nonpositive.
###### Lemma 4.4.
If $`(M_k,p_k)`$ be a sequence of pointed complete Riemannian manifolds of sectional curvatures within $`[a,0]`$ for $`a0`$. Assume that $`(M_k,p_k)`$ converges to $`(M,p)`$ in pointed $`C^{1,\alpha }`$ topology. Then for any $`x_kM_k`$ that converges to $`xM`$, the injectivity radii of $`M_k`$ at $`x_k`$ converge to the injectivity radius of $`M`$ at $`x`$.
###### Proof.
Arguing by contradiction find $`x_k`$ that converges to $`x`$ while $`inj(x_k)`$ does not approach $`inj(x)`$. Pass to a subsequence so that $`inj(x_k)`$ converges to $`cinj(x)`$. Write $`M_k=X_k/\mathrm{\Gamma }_k`$ and $`M=X/\mathrm{\Gamma }`$ and pick preimages $`\stackrel{~}{x}_k`$ of $`x_k`$, and $`\stackrel{~}{x}`$ of $`x`$. Since $`|sec(M_k)|`$ is uniformly bounded, a result in \[Fuk86\] implies that $`(X_k,\stackrel{~}{x}_k,\mathrm{\Gamma }_k)`$ subconverges to $`(X,\stackrel{~}{x},\mathrm{\Gamma })`$ in the equivariant pointed Lipschitz topology. Now since the sectional curvature is nonpositive, $`inj(x)=d_\mathrm{\Gamma }(\stackrel{~}{x})/2`$ and $`inj(x_k)=d_{\mathrm{\Gamma }_k}(\stackrel{~}{x}_k)/2`$ where $`d_\mathrm{\Gamma }(x)`$ is the infimal displacement of the point $`x`$ by the isometries of $`\mathrm{\Gamma }`$. It follows easily from the definition of equivariant Lipschitz convergence that $`d_{\mathrm{\Gamma }_k}(\stackrel{~}{x}_k)`$ converges to $`d_\mathrm{\Gamma }(x)`$ so we get a contradiction. ∎
###### 4.5.
Nikolaev’s smoothing theorem. A useful technical tool is the following smoothing theorem of I. Nikolaev \[Nik91\]. We only state an easy case of $`C^{1,\alpha }`$ metrics even though \[Nik91\] actually applies to any path metric. Let $`(M,g)`$ be a complete $`C^{1,\alpha }`$-Riemannian manifold of Alexandrov curvature $`a`$ and $`b`$. Then there exist a sequence of complete Riemannian metrics $`g_m`$ on $`M`$ with sectional curvatures within $`[a1/m,b+1/m]`$ such that $`(M,g_m)`$ converges to $`(M,g)`$ in (unpointed!) $`C^{1,\alpha }`$-topology. Also $`id:(M,g_m)(M,g)`$ and $`id:(M,g)(M,g_m)`$ are $`2^{1/m}`$-Lipschitz.
###### Proposition 4.6.
Let $`X_k`$ be a sequence of Hadamard $`n`$-manifolds with sectional curvatures in $`[a,b]`$ for $`ab<0`$ and $`n3`$. Let $`\rho _k:\pi \mathrm{Isom}(X_k)`$ be an arbitrary sequence of free and isometric actions such that each $`X_k/\rho _k(\pi )`$ has finite volume. Then after passing to a subsequence, there are points $`p_kX_k`$ such that
(i) $`(X_k,p_k,\rho _k)`$ converges to in the pointwise convergence topology to a free action $`(X,p,\rho )`$, and $`(X_k,p_k,\rho _k(\pi ))`$ converges in the equivariant pointed Lipschitz topology to $`(X,p,\rho (\pi ))`$;
(ii) If $`q_k`$, $`q`$ are projections of $`p_k`$, $`p`$, respectively, then $`(X_k/\rho _k(\pi ),q_k)`$ converges to $`(X/\rho (\pi ),q)`$ in the pointed $`C^{1,\alpha }`$ topology;
(iii) the $`C^{1,\alpha }`$-manifold $`X/\rho (\pi )`$ has finite volume, and is diffeomorphic to the interior of a compact manifold with boundary; the number of connected components of the boundary is the number of maximal parabolic subgroups of $`\pi `$.
###### Proof.
Let $`M_k=X_k/\rho _k(\pi )`$. By 3.5 and \[Bel98a, 2.7, 2.10\] we can assume by passing to a subsequence that $`(X_k,p_k,\rho _k)`$ converges to $`(X,p,\rho )`$ in pointwise convergence topology and $`(X_k,p_k,\rho _k(\pi ))`$ converges to $`(X,p,\mathrm{\Gamma })`$ in equivariant pointed Lipschitz topology. By \[Bel98a, 2.5\] we know that the $`\mathrm{\Gamma }`$-action on $`X`$ is free. By \[Fuk86, And90\] $`(M_k,q_k)`$ converges to $`(X/\mathrm{\Gamma },q)`$ in the pointed $`C^{1,\alpha }`$ topology.
Now $`M=X/\rho (\pi )`$ is a smooth manifold with a $`C^{1,\alpha }`$ Riemannian metric $`g`$ of Alexandrov curvature of $`g`$ is within $`[a,b]`$. Let $`(M,g_m)`$ be the Nikolaev’s smoothing of $`(M,g)`$ \[Nik91\]. Since $`(M,g_m)`$ is a complete Riemannian manifold of pinched negative curvature which is homotopy equivalent to $`M_k`$\[Sch84\] implies that $`(M,g_m)`$ has finite volume so \[BGS85\] implies that $`M`$ is the interior of a compact manifold with boundary and the number of boundary components is the same as the number of ends of $`M_k`$ (or, alternatively, the number of maximal parabolic subgroups of $`\pi `$).
Since $`(M,g_m)`$ converge to $`(M,g)`$ in unpointed $`C^{1,\alpha }`$-topology, $`vol(M,g_m)`$ tends to $`vol(M,g)`$. Also $`id:(M,g_m)(M,g_{m+1})`$ is Lipschitz with Lipschitz constant approaching $`1`$, hence the volumes of $`(M,g_m)`$ are uniformly bounded \[BGS85, p11\]. Thus the volume of $`(M,g)`$ is finite. Being a quotient of $`M`$, the manifold $`X/\mathrm{\Gamma }`$ also has finite volume so $`\rho (\pi )`$ has finite index in $`\mathrm{\Gamma }`$. Hence $`\rho (\pi )=\mathrm{\Gamma }`$ by 4.3. ∎
###### Remark 4.7.
Another obvious application of Nikolaev’s smoothing and the continuity of injectivity radius is that
$$\{inj_{g_k}ϵ2^{1/k}\}\{inj_gϵ\}\{inj_{g_k}ϵ2^{1/k}\}.$$
In particular, $`ϵ`$-thick part $`\{inj_gϵ\}`$ is compact since it is so for $`C^{\mathrm{}}`$ metrics.
###### Lemma 4.8.
The injectivity radius and the diameter of $`ϵ`$-thick part are bounded above and below on $`_{a,b,\pi ,n}`$.
###### Proof.
Arguing by contradiction, let $`M_k_{a,b,\pi ,n}`$ be a sequence of manifolds with the diameter of the $`ϵ`$-thick part going to infinity for some $`ϵ`$. Find points $`p_kM_k`$ such that $`(M_k,p_k)`$ subconverges to $`(M,p)`$. By 4.7 the $`ϵ`$-thick part of $`M`$ is compact hence it lies in the open ball $`B(p,R)`$ for some $`R`$. Pass to subsequence so that diameters of the $`ϵ`$-thick parts of $`M_k`$ are $`>2R+2`$; thus, for all large $`k`$, the $`ϵ`$-thick part of $`M_k`$ contains a point that lies in $`M_kB(p_k,R+1)`$. Furthermore, by continuity of injectivity radius $`B(p_k,R+1)`$ contains a point of injectivity radius $`ϵ`$ for all large $`k`$. Since $`n>2`$, the $`ϵ`$-thick part is connected, so for all large $`k`$, $`B(p_k,2R+1)B(p_k,R+1)`$ contains a point of injectivity radius $`ϵ`$. This point subconverges to a point in $`B(p,2R+2)B(p,R+1/2)`$ of injectivity radius $`ϵ`$ which is a contradiction.
Now a lower bound on the injectivity radius (and hence on the diameter of the $`ϵ`$-thick part) is provided by the Margulis lemma \[BGS85\]. As for an upper bound assume there is a sequence of manifolds $`(M_k,p_k)`$ in with $`M_k_{a,b,\pi ,n}`$ and $`inj_{p_k}>k`$. Pass to a subsequence so that $`(M_k,p_k)`$ converges to $`(M,p)`$. Since the $`ϵ`$-thick part of $`M`$ is compact, the injectivity radius of $`M`$ is bounded above which contradicts the continuity of the injectivity radius. ∎
###### Lemma 4.9.
There exist $`ϵ>0`$ such that for any $`M_{a,b,\pi ,n}`$ the $`ϵ`$-thin part of $`M`$ consists only of cusps.
###### Proof.
Arguing by contradiction find a sequence of manifolds $`M_k=X_k/\rho _k(\pi )_{a,b,\pi ,n}`$ each containing a closed geodesic of length $`<1/k`$. By 4.6 and \[Fuk86\] $`M_k`$ subconverges in the pointed Lipschitz topology. Denote the limiting manifold by $`(M,g)=X/\rho (\pi )`$.
By 4.7 the $`ϵ`$-thick part of $`(M,g)`$ contains the $`2ϵ`$-thick part of $`(M,g_1)`$. For large $`k`$ Lipschitz approximations $`\varphi _k`$ between $`M_k`$ and $`M`$ take the $`ϵ`$-thick part of $`M_k`$ Hausdorff close to the $`ϵ`$-thick part of $`M`$. Hence for all large $`k`$ the $`2ϵ`$-thick part of $`(M,g_1)`$ is contained in the $`\varphi _k`$-image of the $`ϵ`$-thick part of $`M_k`$.
Now let $`ϵ`$ be so small that $`2ϵ`$-thin part of $`(M,g_1)`$ consists only of cusps. Thus, the $`\varphi _k`$-image of any $`ϵ`$-tube of $`M_k`$ lies in a $`2ϵ`$-cusp of $`(M,g_1)`$. Hence if $`\gamma _k`$ be a closed geodesic in a tube of $`M_k`$, then $`\varphi _k(\gamma _k)`$ represents a parabolic element. By algebraic reasons, $`n>2`$ implies that $`\rho _k\rho ^1`$ takes parabolics to parabolics. So $`\gamma _k`$ represents a parabolic element which is a contradiction. ∎
## 5 Group theoretic applications
In this section $`\pi `$ is a group such that $`_{a,b,\pi ,n}`$ is nonempty.
###### Corollary 5.1.
$`Out(\pi )`$ is finite.
###### Proof.
Let $`\rho _k\mathrm{Aut}(\pi )`$ lie in different conjugacy classes. This defines a sequence $`\rho _k`$ of free isometric actions of $`\pi `$ on the universal cover $`X`$ of $`M`$. For a finite generating set $`S`$ of $`\pi `$, let $`D_k(x)=\mathrm{max}_{\gamma S}\{d(x,\rho _k(\gamma )(x))\}`$ and $`D_k=inf\{D_k(x):xX\}`$. As in \[Bel98a, 2.10\], we can choose a sequence of points $`x_kX`$ such that $`D_k(x_k)D_k+1/k`$. If $`D_k\mathrm{}`$, we get an action on an $`𝐑`$-tree hence a splitting which is impossible. So assume $`D_k(x_k)`$ is uniformly bounded.
Let $`F`$ the the Dirichlet fundamental domain for the action of $`\pi `$ on $`X`$. There exists a sequence $`\varphi _k\pi `$ such that $`\varphi _k(x_k)F`$. Then $`(X,\varphi _k(x_k),\varphi _k\rho _k\varphi _k^1)`$ subconverges in the pointwise convergence topology.
Suppose first that $`\varphi _k(x_k)`$ is precompact, so that passing to subsequence we can assume that $`\varphi _k(x_k)`$ converges to $`xX`$. Then $`(X,x,\varphi _k\rho _k\varphi _k^1)`$ also subconverges in the pointwise convergence topology, in other words, passing to subsequence, we deduce that $`\varphi _k\rho _k(\gamma )\varphi _k^1`$ converges in $`\mathrm{Isom}(X)`$ for any $`\gamma \pi `$. Since $`\pi `$ is a closed subgroup the limit lies in $`\pi `$. So $`\varphi _k\rho _k\varphi _k^1`$ converges in $`\mathrm{Hom}(\pi ,\pi )`$. Since the space $`\mathrm{Hom}(\pi ,\pi )`$ is discrete, $`\varphi _k\rho _k\varphi _k^1`$ are all equal for large $`k`$. So $`\rho _k`$ lie in the same conjugacy class for large $`k`$, a contradiction.
If $`\varphi _k(x_k)`$ is not precompact, then passing to subsequence we can assume $`\varphi _k(x_k)`$ converges a parabolic fixed point. (The closure of $`F`$ at infinity is just finitely many parabolic fixed points.) Choose a $`\pi `$-invariant set of mutually disjoint horoballs. Passing to subsequence, we assume that $`\varphi _k(x_k)`$ lies in one horoball $`H`$ for all $`k`$. Since $`\pi `$ is not virtually nilpotent, for each $`k`$ there is a generator $`\gamma _kS`$ such that the horoballs $`H`$ and $`\varphi _k\rho _k\varphi _k^1(\gamma _k)(H)`$ are disjoint. As $`S`$ is finite we can pass to to subsequence so that $`H`$ and $`\varphi _k\rho _k\varphi _k^1(\gamma _k)(H)`$ are disjoint for all $`k`$ and some $`\gamma S`$. Since $`D_k(x_k)`$ is uniformly bounded, the distance between $`\varphi _k(x_k)`$ and $`\varphi _k\rho _k(\gamma )\varphi _k^1(\varphi _k(x_k))`$ is uniformly bounded. On the other hand, this distance has to converge to infinity because it is bounded below by the distance from to $`\varphi _k(x_k)`$ to the horosphere $`H`$. This contradiction completes the proof. ∎
###### Corollary 5.2.
$`\pi `$ is cohopfian.
###### Proof.
Arguing by contradiction assume there exists an injective homomorphism $`\varphi :\pi \pi `$ which is not onto. By \[Sch84\] the manifold $`X/\varphi (\pi )`$ has finite volume hence $`\varphi (\pi )`$ must be of finite index, say $`m>1`$, in $`\pi `$. Iterating $`\varphi `$, we get a sequence of free isometric actions $`\rho _k`$ of $`\pi `$ on the universal cover $`X`$ of $`M`$ such that $`\rho _{k+1}(\pi )`$ is an index $`m`$ subgroup of $`\rho _k(\pi )`$ for each $`k`$.
Same proof as before gives that $`\varphi _k\rho _k\varphi _k^1`$ are all equal for large $`k`$. In particular, $`X/\rho _k(\pi )`$ and $`X/\rho _{k+1}(\pi )`$ are isometric for all large $`k`$. But there exists an $`m`$-sheeted cover $`X/\rho _{k+1}(\pi )X/\rho _k(\pi )`$, hence $`vol(X/\rho _{k+1}(\pi ))=mvol(X/\rho _k(\pi ))`$. Thus, $`m=1`$, a contradiction. ∎
###### Remark 5.3.
There is of course another proof that $`\pi `$ is cohopfian. Namely, by \[Sch84\] $`\varphi (\pi )`$ has finite index $`m`$ in $`\pi `$. Since $`X/\pi `$ and $`X/\varphi (\pi )`$ are properly homotopy equivalent, they have equal simplicial volumes (which are also nonzero \[Gro82\]). On the other hand, simplicial volume is multiplicative under finite covers so $`m=1`$.
## 6 Diffeomorphism finiteness
In this section we prove a theorem that implies 1.1 when combined with 4.6. Main ingredients of the proof are exponential convergence of geodesics and continuity of injectivity radius. We also use in a crucial way that the metrics converge in at least $`C^1`$ topology.
First, we need a better understanding of cusps for manifolds in $`_{a,b,\pi ,n}`$. Fix an $`ϵ(0,\mu _{n,a})`$ where $`\mu _{n,a}`$ is the Margulis constant and fix an $`ϵ`$-cusp of a manifold $`M=X/\pi _{a,b,\pi ,n}`$. Denote by $`inj_ϵ`$ the boundary of the cusp; thus $`inj_ϵ`$ is a compact topological submanifold of $`M`$ of codimension one. Set $`d_ϵ=\mathrm{diam}(inj_ϵ)`$. Also let $`H_ϵ^+`$, $`H_ϵ^{}`$ be the quotients of horospheres such that $`inj_ϵ`$ lies in the region bounded by $`H_ϵ^+`$ and $`H_ϵ^{}`$, and $`H_ϵ^+inj_ϵ`$, $`H_ϵ^{}inj_ϵ`$ are nonempty. We also assume that $`H_ϵ^{}`$ is “closer” to infinity than $`H_ϵ^+`$ (i.e. $`\mathrm{diam}(H_ϵ^{})\mathrm{diam}(H_ϵ^+)`$). Following 2.4 we let $`a=K^2`$, $`b=k^2`$.
###### Lemma 6.1.
Given $`ϵ,\sigma `$ satisfying $`0<\sigma <ϵ<\mathrm{min}\{\mu _{n,a},1\}`$, fix an $`ϵ`$-cusp of a manifold $`M=X/\pi _{a,b,\pi ,n}`$. Then the following holds:
(1) $`\mathrm{diam}(H_ϵ^+)3d_ϵ`$;
(2) $`\mathrm{diam}(inj_\sigma )2d_ϵ+k^1\mathrm{ln}\left(\frac{2ϵ}{\sigma }\right)K^1\mathrm{ln}\left(\frac{ϵ}{2\sigma }\right)`$
(3) $`K^1\mathrm{ln}\left(\frac{ϵ}{2\sigma }\right)d_ϵ\mathrm{dist}(inj_ϵ,inj_\sigma )k^1\mathrm{ln}\left(\frac{2ϵ}{\sigma }\right)+2d_ϵ`$;
(4) if $`\sigma <\frac{ϵ}{2}e^{Kd_ϵ}`$, then $`\mathrm{diam}(H_\sigma ^+)\mathrm{diam}(H_ϵ^{})`$.
###### Proof.
To prove (1) find $`x,yH_ϵ^+`$ with $`dist(x,y)=\mathrm{diam}(H_ϵ^+)`$. Let $`\stackrel{~}{x},\stackrel{~}{y}inj_ϵ`$ be the points obtained by pushing $`x,y`$ along radial geodesics. Then by the triangle inequality $`dist(x,y)dist(x,\stackrel{~}{x})+dist(\stackrel{~}{x},\stackrel{~}{y})+dist(\stackrel{~}{y},y)3d_ϵ`$ where the latter inequality holds because Busemann functions are $`1`$-Lipschitz.
It remains to prove $`(2)(4)`$. Let $`z`$ be a point of the ideal boundary of $`X`$ corresponding to the cusp under consideration. We denote by $`\stackrel{~}{inj}_\delta `$, $`\stackrel{~}{H}_\delta ^+`$, $`\stackrel{~}{H}_\delta ^{}`$ the lifts of $`inj_\delta `$, $`H_\delta ^+`$, $`H_\delta ^{}`$ to the universal cover which are in bounded distance from a horosphere about $`z`$. Let $`\gamma (t)`$ be a geodesic asymptotic to $`z`$ with $`\gamma (0)\stackrel{~}{H}_ϵ^+`$ and assume that $`\gamma (t)`$ intersects $`\stackrel{~}{inj}_\sigma `$ and $`\stackrel{~}{inj}_ϵ`$ in the points $`\gamma (t_\sigma )`$ and $`\gamma (t_ϵ)`$, respectively. Note that $`t_\sigma >t_ϵ`$ and $`t_e[0,d_ϵ]`$.
Since $`inj(\gamma (t_\sigma ))=\sigma `$, one can find $`g\pi `$ such that $`d(g(\gamma (t_\sigma )),\gamma (t_\sigma ))=2\sigma `$. Now $`d(g(\gamma (t_ϵ)),\gamma (t_ϵ))2ϵ`$. Thus, by 2.4
$$\frac{ϵ}{\sigma }=\frac{2ϵ}{2\sigma }\frac{d(g(\gamma (t_ϵ)),\gamma (t_ϵ))}{d(g(\gamma (t_\sigma )),\gamma (t_\sigma ))}2e^{K(t_\sigma t_ϵ)}\mathrm{or}K^1\mathrm{ln}\left(\frac{ϵ}{2\sigma }\right)t_\sigma t_ϵ.$$
Similarly, since $`inj(x_ϵ)=ϵ`$, one can find $`h\pi `$ with $`d(h(\gamma (t_ϵ)),\gamma (t_ϵ))=2ϵ`$. Again, $`2\sigma d(h(\gamma (t_\sigma )),\gamma (t_\sigma ))`$ so
$$\frac{ϵ}{\sigma }\frac{d(h(\gamma (t_ϵ)),\gamma (t_ϵ))}{d(h(\gamma (t_\sigma )),\gamma (t_\sigma ))}e^{k(t_\sigma t_ϵ)}/2\mathrm{or}k^1\mathrm{ln}\left(\frac{2ϵ}{\sigma }\right)t_\sigma t_ϵ.$$
So $`t_\sigma [K^1\mathrm{ln}\left(\frac{ϵ}{2\sigma }\right),d_ϵ+k^1\mathrm{ln}\left(\frac{2ϵ}{\sigma }\right)]`$. Hence, by the triangle inequality
$$\mathrm{diam}(inj_\sigma )2d_ϵ+k^1\mathrm{ln}\left(\frac{2ϵ}{\sigma }\right)K^1\mathrm{ln}\left(\frac{ϵ}{2\sigma }\right)\mathrm{and}$$
$$K^1\mathrm{ln}\left(\frac{ϵ}{2\sigma }\right)d_ϵ\mathrm{dist}(inj_ϵ,inj_\sigma )2d_ϵ+k^1\mathrm{ln}\left(\frac{2ϵ}{\sigma }\right).$$
Furthermore, if $`K^1\mathrm{ln}\left(\frac{ϵ}{2\sigma }\right)>d_ϵ`$, then $`H_\sigma ^+`$ is “closer” to infinity than $`H_ϵ^{}`$ as desired. ∎
###### Corollary 6.2.
The volume function is uniformly bounded above on $`_{a,b,\pi ,n}`$.
###### Proof.
First, show that the volume of the $`\sigma `$-thick part is uniformly bounded above on $`_{a,b,\pi ,n}`$ for any $`\sigma (0,\mu _{a,n})`$. Indeed, observe that the diameter of $`M_{[\sigma ,\mathrm{})}`$ is bounded above by 4.8. Hence, $`M_{[\sigma ,\mathrm{})}`$ is in the image of a ball in $`X`$ of some uniformly bounded above radius. By Bishop-Gromov volume comparison the volume of the ball is uniformly bounded above and the result follows because the projection $`XM`$ is volume non-increasing.
We fix an $`ϵ(0,\mu _{a,n})`$, and an arbitrary $`ϵ`$-cusp of $`M_{a,b,\pi ,n}`$. Now we seek to obtain a uniform upper bound on $`vol(H_ϵ^{})`$. Let $`\sigma =\frac{ϵ}{2}e^{K(d_ϵ+1)}`$ so that $`H_\sigma ^+`$ is “closer” to infinity than $`H_ϵ^{}`$. Let $`T`$ be the distance between between $`H_\sigma ^+`$ and $`H_ϵ^{}`$; note that $`T[d_ϵ+1,(d_ϵ+1)K/k+(\mathrm{ln}4)/k]`$. By above, the volume enclosed between $`inj_ϵ`$ and $`inj_\sigma `$ is bounded above on $`_{a,b,\pi ,n}`$ by a constant $`V`$ depending only on $`\sigma ,a,b,\pi ,n`$. The same is then true for the volume enclosed between $`H_\sigma ^+`$ and $`H_ϵ^{}`$ which is equal to $`_0^Tvol(H_t)𝑑t`$ where $`H_t`$ is the quotient of a horosphere at $`t`$-level, and $`H_0=H_ϵ^{}`$, $`H_T=H_\sigma ^+`$.
By 2.4, pushing along radial vector field gives an $`e^{Kt}`$-Lipschitz map $`H_tH_0`$. Thus, $`vol(H_0)vol(H_t)e^{Knt}`$ (in this proof we always equip $`H_t`$ with the Riemannian metric induced by the inclusion into $`M`$). Hence by the Fubini’s theorem (which applies since the Busemann function is a $`C^2`$-Riemannian submersion) we have
$$vol(H_0)T=_0^Tvol(H_0)𝑑te^{KnT}_0^Tvol(H_t)𝑑te^{KnT}V.$$
Thus $`vol(H_0)=vol(H_ϵ^{})`$ is uniformly bounded above over $`_{a,b,\pi ,n}`$.
Now pushing along radial vector field gives an $`e^{kt}`$-Lipschitz diffeomorphism $`H_0H_t`$ so $`vol(H_t)vol(H_0)e^{ktn}`$. Hence
$$_0^{\mathrm{}}vol(H_t)𝑑tvol(H_0)_0^{\mathrm{}}e^{ktn}𝑑tvol(H_0)/kn$$
and we get a uniform upper bound for the volume of the $`\sigma `$-cusp under consideration. Thus, we get a uniform upper bound on the volume of $`M`$. ∎
###### Theorem 6.3.
Let $`(M_k,g_k)`$ be a sequence of manifolds in $`_{a,b,\pi ,n}`$ such that for some $`p_kM_k`$, $`(M_k,p_k)`$ converges in the pointed $`C^{1,\alpha }`$-topology to $`(M,p)`$ equipped with a $`C^{1,\alpha }`$ Riemannian metric $`g`$. Then for all large $`k`$ there exist diffeomorphisms $`f_k:MM_k`$ such that the pullback metrics $`f_k^{}g_k`$ converge to $`g`$ in $`C^{1,\alpha }`$ topology uniformly on compact subsets.
###### Proof.
Renumerate the sequence so that $`M_k`$ now come with even indices while odd indices correspond to Nikolaev’s smoothings of $`(M,g)`$. We still denote the obtained sequence by $`M_k`$. Since $`(M_k,p_k)`$ converges to $`(M,p)`$ in the pointed $`C^{1,\alpha }`$-topology there are $`1/k`$-Lipschitz smooth embeddings $`\varphi _k:B_{1/k}(p_k)M`$ with $`d(p,\varphi _k(p_k))1/k`$ such that $`\varphi _k`$-pushforward of $`g_k`$ converges to $`g`$ in $`C^{1,\alpha }`$-topology uniformly on compact subsets. By \[Nik91\] we can take $`\varphi _{2k+1}=\mathrm{id}_M`$.
As we shall prove below, for each small enough $`ϵ`$ and large enough $`k`$, there exists a diffeomorphism $`h_{k,ϵ}:M_kM`$ that is equal to $`\varphi _k`$ when restricted to the $`ϵ`$-thick part of $`M_k`$. Note that the continuity of injectivity radius implies that, given $`ϵ`$, the maps $`\varphi _k`$ are defined on the $`1`$-neighborhood of the $`ϵ`$-thick part of $`M_k`$ for all $`kC(ϵ)`$ and some positive integer valued function $`C`$. Now the diffeomorphism $`f_k=h_{k+C(1/k),1/k}`$ enjoys the desired properties. Thus, it remains to construct $`h_{k,ϵ}`$.
Use 4.9 to find an $`ϵ^{}`$ such that $`ϵ^{}`$-thin part of each $`M_k`$ consists of cusps; we assume $`ϵ(0,\mathrm{max}\{1,ϵ^{}/10\})`$. Fix a cusp of $`M`$ and the corresponding cusps of $`M_k`$. Using 6.1, we can make so $`ϵ`$ is so small that
$``$ $`H_{k,ϵ^2}^+`$ is closer to infinity than $`H_{k,ϵ}^{}`$, and $`dist(H_{k,ϵ^2}^+,H_{k,ϵ}^{})>10`$;
$``$ $`H_{k,ϵ^3}^+`$ is closer to infinity than $`H_{k,ϵ^2}^{}`$, and and $`dist(H_{k,ϵ^3}^+,H_{k,ϵ^2}^{})>10`$.
(Here subindex $`k`$ indicates that the quotients of the horospheres lie in a cusp of $`M_k`$.) Let $`R_k`$ be the radial vector field defined on the $`ϵ`$-thin part of $`M_k`$, and let $`D_k`$ be the $`1`$-neighborhood of the region between $`inj_{k,ϵ^2}`$ and $`inj_{k,ϵ^3}`$.
Using exponential convergence of geodesics, one can easily see that for any $`\alpha (0,\pi /2)`$ there exists $`r`$, depending only on $`\alpha ,a,b,n`$ and independent of $`k`$, such that for any $`xD_k`$ and any $`y`$ lying in the same $`ϵ^2`$-cusp of $`M_k`$ with $`d_k(x,y)r`$, the angle at $`x`$ formed by $`R_k`$ and the tangent vector to the geodesic segment $`[y,x]`$ is $`\alpha `$. Fix $`\alpha =\pi /3`$ and fix the corresponding $`r`$.
Assume $`k`$ is large enough so that the embeddings $`\varphi _k`$ are defined on the $`2r`$-neighborhood of $`D_k`$. By continuity of injectivity radius the domains $`\varphi _k(D_k)M`$ converge to some compact set in Hausdorff topology and we can find a smooth domain $`D`$ which is Hausdorff close to the set and satisfies $`D\varphi _k(D_k)`$ for all large $`k`$. Now use $`\psi _k=\varphi _k^1`$ to pullback all metrics to $`D`$.
We want to show that the region bounded by $`\varphi _k(H_{k,ϵ^2}^{})`$ and $`\varphi _{k+1}(H_{k+1,ϵ^3}^+)`$ is diffeomorphic to $`H_{k,ϵ^2}^{}\times [0,1]`$. It suffices to produce a $`C^{\mathrm{}}`$ nowhere vanishing vector field on $`D`$ that is transverse to both $`\varphi _k(H_{k,ϵ^2}^{})`$ and $`\varphi _{k+1}(H_{k+1,ϵ^3}^+)`$. We shall construct such a vector field as a controlled approximation of $`\psi _k^\mathrm{\#}R_k`$, the $`\psi _k`$-pullback of the radial vector field $`R_k`$. Choose a harmonic atlas $`\{B_j\}`$ on $`D`$ as in \[And90\] in which the $`\psi _k`$-pullbacks of $`g_k`$ converge to $`g`$ in $`C^{1,\alpha }`$-topology. Fix a partition of unity associated with the atlas. Let $`yM`$ be a point in the same cusp and $`d(y,D)r+1`$.
Now we construct a $`C^{\mathrm{}}`$ vector field $`X_k`$ on $`D`$ by defining it on each chart neighborhood $`B_j`$ and then gluing via the partition of unity. Look at the $`2r`$-neighborhood of $`D`$ equipped with the metric $`\psi _k^\mathrm{\#}g_k`$. The preimage of $`B_j`$ under $`exp_{\psi _k^\mathrm{\#}g_k}:T_yMM`$ is the disjoint union of copies of $`B_j`$. Pick a copy closest to the origin, join each of its points to the origin by rays, and then project the rays to $`M`$ via $`exp_{\psi _k^\mathrm{\#}g_k}`$. Now the tangent vectors at the endpoints of the obtained geodesic segments joining $`y`$ with points of $`B_j`$ form a vector field on $`B_j`$. Gluing these local data via the partition of unity gives $`X_k`$. Note that by construction $`X_k`$ is a nowhere vanishing vector field such that the angle formed by $`X_k`$ and the exterior normal $`\psi _k^\mathrm{\#}R_k`$ to $`\varphi _k(H_{k,ϵ^2}^{})`$ or $`\varphi _k(H_{k,ϵ^3}^+)`$ is within $`(0,\pi /3]`$.
Now on each chart $`X_k`$ is a solution of the geodesic equation. Since metrics converge in $`C^{1,\alpha }`$ topology, Christoffel symbols converge in at least $`C^0`$ topology, so by standard ODE results \[Rei71, I.5.8\] $`X_k`$ converges in $`C^0`$-topology to some $`C^0`$ vector field $`X`$. So the angle, measured in the metric $`\psi _k^\mathrm{\#}g_k`$, formed by $`X_m`$ and $`\psi _k^\mathrm{\#}R_k`$ is within $`(0,c]`$ for all $`m,k`$ large enough, and some $`c[\pi /3,\pi /2)`$. Now a standard differential topology arguments implies that the region between $`\varphi _k(H_{k,ϵ^2}^{})`$ and $`\varphi _{k+1}(H_{k+1,ϵ^3}^+)`$ is diffeomorphic to $`H_{k,ϵ^2}^{}\times [0,1]`$ as needed.
Finally, we are ready to define $`h_{k,ϵ}`$. Let $`M_{k,\delta }`$ be the compact manifold obtained from $`M_k`$ by chopping off cusps along all surfaces $`H_{k,\delta }^{}`$; we think of $`M_{k,\delta }`$ as a bounded domain in $`M_k`$. Define $`h_{2k+1,ϵ}=\mathrm{id}_M`$, and define $`h_{2k,ϵ}`$ as the following composition. First, map $`M_{2k}`$ diffeomorphically to the interior of $`M_{2k,ϵ^2}`$ by a map which is the identity on the $`1`$-neighborhood of $`M_{2k,ϵ}`$. Then map $`M_{2k,ϵ^2}`$ to $`M`$ by $`\varphi _{2k}`$. Next use the above argument to map $`\varphi _{2k}(M_{2k,ϵ^2})`$ diffeomorphically onto $`\varphi _{2k+1}(M_{2k+1,ϵ^3})=M_{2k+1,ϵ^3}`$ by a map which is the identity on the $`1`$-neighborhood of $`M_{2k+1,ϵ}`$. Last, map $`M_{2k+1,ϵ^3}`$ diffeomorphically to $`M_{2k+1}=M`$, again keeping $`M_{2k+1,ϵ}`$ fixed. ∎
## 7 Pinching
In this section we prove several pinching results that follow from 1.1.
###### Corollary 7.1.
Given a group $`\pi `$, there exists $`ϵ=ϵ(\pi )>0`$ such that any finite volume manifold from $`_{n,1ϵ,1,\pi }`$ is diffeomorphic to a real hyperbolic manifold.
###### Proof.
Arguing by contradiction find a sequence $`M_k`$ of manifolds with fundamental group isomorphic to $`\pi `$ and sectional curvatures within $`[11/k,1]`$. By the main theorem, we can assume that $`(M_k,g_k)`$ converges in $`C^{1,\alpha }`$ topology to a $`C^{1,\alpha }`$-Riemannian manifold $`(M,g)`$. Now the universal cover $`X`$ of $`M`$ is a complete $`C^{1,\alpha }`$-Riemannian manifold of Alexandrov curvature $`1`$. By \[Ale57\] $`X`$ is isometric to the real hyperbolic space and we are done. (Alternatively, one can repeat the argument below appealing to \[Gao90\] to deduce that $`g`$ is a $`C^{\mathrm{}}`$ metric.) ∎
###### Remark 7.2.
It follows from \[Gro82\] that $`ϵ`$ in 7.1 only depends on the simplicial volume of $`M`$ when $`n>3`$. This is formally stronger that dependence on $`\pi _1(M)`$, yet it actually amounts to the same thing because the bounds $`M<V`$, $`asec(M)b<0`$ imply finitely many possibilities for $`\pi _1(M)`$ \[Fuk84\]. Furthermore, if $`n`$ is even and $`>3`$, then $`ϵ`$ in 7.1 only depends on the Euler characteristic $`\chi (M)`$. Indeed, it is well known that there is a positive constant $`C_n`$ such that if $`1C_nsec(M)1`$, then the Gauss-Bonnet integrand $`\chi `$ satisfies $`C_{n,1}\omega \chi C_{n,2}\omega `$ for some constants $`C_{n,1}`$ ,$`C_{n,2}`$ and the Riemannian volume form $`\omega `$. Now the Gauss-Bonnet formula (generalized to finite volume case in \[CG85\]) shows that bounds on $`vol(M)`$ and $`\chi (M)`$ are equivalent.
###### Remark 7.3.
Note that if $`n=4`$, then the Gauss-Bonnet formula (again see \[CG85\] in noncompact case) gives a stronger version of 1.1. Namely, the class of finite volume Riemannian $`4`$-manifolds with sectional curvatures within $`[a,b]`$ where $`ab<0`$ and with uniformly bounded Euler characteristics is precompact in the $`C^{1,\alpha }`$-topology. Indeed, it was shown by Milnor \[Che55\] that in dimension four the Gauss-Bonnet integrand $`\chi `$ satisfies $`\chi 3b^2\omega `$ where $`\omega `$ is the volume form and $`b<0`$ is the upper curvature bound. Hence the Gauss-Bonnet formula implies that, for any closed Riemannian $`4`$-manifold $`M`$ with sectional curvatures within $`[a,b]`$ where $`b<0`$, volume is bounded by the Euler characteristic: $`\chi (M)3b^2vol(M)`$.
Berger’s classification of holonomy groups implies the following possibilities for the restricted holonomy group of a complete negatively curved $`n`$-manifold: $`\mathrm{𝐒𝐎}(n)`$ (generic case), $`𝐔(n/2)`$ (Kähler), $`\mathrm{𝐒𝐩}(1)\mathrm{𝐒𝐩}(n/4)`$ (quaternionic-Kähler), and $`\mathrm{𝐒𝐩𝐢𝐧}(9)`$ (Cayley). Any complete negatively curved manifold with restricted holonomy group $`\mathrm{𝐒𝐩𝐢𝐧}(9)`$ is a quotient of the Cayley hyperbolic plane \[Bes88, 10.96.VI\]. The following uniformization theorem is due to Yeung \[Yeu95\]. Let $`M`$ be a complete finite volume Riemannian manifold that is either Kähler or quaternionic-Kähler. In case $`M`$ is noncompact assume also that sectional curvatures of $`M`$ are within $`[a,b]`$ for some $`ab<0`$ with $`a/b(n1)^4`$. If $`M`$ is homotopy equivalent to a complete pointwise quarter-pinched negatively curved Riemannian manifold, then $`M`$ is locally symmetric.
Combining 1.1, the Yeung’s theorem, some results of Gao \[Gao90\], and some linear algebra of Kähler curvature tensor, we get the following
###### Corollary 7.4.
Given a group $`\pi `$, there exists $`ϵ=ϵ(\pi )>0`$ such that
(1) any finite volume Kähler manifold from $`_{n,4ϵ,1,\pi }`$ is diffeomorphic to a scalar multiple of a complex hyperbolic manifold.
(2) any quaternionic-Kähler manifold from $`_{n,4ϵ,1,\pi }`$ is diffeomorphic to a scalar multiple of a quaternionic hyperbolic manifold.
###### Proof.
Arguing by contradiction find a sequence $`M_k`$ of finite volume Kähler or quaternionic-Kähler manifolds with fundamental group isomorphic to $`\pi `$ and sectional curvatures within $`[41/k,1]`$. By 1.1, we can assume that $`(M_k,g_k)`$ converges in $`C^{1,\alpha }`$ topology to a $`C^{1,\alpha }`$-Riemannian manifold $`(M,g)`$.
First, assume that each $`M_k`$ is quaternionic-Kähler. Any quaternionic-Kähler is Einstein so $`M`$ is smooth Einstein manifold and convergence is in $`C^{\mathrm{}}`$ topology. Hence, sectional curvatures of $`M_k`$ converge to sectional curvatures of $`M`$, so $`M`$ is quarter-pinched, and we are done by \[Yeu95\].
Second, assume that each $`(M_k,g_k)`$ is Kähler. Adapting the argument in \[Ber60\] to negative curvature, we deduce that the holomorphic sectional curvatures of $`M_k`$ converge to $`1`$ uniformly on $`M_k`$. Look at the ”curvature $`4`$-tensor” $`R_k^0`$ for Kähler metric of holomorphic sectional curvature $`1`$ defined in terms of $`g_k`$ and almost complex structure $`J_k`$ of $`M_k`$ \[KN69, IX.7, right before 7.2\]. Since sectional curvature can be written in terms of holomorphic sectional curvature \[BG64\] and the curvature $`4`$-tensor can be written in terms of sectional curvature, the $`4`$-tensor $`R_k`$ of $`g_k`$ is getting close to $`R_k^0`$ uniformly on $`M`$ when $`k\mathrm{}`$. Taking traces we conclude that Ricci tensor of $`g_k`$ is getting close to $`(n+1)g_k/2`$ (see \[KN69, IX.7.5\]). Now $`g_k`$ subconverges to a $`C^{1,\alpha }`$-Riemannian metric on a finite volume manifold $`M`$ hence $`(n+1)g_k/2Ric(g_k)`$ converges to zero. Then the proof of \[Gao90, theorem 0.4\] implies that the limiting metric $`g`$ is a weak solution of the Einstein equation, hence $`g`$ is a $`C^{\mathrm{}}`$ Einstein metric. Also $`(M,g)`$ has Alexandrov curvature within $`[4,1]`$ hence $`sec(M,g)[4,1]`$ and we are done by \[Yeu95\]. ∎
###### Remark 7.5.
It is not clear whether the assumptions of 7.4 are necessary. However, there do exist compact negatively curved Kähler $`4`$-manifolds which are not homotopy equivalent to locally symmetric manifolds \[MS80\]. Also in higher dimensions there are examples of compact almost quarter pinched Riemannian manifolds which are homeomorphic but not diffeomorphic to complex hyperbolic manifolds \[FJ89a\].
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# Random Resistor-Diode Networks and the Crossover from Isotropic to Directed Percolation
## I Introduction
Random resistor-diode networks (RDN) were introduced by Redner . Nevertheless, they were already contained implicitly in the pioneering work of Broadbent and Hammersley on percolation. RDN define a percolation model (for a recent introduction to percolation see Stauffer and Aharony ) on a $`d`$-dimensional hypercubic lattice in which nearest-neighbor sites are connected by a resistor, a positive diode (conducting only in a preferred direction), a negativ diode (conducting only opposite to the preferred direction), or an insulator with respective probabilities $`p`$, $`p_+`$, $`p_{}`$, and $`q=1pp_+p_{}`$. In the three dimensional phase diagram (pictured as a tetrahedron spanned by the four probabilities) one finds a nonpercolating and three percolating phases. The percolating phases are isotropic, positively directed, or negatively directed. Between the phases there are surfaces of continuous transitions. All four phases meet along a multicritical line, where $`0r:=p_+=p_{}1/2`$ and $`p=p_c(r)`$. On the entire multicritical line, i.e., independently of $`r`$, one finds the scaling properties of usual isotropic percolation ($`r=0`$).
About 20 years ago Redner studied the phase diagram sketched in Fig. 1 as well as geometrical properties of RDN in two dimensions. He used real-space renormalization methods and planar lattice duality .
Recently Inui et al have measured the order parameter exponent $`\beta `$ for the special case of a two-dimensional random diode network (with $`p=q=0`$) by Monte Carlo methods combined with series expansions. At the symmetric critical point $`p_+=p_{}=1/2`$ (so called random Manhattan (RM)), they found $`\beta =0.1794\pm 0.008`$, which does not coincide with the known values wether for isotropic percolation (IP), $`\beta _{IP}=5/36`$, nor directed percolation (DP), $`\beta _{DP}=0.27643`$. Therefore, they concluded that the percolation properties of the random diode network constitute a new universality class different from isotropic and directed percolation.
In this paper we study RDN by the methods of renormalized field theory. Contrary to Inui et al. we find that random diode networks at the percolation point as well as RDN at the full multicritical line belong to the universality class of isotropic percolation. The variable $`r=p_+=p_{}`$, which maps out the multicritical line $`p=p_c(r)`$, is a redundant variable in the sense of the classification scheme of scaling variables by Wegner . Thus, all the points of this line are equivalent to the usual isotropic percolation point with $`r=0`$.
The case $`r>0`$, in which positive and negative diodes are distributed with equal probability, leads to a breaking of isotropy and to elongated percolating clusters. However, this symmetry breaking can be easily compensated in the mesoscopic field theoretic formulation by a simple rescaling of the length scale of the preferred direction. The specific length scale is therefore redundant. Wegner showed that the renormalization flow of a redundant variable depends on the particular form of renormalization group used and does not affect the physics. Thus, in the case of Redner’s real space renormalization group, a special fixed point, the so called mixed one, is distinguished.
The situation is different for a symmetry breaking which favors not only an axis but also a direction on that axis. This leads to a relevant variable $`(p_+p_{})`$ and therefore to a new correlation length exponent $`\nu =\varphi \nu _{IP}`$. The crossover exponent $`\varphi `$ describes the beginning of the crossover to directed percolation. It also relates the order parameter exponents, $`\beta =\varphi \beta _{IP}`$. In two dimensions, it has, according to the results of Inui et al, the value $`\varphi =1.29\mathrm{}`$ . Our perturbation calculation yields $`\varphi =1.29\pm 0.05`$ in agreement with the simulations. It should be noted that the interpretation of the multicritical line as a line of equivalent fixed points with the same scaling behavior, as well as the new correlation exponent $`\nu `$, can be already found implicitly and explicitly in the early papers of Redner (his notation is $`\nu _+`$ for the new exponent $`\nu `$).
The organization of the remainder of this paper is as follows. In Sec. II, we first develop a mesoscopic field theoretic model that is capable to describe the crossover from isotropic to directed percolation, which is the basic feature of RDN. We describe the renormalization of the model and calculate the renormalization factors to two-loop order. In Sec. III, we introduce the renormalization group equation for the model and derive the general asymptotic scaling properties. In Sec. IV, we derive an interpolating formula for the crossover exponent from the two-loop $`\epsilon `$-expansion and show that it reproduces the result of Inui et al. In Sec. V we reconsider the theory by Frey, Täuber, and Schwabl for the crossover from isotropic to directed percolation and clear up some minor shortcomings. In Sec. VI, we give some conclusions and summarize our work. In an appendix, we present the two-loop calculation of the renormalization factors.
## II The Field Theoretic Model and its Renormalization
Here we develop a mesoscopic model that is capable to describe the crossover from isotropic to directed percolation. In this paper we are only interested in connectivity properties of the percolating system. In contrast to earlier work on random resistor networks we neglect transport properties as the conductance etc. In other words: all we ask is wether two points on the lattice are connected or not. Formally, we consider the limit of zero resistance of the conducting elements.
Percolating clusters in space and time can be generated by a stochastic spreading process known as the general epidemic process (GEP) . In order to apply field-theoretic methods , it is convenient to use the path-integral representation of the underlying stochastic process $`s(𝐱,t)`$ . With the imaginary-valued response field denoted by $`\stackrel{~}{s}(𝐱,t)`$, the generating functional of the Greens functions, the connected response, and correlation functions takes the form
$$𝒲[H,\stackrel{~}{H}]=\mathrm{ln}𝒟(\stackrel{~}{s},s)\mathrm{exp}\left[𝒥[\stackrel{~}{s},s]+d^dx𝑑t\left(Hs+\stackrel{~}{H}\stackrel{~}{s}\right)\right].$$
(1)
The dynamic functional $`𝒥[\stackrel{~}{s},s]`$ and the functional measure $`𝒟(\stackrel{~}{s},s)`$ ($`𝒟(\stackrel{~}{s},s)`$ is a symbolic notation for $`_{𝐱,t}\left(d\stackrel{~}{s}(𝐱,t)ds(𝐱,t)\right)`$ times a constant) are understood to be defined using a prepoint (Ito) discretization with respect to time . The prepoint discretization leads to the causality rule $`\theta (t0)=0`$ in response functions. This causality rule then forbids response propagator loops in the diagrammatical perturbation expansion.
Using the path-integral formulation a renormalized field-theory of dynamic isotropic percolation can be gained from the GEP . If we add to this model a relevant coupling which breaks isotropy and introduce a preferred direction $`𝐧`$ for the spreading of the disease in the $`d`$-dimensional space, we get the dynamic functional
$$𝒥=d^dx\left\{𝑑t\stackrel{~}{s}\left[_t+\lambda \left(\tau ^2+2v\left(𝐧\right)\right)\frac{\lambda g}{2}\stackrel{~}{s}\right]s+\frac{\lambda ^2g}{2}\stackrel{~}{s}_t\left[^t𝑑t^{}s(t^{})\right]^2\right\}.$$
(2)
This functional corresponds to the Langevin equations
$`_ts(𝐱,t)`$ $`=`$ $`\lambda \left(^22v\left(𝐧\right)\tau gn(𝐱,t)\right)s(𝐱,t)+\zeta (𝐱,t),`$ (3)
$`n(𝐱,t)`$ $`=`$ $`\lambda {\displaystyle _{\mathrm{}}^t}s(𝐱,t^{})𝑑t^{},`$ (4)
$`\zeta (𝐱,t)\zeta (𝐱^{},t^{})`$ $`=`$ $`\lambda gs(𝐱,t)\delta (𝐱𝐱^{})\delta (tt^{}),`$ (5)
for the GEP with the suitable scaled density $`s(𝐱,t)`$ of the infected individuals. Here $`n(𝐱,t)`$ constitutes the density of the immun (or dead) individuals and $`\zeta (𝐱,t)`$ is a Gaussian noise which is zero in spatial regions where the disease is extinguished. The deterministic drift of the disease in space is represented by the flow 2$`\lambda v𝐧s(𝐱,t)`$.
Of course, if the $`d`$-dimensional rotational symmetry is broken to $`(d1)`$ -dimensional isotropy, the diffusion constants for longitudinal (with respect to the direction $`𝐧`$) and transversal spreading are in general different. Thus we have to consider a more general diffusion operator $`\lambda ^2`$ $`\lambda (_{}^2+c^2_{}^2)`$. However, it is easy to see that the new parameter $`c`$ can be absorbed into the definition of the longitudinal length scale by $`cx_{}x_{}`$ followed by an appropriate change of the densities and the coupling constant $`g`$. After that, diffusional spreading looks isotropic again. The parameter $`c`$ depends on the microscopic model. It is a redundant variable in the sense of Wegner and is responsible for the multicritical line in the RDN . From the microscopic RDN-standpoint, the three variables $`\tau `$, $`c`$, $`v`$ are analytical functions of the three probabilities $`p`$, $`p_+`$, $`p_{}`$ for the resistors and diodes near the critical manifolds of RDN and share their spatial symmetries. Thus one has $`\tau (p,p_+,p_{})=\tau (p,p_{},p_+)`$, $`c(p,p_+,p_{})=c(p,p_{},p_+)`$, and $`v(p,p_+,p_{})=v(p,p_{},p_+)`$. In particular we have $`c(p,0,0)=1`$, but in general $`c(p,r,r)1`$ if $`r>0`$. Moreover, $`v(p,r,r)=0`$ holds. Remember that $`p_+=p_{}=r`$ and $`p=p_c(r)`$ defines the multicritical line.
In this paper we are interested only in the static behavior of the process, i.e., in the statistics of the distributions of immunes in space $`n(𝐱,\mathrm{})`$ after a long time when the epidemic is extinguished. These distributions constitute isotropic percolating clusters at the critical point of the GEP, which is given by $`v=v(p_c(r),r,r)=0`$ and $`\tau =\tau (p_c(r),r,r)=\tau _c=0`$ if we neglect fluctuation corrections. The density $`n(𝐱,\mathrm{})`$ is proportional to the Fourier transform of $`s`$ for frequency zero. The statistical weight for these frequency-zero modes can be found by invoking the formal limit
$$\stackrel{~}{s}(𝐱,t)\stackrel{~}{\phi }(𝐱)=\text{const. },\lambda _{\mathrm{}}^{\mathrm{}}𝑑ts(𝐱,t)=n(𝐱,\mathrm{})\phi (𝐱),$$
(6)
for the fields in the dynamic functional $`𝒥`$, Eq. (2). This manipulation can be controlled term by term in a diagrammatical perturbation expansion and leads to the wanted representation of the zero frequency Green’s functions as functional integrals with a weight $`\mathrm{exp}()`$ and a quasistatic Hamiltonian
$$=d^dx\left\{\stackrel{~}{\phi }\left[\tau ^2+2v\left(𝐧\right)+\frac{g}{2}\left(\phi \stackrel{~}{\phi }\right)\right]\phi \right\}.$$
(7)
Remember that closed loops of the response propagators are forbidden in the Feynman diagrams as a consequence of the causality rule. Therefore, the functional $``$ is, up to a rescaling, identical with the statistical functional considered by Frey, Täuber, and Schwabl in their work on the crossover from isotropic to directed percolation.
To absorb ultraviolet divergencies in a perturbational calculation of the Green’s functions with the Hamiltonian $``$ we use, in the case $`v=0`$, the following renormalization scheme :
$`\stackrel{~}{\phi }`$ $``$ $`\stackrel{˚}{\stackrel{~}{\phi }}=Z^{1/2}\stackrel{~}{\phi },\phi \stackrel{˚}{\phi }=Z^{1/2}\phi ,`$ (8)
$`\tau `$ $``$ $`\stackrel{˚}{\tau }=Z^1Z_\tau \tau \mu ^2+\stackrel{˚}{\tau }_c,g^2\stackrel{˚}{g}^2=G_\epsilon ^1Z^3Z_uu\mu ^\epsilon .`$ (9)
Here $`\epsilon =6d`$, $`G_\epsilon =\mathrm{\Gamma }(1+\epsilon /2)/(4\pi )^{d/2}`$, and $`\mu `$ is the usual external momentum scale, which makes the renormalized coupling constant $`u`$ dimensionless. Note that the fields $`\stackrel{~}{\phi }`$ and $`\phi `$ are renormalized by the same $`Z`$-factor as a consequence of the reflection symmetry $`\phi (𝐱_{},x_{})\stackrel{~}{\phi }(𝐱_{},x_{})`$ of $``$, which leads eventually to the equality $`\beta ^{}=\beta `$ between the exponents characterizing the particle density and the percolation probability. The renormalizations are known from percolation field theory up to three-loop order . Using dimensional regularization and minimal subtraction (minimal renormalization) together with the $`\epsilon `$-expansion one finds $`\stackrel{˚}{\tau }_c=0`$ and to two-loop order
$`Z`$ $`=`$ $`1+{\displaystyle \frac{u}{6\epsilon }}+\left(11{\displaystyle \frac{37}{12}}\epsilon \right)\left({\displaystyle \frac{u}{6\epsilon }}\right)^2+O\left(u^3\right),`$ (10)
$`Z_\tau `$ $`=`$ $`1+{\displaystyle \frac{u}{\epsilon }}+\left(9{\displaystyle \frac{47}{12}}\epsilon \right)\left({\displaystyle \frac{u}{2\epsilon }}\right)^2+O\left(u^3\right),`$ (11)
$`Z_u`$ $`=`$ $`1+{\displaystyle \frac{4u}{\epsilon }}+\left(15{\displaystyle \frac{59}{12}}\epsilon \right)\left({\displaystyle \frac{u}{\epsilon }}\right)^2+O\left(u^3\right).`$ (12)
If $`v0`$, further renormalizations are needed. Because $`v\mu `$, this relevant parameter has positive naive dimension like $`\tau `$. Hence, we consider it as a soft variable (which means that the renormalization constants do not depend on $`\tau `$ and $`v`$) as long as $`v`$ is finite and complete the renormalization scheme by
$$\tau \stackrel{˚}{\tau }=Z^1\left(Z_\tau \tau +Y_{v\tau }v^2\right)\mu ^2,v\stackrel{˚}{v}=Z^1Z_vv\mu .$$
(13)
Here we anticipate that $`\tau `$ and $`v^2`$ get mixed under renormalization.
To calculate $`Z_v`$ und $`Y_{v\tau }`$, we need the (unrenormalized) propagator $`G(𝐪)=\phi _𝐪\stackrel{~}{\phi }_𝐪_0^{(trunc)}`$. We define the spatial Fourier transform by $`\phi (𝐱)=_𝐪\phi _𝐪\mathrm{exp}(i𝐪𝐱)`$ (with the abbreviation $`_𝐪=(2\pi )^dd^dq`$) and get
$$G(𝐪)=\frac{1}{\tau +q^2+2i𝐯𝐪}=\frac{1}{\left(\tau +v^2\right)+\left(𝐪+i𝐯\right)^2},$$
(14)
where we used the notation $`𝐯=v𝐧`$. The one-loop contribution to the vertex function $`\mathrm{\Gamma }_{1,1}`$ (the amputated one-particle irreducible Green’s function with one $`s`$\- and one $`\stackrel{~}{s}`$-leg) is given by
$`\mathrm{\Sigma }^{(1)}(𝐪)`$ $`=`$ $`{\displaystyle \frac{g^2}{2}}{\displaystyle _𝐤}G(𝐤+𝐪/2)G(𝐤+𝐪/2)`$ (15)
$`=`$ $`{\displaystyle \frac{g^2}{2}}{\displaystyle _0^{\mathrm{}}}ds_1ds_2{\displaystyle _𝐤}\mathrm{exp}[(s_1+s_2)(\tau +v^2)`$ (17)
$`s_1(𝐤+(𝐪/2+i𝐯))^2s_2(𝐤(𝐪/2+i𝐯))^2]`$
$`=`$ $`{\displaystyle \frac{2G_\epsilon g^2}{(2\epsilon )\epsilon }}\left(\tau +Q\right)^{1\epsilon /2}K_{\epsilon 4}^{(0)}\left(\sqrt{{\displaystyle \frac{v^2Q}{\tau +Q}}}\right)`$ (18)
$`=`$ $`G_\epsilon g^2\tau ^{\epsilon /2}({\displaystyle \frac{2\tau }{(2\epsilon )\epsilon }}K_{\epsilon 4}^{(0)}(v/\sqrt{\tau })+{\displaystyle \frac{Q}{\epsilon }}K_{\epsilon 2}^{(1)}\left(v\sqrt{\tau }\right)`$ (20)
$`{\displaystyle \frac{Q^2}{4}}K_\epsilon ^{(2)}(v/\sqrt{\tau })+O(Q^3)),`$
where we have defined the functions
$$K_\alpha ^{(n)}(p)=_0^1𝑑x\frac{\left(1x^2\right)^n}{\left(1+p^2x^2\right)^{1+\alpha /2}}$$
(21)
and the abbreviation
$$Q=q^2/4+i𝐯𝐪.$$
(22)
The expansion of the yet unrenormalized one-loop selfenergy $`\mathrm{\Sigma }^{(1)}(𝐪)`$ in $`v^2`$ and $`\epsilon `$ yields
$$\mathrm{\Sigma }^{(1)}(𝐪)=\frac{G_\epsilon g^2\tau ^{\epsilon /2}}{\epsilon }\left(\left(1+\frac{\epsilon }{2}\right)\tau +\frac{v^2}{3}+\frac{2i}{3}𝐯𝐪+\frac{1}{6}𝐪^2+O(\epsilon ^2)\right).$$
(23)
Using the renormalization scheme Eqs. (8,9,13 ), the renormalized vertex function $`\mathrm{\Gamma }_{1,1}`$ is found to first order in the renormalized coupling constant $`u`$ as
$$\mathrm{\Gamma }_{1,1}=\left(Z_\tau \tau +Y_{v\tau }v^2\frac{u}{\epsilon }\left(\tau +\frac{v^2}{3}\right)\right)\mu ^2+2i\left(Z_v\frac{u}{3\epsilon }\right)\mu 𝐯𝐪+\left(Z\frac{u}{6\epsilon }\right)𝐪^2+O(\epsilon ^0),$$
(24)
from which the new renormalizations can be gathered to $`O(u)`$ as
$`Z_v`$ $`=`$ $`1+{\displaystyle \frac{u}{3\epsilon }}+\left(23{\displaystyle \frac{73}{12}}\epsilon \right)\left({\displaystyle \frac{u}{6\epsilon }}\right)^2+O\left(u^3\right),`$ (25)
$`Y_{v\tau }`$ $`=`$ $`{\displaystyle \frac{u}{3\epsilon }}+\left(30{\displaystyle \frac{35}{2}}\epsilon \right)\left({\displaystyle \frac{u}{6\epsilon }}\right)^2+O\left(u^3\right).`$ (26)
Here we have included also the two-loop result calculated in the appendix.
## III Renormalization Group Equation and Scaling Behavior
Next we explore the scaling properties of percolation in the RDN system. Scaling properties describe how physical quantities will transform under a change of length scales. By the renormalization, we have introduced the arbitrary mesoscopic length scale $`\mu ^1`$. The freedom to choose $`\mu `$, keeping the unrenormalized fields and bare parameters $`\{\stackrel{˚}{\tau },\stackrel{˚}{v},\stackrel{˚}{g}\}`$, and, in cutoff regularization, the momentum cutoff $`\mathrm{\Lambda }`$ fixed, can be used to derive in a routine fashion the renormalization group (RG) equation for the connected Green’s functions
$$G_{N,\stackrel{~}{N}}(\{𝐱\})=\underset{i=1}{\overset{N}{}}\phi (𝐱_i)\underset{j=N+1}{\overset{N+\stackrel{~}{N}}{}}\stackrel{~}{\phi }(𝐱_j)^{(conn)}.$$
(27)
We denote $`\mu `$-derivatives at fixed bare parameters by $`_\mu |_0`$. From $`\mu _\mu |_0\stackrel{˚}{G}_{N,\stackrel{~}{N}}=0`$ and the renormalization scheme, Eqs. (8,9,13), which lead to $`\stackrel{˚}{G}_{N,\stackrel{~}{N}}=Z^{(N+\stackrel{~}{N})/2}G_{N,\stackrel{~}{N}}`$, we then find the renormalization group (RG) equations
$$\left[𝒟_\mu +\frac{N+\stackrel{~}{N}}{2}\gamma \right]G_{N,\stackrel{~}{N}}=0.$$
(28)
$`𝒟_\mu `$ stands for the renormalization group differential operator
$$𝒟_\mu =\mu _\mu +\beta _u_u+\left(\tau (\kappa _\tau 2)+v^2\kappa _{v\tau }\right)_\tau +v(\kappa _v1)_v.$$
(29)
Here we have introduced the Gell-Mann-Low functions
$`\beta _u`$ $`=`$ $`{\displaystyle \frac{u}{\mathrm{ln}\mu }}|_0=\left(\epsilon +3\gamma \gamma _u\right)u,`$ (30)
$`\tau \kappa _\tau +v^2\kappa _{v\tau }`$ $`=`$ $`{\displaystyle \frac{\tau }{\mathrm{ln}\mu }}|_0=\tau \left(\gamma \gamma _\tau \right)v^2\gamma _{v\tau },`$ (31)
$`v\kappa _v`$ $`=`$ $`{\displaystyle \frac{v}{\mathrm{ln}\mu }}|_0=v\left(\gamma \gamma _v\right),`$ (32)
and the Wilson-functions $`\gamma _{\mathrm{}}=\mathrm{ln}Z_{\mathrm{}}/\mathrm{ln}\mu |_0`$.
The RG equations can be solved in terms of a single flow parameter $`l`$ using the characteristics
$`l{\displaystyle \frac{d}{dl}}\overline{u}(l)`$ $`=`$ $`\beta _u(\overline{u}(l)),\overline{u}(1)=u,`$ (33)
$`l{\displaystyle \frac{d}{dl}}\overline{v}(l)`$ $`=`$ $`\overline{v}(l)\left(\kappa _v(\overline{u}(l))1\right),\overline{v}(1)=v,`$ (34)
$`l{\displaystyle \frac{d}{dl}}\overline{\tau }(l)`$ $`=`$ $`\left(\overline{\tau }\left(\kappa _\tau (\overline{u}(l))2\right)+\overline{v}^2\kappa _{v\tau }(\overline{u}(l))\right),\overline{\tau }(1)=\tau .`$ (35)
With help of these flow equations we recast Eq. (28) as
$$[l\frac{d}{dl}+\frac{N+\stackrel{~}{N}}{2}\gamma (\overline{u}(l))]G_{N,\stackrel{~}{N}}(\{𝐱\},\overline{\tau }(l),\overline{v}(l),\overline{u}(l)\},l\mu )=0.$$
(36)
Equations (33-35) describe how the parameters transform if we change the momentum scale $`\mu `$ according to $`\mu \overline{\mu }(l)=l\mu `$. Being interested in the infrared (IR) behavior of the theory, we must study the limit $`l0`$. According to Eq. (33) we expect that in this IR limit the coupling constant $`\overline{u}(l)`$ flows to a stable fixed point $`u_{}`$ with $`\beta _u=0`$. At the fixed point it is legitimate to diagonalize the part of the RG differential operator containing the relevant parameters $`\tau `$ and $`v`$. Introducing a new parameter $`\sigma `$ instead of $`\tau `$, we find
$$\left(\tau \kappa _\tau +v^2\kappa _{v\tau }\right)_\tau +v\kappa _v_v=\sigma \kappa _\tau _\sigma +v\kappa _v_v,$$
(37)
where
$$\sigma =\tau +a_{}v^2,a_{}=\frac{\kappa _{v\tau }}{\kappa _\tau 2\kappa _v}.$$
(38)
The $`\kappa _i`$ stand for the Gell-Mann-Low functions taken at the fixed point $`u_{}`$.
Using dimensional analysis in conjunction with the flow equations, we readily find the asymptotic behavior of the connected Green’s functions for $`l0`$. Neglecting nonuniversal scale factors we get
$`G_{N,\stackrel{~}{N}}(\{𝐱\},\sigma ,v,u_{},\mu )`$ $`=`$ $`l^{(N+\stackrel{~}{N})\eta /2}G_{N,\stackrel{~}{N}}(\{𝐱\},\sigma /l^{2\kappa _\tau },v/l^{1\kappa _v},u_{},\mu l)`$ (39)
$`=`$ $`(\mu l)^{(N+\stackrel{~}{N})(d2)/2}l^{(N+\stackrel{~}{N})\eta /2}G_{N,\stackrel{~}{N}}(\{l\mu 𝐱\},\sigma /l^{2\kappa _\tau },v/l^{1\kappa _v},u_{},1),`$ (40)
where the Fisher exponent $`\eta `$ is defined by $`\eta =\gamma (u_{})`$. We define the remaining exponents by
$`\nu `$ $`=`$ $`{\displaystyle \frac{1}{1\kappa _v}},\beta =\nu {\displaystyle \frac{d2+\eta }{2}},`$ (41)
$`\nu _{IP}`$ $`=`$ $`{\displaystyle \frac{1}{2\kappa _\tau }},\beta _{IP}=\nu _{IP}{\displaystyle \frac{d2+\eta }{2}},`$ (42)
and the crossover exponent $`\varphi `$ by
$$\varphi =\frac{2\kappa _\tau }{1\kappa _v}=\nu /\nu _{IP}=\beta /\beta _{IP}.$$
(43)
The appropriate choice of the flow parameter $`l`$ in the case $`\left|\sigma \right|\left|v\right|^\varphi `$ is $`l\left|v\right|^\nu `$. In this case the Green’s functions scale as
$$G_{N,\stackrel{~}{N}}(\{𝐱\})=\left|v\right|^{(N+\stackrel{~}{N})\beta }F_{N,\stackrel{~}{N}}(\{𝐱/\xi \},\sigma /\left|v\right|^\varphi )$$
(44)
with a correlation length $`\xi \mu ^1\left|v\right|^\nu `$. Otherwise, in the case $`\left|v\right|\left|\sigma \right|^{1/\varphi }`$, we choose $`l\left|\sigma \right|^\nu `$ and arrive at isotropic percolation scaling
$$G_{N,\stackrel{~}{N}}(\{𝐱\})=\left|\sigma \right|^{(N+\stackrel{~}{N})\beta _{IP}}F_{N,\stackrel{~}{N}}^{(IP)}(\{𝐱/\xi _{IP}\},v/\left|\sigma \right|^{1/\varphi }),$$
(45)
with a correlation length $`\xi _{IP}\mu ^1\left|\sigma \right|^{\nu _{IP}}`$.
## IV Results for the Crossover Exponent from the $`\epsilon `$-Expansion
In this section we derive the scaling indices. Because the $`\epsilon `$-expansions of the usual percolation exponents are well known, we concentrate on the new crossover exponent $`\varphi `$. For this purpose we need the Gell-Mann-Low functions $`\beta _u`$, $`\gamma `$, and $`\kappa _i`$ explicitly. From Eq. (30) we know that $`\beta _u`$ begins with the zero-loop term $`\epsilon u`$, the higher order terms are determined by the Wilson functions. These functions, the logarithmic derivatives of the $`Z`$-factors, are given by $`\gamma _{\mathrm{}}=\mu _\mu |_0\mathrm{ln}Z_{\mathrm{}}=\beta _u_u\mathrm{ln}Z_{\mathrm{}}`$. In minimal renormalization the $`Z`$-factors have a pure Laurent expansion with respect to $`\epsilon `$: $`Z=1+Y^{(1)}/\epsilon +Y^{(2)}/\epsilon ^2+\mathrm{}`$. It follows then recursively in the loop expansion that the Wilson functions have also a pure Laurent expansion and, because they are finite for $`\epsilon 0`$, this expansion reduces to the constant term, i.e., all $`\epsilon `$-poles have to be compensated by the logarithmic derivative. Thus, we get the Wilson functions simply from the formula $`\gamma _{\mathrm{}}=u_uY_{\mathrm{}}^{(1)}`$. Now it is easy to find these functions from the $`Z`$-factors Eqs. (10-12,25, 26). The results are
$`\gamma _u`$ $`=`$ $`4u+{\displaystyle \frac{59}{6}}u^2+O\left(u^3\right),\gamma _\tau =u+{\displaystyle \frac{47}{24}}u^2+O\left(u^3\right),`$ (46)
$`\gamma _v`$ $`=`$ $`{\displaystyle \frac{u}{3}}+{\displaystyle \frac{73}{216}}u^2+O\left(u^3\right),\gamma _{v\tau }={\displaystyle \frac{u}{3}}+{\displaystyle \frac{35}{36}}u^2+O\left(u^3\right),`$ (47)
$`\gamma `$ $`=`$ $`{\displaystyle \frac{u}{6}}+{\displaystyle \frac{37}{216}}u^2+O\left(u^3\right),`$ (48)
from which the Gell-Mann-Low functions follow as
$`\beta _u`$ $`=`$ $`\left(\epsilon +{\displaystyle \frac{7u}{2}}{\displaystyle \frac{671}{72}}u^2+O\left(u^3\right)\right)u,\kappa _\tau ={\displaystyle \frac{5u}{6}}{\displaystyle \frac{193}{108}}u^2+O\left(u^3\right),`$ (49)
$`\kappa _v`$ $`=`$ $`{\displaystyle \frac{u}{6}}{\displaystyle \frac{u^2}{6}}+O\left(u^3\right),\kappa _{v\tau }={\displaystyle \frac{u}{3}}{\displaystyle \frac{35}{36}}u^2+O\left(u^3\right).`$ (50)
From $`\beta _u=0`$ the stable fixed point value $`u_{}=\frac{2}{7}\epsilon +\frac{671}{3087}\epsilon ^2+O(\epsilon ^3)`$ is readily obtained. We finally derive the following $`\epsilon `$-expansion of the crossover exponent:
$$\varphi =2\frac{\epsilon }{7}+\frac{59\epsilon ^2}{221^3}+O\left(\epsilon ^3\right).$$
(51)
Crudely evaluating this expansion for small spatial dimensions, i.e., for $`\epsilon =3`$ or $`4`$, leads inevitably to poor quantitative predictions. Therefore, we improve the $`\epsilon `$-expansion via a rational approximation, which takes into account that for $`d=1`$ the correlation length exponents are, trivially, always equal to one. In addition, we make the hypothesis that the $`\epsilon `$-expansion can be extended up to $`\epsilon =5`$. These considerations lead then to the interpolation formula
$$\varphi =1+\left(1\frac{\epsilon }{5}\right)\left(1+\frac{2\epsilon }{35}+\frac{6767\epsilon ^2}{5021^3}\right).$$
(52)
From Eq. (52) we get values for the crossover and the order parameter exponent for $`d=2`$ to $`d=6`$. These values are summarized in the following table:
| $`d=`$ | $`2`$ | $`3`$ | $`4`$ | $`5`$ | $`6`$ |
| --- | --- | --- | --- | --- | --- |
| $`\varphi =`$ | $`\mathrm{\hspace{0.33em}1.29}`$ | $`\mathrm{\hspace{0.33em}1.52}`$ | $`\mathrm{\hspace{0.33em}1.70}`$ | $`\mathrm{\hspace{0.33em}1.86}`$ | $`\mathrm{\hspace{0.33em}2.00}`$ |
| $`\beta =`$ | $`\mathrm{\hspace{0.33em}0.18}`$ | $`\mathrm{\hspace{0.33em}0.62}`$ | $`\mathrm{\hspace{0.33em}1.09}`$ | $`\mathrm{\hspace{0.33em}1.56}`$ | $`\mathrm{\hspace{0.33em}2.00}`$ |
Of course, the calculated values of $`\varphi `$ depend slightly on the interpolation procedure, i.e., different rational approximands may be used to incorporate $`\varphi (d=1)=1`$. We learn from this numerical sorceries that the displayed numbers may have a failure of roughly $`\pm 0.05`$.
Now we come back to the question which of the order parameter exponents are seen in simulations by Inui et al. From the scaling properties of the Green’s functions $`G_{N,\stackrel{~}{N}}(\left\{𝐱\right\})`$ we have learned that this depends on the relative behavior of the relevant parameters $`\left|\sigma \right|`$ and $`\left|v\right|^\varphi `$ as functions of the microscopic probabilities for the conducting elements. Writing the deviations from the multicritical line as $`\delta p=pp_c(r)`$, $`\delta p_\pm =p_\pm r`$, we have the expansions
$`\sigma `$ $`=`$ $`a_1\left(\delta p_++\delta p_{}\right)+a_2\delta p+a_3\left(\delta p_++\delta p_{}\right)^2+a_4\delta p_+\delta p_{}+a_5\delta p^2+\mathrm{},`$ (53)
$`v`$ $`=`$ $`\left(\delta p_+\delta p_{}\right)\left(b_1+b_2\delta p+b_3\left(\delta p_++\delta p_{}\right)\right)\mathrm{},`$ (54)
where the coefficients $`a_i,b_i`$ depend on the microscopic model. In the simulations of the pure two-dimensional random diode system Inui et al set $`\delta p_+=\delta p_{}=\delta r`$ and $`\delta p=0`$. It follows that $`\sigma \delta r^2`$ and $`v\delta r`$, and consequently $`\left|\sigma \right|/\left|v\right|^\varphi \left|\delta r\right|^{2\varphi }1`$ because $`\varphi <2`$. Thus the exponent $`\beta `$ is found as the simulations clearly state. In their extended model Inui et al. set the variations of the probabilities to $`\delta p_+=\delta p_{}=\delta r`$ and $`\delta p=0`$. In this case we have $`\sigma \delta r`$ and $`v\delta r`$, and consequently $`\left|v\right|/\left|\sigma \right|^{1/\varphi }\left|\delta r\right|^{11/\varphi }1`$ because $`1/\varphi <1`$. Now the isotropic percolation exponent $`\beta _{IP}=5/36`$ is measured.
## V The Crossover from Isotropic to Directed Percolation
### A Preliminaries
The crossover theory of Frey, Täuber, and Schwabl , henceforth called FTS, starts with the quasistatic Hamiltonian $``$ stated in Eq. (7). FTS use the modification of the longitudinal length scale by the parameter $`c`$ which we have discussed in Sec. II. We have seen that $`c`$ is a redundant variable near the isotropic fixed point and hence dropped it from the onset. FTS renormalize $`c`$ also. They find that the anomalous dimension of $`c`$ is zero at the isotropic point. Moreover, they calculate a nonvanishing anomalous dimension of $`c`$ at the directed fixed point. The variable $`c`$ (or better $`1/c^2`$, which couples to the composed field $`\stackrel{~}{\phi }_{}^2\phi `$), being redundant at the isotropic fixed point, changes over to an irrelevant variable at the directed fixed point. Thus, in order to renormalize it, one has to include all the irrelevant operators of equal naive dimension and symmetry which mix under renormalization. This was overlooked in the work of FTS, thus the calculated anomalous dimension of $`1/c^2`$ is meaningless.
The technical problem of the crossover from IP to DP is, as FTS stated, that the two fixed points, which are to connect by the renormalization flow, have different upper critical dimensions, $`6`$ and $`5`$, respectively. Thus, one has from the outset the problem to renormalize the theory for general dimensions below $`6`$ and cannot use the $`\epsilon `$-expansion, which by-passes the problem that the perturbation expansion is ill-defined if one uses critical massless propagators, leading to infrared (IR) divergencies below the upper critical dimension. However, the renormalization can be accomplished in a massive theory, which avoids the IR divergencies, by using normalization conditions for the vertex functions. Asymptotic scaling properties follow from the inhomogeneous Callan-Symanzik equation.
Below the upper critical dimension, the theory is super-renormalizable. Only the vertex function $`\mathrm{\Gamma }_{1,1}`$ is UV divergent and can be renormalized by a mass shift $`\tau \tau +\delta \tau `$, where $`\delta \tau `$ absorbs the UV divergencies. If the theory is regularized dimensionally, these UV divergencies manifest themselves as poles at $`\epsilon =\epsilon _l=2/l`$, $`l=1,2,\mathrm{}`$ , where, at the corresponding spatial dimensions $`d_l=6\epsilon _l`$, only perturbational contributions to $`\mathrm{\Gamma }_{1,1}`$ of loop-order smaller then $`l`$ are superficially UV divergent. These poles can be eliminated, as long as $`v`$ is finite, by a shift $`\delta \tau =g^{4/\epsilon }(\epsilon )`$, where $``$ is meromorphic in $`\epsilon `$. After the mass shift, the theory is UV-finite. Then, however, in those dimensions corresponding to the before mentioned poles, non-analytic logarithmic behavior with respect to the coupling constant arises.
An alternative dimensional regularization formalism was presented some time ago by Schloms and Dohm , henceforth called SD. This formalism circumvents renormalization conditions and resorts to the more convenient minimal renormalization but without using the $`\epsilon `$-expansion. The key feature of that formalism is to use the inverse correlation length itself as the mass-parameter and to define the minimal renormalization at a suitable renormalization point, which introduces the usual external mass scale $`\mu `$. The consequences are relatively simple renormalization factors and a homogeneous RG equation. Hence, the SD formalism fosters higher order calculations and resummation procedures. We will use the SD formalism in the following. FTS use a massive renormalization scheme that is a mixture of renormalization conditions and a mass renormalization at an external mass scale $`\mu `$. It is unclear wether this scheme can be used consistently at higher loop-order.
### B Crossover Renormalization
We revisit the quasistatic Hamiltonian $``$, Eqn. (7), and introduce, in the case $`v0`$, new variables by
$$𝐱_{}=𝐲,x_{}=𝐧𝐱=2v\lambda t,\phi =\left|2v\right|^{1/2}s,\stackrel{~}{\phi }=\left|2v\right|^{1/2}\stackrel{~}{s},g=\left|2v\right|^{1/2}\overline{g}.$$
(55)
Then $``$ appears as
$$\overline{}=𝑑td^{d1}y\left\{\stackrel{~}{s}\left[_t\frac{1}{4\lambda v^2}_t^2+\lambda \left(\tau ^2\right)+\frac{\lambda \overline{g}}{2}\left(s\stackrel{~}{s}\right)\right]s\right\}.$$
(56)
In the limit $`v\mathrm{}`$ we arrive at the dynamic functional which describes directed percolation (DP) . We conclude that the so far unrenormalized Green’s function have the asymptotic property
$$(2v)^{(N+\stackrel{~}{N})/2}G_{N,\stackrel{~}{N}}(\{𝐱\},\tau ,v,g)\overline{G}_{N,\stackrel{~}{N}}(\{𝐱_{},\lambda t\},\tau ,\overline{g}),$$
(57)
where the $`\overline{G}_{N,\stackrel{~}{N}}`$ are the DP Green’s functions. The minimal renormalizations of the DP theory in dimensional regularization are
$`\stackrel{~}{s}`$ $``$ $`\stackrel{˚}{\stackrel{~}{s}}=\overline{Z}_v^{1/2}\stackrel{~}{s},s\stackrel{˚}{s}=\overline{Z}_v^{1/2}s,\lambda \stackrel{˚}{\lambda }=\overline{Z}_v^1\overline{Z}\lambda ,`$ (58)
$`\tau `$ $``$ $`\stackrel{˚}{\tau }=\overline{Z}^1\overline{Z}_\tau \tau \mu ^2+\stackrel{˚}{\overline{\tau }}_c,\overline{g}^2\stackrel{˚}{\overline{g}}^2=\overline{G}_{\overline{\epsilon }}^1\overline{Z}_v^1\overline{Z}^2\overline{Z}_{\overline{u}}\overline{u}\mu ^{\overline{\epsilon }},`$ (59)
where $`\overline{\epsilon }=\epsilon 1=5d=4\overline{d}`$, $`\overline{G}_{\overline{\epsilon }}=\mathrm{\Gamma }(1+\overline{\epsilon }/2)/(4\pi )^{\overline{d}/2}`$, and known as
$`\overline{Z}`$ $`=`$ $`1+{\displaystyle \frac{\overline{u}}{8\overline{\epsilon }}}+{\displaystyle \frac{\overline{u}^2}{128\epsilon }}\left({\displaystyle \frac{13}{\epsilon }}{\displaystyle \frac{31}{4}}+{\displaystyle \frac{35}{2}}\mathrm{ln}{\displaystyle \frac{4}{3}}\right)+O(\overline{u}^3),`$ (60)
$`\overline{Z}_v`$ $`=`$ $`1+{\displaystyle \frac{\overline{u}}{4\overline{\epsilon }}}+{\displaystyle \frac{\overline{u}^2}{32\epsilon }}\left({\displaystyle \frac{7}{\epsilon }}3+{\displaystyle \frac{9}{2}}\mathrm{ln}{\displaystyle \frac{4}{3}}\right)+O(\overline{u}^3),`$ (61)
$`\overline{Z}_\tau `$ $`=`$ $`1+{\displaystyle \frac{\overline{u}}{2\overline{\epsilon }}}+{\displaystyle \frac{\overline{u}^2}{2\epsilon }}\left({\displaystyle \frac{1}{\epsilon }}{\displaystyle \frac{5}{16}}\right)+O(\overline{u}^3),`$ (62)
$`\overline{Z}_u`$ $`=`$ $`1+{\displaystyle \frac{2\overline{u}}{\overline{\epsilon }}}+{\displaystyle \frac{7\overline{u}^2}{2\epsilon }}\left({\displaystyle \frac{1}{\epsilon }}{\displaystyle \frac{1}{4}}\right)+O(\overline{u}^3).`$ (63)
Note that the DP-Hamiltonian $`\overline{}`$, Eq. (56), is non-renormalizable for $`d>5`$. Thus, to study this limit we are restricted to spatial dimensions $`d5`$.
We expect that the theory can be rendered finite for all values of $`v`$, including infinity, by interpolating the renormalizations of $``$. Then it follows from the definitions (55) that the following asymptotic properties hold for $`v\mathrm{}`$:
$$Z_{\mathrm{}}(u,v,\epsilon )\overline{Z}_{\mathrm{}}(\overline{u},\overline{\epsilon }),$$
(64)
with $`\overline{u}=u/2v`$. Of course, as long as $`v`$ is finite the soft renormalizations of $``$ can be used, but for $`v\mathrm{}`$, logarithmic divergencies $`\mathrm{ln}v`$ arise for $`\overline{\epsilon }=\epsilon 1=0`$. These must be included in the renormalization factors to make the theory finite for all $`v`$, especially in the DP limit. Note that in this limit $`v`$ plays the role of a cutoff. In the following we demonstrate that the renormalization procedure suggested by SD, i.e., an extended minimal renormalization scheme without using the $`\epsilon `$-expansion, is feasible and that it allows to extract the crossover behavior of the Green’s functions.
We consider the dimensional regularized bare vertex functions $`\stackrel{˚}{\mathrm{\Gamma }}_{\stackrel{~}{N},N}(\{𝐪\};\stackrel{˚}{\tau },\stackrel{˚}{v},\stackrel{˚}{g};d)`$ as functions of the bare parameters $`\stackrel{˚}{\tau },\stackrel{˚}{v},\stackrel{˚}{g}`$ , and the momenta $`\{𝐪\}`$ at spatial dimension $`d=6\epsilon `$, $`\epsilon >0`$. The critical point $`\stackrel{˚}{\tau }_c`$ is determined by
$$\stackrel{˚}{\mathrm{\Gamma }}_{1,1}(\{\mathrm{𝟎}\};\stackrel{˚}{\tau }_c,\stackrel{˚}{v},\stackrel{˚}{g};d)=0,$$
(65)
which provides an implicit definition of
$$\stackrel{˚}{\tau }_c=\stackrel{˚}{\tau }_c(\stackrel{˚}{v},\stackrel{˚}{g};d)=\stackrel{˚}{g}^{4/\epsilon }S(\epsilon ,\stackrel{˚}{v}/\stackrel{˚}{g}^{2/\epsilon }).$$
(66)
The generalized Symanzik function has pole singularities at $`\epsilon =2/l`$, $`l=1,2,\mathrm{}`$, as long as $`x=\stackrel{˚}{v}/\stackrel{˚}{g}^{2/\epsilon }`$ is finite. An expansion yields $`S(\epsilon ,x)=S_0(\epsilon )+Y(\epsilon )x^2+O(x^4)`$ with a finite $`Y(\epsilon )`$ as long as $`\epsilon >0`$. Taking the limit behavior of the vertex functions for $`x1`$ into account, we find for $`\epsilon >1`$, the asymptotic behavior $`S(\epsilon ,x)x^{2\epsilon }\overline{P}(x^\epsilon ;\epsilon )+x^{2/(1\epsilon )}\overline{S}(\epsilon )`$, where $`\overline{P}`$ denotes an analytic function of $`x^\epsilon =\stackrel{˚}{g}^2/\stackrel{˚}{v}^\epsilon `$. For $`\epsilon >1`$ , the pole-singularities of $`\overline{P}(y;\epsilon )`$ and $`\overline{S}(\epsilon )`$ are found at $`\epsilon 1=\overline{\epsilon }=2/l`$. These pole-singularities combine to yield logarithmic divergencies in $`x`$ instead the poles in $`\overline{\epsilon }`$.
Since the theory is super-renormalizable for $`\epsilon >0`$, we know that all the functions
$$\stackrel{˚}{\mathrm{\Gamma }}_{\stackrel{~}{N},N}^{}(\{𝐪\};\stackrel{˚}{\tau }^{},\stackrel{˚}{v},\stackrel{˚}{g};d)=\stackrel{˚}{\mathrm{\Gamma }}_{\stackrel{~}{N},N}(\{𝐪\};\stackrel{˚}{\tau }^{}+\stackrel{˚}{\tau }_c(\stackrel{˚}{v},\stackrel{˚}{g};d),\stackrel{˚}{v},\stackrel{˚}{g};d)$$
(67)
are finite below six dimensions. However, because of the non-analyticity of $`\stackrel{˚}{\tau }_c`$ with respect to $`\stackrel{˚}{g}`$, they are no more expandable in the coupling constant $`\stackrel{˚}{g}`$. This can be fixed by introducing the inverse transversal correlation length as the new mass. It is defined according to
$$\xi _{}^2=m^2=\frac{\stackrel{˚}{\mathrm{\Gamma }}_{1,1}^{}(\{\mathrm{𝟎}\};\stackrel{˚}{\tau }^{},\stackrel{˚}{v},\stackrel{˚}{g};d)}{_{q_{}^2}\stackrel{˚}{\mathrm{\Gamma }}_{1,1}^{}(\{𝐪\};\stackrel{˚}{\tau }^{},\stackrel{˚}{v},\stackrel{˚}{g};d)|_{𝐪=0}}.$$
(68)
The function $`m(\stackrel{˚}{\tau }^{},\stackrel{˚}{v},\stackrel{˚}{g};d)`$ has poles only at $`d=6`$. It can be inverted to define $`\stackrel{˚}{\tau }^{}`$ as a function of $`m`$,
$$\stackrel{˚}{\tau }^{}=\stackrel{˚}{r}(m,\stackrel{˚}{v},\stackrel{˚}{g};d),$$
(69)
with $`\stackrel{˚}{r}(0,\stackrel{˚}{v},\stackrel{˚}{g};d)=0`$. The function $`\stackrel{˚}{r}`$ can be substituted into $`\stackrel{˚}{\mathrm{\Gamma }}_{\stackrel{~}{N},N}^{}`$ yielding the bare vertex functions $`\stackrel{˚}{\mathrm{\Gamma }}_{\stackrel{~}{N},N}^{\prime \prime }`$ in terms of the mass $`m`$:
$$\stackrel{˚}{\mathrm{\Gamma }}_{\stackrel{~}{N},N}^{}(\{𝐪\};\stackrel{˚}{r}(m,\stackrel{˚}{v},\stackrel{˚}{g};d),\stackrel{˚}{v},\stackrel{˚}{g};d)=\stackrel{˚}{\mathrm{\Gamma }}_{\stackrel{~}{N},N}^{\prime \prime }(\{𝐪\};m,\stackrel{˚}{v},\stackrel{˚}{g};d).$$
(70)
In the following we abbreviate $`\stackrel{˚}{\mathrm{\Gamma }}_{\stackrel{~}{N},N}^{\prime \prime }`$ by $`\stackrel{˚}{\mathrm{\Gamma }}_{\stackrel{~}{N},N}`$ for notational simplicity. The $`\stackrel{˚}{\mathrm{\Gamma }}_{\stackrel{~}{N},N}`$ are now free of dimensional singularities below $`d=6`$, and have an expansion in integer powers of $`\stackrel{˚}{g}`$. In the limit $`v\mathrm{}`$, for $`d<5`$, we deduce from Eq. (57) the asymptotic property
$$\stackrel{˚}{\mathrm{\Gamma }}_{\stackrel{~}{N},N}(\{𝐪\},m,\stackrel{˚}{v},\stackrel{˚}{g};d)(2\stackrel{˚}{v})^{(N+\stackrel{~}{N}2)/2}\stackrel{˚}{\overline{\mathrm{\Gamma }}}_{\stackrel{~}{N},N}(\{𝐪_{},\stackrel{˚}{\lambda }^1\omega \},m,\stackrel{˚}{\overline{g}};d),$$
(71)
with $`\omega =2\stackrel{˚}{\lambda }\stackrel{˚}{v}q_{}`$ and the unrenormalized DP vertex functions $`\stackrel{˚}{\overline{\mathrm{\Gamma }}}_{\stackrel{~}{N},N}`$.
Henceforth we use for simplicity the notation $`q`$ symbolically for all the momenta $`\{𝐪\}`$. It is convenient to write
$$\stackrel{˚}{\mathrm{\Gamma }}_{\stackrel{~}{N},N}(q,m,\stackrel{˚}{v},\stackrel{˚}{g};d)=m^{\delta _{_{\stackrel{~}{N},N}}}\stackrel{˚}{F}_{\stackrel{~}{N},N}(q/m,\stackrel{˚}{v}/m,\stackrel{˚}{g}/m^{\epsilon /2};d)$$
(72)
with dimensionless functions $`\stackrel{˚}{F}_{\stackrel{~}{N},N}`$ and $`\delta _{_{\stackrel{~}{N},N}}=d(d2)\left(\stackrel{~}{N}+N\right)/2`$. Expanding $`\stackrel{˚}{F}_{\stackrel{~}{N},N}`$ in a power series gives
$$\stackrel{˚}{F}_{\stackrel{~}{N},N}(x,y,z;d)=z^\sigma \underset{l,n=0}{\overset{\mathrm{}}{}}f_{\stackrel{~}{N},N}^{(l,n)}(y;d)\left(x\right)^n\left(z^2/\epsilon \right)^l,$$
(73)
where $`\sigma =0`$ if $`\left(\stackrel{~}{N}+N\right)/2`$ is an integer and $`\sigma =1`$ if not. The exponent $`l`$ denotes the loop order. The functions $`f_{\stackrel{~}{N},N}^{(l,n)}(y;d)`$ are finite for $`d6`$ and finite $`y`$. Moreover, they are analytic in $`y`$ and can be expanded to
$$f_{\stackrel{~}{N},N}^{(l,n)}(y;d)=\underset{k=0}{\overset{\mathrm{}}{}}f_{\stackrel{~}{N},N}^{(k,l,n)}\left(d\right)y^k$$
(74)
The contributions of order $`z^l`$ to the scaled vertex functions $`\stackrel{˚}{F}_{\stackrel{~}{N},N}`$ (73) exhibit infrared (IR) divergencies for $`m0`$. SD showed at the instance of the $`\mathrm{\Phi }^4`$-theory, that this problem can be treated by minimal renormalizations
$`\stackrel{˚}{\stackrel{~}{\phi }}`$ $`=`$ $`Z^{1/2}\stackrel{~}{\phi },\stackrel{˚}{\phi }=Z^{1/2}\phi ,`$ (75)
$`\stackrel{˚}{g}^2`$ $`=`$ $`A\left(d\right)G_\epsilon ^1Z^3Z_uu\mu ^\epsilon ,v\stackrel{˚}{v}=Z^1Z_vv\mu ,`$ (76)
where the soft renormalization factors $`Z_{\mathrm{}}(u;d)`$ absorb just the poles at $`d=6`$, followed by the application of the renormalization group. The UV $`\epsilon `$-poles and the IR divergencies are then summed up yielding the correct critical scaling behavior reminiscent of the intimate relation between both. Since the poles of the coefficients $`f_{\stackrel{~}{N},N}^{(k,l,n)}\left(6\epsilon \right)/\epsilon ^l`$ do not depend on any of the parameters, they are identical to the usual soft minimal renormalizations and can be calculated by setting $`m=\mu `$ or in the massless theory with $`\epsilon 0`$. Following SD, we have introduced in Eq. (76) a further amplitude $`A(d)`$ with $`A(0)=1`$. This dimension dependent function can be conveniently defined and is extremely useful for the practical calculation of scaling functions. For simplicity we set $`A(d)=1`$ in the remaining part of this article. Note that from its definition (68), the mass $`m`$ does not need a multiplicative renormalization and that $`m=\stackrel{˚}{m}`$.
Now we consider the DP vertex functions
$$\stackrel{˚}{\overline{\mathrm{\Gamma }}}_{\stackrel{~}{N},N}(q_{},\stackrel{˚}{\lambda }^1\omega ,m,\stackrel{˚}{\overline{g}};d)=m^{\overline{\delta }_{_{\stackrel{~}{N},N}}}\stackrel{˚}{\overline{F}}_{\stackrel{~}{N},N}(q_{}/m,\omega /\stackrel{˚}{\lambda }m^2,\stackrel{˚}{\overline{g}}/m^{\overline{\epsilon }/2};d),$$
(77)
with $`\delta _{_{\stackrel{~}{N},N}}=(d+1)(d1)\left(\stackrel{~}{N}+N\right)/2`$ . The power series expansion of the $`\stackrel{˚}{\overline{F}}_{\stackrel{~}{N},N}`$ reads for $`d<5`$
$$\stackrel{˚}{\overline{F}}_{\stackrel{~}{N},N}(x_{},\overline{x}_{},\overline{z};d)=\overline{z}^\sigma \underset{l,n_{},n_{}=0}{\overset{\mathrm{}}{}}\overline{f}_{\stackrel{~}{N},N}^{(l,n_{},n_{})}\left(\overline{d}\right)\left(x_{}\right)^n_{}\left(x_{}\right)^n_{}\left(\overline{z}^2/\overline{\epsilon }\right)^l,$$
(78)
with $`\overline{\epsilon }=5d=4\overline{d}`$. By the same argumentation as above, one can show that the $`\overline{\epsilon }`$-poles are cancelled and the IR divergencies are summed up by the renormalization group equation based on the minimal renormalizations $`\overline{Z}`$, $`\overline{Z}_v`$, and $`\overline{Z}_u`$, Eqs. (60,61,63), where the renormalization factors $`\overline{Z}_{\mathrm{}}(\overline{u},d)`$ now absorb just the $`\overline{\epsilon }`$ -poles at $`d=5`$. As above, they are identical to the usual soft minimal renormalizations. They can be calculated at the renormalization point $`m=\mu `$ , $`\{𝐪=0\}`$ or in the massless theory with $`\overline{\epsilon }0`$, because the functions $`\overline{f}_{\stackrel{~}{N},N}^{(l,n_{},n_{})}\left(\overline{d}\right)`$ do not depend on any of the parameters.
Combining the asymptotic properties of the vertex functions for $`y\mathrm{}`$, Eq. (71), with the various expansions, we find for $`d<5`$
$$f_{\stackrel{~}{N},N}^{(l,n)}(y;d)(2y)^{(N+\stackrel{~}{N}2\sigma )/2+n_{}l}\left(\epsilon /\overline{\epsilon }\right)^l\overline{f}_{\stackrel{~}{N},N}^{(l,n_{},n_{})}\left(\overline{d}\right).$$
(79)
In case of taking the limit $`d5`$ before $`y\mathrm{}`$, the $`\overline{\epsilon }^l`$-poles are replaced by a polynomial $`P^{(l,n_{},n_{})}(\mathrm{ln}y)`$ of order $`l`$. Therefore, it should be possible to find further finite renormalizations $`\zeta _{\mathrm{}}(u,v;d)`$ which interpolate between the two minimal renormalizations at the renormalization point $`m=\mu `$,
$$Z_{\mathrm{}}(u;d)\zeta _{\mathrm{}}(u,v;d)Z_{\mathrm{}}(u;d)=Z_{\mathrm{}}(u,v;d),$$
(80)
where the $`Z_{\mathrm{}}(u,v;d)`$ tend to the corresponding $`\overline{Z}_{\mathrm{}}(\overline{u};\overline{d})`$ with $`\overline{u}=u/2v`$ in the limit $`v\mathrm{}`$.
Following SD, the remaining problem is to determine the renormalization of the relation between the mass $`m`$ and the “temperature” $`\stackrel{˚}{\tau }=\stackrel{˚}{\tau }(m,\stackrel{˚}{v},\stackrel{˚}{g};d)=m^2\stackrel{˚}{T}(\stackrel{˚}{v}/m,\stackrel{˚}{g}/m^\epsilon ;d)`$. The dimensionless function $`\stackrel{˚}{T}`$ has an expansion
$$\stackrel{˚}{T}(y,z;d)=1+\underset{l=1}{\overset{\mathrm{}}{}}t_l(y;d)\left(z^2/\epsilon \right)^l,$$
(81)
but the $`t_l(y;d)`$ show UV singularities below $`d=6`$. Otherwise, the use of the variable $`\stackrel{˚}{\tau }^{}=\stackrel{˚}{\tau }\stackrel{˚}{\tau }_c=\stackrel{˚}{r}(m,\stackrel{˚}{v},\stackrel{˚}{g};d)=`$ $`m^2\stackrel{˚}{T}(\stackrel{˚}{v}/m,\stackrel{˚}{g}/m^\epsilon ;d)\stackrel{˚}{g}^{4/\epsilon }S(\epsilon ,v/\stackrel{˚}{g}^{2/\epsilon })`$ instead of $`\stackrel{˚}{\tau }`$ eliminates all UV singularities below the upper critical dimension. Then, however, the function $`\stackrel{˚}{r}`$ is not more expandable in $`\stackrel{˚}{g}`$. Following SD, we shall therefore consider the derivative $`\stackrel{˚}{r}/m^2|_{\stackrel{˚}{v},\stackrel{˚}{g}}`$, which is not only expandable but also free of singularities. We define the function
$$\stackrel{˚}{P}(\stackrel{˚}{v}/m,\stackrel{˚}{g}/m^\epsilon ;d)=\frac{\stackrel{˚}{r}}{m^2}|_{\stackrel{˚}{v},\stackrel{˚}{g}}=\frac{\left(m^2\stackrel{˚}{T}\right)}{m^2}|_{\stackrel{˚}{v},\stackrel{˚}{g}}=\frac{\stackrel{˚}{\tau }}{m^2}|_{\stackrel{˚}{v},\stackrel{˚}{g}}$$
(82)
which has the expansion
$$\stackrel{˚}{P}(y,z;d)=1+\underset{l=1}{\overset{\mathrm{}}{}}p_l(y;d)\left(z^2/\epsilon \right)^l.$$
(83)
The functions $`p_l(y;d)`$ are finite for $`d6`$, and $`\stackrel{˚}{P}`$ is minimally renormalized by
$$\stackrel{˚}{P}=Z^1Z_\tau P.$$
(84)
The same argumentation as above shows that it is possible to find a $`Z_\tau (u,v;d),`$ which interpolates between the respective minimal renormalizations of isotropic and directed percolation. Note that $`/m^2|_{\stackrel{˚}{v},\stackrel{˚}{g}}=/m^2|_{v,g,\mu }`$. Thus, the derivative with respect to the mass $`m`$ commutes with the multiplication with $`Z`$-factors.
### C One-Loop Crossover Calculation
In this subsection we derive a minimal crossover renormalization, the corresponding renormalization group equation, and the flow equations, which show the crossover from isotropic to directed percolation. From the self-energy to order $`g^2`$, we get the unrenormalized (for notational simplicity we drop the overcirc) vertex function $`\mathrm{\Gamma }_{1,1}`$ expanded to second order in the momenta:
$`\mathrm{\Gamma }_{1,1}`$ $`=`$ $`\tau \left(1{\displaystyle \frac{2G_\epsilon g^2\tau ^{\epsilon /2}}{(2\epsilon )\epsilon }}K_{\epsilon 4}^{(0)}\left(v/\sqrt{\tau }\right)\right)+2i𝐯𝐪\left(1{\displaystyle \frac{G_\epsilon g^2\tau ^{\epsilon /2}}{2\epsilon }}K_{\epsilon 2}^{(1)}\left(v/\sqrt{\tau }\right)\right)`$ (86)
$`+q^2\left(1{\displaystyle \frac{G_\epsilon g^2\tau ^{\epsilon /2}}{4\epsilon }}K_{\epsilon 2}^{(1)}\left(v/\sqrt{\tau }\right)\right)\left(𝐯𝐪\right)^2{\displaystyle \frac{G_\epsilon g^2\tau ^{\epsilon /2}}{4\tau }}K_\epsilon ^{(2)}\left(v/\sqrt{\tau }\right).`$
The inverse transversal correlation length is given by
$$m^2=\frac{\tau \left(1\frac{2G_\epsilon g^2\tau ^{\epsilon /2}}{(2\epsilon )\epsilon }K_{\epsilon 4}^{(0)}\left(v/\sqrt{\tau }\right)\right)}{\left(1\frac{G_\epsilon g^2\tau ^{\epsilon /2}}{4\epsilon }K_{\epsilon 2}^{(1)}\left(v/\sqrt{\tau }\right)\right)}.$$
(87)
From Eq. (87) we find easily the unrenormalized perturbational temperature as
$$\tau =m^2\left(1+\frac{G_\epsilon g^2m^\epsilon }{4\epsilon }\left(\frac{8}{(2\epsilon )}K_{\epsilon 4}^{(0)}(v/m)K_{\epsilon 2}^{(1)}(v/m)\right)\right).$$
(88)
Taking the derivative with respect to $`m^2`$ while holding $`g`$ and $`v`$ constant and renormalization by multiplication with the factor $`ZZ_\tau ^1`$ yields via the renormalization scheme the function $`P`$ as
$$P=ZZ_\tau ^1+u(\mu /m)^\epsilon \left(\frac{1}{4\epsilon }\left(4K_{\epsilon 2}^{(0)}(v\mu /m)K_{\epsilon 2}^{(1)}(v\mu /m)\right)\frac{1}{8}K_\epsilon ^{(1)}(v\mu /m)\right).$$
(89)
With the formulas
$`p^2K_\alpha ^{(1)}(p)`$ $`=`$ $`\left(1+p^2\right)K_\alpha (p)K_{\alpha 2}(p),`$ (90)
$`(1+\alpha )K_\alpha (p)`$ $`=`$ $`(2+\alpha )K_{\alpha +2}(p)\left(1+p^2\right)^{1\alpha /2},`$ (91)
where we have defined $`K_\alpha :=K_\alpha ^{(0)}`$, we get
$`P`$ $`=`$ $`ZZ_\tau ^1+{\displaystyle \frac{u}{8}}({\displaystyle \frac{20}{3\epsilon }}(1+v^2)^{\epsilon /2}+{\displaystyle \frac{2}{\epsilon 1}}(3v^2)[K_\epsilon (v)(1+v^2)^{\epsilon /2}]`$ (93)
$`+{\displaystyle \frac{1}{\epsilon 3}}(v^2[K_\epsilon (v)(1+v^2)^{\epsilon /2}]{\displaystyle \frac{2}{3}}(1+v^2)^{\epsilon /2})+K_\epsilon (v))`$
at the renormalization point $`m=\mu `$. Note that the poles for $`\epsilon =1`$ and $`\epsilon =3`$ are only fictitious because $`K_1(p)=(1+p^2)^{1/2}`$ and $`K_3(p)=(1+p^2)^{3/2}(1+2p^2/3)`$. Useful properties of the function $`K_\epsilon (p)`$ are
$`K_\epsilon (p)`$ $`=`$ $`1{\displaystyle \frac{2+\epsilon }{6}}p^2+O\left(p^4\right)`$ (94)
$`=`$ $`{\displaystyle \frac{\sqrt{\pi }\mathrm{\Gamma }((1+\epsilon )/2)}{2\mathrm{\Gamma }(1+\epsilon /2)}}p^1{\displaystyle \frac{1}{1+\epsilon }}p^{2\epsilon }+O\left(p^{4\epsilon }\right),`$ (95)
from which we find that the fictitious pole at $`\epsilon =1`$ in Eq. (93) is replaced by a divergence $`\mathrm{ln}v`$, and that $`w=uK_\epsilon (v)\overline{u}/2`$ for $`v\mathrm{}`$. Thus, a minimal crossover renormalization is given by
$$ZZ_\tau ^1=1\frac{5u}{6\epsilon }\left(1+v^2\right)^{\epsilon /2}\frac{3u}{4(\epsilon 1)}\left[K_\epsilon (v)\left(1+v^2\right)^{\epsilon /2}\right]+O\left(u^2\right).$$
(96)
Defining the crossover function
$$C_\epsilon (v)=1\left(\left(1+v^2\right)^{\epsilon /2}K_\epsilon (v)\right)^1,$$
(97)
we render the renormalized temperature derivative
$$P=1+\frac{w}{8}\left(1+\frac{1}{3\epsilon }\left(2\frac{1C_\epsilon (v)}{3}+\frac{5\epsilon }{1\epsilon }v^2C_\epsilon (v)\right)\right)$$
(98)
finite for all $`d5`$ and all $`v`$.
Now we consider the vertex function $`\mathrm{\Gamma }_{1,1}`$, Eq. (86), for $`𝐪0`$. Inserting $`\tau `$, Eq. (88), we get after renormalization
$`\mathrm{\Gamma }_{1,1}`$ $`=`$ $`\left(m^2+q^2\right)\left(Z{\displaystyle \frac{u(\mu /m)^\epsilon }{4\epsilon }}K_{\epsilon 2}^{(1)}\left(v\mu /m\right)\right)`$ (100)
$`+2i\mu 𝐯𝐪\left(Z_v{\displaystyle \frac{u(\mu /m)^\epsilon }{2\epsilon }}K_{\epsilon 2}^{(1)}\left(v\mu /m\right)\right)\left(𝐯𝐪\right)^2{\displaystyle \frac{u(\mu /m)^{2+\epsilon }}{4}}K_\epsilon ^{(2)}\left(v\mu /m\right).`$
The last term is finite for all $`d5`$ and all $`v`$, even in the limit $`v\mathrm{}`$, if one holds $`2𝐯𝐪=\omega /\lambda `$ and $`w=uK_\epsilon (v)`$ constant. The reduction of $`K_{\epsilon 2}^{(1)}(v)`$ with Eqs. (90,91) leads to
$`K_{\epsilon 2}^{(1)}(v)`$ $`=`$ $`{\displaystyle \frac{2}{3\epsilon }}\left(1+v^2\right)^{\epsilon /2}+{\displaystyle \frac{2+v^2}{2(\epsilon 1)}}\left[K_\epsilon (v)\left(1+v^2\right)^{\epsilon /2}\right]`$ (102)
$`{\displaystyle \frac{1}{\epsilon 3}}\left({\displaystyle \frac{1}{2}}v^2\left[K_\epsilon (v)\left(1+v^2\right)^{\epsilon /2}\right]{\displaystyle \frac{1}{3}}\left(1+v^2\right)^{\epsilon /2}\right).`$
We see that the divergencies which have to be absorbed by renormalization arise from the first two terms of the right hand side of Eq. (102). Thus, we define
$`Z`$ $`=`$ $`1+{\displaystyle \frac{u}{6\epsilon }}\left(1+v^2\right)^{\epsilon /2}+{\displaystyle \frac{u}{4(\epsilon 1)}}\left[K_\epsilon (v)\left(1+v^2\right)^{\epsilon /2}\right]+O\left(u^2\right),`$ (103)
$`Z_v`$ $`=`$ $`1+{\displaystyle \frac{u}{3\epsilon }}\left(1+v^2\right)^{\epsilon /2}+{\displaystyle \frac{u}{2(\epsilon 1)}}\left[K_\epsilon (v)\left(1+v^2\right)^{\epsilon /2}\right]+O\left(u^2\right).`$ (104)
Now we turn to the calculation of the one-loop contribution to the vertex function $`\mathrm{\Gamma }_{1,2}`$. To renormalize $`\mathrm{\Gamma }_{1,2}`$ we need it for zero external momenta. We find the contribution
$$V^{(1)}=2g^3_𝐤G(𝐤)^2G(𝐤)=2g\frac{}{\tau }\mathrm{\Sigma }^{(1)}(\mathrm{𝟎})=\frac{2G_\epsilon g^3}{\epsilon }\tau ^{\epsilon /2}K_{\epsilon 2}\left(v/\sqrt{\tau }\right),$$
(105)
which leads after renormalization to
$$\left(\mathrm{\Gamma }_{1,2}\right)^2=G_\epsilon ^1u\mu ^\epsilon \left(Z_u\frac{4u(\mu /m)^\epsilon }{\epsilon }K_{\epsilon 2}\left(v\mu /m\right)\right).$$
(106)
Using once more the reduction formula (91), we finally find
$$Z_u=1+\frac{4u}{\epsilon }\left(1+v^2\right)^{\epsilon /2}+\frac{4u}{(\epsilon 1)}\left[K_\epsilon (v)\left(1+v^2\right)^{\epsilon /2}\right]+O\left(u^2\right).$$
(107)
Now, all the $`Z`$-factors (96,103,104,107 ) are of the form
$$Z_i=1+u\left(\frac{a_i}{\epsilon }\left(1+v^2\right)^{\epsilon /2}+\frac{b_i}{(\epsilon 1)}\left[K_\epsilon (v)\left(1+v^2\right)^{\epsilon /2}\right]\right)+O\left(u^2\right),$$
(108)
where the coefficients are given by
$`a`$ $`=`$ $`1/6,b=1/4,a_\tau =1,b_\tau =1,`$ (109)
$`a_v`$ $`=`$ $`1/3,b_v=1/2,a_u=4,b_u=4.`$ (110)
With help of the derivative of the $`K`$-function,
$$v\frac{K_\epsilon (v)}{v}=\left(1+v^2\right)^{1\epsilon /2}K_\epsilon (v)=K_\epsilon (v)C_\epsilon (v),$$
(111)
and the Gell-Mann-Low functions
$`\beta _u`$ $`=`$ $`u/\mathrm{ln}\mu |_0=(\epsilon +3\gamma \gamma _u)u=\epsilon u+O(u^2),`$ (112)
$`\beta _v`$ $`=`$ $`u/\mathrm{ln}\mu |_0=(1+\gamma \gamma _v)v=(1+O(u))v`$ (113)
we obtain the Wilson-functions $`\gamma _i=\mathrm{ln}Z_i/\mathrm{ln}\mu |_0=(\beta _u_u+\beta _v_v)\mathrm{ln}Z_i|_0`$ easily as
$`\gamma _i`$ $`=`$ $`u\left(a_i\left(1+v^2\right)^{1\epsilon /2}+b_i\left[K_\epsilon (v)\left(1+v^2\right)^{1\epsilon /2}\right]\right)+O\left(u^2\right)`$ (114)
$`=`$ $`w\left(a_i\left(1C_\epsilon (v)\right)+b_iC_\epsilon (v)\right)+O\left(u^2\right).`$ (115)
Here we have used the new coupling constant $`w=uK_\epsilon (v)`$ and the definition (97) of the crossover function $`C_\epsilon (v)`$ . We mention that $`w=u`$ for $`v=0`$ and $`w=\overline{u}/2`$ for $`v=\mathrm{}`$, as well as $`C_\epsilon (0)=0`$ and $`C_\epsilon (\mathrm{})=1`$. $`C_\epsilon (v)`$ crosses over monotonically between this two values. Especially in three dimensions it is a simple rational function of $`v^2`$:
$$C_3(v)=\frac{v^2\left(5+2v^2\right)}{\left(1+v^2\right)\left(3+2v^2\right)}.$$
(116)
The renormalization group equation (28),
$$\left[𝒟_\mu +\frac{N+\stackrel{~}{N}}{2}\gamma \right]G_{N,\stackrel{~}{N}}=0,$$
(117)
follows as usually from the independence of the unrenormalized Green’s functions from the external mass scale $`\mu `$. Here we have to note that the mass $`m`$ is by definition (68) a function of the bare parameters only. Therefore, its Gell-Mann-Low function vanishes. Thus, we have the renormalization group differential operator
$$𝒟_\mu =\mu _\mu +\beta _w_w+\beta _v_v.$$
(118)
To one-loop order we obtain
$`\beta _w`$ $`=`$ $`\left(\epsilon C_\epsilon (v)+{\displaystyle \frac{425C_\epsilon (v)C_\epsilon (v)^2}{12}}w\right)w`$ (119)
$`\beta _v`$ $`=`$ $`\left(1+{\displaystyle \frac{2+C_\epsilon (v)}{12}}w\right)v,`$ (120)
and
$$\gamma =\frac{2+C_\epsilon (v)}{12}w.$$
(121)
The flow equations
$`l{\displaystyle \frac{d}{dl}}\overline{w}(l)`$ $`=`$ $`\beta _w(\overline{w}(l),\overline{v}(l)),\overline{w}(1)=w,`$ (122)
$`l{\displaystyle \frac{d}{dl}}\overline{v}(l)`$ $`=`$ $`\beta _v(\overline{w}(l),\overline{v}(l)),\overline{v}(1)=v,`$ (123)
and their solutions show of course the (unstable) isotropic and the (stable) directed percolation fixed points for $`v_{}=0`$ and $`v_{}=\mathrm{}`$, respectively, and the continuous crossover between both.
Upon defining the flowing amplitude function $`X(l)`$ by
$$l\frac{d}{dl}\mathrm{ln}X(l)=\gamma (\overline{w}(l),\overline{v}(l)),X(1)=1,$$
(124)
the scaling form of the Green’s functions follows from the renormalization group equation (117) as
$`G_{N,\stackrel{~}{N}}(\{𝐱\},m;w,v;\mu )`$ $`=`$ $`\left(m^{d2}X(m/\mu )\right)^{(N+\stackrel{~}{N})/2}G_{N,\stackrel{~}{N}}(\{m𝐱\};\overline{w}(m/\mu ),\overline{v}(m/\mu ))`$ (125)
$``$ $`\left(m^{d2}{\displaystyle \frac{X(m/\mu )}{2\overline{v}(m/\mu )}}\right)^{(N+\stackrel{~}{N})/2}\overline{G}_{N,\stackrel{~}{N}}(\{m𝐱_{},m^2\overline{\lambda }(m/\mu )t\};\overline{w}(m/\mu )),`$ (126)
where the last form holds asymptotically for $`\overline{v}(m/\mu )1`$, and $`\lambda /\overline{\lambda }(l)=l\overline{v}(l)/v`$.
We conclude this section with the determination of the parameter $`m`$ as a function of the temperature $`\tau `$. For the renormalized temperature derivative we have the renormalization group equation
$$\left[𝒟_\mu \kappa \right]P=0$$
(127)
with $`\kappa =\gamma \gamma _\tau `$. Its solution is given by
$$P(m/\mu ;w,v)=Y(m/\mu )P(1;\overline{w}(m/\mu ),\overline{v}(m/\mu )),$$
(128)
where the amplitude function $`Y`$ is determined by the differential equation
$$l\frac{d}{dl}\mathrm{ln}Y(l)=\kappa (\overline{w}(l),\overline{v}(l)),Y(1)=1.$$
(129)
$`P`$ comprises the main (exponential) contribution of the crossover-scaling of $`\tau /m^2|_{w,v,\mu }`$. The amplitude $`P(1;w,v)`$ is given to one-loop order by Eq. (98). The temperature $`\tau `$ results then from the solution of the differential equation $`\tau /m^2|_{w,v,\mu }=P(m/\mu ;w,v)`$. The integration of all the flow-equations has to be done numerically and leads, qualitatively, back to the numerical results of FTS .
## VI Epilogue
Using field theoretic methods, we have analyzed the connectivity behavior of random resistor-diode networks near the percolation critical point. We found that the introduction of positive and negative diodes oriented to a privileged direction in space with unequal probabilities leads to a crossover to the directed percolation problem, whereas a distribution of diodes with equal probabilities results only in elongated isotropic percolation clusters. In the latter case, a simple rescaling of the privileged direction maps the problem to isotropic percolation. A slightly different distribution of the diodes introduces a further relevant variable with a new scaling dimension. We have calculated this scaling dimension to $`O(\epsilon ^2)`$ in an $`\epsilon `$-expansion around six dimensions. An interpolation resulting from this $`\epsilon `$-expansion and an exact value at one dimension leads to a formula that compares very well with a recent simulational result in two dimensions. It would be very interesting to perform simulations also in higher dimensions in order to compare these with our result.
In Sec. V we have reconsidered the theory by Frey, Täuber, and Schwabl for the crossover from isotropic to directed percolation. Some shortcomings are corrected and it is demonstrated, how one can perform crossover calculations consistently by using a type of extended minimal renormalization.
###### Acknowledgements.
This work has been supported in part by the SFB 237 ,,Unordnung und große Fluktuationen“ of the Deutsche Forschungsgemeinschaft.
## Two-Loop Calculation
In this appendix we present briefly the main part of the two-loop calculation of the new renormalization constants $`Z_v`$ and $`Y_{v\tau }`$. The dimensionaly regularized parameter integral
$`I(a,b,c)`$ $`=`$ $`{\displaystyle _{𝐩,𝐪}}{\displaystyle \frac{1}{(a+𝐩^2)(b+𝐪^2)(c+(𝐩+𝐪)^2)}}`$ (130)
$`=`$ $`{\displaystyle \frac{G_\epsilon ^2}{6\epsilon }}(({\displaystyle \frac{1}{\epsilon }}+{\displaystyle \frac{25}{12}})(a^{3\epsilon }+b^{3\epsilon }+c^{3\epsilon })3abc`$ (132)
$`({\displaystyle \frac{3}{\epsilon }}+{\displaystyle \frac{21}{4}})(a^{2\epsilon }(b+c)+b^{2\epsilon }(a+c)+c^{2\epsilon }(b+c))),`$
introduced in Ref. , plays a fundamental role in the calculation. Its derivatives
$$I_{lmn}=\frac{(1)^{l+m+n3}}{(l1)!(m1)!(n1)!}\frac{^{l+m+n3}I(a,b,c)}{a^{l1}b^{m1}c^{n1}}|_{a=b=c=1}$$
(133)
are extensively used in the following. Particularly simple are the integrals
$$I_n=_𝐩\frac{1}{(1+𝐩^2)^n}=\frac{8(1)^{n1}G_\epsilon }{(n1)!(4\epsilon )(2\epsilon )\epsilon }\frac{^{n1}a^{2\epsilon /2}}{a^{n1}}|_{a=0}.$$
(134)
We start with the self-energy diagram displayed in Fig. 2(a). Its value is given by
$`\mathrm{\Sigma }^{(2a)}(𝐪)`$ $`=`$ $`{\displaystyle \frac{g^4}{2}}{\displaystyle _{𝐤,𝐩}}G(𝐪𝐤)G(𝐤)^2G(𝐩)G(𝐤𝐩)`$ (135)
$`=`$ $`{\displaystyle \frac{g^4}{2}}{\displaystyle _{𝐤,𝐩}}{\displaystyle \frac{1}{[\overline{\tau }+(\overline{𝐪}𝐤)^2][\overline{\tau }+𝐤^2]^2[\overline{\tau }+(𝐩+i𝐯)^2][\overline{\tau }+(𝐤𝐩)^2]}},`$ (136)
where we have shifted $`𝐤𝐤i𝐯`$, and defined $`\overline{\tau }=\tau +𝐯^2`$ and $`\overline{𝐪}=𝐪+2i𝐯`$. We expand the last integral in $`𝐪`$ and $`𝐯`$ to second order and find
$`\mathrm{\Sigma }^{(2a)}(𝐪)`$ $`=`$ $`{\displaystyle \frac{g^4}{2}}(I_{113}+({\displaystyle \frac{d4}{d}}I_{123}+{\displaystyle \frac{4\overline{\tau }}{d}}I_{133})𝐯^2`$ (138)
$`+2(I_2I_4+\overline{\tau }I_{124}I_{114}I_{123}){\displaystyle \frac{i𝐯\overline{𝐪}}{d}}({\displaystyle \frac{d4}{d}}I_{114}+{\displaystyle \frac{4\overline{\tau }}{d}}I_{115})\overline{𝐪}^2).`$
In the same fashion we get for the diagram shown in Fig. 2(b)
$`\mathrm{\Sigma }^{(2b)}(𝐪)`$ $`=`$ $`g^4{\displaystyle _{𝐤,𝐩}}G(𝐪𝐤)G(𝐪𝐩)G(𝐤)G(𝐩)G(𝐤𝐩)`$ (139)
$`=`$ $`g^4\tau ^\epsilon (I_{122}+({\displaystyle \frac{d4}{d}}I_{222}+{\displaystyle \frac{4\overline{\tau }}{d}}I_{223})𝐯^2`$ (141)
$`2({\displaystyle \frac{d6}{d}}I_{123}+{\displaystyle \frac{4\overline{\tau }}{d}}I_{124}+{\displaystyle \frac{1}{d}}I_3^2+{\displaystyle \frac{\overline{\tau }}{d}}I_{133})\overline{𝐪}^2).`$
Using Eqs. (133,134) we obtain from Eqs. (138,141) the singular contributions to the two-loop vertex function after renormalization according to the scheme (8,9,13) as
$`\mathrm{\Gamma }_{1,1}^{(2loop)}`$ $`=`$ $`{\displaystyle \frac{u^2}{\epsilon ^2}}\tau ^\epsilon [({\displaystyle \frac{9}{4}}+{\displaystyle \frac{45\epsilon }{16}})\mu ^2\tau +({\displaystyle \frac{5}{2}}+{\displaystyle \frac{25\epsilon }{24}}){\displaystyle \frac{(\mu 𝐯)^2}{3}}`$ (143)
$`+({\displaystyle \frac{11}{6}}+{\displaystyle \frac{7\epsilon }{72}}){\displaystyle \frac{𝐪^2}{6}}+({\displaystyle \frac{23}{12}}+{\displaystyle \frac{13\epsilon }{144}}){\displaystyle \frac{2i\mu 𝐯𝐪}{3}}].`$
Using the renormalization constants to $`O(u)`$, Eqs. (10-12,25,26), we get from Eq. (23), after the renormalization $`\stackrel{˚}{\mathrm{\Gamma }}_{1,1}\mathrm{\Gamma }_{1,1}=Z\stackrel{˚}{\mathrm{\Gamma }}_{1,1}`$, the one-loop self-energy to $`O(u^2)`$ as
$`\mathrm{\Gamma }_{1,1}^{(1loop)}`$ $`=`$ $`{\displaystyle \frac{u}{\epsilon }}\tau ^{\epsilon /2}[(1+({\displaystyle \frac{9}{2}}+{\displaystyle \frac{11\epsilon }{6}}){\displaystyle \frac{u}{\epsilon }})\mu ^2\tau +(1+(5{\displaystyle \frac{5\epsilon }{12}}){\displaystyle \frac{u}{\epsilon }}){\displaystyle \frac{(\mu 𝐯)^2}{3}}`$ (145)
$`+(1+({\displaystyle \frac{11}{3}}{\displaystyle \frac{5\epsilon }{12}}){\displaystyle \frac{u}{\epsilon }}){\displaystyle \frac{𝐪^2}{6}}+(1+({\displaystyle \frac{23}{6}}{\displaystyle \frac{5\epsilon }{12}}){\displaystyle \frac{u}{\epsilon }}){\displaystyle \frac{2i\mu 𝐯𝐪}{3}}].`$
By adding $`\mathrm{\Gamma }_{1,1}^{(1loop)}`$, $`\mathrm{\Gamma }_{1,1}^{(2loop)}`$ and the renormalized zero-loop part, we get
$$\mathrm{\Gamma }_{1,1}=\mu ^2\left(Z_\tau \tau +Y_{v\tau }v^2\right)+2iZ_v\mu 𝐯𝐪+Zq^2+\mathrm{\Gamma }_{1,1}^{(1loop)}+\mathrm{\Gamma }_{1,1}^{(2loop)}+O(\epsilon ^0,u^3).$$
(146)
We see that the nonprimitive divergencies $`\mathrm{ln}\tau `$ of $`\mathrm{\Gamma }_{1,1}^{(1loop)}`$ and $`\mathrm{\Gamma }_{1,1}^{(2loop)}`$ cancel (as a check of a correct calculation) and find finally the renormalization constants cited in Eqs. (10,11,25,26).
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# Possible evidence of quark matter in neutron star X-ray binaries
## 1 Role of quark phase transition
The density and pressure in the interior of neutron stars is high in comparison with nuclear density by a factor of some 5 to 10 depending on the particular models used to estimate it. At such densities it is quite plausible that the quark constituents of hadrons loose their association with particular hadrons—the deconfined quark matter phase replaces the normal phase. In a nonrotating star, the radial boundaries between quark core, mixed phase, and normal hadronic phase would remain fixed. However in a rotating star, because of the centrifugal distortion of the density in the interior, these boundaries will change as the rotational frequency of the star changes with time.
A structural change occurs in such a star with changes in frequency (Glendenning, Pei & Weber (1997)). If there were no change in the nature of matter, the stellar fluid would respond simply under the action of the centrifugal force. However, the compressibility of the normal nuclear matter phase and the deconfined and relatively free Fermi gas of the quark matter phase, are different. The former must be less compressible than the latter. When a ms pulsar spins down, its central density may rise above the critical phase transition density and the central core will then change phase to softer quark matter; it is compressed both by its own gravitational attraction, and by the weight of the overlying nuclear phase. The reverse will be true in spinup due to accretion. In either case, the distribution of mass, radius and moment of inertia are changed by a phase transition beyond those changes that would take place in an immutable fluid under the action of a changing centrifugal force.
These ideas were applied to the spin-down of a ms pulsar (Glendenning, Pei & Weber (1997); Glendenning (1998); Weber (1999)). It was found that as the quark matter core grew in radial extent, the moment of inertia decreased anomalously, and could even introduce an era of spin-up lasting for $`10^7`$ years (Glendenning (1998)). The response of the moment of inertia to changes in spin is very like the so-called “backbending” in nuclei predicted by Mottelson and Valatin (Mottelson & Valatin (1960)) and discovered many years ago (Johnson, Ryde & Hjorth (1972); Stephens & Simon (1972)).
Accreting X-ray neutron stars provide a very interesting contrast to the spin-down of isolated ms pulsars. The X-ray stars are being spun up by the accretion of matter from a low-mass, less-dense white dwarf companion. They are presumably the link between the canonical pulsars with mean period of $`0.7`$ sec and the ms pulsars (Wijnands & van der Klis (1998); Chakrabarty & Morgan (1998)).
If the critical deconfinement density falls within the range spanned by canonical pulsars, quark matter will already exist in them but may be “spun” out of X-ray stars as their frequency increases during accretion. We can anticipate that in a certain frequency range, the changing radial extent of the quark matter phase will actually inhibit changes in frequency because of the increase in moment of inertia occasioned by the gradual disappearance of the quark matter phase. Accreters will tend to spend a greater length of time in the critical frequencies than otherwise. There will be an anomalous number of accreters that appear at or near the same frequency. This is what was found recently with the Rossi X-ray Timing Explorer (RXTE) (van der Klis (2000)).
Presumably, accreters commence their evolution near the death line of active canonical pulsars with frequencies of $`\nu 1`$ Hz and end as ms pulsars with $`\nu 200\mathrm{to}600`$ Hz. The spinup evolution of an accreting star is a more complicated problem than that of the spindown of an isolated ms pulsar of constant baryon number. It is complicated by the accretion of matter ($`\dot{M}>10^{10}M_{}\mathrm{yr}^1`$), a changing magnetic field strength (from $`B10^{12}\mathrm{to}10^8`$ G), and the interaction of the field with the accretion disk.
## 2 Spin-evolution of accreting X-ray neutron stars
The change in moment of inertia as a function of rotational frequency caused by accretion of matter is similar to that described by Glendenning, Pei & Weber (1997) for spindown of a ms pulsar because of magnetic dipole radiation. However, there are additional phenomena as just mentioned. Generally, a canonical pulsar will have evolved from birth with moderate rotational frequency to the deathline. At that time, the usual drag of the dipole radiation will be eclipsed by the accretion phenomena. The spin-up torque of the accreting matter causes a change in the star’s angular momentum according to the relation (Elsner & Lamb (1977); Ghosh, Lamb & Pethick (1977); Lipunov (1992))
$`{\displaystyle \frac{dJ}{dt}}=\dot{M}\stackrel{~}{l}(r_\mathrm{m})N(r_\mathrm{c}).`$ (1)
This can be rewritten as a time evolution equation for the angular velocity $`\mathrm{\Omega }`$ of the accreting star $`(J/I)`$,
$$I(t)\frac{d\mathrm{\Omega }(t)}{dt}=\dot{M}\stackrel{~}{l}(t)\mathrm{\Omega }(t)\frac{dI(t)}{dt}\kappa \mu (t)^2r_\mathrm{c}(t)^3,$$
(2)
with the following definitions: The accretion rate is denoted by $`\dot{M}`$ ($`G=c=1`$) and
$`\stackrel{~}{l}(r_\mathrm{m})=\sqrt{Mr_\mathrm{m}}`$ (3)
is the angular momentum added to the star per unit mass of accreted matter. The quantity $`N`$ stands for the magnetic plus viscous torque term ($`\kappa 0.1`$),
$$N(r_\mathrm{c})=\kappa \mu ^2r_\mathrm{c}^3,$$
(4)
with $`\mu R^3B`$ the star’s magnetic moment. The quantities $`r_\mathrm{m}`$ and $`r_\mathrm{c}`$ denote the radius of the inner edge of the accretion disk and the co-rotating radius, respectively, and are given by $`(\xi 1)`$
$$r_\mathrm{m}=\xi r_\mathrm{A},$$
(5)
and
$`r_\mathrm{c}=\left(M\mathrm{\Omega }^2\right)^{1/3}.`$ (6)
Accretion will be inhibited by a centrifugal barrier if the neutron star’s magnetosphere rotates faster than the Kepler frequency at the magnetosphere. Hence $`r_\mathrm{m}<r_\mathrm{c}`$, otherwise accretion onto the star will cease. The Alfén radius $`r_\mathrm{A}`$, where the magnetic energy density equals the total kinetic energy of the accreting matter, in Eq. (5) is defined by
$$r_\mathrm{A}=\left(\frac{\mu ^4}{2M\dot{M}^2}\right)^{1/7}.$$
(7)
We assume that the magnetic field evolves according to
$$B(t)=B(\mathrm{})+(B(0)B(\mathrm{}))e^{t/t_\mathrm{d}}$$
(8)
with $`t=0`$ at the start of accretion, and where $`B(0)=10^{12}`$ G and $`t_\mathrm{d}=10^6`$ yr. Such a decay to an asymptotic value seems to be a feature of some treatments of the magnetic field evolution (Konar & Bhattacharya (1999)).
It has previously been assumed that the moment of inertia in Eq. 2 does not respond to changes in the centrifugal force, and in that case, the above formula yields a well-known estimate of the period to which a star can be spun up (Bhattacharya & van den Heuvel (1991)). The approximation is true for slow rotation. However, the response of the star to rotation becomes important as the star is spun up by accretion. Not only do changes in the distribution of matter occur but internal changes in composition occur also because of changes induced in the central density by centrifugal dilution (Glendenning, Pei & Weber (1997)); both changes effect the moment of inertia and hence the response of the star to accretion. In this Letter we wish to follow the spin-evolution of the star, and so, must take such refinements into account.
The moment of inertia of ms pulsars or of neutron star accreters has to be computed in GR without making the usual assumption of slow rotation (Hartle (1967); Hartle & Thorne (1968)). Fortunately, we have previously obtained an expression for the moment of inertia of a rotating star (Glendenning & Weber (1992)). The expression is too cumbersome to reproduce here. Stars that are spun up to high frequencies close to the breakup limit (Kepler frequency) undergo dramatic interior changes; the central density may change by a factor of four or so over that of a slowly rotating star if a phase change occurs during spin-up (Glendenning (1998); Weber (1999)).
Figure 1 shows how the moment of inertia changes for neutron stars in binary systems that are spun up by mass accretion according to Eq. (2) until $`0.4M_{}`$ has been accreted. The neutron star models are fully described in (Glendenning, Pei & Weber (1997)) and references therein, and the initial mass of the star in our examples is $`1.42M_{}`$. Confined nuclear matter is described by a covariant Lagrangian describing the interaction of the members of the baryon octet with scalar, vector and vector-isovector mesons and solved in the meanfield approximation. Quark matter is described by the MIT bag model. In one case, it is assumed that a phase transition to quark matter occurs, and in the other that it does not. This accounts for the different initial moments of inertia, and also, as we see, the response to spinup. Three accretion rates are assumed, which range from $`\dot{M}_{10}=1`$ to 100. These rates are in accord with observations made on low-mass X-ray binaries (LMXBs) observed with the Rossi X-ray Timing Explorer (van der Klis (2000)). The observed objects, which are divided into Z sources and A(toll) sources, appear to accrete at rates of $`\dot{M}_{10}200`$ and $`\dot{M}_{10}2`$, respectively.
Figure 2 shows the spin evolution of accreting neutron stars as determined by the changing moment of inertia and the spin evolution equation (2). Neutron stars without quark matter in their centers are spun up along the dashed lines to equilibrium frequencies between about 600 Hz and 850 Hz, depending on accretion rate and magnetic field. The $`dI/dt`$ term for these sequences manifests itself only insofar as it limits the equilibrium periods to values smaller than the Kepler frequency, $`\nu _\mathrm{K}`$. In both Figs. 1 and 2 we assume that $`0.4M_{}`$ is accreted. Otherwise the maximum frequency attained is less.
The spin-up scenario is dramatically different for neutron stars in which quark deconfinement occurs. In this case, as known from Fig. 1, the temporal conversion of quark matter into its mixed phase of quarks and confined hadrons is accompanied by a pronounced increase of the stellar moment of inertia. This increase contributes so significantly to the torque term $`N(r_\mathrm{c})`$ in Eq. (2) that the spin-up rate $`d\mathrm{\Omega }/dt`$ is driven to saturation around those frequencies at which the pure quark matter core in the center of the neutron star gives way to the mixed phase of confined hadronic matter and quark matter. The star resumes ordinary spin-up if this transition is completed. The epoch during which the spin rates are saturated are determined by attributes like the accretion rate, magnetic field, and its assumed decay time. The epoch lasts between $`10^7\mathrm{and}10^9`$ yr depending on the accretion rate at the values taken for the other factors.
We can translate the information in Fig. 2 into a frequency distribution of X-ray stars by assuming that neutron stars begin their accretion evolution at the average rate of one per million years. A different rate will only shift some neutron stars from one bin to an adjacent one, but will not change the basic form of the distribution. The result is shown in Fig. 3. The result is striking. Spinout of the quark matter core as the neutron star spins up is signalled by a spike in the frequency distribution. The concentration in frequency is centered around 200 Hz, about 100 Hz lower than the observed spinup anomaly. This discrepancy is not surprising given the crude representation of confinement by the MIT bag model while the physics underlying the effect of a phase transition on spin rate is robust.
We address now the bump in the histogram at high frequencies. Certainly there are high frequency pulsars. However, if a histogram of ms pulsar frequencies is made from Ref. (Princeton data base ), a spike is found near 300 Hz, and a strong attenuation in number of high frequency pulsars above the spike. So both the (sparse) data on X-ray objects and on ms pulsars agree on a spike and on attenuation at high frequency. Why the high frequency bump (containing about 9 X-ray objects) in our Fig. 3? The actual white dwarf masses in these low-mass binaries is believed to be $`0.1`$ to $`0.4M_{}`$. We computed the frequency distribution for donor masses in this range in steps of $`0.1M_{}`$ and until all mass has been transferred at the chosen rate. (The result has little sensitivity to the accretion rate.) Since we do not know the mass distribution of donors, we averaged the results. The highest frequencies are attributable to the $`0.4M_{}`$ mass donors. These presumably account for the high frequency tail of the ms pulsar distribution. Of course it is only an assumption that the mass distribution of donors is flat and that the donor is entirely consumed. Each binary represents a unique combination of neutron star and companion masses, magnetic fields, and accretion rates, but with unknown weight for these differences. We simply do not know how the high frequency end is attenuated, but it surely is, as observation tells us, though not so severely as the few data on LMXBs would suggest, inasmuch as they are believed to be the pathway to ms pulsars, several of which have frequencies as high as $`650`$ Hz.
## 3 Summary
We have traced the time evolution of the moment of inertia and rotational frequency for a neutron star accreting matter from a low-mass companion, under various assumptions about the accretion rate and for two stellar models, one an ordinary neutron star populated by nucleons, hyperons and leptons, and one in which phase equilibrium between ordinary and quark deconfined matter occurs within the density range found in canonical pulsars. In the second case the computed frequency distribution of X-ray neutron stars shows a spike, much as is observed in a recent compilation of data (van der Klis (2000)). There are various suggestions as to the cause of the spike, several of which we cite (c.f. Bildsten (1998); Andersson et al (2000); Levin (1999)). A possible contributing mechanism which causes some accreters of suitable mass to resist spinup for a lengthy era is that discussed in this paper— the ongoing reduction of quark matter cores in the centers of neutron stars as they are spun up. This occurs because, with increasing spin, the density of the inner region is centrifugally diluted until it falls below the threshold density at which quark matter can exist, first in the center, and then in an expanding region. As explained in the introduction, the conversion of quark mattter to confined hadronic matter manifests itself in an expansion of the star and a significant increase in its moment of inertia. As a consequence, the angular momentum added to a neutron star during this phase of evolution is then consumed by the star’s expansion, inhibiting a further spin-up until the quark matter has been converted into a mixed phase of matter made up of hadrons and quarks.
###### Acknowledgements.
This work was supported by the Deutsche Forschungsgemeinschaft (DFG), and by the Director, Office of Science, Office of High Energy and Nuclear Physics, Division of Nuclear Physics, of the U.S. Department of Energy under Contract DE-AC03-76SF00098.
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# On the Relationship of Gravitational Constants in KK Reduction
## Abstract
In this short note, we try to clarify a seemly trivial but often confusing question in relating a higher-dimensional physical gravitational constant to its lower-dimensional correspondence in Kaluza-Klein reduction. In particular, we re-derive the low-energy M-theory gravitational constant in terms of type IIA string coupling $`g_s`$ and constant $`\alpha ^{}`$ through the metric relation between the two theories.
preprint: UMTH-00-06 hep-th/0003132
The proper determination of eleven dimensional M-theory gravitational constant (therefore, the eleven-dimensional Planck constant), in terms of type IIA string coupling $`g_s`$ and constant $`\alpha ^{}`$, is important, for example, for the BFSS matrix proposal of M-theory. It is also important for whether brane modes can possibly decouple from bulk gravity modes in the so-called decoupling limit. Given the string constant $`\alpha ^{}`$ ( therefore the units in type IIA string theory) and the relationship between 11-D M-theory and type IIA string theory, the 11-D M-theory physical gravitational constant as well as the units for M-theory are also given. We therefore expect a precise expression for the M-theory gravitational constant in terms of type IIA string coupling $`g_s`$ and constant $`\alpha ^{}`$. However, there exist no unique answers in the literature for this constant. We try to clarify, in this short note, possible confusion about the derivations of this constant.
Let us begin with a general discussion in relating a higher-dimensional physical gravitational constant to its lower-dimensional correspondence in dimension redection. We start with the following gravity action in dimension D
$$I_D=\frac{1}{2\overline{\kappa }_D^2}d^Dx\sqrt{det\widehat{G}}\left[\widehat{R}+\mathrm{}\right],$$
(1)
where $`\widehat{G}_{MN}`$ is the metric, $`\widehat{R}`$ is the Ricci scalar, the constant $`\overline{\kappa }_D^2`$ is usually called gravitational constant (the Newton constant $`8\pi G_N^D\kappa _D^2`$)<sup>*</sup><sup>*</sup>*The $`\kappa `$ without a bar corresponds to the physical gravitional constant while the one with a bar is not necesarily a physical one, see the explanation given in the text., and $`\mathrm{}`$ in the above action represents other possible fieldsFor the purpose of this paper, we need to consider only the first term in the action..
Now we wish to compactify the above action to dimension d ($`<D`$). For our purpose, we need to consider only the massless graviton whose effective action is
$$I_d=\frac{1}{2\overline{\kappa }_d^2}d^dx\sqrt{detg}\left[R+\mathrm{}\right].$$
(2)
In obtaining the above action from Eq. (1), we made the split for the higher-dimensional coordinate $`x^M=(x^\mu ,y^i)`$ with $`M=0,1\mathrm{},D1;\mu =0,1,\mathrm{}d1`$ and $`i=1,\mathrm{}Dd`$. Here $`x^\mu `$ are the coordinates of the lower-dimensional spacetime and $`y^i`$ are the compactified coordinates. We therefore have the same units for both the D-dimensional theory and the compactified d-dimensional theory. The massless sector of the lower-dimensional theory can be obtained by assuming the higher-dimensional fields to be independent of $`y^i`$. We then simply integrate out the $`y^i`$ from the action (1). By comparing the resulted action with Eq. (2), we have the relation for the gravitational constants and the compactified volume $`V_{Dd}`$ asI will make this precise when we discuss on how to reduce the low energy M-theory to the low energy type IIA string theory.
$$\kappa _D^2=\kappa _d^2V_{Dd}.$$
(3)
We often say that the above equation, relating the higher-dimensional physical gravitational constant to its lower-dimensional correspondence through the physical volume measured with respect to the lower-dimensional metric, is independent of the actual metric relation between the higher-dimensional and the lower-dimensional theories. This is true indeed.
However, what is often ignored in practice is the implicit assumption used in deriving Eq. (3) that we choose the asymptotic metric (i.e., the underlying vacuum) for the higher-dimensional theory to be same as that for the compactified lower-dimensional theory<sup>§</sup><sup>§</sup>§Usually we choose the asymptotic metric for the lower-dimensional theory to be Minkowskian, i.e., flat Minkowski metric $`\eta _{\mu \nu }=(,+,\mathrm{},+)`$, see the definition given in . This is the metric used in defining the physical gravitational constant. This is also the metric used in perturbative string theory in defining the string tension $`T_f=1/(2\pi \alpha ^{})`$ or the string constant $`\alpha ^{}`$.. Only in this case, we can take the constant $`\overline{\kappa }^2`$ in front of the respective action as the physical gravitational constantIn this case, for example, the metric $`g_{\mu \nu }=\eta _{\mu \nu }+\kappa h_{\mu \nu }`$ for small graviton fluctuation around flat Minkowski spacetime and the Einstein-Hilbert action reduces to the cannonical form $`(h)^2+\kappa (h)^2h`$.. In general, however, the asymptotic metrics for both the higher-dimensional theory and the compactified lower-dimensional theory are not necessarily the same because the scalars due to the compactification develop VEV Of course, one can always define both of the higher-dimensional metric and the lower-dimensional one to be same asymptotically by absorbing the possible constant factor due to the VEV of scalars into the $`\overline{\kappa }^2`$ in front of the action. For the higher-dimensional theory, the compactified coordinates $`y^i`$ have to be rescaled properly with respect to the lower-dimensional metric, see the example given later in relating M-theory to type IIA string. Then the resulting $`\kappa ^2`$ is the physical one. This is what Maldacena did in for obtaining masses properly for BPS states in string theory using U-duality.. If this happens, we cannot take both $`\overline{\kappa }_D^2`$ and $`\overline{\kappa }_d^2`$ in front of the respective action in the above as physical<sup>\**</sup><sup>\**</sup>\**Polchinski in chose both the units and metric for M-theory the same as those for string theory. By definition, his $`\kappa _{11}^2`$ and the compactified radius are both physical, not the usual $`\overline{\kappa }_{11}^2`$ and the coordiante radius $`r`$.. The ignorance of this fact is often the source of confusion in the literature. For example, the low-energy M-theory physical gravitational constant has been given correctly in as $`2\kappa _{11}^2=(2\pi )^8g_s^3\alpha ^{9/2}`$ in terms of type IIA string coupling $`g_s`$ and constant $`\alpha ^{}`$. But this constant has also been given incorrectly in the literature precisely because $`\overline{\kappa }`$ is mistaken as $`\kappa `$.
In the remainder of this note, I will focus, as an example, on the reduction of 11-dimensional low-energy M-theory on a circle $`S^1`$ to give the low energy theory of type IIA string. The 11-D M-theory metric is related to the type IIA string metric as
$$ds_{11}^2=e^{2\varphi /3}ds_{10}^2+e^{4\varphi /3}(dx^{11})^2,$$
(4)
where $`\varphi `$ is the dilaton in type IIA string theory, 11-th coordinate $`x^{11}`$ has a period $`2\pi r`$ with the coordinate radius $`r`$. In the above we have dropped the KK vector field $`A_\mu `$. For our purpose, $`A_\mu `$ is irrelevant. We take type IIA string metric $`ds_{10}^2`$ as asympotically Minkowski since that is where we quantize the perturbative type IIA string. That is the metric used in defining the fundamental string tension $`T_f=1/(2\pi \alpha ^{})`$. So are the tensions for D-branes and NSNS branes. With Eq. (4) and taking $`D=11`$ in Eq. (1), we have the low energy action of type IIA string as
$$I_{10}=\frac{2\pi r}{\overline{\kappa }_{11}^2}d^{10}x\sqrt{detg}e^{2\varphi }\left[R+\mathrm{}\right],$$
(5)
where $`g_{\mu \nu }`$ is the string metric. By definition, we have the following relation
$$\overline{\kappa }_{11}^2=2\pi re^{2\varphi _0}\kappa _{10}^2,$$
(6)
where $`\kappa _{10}^2\overline{\kappa }_{10}^2e^{2\varphi _0}`$ is the physical gravitational constant in D = 10. $`\varphi _0`$ is the VEV of dilaton or the asymptotic value of the dilaton and is related to the string coupling as $`g_s=e^{\varphi _0}`$.
As we stress above that Eq. (3) holds true always. In the present context, it is
$$\kappa _{11}^2=2\pi \rho \kappa _{10}^2,$$
(7)
with $`\rho `$ the physical radius. Let me explain why $`2\kappa _{11}^2=(2\pi )^8g_s^3\alpha ^{9/2}`$ given in must be correct. As I mentioned above, Eq. (7) should hold always true. The physical gravitational constant $`2\kappa _{10}^2=(2\pi )^7g_s^2\alpha ^4`$ was given in . As we now know that the strong coupling of type IIA string is just M-theory compactified on a big circle. In order for this to be true, one needs to identify the spectrum of D0 branes with that of momentum (Kaluza-Klein) states. This implies that the physical radius of the circle measured in string metric is given as the inverse of mass of a single D0 brane, i.e., $`\rho =g_s\alpha ^{1/2}`$. Then we have the 11-D physical gravitational constant from Eq. (7) as given above.
From Eqs. (6) and (7), we have
$$\frac{\overline{\kappa }_{11}^2}{re^{2\varphi _0}}=\frac{\kappa _{11}^2}{\rho }.$$
(8)
We intend to determine the relation between $`\kappa _{11}`$ and $`\overline{\kappa }_{11}`$ and the relation between $`r`$ and $`\rho `$ unambiguously. To my knowledge, no explicit derivations for these two relations have been given in the literature. We dare to present them here.
For our purpose, we need to consider only the asymptotic metric relation in Eq. (4), i.e.,
$`(ds_0)_{11}^2`$ $`=`$ $`e^{2\varphi _0/3}(ds_0)_{10}^2+e^{4\varphi _0/3}(dx^{11})^2,`$ (9)
$`=`$ $`e^{2\varphi _0/3}\left[(ds_0)_{10}^2+(d\stackrel{~}{x}^{11})^2\right],`$ (10)
where the asymptotic string metric $`(ds_0)_{10}^2`$ is actually Minkowskian and the rescaled 11-th coordinate $`\stackrel{~}{x}^{11}=e^{\varphi _0}x^{11}`$ with its radius $`\stackrel{~}{r}=e^{\varphi _0}r`$. The first line in the above equation indicates clearly that the 11-D metric cannot be asymptotically Minkowskian if we insist $`ds_{10}^2`$ be so<sup>††</sup><sup>††</sup>††We can no longer rescale $`(ds_0)_{10}^2`$ since that is the metric used in defining the string constant $`\alpha ^{}`$.. The second line says that the 11-D metric can be made asymptotically Minkowskian up to a constant scaling factor $`e^{2\varphi _0/3}`$ if we rescale $`x^{11}`$ to $`\stackrel{~}{x}^{11}`$ as given above. In other words, the scaled radius $`\stackrel{~}{r}`$ is measured with respect to the string metric. Because the string constant $`\alpha ^{}`$ is defined with respect to the string metric, the second line in Eq. (10) should be used in the following equation.
By definition, from the above asymptotic metric relation, we have
$$\kappa _{11}^2=\frac{\overline{\kappa }_{11}^2}{\sqrt{det\widehat{G}_0}\widehat{G}_{0}^{}{}_{}{}^{1}}=e^{3\varphi _0}\overline{\kappa }_{11}^2,$$
(11)
where $`\widehat{G}_0`$ denotes the 11-D asymptotic metric given in the second line of Eq. (10). Using the above and Eq. (8), we derive $`r=\alpha ^{1/2}`$. Then $`\stackrel{~}{r}=e^{\varphi _0}\alpha ^{1/2}=\rho `$ is the physical radius measured in the string metric.
Let us provide an independent check of Eq. (11). For simplicity, we consider the reduction of a scalar field $`\mathrm{\Phi }(x^M)`$ from 11-D to 10-D on a circle $`S^1`$. The usual KK reduction says
$$\mathrm{\Phi }(x^\mu ,x^{11})=\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\mathrm{\Phi }_n(x^\mu )e^{inx^{11}/r},$$
(12)
or
$$\mathrm{\Phi }(x^\mu ,\stackrel{~}{x}^{11})=\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\mathrm{\Phi }_n(x^\mu )e^{in\stackrel{~}{x}^{11}/\stackrel{~}{r}},$$
(13)
where $`x^{11}`$ or $`\stackrel{~}{x}^{11}`$ is the compactified coordinate defined earlier. The 11-D wave equation $`_M^M\mathrm{\Phi }=0`$ in the asymptotic region (or around the Minkowski vacuum) becomes
$$\eta ^{\mu \nu }_\mu _\nu \mathrm{\Phi }_n(x^\mu )=\frac{n^2}{r^2e^{2\varphi _0}}\mathrm{\Phi }_n(x^\mu ),$$
(14)
where we have used the first line in Eq. (10), or we have
$$\eta ^{\mu \nu }_\mu _\nu \mathrm{\Phi }_n(x^\mu )=\frac{n^2}{\stackrel{~}{r}^2}\mathrm{\Phi }_n(x^\mu ),$$
(15)
where we have used the second line in Eq. (10). The mass spectrum with respect to the 10-D Minkowski vacuum in string frame can be obtained from either of the above equations as
$$M_n^2=p_\mu p^\mu =\frac{n^2}{r^2e^{2\varphi _0}}=\frac{n^2}{\stackrel{~}{r}^2}.$$
(16)
From the above, we should identify $`\stackrel{~}{r}=re^{\varphi _0}`$ as the physical radius with respect to the string metric. The mass $`M_n`$ should be identified with that of n D0 brane for the reason mentioned earlier. We therefore have $`\stackrel{~}{r}=g_s\alpha ^{1/2}`$. So we have $`r=\alpha ^{1/2}`$. Using Eq. (8) and $`\rho =\stackrel{~}{r}`$, we obtain also Eq. (11). Our discussion gives $`2\overline{\kappa }_{11}^2=(2\pi )^8\alpha ^{9/2}`$.
In summary, we explain how to obtain the physical gravitational constant for the original higher dimensional theory if we know the physical gravitational constant in the compactified lower-dimensional theory. In particular, we derive the relation $`\kappa _{11}^2=g_s^3\overline{\kappa }_{11}^2`$ and determine the radius $`r=\alpha ^{1/2}`$ (or $`\stackrel{~}{r}=g_s\alpha ^{1/2}`$) unambiguously.
###### Acknowledgements.
The author would like to thank Dan Chung for discussion and to acknowledge the support of U. S. Department of Energy.
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# Adaptive filtering techniques for gravitational wave interferometric data: Removing long-term sinusoidal disturbances and oscillatory transients.
## I Introduction
Over the next decade, several large-scale interferometric gravitational wave detectors will come on-line. These include LIGO, composed of two Laser Interferometer Gravitational-wave Observatories situated in the U.S. , VIRGO, a French/Italian project located near Pisa , GEO600, a German/British interferometer under construction near Hannover , TAMA in Japan, a medium-scale laser interferometer , and with funding approval AIGO500, the proposed 500 meter project sponsored by ACIGA. There are also separate proposals for space-based detectors which could be operational twenty-five years from now (e.g., LISA: the Laser Interferometer Space Antenna, a cornerstone project of the European Space Agency ). In the meantime, a number of existing resonant bar detectors will have had their sensitivities further enhanced.
The key to gravitational wave detection is the very precise measurement of small changes in distance. For laser interferometers, this is the distance between pairs of mirrors hanging at either end of two long, mutually perpendicular vacuum chambers. Gravitational waves passing through the instrument will shorten one arm while lengthening the other. By using an interferometer design, the relative change in length of the two arms can be measured, thus signaling the passage of a gravitational wave at the detector site. Long arm lengths, high laser power, and extremely well-controlled laser stability are essential to reach the requisite sensitivity, since the gravitational waves will be faint and will modify only weakly the structure of space-time in the detector’s arms (see e.g., ).
Gravitational wave detectors produce an enormous volume of output (e.g., of the order of 16 MB/sec for the LIGO instruments) consisting mainly of noise from a host of sources both environmental and intrinsic to the apparatus. Buried in this noise will be the gravitational wave signature. Sophisticated data analysis techniques will have to be developed to optimally extract astrophysical data. Many of the techniques developed so far are based on matched filtering and assume stationary Gaussian noise.
However, the real data stream from the detectors is not expected to satisfy the stationary and Gaussian assumptions. In fact, the data from the Caltech 40 meter proto-type interferometer has the expected broadband noise spectrum, but superposed on this are several other noise features ; such as long-term sinusoidal disturbances emanating from suspensions and electric main harmonics and also transients occurring occasionally, typically due to servo-controls instabilities or mechanical relaxation in suspension system etc. While no precise a priori model can be given for this noise until the detector is completed and fully tested, matched filtering techniques cannot be used to locate/remove these noisy signals.
This disparity between standard Gaussian assumptions and real data characteristics poses a major problem to the direct application of matched filtering techniques. This is true when searching for burst sources such as blackhole binary quasinormal ringings . This is also the case for the inspiral searches in Caltech 40meter data, where one has to introduce a veto on the decision taken with the matched filter to ensure that the detected signal is actually the one we are looking for.
It is possible that in the future, improved experimental techniques and greater experience, will reduce or even completely eliminate some of these nonstationary and non-Gaussian features. Nevertheless, it will take probably some time to reach such acceptable and high quality of data. Therefore, it is necessary and desirable to somehow combat this noise. Since such noise features defy modeling, a novel approach to the problem is called for.
We propose a denoising method based on LMS adaptive linear prediction techniques which does not require any precise a priori information about the noise characteristics. Although our method does not pretend to optimality, we believe that its simplicity makes it useful for data preparation and for the understanding of the first data.
In the following, we present the principles of LMS adaptive denoising (Sect. II), a characterization of its behavior on a simple model of the noise from the interferometer (Sect. III), the precise structure of the denoising algorithm (Sect. IV) and results (Sect. V) obtained with simulated data and also with real data taken from the Caltech 40 meter proto-type interferometer .
This work here is preliminary; its goal is to explore how effectively adaptive filtering techniques perform on the problem we address. It is a first step towards a more complete statistical evaluation of the algorithm.
## II Methods
### A From hypothesis to method
We assume that the noise consists of broadband Gaussian noise plus large amplitude oscillating interference signals. The model does not include any a priori knowledge of the signal such as its exact frequency or shape of the envelope. The only assumption we make is that its autocorrelation over a small time-lag – the time-lag chosen greater than the decorrelation time scale of the broadband noise – is appreciable, while for the broadband noise it is essentially zero. This difference can be used to advantage to discriminate between the narrow band interferences and the broadband noise.
The idea is to predict the current signal sample given the previous samples of the data. This is possible, only if the target sample shares enough information with (i.e., is sufficiently correlated to) the previous samples. In other words, the only predictable part of the signal is the one whose correlation length is sufficiently large (i.e., long-term sinusoids or ringdowns). Conversely, broadband noise cannot be predicted, as it is not possible to guess the next value in this way. It is this crucial underlying idea we use to discriminate between the two noise signals.
### B Mean square linear prediction
Let us recall some standard principles to design an optimal linear predictor. The question to address is to optimally predict the data sample $`x_k`$ with a collection of past samples $`𝒙_k:=(x_{kdm},m=0,1,\mathrm{},N1)^t`$ given the delay $`d1`$ (the quantity $`d`$ is also referred to as prediction depth). The prediction is obtained by linearly combining these data samples weighted by the $`N`$ corresponding coefficients<sup>*</sup><sup>*</sup>*We put brackets around indices of vectors and matrices in order to distinguish them from the time index. $`w^{(m)}`$, forming the tap-weight vector $`𝒘:=(w^{(m)},m=0,1,\mathrm{},N1)^t`$, where the superscript ‘t’ denotes the transpose of the vector. Therefore, the prediction $`y_k`$ of $`x_k`$ reads,
$$y_k:=𝒘^t𝒙_k.$$
(1)
The predictor is optimal in the mean square sense when the variance of the prediction error $`e_k=x_ky_k`$ is minimum. Therefore, the problem is to find the set of weight coefficients which minimizes
$$J_k(𝒘):=\text{E}[e_k^2]=\text{E}[(x_k𝒘^t𝒙_k)^2],$$
(2)
where E denotes the expectation value operator.
This leads to the minimization of the following quadratic form
$$J_k(𝒘)=\sigma _k^22𝒘^t𝒑_k+𝒘^t𝑹_k𝒘,$$
(3)
where $`\sigma _k^2:=\text{E}[x_k^2]`$, $`𝒑_k:=\text{E}[x_k𝒙_k]`$ and $`𝑹_k:=\text{E}[𝒙_k^t𝒙_k]`$. There exists only one solution $`𝒘_k^{}`$, obtained when the gradient of $`J_k`$ vanishes. This situation is realized when
$$𝑹_k𝒘_k^{}=𝒑_k.$$
(4)
When the signal is stationary, $`𝑹_k=𝑹`$ and $`𝒑_k=𝒑`$ are constant (independent of $`k`$). In this case, $`𝑹`$ defines the autocorrelation matrix of the signal $`x_k`$ and the solution of (4) is referred to as the Wiener filter.
### C Linear prediction and LMS method
Eq. (4) requires the computationally expensive inversion of the matrix $`𝑹_k`$. An alternative and more efficient solution for finding the minimum of the (3) consists in starting from an arbitrary initial value $`𝒘_{k,0}`$, and iterate the tap-weight vector along the steepest descent direction,
$$𝒘_{k,n+1}=𝒘_{k,n}\mu _𝒘J_k(𝒘_{k,n}),$$
(5)
given by the gradient
$$_𝒘J_k(𝒘)=2(𝑹_k𝒘𝒑_k).$$
(6)
For a sufficiently small gain $`\mu `$, the weight vectors will eventually converge to the optimal predictor filter $`𝒘_k^{}`$. This procedure requires the second order statistics (namely $`𝑹_k`$ and $`𝒑_k`$) of the signal. In our case, this information is not available and one has therefore to estimate these quantities. Instead of estimating directly $`𝑹_k`$ and $`𝒑_k`$ and combining them with (6), a more efficient solution is to estimate the gradient. From the derivation of (2), one can rewrite the gradient as
$$_𝒘J_k(𝒘)=2\text{E}[e_k𝒙_k].$$
(7)
A simple and natural way to obtain an estimator of this quantity is to omit the expectation operator :
$$\widehat{_𝒘J_k}=2e_k𝒙_k.$$
(8)
Because the noise perturbs this estimate, the algorithm may iterate in a direction which does not lie along the direction of steepest descent, thus preventing the filter from converging to the Wiener filter. For this purpose, we stabilize the estimation above by setting the algorithm time index $`n`$ equal to signal time index $`k`$ in the Eq. (5). The final evolution equation for the tap-weight vector finally reads:
$$𝒘_{k+1}=𝒘_k+2\mu e_k𝒙_k.$$
(9)
At a fixed time $`k`$, the weight vector evolves along the crude estimate of the steepest descent direction. But on a longer duration, the direction followed by the tap-weight vector is governed by the sum of the successive gradient estimates obtained with different noise samples. In other words, we have replaced an ensemble average in (7) by a time average. It also implies that we have implicitly called for further assumptions on the signal $`x_k`$: first its local stationarity (more precisely, the second order statistics are supposed to be constant during the convergence time of the algorithm) and second, its ergodicity.
Summarizing, the method we propose consists in linearly filtering the data to extract the part of the signal with a long correlation time. As illustrated with the block diagram in Fig. 1, the finite impulse response filter (given by $`𝒘_k`$) is modified at each iteration according to the relation (9) with the final goal to minimize the mean square error. Once the filter has converged (i.e., $`𝒘_k`$ is stable in time), we reject the predicted part of the signal (corresponding to the long-term sinusoidal or the ringdown signals) and we send the rest of the signal for further analysis for detection.
### D Properties of the LMS method
The method we described above is referred to as adaptive line enhancer (ALE). It is a special case of the LMS algorithm. Both, ALE and LMS algorithms have been first introduced by Widrow and Hoff in the 1960’s.
The acronym LMS (Least Mean Square) designates a general scheme to design signal processing methods where a minimization (in a statistical sense) of a definite positive quadratic cost function (usually related to some mean quadratic error) is needed. Its central idea is the use of the estimate of the gradient of this function given in Eq. (8). The LMS technique has been extensively used for the last 30 years in communications problems such as echo cancellation, channel equalization, antenna processing, etc. The main advantages to be gained by applying the LMS technique are (i) adaptivity, (ii) robustness, (iii) simplicity.
In this context, the term “adaptivity” has two different meanings. First, it means that the LMS technique will automatically modify its parameters to reach for the best setup for a problem which has not been initially precisely defined. Second, it is also able to follow changes in the characteristics of the data being processed in the event that they occur. The latter property also shows that the method is robust. In fact, this method has been proved to be robust according to specific statistical criterion such as the minimax criterion .
The ALE is an adaptive prediction algorithm using the LMS technique. We have seen that the signal is predicted from a reference signal which is the signal itself. In some other applications, although the same principles are applied, the reference signal can be another signal, e.g. echo cancellation or denoising. In such cases, the quantity of interest might not be the prediction output but the linear filter used to compute it, e.g., deconvolution.
## III Adapting ALE filter to canceling noise in GW data
In this section we essentially describe a model for understanding the behaviour of the ALE algorithm. The model we assume consists of a high amplitude narrowband signal superposed on broadband noise. For simplicity, we assume the broadband noise to be white and Gaussian and the narrowband signals are sinusoids of constant envelope. The results we obtain hold for more realistic signals when the evolution of their amplitude and/or instantaneous frequency occurs adiabatically, i.e., the change is small over the period of the sinusoid.
The assumption of white noise is not too restrictive because this is equivalent to choosing the noise correlation time to be zero and therefore we are free to choose the prediction depth (i.e., the time delay between the current predicted data sample and the reference signal to the LMS filter) to be arbitrarily small. In a real situation, we must fix the delay to be greater than the correlation time of the broadband noise. We first analyze the case of the sinusoid because it is easier to investigate and provides invaluable insights into the workings of the LMS algorithm.
It may be remarked that the denoising of sinusoids in white noise has been treated in the literature with great detail (see for a review). We give here only pertinent results (with a short proof) for introducing the structure of the algorithm, which we present later in the text.
### A Optimal filter
We consider the data to be of the following form,
$$x_k:=\mathrm{cos}(2\pi f_0t_k+\mathrm{\Phi })+n_k,$$
(10)
where $`t_k:=k\delta `$, $`\delta :=1/f_s`$ being the sampling interval and $`\mathrm{\Phi }`$ is a random phase (at the origin) with uniform probability density function between $`\pi `$ and $`\pi `$. The sinusoid has frequency $`f_0`$ and the units are so chosen that it is of unit amplitude. The additive white noise $`n_k`$ with variance $`\sigma ^2`$ satisfies the relation,
$$\text{E}[n_kn_m]:=\sigma ^2\delta _{km},$$
(11)
where $`\delta _{km}`$ is the Kronecker delta.
The reference signal to the adaptive filter is just the delayed data by the amount $`d\delta `$, where $`d`$ is the number of time samples. We choose $`N`$ weights $`w_kW`$ ($`𝒘`$ can be thought of as a column vector) for the length of our filter, then the “reference vector” $`𝒙_k`$ at the $`k`$th time instant $`t_k`$ has the components $`x_{kdn}`$, $`n=0,1,\mathrm{},N1`$. The components of the autocorrelation matrix $`𝑹`$ and the vector $`𝒑`$ in Eq.(3) are given by
$`𝑹^{(mn)}`$ $`=1/2\mathrm{cos}(mn)\delta \varphi +\sigma ^2\delta _{mn}`$ (12)
$`𝒑^{(m)}`$ $`=1/2\mathrm{cos}(m+d)\delta \varphi ,`$ (13)
where $`(m,n)=0,1,\mathrm{},N1`$ and $`\delta \varphi =2\pi f_0\delta `$. Note that we have dropped the index $`k`$ because the autocorrelation $`𝑹`$ does not depend upon $`k`$, since we are dealing with a stationary signal.
From the above expressions of $`𝑹`$ and $`𝒑`$ and solving Eq.(4), we obtain the optimum Wiener filter
$$𝒘^{(m)}=\frac{2}{N+4\sigma ^2}\mathrm{cos}(m+d)\delta \varphi ,m=0,1,\mathrm{},N1.$$
(14)
where we have chosen the length of the filter to be half-integral number of cycles for reasons of simplicity, i.e. $`N\delta \varphi =l\pi `$, where $`l`$ is an integer.
In other words, the optimum linear predictor is nothing but a copy of the expected signal itself. The filter in Eq. (14) is also referred to as matched filter. In our situation, in practise, $`N\sigma ^2`$ and the term $`4\sigma ^2`$ can be omitted from the amplitude of $`𝒘^{}`$.
For the reasons detailed before, we propose to use the ALE algorithm in order to find a good approximation of $`𝒘^{}`$. Starting from an arbitrary initial tap-weight vector, we iterate the weights $`𝒘_k`$ according to Eq. (9) to converge to $`𝒘^{}`$. Once the filter is “close” enough to the optimal solution (the word “close” will be defined later in the text), we then say that the filter has locked on to the signal.
### B Approach to locking
#### a Continuous time approximation of the locking trajectory —
We may analyze the approach to locking by deriving a difference equation for the averaged evolution of the weights and then investigating this equation. It is impossible to obtain the average evolution of the weights by using the standard definition of the expectation operator E because of the nonlinearity and the recursive scheme involved in evolving the weights. We therefore adopt the time-average over successive data points as the operational definition of E.
Shifting the origin to $`𝒘^{}`$ by defining $`𝒗_k:=𝒘_k𝒘^{}`$, we may write the LMS evolution equation (9) in the following form :
$$𝒗_{k+1}𝒗_k=2\mu (𝒙_k𝒙_k^t)𝒗_k+2\mu e_k^{}𝒙_k,$$
(15)
where $`e_k^{}:=x_k𝒘^t𝒙_k`$ is the prediction error produced when using the optimal filter.
During the locking phase, the filter is far apart from the optimal location (i.e., $`𝒗_k`$ has a large modulus). The homogeneous term dominates the forcing term in the difference equation (15) which then can be approximated by :
$$𝒗_{k+1}𝒗_k=2\mu (𝒙_k𝒙_k^t)𝒗_k.$$
(16)
In the situation where the step gain parameter $`\mu `$ is chosen to be very small so that the weight coefficients are almost constant over a given time interval, the recursivity eventually acts as an averaging operation on both sides of the equation above. This leads to the difference equation which we use to describe the tap-weight trajectory in the space of weight coefficients, we denote $`𝒲`$ :
$$𝒗_{k+1}𝒗_k=2\mu 𝑹𝒗_k.$$
(17)
Let $`𝑸`$ be the transformation which diagonalizes $`𝑹`$. The above difference equation is best analyzed by changing the frame in $`𝒲`$ to the principal axis
$$\stackrel{~}{𝒗}_{k+1}\stackrel{~}{𝒗}_k=2\mu \stackrel{~}{𝑹}\stackrel{~}{𝒗}_k,$$
(18)
where $`\stackrel{~}{𝑹}:=𝑸𝑹𝑸^1=\mathrm{diag}(\lambda ^{(0)},\mathrm{},\lambda ^{(N1)})`$ and $`\stackrel{~}{𝒗}_k:=𝑸𝒗_k`$. Eq. (18) gives decoupled difference equations for the components $`\stackrel{~}{𝒗}_k^{(m)},m=0,\mathrm{},N1`$ of $`\stackrel{~}{𝒗}`$ which can be solved given the initial weight vector $`\stackrel{~}{𝒗}_0^{(m)}`$ :
$$\stackrel{~}{𝒗}_k^{(m)}=\stackrel{~}{𝒗}_0^{(m)}(12\mu \lambda ^{(m)})^k.$$
(19)
#### b Eigenvalues and eigenvectors of $`𝑹`$
We need to compute the eigenvalues $`\lambda ^{(m)}`$ of $`𝑹`$. This can be conveniently implemented by splitting $`𝑹`$ into the noise part, which is just $`\sigma ^2`$ times the identity plus the signal part which we denote by $`𝑺/2`$, and thus,
$$𝑹=\sigma ^2𝑰+𝑺/2,$$
(20)
where $`𝑺^{(mn)}:=\mathrm{cos}(mn)\delta \varphi `$. It is easily verified that the eigenvectors of $`𝑹`$ and $`𝑺`$ are identical and the eigenvalues of $`𝑹`$ are obtained from those of $`𝑺`$ by first halving them and then adding $`\sigma ^2`$ to the result. It remains, therefore, to compute the eigenvalues and eigenvectors of $`𝑺`$. We do this by observing that we can write $`𝑺`$ as follows By definition, the vectors $`\overline{𝒙}`$ and $`𝒙^{}:=\overline{𝒙}^t`$ denote respectively the complex conjugate and the hermitian transpose of $`𝒙`$.:
$$𝑺=(𝝂𝝂^{}+\overline{𝝂}\overline{𝝂}^{})/2,$$
(21)
where $`𝝂:=(1,\mathrm{exp}(i\delta \varphi ),\mathrm{exp}(2i\delta \varphi ),\mathrm{},\mathrm{exp}((N1)i\delta \varphi ))^t`$.
Since the matrix $`𝑺`$ is real and is essentially made out of two external products of $`𝝂`$ and $`\overline{𝝂}`$, its rank equals $`2`$ ($`𝑺`$ has $`N2`$ degenerate eigendirections in $`𝒲`$ with eigenvalue zero) with two non-zero real eigenvalues. Let $`𝒗`$ be an eigenvector associated to one of the non-trivial eigenvalues. According to the structure of $`𝑺`$, the vector $`𝒗`$ can be written without loss of generality as the following linear combination,
$$𝒗=𝝂\mathrm{exp}(i\alpha )+\overline{𝝂}\mathrm{exp}(i\alpha ),$$
(22)
where the coefficients have been chosen to have unit modulus arbitrarily.
Using the two scalar products $`𝝂^{}𝝂=N`$ and
$`𝝂^t𝝂`$ $`=1+\mathrm{exp}(2i\delta \varphi )+\mathrm{}+\mathrm{exp}(2(N1)i\delta \varphi )`$ (23)
$`:=\beta \mathrm{exp}i\gamma ,`$ (24)
where the geometric series can be summed up and modulus and phase ascertained :
$`\beta `$ $`={\displaystyle \frac{\mathrm{sin}(N\delta \varphi )}{\mathrm{sin}\delta \varphi }}`$ (25)
$`\gamma `$ $`=(N1)\delta \varphi ,`$ (26)
we obtain the effect of the matrix $`𝑺`$ on the vector $`𝒗`$, given by
$$𝑺𝒗=𝝂\mathrm{exp}(i\alpha )(N+\beta \mathrm{exp}(i\gamma +2i\alpha ))/2+c.c.$$
(27)
where c.c. denotes complex conjugate.
This expression has to be compared to the second term of the eigenvalue equation $`𝑺𝒗=\lambda 𝒗`$, leading to two solutions for $`\alpha `$, namely, $`\alpha =\gamma /2`$ and $`\alpha =(\gamma \pi )/2`$. These yield the eigenvectors $`𝒗_\pm `$ and the corresponding eigenvalues $`\lambda _\pm `$ :
$$\begin{array}{ccc}\hfill 𝒗_+& =& 𝝂\mathrm{exp}(i\gamma /2)+\overline{𝝂}\mathrm{exp}(i\gamma /2),\hfill \\ \hfill 𝒗_{}& =& i𝝂\mathrm{exp}(i\gamma /2)i\overline{𝝂}\mathrm{exp}(i\gamma /2),\hfill \end{array}$$
(28)
$$\lambda _\pm =(N\pm \beta )/2.$$
(29)
If we choose $`N`$ large enough and $`N\delta \varphi =m\pi `$, where $`m`$ is an integer then the analysis becomes simpler. This amounts to choosing the length of the filter to have half-integral number of cycles: we have $`\beta =0`$ and $`\lambda _\pm =N/2`$. (Geometrically, this means that the eigenvalue problem is degenerate with respect to the two signal eigenvectors: there is a two dimensional eigenspace belonging to the eigenvalue $`N/2`$. The weights thus evolve non-preferentially with respect to the signal eigendirections.) Since typical cases imply generally $`N\beta `$, we will assume this simplification in the rest of the paper.
In this situation, the spectrum of $`𝑹`$
$$\begin{array}{c}\text{sp}(𝑹)=\{\lambda ^{(0)}=\lambda ^{(1)}=N/4+\sigma ^2\text{ and }\hfill \\ \hfill \lambda ^{(m)}=\sigma ^2,m=2,\mathrm{},N1\},\end{array}$$
(30)
consist of two sets of eigenvalues : the first two correspond to directions in the signal space associated to “signal+noise” (or “signal”, for short) whereas the remaining $`N2`$ characterize “noise” directions.
According to Eq. (19), the weight vector will converge more rapidly in directions associated with the largest eigenvalues, which are the signal eigenvalues. The other noise eigenvalues are unimportant in this consideration. The eigenvectors pertaining to the signal provide preferred directions in $`𝒲`$: it is along these directions that the slope of the performance surface is steep and hence promotes faster convergence.
### C Steady state evaluation
If the step gain factor is sufficiently small, the tap-weight coefficients eventually converge and stabilize in a neighbourhood of the optimal value. At this stage, the assumptions made in obtaining the approximated evolution equation (17) do not hold anymore. In contrast with the case of “the approach to locking”, the right hand side of the difference equation (15) is now dominated by the forcing term :
$$𝒗_{k+1}𝒗_k=2\mu e_k^{}𝒙_k.$$
(31)
Roughly speaking, the trajectory of the vector $`𝒘_k`$ during the steady state can be viewed as a random walk centered around $`𝒘^{}`$ lying within a region of $`𝒲`$ space whose extent is determined by two factors, namely, $`\mu `$ and the intrinsic geometry of $`𝒲`$ in the vicinity of $`𝒘^{}`$.
The misalignment between the actual ALE filter $`𝒘_k`$ and the optimal one $`𝒘^{}`$ creates an additional error in the output. In fact, a direct calculation from Eq. (3) shows that the total mean square error may decomposed as :
$$J_k(𝒘_k)=\xi _{min}+\xi _k,$$
(32)
where (i) $`\xi _{min}:=J_k(𝒘^{})`$ is the minimum mean square error arising from the fraction of the input noise which still remains in the output, assuming that the ALE filter has reached exact optimality and (ii) $`\xi _k:=𝒗_k^t𝑹𝒗_k`$ is the excess mean square error (EMSE) due to the misalignment between the ALE filter and the Wiener filter.
One can verify that the EMSE vanishes when reaching optimality i.e., when $`𝒗_k=\mathrm{𝟎}`$. In other words, this term quantifies the non optimality of the current filter in use. We can imagine $`\xi _k`$ as the square of a natural distance in $`𝒲`$ and $`𝑹`$ as a intrinsic metric over $`𝒲`$.
A good approximation of $`\xi _{min}`$ can be found for large number of weight coefficients for the specific case of sinusoidal signals with high SNRs. Using Eqs. (3) and (4), we may write $`\xi _{min}=\text{E}[x_k^2]𝒑^t𝒘^{}`$. When $`N\mathrm{}`$, a direct calculation shows that the second term $`𝒑^t𝒘^{}`$ tends to the energy of the sinusoid, which means that the remaining energy is that of the noise: $`\xi _{min}\sigma ^2`$.
We complete the characterization of the mean square error (32) with the evaluation of the average value of the EMSE, which we denote by $`\xi ^{(st)}`$. Firstly, noticing that the EMSE is invariant under the principal axis transformation
$$\xi _k=\stackrel{~}{𝒗}_k^t\stackrel{~}{𝑹}\stackrel{~}{𝒗}_k=\underset{m=0}{\overset{N1}{}}\lambda ^{(m)}\left(\stackrel{~}{𝒗}_k^{(m)}\right)^2,$$
(33)
and secondly, using the approximation $`\text{E}[\stackrel{~}{𝒗}_k\stackrel{~}{𝒗}_k^t]\mu \xi _{min}𝑰`$ proposed in to obtain the typical value for $`(\stackrel{~}{𝒗}_k^{(m)})^2`$ yields
$$\xi ^{(st)}\mu \underset{m=0}{\overset{N1}{}}\lambda ^{(m)}\xi _{min}.$$
(34)
Since the signal is of much larger amplitude than the broadband noise, the trace of $`\stackrel{~}{𝑹}`$ is essentially due to the signal eigenvalues (see Eq. 30). Combining with the expression of $`\xi _{min}`$ above, this leads to :
$$\xi ^{(st)}\mu N\sigma ^2/2.$$
(35)
A better estimate of $`\xi ^{(st)}`$ can be obtained starting with more realistic hypotheses and using more sophisticated approximations :
$$\xi ^{(st)}\underset{m=0}{\overset{N1}{}}\frac{\mu \lambda ^{(m)}}{1\mu \lambda ^{(m)}}\xi _{min}.$$
(36)
In the limit of small step size, this approximation tends to the simpler one in Eq. (35).
### D Convergence time
In the expression of the EMSE in Eq. (33), we separate the sum into two parts: the first, $`\xi _k^{(n)}`$ associated with the noise (i.e., consisting of terms involving noise eigenvalues and vectors), the second, $`\xi _k^{(s)}`$ with the signal. Because the signal eigenvalues are much larger than those of the noise, the sum in Eq. (33) is essentially dominated in the beginning (for small $`k`$) by $`\xi _k^{(s)}`$. These two errors decrease during the locking phase until reaching a steady state value. The locking time (i.e., time at which the steady state is reached) is defined to be that, when $`\xi _k^{(s)}`$ is of the order of the total EMSE expected in the steady state.
From Eqs. (19), (29) and (33) we obtain,
$`\xi _k^{(s)}`$ $`:=\lambda ^{(0)}\left(\stackrel{~}{𝒗}_k^{(0)}\right)^2+\lambda ^{(1)}\left(\stackrel{~}{𝒗}_k^{(1)}\right)^2`$ (37)
$`{\displaystyle \frac{N}{4}}\left(\left(\stackrel{~}{𝒗}_0^{(0)}\right)^2+\left(\stackrel{~}{𝒗}_0^{(1)}\right)^2\right)\left(1{\displaystyle \frac{\mu N}{2}}\right)^{2k},`$ (38)
where we have assumed $`N\beta `$.
We set the starting point $`𝒘_0`$ in $`𝒲`$ to be $`\mathrm{𝟎}`$. It corresponds to the initial value $`\stackrel{~}{𝒗}_0`$ in the eigenspace which is given by $`\stackrel{~}{𝒗}_0=𝑸𝒘^{}`$. The first two coordinates of $`\stackrel{~}{𝒗}_0`$ can be directly obtained since the first two row vectors of $`𝑸`$ are just the normalized signal eigenvectors of the $`𝑹`$ matrix, leading to, for half integral wavelength filters and $`N\sigma ^2`$,
$`\stackrel{~}{𝒗}_0^{(0)}`$ $`=\sqrt{2/N}\mathrm{cos}(d\delta \varphi +\gamma /2)`$ (39)
$`\stackrel{~}{𝒗}_0^{(1)}`$ $`=\sqrt{2/N}\mathrm{sin}(d\delta \varphi +\gamma /2).`$ (40)
These considerations yield
$$\xi _k^{(s)}\frac{1}{2}\left(1\frac{\mu N}{2}\right)^{2k}.$$
(41)
This error must now be compared with the averaged EMSE in Eq. (35) in order to find the time $`t_{\mathrm{lock}}`$ at which $`\xi ^{(s)}`$ and $`\xi ^{(st)}`$ are equal :
$$t_{\mathrm{lock}}\delta \frac{\mathrm{ln}(\mu N\sigma ^2)}{2\mathrm{ln}(1\mu N/2)}.$$
(42)
It is important to mention that, when the product $`\mu N/2`$ tends to $`1`$, the convergence time diverges to infinity meaning that the weights do not converge toward $`𝒘^{}`$ anymore. In order to ensure the stability of the algorithm, the parameters will have to satisfy the stability condition $`0<\mu N/2<1`$. However, we have observed in our simulations that when $`1/2\mu N/2<1`$, the convergence is slowed down, because of the presence of oscillatory terms in the gradient which do not average to zero anymore. In practise, it is advisable to choose the parameters so that $`\mu N/2<1/2`$.
For a sinusoid of amplitude $`A`$ instead of unity as we have considered before, the condition for stability can be simply obtained by replacing the parameter $`\mu N/2`$ by $`\rho :=\mu NA^2/2`$ leading to $`0<\rho <1`$.
We illustrate in Figs. 2 and 3 with an example the results of this Section pertaining to the approach to locking and steady state analysis.
## IV The ALE in practice
In the previous Sections we have characterized the behaviour of the ALE in cases of interest. We will now elaborate on how this algorithm can be adapted to the interferometric data.
In the scheme we present here, we first decompose the signal in $`p`$ frequency subbands to which we apply the ALE twice with different sets of parameters. In the first stage, the parameters are tuned to best remove long-term sinusoidal components of the noise; whereas in the second stage, the target consists of shorter oscillatory transients.
### A Subband decomposition
Interferences such as mains power and violin mode harmonics are distributed over a large dynamic scale (the first harmonics are of much larger amplitude than those of high order). But, since the interferometer noise curve also decreases at low frequencies, their relative amplitude as compared to the background noise power spectrum at the same frequency remains large. Therefore, the model introduced previously, namely that of large amplitude sinusoidal signals embedded in broadband noise, is a reasonable approximation within the relevant small bandwidth of frequencies.
For this reason, we divide the frequency axis in $`p`$ disjoint frequency subbands of the same size. The $`p`$ signals lying in each of the subbands are heterodyned and decimated to the sampling frequency $`f_s^{\text{band}}:=f_s/p`$.
The tiling has the advantage that, if $`p`$ is sufficiently large, we can consider the interferometer background noise almost white within a subband, which implies that the noise has vanishing correlation time. The prediction depth $`d`$ which has to be larger than the correlation time, can be then simply fixed to any value greater than $`1`$ sample period in each of the subbands.
### B Long-term sinusoid removal
Certain parts of the spectrum may not contain any long-term periodic interferences. We apply a preliminary test to exclude subbands which may not require the first denoising step. The test is crudely done by estimating the amplitude $`A`$ of the sinusoid from the largest peak of the power spectrum (Welch estimate) and comparing it to the variance $`\sigma ^2`$ of the broadband noise (also estimated from the power spectrum). If it is found that $`A>\sigma `$, we decide that there exists a long-term sinusoidal signal of sufficient amplitude in the band which needs to be removed, otherwise we proceed directly to the second step.
We apply the ALE in each of the selected subbands choosing parameters as follows:
* Number of tap-weight coefficients $`N`$
The number of tap-weight coefficients is fixed by prescribing an upperbound $`0<\eta _{noise}<1`$ to the ratio between the noise power corrupting the filtered output $`y_k`$ of the optimal filter $`𝒘^{\mathbf{}}`$ and the input noise power. Let $`𝒏_k:=(n_{kdm},m=0,1,\mathrm{},N1)^t`$ be a collection of noise samples, then the above condition reads ,
$$\text{E}[(𝒘^t𝒏_k)^2]\eta _{noise}\text{E}[(n_k)^2],$$
(43)
which, with the stationarity and whiteness of the background noise $`n_k`$, results in bound on the optimal filter gain :
$$𝒘^t𝒘^{}\eta _{noise}.$$
(44)
The $`L^2`$ norm $`𝒘^{}_2^2=(2/N^2)(N+\beta \mathrm{cos}(\gamma +2d\delta \varphi ))`$ is obtained by squaring and summing the Eq. (14) for the optimal filter. Since in typical cases $`\beta N`$, this leads to simpler expression $`𝒘^{}_2^22/N`$.
Consequently, the number of tap-weight coefficients $`N`$ has to be chosen so that,
$$N2/\eta _{noise}.$$
(45)
* Step gain parameter $`\mu `$
We fix the step gain parameter by imposing to the distance of the ALE filter from optimality in the steady state to be smaller than a given threshold on average. As we have seen in Sect. III C, this can be done naturally by imposing an upperbound $`0<\eta _{sig}<1`$ on the excess square mean error as compared to the signal power $`E_s=A^2/2`$ :
$$\xi ^{(st)}\eta _{sig}E_s.$$
(46)
Using the expression obtained in the steady state analysis in Eq. (35) for the EMSE, this condition reduces to :
$$\mu \eta _{sig}/(N\sigma ^2).$$
(47)
Generally, this equation leads to small values of $`\mu `$ which prevent the convergence of the ALE filter from its initial state (i.e., all tap-weight coefficients are fixed to $`0`$) in a reasonable time (convergence faster than a tenth of second, which is the duration of the chunk of data). We solve this problem by first applying the ALE on a sequence of training data, the step gain parameter being set at the beginning to a large value (for fast convergence) and decreased gradually to the value given in Eq. (47). The filter obtained after the completion of this training is close to the objective (i.e., the Wiener filter). We then start the longterm sinusoid removal using this prepared filter.
We remark here that although $`\mu `$ is small, it is non-zero thus giving the ALE filter some flexibility of adapting to changes (non-stationarities) in the signal such as slow drifts in frequency and amplitude modulation. This property, however, needs to be investigated more in detail.
### C Ringdown removal
The aim of the second step of the algorithm is to remove oscillatory transients (ringdowns) of large amplitude. These transients are either frequency bands excited from time to time (caused by dysfunctions in the interferometer) or relics from the previous step (when the envelope of a long-term sinusoid possesses fast variations to which the algorithm cannot adapt or converge to during the first step of removal).
The cleaning procedure consists in applying ALE the second time to each of the subbands but now, the parameters are so adjusted that, (i) they select features with a larger bandwidth than in the previous step, and (ii) converge rapidly onto an oscillatory noisy signal that may appear.
* Number of tap-weight coefficients $`N`$
The impulse response duration and frequency selectivity (i.e., the filter bandwidth $`\mathrm{\Delta }f`$) of the transfer function are dual in character. This follows from the uncertainty relation. The rough approximate relation between these quantities is given by,
$$N=f_s/(p\mathrm{\Delta }f),$$
(48)
where $`f_s`$ is the sampling frequency. We choose the number of tap-weight coefficients $`N`$ by imposing a minimum bandwidth $`\mathrm{\Delta }f_{min}`$ to the filter and using the above equation.
* Step gain parameter $`\mu `$
Assuming that the ringdown can be locally approximated by a sinusoid, we choose the step gain parameter by imposing a convergence time of the order of a typical transient duration (i.e., $`t_{lock}N\delta `$). More concretely, setting $`\rho :=\mu NA^2/2`$ in the unnormalized form of Eq. (42) (i.e., for arbitrary ringdown amplitude $`A`$), we solve for
$$\frac{\mathrm{ln}(2\sigma ^2\rho )}{2\mathrm{ln}(1\rho )}=N.$$
(49)
Using the crude estimate $`A^2/2𝒙_k_2^2/N`$ for the ringdown amplitude, the step gain parameter is finally obtained as $`\mu =\rho /𝒙_k_2^2`$.
Since the ringdown signals are of short duration and can occur with large time gaps, the ALE does not need to operate on each data segment. Accordingly, we have added a supervision test which decides whether or not the denoising algorithm should be applied to a given data segment. The test consists of observing the Gaussianity of the filtered output $`y_k=𝒘_k^t𝒙_k`$. If the input signal $`x_k`$ is a zero-mean white Gaussian process of variance $`\sigma ^2`$, then the output of the filter $`y_k`$ shares the same characteristics, except that the variance gets multiplied by the filter gain : $`\text{var }y_k=𝒘_k_2^2\sigma ^2`$. Furthermore, under this hypothesis, the envelope $`𝒴_k=|(y)_k|^2`$ ($``$ denotes the discrete Hilbert transform The discrete Hilbert transform $`y_n=(x)_n`$ of a signal $`x_n`$ is essentially obtained by cancelling its negative frequencies; more precisely, $`Y(f):=2U(f)X(f)`$, with $`U(f)=1`$ when $`f[0,1/2]`$ and $`0`$ when $`f]1/2,0[`$ and where $`X(f)`$ (and $`Y(f)`$) denotes the Fourier transform of the corresponding signal $`X(f)=_{n=0}^Nx_ne^{2\pi inf}`$. of $`y_k`$) follows by definition a chi-square distribution with $`1`$ degree of freedom.
This implies that, up to an arbitrary probability $`P_0`$, the envelope $`𝒴_k`$ does not exceed the threshold given below:
$$𝒴_k<\kappa (P_0)\mathrm{\hspace{0.17em}2}f_s^{band}𝒘_k_2^2\sigma ^2,$$
(50)
where $`\kappa ()`$ is the inverse function of the (unit variance) $`\chi ^2`$ cumulative distribution function (cdf).
If Eq. (50) is satisfied, we conclude that the filtered output is essentially due to a Gaussian background noise and we leave the input signal as it is. Otherwise, we conclude that the filtered output carries a ringdown signal and decide to remove it from the input data.
The functioning of the second step of the denoising algorithm could be interpreted as follows : it removes from the input data, regions in the time-frequency plane presumably associated with transients, whose support is defined along the frequency axis by the ALE filter, and along the time axis by the supervision criterion (50).
After completing these two steps, we recombine the signal in all the subbands together to retrieve a single strain signal.
## V Numerical results
### A Simulated data: test of the ringdown removal
In this section, the goal is to test how effectively the second stage of the denoising algorithm (i.e., the ringdown removal) described in Sect. IV operates on a simple signal. The test signal is composed of three ringdown signals (of fixed amplitude and frequency) occurring successively in the data stream and embedded in a additive Gaussian white noise. This model may be used to represent ringdown disturbances originating from the same underlying physical mechanism.
Each of these ringdown signals is a sinusoidal waveform, similar to Eq. (10) (with $`A=1`$, $`f_0=50`$ Hz and sampling frequency $`f_s=200`$ Hz), whose support is limited in time by a Gaussian envelope :
$$r_k:=A\mathrm{exp}\left(\pi (t_kt_c)^2/T^2\right)\mathrm{cos}(2\pi f_0t_k+\mathrm{\Phi }),$$
(51)
where three different reference times $`t_c`$ are given and the equivalent time duration is $`T=200`$ms (giving a frequency bandwidth of $`\mathrm{\Delta }f1/T=5`$ Hz and $`Q:=f_0T10`$ cycles).
Figure 4 describes the application of the denoising algorithm configured with $`d=5`$ sampling periods (equal to $`25`$ ms) $`\mathrm{\Delta }f_{min}=3`$ Hz, and $`P_0=0.01`$. It can be seen that the algorithm operates better on the transient encountered later in the data train than its predecessor. The explanation is that a transient duration is too short for the filter to reach the steady state but, when it encounters the next transient, the filter benefits from the distance to $`𝒘^{}`$ previously covered, thus improving the convergence towards optimality.
This can be verified with a time-frequency representation of the output signal such as Fig. 4, where we have chosen the spectrogram $`S_x^h[n,m]:=|F_x^h[n,m]|^2`$ defined as the squared modulus of the short-time Fourier transform :
$$F_x^h[n,m]:=\underset{k}{}x_nh_{kn}e^{2\pi inm},$$
(52)
where $`n[1,2,\mathrm{},N]`$, $`m]1/2\mathrm{}1/2]`$ and $`h_k`$ is an arbitrary window (a Gaussian window here).
Notice that real time and frequency coordinates can be retrieved through the relations : $`t=n/f_s`$ and $`f=mf_s`$.
### B Results on Caltech 40m proto-type data
Here we have applied the algorithm to the Caltech 40meter proto-type data taken in October 1994 . This data was recorded with a sampling frequency of $`f_s=9.86`$kHz. We have used the calibrated strain signal (relative arm length measurement) for applying our algorithm.
We tile the complete spectrum into $`p=32`$ frequency subbands of approximately $`154`$ Hz each. Each subband encounters typically one or two long-term sinusoidal interferences.
We have chosen the prediction depth to be $`d=5`$ sampling periods, which corresponds to a delay of $`pd/f_s16`$ ms in real time. The correlation time of the broadband noise is effectively smaller in each subband except at the extremities of the spectrum where the steep slope of the spectrum does not allow us to assume the background noise to be locally white. It only affects the first and last subbands which are not too important for detection purposes.
In the first stage, we have chosen $`\eta _{noise}=0.01`$ (giving $`N=200`$ according to Eq. (45)) and $`\eta _{sig}=0.01`$. In the second stage of ringdown removal, the minimum filter bandwidth has been fixed to $`\mathrm{\Delta }f_{min}=3`$ Hz, which gives a filter with $`N=100`$ tap-weight coefficients (see Eq. (48)) and we have set $`P_0=0.01`$ for the Gaussianity test.
We have performed two types of simulations:
* a “Caltech signal only” simulation to measure improvements after denoising : we check firstly, whether the frequency peaks are removed from the noise power spectrum and secondly, whether the noise statistics is closer to Gaussian than before denoising,
* a “Caltech+inspiral” simulation to evaluate the consequences of the denoising algorithm on gravitational wave detection; specifically, for the case of the inspiralling compact binary signal. The question here is to check whether the denoising operation has removed a significant part or even whole of the inspiral signal.
Caltech signal only — Eleven of the thirty-two frequency subbands (# $`1`$$`9`$, $`11`$ and $`17`$) are selected and sent to the first cleaning step of the algorithm. In these subbands, we obtained the following mean values for $`A1.5\times 10^{16}`$ and $`\sigma 3.6\times 10^{17}`$ (the sinusoid amplitude $`A`$ equals approximately $`1`$ to at most $`5`$ times the noise standard deviation $`\sigma `$) leading to typical values for the signal-to-noise ratio of about $`\text{SNR}=A^2/(2\sigma ^2)8.7(9.4\mathrm{𝖽𝖡})`$ and for the step gain parameter (see Eq. (47)) of $`\mu NA^2/20.04`$ (spanning from $`0.01`$ to $`0.14`$).
The complete set of subband signals is processed in the second step. The typical noise variance estimate is $`\sigma =1.35\times 10^{17}`$ (from $`4.8\times 10^{18}`$ to $`10^{16}`$) leading according to Eq. (49) to values of $`\mu N\sigma ^2`$ which span the range of values from $`0.07`$ to $`10^4`$.
Figure 5 illustrates how the algorithm operates in the fifth frequency subband (from $`617`$ Hz to $`771`$ Hz) among the $`p=32`$ ones being processed. This frequency band contains two power line harmonics (the $`11`$th at $`660`$ Hz and the $`12`$th at $`720`$ Hz).
Figures 6 and 7 show respectively comparisons between the power spectra and histograms of the signal before and after denoising. We observe that after denoising, the frequency peaks have been removed from the input signal and the histogram appears much closer to the Gaussian bell curve.
Caltech signal + inspiral waveform — The purpose of this test is to evaluate how the cleaning operation affects gravitational wave detection and in particular to make sure whether a significant part of the gravitational signature could be removed from data. Answering this question by analytical means is difficult, however a qualitative rational in the case of inspiral binaries can be made and verified with simulations.
The theory predicts that the gravitational waves emitted from inspiralling binaries of neutron stars are oscillating waveforms whose frequency evolves in time in a prescribed manner and scans the interferometer bandwidth from lower end to the higher.
Their weak amplitude and short time duration within a single subband (in the case we have considered, less than a second) make them “invisible” to the ALE filter. The amplitude and the duration of the gravitational wave signal are simply not large enough for the ALE coefficients to converge onto the gravitational wave instantaneous frequency.
We have checked the validity of this argument by adding to the Caltech signal the inspiralling ‘chirp’ waveform in the Newtonian approximation of two neutron star binaries each having a mass of $`M=1.4`$ solar masses, and located at a distance of $`r=7`$kpc from the Earth.
Figure 8 depicts a comparison of matched filter detector response on the same signal with and without denoising. The detector output displays a peak of the same height and at the correct instant, showing that the cleaning algorithm has not removed the inspiral signal from the data. This can be crosschecked in Fig. 9 showing a zoomed view of the same signal after denoising.
## VI Concluding remarks
The originality of the idea of the proposed denoising algorithm lies in its wide applicability, so that both types of disturbances, long-term sinusoidal and oscillatory transients (the type of noise which has been ignored till now) can be treated. Although the question of the computational burden in applying this algorithm has not quite been addressed here, it appears from the simplicity of the operations involved (e.g., no requirement such as long-term FFTs) that the total computational cost should be within acceptable limits, so that the algorithm can be operated in real time. Furthermore, the structure of the algorithm already implemented with Matlab can be easily translated into a parallel code (each processing node can be associated with one frequency subband and the processing can be done independently).
As part of future extensions to the present work, some improvements to the current code might be needed : in order to limit the finite size effects in the subband decomposition and reconstruction, a reversible filter bank (e.g., a Gabor transform) would be preferable than the crude method used here.
The key idea (i.e., looking for correlation between the current sample of the strain signal and a reference signal, namely a set of past samples) can be also extended to investigate correlations of the detector output with other environmental channels by simply using them as a reference rather than the strain signal itself. Similarly to the cross-talk removal in but with adaptive methods, such an algorithm would provide an estimation of any poorly known (linear) transfer functions relating noise sources to their final leaking in the detector output and of the environmental contamination that must be subtracted from the data, if so desired.
### Acknowledgments
We would like to thank B. F. Schutz for suggesting the idea of adaptive methods and also for fruitful conversations and the LIGO collaboration for providing us the Caltech 40meter proto-type data. E. C.-M. would like to thank W. Anderson, R. Balasubramian, J. Creighton and S. Mohanty for their useful comments and suggestions.
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# Local measurements of nonlocal observables and the relativistic reduction process
## 1 Introduction
As it is well known, standard quantum mechanics is characterized by irreducible stochastic features which enter into play and give rise to puzzling situations in connection with the measurement process or, more appropriately and more generally, with the so called macro-objectification problem. It is useful to stress that it is precisely with reference to such processes which, when one takes further into account the unavoidable nonlocal character of the theory, the problem of the causal relations between events (such as the measurement outcome in one region and the transition from potential to actual physical situations in another) emerges as a central one. It goes without saying that an adequate discussion of such questions requires a relativistic approach. Our main concern here will be the analysis of statevector reduction in a relativistic quantum context and the identification of the basic features which any relativistic reduction mechanism must exhibit. In particular we will show that recent attempts in this direction -, even though not fully satisfactory due to some technical problems (such as the appearence of untractable divergencies) have given clear indications about the line of thought one should follow to account for the nonlinear and stochastic reduction process meeting the strict demands of special relativity<sup>2</sup><sup>2</sup>2Some of the arguments we will discuss in this paper have been already presented in refs. and . However, the new approach to the problem of local measurements of nonlocal observables which is presented in section 3 will allow us to develope in a clearer and more convincing way our arguments and to draw the relevant conclusions of the final section..
In section 2 we reconsider the crucial aspects of the problem under investigation by following the recent lucid presentation of the matter by Breuer and Petruccione . Particular attention will be devoted to the fundamental contributions to the subject by Aharonov and Albert , and . These authors have identified the necessary formal features that any acceptable relativistic reduction mechanism must exhibit and have shown that all previous attempts in this direction were fundamentally unsatisfactory. The key argument which has led Aharonov and Albert to draw such a conclusion is their proof that one can measure nonlocal observables by resorting to local interactions between appropriately (and smartly) chosen quantum systems, followed by local detection processes.
There are, however, two reasons which require to deepen and to improve the line of thought of these authors. First, their explicit example of a local measurement of nonlocal observables is quite formal and rests on the consideration of nonnormalizable states. Secondly, in spite of the fact that they have given clear indications about the general formal features of any acceptable relativistic dynamical reduction mechanism, they have made no attempt to work out an explicit example of such a dynamics. On the other hand, as already mentioned, in recent years precise relativistic models of dynamical reduction have been formulated -. It is one of the aims of this paper to investiagte critically such models from the point of view of the analysis of Aharonov and Albert.
To this purpose it is useful to work out, first of all, a remarkably simplified version of the measurement procedure suggested by the above authors in ref.. This is achieved in section 3 by relating the relevant physical process to degrees of freedom whose associated Hilbert space is finite–dimensional. In this way the treatment turns out to be extremely simple and intuitive and the procedure easy to implement in practice. This simplified version of the proposal of ref. will allow us to focus the essential features of any relativistic reduction theory and to show that such features are essential and natural ingredients of the theoretical framework presented in refs.-. Accordingly, in the rest of the paper we will discuss how such precise theoretical schemes account for the relativistic macro-objectification process and we will show that they lead to a perfectly coherent view which allows one to analyze basic issues such as the attribution of objective properties to individual physical systems, the necessary generalization of the concept of event and the appropriate way to resort to counterfactual arguments within a relativistic and nonlocal context.
## 2 Relativistic reduction processes
This section is devoted to reviewing some of the issues of the central theme of the paper, i.e., the relativistic aspects of the reduction process.
### 2.1 Local relativistic reduction processes
The main problems which one meets in trying to generalize the nonrelativistic process of statevector reduction derive from the assumed instantaneity of such a process. As lucidly described in ref. one can consider the case in which one has a system whose associated wavefunction has an appreciable spatial extension and, at time $`t=0,`$ is found at $`x=0`$ by a detector which is placed there. The problem one has to face is quite obvious: even if one were able to account for the local position measurement in terms of a local covariant interaction between the measured object and the measuring device, the ensuing picture would obviously turn out not to be covariant for the very simple reason that in any other reference frame the space-like surface $`t=0`$ is not an equal time surface; consequently, the reduction cannot be instantaneous for any observer in motion with respect to the original one.
Hellwig and Kraus have proposed to circumvent the above difficulty by postulating that in a local measurement process statevector reduction takes place along the past light cone originating from the space-time point at which the covariant system-apparatus interaction occurs. In contrast with the case mentioned above it is obvious that the proposed prescription is manifestly covariant since the light cone from a given point is the same in all reference frames. These authors have also shown that their prescription leads to the correct quantum predictions concerning the probabilities of the outcomes of local measurements of local observables. In the early stages of the debate about relativistic reduction processes the expression “local” has been used in a quite loose way which we will follow in this subsection with the purpose of making quite intuitive the problems one meets. Suppose our system, let us say an elementary particle, prior to the local position measurement at $`t=0,`$ $`x=0,`$ is described, in a given reference frame, by a wavefunction which has an appreciable extension in space (the limiting case would be the one of a system in a state of definite momentum). A measurement in which the particle is found at $`x=0`$ will, according to the Hellwig and Kraus prescription, lead to a statevector which is different from zero only within the past light cone from the origin, contradicting the fact that at a time prior to $`t=0`$ the wavefunction extends over a much larger region as, e.g., in the limiting case of a momentum eigenstate, which, if subjected to a momentum measurement, would remain unaltered. According to Hellwig and Kraus the way out from such a difficulty derives from assuming that, e.g., the alleged momentum measurement cannot be performed by resorting to local interactions. Actually, already in 1931, Landau and Peierls had suggested that all nonlocal quantities, like the momentum operators, cannot be considered as observables in relativistic quantum theories. Hellwig and Kraus adopt this point of wiev and, consequently, they dismiss as non pertinent any criticism concerning their proposal. Here comes the fundamental contribution by Aharonov and Albert . They have shown that it is actually possible to measure nonlocal observables and, even more important, that this can be done by local interactions and local detections.
Before going on, it is appropriate to stress a point of great conceptual relevance for the subject of this paper. In the debate we have just reviewed it was always assumed (more or less tacitly) that the wavefunction $`\mathrm{\Psi }(x,t),`$ in the relativistic case must be considered as a function on the space-time continuum. As we will see, the analysis of Aharonov and Albert requires a radical change of perspective about this fact.
### 2.2 Nonlocal measurements in a relativistic context and their difficulties
As already anticipated, within a relativistic context, nonlocal observables raise particular problems in connection with measurement processes. To better focus this point, in place of the rather vague considerations of the previous subsection related to momentum measurements, we present a slightly modified version of a simple and nice example by Breuer and Petruccione .
Let us consider (Fig.1) a given inertial reference frame and a pair of particles having space plus internal degrees of freedom. For simplicity, since these assuptions do not change in any way whatsoever the conceptually relevant aspects of the analysis, we will disregard the spreading of the wavepackets of the two particles and we will assume that the internal degrees of freedom refer to quantities, like the isospin, which do not transform under Lorentz transformations. Suppose the two particles are created, e.g., in a decay process at a given objective space-time point and that they propagate in opposite directions. Concerning the internal degrees of freedom, let us assume that the two particles have isospin $`1/2`$ and are in the common eigenstate of $`T^2`$ and $`T_z`$ belonging to the eigenvalues (in the appropriate units) $`2`$ and $`1`$, respectively, i.e., in the state we will denote (in analogy with the case of spin variables) as $`|_{1z},_{2z}>.`$ Suppose that, along the world line of particle 1 (the one propagating at left), at a precise space-time point $`L(x_1,t_1ϵ)`$, there is a device which can be switched on and off at the experimenter’s whim, whose only (local) effect is that of inducing a isospin flip of the particle interacting with it when it is on, while no change occurs when it is off. Suppose also that, at time $`t_1`$, a measurement of the nonlocal observable $`T^2`$ is performed on the composite system. We stress that since the two particles are in different space regions the measurement of $`T^2`$ is unavoidably nonlocal. Finally let us assume that at time $`t_1+ϵ,`$ at a point $`R(x_2,t_1+ϵ)`$ on the world line of particle 2 there is an apparatus measuring $`T_{2z},`$ the $`z`$ component of the isospin of particle 2.
Given these premises one can argue as follows:
* Suppose the apparatus at $`L`$ is off. Then the state of the composite system at $`t_1`$ is $`|_{1z},_{2z}>,`$ which, being an eigenstate of $`T^2,`$ is left unchanged by the measurement of this observable. There follows that in the subsequent measurement of $`T_{2z}`$ the outcome +1 is obtained with certainity.
* Suppose now the apparatus at $`L`$ is on. Then the state of the composite system at $`t_1ϵ`$ becomes $`|_{1z},_{2z}>,`$ which is the superposition with equal squared amplitudes of the eigenstates of $`T^2`$ belonging to the eigenvalues 0 and 2, respectively. The measurement taking place at $`t_1`$ will then reduce the state to one of such eigenstates, i.e., $`\frac{1}{\sqrt{2}}[|_{1z},_{2z}>`$ $`|_{1z},_{2z}>],`$ which assign probability $`1/2`$ both to the outcome +1 and to –1 in the measurement of $`T_{2z}`$ at $`t_1+ϵ`$. If the observer at $`R`$ is informed about the experimental set-up, in the cases in which he gets the outcome –1 he can infer that the observer at $`L`$ has actually decided to induce the spin flip of particle 1. This argument shows quite nicely that nonlocal measurements can lead to a violation of causality allowing superluminal communication and makes clear while such measurements must be excluded if one wants to take the position of Hellwig and Kraus about statevector reduction.
### 2.3 The Aharonov and Albert procedure
The simplest example of a measurement of a nonlocal observable proposed by these authors is, in our language, a measurement of the $`z`$component of the total isospin $`T_z=T_{1z}+T_{2z}`$ of the previous system. The smart trick consists in considering an apparatus constituted by two subsystems (probes) whose world lines intersect the world lines of particles 1 and 2 (once more we disregard the spreading associated with the free motion of the probes) and considering “internal generalized coordinates $`\widehat{q}_1`$ and $`\widehat{q}_2`$”of the probes. The local interactions of the probes with the particles are described by the hamiltonian:
$$H_I=g(t)\left[\widehat{q}_1T_{1z}+\widehat{q}_2T_{2z}\right],$$
(2.1)
where $`g(t)`$ is a time dependent coupling constant vanishing outside a small time interval $`(t_1,t_2)`$. Finally, one assumes that the probes are, at any time $`t_0`$ preceeding their interactions with the particles, in an entangled state of the kind of the one considered in the celebrated EPR paper, i.e., one in which both the coordinates $`\widehat{q}_1`$ and $`\widehat{q}_2`$ and the conjugated momenta $`\widehat{\pi }_1`$ and $`\widehat{\pi }_2`$ are perfectly correlated:
$$\widehat{\pi }_1(t_0)+\widehat{\pi }_2(t_0)=0$$
$$\widehat{q}_1(t_0)\widehat{q}_2(t_0)=0.$$
(2.2)
As in the standard measurement scheme by von Neumann the local interactions described by eq.(2.1) imply the following equation of motion for the total canonical momentum $`\pi (t)`$:
$$\frac{\pi (t)}{t}=i[H_I,\pi (t)]=g(t)\left[T_{z1}+T_{z2}\right]=g(t)T_{z,}$$
(2.3)
so that one can infer the value of $`T_z`$ from the change of the total momentum:
$$T_z=\frac{\pi (t>t_2)}{_{t_1}^{t_2}𝑑tg(t)},$$
(2.4)
where we have taken into account that $`\pi (t<t_1)=0,`$ as shown by the first of eqs.(2.2). To make eq.(2.4) more understandable to the reader, we remark that the local interaction of one probe with the corresponding particle decreases (increases) the value of the associated momentum according whether the particle is in the eigenstate $`T_{iz}=+1`$ $`(1),`$ respectively. Accordingly, the total momentum changes when the two particles have parallel $`z`$isospin components and remains unchanged when their projections are opposite.
It is extremely important to realize that the measurement is a genuinely nonlocal measurement performed by means of local interactions and detections. In fact it is easy to convince oneselves that the knowledge, after the measurement, of the value of the momentum of only one of the probes does not allow to draw any conclusion referring to the isospin component of the particle which has interacted with the probe, or to the component of the total isospin. This point will become much more clear when we will consider the simplified version of a measurement of this kind which is the subject of the following Section.
Some remarks are at order:
* The two probes considered in refs. have translational degrees of freedom related to their propagation and interactions with the particles and further internal degrees of freedom which, however, are associated to continuous generalized coordinates and momenta. Moreover, the state of the probes before the measurement exhibits perfect “internal” momentum and position correlations, just as the state of the original EPR argument. Leaving aside the problem of the difficulty of preparing such an entangled state in these continuous variables, the resulting state turns out to be nonnormalizable. It should be clear to the reader that the very cute proposal by Aharonov and Albert corresponds to a quite idealized situation. Just in the same way in which Bohm has felt the necessity of rephrasing the EPR argument resorting to the consideration of spin degrees of freedom, transforming it from a gedanken to a feasible experiment, we consider it useful to perform an analogous step with reference to the process under discussion.
* The above procedure has been generalized by the authors of refs. to measure the square $`T^2`$ of the total isospin of the composite system. Such a step involves the use of three pairs of entangled probes which interact with particles 1 and 2. Only the comparison of the outcomes obtained at the two wings of the apparatuses devised to measure the momenta of the three probes allows, in the case in which all pairs of momenta sum up to zero, to assert that the system has been found in the singlet state. The measurement is fundamentally a quantum nondemolition measurement for the singlet state, while it is not a measurement of the standard type for the states of the three dimensional manifold $`T^2=2.`$ Once more these subtle questions will become clear in the modified version of a local measurementof nonlocal observables we are going to discuss.
## 3 A new procedure for a local measurement of nonlocal observables
In this Section we present a possible experimental setup for the measurement of the observables $`T_z`$ and $`T^2`$ of the system S of two particles of isospin 1/2 which is based on local interactions between the particles and probe particles whose internal degrees of freedom are accounted by the states of a three-dimensional Hilbert space (Fig.2).
To begin, we consider, as before, two particles 1 and 2 which are produced at a given space-time point and propagate along two world lines. Since we disregard the spreading of the wavefunctions of the particles, their quantum features are entirely described by the vectors of the four-dimensional Hilbert spaced spanned by the states (with obvious meaning of the symbols):
$$|_{1z},_{2z}>,|_{1z},_{2z}>,|_{1z},_{2z}>,|_{1z},_{2z}>.$$
(3.1)
We consider also a system $`\mathrm{\Sigma }`$ of two probe particles (which for reasons which will become clear in what follows we identify as 3 and 6) whose world lines intersect, at the space time points $`P1`$ and $`P2`$, the world lines of the system particles. The probe particles, as already stated, besides their configurational degrees of freedom (which we treat once more as classical), have internal degrees of freedom associated to a three-dimensional Hilbert space. We denote as
$$|1>_i,|0>_i,|+1>_i,(i=3,6),$$
(3.2)
the complete set of the eigenstates of an observable $`\mathrm{\Omega }^{(i)}`$ of the Hilbert space of the i-th probe and by
$$|a,b>_{3,6}|a>_3|b>_6,(a,b=1,0,+1)$$
(3.3)
the elements of the corresponding basis in the nine-dimensional Hilbert space of $`\mathrm{\Sigma }`$. The initial state of the probe is assumed to be the normalized entangled state of particles 3 and 6:
$$|\mathrm{\Phi }>_{3,6}=\frac{1}{\sqrt{3}}\left[|0,0>_{3,6}+|+1,1>_{3,6}+|1,+1>_{3,6}\right].$$
(3.4)
The system and probe particles interact locally (constituent 1 (2) of $`S`$ with constituent 3 (6) of $`\mathrm{\Sigma }`$) for a given time interval, the effect of the interaction being described by the unitary operator $`U_z:`$
$$U_z=\left[P_{z+}^{(1)}P_L^{(3)}+P_z^{(1)}P_R^{(3)}\right]\left[P_{z+}^{(2)}P_L^{(6)}+P_z^{(2)}P_R^{(6)}\right],$$
(3.5)
where the operators $`P_{z\pm }^{(i)}`$ are the projection operators on the eigenmanifolds of $`T_{iz}`$:
$$P_{z+}^{(i)}=|_{iz}><_{iz}|,P_z^{(i)}=|_{iz}><_{iz}|,(i=1,2)$$
and the operators $`P_{R,L}^{(j)}`$ act on the probe space and have, in the basis (3.2), the representation indicated below:
$$P_L^{(j)}=\left(\begin{array}{ccc}0\hfill & 0\hfill & 1\hfill \\ 1\hfill & 0\hfill & 0\hfill \\ 0\hfill & 1\hfill & 0\hfill \end{array}\right)_j,P_R^{(j)}=\left(\begin{array}{ccc}0\hfill & 1\hfill & 0\hfill \\ 0\hfill & 0\hfill & 1\hfill \\ 1\hfill & 0\hfill & 0\hfill \end{array}\right)_j,(j=3,6).$$
(3.6)
We note that the operators $`P_L^{(j)}`$ and $`P_R^{(j)}`$ act in the following way on the states of the base of probe j:
$$P_L^{(j)}|1>_j=|+1>_j,P_L^{(j)}|0>_j=|1>_j,P_L^{(j)}|+1>_j=|0>_j,$$
(3.7)
$$P_R^{(j)}|1>_j=|0>_j,P_R^{(j)}|0>_j=|+1>_jP_R^{(j)}|+1>_j=|1>_j,$$
(3.8)
so that they produce an anticlockwise and a clockwise permutation of the three states ($`|1>,|0>,|+1>`$) of the basis, respectively. We also note the following properties of these operators:
$$P_R^{(j)}=P_L^{(j)},P_L^{(j)}=P_R^{(j)},P_R^{(j)}P_L^{(j)}=P_L^{(j)}P_R^{(j)}=1^{(j)}.$$
(3.9)
The above relations imply:
$`U_z^{}U_z`$ $`=`$ $`[P_{z+}^{(1)}P_R^{(3)}+P_z^{(1)}P_L^{(3)}][P_{z+}^{(2)}P_R^{(6)}+P_z^{(2)}P_L^{(6)}]`$
$`\left[P_{z+}^{(1)}P_L^{(3)}+P_z^{(1)}P_R^{(3)}\right]\left[P_{z+}^{(2)}P_L^{(6)}+P_z^{(2)}P_R^{(6)}\right]`$
$`=`$ $`\left[P_{z+}^{(1)}1^{(3)}+P_z^{(1)}1^{(3)}\right]\left[P_{z+}^{(2)}1^{(6)}+P_z^{(2)}1^{(6)}\right]=1^{(1)}1^{(2)}1^{(3)}1^{(6)}=1,`$
so that $`U_z`$ is a unitary operator.
Given the above equations it is immediate to evaluate the effect of applying $`U_z`$ to the direct product of any one of the states (3.1) times the initial state of the probe $`|\mathrm{\Phi }>_{3,6}.`$ To this purpose it is useful to introduce the following normalized states of the probe system (with obvious meaning of the symbols):
$`|\mathrm{\Pi }(1,2)>_{3,6}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{3}}}\left[|1,1>_{3,6}+|0,+1>_{3,6}+|+1,0>_{3,6}\right]`$ (3.11)
$`|\mathrm{\Pi }(2,1)>_{3,6}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{3}}}\left[|+1,+1>_{3,6}+|0,1>_{3,6}+|1,0>_{3,6}\right].`$
As easily checked we then have:
$`U_z[|_{1z},_{2z}>|\mathrm{\Phi }>_{3,6}]`$ $`=`$ $`|_{1z},_{2z}>|\mathrm{\Pi }(1,2)>_{3,6}`$ (3.12)
$`U_z[|_{1z},_{2z}>|\mathrm{\Phi }>_{3,6}]`$ $`=`$ $`|_{1z},_{2z}>|\mathrm{\Pi }(2,1)>_{3,6}`$
$`U_z[|_{1z},_{2z}>|\mathrm{\Phi }>_{3,6}]`$ $`=`$ $`|_{1z},_{2z}>|\mathrm{\Phi }>_{3,6}`$
$`U_z[|_{1z},_{2z}>|\mathrm{\Phi }>_{3,6}]`$ $`=`$ $`|_{1z},_{2z}>|\mathrm{\Phi }>_{3,6}`$
The model we have presented, in complete analogy with the one proposed by Aharonov and Albert, describes a nonlocal measurement procedure (in terms of local interactions and detections) of the nonlocal observable $`T_z,`$ the z-component of the total isospin of the composite system $`S`$. In fact, in order to complete the measurement, we have only to put on the world lines of the probe particles two apparata measuring the observables $`\mathrm{\Omega }^{(i)}`$(i=3,6), to register the obtained outcomes and to sum them. If they sum up to 0, then reduction takes place to a state for which $`T_z=0,`$ if they sum up to either 1 or –2 then reduction takes place to the state for which $`T_z=1,`$ finally, if they sum up to either 2 or –1, reduction takes place to the state for which $`T_z=1.`$
It is of great relevance to remark that for each of the states of the probes appearing at the right hand side of eq.(3.12) all possible eigenvalues ($`\omega _k^{(i)}=1,0,1)`$ of the observable $`\mathrm{\Omega }^{(i)}`$ referring to each probe particle appear with equal weights. Thus, knowledge of the outcome, e.g., of the detector at the right wing of the apparatus, does not give any information whatsoever about the isospin component of the particle with which the probe particle has interacted locally, nor about the component of the total isospin. In this precise sense the measurement is genuinely nonlocal (as it must be, since the two particles lie in far apart regions).
We can now proceed to generalize the previous procedure to a measurement process for $`T^2.`$ To this purpose we have to resort to a more complicated device corresponding to the production of three entangled pairs of probe particles which interact locally with the particles whose total isospin is being measured. So, besides the pair (3-6) of probe particles we will consider also the pairs<sup>3</sup><sup>3</sup>3We have decided to use the indices (3,6), (2\*,4\*) and (3\*,6\*) to refer to the three pairs of probe particle for obvious reasons: the specification 3 and 6 (2 and 4) refer to degrees of freedom which are locally coupled to the third, i.e. the z, (the second, i.e. the y) component of the isospin of the measured particles. The stars are used to avoid confusion with the indices 1,2 of the measured particles or to stress that the last pair (3\*,6\*) of probe particles is different from the first one. (2\*-4\*) and (3\*-6\*). The initial state of the whole probe will simply be the direct product of three states like (3.4):
$$|\mathrm{\Phi }>=|\mathrm{\Phi }>_{3,6}|\mathrm{\Phi }>_{2,4}|\mathrm{\Phi }>_{3,6}$$
(3.13)
and the unitary evolution operator $`U`$ for the whole system is assumed to be the product of three operators: $`\stackrel{~}{U}_zU_yU_z`$. In turn, $`U_z`$ is the operator defined in (3.5), $`U_y`$ has exactly the same form of $`U_z`$ with the operators $`P_{y+}^{(1)},P_y^{(1)},P_{y+}^{(2)},P_y^{(2)}`$ replacing the corresponding ones with the index $`z`$, and the operators $`P_L^{(2)},P_R^{(2)},P_L^{(4)},P_R^{(4)}`$ replacing the corresponding ones with the indices $`3`$ and $`6`$. Finally $`\stackrel{~}{U}_z`$ coincides with $`U_z`$ with the replacement of the indices 3 and 6 by the corresponding indices 3\* and 6\*.
To see how the mechanism works we have, first of all, to express the states (3.1) in terms of the corresponding states of the type $`|_{1y},_{2y}>`$ etc., in order to evaluate the effetct of applying the operator $`U_y.`$ Moreover, since now we are interested in identifying the eigenstates of $`T^2`$, it is useful to replace the states $`|_{1z,}_{2z}>`$ and $`|_{1z,}_{2z}>`$ by their simmetrical and skew-symmetrical combinations:
$`|Triplet>_z`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}[|_{1z},_{2z}>+|_{1z},_{2z}>],`$ (3.14)
$`|Singlet>_z`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}[|_{1z},_{2z}>|_{1z},_{2z}>]`$
and analogous ones for the states<sup>4</sup><sup>4</sup>4We remind that while no specification is necessary for the singlet state, in the case of the triplet state the specifications z or y are essential since they are genuinely different states. $`|_{1y,}_{2y}>`$ and $`|_{1y,}_{2y}>`$.
With these premises an elementary (but tedious) calculation shows the effect of applying the unitary operator $`U`$ to the isospin singlet and triplet states:
$$U|Singlet>|\mathrm{\Phi }>=|Singlet>|\mathrm{\Phi }>$$
(3.15)
$`U|_{1z},_{2z}>|\mathrm{\Phi }>=`$
$`+\{{\displaystyle \frac{1}{4}}|\mathrm{\Pi }(1,2)>_{3,6}[|\mathrm{\Pi }(1,2)>_{2,4}+|\mathrm{\Pi }(2,1)>_{2,4}+2|\mathrm{\Phi }>_{2,4}]|_{1z},_{2z}>`$
$`{\displaystyle \frac{1}{4}}|\mathrm{\Pi }(2,1)>_{3,6}[|\mathrm{\Pi }(1,2)>_{2,4}+|\mathrm{\Pi }(2,1)>_{2,4}2|\mathrm{\Phi }>_{2,4}]|_{1z},_{2z}>`$
$`+{\displaystyle \frac{i}{2\sqrt{2}}}|\mathrm{\Phi }>_{3,6}[|\mathrm{\Pi }(1,2)>_{2,4}|\mathrm{\Pi }(2,1)>_{2,4}]|Triplet>_z\}|\mathrm{\Pi }(1,2)>_{3,6};`$
$`U|_{1z},_{2z}>|\mathrm{\Phi }>=`$
$`\{{\displaystyle \frac{1}{4}}|\mathrm{\Pi }(1,2)>_{3,6}[|\mathrm{\Pi }(1,2)>_{2,4}+|\mathrm{\Pi }(2,1)>_{2,4}2|\mathrm{\Phi }>_{2,4}]|_{1z},_{2z}>`$
$`+{\displaystyle \frac{1}{4}}|\mathrm{\Pi }(2,1)>_{3,6}[|\mathrm{\Pi }(1,2)>_{2,4}+|\mathrm{\Pi }(2,1)>_{2,4}+2|\mathrm{\Phi }>_{2,4}]|_{1z},_{2z}>`$
$`{\displaystyle \frac{i}{2\sqrt{2}}}|\mathrm{\Phi }>_{3,6}[|\mathrm{\Pi }(1,2)>_{2,4}|\mathrm{\Pi }(2,1)>_{2,4}]|Triplet>_z\}|\mathrm{\Pi }(2,1)>_{3,6};`$
$`U|Triplet>_z|\mathrm{\Phi }>=`$
$`\{{\displaystyle \frac{i}{2\sqrt{2}}}|\mathrm{\Pi }(1,2)>_{3,6}[|\mathrm{\Pi }(1,2)>_{2,4}|\mathrm{\Pi }(2,1)>_{2,4}]|_{1z},_{2z}>`$
$`+{\displaystyle \frac{i}{2\sqrt{2}}}|\mathrm{\Pi }(2,1)>_{3,6}[|\mathrm{\Pi }(1,2)>_{2,4}|\mathrm{\Pi }(2,1)>_{2,4}]|_{1z},_{2z}>`$
$`+{\displaystyle \frac{1}{2}}|\mathrm{\Phi }>_{3,6}[|\mathrm{\Pi }(1,2)>_{2,4}+|\mathrm{\Pi }(2,1)>_{2,4}]|Triplet>_z\}|\mathrm{\Phi }>_{3,6}.`$
Suppose now that there are six local detectors (three at each wing) devised to measure the observables $`\mathrm{\Omega }^{(s)},`$ $`s=`$ 3,6,2\*,4\*,3\*,6\* and denote as $`\omega ^{(s)}`$ the obtained outcomes. The above relations exhibit some interesting features:
i). If $`\omega ^{(3)}+\omega ^{(6)}=\omega ^{(2)}+\omega ^{(4)}=\omega ^{(3)}+\omega ^{(6)}=0,`$ then reduction has taken place to $`|Singlet>.`$
ii). If at least one of the relations under i) is not satisfied, reduction has taken place to a state belonging to the three-dimensional $`T^2=2`$ eigenmanifold.
This second case can be further analyzed according to the following table:
$$\begin{array}{cccccc}\text{If}& \omega ^{(3)}+\omega ^{(6)}\hfill & =& 1\text{ or }2\hfill & \text{reduction has taken place to}& |_{1z},_{1z}>\hfill \\ \text{If}& \omega ^{(3)}+\omega ^{(6)}\hfill & =& 2\text{ or }1\hfill & \text{reduction has taken place to}& |_{1z},_{1z}>\hfill \\ \text{If}& \omega ^{(3)}+\omega ^{(6)}\hfill & =& 0\text{ and }\omega ^{(2)}+\omega ^{(4)}0\hfill & \text{reduction has taken place to}& |Triplet>_z.\hfill \end{array}$$
We can consider now an arbitrary state of the system, i.e. a linear superposition of the four states (3.1),
$$|\mathrm{\Psi }>=\alpha |_{1z},_{2z}>+\beta |_{1z},_{2z}>+\gamma |_{1z},_{2z}>+\delta |_{1z},_{2z}>$$
(3.19)
and use it to trigger our measuring devices. The resulting state $`U|\mathrm{\Psi }>|\mathrm{\Phi }>`$ will be the linear combination with the same coefficients of the states at the right hand side of eqs.(3.15-3.18). Let us write, for simplicity such a state as:
$$U|\mathrm{\Psi }>|\mathrm{\Phi }>=\eta |Singlet>|\mathrm{\Phi }>+|\mathrm{\Gamma }>.$$
(3.20)
As already remarked, if the detection procedure of the probe particles indicates that $`\omega ^{(3)}+\omega ^{(6)}=\omega ^{(2)}+\omega ^{(4)}=\omega ^{(3)}+\omega ^{(6)}=0`$ (and this occurs with probability $`|\eta |^2`$) then reduction takes place to the state $`|Singlet>.`$ In all other cases, even though reduction leads to a state of the eigemanifold $`T^2=2,`$ the measurement process does not respect the request of being “moral”, i.e., the condition that the reduced state be the (normalized) projection of the state prior to the measurement on the considered eigenmanifold. Said differently, the measurement, when it yields the outcome $`T^2=2,`$ turns out to be distorting for $`T_z`$. This is immediately seen by choosing in (3.19) $`\beta =\gamma =\delta =0`$ and by looking at eq.(3.16) which shows that there is a nonzero probability that the measurement leads to reduction on one of the states $`|_{z1},_{z2}>`$ and $`|Triplet>_z.`$ This is naturally related to the fact that, in a loose sense, one could state that in order to measure $`T^2`$ one has to measure simultaneously incompatible observables, such as $`T_{1z}`$ and $`T_{1x,}`$ of one of the particles. But for our purposes this nonideality of the measurement for the states of the eigenmanifold $`T^2=2,`$ is not relevant.
Actually, the feature of the measurement process we have just analyzed, besides being unavoidable, is an extremely positive one. In fact, if consideration is given to the situation discussed in Subsection 2.2, one easily sees that the distorting nature of the measurement eliminates the possibility of faster than light effects. From eq.(3.16) one sees that if a measurement of $`T^2`$ is performed on the state $`|_{1z},_{2z}>`$, the probability of getting in a subsequent measurement the outcome $`T_{2z}=+1`$ is no more equal to 1, but precisely to $`1/2`$. This agrees with the probability of getting such an outcome if one performs the measurement of $`T^2`$ on the state which has been subjected to a spin flip of particle 1, i.e. the state $`|_{1z},_{2z}>`$, as one derives immediately from a combined use of eqs.(3.15) and (3.18). On the contrary, eqs.(3.12) show that a local measurement of the nonlocal observable $`T_z`$ does not alter the fact that the probability of getting the outcome $`T_{2z}=+1`$ remains unaltered (and equal to 1) when the spin of particle 1 is flipped.
It is time to come to discuss the conceptual implications of the actual possibility of measuring nonlocal observables identified by Aharonov and Albert and reformulated by us in this Section.
## 4 The implications of the request of a covariant reduction
The analysis of ref. and of the previous section rules out the Hellwig and Kraus proposal for a covariant reduction mechanism and requires to consider an alternative approach to the problem. This has been suggested in another important paper by Aharonov and Albert. The proposal is quite simple and natural, but it has far reaching consequences. One could describe it by stating that, in a certain sense, they assume that reduction takes place instantaneously in all inertial frames. To be more clear, let us consider an objective space-time point $`P`$ where a local measurement occurs and the set of all space-like hyperplanes<sup>5</sup><sup>5</sup>5Actually the above mentioned authors have considered more generally and more appropriately arbitrary space like surfaces through $`P`$ and have related each of them to the inertial frame in which the normal to the surface at the considered point coincides with the time axis of the frame. This choice is not only useful but necessary in a perspective like the one of dynamical reduction models in which the dynamical equations, given the statevector on an initial space-like surface, determine (stochastically) the statevector to be attached to any space-like surface lying entirely in the future of the initial one. Here, however, to allow the reader to grasp intuitively the meaning of the proposal, we will resort to the consideration of space-like hyperplanes. through $`P.`$ Each such hyperplane is a $`t=const`$ surface for an appropriate observer. The proposal of ref. can then be formulated by stating that, for the considered reference frame, reduction takes place precisely along the $`t=const`$ hyperplane through $`P.`$
Note that this rule is manifestly covariant and it is quite natural, since the scope of statevector reduction is to allow one to exploit the information gained in a measurement for evaluating the probabilities of future (for him) events. However, taking such a position has far reaching consequences and requires a radical revision of the meaning of the wavefunction. Namely, the wavefunction cannot any longer be seen as a function on the space-time continuum but it becomes a function on the set of space-like hypersurfaces: the value of $`\mathrm{\Psi }`$ at a space-time point depends, in general, on the particular space-time surface crossing the point we take into account. We should then write, in place of $`\mathrm{\Psi }(x,t),`$ $`\mathrm{\Psi }=\mathrm{\Psi }(\sigma ,x),`$ where $`x`$ runs over the points of the hypersurface $`\sigma .`$
An intuitive understanding of this fact can be obtained by considering (Fig.3) an objective space-time point $`P`$ at which a local position measurement occurs for a particle in a state which is different from zero along two distinct world lines ($`l_1`$ and $`l_2)`$.
For a given reference frame, we can consider the set of $`t=const`$ hyperplanes. Then, for one of such hyperplanes $`\sigma (<)`$ associated to a time prior to the one characterizing the hyperplane through $`P`$, the wavefunction has a value different from zero at the point $`Q`$ in which $`\sigma (<)`$ intersects $`l_2:`$ $`\mathrm{\Psi }(\sigma (<),Q)0.`$ We can now consider another reference frame such that in it, the $`t=const`$ hyperplane $`\stackrel{~}{\sigma }(Q)`$ through $`Q`$ is characterized by a time label greater than the one characterizing the hyperplane $`\stackrel{~}{\sigma }(P)`$ of the same family going through $`P.`$ If we suppose that in the local measurement the particle is detected at $`P`$ (note that this is an objective, i.e., reference frame independent statement), then we have: $`\mathrm{\Psi }(\stackrel{~}{\sigma }(Q),Q)=0.`$ Thus, according to the Aharonov and Albert proposal, the wavefunction must take different values at the same space-time point according to which space-like surface passing through it we take into account.
We conclude this section by remarking that the idea that to deal with relativistic quantum systems one must associate different statevectors to different space-like hypersurfaces is a quite hold one, going back to the fundamental papers by Tomonaga and Schwinger . However, within the formalism proposed by these authors, the expectation value of any local observable having its compact support contained in the common part of two hypersurfaces (and thus in particular the value of the wavefunction at a point in which two such surfaces intersect) does not depend on which of the surfaces one is taking into account. This is not surprising: the Tomonaga-Schwinger theory is intended to describe (within a relativistic context) the linear and deterministic evolution of statevectors and does not pretend to account for nonlinear and stochastic processes like those we are considering here, i.e., the reduction processes.
## 5 Relativistic dynamical reduction theories
As repeatedly stressed in the previous sections, Aharonov and Albert, in their fundamental papers, have identified the crucial problems of relativistic reduction processes and have given clear hints about the features of any theory which should account for their occurrence. But they have not proposed any explicit and consistent dynamical mechanism specifying what occurs during a measurement process. Identifying such a new (universal) mechanism is precisely the aim of dynamical reduction theories.
### 5.1 The GRW theory
The first consistent and precise proposal of a dynamical theory accounting, at the nonrelativistic level, for the nonlinear and stochastic reduction process is the so called GRW theory . As pointed out by Bell , this approach corresponds to accepting that Schrödinger’s equation is not always right and taking into consideration stochastic and nonlinear modifications of it which allow a unified treatment of all natural processes including the objectification of macroscopic properties. The model is based on the assumption that besides the standard quantum evolution, physical systems are subjected to spontaneous localizations occurring at random times with an appropriate average frequency and affecting their elementary constituents. Such processes are formally described in the following way.
Let us consider a system of $`N`$ particles. When the $`i`$-th particle of the system suffers a localization, the wavefunction changes according to
$$\mathrm{\Psi }(𝐫_1,\mathrm{},𝐫_N)\mathrm{\Psi }_𝐱(𝐫_1,\mathrm{},𝐫_N)=\frac{\mathrm{\Phi }_𝐱(𝐫_1,\mathrm{},𝐫_N)}{\mathrm{\Phi }_𝐱}$$
(5.1)
$$\mathrm{\Phi }_𝐱(𝐫_1,\mathrm{},𝐫_N)=\left(\frac{\alpha }{\pi }\right)^{\frac{3}{4}}e^{\frac{\alpha }{2}(𝐫_i𝐱)^2}\mathrm{\Psi }(𝐫_1,\mathrm{},𝐫_N).$$
(5.2)
The probability density of the process occurring at point x is given by $`\mathrm{\Phi }_𝐱^2.`$ For what concerns the temporal aspects of the process we assume that the localizations of the various constituents (particles) occur independently at randomly distributed times with a mean frequency $`\lambda _m`$ which depends on their mass. We choose $`\lambda _m=\lambda (m/m_0)`$, where $`m`$ is the mass of the particle, $`m_0`$ the nucleon mass and $`\lambda `$ is of the order of $`10^{16}\mathrm{sec}^1.`$ The localization parameter $`1/\sqrt{\alpha }`$ is assumed to take the value $`10^5cm`$.
As the reader can easily grasp the model does not entail any appreciable deviation from standard quantum mechanics for microsystems since such systems suffer, on the average, one localization every $`10^9`$ years. The most appealing feature of the model derives from its trigger (or amplification) mechanism: in the case of a macroscopic system the frequency of the localizations is amplified with the number of its constituents and, in the case of an almost rigid body, each localization process amounts to a localization of the centre of mass, so that superpositions corresponding to different locations of a macroscopic object are suppressed in about $`10^7\mathrm{sec}`$.
The physically relevant features of the just described model for what concerns the reduction process (the one in which we are primarily interested in here) should be clear: if one wants to get information about an observable (a property) of a microsystem, one has to use it to trigger a macroscopic change. Different outcomes of the measurement are then correlated to different positions of some macroscopic system (to be precise, to macroscopically different mass distributions). But the theory does not tolerate the formation of superpositions of such macroscopically different states, leading, just as a consequence of the universal dynamics ruling all natural processes, to a definite outcome. Summarizing: the interaction between the measured system and the measuring apparatus strives to create superpositions of macroscopically different states, but the dynamics forbids such processes - measurements have outcomes in extremely short times.
### 5.2 The CSL theory
The model of the previous subsection, even though it contains all the essential elements allowing to overcome the problems affecting the quantum theory of measurement, has a drawback: it does not preserve the symmetry requirements of quantum mechanics for identical constituents. One could easily circumvent this difficulty by a theoretical scheme quite similar to the one presented above, but a more elegant (even though physically equivalent) formalism (CSL) based on a stochastic evolution equation for the statevector has been worked out , . Let us list its basic features. The evolution equation is:
$$\frac{d|\mathrm{\Psi }_w(t)}{dt}=\left[\frac{i}{\mathrm{}}H+\underset{i}{}A_iw_i(t)\gamma \underset{i}{}A_i^2\right]|\mathrm{\Psi }_w(t).$$
(5.3)
In equation (5.3) $`H`$ is the hamiltonian, the quantities $`A_i`$ are a set of commuting self-adjoint operators while $`w_i(t)`$ are c-number Gaussian stochastic processes satisfying:
$$w_i(t)=0,w_i(t)w_j(t^{})=\gamma \delta _{ij}\delta (tt^{}).$$
(5.4)
For the moment, let us assume that the operators $`A_i`$ have a purely discrete spectrum and let us denote as $`P_\sigma `$ the projection operators on their common eigenmanifolds. The precise way in which the model works is defined by the following prescription: if a homogeneous ensemble (pure case) is associated at the initial time $`t=0`$ to the statevector $`|\mathrm{\Psi }(0),`$ then the ensemble at time $`t`$ is the union of homogeneous ensembles associated with the normalized statevectors $`|\mathrm{\Psi }_w(t)/|\mathrm{\Psi }_w(t),`$ where $`|\mathrm{\Psi }_w(t)`$ is the solution of the evolution equation with the assigned initial condition and for the specific stochastic process $`w`$ which has occurred in the interval $`(0,t).`$ The probability density for such a subensemble is the cooked one, given by:
$$P_{cooked}[w]=P_{Raw}[w]|\mathrm{\Psi }_w(t)^2,$$
(5.5)
where
$$P_{raw}[w]=\frac{1}{N}e^{\frac{1}{2\gamma }_i_0^t𝑑\tau w_i^2(\tau )},$$
(5.6)
$`N`$ being a normalization factor. It is easy to show that the just described dynamical process, when the non-hamiltonian part dominates the hamiltonian one, drives each statevector within one of the eigenmanifolds associated to the projection operators $`P_\sigma ,`$ with the appropriate probabilities.
To get a theory which leads to the desired reductions (i.e. which works like the GRW theory) one has to make a choice for the operators $`A_i`$ which is directly suggested by the GRW theory itself. It is obtained by identifying the discrete index $`i`$ with the continuos index r and the above operators with an appropriately averaged mass density operator $`M(𝐫):`$
$$M(𝐫)=\underset{k}{}m^{(k)}N^{(k)}(𝐫),$$
(5.7)
$$N^{(k)}(𝐫)=\left[\frac{\alpha }{2\pi }\right]^{\frac{3}{2}}\underset{s}{}𝑑𝐪e^{\frac{\alpha }{2}(𝐪𝐫)^2}a_{(k)}^{}(𝐪,s)a_{(k)}(𝐪,s),$$
(5.8)
where $`a_{(k)}^{}(𝐪,s)`$ and $`a_{(k)}(𝐪,s)`$ are the creation and annihilation operators of a particle of type $`k`$ ($`k=`$ proton, neutron, electron,…) at a point q, with spin component $`s.`$ The equations replacing (5.3) and (5.4) are then:
$$\frac{d|\mathrm{\Psi }_w(t)>}{dt}=\left[\frac{i}{\mathrm{}}H+𝑑𝐫M(𝐫)\stackrel{~}{w}(𝐫,t)\frac{\gamma }{m_0^2}𝑑𝐫M^2(𝐫)\right]|\mathrm{\Psi }_w(t)>$$
(5.9)
and
$$\stackrel{~}{w}(𝐫,t)=0,\stackrel{~}{w}(𝐫,t)\stackrel{~}{w}(𝐫^{},t^{})=\frac{\gamma }{m_0^2}\delta (𝐫𝐫^{})\delta (tt^{}).$$
(5.10)
With these choices, when the parameter $`\gamma `$ is related to those of the GRW theory according to $`\gamma =\lambda \left(4\pi /\alpha \right)^{\frac{3}{2}}`$, the theory implies that any macroscopic object is always extremely well localized in space and that if any interaction leads, as a consequence of the hamiltonian dynamics, to a superposition of differently located states of a macroscopic system, then reduction takes place almost immediately to one of them, with the appropriate probability.
### 5.3 Relativistic CSL
The first attempt to get a relativistic generalization of the CSL theory has been performed by P. Pearle. A detailed investigation of the model proposed by him together with a discussion of all its most relevant features has been presented in ref.. In particular, in this last paper the relativistic stochastic invariance of the model, as well as the integrability of its evolution equation have been analyzed in all details. For the moment, let us simply recall the general scheme. One adopts a Tomonaga-Schwinger approach within a quantum field theoretical framework and one accounts for the stochasticity of the evolution by an appropriate stochastic interaction term. Accordingly, one considers a Lagrangian density
$$L(x)=L_0(x)+L_I(x)V(x)$$
(5.11)
where $`L_0(x)`$ and $`L_I(x)`$ are scalar functions of the fields, $`L_I(x)`$ does not contain derivative couplings, and $`V(x)`$ is a c-number stochastic process which is a scalar for Poincaré transformations. One chooses for $`V(x)`$ a Gaussian noise with mean zero. In order to have relativistic stochastic invariance, its covariance must be an invariant function:
$$V(x)V(x^{})=A(xx^{}),A(\mathrm{\Lambda }^1x)=A(x).$$
(5.12)
In refs. and the choice
$$A(x)=\lambda \delta (x)$$
(5.13)
has been made. It has to be stressed that such a choice, due to its white nature in time, gives rise to specific problems related to the appearence of untractable divergences. On the other hand, if one would not require the function $`A(x)`$ to be white in time various unacceptable consequences would emerge . Thus, we plainly accept that the program of a relativistic generalization of dynamical reduction theories is still an open one. However, the just mentioned problems are of technical nature, and one can hope to succeed in overcoming them. On the other hand, as we will show, the formal structure of the theory is perfectly satisfactory. Accordingly, we will disregard here this kind of difficulties and we will concentrate our attention on the general features of the theory.
We still have to define the precise dynamics of the model. This can be summarized in the following terms:
* The fields are solutions of the Heisemberg equations associated to $`L_0(x).`$
* The statevector obeys the evolution equation
$$\frac{\delta |\mathrm{\Psi }_V(\sigma )>}{\delta \sigma (x)}=\left[L_I(x)V(x)L_I^2(x)\right]|\mathrm{\Psi }_V(\sigma )>.$$
(5.14)
Note the skew-hermitian character of the coupling to the stochastic field. This equation, just as those of nonrelativistic CSL, does not preserve the norm of the statevector but it preserves its average value. Accordingly, one has to introduce an appropriate “cooking” procedure for the probability of occurrence of a given potential which parallels the one (5.5) of CSL.
In the specific case of refs. and consideration has been given to a hermitian scalar meson field $`\mathrm{\Phi }(x)`$ coupled to a fermion field $`\mathrm{\Psi }(x)`$, and the following choices have been made for the Lagrangian densities:
$$L_0(x)=\frac{1}{2}\left[_\mu \mathrm{\Phi }(x)^\mu \mathrm{\Phi }(x)m^2\mathrm{\Phi }^2(x)\right]+\overline{\mathrm{\Psi }}(x)\left[i\gamma ^\mu _\mu M\right]\mathrm{\Psi }(x)+\eta \overline{\mathrm{\Psi }}(x)\mathrm{\Psi }(x)\mathrm{\Phi }(x)$$
(5.15)
$$L_I(x)=\mathrm{\Phi }(x).$$
(5.16)
The physically relevant features of the model can be easily understood. The nonlinear and stochastic nature of eq.(5.14) implies that the dynamics leads to the suppression of superpositions of different states of the meson field. This, if one takes into account that fermions which are differently located are associated with different mesonic clouds, makes clear that the dynamics leads (indirectly) to a (very unfrequent) localization of the fermions. For lack of space we cannot analyze the theory in all its details and we refer the reader to ref. for an exhaustive discussion of all its features.
What matters for our analysis is just the characteristic of the theory of leading to the localization of fermions, in particular of nucleons. Actually, by appropriately choosing the constants of the model, one shows (ignoring the problem of the divergences) that in the nonrelativistic limit the theory exhibits features which are extremely similar to those of the CSL model.
Two remarks seem appropriate:
* Due to the nonhermitian structure of the evolution operator, the states associated to different hypersurfaces can take different values at a given objective space-time point (equivalently, the theory attaches different expectation values to operators having a common compact support). This fact parallels strictly the features which have been identified by Aharonov and Albert and which have been discussed in section 4.
* The theory, just as the GRW and CSL models, does not exhibit stochastic time reversal invariance: only the forward time translation Poincaré semigroup is represented. Accordingly, the initial conditions must be given on a precise, objective space-like hypersurface. The transformation from one observer to another must be discussed by adopting the passive point of view.
We conclude this subsection by discussing an explicit example of a measurement–like process. To this purpose it is sufficient to take into account once more a microscopic system possessing an internal degree of freedom which will be treated quantum mechanically, while one treats its space degrees of freedom as classical (in particular one disregards the spreading of wavepackets and one speaks of the world lines of the system). One also assumes that a macroscopic system enters into the game and that it also can be treated, for what concerns its spatial degrees of freedom, in classical terms. The world lines of the micro and the macroscopic systems are assumed to intersect at a precise objective space-time point. The macroscopic system mimics an apparatus measuring an observable of the internal Hilbert space of the microscopic system. The system-apparatus interaction is assumed to induce, according to the internal microstates of the system, different displacements of a macroscopic part (the pointer) of the apparatus. Such a mobile part of the device contains, obviously, an enormous number of nucleons (of the order of Avogadro’s number). They would end up in a superposition of states corresponding to different positions if no reducing dynamics would be effective. But, as we have made plausible, different locations of a nucleon are (very seldom) suppressed by the nonlinear evolution which does not tolerate superpositions of the associated different mesonic clouds. An amplification mechanism mirroring precisely the one of the GRW and CSL theories is present also in the relativistic model we are discussing. This means that the relativistic dynamics leads to the suppression of all but one the possible final macroscopically different configurations in an extremely short interval of the proper time of the apparatus.
We have described the process in physical terms. Let us now look at it in mathematical terms and in a language which does not mention different reference frames but only objective points or surfaces of the space-time continuum. We have a space-like surface $`\sigma _0`$ (to be identified with the one chosen for accounting of the big bang?) on which the initial conditions for the system and the apparatus are given. The dynamical equation determines in a unique way the statevector on any space-like surface lying entirely in the future of $`\sigma _0.`$ The previous discussion should have made clear that if we change the space-like surface we are interested in (e.g. by considering smooth continuous variations of it), it is just when we cross the point in which the system-apparatus interaction takes place that the nonlinear dynamics becomes effective and leads to a macroscopic definite state of the pointer (and, obviously, to a corresponding precise microstate of the system). Note that what matters for the reduction process is the fact that the considered space-like surface crosses an objective space-time point and not the precise way in which it is changed. If we cross the point by modifying the whole space-like surface we are considering (for instance by translating it rigidly in the direction of the time axis) or if we keep a far-away part of the surface fixed and we modify it only around the space-time point at which the system-apparatus interaction occurs, the change of the statevector is (objectively) the same<sup>6</sup><sup>6</sup>6Obviously, the same objective situation of the pointer as well as the observable which becomes definite for the system will be described in different terms by different observers. For instance if one observer adopts a certain reference frame he could claim that the pointer is aligned with his z-axis, while a rotated observer will claim that it is aligned with his x-axis. But this is totally irrelevant: we are not making reference to the language used by the observers, to the labels they attach to different space-time points and so on, we are specifying what happens in a covariant language, such as the one making reference to the objective point in which the system apparatus interaction takes place or to the fact that such a point belongs or does not belong to the volume between the initial space-like surface and the one we are interested in.: the pointer, in an extremely short time, acquires a precise final position differing from the initial ready to measure one, a position which is correlated to a precise final state of the measured system, i.e. the one corresponding to the measurement outcome.
At this point the reader should have perfectly clear that the basic features of the model under discussion are precisely those which have been identified in refs. and , in particular that different statevectors are attached to different hypersurfaces, that the value of the wavefunction can differ at the same space-time point according to the space-like hypersurface going through that point one takes into account, and that, in a quite precise sense, reduction occurs, for all observers, along the hypersurfaces going through the point at which the measurement process takes place.
In what follows we will reconsider the reduction process induced by the theory with reference to an oversimplified model of two correlated microsystems subjected to measurements in space-like separated regions, an analysis which will allow us to clarify some relevant features connected with quantum nonlocality and with the use of counterfactual arguments within a relativistic and nonlocal context. Now it is time to devote particular attention to the problem of the properties objectively possessed by individual physical systems within a theoretical framework like the one we have described in this section.
## 6 Properties and events in the relativistic context
We are ready to tackle a problem of great conceptual relevance, i.e., the one of the emergence of definite properties of an individual physical system from the sea of the potentialities which characterize it within the quantum framework. Obviously, this problem has extremely close connections with the measurement or the macro-objectification problem. For these reasons it is obvious that taking a precise position about the reduction process, such as the one entailed by relativistic CSL, puts precise limitations concerning the situations in which one can consistently speak of properties objectively possessed by an individual physical system at a give space-time point of its world line.
Before proceeding it is useful to reconsider shortly the same problem within nonrelativistic quantum mechanics. In such a context there is an absolute time, so that one can discuss the question of the properties objectively possessed by a physical system at a given instant $`t.`$ Suppose we have a system $`S`$ (elementary or composed of various constituents), whose state is represented, at the considered time, by the statevector $`|\mathrm{\Psi }(t)>.`$ The standard wisdom about properties (or in Einstein’s language about the elements of physical reality) leads to a quite natural assumption:
> If consideration is given to an observable whose associated self-adjoint operator is $`\mathrm{\Omega },`$ we claim that the events “the considered observable is definite” and “the system $`S`$ possesses the objective property $`\mathrm{\Omega }=\omega _k`$” occur iff $`|\mathrm{\Psi }(t)>`$ is an eigenstate of $`\mathrm{\Omega }`$ belonging to the indicated eigenvalue. If this is not the case, we claim that the event “the considered observable is indefinite” occurs.
Some remarks are appropriate.
* If one accepts, as it is usual within textbook quantum mechanics, that all (bounded) self-adjoint operators correspond to physically measurable quantities, then any system S, considered as a whole, possesses always some properties. In fact in any case there is at least one self-adjoint operator such that $`|\mathrm{\Psi }(t)>`$ is one of its eigenstates (the most trivial example being the projection operator $`|\mathrm{\Psi }(t)><\mathrm{\Psi }(t)|`$ itself).
* The above position concerning the possibility of speaking of objective properties matches the well known fact that the theory allows one to make the counterfactual statement: if at time $`t`$ a measurement process of $`\mathrm{\Omega }`$ were performed it would yield with certainity the outcome $`\omega _k.`$ For an analysis of this point, see the remarks below.
* If the assumption under the first item above is satisfied, then there are certainly physical observables which do not have a definite value, in the sense that the theory attaches genuinely nonepistemic nonvanishing probabilities to different outcomes in a prospective measurement of such observables at the considered time.
To analyze in greater details the just discussed situation let us consider the case of one particle which is in the (improper) superposition of two position eigenstates <sup>7</sup><sup>7</sup>7Obviously, to be correct, one should consider in place of the states $`|x_1>`$ and $`|x_2>`$ normalized wavefunctions different from zero only in extremely small intervals around the indicated positions and correspondingly a detector whose acceptance window is larger than the extension of the wavepacket around $`x_2`$.:
$$|\mathrm{\Psi }>=\frac{1}{\sqrt{2}}\left[|x_1>+|x_2>\right].$$
(6.1)
For such a state there is no matter of fact about the location of the particle: its position is indefinite. Suppose now that at time $`t`$ a detector is placed at point $`x_2`$ and it does not detect the particle. Immediately after $`t`$ the state, in accordance with the reduction postulate, becomes $`|x_1>`$ so that we can legitimately claim that the definite objective property of the particle being at $`x_1`$ has emerged as a consequence of the measurement at $`x_2.`$ Note that this statement is perfectly legitimate both from the point of view of the above criterion concerning possessed properties as well as from the one of counterfactual reasoning. In fact, the formalization of such reasonings requires to define the accessibility sphere from the actual world, and it is common practice in physics to consider as accessible (i.e. as those nearest to the actual one) those worlds in which the physical laws are the same as those of the actual world, and which coincide, up to the time one is interested in, with the actual world itself. Accordingly, in our case the accessibility sphere is represented by those worlds in which quantum mechanics holds and for which the state of the physical system we are interested in coincides up to a time following immediately $`t`$ with the one describing it in the actual world. In particular in all accessible worlds the premises: the particle has not been detected at $`x_2`$ and, accordingly, its state is $`|x_1>`$, are true.
As the reader should have clear, the situation changes radically in a relativistic context. In fact, as repeatedly stressed, the value of the wavefunction at the space-time point $`(x_1,t)`$ (i.e. in our case the fact that it differs from zero either only around point $`x_1`$ or around both points) depends not only on the considered (objective) space-time point, but also on the space-like surface through it which we take into account. How should one proceed to make statements about the property related to the location of the particle at a given time? The first important thing which has to be remarked is that in the spirit of an approach like the one of CSL bearing on reality as opposed to intersubjective appearences, the prescription which should lead to the legitimate conclusion that a property is definite and possessed or it is indefinite, must be covariant, i.e., independent from the reference frame one is considering. If one takes into account that the theory attaches a precise statevector to any space-like surface and that one is interested in a statement concerning an objective space-time point $`P`$ on the world line of the system we are considering, there are only two natural covariant prescriptions satisfying the above conditions, and they make reference to the past or to the future light cone from $`P`$, respectively. For reasons which are easily understandable (more about this in what follows) it is appropriate to resort to the criterion which takes into account the past light cone. Thus we replace the previous assumption about the attribution of objective properties to an individual physical system (at a point $`P`$ of its world line) by the following one:
> If consideration is given to a space-time point $`P`$ and to a given observable whose associated self-adjoint operator is $`\mathrm{\Omega },`$ we consider, first of all, the space-like surface $`\sigma (P)`$, constituted by the past light cone from $`P`$ and the part outside it of the initial surface $`\sigma _0`$. The theory assignes a precise statevector $`|\mathrm{\Psi }(\sigma (P))>`$ to such a surface. Then we proceed as before, i.e., we claim that the system $`S`$ possesses the objective property “$`\mathrm{\Omega }=\omega _k\mathrm{"}`$ iff $`|\mathrm{\Psi }(\sigma (P))>`$ is an eigenstate of $`\mathrm{\Omega }`$ belonging to the indicated eigenvalue, and, if this is not the case, we claim that the considered observable is indefinite.
Once more some remarks are appropriate:
* The above criterion implies that the statement “the particle has the definite property $`\mathrm{\Omega }=\omega _k`$” is correct when, in the past light cone from $`P`$ an appropriate preparation procedure or an interaction corresponding to the measurement of $`\mathrm{\Omega }`$ has occurred. Note that for an elementary particle this means that there is a point along its world line in which it has interacted with an appropriate device, while in more general cases like the one of a pair of far-away correlated particles in an entangled state, the property of one particle becomes definite also when in its past light cone there is a point in which a measurement of the relevant observable has been performed on the correlated particle.
* The above position concerning the possibility of speaking of objective properties matches the fact that the theory allows one to make the counterfactual statement “if at the space-time point $`P`$ a measurement process of the observable $`\mathrm{\Omega }`$ were performed it would yield the outcome $`\mathrm{\Omega }_k`$”, provided one makes the perfectly reasonable and covariant assumption that the accessible worlds from the actual one are those in which the laws of nature are the same as those of the actual world and the physical situation matches exactly the actual one within the past light cone from $`P`$.
## 7 Relativistic reduction and nonlocality
The fundamental issues discussed in the previous section become particularly interesting in connection with quantum nonlocality. To investigate this subtle point it is particularly useful to resort to an oversimplified model which exhibits all relevant features of relativistic CSL (in particular it is relativistically invariant, stochastic, nonlinear and nonlocal) but it is remarkably more simple. Such a model has been already considered in the previous sections, but, before proceeding, it is appropriate to make it absolutely precise.
### 7.1 Preliminary considerations: a relativistic toy model
We will deal with one or two particles, each having (as in the previous sections) space degrees of freedom obeying a classical relativistic dynamics, plus a quantum internal degree of freedom which behaves like a scalar under Lorentz transformations and which we will identify, for simplicity, with the isospin space of an isospin $`1/2`$ particle. We will consider an operator of this space with eigenvalues +1 and –1. When we will deal with two particles we will correspondingly consider two such operators , $`\mathrm{\Theta }^{(i)}`$ (i=1,2) and we will denote by $`|i,+>,`$ $`|i,>`$ the corresponding eigenvectors. Obviously, in the internal space of a particle one can consider any two by two hermitian matrix and in the case of two particles the full algebra of hermitian operators in the four dimensional internal space. However, for our purposes (as in the modern versions of the EPR argument) there will never be the need to resort to noncommuting observables referring to a particle, so that we will always deal with the operators $`\mathrm{\Theta }^{(i)}`$ .
In the theory, besides the two particles (1,2), there are objects (A,B, … ), simulating apparata measuring the observables $`\mathrm{\Theta }^{(i)}`$, which are characterized by space degrees of freedom obeying a classical relativistic dynamics and by three possible internal states, $`(r,+,)`$. Since they will always be in one of these three states it is not relevant, for the present analysis, to be precise about the nature of their internal space. In particular, one could consider the states $`|r>`$, $`|+>`$ and $`|>`$ as three ortogonal vectors in a three-dimensional Hilbert space or as three classical labels, which are Lorentz scalars. The objects, even though representing measuring devices, are supposed to be point-like, so that we can consider their world lines (which we will not draw in the figures) and the space-time points (which will be represented by small black squares) in which such world lines intersect the particles’ world lines. Moreover the objects (A,B, … ) are characterized by parameters $`g_A,g_B,`$ etc., which can take one of two values {0,1} (chosen at free-will by an experimenter) corresponding to the apparatus being “switched on” or “switched off” respectively.
### 7.2 The one-particle case
To warm up we begin by discussing the case of one particle. We consider its world line originating from the space-like surface $`\sigma _0`$ on which the initial conditions are given and, with reference to the internal degree of freedom, we assign the statevector on this surface by expressing it as a linear superposition of the eigenstates of the operator $`\mathrm{\Theta }`$ according to:
$$|\mathrm{\Psi }(\sigma _0)>=\alpha |+>+\beta |>.$$
(7.1)
\- The completeness assumption is embodied in the assertion that the assignement of the initial state (7.1) (besides the relativistic classical dynamics for the propagation of the free particle) represents the maximum of information we can have about the particle itself and determines all what we can know about it.
\- The experimental context: along the world line of the particle, at the space-time point $`R`$, there is an apparatus $`A`$ devised to measure $`\mathrm{\Theta }`$, which can be switched off or on.
\- Dynamics: it is nonlinear and stochastic. The theory associates to any space-like surface $`\sigma `$ in the future of $`\sigma _0`$ a statevector according to the following rules<sup>8</sup><sup>8</sup>8As already remarked, we will never consider other observables besides $`\mathrm{\Theta }`$. However, an exhaustive theory should deal with all hermitian operators in the internal space. The reader will have no difficulty in generalizing the rules to cover such a case. The whole procedure requires only to express the initial state as a linear combination of the eigenstates of the observables one is interested in.:
Denote by $`V(\sigma ,\sigma _0)`$ the space-time volume enclosed by the two indicated surfaces.
a). If
$$\left[RV(\sigma ,\sigma _0)\right]\left[g_A=0\right],$$
(7.2)
then the state $`|\mathrm{\Psi }(\sigma )>`$ associated to the surface is
$$|\mathrm{\Psi }(\sigma )>=|\mathrm{\Psi }(\sigma _0)>,$$
(7.3)
while
b). if
$$\left[RV(\sigma ,\sigma _0)\right]\left[g_A=1\right],$$
(7.4)
the state is
$$either|\mathrm{\Psi }(\sigma )>=|+>or|\mathrm{\Psi }(\sigma )>=|>,$$
(7.5)
the two alternatives occurring randomly with probabilities $`\left|\alpha \right|^2`$ and $`\left|\beta \right|^2`$ , respectively. Thus, when a spacelike surface crosses the region in which an apparatus is switched on a real dice-playing leading to one among two possible states takes place: the probabilities governing the process have a nonepistemic status.
In equations (7.1), (7.3) and (7.5) we have skipped the indication of the apparatus state but it is understood that the apparatus will be in one of the states (r,+,-) matching the one of the system.
A final comment is appropriate. The dynamics satisfies the consistency requirement that considering the evolution leading from $`\sigma _0`$ to $`\sigma `$ and then the one leading from $`\sigma `$ to $`\sigma _1`$ is the same as going directly from $`\sigma _0`$ to $`\sigma _1.`$ Thus, if one considers the case b) and assumes that along the particle’s world line there is another apparatus which is switched on and devised to measure the same observable at a point $`\stackrel{~}{R}`$ in the future of $`R`$, then the statevector to be assigned to surfaces $`\sigma `$ such that $`\stackrel{~}{R}V(\sigma ,\sigma _0)`$ would be the same as the one appearing in (7.5).
### 7.3 Events in the toy model: the one-particle case
In the one-particle case of subsection (7.2) let us consider a point P preceeding, on the world line of the particle, the point R. In such a case, since the precise rules of the model tell us that the state on the space-like surface $`\sigma (P)`$ is not an eigenstate of $`\mathrm{\Theta }`$ , the specific event “the associated property is indefinite” occurs.
Obviously, if consideration is given to a point $`\stackrel{~}{P}`$ followingR when $`g_A=1`$, then the surface $`\sigma (\stackrel{~}{P})`$ is such that for it condition (7.4) holds. Accordingly, as shown by Eqs.(7.5), the statevector is an eigenstate of $`\mathrm{\Theta }`$ (which one between $`|+>`$ and $`|>`$ is determined by the genuinely stochastic process taking place at $`R`$) and the the specific event “the property associated to $`\mathrm{\Theta }`$ is definite and equals + (-)” occurs.
The case under discussion, since it does not involve space-like separated events and nonlocal effects, does not raise any specific problem differing significantly from those which characterize also nonrelativistic quantum mechanics, i.e. the fact that one cannot avoid to take into account the event “there is no property referring to $`\mathrm{\Theta }`$ ”. The situation changes remarkably in the two particle case we are going to discuss.
### 7.4 The two-particle case
We start by generalizing the toy dynamics to the case in which there are two particles in place of one. First of all we have to specify the world lines describing the classical motion of the two particles and we must assign the statevector referring to the internal degrees of freedom on the initial space-like surface. To be general we should express it as an arbitrary normalized linear combination of the four ortonormal states $`|i,s>|j,t>`$ ($`i,j=1,2;s,t=+,`$) which are the common eigenvectors of the two commuting operators $`\mathrm{\Theta }^{(k)}`$. Similarly we should take into account the possibility of measuring, for each particle, operators which do not commute with those considered above. Moreover also the case of correlation measurements involving different operators for the two particles should be taken into account. However, for our purposes we can limit our considerations to a very specific initial state, i.e. to the state
$$|\mathrm{\Psi }(\sigma _0)>=\frac{1}{\sqrt{2}}\left[|1+,2>|1,2+>\right]$$
(7.6)
and to the operators $`\mathrm{\Theta }^{(i)}`$.
The reader could easily generalize our rules to arbitrary initial states and arbitrary measurement processes. Once more we specify the rules of the game (see Figs. 4 and 5):
\- Completeness: the assignement of the initial state (7.6) represents the maximum information one can have about the system.
\- The experimental context: along the world lines of the particles, at two space-time points R (at right for particle 2) and L (at left for particle 1) there are two apparatuses A and B (each characterized by three possible internal states r,+,-) devised to measure $`\mathrm{\Theta }^{(2)}`$ and $`\mathrm{\Theta }^{(1)},`$ respectively . Each apparatus can be switched on or off at the experimenter’s free will.
\- Dynamics: once more it is nonlinear and stochastic and associates to any space-like surface in the future of $`\sigma _0`$ a precise statevector according to rules which are the natural generalization of those of Section (7.2) (as before we denote by $`V(\sigma ,\sigma _0)`$ the space-time volume enclosed by the two indicated surfaces):
i. If
$$\left\{\left[RV(\sigma ,\sigma _0)\right]\left[g_A=0\right]\right\}\left\{\left[LV(\sigma ,\sigma _0)\right]\left[g_B=0\right]\right\},$$
(7.7)
then the state $`|\mathrm{\Psi }(\sigma )>`$ associated to the surface $`\sigma `$ is
$$|\mathrm{\Psi }(\sigma )>=|\mathrm{\Psi }(\sigma _0)>.$$
(7.8)
This situation occurs for the space-like surfaces represented by continuous lines in Fig.5.
ii. If
$$\left\{\left[RV(\sigma ,\sigma _0)\right]\left[g_A=0\right]\right\}\left\{\left[LV(\sigma ,\sigma _0)\right]\left[g_B=1\right]\right\},$$
(7.9)
then the state is
$$either|\mathrm{\Psi }(\sigma )>=|1+,2>or|\mathrm{\Psi }(\sigma )>=|1,2+>,$$
(7.10)
the two alternatives occurring at random with equal probabilities.
iii. If
$$\left\{\left[RV(\sigma ,\sigma _0)\right]\left[g_A=1\right]\right\}\left\{\left[LV(\sigma ,\sigma _0)\right]\left[g_B=0\right]\right\},$$
(7.11)
then the state is
$$either|\mathrm{\Psi }(\sigma )>=|1+,2>or|\mathrm{\Psi }(\sigma )>=|1,2+>,$$
(7.12)
the two alternatives occurring at random with equal probabilities. This situation occurs for the space-like surfaces represented by dashed lines in Fig.4.
iv. Finally if
$$\left\{\left[RV(\sigma ,\sigma _0)\right]\left[g_A=1\right]\right\}\left\{\left[LV(\sigma ,\sigma _0)\right]\left[g_B=1\right]\right\},$$
(7.13)
then the state is
$$either|\mathrm{\Psi }(\sigma )>=|1+,2>or|\mathrm{\Psi }(\sigma )>=|1,2+>,$$
(7.14)
the two alternatives occurring at random with equal probabilities. This case is represented by the gray dashed lines of Fig.5.
We have depicted in Fig. 4 the two alternatives corresponding to case iii, and in Fig. 5 those corresponding to case iv. It has to be stressed that the two occurrences in cases ii. and iii. (when only one apparatus is on) have no relations with the corresponding ones of case iv. In fact, leaving aside the case in which no measurement occurs, it has to be stressed that since in the actual world either one or both apparatuses are switched on, and the outcomes are genuinely stochastic, there is no definite relation between the two cases of Fig.4 and of Fig.5 . On the other hand, in case iv., if one considers a space-like surface, like the black dashed lines of Fig.5, passing below one of the two points where there is an apparatus and above the other one, and one supposes that one of the two alternatives has occurred, then the subsequent evolution must be consistent with the chosen alternative, i.e., for all space-like surfaces in the future of both A and B, the statevector of the system remains the same and the previously untriggered apparatus simply registers the property possessed by the microsystem. In this way the necessary requirement that in any case one can consistently describe the evolution from $`\sigma _0`$ to $`\sigma _1`$ and then the one from $`\sigma _1`$ to $`\sigma _2`$ is satisfied.
### 7.5 Some features of the two-particle model
The model we have just introduced has many interesting features. It contains precise dynamical rules for assigning to each space-like surface in the future of the surface defining the initial conditions a definite statevector. The model is fundamentally stochastic so that, when various alternatives can occur, they occur genuinely at random but in accordance with precise probabilistic laws. The macroscopic apparata are always in definite states (i.e. they always possess definite macroscopic properties) and, for those apparata which are switched on and for space-time points following (on their world lines) the event ”the microsystem triggers the apparatus”, they match the eigenvalues of the observables they are devised to measure. In all other instances, they correspond to their initial untriggered states.
The physical implications of the model obviously agree with those of SQM. In fact, since taken any objective space-time point (after the system-apparatus interaction) on the world line of an apparatus<sup>9</sup><sup>9</sup>9These world lines are not shown in the figures, but they can be simply thought as vertical lines in the reference frame in which the figures are drawn, corresponding to the fact that they are at rest in this frame. the apparatus state is precisely defined, we can make reference to these states to “read” the outcomes of the process.
Concerning its formal structure it has to be stressed that the model is entirely formulated in a coordinate-free language and thus it satisfies the relativistic requirements of a stochastically Lorentz invariant theory. In fact the statement that an objective space-time point (the one in which there is an apparatus which is on) belongs or does not belong to a precisely defined space-time volume is frame independent and the internal degrees of freedom are assumed to be Lorentz scalars. If we consider the correlations between outcomes, we see that when both apparatuses are switched on they register either (A+) (B-) or (A-) (B+) with equal probabilities and they never register the same outcome. Consequently the model reproduces the perfect correlations of SQM for isospin measurements along the same direction in the isospin singlet state.
The model satisfies the completeness requirement by assumption: there is no better specification of the initial state than the one given by $`|\mathrm{\Psi }(\sigma _0)>,`$ and its knowledge specifies everything about the future of the system exception made for the actual outcomes of processes whose probability of occurence is fundamentally nonepistemic.
Due to the fact that the model guarantees the perfect correlations of outcomes at the two wings of the apparatus it violates Bell’s locality requirement. It is quite important to stress that:
a). The model exhibits Parameter Independence. In fact, denoting by $`P_S(+1|g_R=\alpha )`$ the probability that the outcome at $`S`$ (taking the values $`L`$,$`R`$) be $`i`$ (taking the values +1,–1) when the apparatus at $`S`$ (taking the value $`R`$ when $`S=L`$ and $`L`$ when $`S=R`$) is switched off ($`g_S=0`$ ) or it is on ($`g_S=1`$ ), we have:
$`P_L(+1|g_R`$ $`=`$ $`1)=P_L(+1|g_R=0)={\displaystyle \frac{1}{2}}`$ (7.15)
$`P_L(1|g_R`$ $`=`$ $`1)=P_L(1|g_R=0)={\displaystyle \frac{1}{2}}`$
and, analogously:
$`P_R(+1|g_R`$ $`=`$ $`1)=P_R(+1|g_R=0)={\displaystyle \frac{1}{2}}`$ (7.16)
$`P_R(1|g_R`$ $`=`$ $`1)=P_R(1|g_R=0)={\displaystyle \frac{1}{2}}`$
b). The model violates Outcome Independence since the outcomes are perfectly correlated in spite of the fact that they have probability 1/2 of being +1 or -1.
### 7.6 Events in the two-particle case
At this point the reader should already have perfectly clear all the implications of the model. Suppose one is interested in an event concerning a space-time point $`P`$ of the world line of the i-th microcostituent of the composite system. The situation can be summarized as follows:
* No one of the space-time points $`R`$ and/or $`L`$ at which an apparatus is switched on belongs to the volume $`\stackrel{~}{V}(\sigma (P),\sigma _0)`$ lying between $`\sigma _0`$ and $`\sigma (P)`$ (this last surface being the one defined in section 6). Then the event “the observable $`\mathrm{\Theta }^{(i)}`$ is indefinite” is true.
* If any one of the space-time points $`R`$ and/or $`L`$ at which an apparatus is switched on belongs to $`\stackrel{~}{V}(\sigma (P),\sigma _0)`$ then the corrisponding event is “the microproperty related (or anticorrelated) to the outcome of the isospin component which has been measured” is definite. The probability of its value is precisely determined by the theory, the actual occurence of one of the possible outcomes is a genuinely random event.
Note that, in accordance with the above statements and when the apparatus at $`R`$ is on, the assertion “the observable $`\mathrm{\Theta }^{(1)}`$ of particle $`1`$ is indefinite” holds for all space-time points of the world line of microsystem 1 preceeding the point in which the future light cone from $`R`$ intersects its world line. The event “the property of microcostituent 1 is definite and it is opposite to the outcome of the measurement at $`R`$ on system 2” emerges when system 1 reaches the future light cone of the measurement event.
Suppose now one is interested in the event characterizing a precise space-time point of the world line of a macroscopic measuring apparatus. As already remarked and as it should be evident by our argument the event for such a system is always precisely defined and it corresponds to one of the alternatives “the pointer points to the $`r`$ (ready) position, it points to +, it points to -”. It is important to stress that this holds for both world lines of the apparata independently of the fact that they are switched on or off (in which case they are always in the $`r`$ state) and independently of the fact that only one or both of them are switched on. In spite of this fact, in the case in which both apparatuses are on, the “definite events” referring to space-time points following, on their world lines, the objective space-time points at which the system-apparatus interactions take place are always opposite, i.e., the perfect anticorrelations characterizing SQM predictions are respected.
### 7.7 Counterfactuals and nonlocality
We consider it appropriate to call the attention of the reader on the extremely relevant implications of the nonlocal nature of quantum theory for counterfactual arguments within a relativistic context. We recall, first of all, that we have related the possibility of making counterfactual assertions about an objective space-time point to the consideration of the past light cone from the considered point. This is unavoidable within a context like the present one in which the dynamics is fundamentally irreversible, so that the absolute past plays a basic role for any consideration concerning the absolute future.
In this subsection we want to analyze in greater details the problems which are specifically related to nolocality, to have the opportunity of stressing some subtle points. We start by considering a quite simple objection to our way of dealing with counterfactual assertions which could be raised by a naive reader: in the two particle case of subsections 7.4 and 7.6, why an observer who is on the world line of particle 2 at a point Q in the immediate future of the point $`R`$ at which the isospin component $`T_{2z}`$ of this particle has been measured (and found, e.g. to have the value –1) is not allowed to make the statement “if an apparatus were switched on on the world line of particle 1 at a point L which is space-like with respect to both $`R`$ and $`Q,`$ such an apparatus would register with certainity the outcome +1”? Here is where nonlocality enters. To claim that the above statement is appropriate means to assume that the accessibility sphere from the actual world is represented by all those worlds in which the antecedent, i.e., the fact that the outcome at R has been –1, is true. If the theory were local, i.e. if the outcome at a given point were totally independent from all what is going on at space-like separations, then the argument would be perfectly correct. But, as we know, this is not the case. To allow the reader to grasp this subtle point as well as the argument we will present below, we invite him to consider the three following situations and the related statements:
* We consider an actual world in which everything is like in the two particle case of subsection 7.4 and, moreover, both apparatuses are on. We also consider an observer along the world line of particle 2 immediately after the point $`R,`$ who is aware of the outcome of the measurement and is also aware of the fact that the apparatus at $`L`$ is on. Then he can argue in the following way: I know that the apparatus $`A`$ (the one at $`R`$) has registered “the isospin has the value –1”, I also know that another apparatus is on at $`L`$ and that the theory guarantees that the final outcomes are always anticorrelated. I can then claim that the apparatus $`B`$ at $`L`$, at any instant following (on its world line) the one of its interaction with particle 1, registers the value +1.
Note that the above is not a counterfactual argument since it makes exclusive reference to the actual world. Moreover, as it should be clear from our analysis in the previous sections, such an argument is perfectly legitimate and correct.
* We consider now the same situation, but we assume that, in the actual world, the apparatus at $`L`$ is off (note that this is a statement about a precise space-time event, meaning that the apparatus is off when the particle might trigger it). Let us consider once more our observer and his reasoning: I know that the apparatus $`A`$ at $`R`$ has registered “the isospin has the value –1” and I also know that the theory guarantees that if two apparata were present and switched on at L and R their outcomes would be perfectly anticorrelated. I can then claim that if the apparatus at $`L,`$ were on, at any instant following (on its world line) the one of its interaction with particle 1, it would register the value +1.
This, as the reader has certainly clear, is a genuinely counterfactual argument and, according to our criterion (but as we will show, completely in general, due to the nonlocal character of the theory) is definitely illegitimate.
* The final case is exactly the same as the previous one, and, in particular, the apparatus at $`L`$ is switched off. The only difference consists in the statement made by the observer. Such statement does not refer to a point which is space-like with respect to him, but to one, let us say $`T,`$ along the world line of particle 1 lying in the future light cone of the observer. The he can claim: I know that the apparatus at $`R`$ has registered “the isospin has the value –1” and I also know that the theory guarantees that the outcomes of pairs of apparata at R and T (which are time-like separated) are always anticorrelated. So, in spite of the fact that in the actual world no apparatus is on along the world line of particle 1, I can claim (which, for time–like separations also means predict) that if such an apparatus were on at $`T,`$ then it would register the outcome “the isospin equals $`+1`$”.
We stress that this is a genuinely counterfactual argument and we also stress that, according to our criterion and to standard wisdom, is a perfectly legitimate one.
To further clarify why the position analyzed in the second of the above cases is inappropriate and to show that the fundamental reason for this derives from the basically nonlocal nature of physical processes which has been so appropriately brought to the attention of the scientific community by the analysis of J.S. Bell , we will consider now an hypothetical (and certainly possible) deterministic completion of our toy model, i.e. a nonlocal deterministic hidden variable theory equivalent to it. The theory will be characterized by hidden variables which can take values from a set $`\mathrm{\Lambda }`$ and whose knowledge would unambiguously determine all outcomes of all conceivable measurement processes.
Within the set $`\mathrm{\Lambda }`$ we identify (Fig.6) two subsets $`\mathrm{\Lambda }_1(2,+)`$ and $`\mathrm{\Lambda }_1(2,)`$ such that in the case in which only the apparatus at $`A`$ is switched on, if the actual value $`\lambda `$ of the hidden variable belongs to $`\mathrm{\Lambda }_1(2,+)`$ \[$`\mathrm{\Lambda }_1(2,)]`$ then the outcome of the measurement is $`+1`$ ($`1).`$ Let us now consider the case in which both apparatuses are switched on. Then the fact that locality is violated implies that there exist a non empty subset $`\mathrm{\Lambda }_{12}`$ of $`\mathrm{\Lambda }_1(2,+)`$ such that, for $`\lambda \mathrm{\Lambda }_{12}`$ the outcome at $`A`$ is –1 and the one at $`B`$ is +1 <sup>10</sup><sup>10</sup>10Actually nonlocality implies that this should occur for at least some of the pairs of perfectly correlated observables of the constituents. For simplicity we assume here that this actually happens for the pair we are interested in. This does not change in any way the conceptula implications of our analysis..
To judge of the validity of a counterfactual statement like “if besides the apparatus at $`R`$ also the apparatus at $`L`$ were on then the outcome at $`L`$ would be …” one has to identify the accessibility sphere from the actual world. Which is the appropriate criterion to characterize the worlds which are nearest to the actual one? If one claimed that they are those characterized by the same value of the hidden variables, then for the subset $`\lambda \mathrm{\Lambda }_{12}`$ of $`\mathrm{\Lambda }_1(2,+)`$ the appropriate specification which must replace the dots in the previous sentence would be “$`+1`$” (while the outcome at $`R`$ should be characterized by the value opposite to the one which occurs with certainty when only such apparatus is on). On the contrary, if one took the position that the worlds which are nearest to the actual one are those in which the outcome at $`R`$ is the same as the one obtained in the actual world, then one would be including alternative worlds characterized by a value of the hidden variables belonging to an appropriate subset of $`\mathrm{\Lambda }_1(2,)`$, i.e. worlds such that if only the apparatus at $`R`$ were on would give the outcome opposite to the one which has occurred in the actual world.
This elementary example should have made clear how delicate is the handling of counterfactual statements in a relativistic nonlocal context and why the prescription we have adopted is the only consistent and logical one. For a further discussion of this delicate point we refer the reader to ref..
## 8 Relativistic dynamical reduction and local measurements of nonlocal observables
In this section, with reference to our toy model, we will reconsider the experimental set-up of section 3 devised to measure nonlocal observables of a two particle system by resorting to local interactions and detections. The first important remark is that the dynamical evolution has to be enriched (with respect to the examples analyzed in the previous section) to take into account the further interaction processes between the microprobes (, \[2\*\], \[4\*\] and , \[4\*\], \[6\*\]) and the microscopic constituents ((1) and (2)). Just because these interactions involve only microscopic systems they do not give rise to reduction processes and are accounted by the linear and deterministic evolution summarized in eqs. (3.15)-(3.18). Obviously, also such interactions occur at precise, objective space-time points, so that the whole processs is perfectly covariant.
The first appropriate step is to analyze the evolution from the initial space-like surface $`\sigma _0`$ to a surface $`\sigma _1`$ which has ”crossed” a region in which the interactions between the probes and the particles occur (see Fig. 7, where the considered region is the one at right, lying on the world line of particle 2).
We still assume that the initial state is the singlet isospin state. In order to evaluate the evolution due to the interactions of particle 2 with the probe particles , \[4\*\] and \[6\*\] we remark that such an evolution is governed by the appropriate part of the operator $`U=\stackrel{~}{U_z}U_yU_z`$ defined after eq.(3.13). Such an operator can be written as the direct product of an operator acting on particle 2 and one acting on particle 1 (this expresses the local nature of the interactions with the probes):
$`U=U^{(1)}U^{(2)},`$ (8.1)
$`U^{(2)}=[P_{z+}^{(2)}P_L^{(6)}+P_z^{(2)}P_R^{(6)}][P_{y+}^{(2)}P_L^{(4^{})}+P_y^{(2)}P_R^{(4^{})}][P_{z+}^{(2)}P_R^{(6^{})}+P_z^{(2)}P_L^{(6^{})}].`$
According to eq. (8.1), the effect of applying $`U^{(2)}`$ to the singlet state is (as one proves by a rather involved calculation):
$`|\mathrm{\Psi }(\sigma _1)>U^{(2)}|Singlet>|\varphi >`$ (8.2)
$`={\displaystyle \frac{i}{2\sqrt{2}}}(P_R^{(6)}|\varphi >_{3,6})[(P_R^{(4^{})}P_L^{(4^{})})|\varphi >_{2^{},4^{}}](P_L^{(6^{})}|\varphi >_{3^{},6^{}})|_{1z},_{2z}>`$
$`+{\displaystyle \frac{1}{2\sqrt{2}}}(P_R^{(6)}|\varphi >_{3,6})[(P_R^{(4^{})}+P_L^{(4^{})})|\varphi >_{2^{},4^{}}](P_R^{(6^{})}|\varphi >_{3^{},6^{}})|_{1z},_{2z}>`$
$`{\displaystyle \frac{1}{2\sqrt{2}}}(P_L^{(6)}|\varphi >_{3,6})[(P_L^{(4^{})}+P_R^{(4^{})})|\varphi >_{2^{},4^{}}](P_L^{(6^{})}|\varphi >_{3^{},6^{}})|_{1z},_{2z}>`$
$`{\displaystyle \frac{i}{2\sqrt{2}}}(P_L^{(6)}|\varphi >_{3,6})[(P_L^{(4^{})}P_R^{(4^{})})|\varphi >_{2^{},4^{}}](P_R^{(6^{})}|\varphi >_{3^{},6^{}})|_{1z},_{2z}>.`$
Recalling the effect of the operators $`P_L^{(i)}`$ and $`P_R^{(i)}`$ on the states $`|\varphi >_{r,i}:`$
$$P_L^{(i)}|\varphi >_{r,i}=\frac{1}{\sqrt{3}}\left[|0,1>_{r,i}+|1,1>_{r,i}+|1,0>_{r,i}\right]$$
(8.3)
$$P_R^{(i)}|\varphi >_{r,i}=\frac{1}{\sqrt{3}}\left[|0,1>_{r,i}+|1,0>_{r,i}+|1,1>_{r,i}\right]$$
we can rewrite the above state (8.2) as:
$`|\mathrm{\Psi }(\sigma _1)U^{(2)}|Singlet>|\varphi >=`$ (8.4)
$`{\displaystyle \frac{1}{6\sqrt{6}}}[i|A>|_{1z},_{2z}>+|B>|_{1z},_{2z}>|C>|_{1z},_{2z}>i|D>|_{1z},_{2z}>],`$
where the probe states $`|A>,`$ $`|B>,`$ $`|C>`$ and $`|D>`$ are given by:
$`|A>`$ $`=`$ $`(|0,1>_{3,6}+|1,0>_{3,6}+|1,1>_{3,6})`$
$`[|0,1>_{2^{},4^{}}+|1,0>_{2^{},4^{}}+|1,1>_{2^{},4^{}}|0,1>_{2^{},4^{}}|1,1>_{2^{},4^{}}|1,0>_{2^{},4^{}}]`$
$`\left(|0,1>_{3^{},6^{}}+|1,1>_{3^{},6^{}}+|1,0>_{3^{},6^{}}\right),`$
$`|B>`$ $`=`$ $`(|0,1>_{3,6}+|1,0>_{3,6}+|1,1>_{3,6})`$
$`[|0,1>_{2^{},4^{}}+|1,0>_{2^{},4^{}}+|1,1>_{2^{},4^{}}+|0,1>_{2^{},4^{}}+|1,1>_{2^{},4^{}}+|1,0>_{2^{},4^{}}]`$
$`\left(|0,1>_{3^{},6^{}}+|1,0>_{3^{},6^{}}+|1,1>_{3^{},6^{}}\right),`$
$`|C>`$ $`=`$ $`(|0,1>_{3,6}+|1,1>_{3,6}+|1,0>_{3,6})`$
$`[|0,1>_{2^{},4^{}}+|1,0>_{2^{},4^{}}+|1,1>_{2^{},4^{}}+|0,1>_{2^{},4^{}}+|1,1>_{2^{},4^{}}+|1,0>_{2^{},4^{}}]`$
$`\left(|0,1>_{3^{},6^{}}+|1,1>_{3^{},6^{}}+|1,0>_{3^{},6^{}}\right),`$
$`|D>`$ $`=`$ $`(|0,1>_{3,6}+|1,1>_{3,6}+|1,0>_{3,6})`$
$`[|0,1>_{2^{},4^{}}+|1,1>_{2^{},4^{}}+|1,0>_{2^{},4^{}}|0,1>_{2^{},4^{}}|1,0>_{2^{},4^{}}|1,1>_{2^{},4^{}}]`$
$`\left(|0,1>_{3^{},6^{}}+|1,0>_{3^{},6^{}}+|1,1>_{3^{},6^{}}\right).`$
We consider now the further evolution from $`\sigma _1`$ to $`\sigma _2.`$ In this process a micro-macro interaction between the probe particles and the detectors $`\mathrm{\Omega }^{[6]},`$ $`\mathrm{\Omega }^{[4^{}]}`$ and $`\mathrm{\Omega }^{[6^{}]}`$ takes place and leads almost immediately to precise locations of their macroscopic pointers. From the above equations we see that all possible triplets of oucomes can occur. Just to fix our mind, let us suppose that reduction has lead to the three outcomes $`\mathrm{\Omega }^{[6]}=0,`$ $`\mathrm{\Omega }^{[4^{}]}=0`$ and $`\mathrm{\Omega }^{[6^{}]}=0.`$ The state after the completion of this process is:
$`|\mathrm{\Psi }(\sigma _2)>=`$ (8.9)
$`{\displaystyle \frac{N}{6\sqrt{6}}}\{i[|1,0>_{3,6}(|1,0>_{2^{},4^{}}|1,0>_{2^{},4^{}})|1,0>_{3^{},6^{}}]|_{1z},_{2z}>`$
$`+[|1,0>_{3,6}(|1,0>_{2^{},4^{}}+|1,0>_{2^{},4^{}})|1,0>_{3^{},6^{}}]|_{1z},_{2z}>`$
$`[|1,0>_{3,6}(|1,0>_{2^{},4^{}}+|1,0>_{2^{},4^{}})|1,0>_{3^{},6^{}}]|_{1z},_{2z}>`$
$`i[|1,0>_{3,6}(|1,0>_{2^{},4^{}}|1,0>_{2^{},4^{}})|1,0>_{3^{},6^{}}]|_{1z},_{2z}>,`$
$`N`$ being a normalization factor.
And now comes the final step. We consider the subsequent evolution to a final state which is associated to a space-like surface in the future of $`\sigma _2,`$ obtained by crossing the region of Fig. 7 where the interactions of particle 1 with the probe particles takes place (subsequently the detectors at left will simply register the definite situation of the probes). Such evolution is governed by the unitary operator $`U^{(1)}`$ which has the same form as $`U^{(2)}`$ of eq. (8.1) with the apices , \[2\*\] and \[3\*\] replacing , \[4\*\] and \[6\*\]. With a rather tedious calculations one gets the result:
$$|\mathrm{\Psi }(\sigma _{final})>=\left[|0,0>_{3,6}|0,0>|0,0>_{3^{},6^{}}\right]|Singlet>.$$
(8.10)
Obviously, one can perform the calculation by assuming any other possible set of outcomes of the detectors at the right wing. For instance, if one chooses the set $`\mathrm{\Omega }^{[6]}=1,`$ $`\mathrm{\Omega }^{[4^{}]}=0`$ and $`\mathrm{\Omega }^{[6^{}]}=0`$ and goes through precisely the same calculation, one gets, in place of eq. (8.10) the following one:
$$|\mathrm{\Psi }(\sigma _{final})>=\left[|1,1>_{3,6}|0,0>|0,0>_{3^{},6^{}}\right]|Singlet>.$$
(8.11)
As expected, the final state of the system is always the singlet and the states of the various pairs of apparata appearing in it are such that the sum of their eigenvalues is always zero, as it must be in order that the formalism be internally consistent. The same happens for different initial states and for any possible choice of the outcomes at the right wing. We believe that this simple exercise should have made clear to the reader the elegance and the logical coherence of the formalism.
Obviously, mastering completely the theoretical scheme allows us to tackle, also in the more general case of nonlocal measurements, the problem of property attribution to the system. Since such measurements involve a pair of space-time points one has to identify the space-like hypersurface which has to be considered for drawing the relevant conclusions. According to our criterion, such a hypersurface is the boundary of the two past light cones from the pair of points one is taking into account on the world lines of the two particles. It is then obvious that, if consideration is given to a pair of points which preceed the two interaction regions (both at right and at left) of the particles with the probes, the theory associates to such a surface the singlet state and one can claim that “the composite system possesses the objective property of having total isospin equal to 0”. Such an assertion is confirmed by the subsequent measurement procedure and the communication between the two obserbers concerning the outcomes they have registered. On the contrary, if one considers a pair of points such that one follows and the other one preceeds the particle-probe interactions, i.e. two points like those in which the hypersurface $`\sigma _1`$ of Fig. 7 crosses the two world lines, the state describing the overall situation is the one of eq.(8.2) for which the square of the total isospin is undefined. One can go on reasoning in the just outlined way. We have no space to discuss all aspects of such a procedure in detail, we simply point out that, with the chosen criterion the unfolding of the process and the statements one is led to make are quite natural.
The last comment which is worth making concerns the other possible (in principle) choice for the space-like surface one uses to identify the possessed properties. As already remarked, one could make reference to the future rather than to the past light cone without meeting any logical contradiction. However, we cannot avoid stressing that the ensuing picture is rather peculiar. The reasons should be quite obvious: in order to speak of an objective situation at an objective space-time point it is much more natural to make reference to its absolute past than to the past plus the whole region which is space-like with respect to it. In particular, we point out that when the future light cone criterion is adopted, in the situation of section 7.4 and in the case in which only the apparatus at $`R`$ is switched on, the isospin of particle 1 becomes definite at the space-time point at which the past light cone from $`R`$ intersects its world line. In our opinion, this sort of (luminal) backward causation is less natural than the forward one implied by the alternative criterion. Even more peculiar aspects emerge when one takes into account nonlocal measurements. These reasons make our choice, whith its precise physical implications, definitely preferable.
## 9 Conclusions
The detailed analysis we have performed should, in our opinion, have made clear the conceptually most relevant features that any relativistic model of dynamical reduction must satisfy. It should also have made clear that the relativistic generalizations of CSL satisfy, in principle, all requirements we have identified. As already mentioned, such models encounter some technical difficulties. But we believe that the present analysis makes plausible that the main physical ideas at the basis of the considered approaches must be essential ingredients of any attempt to account for the macro-objectification process respecting the relativistic requirements and emboding the fundamental nonlocal aspects of natural processes.
## Acknowledgements
We thank Dr. F. Camana for a careful checking of some of the cumbersome calculations and Drs. A. Bassi and L. Marinatto for stimulating exchanges of views and for help in preparing the LaTex version of the manuscript.
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# POLARIZED MØLLER SCATTERING ASYMMETRIES**footnote *Talk given at 𝑒⁻𝑒⁻99: 3rd International Workshop on Electron-Electron Interactions at TeV Energies, Santa Cruz, December 1999.
## 1 Polarization and Precision Measurements
Polarized beams provide powerful tools for testing the Standard Model and probing “new physics” effects. They can be used to enhance signals, suppress backgrounds, study particle properties, and carry out precision measurements. A beautiful illustration of the last possibility is provided by the SLD measurement of $`A_{LR}`$ at the $`Z`$ pole
$`A_{\mathrm{LR}}{\displaystyle \frac{\sigma \left(e^+e_L^{}\mathrm{hadrons}\right)\sigma \left(e^+e_R^{}\mathrm{hadrons}\right)}{\sigma \left(e^+e_L^{}\mathrm{hadrons}\right)+\sigma \left(e^+e_R^{}\mathrm{hadrons}\right)}}.`$ (1)
That quantity is very sensitive to $`\mathrm{sin}^2\theta _W`$
$`A_{\mathrm{LR}}={\displaystyle \frac{2(14\mathrm{sin}^2\theta _W)}{1+(14\mathrm{sin}^2\theta _W)^2}}\text{(Tree level)}.`$ (2)
In fact, for $`\mathrm{sin}^2\theta _W0.23`$, one finds $`\mathrm{\Delta }\mathrm{sin}^2\theta _W/\mathrm{sin}^2\theta _W\frac{1}{10}\mathrm{\Delta }A_{\mathrm{LR}}/A_{\mathrm{LR}}`$. Hence, a $`\pm 1\%`$ measurement of $`A_{\mathrm{LR}}`$ determines $`\mathrm{sin}^2\theta _W`$ at the $`\pm 0.1\%`$ level.
Based on about 500 thousand $`Z`$ decays and employing a polarized $`e^{}`$ beam with polarization reaching $`P_e^{}77\%`$, the SLD collaboration has reported<sup>?</sup> the single best measurement of the weak mixing angle (defined here by modified minimal subtraction)
$`\mathrm{sin}^2\theta _W(m_Z)_{\overline{MS}}=0.23073\pm 0.00028,`$ (3)
which weighs heavily in the (leptonic) $`Z`$ pole average (from SLD and LEP)
$`\mathrm{sin}^2\theta _W(m_Z)_{\overline{MS}}=0.23091\pm 0.00021.`$ (4)
Taken on their own, the quantities in eqs. (3) and (4) are merely precise numbers. They become interesting when interpreted in the context of a complete (renormalizable) theory such as the $`SU(3)_C\times SU(2)_L\times U(1)_Y`$ Standard Model or its various extensions. Then, symmetries provide natural relationships among couplings and masses which can be tested by comparing different precision measurements. For example, the fine structure constant, Fermi constant, and $`Z`$ mass
$`\alpha ^1`$ $`=`$ $`137.03599959(40)`$
$`G_\mu `$ $`=`$ $`1.16637(1)\times 10^5\text{GeV}^2`$
$`m_Z`$ $`=`$ $`91.1871(21)\text{GeV}`$ (5)
can be compared with the weak mixing angle via
$`\mathrm{sin}^22\theta _W(m_Z)_{\overline{MS}}={\displaystyle \frac{4\pi \alpha }{\sqrt{2}G_\mu m_Z^2\left[1\mathrm{\Delta }\widehat{r}(m_t,m_h)\right]}}`$ (6)
where $`\mathrm{\Delta }\widehat{r}`$ represents finite, calculable quantum loop effects which depend on the top quark and Higgs scalar masses. Taking $`m_t=174.3\pm 5.1`$ GeV and $`m_h100`$ GeV leads to $`\mathrm{\Delta }\widehat{r}=0.05940\pm 0.0005\pm 0.0002`$, where the errors correspond to $`\mathrm{\Delta }m_t`$ and hadronic loop uncertainties.
Leaving $`m_h`$ arbitrary, eq. (6) leads to the prediction<sup>?</sup>
$`\mathrm{sin}^2\theta _W(m_Z)_{\overline{MS}}=(0.23112\pm 0.00016\pm 0.00006)\left(1+0.00226\mathrm{ln}{\displaystyle \frac{m_h}{100\text{GeV}}}\right).`$ (7)
Comparing that prediction with the world average in eq. (4) suggests a relatively light Higgs,
$`m_h65\begin{array}{ccc}\hfill +35& \hfill +28& \hfill +9\\ \hfill 20& \hfill 21& \hfill 8\end{array}\text{GeV},`$ (10)
which is centered somewhat below the LEP II direct search bound<sup>?</sup>
$`m_h>106\mathrm{GeV}\text{(95\% C.L.)}.`$ (11)
In fact, the SLD value in eq. (3) favors an even smaller $`m_h`$. If the Higgs mass turns out to be well outside the range in eq. (10), then one must append “new physics” to the Standard Model either through loop effects or small tree level contributions.
It would be nice to push the current $`\pm 0.1\%`$ test in eq. (6) as far as possible. Indeed, $`\alpha `$, $`G_\mu `$, and $`m_Z`$ are all already known to much better than $`\pm 0.01\%`$ (and will be or can be further improved). Can one reduce the uncertainty in $`\mathrm{sin}^2\theta _W(m_Z)_{\overline{\mathrm{MS}}}`$ from its current $`\pm 0.1\%`$ to $`\pm 0.01\%`$? If so, it would provide a sensitivity to $`m_h`$ at the incredible $`\pm 5\%`$ level (assuming $`m_t`$ and hadronic loop uncertainties are also improved).
The only known way to improve $`\mathrm{sin}^2\theta _W(m_Z)_{\overline{\mathrm{MS}}}`$ is to carry out a clean high statistics study of asymmetries such as $`A_{\mathrm{LR}}`$. In that regard, the NLC (Next Linear Collider) will be capable at an early stage of sitting at the $`Z`$ resonance and collecting $`10^810^9`$ $`Z`$ decays in a relatively short time. With such statistics, $`\mathrm{sin}^2\theta _W(m_Z)_{\overline{\mathrm{MS}}}`$ can, in principle, be obtained via $`A_{\mathrm{LR}}`$ to better than $`\pm 0.01\%`$. Systematics then become the issue. The dominant systematic uncertainty at the SLD was a $`\pm 0.5\%`$ polarization error which contributes to $`\mathrm{\Delta }\mathrm{sin}^2\theta _W`$ at the $`\pm 0.0001`$ level. One would need to reduce the polarization uncertainty to $`\pm 0.1\%`$ to reach $`\pm 0.01\%`$ in $`\mathrm{sin}^2\theta _W(m_Z)_{\overline{\mathrm{MS}}}`$. Such a reduction would be possible if both the $`e^+`$ and $`e^{}`$ beams were polarized. Then, the effective polarization (they add like relativistic velocities)
$`P_{\mathrm{eff}}={\displaystyle \frac{P_e^{}P_{e^+}}{1P_e^{}P_{e^+}}}`$ (12)
enters
$`{\displaystyle \frac{N_{LR}N_{RL}}{N_{LR}+N_{RL}}}=P_{\mathrm{eff}}A_{\mathrm{LR}},`$ (13)
where $`N_{LR}`$ denotes the number of $`e_L^{}e_R^+`$ induced hadronic $`Z`$ decays. For $`|P_e^{}|=0.9000\pm 0.0045`$ and $`|P_{e^+}|=0.6500\pm 0.0065`$ (i.e. $`\pm 1\%`$ $`e^+`$ polarization), one finds $`P_{\mathrm{eff}}=0.9779\pm 0.0012`$ as required for a $`\pm 0.01\%`$ determination of $`\mathrm{sin}^2\theta _W(m_Z)_{\overline{\mathrm{MS}}}`$.
Improving the direct measurement of $`\mathrm{sin}^2\theta _W(m_Z)_{\overline{\mathrm{MS}}}`$ can have other applications. The $`Z`$ pole determination is relatively pure and free of “new physics.” Below, we demonstrate its utility for comparison with polarized Møller scattering asymmetries which could exhibit effects from “new physics” beyond the Standard Model.
## 2 Polarized Møller Scattering – Fixed Target
Møller scattering $`e^{}e^{}e^{}e^{}`$ has been a well studied, classic low energy reaction.<sup>?</sup> Employing polarized electrons, one can, in principle, measure parity violating weak interaction asymmetries.<sup>?</sup> At tree level, the $`A_{\mathrm{LR}}`$ in Møller scattering comes from an interference among the diagrams in Fig. 2. For a single polarized $`e^{}`$, the asymmetry corresponds to
$`A_{\mathrm{LR}}^{(1)}{\displaystyle \frac{\mathrm{d}\sigma _{\mathrm{LL}}+\mathrm{d}\sigma _{\mathrm{LR}}\mathrm{d}\sigma _{\mathrm{RL}}\mathrm{d}\sigma _{\mathrm{RR}}}{\mathrm{d}\sigma _{\mathrm{LL}}+\mathrm{d}\sigma _{\mathrm{LR}}+\mathrm{d}\sigma _{\mathrm{RL}}+\mathrm{d}\sigma _{\mathrm{RR}}}}`$ (14)
while in the case of both $`e^{}`$ polarized, a second asymmetry becomes possible<sup>?</sup>
$`A_{\mathrm{LR}}^{(2)}{\displaystyle \frac{\mathrm{d}\sigma _{\mathrm{LL}}\mathrm{d}\sigma _{\mathrm{RR}}}{\mathrm{d}\sigma _{\mathrm{LL}}+\mathrm{d}\sigma _{\mathrm{RR}}}}.`$ (15)
The subscripts denote the initial $`e^{}e^{}`$ states’ polarizations. As we subsequently show, both asymmetries would be measurable at a high energy $`e^{}e^{}`$ collider where polarizations of $`0.90`$ for each beam are likely. Since d$`\sigma _{LR}=\mathrm{d}\sigma _{RL}`$ by rotational invariance, they differ only in their denominators.
Let us begin by considering a fixed target scenario in which a 50 GeV polarized electron beam scatters off a fixed target of electrons. That case will be addressed in the near future by SLAC experiment E158.<sup>?</sup>
In the center-of-mass frame, the differential cross section is characterized by the scattering angle $`\theta `$ with respect to the beam axis or
$`y={\displaystyle \frac{1\mathrm{cos}\theta }{2}},0\theta \pi .`$ (16)
The variable $`y`$ relates the momentum transfer $`Q^2=q^2`$ and center-of-mass energy $`\sqrt{s}`$ via
$`Q^2=ys,0y1.`$ (17)
Since the cross section grows as $`1/y^2(1y)^2s`$, very high statistics are possible at small angle and/or small $`s`$. However, the asymmetry grows with $`s`$. All things considered, it is generally better to measure $`A_{\mathrm{LR}}`$ at high $`s`$, but lower energy fixed target facilities can compensate by having very large effective luminosities. For example, E158 at SLAC will have $`s0.05\mathrm{GeV}^2`$ and aims to measure (with high precision) a very small asymmetry $`A_{\mathrm{LR}}1.5\times 10^7`$. That is only possible because their luminosity will be $`4\times 10^{38}\mathrm{cm}^2/s`$.
At small $`Q^2=ysm_Z^2`$, the left-right polarization asymmetry in Møller scattering is given by (at tree level)<sup>?</sup>
$`A_{\mathrm{LR}}^{(1)}(e^{}e^{}e^{}e^{})={\displaystyle \frac{G_\mu s}{\sqrt{2}\pi \alpha }}{\displaystyle \frac{y(1y)}{1+y^4+(1y)^4}}(14\mathrm{sin}^2\theta _W),`$ (18)
or for comparison with the $`Z`$ pole asymmetry
$`A_{\mathrm{LR}}^{(1)}(e^{}e^{}e^{}e^{})`$ (19)
$`={\displaystyle \frac{12}{\alpha }}{\displaystyle \frac{y(1y)}{1+y^4+(1y)^4}}{\displaystyle \frac{s\mathrm{\Gamma }(Ze^+e^{})}{m_Z^3}}A_{\mathrm{LR}}(e^+e^{}Z\mathrm{hadrons}).`$
To be at all competitive with the $`\pm 0.00028`$ uncertainty in $`\mathrm{sin}^2\theta _W`$ found by SLD, very high statistics are required or equivalently, a very good determination of $`A_{\mathrm{LR}}`$,
$`{\displaystyle \frac{\delta \mathrm{sin}^2\theta _W}{\mathrm{sin}^2\theta _W}}{\displaystyle \frac{14\mathrm{sin}^2\theta _W}{4\mathrm{sin}^2\theta _W}}{\displaystyle \frac{\delta A_{\mathrm{LR}}}{A_{\mathrm{LR}}}}.`$ (20)
Again, one sees the enhanced sensitivity to small changes in $`\mathrm{sin}^2\theta _W`$. E158 aims for a $`\pm 0.0007`$ to $`\pm 0.0004`$ measurement of $`\mathrm{sin}^2\theta _W`$ which will make it the best low energy determination of that quantity. As we subsequently illustrate, it will be sensitive to the running of the weak mixing angle as well as “new physics” effects.
## 3 Polarized Møller Scattering at Collider Energies
Møller scattering, $`e^{}e^{}e^{}e^{}`$, at the NLC can also be used for precision tests of the Standard Model as well as direct and indirect searches for “new physics.”<sup>?,?</sup> Indeed, in some cases it can provide a more powerful probe than $`e^+e^{}`$. One can assume with some confidence that both $`e^{}`$ beams will be polarized with $`|P_1|=|P_2|=0.9`$ and about $`\pm 0.5`$% uncertainty each. The effective polarization will therefore be (with like sign $`P_1`$ and $`P_2`$)
$`P_{\mathrm{eff}}={\displaystyle \frac{P_1+P_2}{1+P_1P_2}}=0.9945\pm 0.0004.`$ (21)
We see that $`P_{\mathrm{eff}}`$ will be very large and has essentially negligible uncertainty compared to $`P_1`$ and $`P_2`$.
The differential cross section in high energy collider Møller scattering is also characterized by a single parameter, the scattering angle $`\theta `$ with respect to the beam axis or
$`y={\displaystyle \frac{1\mathrm{cos}\theta }{2}},0\theta \pi .`$ (22)
The cross section grows as $`1/y^2`$ for small angle scattering. Hence, very high statistics are possible in the small angle region. Good angular coverage is therefore important for precision measurements. As before, the variable $`y`$ relates $`s`$ and the momentum transfer $`Q^2=q^2`$ via $`Q^2=ys,0y1`$. Note, that $`y`$ and $`1y`$ correspond to indistinguishable events. Very forward (small angle) $`e^{}e^{}`$ events will therefore be composed of high and low $`Q^2`$ contributions.
As previously noted, one can consider two distinct but similar parity violating Møller asymmetries: the single spin asymmetry $`A_{\mathrm{LR}}^{(1)}`$ defined in eq. (14) and double spin asymmetry $`A_{\mathrm{LR}}^{(2)}`$ in eq. (15).
Experimentally, one can and probably will flip the individual polarizations (pulse by pulse) and measure $`N_{\mathrm{LL}}`$, $`N_{\mathrm{LR}}`$, $`N_{\mathrm{RL}}`$, and $`N_{\mathrm{RR}}`$ (the number of events in each mode) for fixed luminosity and polarization. From those measurements, the polarizations and $`A_{\mathrm{LR}}^{(2)}(y)`$ can be simultaneously determined using<sup>?,?</sup>
$`{\displaystyle \frac{N_{\mathrm{LL}}+N_{\mathrm{LR}}N_{\mathrm{RL}}N_{\mathrm{RR}}}{N_{\mathrm{LL}}+N_{\mathrm{LR}}+N_{\mathrm{RL}}+N_{\mathrm{RR}}}}`$ $`=`$ $`P_1A_{\mathrm{LR}}^{(1)}(y),`$ (23)
$`{\displaystyle \frac{N_{\mathrm{RR}}+N_{\mathrm{LR}}N_{\mathrm{RL}}N_{\mathrm{LL}}}{N_{\mathrm{RR}}+N_{\mathrm{LR}}+N_{\mathrm{RL}}+N_{\mathrm{LL}}}}`$ $`=`$ $`P_2A_{\mathrm{LR}}^{(1)}(y),`$ (24)
$`{\displaystyle \frac{N_{\mathrm{LL}}N_{\mathrm{RR}}}{N_{\mathrm{LL}}+N_{\mathrm{RR}}}}`$ $`=`$ $`P_{\mathrm{eff}}A_{\mathrm{LR}}^{(2)}(y)\left({\displaystyle \frac{1}{1+\frac{1P_1P_2}{1+P_1P_2}\frac{\sigma _{\mathrm{LR}}+\sigma _{\mathrm{RL}}}{\sigma _{\mathrm{LL}}+\sigma _{\mathrm{RR}}}}}\right),`$ (25)
$`P_{\mathrm{eff}}`$ $`=`$ $`{\displaystyle \frac{P_1+P_2}{1+P_1P_2}}.`$
For $`P_1=P_2=0.9`$, the correction term in parentheses of Eq. (25) is small but must be accounted for. Using Eq. (25), $`A_{\mathrm{LR}}^{(2)}`$ (which depends on $`\mathrm{sin}^2\theta _W`$) can be extracted from data and compared with the Standard Model prediction. A deviation from expectations would signal “new physics.”
In general the d$`\sigma _{ij}`$ for Møller scattering are somewhat lengthy expressions<sup>?</sup> with contributions from direct and crossed $`\gamma `$ and $`Z`$ exchange amplitudes (see Fig. 1). To simplify our discussion, we consider for illustration the case $`ys`$ and $`(1y)sm_Z^2`$; so, terms of relative order $`m_Z^2/ys`$ and $`m_Z^2/(1y)s`$ can be neglected. In that limit, one finds at tree level<sup>?</sup>
$`{\displaystyle \frac{\mathrm{d}\sigma _{\mathrm{LL}}}{\mathrm{d}y}}`$ $`=`$ $`\sigma _0{\displaystyle \frac{1}{y^2(1y)^2}}{\displaystyle \frac{1}{16\mathrm{sin}^4\theta _W}},`$
$`{\displaystyle \frac{\mathrm{d}\sigma _{\mathrm{RR}}}{\mathrm{d}y}}`$ $`=`$ $`\sigma _0{\displaystyle \frac{1}{y^2(1y)^2}},`$
$`{\displaystyle \frac{\mathrm{d}\sigma _{\mathrm{LR}}}{\mathrm{d}y}}`$ $`=`$ $`{\displaystyle \frac{\mathrm{d}\sigma _{\mathrm{RL}}}{\mathrm{d}y}}=\sigma _0{\displaystyle \frac{y^4+(1y)^4}{y^2(1y)^2}}{\displaystyle \frac{1}{4}},`$ (26)
and the asymmetries become
$`A_{\mathrm{LR}}^{(1)}(y)`$ $`=`$ $`{\displaystyle \frac{(14s_W^2)(1+4s_W^2)}{1+16s_W^4+8\left[y^4+(1y)^4\right]s_W^4}},`$ (27)
$`A_{\mathrm{LR}}^{(2)}(y)`$ $`=`$ $`{\displaystyle \frac{(14s_W^2)(1+4s_W^2)}{1+16s_W^4}}.`$ (28)
Expanding about $`\mathrm{sin}^2\theta _W=1/4`$, Eq. (28) becomes
$`A_{\mathrm{LR}}^{(2)}(y)=(14\mathrm{sin}^2\theta _W)+𝒪[(14\mathrm{sin}^2\theta _W)^2].`$ (29)
For arbitrary $`s`$, the asymmetries are maximal at $`y=1/2`$. There we find, up to terms of $`𝒪[(14\mathrm{sin}^2\theta _W)^2]`$,
$`A_{\mathrm{LR}}^{(1)}(y=1/2)`$ $``$ $`(14\mathrm{sin}^2\theta _W){\displaystyle \frac{16x\left(3+2x\right)}{3\left(27+34x+11x^2\right)}},`$
$`A_{\mathrm{LR}}^{(2)}(y=1/2)`$ $``$ $`(14\mathrm{sin}^2\theta _W){\displaystyle \frac{2x}{3+2x}},x{\displaystyle \frac{s}{m_Z^2}}.`$ (30)
Because of the $`(14\mathrm{sin}^2\theta _W)`$ dependence of $`A_{\mathrm{LR}}(e^{}e^{})`$, even with relatively modest angular coverage limited to $`0.1y0.9`$, Møller scattering can be used to measure $`\mathrm{sin}^2\theta _W`$ rather precisely, to about $`\pm 0.0003`$ at $`\sqrt{s}1`$ TeV. Although not likely to compete with future potential very high statistics $`Z`$ pole measurements, it will be competitive with present day measurements. In addition, Møller scattering can be used as a powerful probe for “new physics” effects. Indeed, for electron composite effects parametrized by the four fermion interaction<sup>?</sup> $`\frac{2\pi }{\mathrm{\Lambda }^2}\overline{e}_L\gamma _\mu e_L\overline{e}_L\gamma ^\mu e_L`$ one finds $`\mathrm{\Delta }A_{\mathrm{LR}}sy(1y)c_W^2/\alpha \mathrm{\Lambda }^2`$ for $`e^{}e^{}`$ Møller scattering. It can, therefore, be more sensitive than $`e^+e^{}e^+e^{}`$ (about 50% better) and could probe $`\mathrm{\Lambda }150`$ TeV.
If one is interested in an even more precise determination of $`\mathrm{sin}^2\theta _W`$ via Møller scattering, extremely forward events must be detected. For example, assuming detector acceptance down to about $`5^{}`$ ($`y=0.0019`$), Cuypers and Gambino<sup>?</sup> have shown that $`\mathrm{\Delta }\mathrm{sin}^2\theta _W\pm 0.0001`$ may be possible at a $`\sqrt{s}=2`$ TeV $`e^{}e^{}`$ collider with $`P_1=P_2=90\%`$.
## 4 Radiative Corrections and $`\mathrm{𝐬𝐢𝐧}^\mathrm{𝟐}𝜽_𝑾\mathbf{(}𝑸^\mathrm{𝟐}\mathbf{)}`$
The tree level $`A_{\mathrm{LR}}`$ for both E158 and future $`e^{}e^{}`$ collider studies are proportional to $`14\mathrm{sin}^2\theta _W`$ and hence suppressed because $`\mathrm{sin}^2\theta _W0.23`$. Since some electroweak radiative corrections are not suppressed by $`14\mathrm{sin}^2\theta _W`$, they can be potentially very large. A complete calculation has been carried out<sup>?</sup> for small $`s`$ as appropriate to E158. There it was shown that such effects reduce $`A_{\mathrm{LR}}`$ by 40% and must be included in any detailed study. Here, we comment on the primary sources of those large corrections and show how much of the effect can be incorporated into a running $`\mathrm{sin}^2\theta _W(Q^2)`$. We also discuss how those large effects carry over to collider energies. For a complete study of radiative corrections to Møller scattering at high energies, see ref. ?, ?, ?.
The largest radiative corrections to $`A_{\mathrm{LR}}`$ at low energies come from three sources:
1. $`WW`$ box diagrams,
2. Photonic vertex and box diagrams,
3. $`\gamma Z`$ mixing and the anapole moment.
The first two are of order $`+4\%`$ and $`6\%`$ respectively.<sup>?</sup> $`\gamma Z`$ mixing along with the anapole moment in Fig. 3 is the largest effect. It effectively replaces the tree level $`14\mathrm{sin}^2\theta _W`$ in $`A_{\mathrm{LR}}`$ by<sup>?</sup>
$`14\kappa (0)\mathrm{sin}^2\theta _W(m_Z)_{\overline{\mathrm{MS}}}`$ (31)
where
$`\kappa (0)=1.0301\pm 0.0025`$ (32)
represents a 3% shift in the effective $`\mathrm{sin}^2\theta _W`$ due to loop effects illustrated in Fig. 3. That $`+3\%`$ increase in the effective $`\mathrm{sin}^2\theta _W`$ appropriate for low $`Q^2`$ gives rise to a $`38\%`$ reduction in $`A_{\mathrm{LR}}`$. Interestingly, that reduction actually makes E158 more sensitive to $`\mathrm{sin}^2\theta _W(m_Z)_{\overline{\mathrm{MS}}}`$ as well as “new physics.”
In the case of very large $`Q^2`$, appropriate for $`e^{}e^{}`$ colliders, the electroweak radiative corrections will change and must be reevaluated. In particular, the $`WW`$ box diagram gives a large negative contribution to $`A_{\mathrm{LR}}`$. The effects of $`\gamma Z`$ mixing and anapole moment can also be very large, but they are easy to obtain from the loops in Fig. 3. One finds for arbitrary $`Q^2`$ that they replace $`14\mathrm{sin}^2\theta _W`$ in the tree level asymmetry by
$`14\kappa (Q^2)\mathrm{sin}^2\theta _W(m_Z)_{\overline{\mathrm{MS}}}`$ $``$ $`14\mathrm{sin}^2\theta _W(Q^2),`$
$`\kappa (Q^2)=\kappa _f(Q^2)+\kappa _b(Q^2),`$ (33)
where the subscripts $`f`$ and $`b`$ denote fermion and boson loops, and $`\mathrm{sin}^2\theta _W(Q^2)`$ is a running effective parameter. In perturbation theory (i.e. without QCD dressing)
$`\kappa _f(Q^2)`$ $`=`$ $`1{\displaystyle \frac{\alpha }{2\pi \mathrm{sin}^2\theta _W}}\{{\displaystyle \frac{1}{3}}{\displaystyle \underset{f}{}}(T_{3f}Q_f2\mathrm{sin}^2\theta _WQ_f^2)`$
$`\times [\mathrm{ln}{\displaystyle \frac{m_f^2}{m_Z^2}}{\displaystyle \frac{5}{3}}+4z_f+(12z_f)p_f\mathrm{ln}{\displaystyle \frac{p_f+1}{p_f1}}]\},`$
$`z_f`$ $``$ $`{\displaystyle \frac{m_f^2}{Q^2}},p_f\sqrt{1+4z_f},`$ (34)
with $`T_{3f}=\pm 1/2`$, $`Q_f=`$ fermion charge, and the sum is over all fermions;
$`\kappa _b(Q^2)`$ $`=`$ $`1{\displaystyle \frac{\alpha }{2\pi \mathrm{sin}^2\theta _W}}\{{\displaystyle \frac{42\mathrm{cos}^2\theta _W+1}{12}}\mathrm{ln}\mathrm{cos}^2\theta _W+{\displaystyle \frac{1}{18}}`$
$`\left({\displaystyle \frac{p}{2}}\mathrm{ln}{\displaystyle \frac{p+1}{p1}}1\right)\left[(74z)\mathrm{cos}^2\theta _W+{\displaystyle \frac{1}{6}}(1+4z)\right]`$
$`z[{\displaystyle \frac{3}{4}}z+(z{\displaystyle \frac{3}{2}})p\mathrm{ln}{\displaystyle \frac{p+1}{p1}}+z(2z)\mathrm{ln}^2{\displaystyle \frac{p+1}{p1}}]\},`$
$`z`$ $``$ $`{\displaystyle \frac{m_W^2}{Q^2}},p\sqrt{1+4z}.`$ (35)
(Note, Eqs. (28) and (29) of ref. ? contain misprints in the $`\kappa (Q^2)`$ expressions.)
In Fig. 4 we illustrate the expected dependence of $`\mathrm{sin}^2\theta _W(Q^2)`$ on $`Q`$ and show how well it has already been measured for several $`Q^2`$. We also illustrate the approximate potential of E158 and future $`e^{}e^{}`$ and $`e^+e^{}`$ collider measurements at $`\sqrt{s}=1`$ TeV. One notices a $`2\sigma `$ discrepancy in the atomic parity violation result as compared with Standard Model expectations. That issue could be resolved or made even more interesting by results from E158 at SLAC.
In the case of $`e^{}e^{}`$ collider studies, one can actually map out the variation in $`\mathrm{sin}^2\theta _W(Q^2)`$ in a single experiment through measurements at different $`\theta `$. We illustrate in Fig. 4 the type of running that one is predicted to find at a $`\sqrt{s}=1`$ TeV $`e^{}e^{}`$ collider. Notice, that by going to small angles (low $`Q^2`$), one can obtain very high precision. Of course, within the Standard Model, the measurements at different $`Q^2`$ would be radiatively corrected to provide a single precise determination of $`\mathrm{sin}^2\theta _W(m_Z)_{\overline{\mathrm{MS}}}`$. However, demonstrating the running of $`\mathrm{sin}^2\theta _W(Q^2)`$ over a large range in $`Q^2`$ in a single experiment will be an added bonus.
## 5 “New Physics” Effects
The real utility of high precision $`A_{\mathrm{LR}}`$ measurements away from the $`Z`$ pole is to search for or constrain “new physics.” A disagreement with the extracted $`\mathrm{sin}^2\theta _W(m_Z)_{\overline{\mathrm{MS}}}`$ value from $`Z`$ pole determinations could signal the presence of additional tree or loop level neutral current effects. Examples that have been considered include $`Z^{}`$ bosons, compositeness, anomalous anapole moment effects, doubly charged scalars $`\mathrm{\Delta }^{}`$, extra dimensions, etc. For example, if E158 meets its phase one goal of a $`\pm 0.0007`$ determination of $`\mathrm{sin}^2\theta _W(m_Z)_{\overline{\mathrm{MS}}}`$, it will probe the $`m_{Z_\chi }`$ of SO(10) at about the 800 GeV level, compositeness at the 10–15 TeV scale, the anapole moment at $`10^{17}`$ cm (or the $`X`$ parameter<sup>?</sup> at $`\pm 0.15`$), and $`g^2/m_\mathrm{\Delta }^{}^20.01G_\mu `$.
At an $`e^{}e^{}`$ collider, the larger value of $`s`$ would significantly improve the “new physics” reach. Roughly, at $`\sqrt{s}500`$ GeV one could do a factor of 10 better in $`m_{Z_\chi }`$ and $`\mathrm{\Lambda }_{\mathrm{comp}}`$ than E158. In the case of the doubly charged Higgs, $`g^2/m_\mathrm{\Delta }^{}^25\times 10^5G_\mu `$ would be probed. Of course, the sensitivity would further improve as higher $`\sqrt{s}`$ values are reached.
Parity violating left-right asymmetries have played key roles in establishing the validity of the Standard Model. From the classic SLAC polarized $`eD`$ measurement to the $`Z`$ pole asymmetry, polarized electron beams have proved their worth. They will continue to provide valuable tools during the NLC era both in the $`e^+e^{}`$ and $`e^{}e^{}`$ modes. In the case of precision studies of parity violating left-right scattering asymmetries, short–distance physics up to $`𝒪`$(150 TeV) will be indirectly explored. Even more exciting is the possible direct detection of new phenomena such as supersymmetry at these high energy facilities. If “new physics” is uncovered, polarization will help sort out its properties and decipher its place in nature.
Acknowledgments
This work was supported by the DOE under grant number DE-AC02-76CH00016.
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# Relativistic Corrections for Polarized 𝐽/Ψ-Production in b-decay
## Abstract
We analyse the structure of relativistic corrections in the inclusive production of polarized $`J/\mathrm{\Psi }`$ from b-quark decay. The analysis is performed not only for the production channel in which the $`c\overline{c}`$ pair is a color sinlet, but also for the channels in which the $`c\overline{c}`$ is a color octet. We find that the correction in the color-singlet channel at the tree-level is completely determined by the decay constant of $`J/\mathrm{\Psi }`$, while in the color-octet channels the corrections are characterized by three matrix elements defined in NRQCD, whose numerical values are unknown. We discuss the impact of these corrections on the polarized $`J/\mathrm{\Psi }`$-production, and the impact is so significant that the predictions based on the analysis for the considered process may be unreliable. Finally, we propose an integrated spin observable to measure the polarization of $`J/\mathrm{\Psi }`$.
PACS numbers: 13.25.Hw, 14.40.Gx, 12.38.Bx, 13.85.Ni
preprint: AS-ITP-2000-03 preprint: AS-ITP-2000-03
Quarkonium systems generally are thought as simpler than light hadrons and it may be easier to handle them in the framework of QCD. However, only recently we have been able to treat their decays and productions rigorously, based on a factorization with non-relativistic QCD(NRQCD). In the framework of the NRQCD factorization the effect of short distance is handled with perturbative QCD and the effect of long distance is parameterized with NRQCD matrix elements. The factorization is performed by utilizing the fact that a heavy quark inside a quarkonium moves with a small velocity $`v`$ in the quarkonium rest frame and an expansion in $`v`$ can be employed. In this factorization, inclusive productions of a quarkonium, e.g., like $`J/\mathrm{\Psi }`$, can be imagined at the leading order of $`v`$ as the following: The $`c`$\- and $`\overline{c}`$\- quark are produced at a same space-time point, and then this pair is transformed into the $`J/\mathrm{\Psi }`$. The production of the $`c\overline{c}`$ pair can be calculated with perturbative QCD, while the transformation is a nonperturbative process, which can be described by matrix elements defined in NRQCD. At higher order of $`v`$ various effects are taken into account, e.g., relativistic effect, and the effect of the produced $`c\overline{c}`$ pair which does not have the same quantum numbers as those of $`J/\mathrm{\Psi }`$. For $`J/\mathrm{\Psi }`$ the transition of a $`c\overline{c}`$ pair in a color-octet state happens at higher orders in the small velocity expansion, but in some cases, the produced $`c\overline{c}`$ pair is more likely in a color-octet state than in a color-singlet state, hence the production of $`J/\mathrm{\Psi }`$ through a color-octet $`c\overline{c}`$ pair is not negligible, and has an important contribution to the total production rate. Including this contribution one is able to fitting experimental results of the total production rate obtained in Tevatron.
With the NRQCD factorization, not only the total production rate of a quakonium can be predicted, but also the polarization of the quarkonium, if it has a spin. Recent preliminary measurement at Tevatron by CDF shows that the produced $`J/\mathrm{\Psi }`$ is polarized in the way which is unexpected from theoretical predictions. Several attempts to explain the discrepancy are made. It should be kept in mind that these predictions are only based the theoretical analysison at the leading order in $`\alpha _s`$ and at the leading order in $`v`$. For charmonia the velocity is not very small, typically $`v^20.3`$, and one-loop effects are generally substantial. Adding effects at higher orders of $`\alpha _s`$ and of $`v`$, the predictions may be changed significantly. In this work we study the effects of higher orders in $`v`$ on the polarization of $`J/\mathrm{\Psi }`$ produced in b-quark decay, to see how important the effects are.
In the framework of heavy quark effective theory inclusive decays of a b-flavored hadron can be thought approximately as inclusive decays of a free b-quark. The inclusive $`b`$-quark decay into a $`J/\mathrm{\Psi }`$ were studied at the leading order in $`v`$ in , in the one-loop correction from QCD for the unpolarized decay was studied. In these studies it is already shown that the color-octet $`c\overline{c}`$ pair plays an important role. In this work we will not only considered the color-singlet contribution but also the color-octet contribution. With the NRQCD factorization, the polarized decay width for the process
$$bJ/\mathrm{\Psi }+X$$
(1)
can be written as:
$$\mathrm{\Gamma }_\lambda (J/\mathrm{\Psi })=\underset{n}{}C_n(bc\overline{c}[n]+X)0|O^{J/\mathrm{\Psi }}[n]|0,$$
(2)
in which the coefficients $`C_n(bc\overline{c}[n]+X)`$ describe the production of a $`c\overline{c}`$ pair in a state $`n`$, the matrix elements $`0|O^{J/\mathrm{\Psi }}[n]|0`$ characterize the transition of a $`c\overline{c}`$ pair in $`n`$-state into the $`J/\mathrm{\Psi }`$. The coefficients can be calculated with perturbative QCD because the production of a $`c\overline{c}`$ pair is a short-distance process, while the matrix elements represent nonperturbative effects and they are defined with operators in NRQCD. These matrix elements are scaled by the power of $`v`$ with the rule of power counting in $`v`$. The index $`\lambda `$ stands for the helicity of the $`J/\mathrm{\Psi }`$, $`\lambda =L`$ is for the longitudinal polarization, $`\lambda =T`$ is for the transversal polarization. We will use the matching procedure proposed in to identify the operators at the next-to-leading order in $`v`$ and to calculate the corresponding coefficients.
To do the matching we need to consider the process
$$b(p_b)c(p_1)+\overline{c}_(p_2)+X,$$
(3)
where the momenta are given in the brackets. With a Lorentz transformation we can boost the $`c\overline{c}`$ pair into its rest-frame and denote in the rest-frame the three- momentum of $`c`$ and of $`\overline{c}`$ as $`𝐪`$ and $`𝐪`$ respectively. The total decay width for the process in Eq.(3) can be written
$$\mathrm{\Gamma }(bc\overline{c}+X)=\frac{d^3q}{(2\pi )^3}\widehat{\mathrm{\Gamma }}(𝐪)$$
(4)
The matching condition reads:
$$\widehat{\mathrm{\Gamma }}(𝐪)=\underset{n}{}C_n(bc\overline{c}[n]+X)0|O^{c\overline{c}}[n]|0.$$
(5)
In the above equation, the $`𝐪`$-dependence in the right hand side is only contained in the matrix elements, in which the hadronic state is replaced by the partonic $`c\overline{c}`$ state. With Eq.(5) one can identify the matrix elements appearing in Eq.(2) and can determine the corresponding coefficients.
The effective weak Hamiltonian for $`b`$-quark decay is:
$`H_{\mathrm{eff}}`$ $`=`$ $`{\displaystyle \frac{G_F}{\sqrt{2}}}{\displaystyle \underset{q=s,d}{}}\{V_{cb}V_{cq}^{}[{\displaystyle \frac{1}{3}}C_1(\mu )\overline{c}\gamma ^\mu (1\gamma _5)c\overline{q}\gamma _\mu (1\gamma _5)b`$ (6)
$`+C_8(\mu )\overline{c}T^a\gamma ^\mu (1\gamma _5)c\overline{q}T^a\gamma _\mu (1\gamma _5)b]\}.`$ (7)
We neglected the contributions of QCD penguin operators in $`H_{\mathrm{eff}}`$. $`T^a`$($`a=1,\mathrm{}8`$) is SU(3) color-matrix. The coefficients $`C_1`$ and $`C_8`$ are related to the usual $`C_\pm `$ by
$$C_1(\mu )=2C_+(\mu )C_{}(\mu ),C_8(\mu )=C_+(\mu )+C_{}(\mu ).$$
(8)
With the one-loop evolution one can obtain
$$\frac{C_1(m_b)}{C_8(m_b)}0.18$$
(9)
This indicates that the $`c\overline{c}`$ pair is produced more likely in a color-octet state than in a color-singlet state. Hence the color-octet contribution is significant and should be included in the decay width. With the effective Hamiltonia it is straightforward to calculate $`\widehat{\mathrm{\Gamma }}(𝐪)`$ and to determine the coefficients in Eq.(5). However, the power counting of $`v`$ for matrix elements with partonic states is different than that for matrix elements with hadronic states. Before presenting our results, we briefly discuss the identification of operators in Eq.(2). In general the matrix elements in Eq.(2) take the form
$$0|O^H[n]|0=\underset{X}{}0|𝒦_n|H+XH+X|𝒦_n^{}|0,$$
(10)
where $`H`$ denotes a quarkonium. $`𝒦_n`$ and $`𝒦_n^{}`$ are operators defined in NRQCD. The general constraints can be obtained by considering symmetries of QCD or NRQCD. Because QCD respects charge conjugation- and parity symmetry, the operators must be transformed in the same way under C- and P transformation. The operators $`𝒦_n`$ and $`𝒦_n^{}`$ can be classified with the weight $`j`$ of $`SU(2)`$ representations. Since the total angle momentum is conserved, the operators must have the same weight. By considering the approximated symmetry, spin-symmetry of NRQCD, further constraints can be obtained. With these constraints, one can identify the operators appearing in Eq.(2).
Now we give our results for $`\mathrm{\Gamma }_\lambda (J/\mathrm{\Psi })`$. We write
$$\mathrm{\Gamma }_\lambda (J/\mathrm{\Psi })=\mathrm{\Gamma }_\lambda ^{(1)}+\mathrm{\Gamma }_\lambda ^{(8)},$$
(11)
where the index 1 or 8 stands for color singlet contributions or for color-octet contributions, respectively. The leading order in $`v`$ for color singlet contributions is $`v^0`$, while the leading order for color-octet contributions is $`v^4`$. We expand the quantities:
$`\mathrm{\Gamma }_\lambda ^{(1)}`$ $`=`$ $`\mathrm{\Gamma }_\lambda ^{(1,v^0)}+\mathrm{\Gamma }^{(1,v^2)}+𝒪(v^4),`$ (12)
$`\mathrm{\Gamma }_\lambda ^{(8)}`$ $`=`$ $`\mathrm{\Gamma }_\lambda ^{(8,v^4)}+\mathrm{\Gamma }^{(8,v^6)}+𝒪(v^8).`$ (13)
For color-singlet contributions at the leading order, there is only one matrix element representing the transformation of a $`c\overline{c}`$ pair in $`{}_{}{}^{3}S_{1}^{}`$ state into the $`J/\mathrm{\Psi }`$. At this order we have:
$`\mathrm{\Gamma }_T^{(1,v^0)}`$ $`=`$ $`{\displaystyle \frac{C_1^2G_F^2|V_{cb}|^2}{432\pi m_b^3m_c}}(m_b^24m_c^2)^20|O_1^{J/\mathrm{\Psi }}(^3S_1)|04m_c^2,`$ (14)
$`\mathrm{\Gamma }_L^{(1,v^0)}`$ $`=`$ $`{\displaystyle \frac{C_1^2G_F^2|V_{cb}|^2}{432\pi m_b^3m_c}}(m_b^24m_c^2)^20|O_1^{J/\mathrm{\Psi }}(^3S_1)|0m_b^2,`$ (15)
where we used $`|V_{cd}|^2+|V_{cs}|^21`$. At order of $`v^2`$, there is also one matrix element which represents the effect of the relative movement of the $`c\overline{c}`$ pair inside $`J/\mathrm{\Psi }`$. The results reads:
$$\mathrm{\Gamma }_\lambda ^{(1,v^2)}=\frac{C_1^2G_F^2|V_{cb}^2}{432\pi m_c}m_b(m_b^24m_c^2)\frac{0|P_1^{J/\mathrm{\Psi }}(^3S_1)|0}{m_c^2}G_\lambda (\frac{m_c^2}{m_b^2})$$
(16)
with
$`G_L(y)`$ $`=`$ $`{\displaystyle \frac{3+92y+208y^2+576y^3}{6(1+4y)^2}},`$ (17)
$`G_T(y)`$ $`=`$ $`{\displaystyle \frac{2y(368y304y^2960y^3)}{3(1+4y)^2}}.`$ (18)
The definition of the matrix elements $`0|O_1^{J/\mathrm{\Psi }}(^3S_1)|0`$ and $`0|P_1^{J/\mathrm{\Psi }}(^3S_1)|0`$ can be found in . These matrix elements can only be calculated nonperturbatively and they must be calculated with the same accuracy in order of $`v`$ if one uses both to make predictions.
For color octet contributions at the leading order in $`v`$, the $`c\overline{c}`$ pair can be in $`{}_{}{}^{1}S_{0}^{}`$, $`{}_{}{}^{3}S_{1}^{}`$ and $`{}_{}{}^{3}P_{1}^{}`$ states, these states can be transformed into a $`J/\mathrm{\Psi }`$ through emission or absorption of soft gluons. Correspondingly, we have in the matching three types of matrix elements, these matrix elements can be reduced to three matrix elements with the symmetries mentioned before. They are:
$`{\displaystyle \underset{X}{}}0|\chi ^{}T^a\psi |J/\mathrm{\Psi }+XJ/\mathrm{\Psi }+X|\psi ^{}T^a\chi |0`$ (19)
$`={\displaystyle \frac{1}{3}}\epsilon _i(\lambda )\epsilon _i^{}(\lambda )0|O_8^{J/\mathrm{\Psi }}(^1S_0)|0`$ (20)
$`{\displaystyle \underset{X}{}}0|\chi ^{}T^a\sigma ^i\psi |J/\mathrm{\Psi }+XJ/\mathrm{\Psi }+X|\psi ^{}T^a\sigma ^j\chi |0`$ (21)
$`={\displaystyle \frac{1}{3}}\epsilon _i(\lambda )\epsilon _j^{}(\lambda )0|O_8^{J/\mathrm{\Psi }}(^3S_1)|0`$ (22)
$`{\displaystyle \underset{X}{}}0|\chi ^{}T^a(\sigma \times ({\displaystyle \frac{i}{2}}\stackrel{}{𝐃}))^i\psi |J/\mathrm{\Psi }+XJ/\mathrm{\Psi }+X|\psi ^{}T^a\sigma \times ({\displaystyle \frac{i}{2}}\stackrel{}{𝐃}))^j\chi |0`$ (23)
$`={\displaystyle \frac{1}{3}}(\delta _{ij}\epsilon (\lambda )\epsilon ^{}(\lambda )\epsilon _i(\lambda )\epsilon _j^{}(\lambda )0|O_8^{J/\mathrm{\Psi }}(^3P_1)|0`$ (24)
In the above equations $`\epsilon _i(\lambda )(i=1,2,3)`$ denotes the polarization vector of $`J/\mathrm{\Psi }`$ in its rest-frame. The field $`\psi (\chi )`$ is the field in NRQCD for the $`c(\overline{c})`$ quark. $`𝐃`$ is the spacial part of the covariant derivative $`D^\mu `$. The definition of the operator $`O_8^{J/\mathrm{\Psi }}(^1S_0)`$, $`O_8^{J/\mathrm{\Psi }}(^3S_1)`$ and $`O_8^{J/\mathrm{\Psi }}(^3P_1)`$ can be found in . The rotation invariance leads to Eq.(15), while spin symmetry is used for Eq.(16) and Eq.(17). The leading order of these matrices is at $`v^4`$. It should be noted that Eq.(16) and Eq.(17) also hold at order of $`v^6`$, because to keep the correct spin configuration the spin of the $`c\overline{c}`$ pair must be flipped twice through the interaction which violates the spin symmetry, this can only happen at order of $`v^8`$. With the above identification we obtain the color-octet contribution at the leading order:
$`\mathrm{\Gamma }_L^{(8,v^4)}`$ $`=`$ $`{\displaystyle \frac{C_8^2G_F^2|V_{cb}|^2}{288\pi m_b^3m_c}}(m_b^24m_c^2)^2`$ (25)
$``$ $`\{m_b^20|O_8^{J/\mathrm{\Psi }}(^1S_0)|0+m_b^20|O_8^{J/\mathrm{\Psi }}(^3S_1)|0+80|O_8^{J/\mathrm{\Psi }}(^3P_1)|0\}`$ (26)
$`\mathrm{\Gamma }_T^{(8,v^4)}`$ $`=`$ $`{\displaystyle \frac{C_8^2G_F^2|V_{cb}|^2}{288\pi m_b^3m_c}}(m_b^24m_c^2)^2`$ (27)
$``$ $`\{m_b^20|O_8^{J/\mathrm{\Psi }}(^1S_0)|0+4m_c^20|O_8^{J/\mathrm{\Psi }}(^3S_1)|0+{\displaystyle \frac{m_b^2+4m_c^2}{m_c^2}}0|O_8^{J/\mathrm{\Psi }}(^3P_1)|0\}.`$ (28)
At order of $`v^6`$ for the color-octet contributions various matrix elements appear at the first look. For example, the matrix element
$$\underset{X}{}0|\chi ^{}T^a\sigma ^l(\frac{i}{2})^2\stackrel{}{𝐃}^{(l}\stackrel{}{𝐃}^{i)}\psi |J/\mathrm{\Psi }+XJ/\mathrm{\Psi }+X|\psi ^{}T^a\sigma ^j\chi |0+h.c..$$
(29)
To take a close look at the matrix element we consider
$$B^{ijlk}=\underset{X}{}0|\chi ^{}T^a\sigma ^i(\frac{i}{2})^2\stackrel{}{𝐃}^{(l}\stackrel{}{𝐃}^{k)}\psi |J/\mathrm{\Psi }+XJ/\mathrm{\Psi }+X|\psi ^{}T^a\sigma ^j\chi |0+h.c..$$
(30)
At order of $`v^6`$ the spin symmetry still can be used for the matrix element, it leads:
$$B^{ijlk}=\epsilon _i(\lambda )\epsilon _j^{}(\lambda )A_{lk}.$$
(31)
By rotation invariance, $`A_{lk}`$ is proportional to $`\delta _{lk}`$. Because the tensor $`\stackrel{}{𝐃}^{(l}\stackrel{}{𝐃}^{k)}`$ is symmetric and trace-less, we conclude that $`B^{ijlk}`$ is zero at this order, hence also the matrix element in Eq.(19). This can also be understood as the following: Because the spin symmetry holds, the total orbit angle moment is conserved. The $`X`$-state in the bra in Eq.(20) is a $`l=0`$ state because of the conservation, while the $`X`$-state in the ket is a $`l=2`$ state. Therefore, the matrix element is excluded by the definition of the sum over $`X`$ state. By considering all constraints, only three types of matrix elements remain, and they can be further reduced to three matrix elements:
$`{\displaystyle \underset{X}{}}0|\chi ^{}T^a({\displaystyle \frac{i}{2}}\stackrel{}{𝐃})^2\psi |J/\mathrm{\Psi }+XJ/\mathrm{\Psi }+X|\psi ^{}T^a\chi |0+h.c.`$ (32)
$`={\displaystyle \frac{2}{3}}\epsilon _i(\lambda )\epsilon _i^{}(\lambda )0|P_8^{J/\mathrm{\Psi }}(^1S_0)|0`$ (33)
$`{\displaystyle \underset{X}{}}0|\chi ^{}T^a\sigma ^i({\displaystyle \frac{i}{2}}\stackrel{}{𝐃})^2\psi |J/\mathrm{\Psi }+XJ/\mathrm{\Psi }+X|\psi ^{}T^a\sigma ^j\chi |0+h.c.`$ (34)
$`={\displaystyle \frac{2}{3}}\epsilon _i(\lambda )\epsilon _j^{}(\lambda )0|P_8^{J/\mathrm{\Psi }}(^3S_1)|0`$ (35)
$`{\displaystyle \underset{X}{}}0|\chi ^{}T^a(\sigma \times ({\displaystyle \frac{i}{2}}\stackrel{}{𝐃}))^i({\displaystyle \frac{i}{2}}\stackrel{}{𝐃})^2\psi |J/\mathrm{\Psi }+XJ/\mathrm{\Psi }+X|\psi ^{}T^a\sigma \times ({\displaystyle \frac{i}{2}}\stackrel{}{𝐃}))^j\chi |0+h.c.`$ (36)
$`={\displaystyle \frac{2}{3}}(\delta _{ij}\epsilon (\lambda )\epsilon ^{}(\lambda )\epsilon _i(\lambda )\epsilon _j^{}(\lambda )0|P_8^{J/\mathrm{\Psi }}(^3P_1)|0.`$ (37)
The matrix elements $`0|P_8^{J/\mathrm{\Psi }}(^1S_0)|0`$, $`0|P_8^{J/\mathrm{\Psi }}(^3S_1)|0`$ and $`0|P_8^{J/\mathrm{\Psi }}(^3P_1)|0`$ can be obtained by summing over the helicity and by the contraction over the indices in the above equations. With these matrix elements the results for the color-octet contribution at order of $`v^6`$ are:
$`\mathrm{\Gamma }_\lambda ^{(8,v^6)}`$ $`=`$ $`{\displaystyle \frac{C_8^2G_F^2|V_{cb}|^2}{288\pi m_c^3}}m_b(m_b^24m_c^2)`$ (38)
$``$ $`\{F({\displaystyle \frac{m_c^2}{m_b^2}})0|P_8^{J/\mathrm{\Psi }}(^1S_0)|0+G_\lambda ({\displaystyle \frac{m_c^2}{m_b^2}})0|P_8^{J/\mathrm{\Psi }}(^3S_1)|0`$ (40)
$`+H_\lambda ({\displaystyle \frac{m_c^2}{m_b^2}}){\displaystyle \frac{1}{m_c^2}}0|P_8^{J/\mathrm{\Psi }}(^3P_1)|0\}`$
with
$`F(y)`$ $`=`$ $`{\displaystyle \frac{7+108y+144y^2+320y^3}{6(1+4y)^2}},`$ (41)
$`H_L(y)`$ $`=`$ $`{\displaystyle \frac{4y(1+84y+240y^2+704y^3)}{6(1+4y)^2}},`$ (42)
$`H_T(y)`$ $`=`$ $`{\displaystyle \frac{7+84y+144y^2+704y^3}{6(1+4y)^2}}.`$ (43)
The function $`G_\lambda (y)`$ is given in Eq.(14). The above results indicate that the correction in the color-octet contributions is to take the effect of the relative movement of the $`c\overline{c}`$ pair inside the $`J/\mathrm{\Psi }`$ into account. This is similar as the case with the color-singlet contribution. Hence, the total corrections are just relativistic corrections characterized by four matrix elements.
With the results given above one may predict the polarization of the produced $`J/\mathrm{\Psi }`$ and the total decay width $`\mathrm{\Gamma }=2\mathrm{\Gamma }_T+\mathrm{\Gamma }_L`$, if one knows the numerical values of the eight matrix elements. Among the four matrix elements at the leading order in $`v`$ the best known matrix element is $`0|O_1^{J/\mathrm{\Psi }}(^3S_1)|0`$, which is calculated with potential models and with lattice QCD, while the other three are extracted from experimental data. The value of $`0|O_8^{J/\mathrm{\Psi }}(^3S_1)|0`$ is rather well determined by direct $`J/\mathrm{\Psi }`$ production at large transverse momentum in $`p\overline{p}`$ collisions. However the uncertainty for this determination is large and it can be at the level of $`100\%`$. The other two are not well determined, only certain combination of them is know with a large uncertainty. In a re-analysis of experimental data from Tevatron and from Hera is performed, in which some higher-orders effects due to multiple-gluon initial-states radiation are considered. It is shown that the values of these matrix elements can be changed substantially by including these effects. Based on the leading order results an analysis for predictions of the polarized $`J/\mathrm{\Psi }`$ is given. The four matrix elements at the next-to-leading order in $`v`$ are completely unknown. With the power-counting in $`v`$, we only know that they are suppressed by $`v^2`$ relatively to the corresponding matrix elements at the leading order in $`v`$. All of these prevents us from a detailed prediction for the polarization and for the total decay width. Nevertheless, one can still see the impact of these corrections. For the color-singlet contributions, if one neglects the one-loop QCD correction and uses the vacuum saturation, then with $`H_{\mathrm{eff}}`$ in Eq.(6) one has only one constant representing the nonperturbative effect in the process in Eq.(1):
$$J/\mathrm{\Psi }|\overline{c}\gamma ^\mu c|0=if_{J/\mathrm{\Psi }}M_{J/\mathrm{\Psi }}(\epsilon ^\mu (\lambda ))^{},$$
(44)
where $`f_{J/\mathrm{\Psi }}`$ is the leptonic decay constant of $`J/\mathrm{\Psi }`$ and is related to the leptonic decay width:
$$\mathrm{\Gamma }(J/\mathrm{\Psi }\mathrm{}^+\mathrm{}^{})=\frac{16\pi \alpha _{\mathrm{em}}^2}{27M_{J/\mathrm{\Psi }}}f_{J/\mathrm{\Psi }}^2.$$
(45)
In the above equation, the only approximation is to neglect effects of higher orders in $`\alpha _{\mathrm{em}}`$ and lepton masses, relativistic corrections are automatically included. In NRQCD the vacuum saturation brings uncertainty at order of $`v^4`$. Therefore the color-singlet contributions can be written as:
$`\mathrm{\Gamma }_T^{(1)}`$ $`=`$ $`{\displaystyle \frac{C_1^2G_F^2|V_{cb}|^2}{144\pi m_b^3}}(m_b^2M_{J/\mathrm{\Psi }}^2)^2M_{J/\mathrm{\Psi }}^2f_{J/\mathrm{\Psi }}^2+𝒪(\alpha _s^2)+𝒪(v^4),`$ (46)
$`\mathrm{\Gamma }_L^{(1)}`$ $`=`$ $`{\displaystyle \frac{C_1^2G_F^2|V_{cb}|^2}{144\pi m_b^3}}(m_b^2M_{J/\mathrm{\Psi }}^2)^2m_b^2f_{J/\mathrm{\Psi }}^2+𝒪(\alpha _s^2)+𝒪(v^4),`$ (47)
where the relativistic correction calculated before is contained in the $`J/\mathrm{\Psi }`$-mass $`M_{J/\mathrm{\Psi }}`$ and in the decay constant. Comparing the results in Eq.(29) with that in Eq.(12) we may see how large the relativistic correction is. For this we define:
$$R_\lambda ^{(1)}=\frac{\mathrm{\Gamma }_\lambda ^{(1,v^0)}}{\mathrm{\Gamma }_\lambda ^{(1)}},$$
(48)
where $`\mathrm{\Gamma }_\lambda ^{(1,v^0)}`$ is given in Eq.(12) and the results in Eq.(29) are used for $`\mathrm{\Gamma }_\lambda ^{(1)}`$. To obtain numerical values of $`R_\lambda ^{(1)}`$, we take $`m_b=4.7`$GeV and $`m_c=1.5`$GeV. $`f_{J/\mathrm{\Psi }}`$ is determined by the leptonic decay to be 405MeV. For the matrix element $`0|O_1^{J/\mathrm{\Psi }}(^3S_1)|0`$, its value is determined based on a potential model in to be 1.16(GeV$`)^3`$, which is in agreement with a calculation of lattice QCD in . With these values we obtain:
$$R_L^{(1)}R_T^{(1)}1.6.$$
(49)
This indicates that the correction is negative and very large. Another way to see the impact of the correction is to compare the coefficients in the front of the matrix element. We take the above numerical values for quark masses and obtain:
$`\mathrm{\Gamma }_L^{(1)}`$ $`=`$ $`{\displaystyle \frac{C_1^2G_F^2|V_{cb}|^2}{432\pi }}\{24.30|O_1^{J/\mathrm{\Psi }}(^3S_1)|052.4{\displaystyle \frac{1}{m_c^2}}0|P_1^{J/\mathrm{\Psi }}(^3S_1)|0\}+𝒪(v^4),`$ (50)
$`\mathrm{\Gamma }_T^{(1)}`$ $`=`$ $`{\displaystyle \frac{C_1^2G_F^2|V_{cb}|^2}{432\pi }}\{9.900|O_1^{J/\mathrm{\Psi }}(^3S_1)|011.4{\displaystyle \frac{1}{m_c^2}}0|P_1^{J/\mathrm{\Psi }}(^3S_1)|0\}+𝒪(v^4).`$ (51)
With these results the relativistic correction can be at the level of $`60\%`$ for $`\mathrm{\Gamma }_L^{(1)}`$ and at the level of $`35\%`$ for $`\mathrm{\Gamma }_T^{(1)}`$ if one takes
$$\frac{1}{m_c^2}0|P_1^{J/\mathrm{\Psi }}(^3S_1)|0v^20|O_1^{J/\mathrm{\Psi }}(^3S_1)|0,v^20.3.$$
(52)
Similarly we obtain for the color-octet contributions:
$`\mathrm{\Gamma }_L^{(8)}`$ $`=`$ $`{\displaystyle \frac{C_8^2G_F^2|V_{cb}|^2}{288\pi }}\{24.40|O_8^{J/\mathrm{\Psi }}(^1S_0)|068.3{\displaystyle \frac{1}{m_c^2}}0|P_8^{J/\mathrm{\Psi }}(^1S_0)|0`$ (55)
$`+24.40|O_8^{J/\mathrm{\Psi }}(^3S_1)|052.4{\displaystyle \frac{1}{m_c^2}}0|P_8^{J/\mathrm{\Psi }}(^3S_1)|0`$
$`+19.8{\displaystyle \frac{1}{m_c^2}}0|O_8^{J/\mathrm{\Psi }}(^3P_1)|035.8{\displaystyle \frac{1}{m_c^4}}0|P_8^{J/\mathrm{\Psi }}(^3P_1)|0\}+𝒪(v^8),`$
$`\mathrm{\Gamma }_T^{(8)}`$ $`=`$ $`{\displaystyle \frac{C_8^2G_F^2|V_{cb}|^2}{288\pi }}\{24.40|O_8^{J/\mathrm{\Psi }}(^1S_0)|068.3{\displaystyle \frac{1}{m_c^2}}0|P_8^{J/\mathrm{\Psi }}(^1S_0)|0`$ (58)
$`+9.90|O_8^{J/\mathrm{\Psi }}(^3S_1)|011.4{\displaystyle \frac{1}{m_c^2}}0|P_8^{J/\mathrm{\Psi }}(^3S_1)|0`$
$`+34.3{\displaystyle \frac{1}{m_c^2}}0|O_8^{J/\mathrm{\Psi }}(^3P_1)|086.5{\displaystyle \frac{1}{m_c^4}}0|P_8^{J/\mathrm{\Psi }}(^3P_1)|0\}+𝒪(v^8).`$
If one assumes that there are similar relations among the color-octet matrix elements like that in Eq.(33), then the correction is very large. For example, the correction for the $`{}_{}{}^{1}S_{0}^{}`$ production channel can be at the level of $`80\%`$, and the correction for the $`{}_{}{}^{3}P_{1}^{}`$ production channel can be at the level of $`70\%`$. With the above discussions one may conclude that the predictions based on the leading- and next-to-leading order in the small velocity expansion for the process are unreliable.
The last subject of this work is to propose an integrated observable to measure the spin of the produced $`J/\mathrm{\Psi }`$. Usually, one looks at the leptonic decay of $`J/\mathrm{\Psi }`$ to measure the spin. At a fixed momentum $`𝐏`$ of $`J/\mathrm{\Psi }`$ one measures the distribution:
$$\frac{d\mathrm{\Gamma }}{d\mathrm{cos}\theta }(J/\mathrm{\Psi }\mathrm{}^+\mathrm{}^{})1+\alpha \mathrm{cos}^2\theta ,$$
(59)
where $`\theta `$ is the angle between $`𝐏`$ and $`𝐤`$, $`𝐤`$ is the momentum of the lepton in the $`J/\mathrm{\Psi }`$-rest frame. The parameter $`\alpha `$ is predicted as:
$$\alpha =\frac{\mathrm{\Gamma }_T\mathrm{\Gamma }_L}{\mathrm{\Gamma }_T+\mathrm{\Gamma }_L}.$$
(60)
This may have a disadvantage that it will be hard to determine the distribution if the number of events with a fixed $`𝐏`$ is small, hence the parameter $`\alpha `$. We propose to use an integrated spin observable to overcome this disadvantage. For this purpose we define the density matrix $`R_{ij}`$ of the produce $`J/\mathrm{\Psi }`$. For arbitrary polarization the decay width can be written:
$$\mathrm{\Gamma }(bJ/\mathrm{\Psi }+X)=\frac{1}{4\pi }𝑑\mathrm{\Omega }\epsilon _iR_{ij}(\widehat{𝐏})\epsilon _j,$$
(61)
where $`\mathrm{\Omega }`$ is the solid angle of $`𝐏`$, $`\widehat{𝐏}`$ denotes the direction of $`𝐏`$. Similarly we can define the density matrix $`\rho _{ij}`$ for the leptonic decay. Any observable $`O`$ can be then predicted by
$$O=\frac{1}{N}\frac{d\mathrm{\Omega }}{4\pi }\frac{d\mathrm{\Omega }_k}{4\pi }O\rho _{ij}(\widehat{𝐤})R_{ji}(\widehat{𝐏}),$$
(62)
where $`N`$ is a normalization factor so that $`1=1`$. A simple calculation leads:
$$\rho _{ij}(\widehat{𝐤})=\frac{1}{3}\delta _{ij}\frac{1}{2}(\widehat{k}_i\widehat{k}_j\frac{1}{3}\delta _{ij}),$$
(63)
where $`\rho _{ij}`$ is normalized. $`R_{ij}`$ can be written:
$`R_{ij}(\widehat{𝐏})`$ $`=`$ $`a\delta _{ij}+ib\epsilon _{ijk}\widehat{P}_k+c\widehat{P}_i\widehat{P}_j,`$ (64)
$`a`$ $`=`$ $`\mathrm{\Gamma }_T,c=\mathrm{\Gamma }_L\mathrm{\Gamma }_T.`$ (65)
The anti-symmetric part exits due to that parity is violated, and it is irrelevant here because the parity is conserved in the leptonic decay. With the tensor structure one can construct the integrated spin observable $`O_P`$, and predict its value with Eq.(38):
$$O_P=(\widehat{𝐏}\widehat{𝐤})^2\frac{1}{3},O_P=\frac{2(\mathrm{\Gamma }_T\mathrm{\Gamma }_L)}{15\mathrm{\Gamma }}.$$
(66)
If $`O_P=0`$ the $`J/\mathrm{\Psi }`$ is unpolarized. If one know the values of the matrix elements discussed before, one can obtain the value for $`O_P`$ to compare with the measured in experiment. One may also obtain the invariance $`O_P^2`$ to determine the statistical error of $`O_P`$.
We summarize our work: We have analyzed the relativistic corrections for the polarized $`J/\mathrm{\Psi }`$-production in $`b`$-quark decay. We have calculated the perturbative coefficients and identified the matrix elements at the next-to-leading order in $`v`$. We find that the corrections can be very large in the color-singlet production channel and as well as in the color-octet production channels. For the color-singlet production channel the correction is determined by the leptonic decay constant of $`J/\mathrm{\Psi }`$. These corrections are so large that the predictions for the considered process may be unreliable, if one only keeps several leading terms in the small velocity expansion. An integrated spin observable is proposed to measure the $`J/\mathrm{\Psi }`$ polarization. If detailed information of the eight matrix elements is known, numerical values of polarized decay widths and the observable may be obtained.
Acknowledgment: The author would like to thank Prof. K.T. Chao and Prof. Y.Q. Chen for discussions. This work is supported by National Science Foundation of P.R. China and by the Hundred Yonng Scientist Program of Sinica Academia of P.R.China.
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# Indefinite theta series of signature (1,1) from the point of view of homological mirror symmetry
## 1. Indefinite theta series of signature $`(1,1)`$
Let $`\mathrm{\Lambda }`$ be a rank 2 lattice equipped with a $``$-valued quadratic form $`Q`$, $`\tau `$ be an element in the upper-half plane. We assume that $`Q`$ has signature $`(1,1)`$ and fix a rational open cone $`C\mathrm{\Lambda }_{}`$ such that $`Q|_C>0`$. We will use the following notation for indefinite theta series with characteristics associated with $`(\mathrm{\Lambda },Q,C)`$: for an element $`𝐜\mathrm{\Lambda }_{}/\mathrm{\Lambda }`$ we set
(1.1)
$$\mathrm{\Theta }_{\mathrm{\Lambda },Q,C;𝐜}(𝐳,\tau )=\underset{𝐧𝐜+\mathrm{\Lambda }:𝐧+\alpha (𝐳)C}{}\mathrm{sign}(𝐧+\alpha (𝐳))\mathrm{exp}(\pi i\tau Q(𝐧)+2\pi i𝐧𝐳)$$
In other words, we have
$$\mathrm{\Theta }_{\mathrm{\Lambda },Q,C;𝐜}(𝐳,\tau )=\mathrm{exp}(\pi i\tau Q(𝐜)+2\pi i𝐜𝐳)\mathrm{\Theta }_{\mathrm{\Lambda },Q,C}(𝐳+\tau 𝐜,\tau ).$$
The following identities follow immediately from the definition:
$$\mathrm{\Theta }_{N\mathrm{\Lambda },Q,C,𝐜}(𝐳,\tau )=\mathrm{\Theta }_{\mathrm{\Lambda },Q,C,\frac{𝐜}{N}}(N𝐳,N^2\tau ),$$
$$\mathrm{\Theta }_{\mathrm{\Lambda },NQ,C,𝐜}(𝐳,\tau )=\mathrm{\Theta }_{\mathrm{\Lambda },Q,C,𝐜}(N𝐳,N\tau )$$
where $`N>0`$ in an integer,
$$\mathrm{\Theta }_{\mathrm{\Lambda },Q,C,𝐜}(𝐳,\tau )=\underset{𝐧\mathrm{\Lambda }/\mathrm{\Lambda }^{}}{}\mathrm{\Theta }_{\mathrm{\Lambda }^{},Q,C;𝐜+𝐧}(z,\tau )$$
for any sublattice $`\mathrm{\Lambda }^{}\mathrm{\Lambda }`$. Since $`N\mathrm{\Lambda }\mathrm{\Lambda }^{}`$ for some $`N`$ we can also use these formulas to express $`\mathrm{\Theta }_{\mathrm{\Lambda }^{},Q,C}`$ in terms of $`\mathrm{\Theta }_{\mathrm{\Lambda },Q,C}`$.
On the other hand, since the cone $`C`$ is rational we can choose coordinates in such a way that $`C=\{(x_1,x_2)^2:x_1x_2>0\}`$, $`x_i>0`$ in $`C^+`$ and $`\mathrm{\Lambda }`$ is a lattice in $`^2`$ commensurable with $`^2`$. Now the condition $`Q|_C>0`$ and the requirement that the signature of $`Q`$ is $`(1,1)`$ mean that $`Q(x_1,x_2)=a_{11}x_1^2+2a_{12}x_1x_2+a_{22}x_2^2`$ where $`a_{ii}0`$ and $`D=a_{12}^2a_{11}a_{22}>0`$.
Now let us consider the case when $`Q`$ splits into a product of linear forms over $``$. Then by additivity of $`\mathrm{\Theta }`$ in $`C`$ it suffices to consider the case when $`Q`$ vanishes on one of the lines forming the boundary of $`C`$. Then we can choose coordinates in such a way that $`C=\{(x_1,x_2)^2:x_1x_2>0\}`$, $`x_i>0`$ in $`C^+`$, $`Q(n_1,n_2)=an_1(n_1+2n_2)`$ for some $`a`$, $`\mathrm{\Lambda }`$ is a lattice commensurable with $`^2`$. Rescaling $`\tau `$ and $`𝐳`$ (rationally) we can assume that $`a=1`$. Also it suffices to consider the case $`\mathrm{\Lambda }=^2`$. Then for any $`𝐜=(c_1,c_2)`$ we have
$`\mathrm{\Theta }_{^2,Q,C,𝐜}(𝐳,\tau )={\displaystyle \underset{𝐧^2+𝐜,(n_1+\alpha (z_1))(n_2+\alpha (z_2))>0}{}}\mathrm{sign}(n_1+\alpha (z_1))\times `$
$`\mathrm{exp}(\pi i\tau n_1(n_1+2n_2)+2\pi i(n_1z_2+n_2z_1)+2\pi in_1z_1).`$
We can split this sum in two pieces and sum the geometric progression in $`n_2`$ in each of them:
$`\mathrm{\Theta }_{^2,Q,C,𝐜}(𝐳,\tau )=`$
$`{\displaystyle \underset{n_1+c_1,n_1+\alpha (z_1)>0}{}}\mathrm{exp}(\pi i\tau n_1^2+2\pi in_1(z_1+z_2)){\displaystyle \underset{n_2_0+n_2^0}{}}\mathrm{exp}(2\pi in_2(\tau n_1+z_1))`$
$`{\displaystyle \underset{n_1+c_1,n_1+\alpha (z_1)<0}{}}\mathrm{exp}(\pi i\tau n_1^2+2\pi in_1(z_1+z_2)){\displaystyle \underset{n_2_0+n_2^01}{}}\mathrm{exp}(2\pi in_2(\tau n_1+z_1))=`$
$`{\displaystyle \underset{n_1+c_1,n_1+\alpha (z_1)>0}{}}{\displaystyle \frac{\mathrm{exp}(\pi i\tau n_1^2+2\pi in_1(z_1+z_2)+2\pi in_2^0(\tau n_1+z_1))}{1\mathrm{exp}(2\pi i(\tau n_1+z_1))}}`$
$`{\displaystyle \underset{n_1+c_1,n_1+\alpha (z_1)<0}{}}{\displaystyle \frac{\mathrm{exp}(\pi i\tau n_1^2+2\pi in_1(z_1+z_2)+2\pi i(n_2^01)(\tau n_1+z_1))}{1\mathrm{exp}(2\pi i(\tau n_1+z_1))}}=`$
$`{\displaystyle \underset{n_1+c_1}{}}{\displaystyle \frac{\mathrm{exp}(\pi i\tau n_1^2+2\pi in_1(z_1+z_2)+2\pi in_2^0(\tau n_1+z_1))}{1\mathrm{exp}(2\pi i(\tau n_1+z_1))}}`$
where $`n_2^0`$ is the minimal $`n_2+c_2`$ such that $`n_2+\alpha (z_2)>0`$. Hence, we derive the following formula:
(1.2)
$$\mathrm{\Theta }_{^2,Q,C,𝐜}(𝐳,\tau )=\mathrm{exp}(2\pi in_2^0z_1)\kappa _{c_1}(z_1,(1n_2^0)\tau z_1z_2;\tau ).$$
## 2. Fukaya category of a torus
### 2.1. Definition
Let us recall the definition of the Fukaya $`A_{\mathrm{}}`$-category of the torus $`^2/^2`$ with the (complexified) symplectic form $`2\pi i\tau dxdy`$ where $`\tau `$ is an element of the upper half-plane (for more details see ). More precisely, it is not quite an $`A_{\mathrm{}}`$-category since morphisms are only defined for transversal configurations of objects, however, the axiomatics can be changed appropriately (see , sec. 4.3) Also, we will need only the subcategory $`_s`$ which is described as follows. The objects of $`_s`$ are pairs $`(L,t)`$ where $`L^2`$ is a non-vertical line with rational slope considered modulo translations by $`^2`$, $`t`$ is a real number. Morphisms between two such objects $`(L_1,t_1)`$ and $`(L_2,t_2)`$ are defined only if $`L_1L_2mod^2`$. In this case $`\mathrm{Hom}((L_1,t_1),(L_2,t_2))=\mathrm{Hom}(L_1,L_2)`$ is a $``$-vector space with the basis $`[P]`$ enumerated by points $`P(L_1+^2)(L_2+^2)`$ modulo $`^2`$ (the numbers $`t_i`$ will play a role only in the definition of compositions). Let $`\lambda _i`$ be the slope of the line $`L_i`$ ($`i=1,2`$). Then $`\mathrm{Hom}(L_1,L_2)0`$ only if $`\lambda _1\lambda _2`$. This space has grading $`0`$ if $`\lambda _1<\lambda _2`$ and grading $`1`$ if $`\lambda _1>\lambda _2`$. By definition the differential $`m_1`$ is zero. The compositions $`m_k`$ for $`k2`$ are (partially) defined as follows. Let $`L_0,L_1,\mathrm{},L_k`$ be the set of lines in $`^2`$ with slopes $`\lambda _0,\lambda _1,\mathrm{},\lambda _k`$. Assume that the images of $`L_i`$ in $`^2/^2`$ form a transversal configuration, i.e., no three of them intersect in one point. For every $`i=0,\mathrm{},k1`$ let $`d_i`$ be the grading of $`\mathrm{Hom}(L_i,L_{i+1})`$. The composition
$$m_k:\mathrm{Hom}(L_0,L_1)\mathrm{}\mathrm{Hom}(L_{k1},L_k)\mathrm{Hom}(L_0,L_k)$$
is non-zero only if $`_{i=0}^{k1}d_ik+2`$ is equal to the degree of $`\mathrm{Hom}(L_0,L_k)`$. Let $`P_{i,i+1}`$ be some intersection points of $`L_i`$ and $`L_{i+1}`$ modulo $`^2`$. Then
$`m_k([P_{0,1}],[P_{1,2}],\mathrm{},[P_{k1,k}])=`$
$`{\displaystyle \underset{P_{0,k},\mathrm{\Delta }}{}}\pm \mathrm{exp}\left(2\pi i\tau {\displaystyle _\mathrm{\Delta }}𝑑xdy+2\pi i{\displaystyle \underset{j/(k+1)}{}}(x(p_j)x(p_{j1}))t_j\right)[P_{0,k}]`$
where the sum is taken over points of intersections $`P_{0,k}`$ of $`L_0`$ with $`L_k`$ modulo $`^2`$ and over all $`(k+1)`$-gons $`\mathrm{\Delta }`$ (considered up to traslation by $`^2`$) with vertices $`p_iP_{i,i+1}mod^2`$, $`i/(k+1)`$, such that the edge $`[p_{i1},p_i]`$ belongs to $`L_i+^2`$. We also require that the path formed by the edges $`[p_0,p_1],[p_1,p_2],\mathrm{},[p_k,p_0]`$ goes in the clockwise direction. The sign in the RHS is “plus” if $`k`$ is even and is equal to the sign of $`x(p_0)x(p_k)`$ if $`k`$ is odd.
The $`A_{\mathrm{}}`$-constraint we are going to use is
$`m_3(m_2(a_1,a_2),a_3,a_4)m_3(a_1,m_2(a_2,a_3),a_4)+m_3(a_1,a_2,m_2(a_3,a_4))=`$
(2.1) $`=m_2(m_3(a_1,a_2,a_3),a_4)+(1)^{\mathrm{deg}(a_1)}m_2(a_1,m_3(a_2,a_3,a_4)),`$
where $`a_1,\mathrm{},a_4`$ are composable morphisms between $`5`$ objects in $`_s`$ forming a transversal configuration. Below we will often abbreviate $`m_2(a,b)`$ to $`ab`$.
### 2.2. Double products and vector bundles on elliptic curves
Since $`m_1=0`$ the composition $`m_2`$ is associative, so we can consider the category $`_s`$ with $`m_2`$ as a usual category. It was shown in that the obtained category is equivalent to the category of stable vector bundles on the elliptic curve $`E=/+\tau `$, where morphisms between vector bundles $`V_1`$ and $`V_2`$ are elements of the graded vector space $`_i\mathrm{Ext}^i(V_1,V_2)`$ (in fact, we showed in how to extend this equivalence to the whole derived category of coherent sheaves on elliptic curve, but we don’t need this extension here). The construction of this equivalence (which we recall below) is based on the observation due to M. Kontsevich that the Fukaya product $`m_2`$ on a torus is given essentially by theta functions. Here is a more precise statement.
For a pair $`(\lambda ,y)`$ where $`\lambda `$, $`y`$, let us denote by $`L(\lambda ,y)`$ the line in $`^2`$ given by
$$L(\lambda ,y)=\{(t,\lambda ty),t\}.$$
Now let $`L_i=L(\lambda _i,y_i)`$, $`i=0,1,2`$, be lines in $`^2`$ with distinct slopes $`\lambda _i`$. Let us denote
$`y_{ij}={\displaystyle \frac{y_jy_i}{\lambda _j\lambda _i}},`$
$`y_{ij}^{}={\displaystyle \frac{\lambda _iy_j\lambda _jy_i}{\lambda _j\lambda _i}}.`$
The lines $`L_i`$ and $`L_j`$ intersect at the point
$$P_{ij}(y_i,y_j)=(y_{ij},y_{ij}^{}).$$
Note that if we change $`y_j`$ by $`y_j+m\lambda _j+n`$ where $`m,n`$, the new line $`L(\lambda _j,y_j+m\lambda _j+n)`$ is a shift of $`L_j`$ by an integer vector. Thus, the new point of intersection $`P_{ij}(y_i,y_j+m\lambda _j+n)`$ still belongs to $`L_i(L_j+^2)`$. One has $`P_{ij}(y_i,y_j+m\lambda _j+n)P_{ij}(y_i,y_j)mod^2`$ if and only if $`(m,n)\mathrm{\Lambda }(\lambda _i,\lambda _j)`$ where
$$\mathrm{\Lambda }(\lambda _i,\lambda _j)=\{(m,n)^2:\frac{m\lambda _j+n}{\lambda _j\lambda _i}I_{\lambda _i}\}$$
where for every $`\lambda `$ we denote
(2.2)
$$I_\lambda =\{n:\lambda n\}.$$
Thus, we have the following basis in $`\mathrm{Hom}(L_i,L_j)`$:
$$[P_{ij}(y_i,y_j+m\lambda _j+n)],(m,n)^2/\mathrm{\Lambda }(\lambda _i,\lambda _j).$$
Instead of shifting $`y_j`$ we could also shift $`y_i`$ and get a different indexing of intersection points modulo $`^2`$. However, this indexing is related to the previous one by the formula
$$P_{ij}(y_im\lambda _in,y_j)=P_{ij}(y_i,y_j+m\lambda _j+n).$$
Note also that we have $`\mathrm{\Lambda }(\lambda _i,\lambda _j)=\mathrm{\Lambda }(\lambda _j,\lambda _i)`$, so changing the order of lines we would get essentially the same indexing.
Assume that $`\mathrm{deg}\mathrm{Hom}(L_0,L_1)+\mathrm{deg}\mathrm{Hom}(L_1,L_2)=\mathrm{deg}\mathrm{Hom}(L_0,L_2)`$. Let $`t_i`$, $`i=0,1,2`$, be some real numbers. Then we can consider objects $`(L_i,t_i)`$ in Fukaya category. An easy computation shows that
$`m_2([P_{01}(y_0,y_1)],[P_{12}(y_1,y_2)])=`$
(2.3) $`{\displaystyle \underset{nI_{\lambda _1}}{}}\mathrm{exp}(\pi i\tau p(v_1+n)^22\pi ip(v_1+n)w_1)[P_{02}(y_0,y_2+n\lambda _2n\lambda _1)]`$
where
$$p=p(\lambda _0,\lambda _1,\lambda _2)=\frac{(\lambda _2\lambda _1)(\lambda _1\lambda _0)}{(\lambda _2\lambda _0)},$$
$`v_1=y_{12}y_{01}`$, $`w_1=t_{12}t_{01}`$,
$$t_{ij}=\frac{t_jt_i}{\lambda _j\lambda _i},$$
$`I_{\lambda _1}`$ is defined by (2.2). Note that the class of the point $`[P_{02}(y_0,y_2+n\lambda _2n\lambda _1)]`$ modulo $`^2`$ depends only on the class of $`n`$ modulo the following subgroup
$$I_{\lambda _0,\lambda _1,\lambda _2}=I_{\lambda _1}\frac{\lambda _2\lambda _0}{\lambda _2\lambda _1}I_{\lambda _0}.$$
The matrix coefficients of the above product are given by values of elliptic functions at $`(\tau v_1w_1,\tau )`$ times the non-holomorphic factor $`\mathrm{exp}(\pi i\tau pv_1^22\pi ipv_1w_1)`$. One can get rid of this factor by rescaling the bases in $`\mathrm{Hom}(L_i,L_j)`$ appropriately. Namely, we set
(2.4)
$$e_{ij}(m,n)=e_{y_i,y_j}(m,n)=\mathrm{exp}(\pi i\tau (\lambda _i\lambda _j)y_{ij}^22\pi i(t_it_j)y_{ij})[P_{ij}(y_i,y_j+m\lambda _j+n)]$$
where $`(m,n)^2/\mathrm{\Lambda }(\lambda _i,\lambda _j)`$. Then the above formula is equivalent to
$$m_2(e_{01}(0,0),e_{12}(0,0))=\underset{nI_{\lambda _1}/I_{\lambda _0,\lambda _1,\lambda _2}}{}\theta _{I_{\lambda _0,\lambda _1,\lambda _2},n}(p(v_1\tau w_1),p\tau )e_{02}(n,n\lambda _1)$$
where we use the notation $`\theta _{I,c}`$ from the introduction. Changing $`y_i`$’s appropriately in the formula (2.2) we derive a more general formula
$`m_2(e_{01}(a,b),e_{12}(c,d))=`$
(2.5) $`{\displaystyle \underset{nI_{\lambda _1}/I_{\lambda _0,\lambda _1,\lambda _2}}{}}\theta _{I_{\lambda _0,\lambda _1,\lambda _2},u+n}(p(v_1\tau w_1),p\tau )e_{02}(a+c+n,b+dn\lambda _1)`$
where
$$u=\frac{c\lambda _2+d}{\lambda _2\lambda _1}\frac{a\lambda _0+b}{\lambda _1\lambda _0}.$$
The corresponding coefficients will be holomorphic in $`v_1\tau w_1`$. The associativity condition for $`m_2`$ is equivalent to the classical addition formulas for theta-functions.
The equivalence with the category of stable bundles on elliptic curve $`E_\tau =/+\tau `$ is constructed in as follows. First let us consider the subcategory in $`_s`$ formed by lines with integer slopes. To an object of this subcategory $`(L(\lambda ,y),t)`$ where $`\lambda `$, $`y,t`$, we associate the line bundle $`t_{y\tau t}^{}^{(\lambda 1)}`$ where $`=_\tau `$ is the line bundle of degree $`1`$ on $`E_\tau `$ such that $`\theta (z,\tau )`$ is a holomorphic section of $`_\tau `$, $`t_z:E_\tau E_\tau `$ denotes the translation by $`z`$. Assume that we have two such objects $`(L_i,t_i)`$, $`i=1,2`$, where $`L_i=L(\lambda _i,y_i)`$, $`\lambda _i`$, $`\lambda _1<\lambda _2`$. Then we identify $`\mathrm{Hom}(L_1,L_2)`$ with the space of morphisms between the corresponding line bundles by sending the basis elements $`e_{12}(0,k)`$, $`k/(\lambda _2\lambda _1)`$, defined by (2.4) to the functions
$$\theta _{(\lambda _2\lambda _1),k}(z+y_{12}\tau t_{12},\frac{\tau }{\lambda _2\lambda _1})$$
regarded as holomorphic sections of the line bundle
$$(t_{y_1\tau t_1}^{}^{(\lambda _11)})^{}(t_{y_2\tau t_2}^{}^{(\lambda _21)})t_{y_{12}\tau t_{12}}^{}^{(\lambda _2\lambda _1)}.$$
The fact that this map respects $`m_2`$ follows from addition formulas for theta functions. To extend this equivalence to all lines and all stable bundles we use isogenies. For every positive integer $`r`$ consider the natural isogeny of degree $`r`$
$$\pi _r:E_{r\tau }E_\tau .$$
Then we have the natural functors $`\pi _r`$ and $`\pi _r^{}`$ between the categories of bundles on $`E_\tau `$ and $`E_{r\tau }`$. We complete the construction of our equivalence by requiring that these functors correspond to the obvious functors $`\pi _r`$ and $`\pi _r^{}`$ between the corresponding Fukaya categories (see for details). One also has to identify morphisms of degree $`1`$ in both categories. For this one has to fix a non-zero holomorphic $`1`$-form on $`E_\tau `$ and use the isomorphisms $`\mathrm{Hom}(V_1,V_2)^{}\mathrm{Ext}^1(V_2,V_1)`$ (where $`V_1`$ and $`V_2`$ are vector bundles on $`E_\tau `$) induced by Serre duality together with the obvious isomorphisms $`\mathrm{Hom}^0(L_1,L_2)^{}\mathrm{Hom}^1(L_2,L_1)`$ in the Fukaya category.
### 2.3. Triple products and indefinite theta series
Consider $`4`$ lines $`(L_i=L(\lambda _i,y_i),i/4)`$ where $`\lambda _i`$, $`y_i`$. As before, we assume that the corresponding circles in $`^2/^2`$ form a transversal configuration, in particular, the lines $`L_i`$ are distinct modulo $`^2`$ and $`\lambda _i\lambda _{i+1}`$ for $`i/4`$. Let $`t_i,i/4`$ be some real numbers, then $`(L_i,t_i)`$ are objects of the Fukaya category. We are going to compute the Fukaya triple product $`m_3([P_{01}],[P_{12}],[P_{23}])`$, where $`P_{i,i+1}:=P_{i,i+1}(y_i,y_{i+1})`$ for $`i=0,1,2`$. This product is zero unless the following equality is satisfied:
(2.6)
$$\underset{i=0}{\overset{2}{}}\mathrm{deg}\mathrm{Hom}(L_i,L_{i+1})=\mathrm{deg}\mathrm{Hom}(L_0,L_3)+1.$$
Following the definition we have to consider all quadrangles (up to translation by $`^2`$) with vertices $`p_i`$ such that for every $`i`$ the vector $`p_ip_{i1}`$ has slope $`\lambda _i`$, $`p_iP_{i,i+1}mod^2`$, for $`i=0,1,2`$, and the piecewise linear path $`[p_0,p_1,p_2,p_3]`$ goes in the clockwise direction. First of all, it is easy to check that the condition (2.6) implies that all such quadrangles are convex. Secondly, the condition on the orientation of the path is equivalent to the system of inequalities
(2.7)
$$det(p_{i+1}p_i,p_ip_{i1})>0$$
where $`i/4`$. These quadrangles (considered up to translations by $`^2`$) can be parametrized by elements of a rank-$`2`$ lattice. Namely, let
$$\mathrm{\Lambda }=\mathrm{\Lambda }(\lambda _0,\mathrm{},\lambda _3)=\{𝐧=(n_0,\mathrm{},n_3)^4:n_i=\lambda _in_i=0,n_1I_{\lambda _1},n_2I_{\lambda _2}\}.$$
Then writing
$$p_ip_{i1}=x_i(1,\lambda _i)$$
for $`i/4`$ we obtain the vector $`𝐱=(x_0,\mathrm{},x_3)`$ in $`\mathrm{\Lambda }_{}`$. The inequalities (2.7) become
(2.8)
$$(\lambda _i\lambda _{i+1})x_ix_{i+1}>0$$
for $`i/4`$. On the other hand, setting
$$P_{i,i+1}P_{i1,i}=v_i(1,\lambda _i)$$
we obtain the vector $`𝐯=(v_0,\mathrm{},v_i)\mathrm{\Lambda }_{}`$ (note that $`v_i=y_{i,i+1}y_{i1,i}`$). Now the conditions $`p_iP_{i,i+1}mod^2`$ for $`i=0,1,2`$ imply that $`𝐱𝐯`$ belongs to $`\mathrm{\Lambda }`$. Conversely, given an element $`𝐧=(n_0,\mathrm{},n_3)\mathrm{\Lambda }`$ we have the corresponding quadrangle $`\mathrm{\Delta }(𝐧)`$ with vertices $`p_i`$ such that $`p_0=P_{0,1}`$ and $`p_ip_{i1}=(v_i+n_i)(1,\lambda _i)`$. Fixing the fourth vertex $`p_3`$ modulo $`^2`$ is equivalent to choosing $`𝐧`$ in a fixed coset modulo the sublattice $`\mathrm{\Lambda }^+\mathrm{\Lambda }`$ defined as follows:
$$\mathrm{\Lambda }^+=\mathrm{\Lambda }^+(\lambda _0,\mathrm{},\lambda _3)=\{𝐧=(n_0,\mathrm{},n_3)^4:n_i=\lambda _in_i=0,\lambda _in_i\}.$$
Thus, the sums in the definition of the Fukaya coefficients are taken over all elements $`𝐧`$ of a coset of $`\mathrm{\Lambda }^+`$ in $`\mathrm{\Lambda }`$, such that $`𝐯+𝐧C`$, where $`C\mathrm{\Lambda }_{}`$ is an open subset defined by inequalities (2.8). It is easy to see that the condition (2.6) implies that $`C`$ is a non-empty open cone. The relation between the vertex $`p_3`$ and the coset $`𝐧\mathrm{\Lambda }/\mathrm{\Lambda }^+`$ can be found explicitly as follows. We know that $`p_0p_3=(v_0+n_0)(1,\lambda _0)`$, and that $`p_0P_{01}`$. Hence, $`p_3P_{03}(y_0,y_3+a\lambda _3+b)`$ where $`a`$ and $`b`$ are integers satisfying
$$\frac{a\lambda _3+b}{\lambda _3\lambda _0}n_0mod(I_{\lambda _0}).$$
It follows that
$$(a,b)(n_1+n_2,\lambda _1n_1\lambda _2n_2)mod\mathrm{\Lambda }(\lambda _0,\lambda _3).$$
The area of $`\mathrm{\Delta }(𝐧)`$ is given by
$$_{\mathrm{\Delta }(𝐧)}𝑑xdy=\frac{1}{2}\left(det(p_1p_0,p_0p_3)+det(p_3p_2,p_2p_1)\right)=\frac{1}{2}Q(𝐯+𝐧)$$
where $`Q`$ is the quadratic form on $`\mathrm{\Lambda }_{}`$ given by
$$Q(𝐱)=(\lambda _0\lambda _1)x_0x_1+(\lambda _2\lambda _3)x_2x_3.$$
Finally, we have
$$\underset{i/4}{}(x(p_i)x(p_{i1}))t_i=\underset{i}{}t_ix_i=𝐰𝐱$$
where $`𝐰=(w_0,\mathrm{},w_3)\mathrm{\Lambda }_{}`$, $`w_i=t_{i,i+1}t_{i1,i}`$, $`t_{i,j}=\frac{t_jt_i}{\lambda _j\lambda _i}`$, $`𝐱𝐱^{}`$ is the symmetric pairing induced by $`Q`$ (so that $`Q(𝐱)=𝐱𝐱`$).
Thus, we have
$`m_3([P_{01}(y_0,y_1)],[P_{12}(y_1,y_2)],[P_{23}(y_2,y_3)])=`$
$`{\displaystyle \underset{𝐧\mathrm{\Lambda },𝐯+𝐧C}{}}\pm \mathrm{exp}(\pi i\tau Q(𝐯+𝐧)2\pi i𝐰(𝐯+𝐧))[P_{03}(y_0,y_3+(n_1+n_2)\lambda _3\lambda _1n_1\lambda _2n_2)]`$
where the sign is equal to the sign of $`v_0+n_0`$. Let us choose $`C^+`$ to be the component of $`C`$ where $`x_0>0`$. Then using the bases $`e_{ij}(m,n)`$ defined by (2.4) we can rewrite the above formula as follows:
$`m_3(e_{01}(0,0),e_{12}(0,0),e_{23}(0,0))=`$
$`{\displaystyle \underset{𝐧\mathrm{\Lambda }/\mathrm{\Lambda }^+}{}}\mathrm{\Theta }_{\mathrm{\Lambda }^+,Q,C;𝐧}(\tau 𝐯𝐰,\tau )e_{03}(n_1+n_2,\lambda _1n_1\lambda _2n_2)`$
where $`\mathrm{\Theta }_{\mathrm{\Lambda }^+,Q,C;𝐧}`$ is the indefinite theta series with characteristic $`𝐧`$ defined by (1.1). Similarly we can compute products of all basis elements:
$`m_3(e_{01}(a,b),e_{12}(c,d),e_{23}(f,g))=`$
(2.9) $`{\displaystyle \underset{𝐧\mathrm{\Lambda }/\mathrm{\Lambda }^+}{}}\mathrm{\Theta }_{\mathrm{\Lambda }^+,Q,C;𝐮+𝐧}(\tau 𝐯𝐰,\tau )e_{03}(a+c+f+n_1+n_2,b+d+g\lambda _1n_1\lambda _2n_2)`$
where $`𝐮=(u_0,u_1,u_2,u_3)\mathrm{\Lambda }_{}`$ is the following vector:
$`𝐮=({\displaystyle \frac{a\lambda _1+b}{\lambda _1\lambda _0}}{\displaystyle \frac{(a+c+f)\lambda _3+b+d+g}{\lambda _3\lambda _0}},{\displaystyle \frac{c\lambda _2+d}{\lambda _2\lambda _1}}{\displaystyle \frac{a\lambda _0+b}{\lambda _1\lambda _0}},{\displaystyle \frac{f\lambda _3+g}{\lambda _3\lambda _2}}{\displaystyle \frac{c\lambda _1+d}{\lambda _2\lambda _1}},`$
$`{\displaystyle \frac{(a+c)\lambda _0+f\lambda _3+b+d+g}{\lambda _3\lambda _0}}{\displaystyle \frac{f\lambda _3+g}{\lambda _3\lambda _2}}).`$
The $`A_{\mathrm{}}`$-axiom (2.1) can be converted into a certain identity for indefinite theta series (and the usual theta functions which appear from $`m_2`$). The explicit formula can be found in the last section of .
It is convenient for explicit computations to choose a pair of components $`(x_i,x_j)`$ as coordinates on $`\mathrm{\Lambda }_{}`$ in such a way that the cone $`C`$ defined by inequalities (2.8) coincides with the cone $`x_ix_j>0`$. We will do this in two particular cases.
1. First assume that $`\lambda _1<\lambda _0\lambda _2<\lambda _3`$. Choose $`(x_0,x_1)`$ as coordinates in $`\mathrm{\Lambda }_{}`$. Then we have $`C=\{𝐱:x_0x_1>0\}`$. The quadratic form in these coordinates can be written as
$$Q(𝐱)=ax_0^2+2bx_0x_1+cx_1^2,$$
where
$$a=\frac{(\lambda _2\lambda _0)(\lambda _3\lambda _0)}{\lambda _3\lambda _2},$$
$$b=\frac{(\lambda _2\lambda _1)(\lambda _3\lambda _0)}{\lambda _3\lambda _2},$$
$$c=\frac{(\lambda _2\lambda _1)(\lambda _3\lambda _1)}{\lambda _3\lambda _2}.$$
If $`\lambda _0=\lambda _2`$ (in this case we say that this is a trapezoid triple product) then $`a=0`$, so the form $`Q`$ splits over $``$. On the other hand, if we set $`\lambda _0=0`$ then the corresponding transformation $`(\lambda _1,\lambda _2,\lambda _3)(a,b,c)`$ is birational. More precisely, the inverse transformation is given by
$$\lambda _1=\frac{D}{b},$$
$$\lambda _2=\frac{aD}{b(ba)},$$
$$\lambda _3=\frac{D}{cb}$$
where $`D=b^2ac`$. The inequalities satisfied by $`\lambda _i`$ are equivalent to the following inequalities for $`a,b,c`$:
$$D<0a<b<c.$$
2. Now assume that $`\lambda _2<\lambda _0<\lambda _3<\lambda _1`$. Choose $`(x_0,x_3)`$ as coordinates in $`\mathrm{\Lambda }_{}`$. Then $`C=\{𝐱:x_0x_3>0\}`$ and
$$Q(𝐱)=ax_0^2+2bx_0x_3+cx_3^2,$$
where
$$a=\frac{(\lambda _1\lambda _0)(\lambda _0\lambda _2)}{\lambda _1\lambda _2},$$
$$b=\frac{(\lambda _1\lambda _0)(\lambda _3\lambda _2)}{\lambda _1\lambda _2},$$
$$c=\frac{(\lambda _1\lambda _3)(\lambda _3\lambda _2)}{\lambda _1\lambda _2}.$$
Setting $`\lambda _0=0`$ we get a birational transformation coinciding with the previous one up to permutation of variables and signs. So the inverse map is given by $`\lambda _1=D/(bc)`$, $`\lambda _2=aD/b(ab)`$, $`\lambda _3=D/b`$. The inequalities for $`\lambda _i`$ are equivalent to the following inequalities:
$$D<0<a,c<b.$$
It follows that every indefinite theta series associated with rational quadratic form of signature $`(1,1)`$ appears as a coefficient of certain Fukaya triple product. Indeed, let us consider the $``$-valued quadratic form $`Q(x,y)=ax^2+2bxy+cy^2`$ on $`^2`$ such that $`D=b^2ac>0`$ and $`Q`$ is positive on the cone $`C_0=\{(x,y)^2:xy>0\}`$. Assume first that $`ac0`$ and $`bc`$. Then we necessarily have that $`a`$, $`b`$ and $`c`$ are positive and either $`a<b`$ or $`c<b`$. Permuting the coordinates if necessary we can assume that $`a<b`$. Then either $`c>b`$ and we are in the situation of the case 1 above or $`c<b`$ so we can apply the case 2. Note that the lattice coming from the configuration of lines will be commensurable with $`^2`$, so we can use formulas of section 1 to relate the indefinite theta series associated with $`(^2,Q,C_0)`$ to the corresponding Fukaya triple product. Furthermore, by rescaling the coordinates $`x`$, $`y`$ we can always achieve that we are in the situation of a given case. For example, for the proof of Theorem 1 we will use a rescaling which leads to the case 1. If $`a=0`$ and $`bc`$ we can still apply the formulas of either case 1 or case 2. If $`b=c`$ then one of the slopes will be infinite. The only reason why we didn’t include vertical lines in our category was because they correspond to torsion sheaves on elliptic curves while we want to deal only with bundles. However, the Fukaya compositions are well-defined for all lines including vertical, so a slight modification of the above computation will work in this case. Alternatively, we can always rescale the coordinates in such a way that $`bc`$ and then apply the above formulas. On the other hand, we notice that replacing the form $`Q`$ by $`NQ`$ for some $`N>0`$ we can always achive that all the slopes $`\lambda _i`$ are integers.
Remark. For computations with the form $`Q`$ on $`\mathrm{\Lambda }_{}`$ defined above the following formula is useful:
$$𝐱𝐲=(\lambda _i\lambda _{i+1})x_iy_{i+1}+(\lambda _{i+2}\lambda _{i+3})y_{i+2}x_{i+3},$$
for any $`i/4`$, $`𝐱,𝐲\mathrm{\Lambda }_{}`$.
## 3. Morphisms of vector bundles on elliptic curves
We identify an elliptic curve $`E`$ with its dual by associating to a point $`xE`$ the line bundle $`𝒫_x=𝒪_E(xe)`$ of degree zero on $`E`$, where $`eE`$ is the neutral element of the group law. For every integer $`d`$ we denote by $`E_d`$ the kernel of the homomorphism $`[d]:EE:xdx`$.
###### Proposition 3.1.
Let $`V_1`$, $`V_2`$ and $`V_3`$ be stable vector bundles on an elliptic curve $`E`$. Assume that $`V_i`$ has rank $`r_i`$ and degree $`d_i`$ and that the slopes $`\mu _i=d_i/r_i`$ satisfy $`\mu _1<\mu _2<\mu _3`$. Then there exists an integer $`d`$ depending only on $`(d_i)`$ and $`(r_i)`$ such that the following natural map is surjective
$$_{xE_d}\mathrm{Hom}(V_1,V_2𝒫_x)\mathrm{Hom}(V_2𝒫_x,V_3)\mathrm{Hom}(V_1,V_3).$$
Proof. Consider the action of $`E\times E`$ on the category of vector bundles on $`E`$, such that a point $`(x,y)`$ acts as the functor $`T_{(x,y)}:Ft_x^{}F𝒫_y`$. Then the statement of the theorem can be reformulated as follows: there exists a finite subgroup $`SE\times E`$ such that the map
(3.1)
$$_{sS}\mathrm{Hom}(V_1,T_s(V_2))\mathrm{Hom}(T_s(V_2),V_3)\mathrm{Hom}(V_1,V_3)$$
is surjective. Indeed, this follows from the fact that for any $`xE`$ one has
$$t_{r_2x}^{}V_2V_2𝒫_{d_2x}.$$
Now we claim that in proving the surjectivity of (3.1) we can replace the bundles $`V_i`$ by $`V_iL`$, where $`L`$ is a line bundle, or by $`\mathrm{SS}(V_i)`$ provided that $`d_1>0`$, where $`\mathrm{SS}`$ is the Fourier-Mukai transform (see ). Indeed, in the first case this is clear. In the case of the Fourier-Mukai transform this follows from the fact that $`\mathrm{SS}`$ interchanges translations with tensoring by line bundles of degree zero.
Using these two operations (tensoring with a line bundle and the Fourier-Mukai transform) we can make $`V_1=𝒪_E`$. Next we want to reduce the proof to the case when $`V_2`$ is a line bundle. Indeed, assume that in this case the assertion is true. Then consider an isogeny $`\pi :E^{}E`$ of degree $`r_2`$ and a line bundle $`L`$ on $`E^{}`$ such that $`\pi _{}LV_2`$ (such $`\pi `$ and $`L`$ always exist). By assumption the statement is true for the triple $`(𝒪_E^{},L,\pi ^{}V_3)`$ on $`E^{}`$ (since $`\pi ^{}V_3`$ is a direct sum of stable bundles), hence there exists $`d`$ such that the map
(3.2)
$$_{xE_d^{}}H^0(E^{},L𝒫_x)\mathrm{Hom}(L𝒫_x,\pi ^{}V_3)H^0(E^{},\pi ^{}V_3)$$
is surjective. Now we notice that for every $`xE^{}`$ there is a natural commutative diagram
(3.3)
$$\begin{array}{ccc}H^0(E^{},L𝒫_x)\mathrm{Hom}(L𝒫_x,\pi ^{}V_3)& \text{}& H^0(E^{},\pi ^{}V_3)\\ \text{}& & \text{}\\ H^0(E,\pi _{}(L𝒫_x))\mathrm{Hom}(\pi _{}(L𝒫_x),V_3)& \text{}& H^0(E,V_3)\end{array}$$
in which the right vertical arrow is the following composition of natural morphisms:
$$H^0(E^{},\pi ^{}V_3)\stackrel{~}{}H^0(E,\pi _{}\pi ^{}V_3)H^0(E,V_3).$$
Since $`V_3`$ is a direct summand in $`\pi _{}\pi ^{}V_3`$, this map is surjective. Also for every $`yE`$ we have an isomorphism
$$\pi _{}(L𝒫_{\widehat{\pi }(x)})(\pi _{}L)𝒫_x$$
where $`\widehat{\pi }:EE^{}`$ is the isogeny dual to $`\pi `$. Hence, the surjectivity of (3.2) implies the surjectivity of the following map
$$_{y\widehat{\pi }^1(E_d^{})}H^0(E,(\pi _{}L)𝒫_y)\mathrm{Hom}((\pi _{}L)𝒫_y,V_3)H^0(E,V_3)$$
as required.
It remains to prove the statement for a triple $`(𝒪_E,L,V)`$ where $`L`$ is a line bundle, $`V`$ is a stable bundle, such that $`0<\mathrm{deg}(L)<\mu (V)`$. Note that $`H^0(E,V)`$ is an irreducible representation of the Heisenberg group $`H`$ which is an extension of $`E_d`$ by $`𝔾_m`$, where $`d=\mathrm{deg}V`$. More precisely, $`H`$ is the group of pairs $`(x,\varphi )`$ where $`xE_d`$, $`\varphi :Vt_x^{}V`$. It follows that the image of the natural map
$$_{xE_d}H^0(E,t_x^{}L)\mathrm{Hom}(t_x^{}L,V)H^0(E,V)$$
is invariant under the $`H`$-action. Therefore, it suffices to prove that this map is not zero. Let $`f:LV`$ be a non-zero morphism. Then it is an injection of sheaves, hence the induced morphism $`H^0(E,L)H^0(E,V)`$ is injective which finishes the proof. ∎
Using the equivalence of the Fukaya category of $`^2/^2`$ (without higher products) with the category of vector bundles on $`E`$ we deduce the following corollary.
###### Corollary 3.2.
Let $`(L,t)`$ (resp. $`(L^{},t^{})`$) be an object in the Fukaya category, where $`L`$ (resp. $`L^{}`$) is a line of slope $`\lambda `$ (resp. $`\lambda ^{}`$). Assume that $`\lambda <\mu <\lambda ^{}`$, where $`\mu `$. Then there exists a finite number of objects $`(M_i,t_i)`$ in Fukaya category, where $`M_i`$ are lines of slope $`\mu `$, such that the composition
$$m_2:_i\mathrm{Hom}((L,t),(M_i,t_i))\mathrm{Hom}((M_i,t_i),(L^{},t^{}))\mathrm{Hom}((L,t),(L^{},t^{}))$$
is surjective. Furthermore, one can choose $`(M_i,t_i)`$ in a generic position (i.e. in such a way that $`M_i`$ do not pass through a finite number of given points).
## 4. Expression of all triple products via trapezoid ones
In this section we use the $`A_{\mathrm{}}`$-axiom (2.1) to get a simple expression of an arbitrary triple product in $`_s`$ (corresponding to a transversal configuration of lines) in terms of trapezoid triple products and all double products.
First, consider the triple product $`m_3(r,s,t)`$ where $`r\mathrm{Hom}^1(L_0,L_1)`$, $`s\mathrm{Hom}^0(L_1,L_2)`$, $`t\mathrm{Hom}^0(L_2,L_3)`$, $`\lambda _0>\lambda _1<\lambda _2<\lambda _3`$, $`\lambda _0<\lambda _3`$. If $`\lambda _2=\lambda _0`$ then this is a trapezoid product. Otherwise, there are two possibilities:
a) $`\lambda _2>\lambda _0`$. In this case using Corollary 3.2 we can write $`s`$ as a linear combination of products $`s^{}s^{\prime \prime }`$ where $`s^{}\mathrm{Hom}^0(L_1,L_0^{})`$, $`s^{\prime \prime }\mathrm{Hom}^0(L_0^{},L_2)`$, $`L_0^{}`$ is a line of slope $`\lambda _0`$, $`L_0^{}L_0mod(^2)`$. Then applying the $`A_{\mathrm{}}`$-constraint to the quadruple $`r,s^{},s^{\prime \prime },t`$ we obtain
$$m_3(r,s^{}s^{\prime \prime },t)=m_3(r,s^{},s^{\prime \prime })t+m_3(r,s^{},s^{\prime \prime }t)$$
(note that $`rs^{}=0`$ by assumption while $`m_3(s^{},s^{\prime \prime },t)=0`$ as an element of $`\mathrm{Hom}^1(L_1,L_3)`$).
b) $`\lambda _2<\lambda _0`$. In this case we write $`t`$ as a linear combination of products $`t^{}t^{\prime \prime }`$ where $`t^{}\mathrm{Hom}^0(L_2,L_0^{})`$, $`t^{\prime \prime }\mathrm{Hom}^0(L_0^{},L_3)`$, $`L_0^{}`$ is a line of slope $`\lambda _0`$, $`L_0^{}L_0mod(^2)`$. Then we have $`m_3(r,s,t^{})\mathrm{Hom}^0(L_0,L_0^{})=0`$. Hence, applying the $`A_{\mathrm{}}`$-constraint to the quadruple $`r,s,t^{},t^{\prime \prime }`$ we get
$$m_3(r,s,t^{}t^{\prime \prime })=m_3(rs,t^{},t^{\prime \prime })+m_3(r,st^{},t^{\prime \prime })$$
(note that $`m_3(s,t^{},t^{\prime \prime })=0`$).
One deals similarly with products of the type
$$\mathrm{Hom}^0(L_0,L_1)\mathrm{Hom}^0(L_1,L_2)\mathrm{Hom}^1(L_2,L_3)\mathrm{Hom}^0(L_0,L_3)$$
where $`\lambda _0<\lambda _1<\lambda _2>\lambda _3`$.
Now let us consider $`m_3(r,s,t)`$ where $`r\mathrm{Hom}^0(L_0,L_1)`$, $`s\mathrm{Hom}^1(L_1,L_2)`$, $`t\mathrm{Hom}^0(L_2,L_3)`$, $`\lambda _0<\lambda _1>\lambda _2<\lambda _3`$, $`\lambda _0<\lambda _3`$. If $`\lambda _1=\lambda _3`$ then this is a trapezoid product. Otherwise, there are two possibilities:
a) $`\lambda _1>\lambda _3`$. Then we can write $`r`$ as a linear combination of products $`r^{}r^{\prime \prime }`$, where $`r^{}\mathrm{Hom}^0(L_0,L_3^{})`$, $`r^{\prime \prime }\mathrm{Hom}^0(L_3^{},L_1)`$, $`L_3^{}`$ is a line of slope $`\lambda _3`$, $`L_3^{}L_3mod(^2)`$. Applying $`A_{\mathrm{}}`$-constraint to $`(r^{},r^{\prime \prime },s,t)`$ we get
$$m_3(r^{}r^{\prime \prime },s,t)=m_3(r^{},r^{\prime \prime }s,t)m_3(r^{},r^{\prime \prime },st)+m_3(r^{},r^{\prime \prime },s)t$$
(since $`m_3(r^{\prime \prime },s,t)\mathrm{Hom}^0(L_3^{},L_3)=0`$). The first two terms in the RHS are trapezoid, while the product $`m_3(r^{},r^{\prime \prime },s)`$ is of the form considered before.
a) $`\lambda _1<\lambda _3`$. In this case we can write $`t`$ as a linear combination of products $`t^{}t^{\prime \prime }`$, where $`t^{}\mathrm{Hom}^0(L_2,L_1^{})`$, $`t^{\prime \prime }\mathrm{Hom}^0(L_1^{},L_3)`$, $`L_1^{}`$ is a line of slope $`\lambda _1`$, $`L_1^{}L_1mod(^2)`$. Applying $`A_{\mathrm{}}`$-constraint to $`(r,s,t^{},t^{\prime \prime })`$ we get
$$m_3(r,s,t^{}t^{\prime \prime })=m_3(rs,t^{},t^{\prime \prime })+rm_3(s,t^{},t^{\prime \prime })+m_3(r,s,t^{})t^{\prime \prime }$$
(since $`st^{}\mathrm{Hom}^1(L_1,L_1^{})=0`$). The products $`m_3(s,t^{},t^{\prime \prime })`$ and $`m_3(r,s,t^{})`$ are trapezoid, while the product $`m_3(rs,t^{},t)`$ is of the form considered before.
Finally, using the cyclic symmetry of $`m_3`$ we can reduce all non-zero transversal higher products $`m_3`$ to the ones considered above. For example, the product
$$\mathrm{Hom}^1(L_0,L_1)\mathrm{Hom}^1(L_1,L_2)\mathrm{Hom}^0(L_2,L_3)\mathrm{Hom}^1(L_0,L_3)$$
is equivalent to the product
$$\mathrm{Hom}^1(L_1,L_2)\mathrm{Hom}^0(L_2,L_3)\mathrm{Hom}^0(L_3,L_0)\mathrm{Hom}^0(L_1,L_0).$$
Proof of Theorem 1.
First, let us introduce some notation. We assume that for all objects $`(L,t)`$ of the Fukaya category that appear below a representative of the line $`L`$ modulo $`^2`$-translations is fixed. Then for every pair of objects $`(L_1,t_1)`$ and $`(L_2,t_2)`$ and every intersection point $`P(L_1+^2)(L_2+^2)/^2`$ we represent $`P`$ in the form $`P_{12}(y_1,y_2+m\lambda _2+n)`$, where $`L_i=L(\lambda _i,y_i)`$, and set
$$e(P)=e_{12}(m,n)$$
where $`e_{12}(m,n)`$ is defined by formula (2.4).
As was explained in section 2.3 (after passing to a commensurable lattice $`\mathrm{\Lambda }`$) we can assume that the lattice $`\mathrm{\Lambda }`$, the quadratic form $`Q`$, and the cone $`C`$ come from a quadruple of rational numbers $`\lambda _0,\mathrm{},\lambda _3`$ such that $`\lambda _1<\lambda _0<\lambda _2<\lambda _3`$. Let us represent the variable $`𝐳\mathrm{\Lambda }_{}`$ in the form $`𝐳=\tau 𝐯𝐰`$ with $`𝐯,𝐰\mathrm{\Lambda }_{}`$. Then the value of $`\mathrm{\Theta }_{\mathrm{\Lambda },Q,C}(𝐳,\tau )`$ appears as the coefficient with $`e(P_{03})`$ of the triple product $`m_3(e(P_{01}),e(P_{12}),e(P_{23}))`$ for a quadruple of objects $`(L_i,t_i)`$, $`i=0,\mathrm{},3`$, and intersection points $`P_{i,i+1}L_iL_{i+1}`$, where $`L_0=L(\lambda _0,0)`$, $`L_1=L(\lambda _1,0)`$, $`L_2=(\lambda _2,(\lambda _2\lambda _1)v_1)`$, $`L_3=(\lambda _3,(\lambda _0\lambda _3)v_0)`$, $`t_0=t_1=0`$, $`t_2=(\lambda _2\lambda _1)w_1`$, $`t_3=(\lambda _0\lambda _3)w_0`$. Now we can apply the above procedure of expressing this triple product in terms of the trapezoid ones. More precisely, due to the inequalities $`\lambda _1<\lambda _0<\lambda _2<\lambda _3`$ we apply the very first case of the above argument. This means that we choose a finite number of (not necessarily distinct) objects $`(M_j,t_j)`$, where $`M_j`$ are lines of slope $`\lambda _0`$ different from $`L_0`$, and intersection points $`Q_jL_1(M_j+^2)`$, $`R_jM_j(L_2+^2)`$ such that $`m_2(e(Q_j),e(R_j))`$ form a basis in $`\mathrm{Hom}((L_1,t_1),(L_2,t_2))`$. The transition matrix from this basis to the standard basis of intersection points of $`L_1`$ and $`L_2`$ modulo $`^2`$ is given by elliptic functions of $`(z_1,\tau )`$. Thus, we can write
$$e(P_{12})=\underset{j}{}\varphi _jm_2(e(Q_j),e(R_j))$$
where $`\varphi _j`$ are meromorphic elliptic functions of $`(z_1,\tau )`$. Applying the $`A_{\mathrm{}}`$-identity (2.1) we get
$`m_3(e(P_{01}),e(P_{12}),e(P_{23}))=`$
$`{\displaystyle \underset{j}{}}\rho _j(m_3(e(P_{01}),e(Q_j),e(R_j))e(P_{23})+m_3(e(P_{0,1}),e(Q_j),e(R_j)e(P_{23}))).`$
Now by the results of section 2.3 and by formula (1.2) the coefficients of the trapezoid product $`m_3(e(P_{0,1}),e(Q_j),e(R_j))`$ are given (up to factors of the form $`\mathrm{exp}(\pi ia\tau )`$ with $`a_{>0}`$) by the functions of the form $`\kappa _e(a\tau +b,gz_1+c\tau +d;h\tau )`$, where $`a,b,c,d,e,g`$, $`h_{\mathrm{\hspace{0.25em}0}}`$, Multiplying the result with $`e(P_{2,3})`$ means that we get some linear combination of the above functions with coefficients which are elliptic functions of $`(s(𝐳),\tau )`$ for some linear functional $`s`$. Finally the products $`m_3(e(P_{0,1}),e(Q_j),e(R_j)e(P_{2,3}))`$ are expressed via elliptic functions of $`(s(𝐳),\tau )`$ and the functions of the form $`\kappa _e(a\tau +b,gz_2+c\tau +d;h\tau )`$. ∎
Remarks. 1. Since after rescaling $`Q`$ the slopes $`\lambda _i`$ can always be chosen to be integers one can replace the reference to Proposition 3.1 in the above proof by the well-known surjectivity statement for morphisms between line bundles.
2. It may seem strange that in Theorem 1 we substitute only constants in the second argument of $`\kappa `$. However, the function $`\kappa `$ satisfies some identities (see , p. 481, formula (45) and the next one, or , formula (3.4.3)) which imply that one can express $`\kappa (y,x;\tau )`$ in terms of $`\tau `$-elliptic functions and the function $`\kappa (c,x+yc;\tau )`$ for any $`c+\tau `$.
## 5. Example
In this section we will give an example of identity produced by Theorem 1. Let us fix $`a`$ such that $`a2`$ and consider the quadratic form $`Q`$ on $`^2`$ given by
$$Q(n_0,n_1)=an_0^2+4an_0n_1+(4a2)n_1^2.$$
Let us also consider the following split quadratic forms:
$$Q^1(n_0,n_1)=2(n_0+n_1)n_1$$
on $`^2`$ and
$$Q^2(n_0,n_1)=(2n_0+\frac{2a1}{a}n_1)n_1$$
on the lattice $`\mathrm{\Lambda }^2=\{(n_0,n_1):n_0,n_1a\}`$. As a cone $`C`$ in all three cases we choose $`x_0x_1>0`$ (with $`x_0>0`$ in $`C^+`$) and set
$$\mathrm{\Theta }(z_0,z_1)=\mathrm{\Theta }_{^2,Q,C}(z_0,z_1;\tau ),$$
$$\mathrm{\Theta }^1(z_0,z_1)=\mathrm{\Theta }_{^2,Q^1,C}(z_0,z_1;\tau ),$$
$$\mathrm{\Theta }_{c_0,c_1}^2(z_0,z_1)=\mathrm{\Theta }_{\mathrm{\Lambda }^2,Q^2,C;(c_0,c_1)}(z_0,z_1;\tau )$$
(we omit the variable $`\tau `$ in notation for brevity). Let us also denote by $`\mathrm{\Delta }(z)`$ the determinant of the $`2\times 2`$ matrix
$$(\theta _{2,i}(z+\frac{j}{2},\frac{\tau }{2}))_{i/2,j/2}.$$
Then we have the following identity:
$`\mathrm{\Delta }(z_1{\displaystyle \frac{\tau }{2}})\mathrm{\Theta }(z_0,z_1)=\theta _{2,1}(z_1{\displaystyle \frac{\tau 1}{2}},{\displaystyle \frac{\tau }{2}})\times `$
$`\{\theta _a(z_0+2z_1,{\displaystyle \frac{\tau }{a}})\mathrm{\Theta }^1(2z_1,{\displaystyle \frac{\tau }{2}})+{\displaystyle \underset{l/a}{}}\theta _{a,l}(z_02z_1+{\displaystyle \frac{\tau }{2a}},{\displaystyle \frac{\tau }{a}})\mathrm{\Theta }_{\frac{a1}{a}l,l}^2(z_0,{\displaystyle \frac{\tau }{2}})\}`$
$`\theta _{2,1}(z_1{\displaystyle \frac{\tau }{2}},{\displaystyle \frac{\tau }{2}})\times `$
$`\{\theta _a(z_0+2z_1,{\displaystyle \frac{\tau }{a}})\mathrm{\Theta }^1(2z_1,{\displaystyle \frac{\tau 1}{2}})+{\displaystyle \underset{l/a}{}}\theta _{a,l}(z_02z_1+{\displaystyle \frac{\tau 1}{2a}},{\displaystyle \frac{\tau }{a}})\mathrm{\Theta }_{\frac{a1}{a}l,l}^2(z_0,{\displaystyle \frac{\tau 1}{2}})\}.`$
Let us write the variables in the form $`z_i=v_i\tau w_i`$, where $`i=0,1`$, and consider the following objects in the Fukaya category: $`(L_0=L(0,0),t_0=0)`$, $`(L_1=L(1,0),t_1=0)`$, $`(L_2=L(1,2v_1),t_2=2w_1)`$, $`(L_3=L(\frac{a}{a1},\frac{a}{a1}v_0),t_3=\frac{a}{a1}w_0)`$. Then the series $`\mathrm{\Theta }(z_0,z_1)`$ is equal to the coefficient with $`e_{03}`$ in the triple product $`m_3(e_{01},e_{12},e_{23})`$ where we set $`e_{ij}=e_{ij}(0,0)`$. Now we consider two auxiliary objects in the Fukaya category: $`(M_0=L(0,\frac{1}{2}),0)`$ and $`(M_1=L(0,\frac{1}{2}),\frac{1}{2}dx)`$. There are unique points of intersection $`Q_iL_1M_i`$, $`R_iM_iL_2`$, $`i=0,1`$. Note that the points $`Q_0`$ and $`Q_1`$ (resp. $`R_0`$ and $`R_1`$) coincide but we denote them differently since they belong to morphism spaces between different objects in the Fukaya category. The first step is to represent $`e_{12}`$ as a linear combination of $`e(Q_0)e(R_0)`$ and $`e(Q_1)e(R_1)`$. We have
$$e(Q_i)e(R_i)=\theta _2(z_1\frac{\tau i}{2},\frac{\tau }{2})e_{12}+\theta _{2,1}(z_1\frac{\tau }{2},\frac{\tau i}{2})e_{12}(0,1)$$
for $`i=0,1`$. Hence,
$$\mathrm{\Delta }(z_1\frac{\tau }{2})e_{12}=\theta _{2,1}(z_1\frac{\tau 1}{2},\frac{\tau }{2})e(Q_0)e(R_0)\theta _{2,1}(z_1\frac{\tau }{2},\frac{\tau }{2})e(Q_1)e(R_1).$$
Next we use the formula
$$m_3(e_{01},e(Q_i)e(R_i),e_{23})=m_3(e_{01},e(Q_i),e(R_i))e_{23}+m_3(e_{01},e(Q_i),e(R_i)e_{23}).$$
We have
$$m_3(e_{01},e(Q_i),e(R_i))=\mathrm{\Theta }^1(2z_1,\frac{\tau i}{2})e_{02},$$
$$e(R_i)e_{23}=\underset{l/a}{}\theta _{a,l}(z_02z_1+\frac{\tau i}{2a},\frac{\tau }{a})e_{03}^i(0,l)$$
where $`i=1,2`$, $`e_{03}^i(0,l)`$ are the basis elements in $`\mathrm{Hom}(M_i,L_3)`$ defined by (2.4). The coefficient with $`e_{03}`$ in the product $`e_{02}e_{03}`$ is equal to
$$\theta _a(z_0+2z_1,\frac{\tau }{a}).$$
One more computation shows that the coefficient with $`e_{03}`$ in the product $`m_3(e_{01},e(Q_i),e_{03}^i(0,l))`$ (where $`i=1,2`$) is equal to
$$\mathrm{\Theta }_{\frac{a1}{a}l,l}^2(z_0,\frac{\tau i}{2}).$$
Combining all these calculations we get the identity above.
## 6. Massey products
In this section we consider a family of well-defined univalued triple Massey products which define global sections of certain line bundles on the second cartesian power of the universal curve over the moduli stack of elliptic curves with some level structure.
### 6.1. Definition of triple Massey products
Let $`V_i`$, $`0i3`$ be holomorphic vector bundles on a complex manifold. Let $`\alpha _1\mathrm{Hom}(V_0,V_1)`$, $`\alpha _2\mathrm{Ext}^1(V_1,V_2)`$, $`\alpha _3\mathrm{Hom}(V_2,V_3)`$ be elements satisfying $`\alpha _2\alpha _1=0`$, $`\alpha _3\alpha _2=0`$. Below we recall two equivalent constructions of the triple Massey product $`MP(\alpha _1,\alpha _2,\alpha _3)`$ which belongs to the cokernel of the morphism
(6.1)
$$\mathrm{Hom}(V_0,V_2)\mathrm{Hom}(V_1,V_3)\mathrm{Hom}(V_0,V_3):(\beta _1,\beta _2)\alpha _3\beta _1+\beta _2\alpha _1.$$
Let us represent $`\alpha _2`$ by a $`\overline{}`$-closed $`(0,1)`$-form $`\stackrel{~}{\alpha }_2`$ with values in $`V_1^{}V_2`$. Then by our assumption we have
$$\stackrel{~}{\alpha }_2\alpha _1=\overline{}(\alpha _{12}),$$
$$\alpha _3\stackrel{~}{\alpha }_2=\overline{}(\alpha _{23})$$
for some sections $`\alpha _{12}C^{\mathrm{}}(V_1^{}V_2)`$, $`\alpha _{23}C^{\mathrm{}}(V_2^{}V_3)`$. Now we set
$$MP(\alpha _1,\alpha _2,\alpha _3)=\alpha _3\alpha _{12}\alpha _{23}\alpha _1.$$
The ambiguity in a choice of $`\alpha _{12}`$ and $`\alpha _{23}`$ precisely means that $`MP(\alpha _1,\alpha _2,\alpha _3)`$ is correctly defined modulo the image of the map (6.1).
In the second definition <sup>2</sup><sup>2</sup>2This definition is a particular case of the general construction of Massey products in triangulated categories, cf. ,IV.2 we consider an extension
$$0V_2\stackrel{i}{}V\stackrel{p}{}V_10$$
with the class $`\alpha _2`$. By our assumption there exist morphisms $`\alpha _1^{}:V_0V`$ and $`\alpha _3^{}:VV_3`$ such that
$$\alpha _1=p\alpha _1^{},$$
$$\alpha _3=\alpha _3^{}i.$$
Now the composition $`\alpha _3^{}\alpha _1^{}\mathrm{Hom}(V_0,V_3)`$ is well-defined modulo the image of (6.1).
###### Proposition 6.1.
One has
$$MP(\alpha _1,\alpha _2,\alpha _3)=\alpha _3^{}\alpha _1^{}$$
in the cokernel of the map (6.1).
Proof. A choice of a closed $`(0,1)`$-form $`\stackrel{~}{\alpha }_2`$ representing $`\alpha _2`$ leads to the choice of $`V`$ as follows. We set $`V=V_2V_1`$ as a $`C^{\mathrm{}}`$-bundle and define a holomorphic structure on it by the following $`\overline{}`$-operator:
$$\overline{}_V=\left(\begin{array}{cc}\overline{}_{V_2}& \stackrel{~}{\alpha }_2\\ 0& \overline{}_{V_1}\end{array}\right).$$
Let $`\sigma :V_1V`$ be the natural $`C^{\mathrm{}}`$-splitting arising from this description (so $`p\sigma =\mathrm{id}_{V_1}`$). Then $`\overline{}(\sigma )=i\stackrel{~}{\alpha }_2`$. Define $`\rho :VV_2`$ by the condition
$$i\rho =\sigma p\mathrm{id}_V.$$
Then we have
$$\overline{}(\rho )=\stackrel{~}{\alpha }_2p.$$
Thus, we can choose
$$\alpha _{12}=\rho \alpha _1^{},$$
$$\alpha _{23}=\alpha _3^{}\sigma .$$
Hence,
$$\alpha _3\alpha _{12}\alpha _{23}\alpha _1=\alpha _3^{}i\rho \alpha _1^{}\alpha _3^{}\sigma p\alpha _1^{}=\alpha _3^{}\alpha _1^{}.$$
We will be interested in a particular case when $`\mathrm{Hom}(V_0,V_2)=\mathrm{Hom}(V_1,V_3)=0`$. In this case the Massey product $`MP(\alpha _1,\alpha _2,\alpha _3)`$ is an element of $`\mathrm{Hom}(V_0,V_3)`$. Homological mirror conjecture for elliptic curve $`E=E_\tau `$ (proven in for transversal products) implies that $`MP(\alpha _1,\alpha _2,\alpha _3)`$ is equal to the corresponding triple Fukaya product $`m_3(\alpha _1,\alpha _2,\alpha _3)`$. Indeed, the corresponding products are homotopic but in our case the homotopy takes values in zero spaces (see , sec.1.1, for a more detailed explanation).
Remark. The notion of transversality considered in has to be strengthened, since the definition of Fukaya products $`m_k`$ given in section 2.1 requires that no three of the corresponding circles intersect in one point (this was overlooked in ). However, all the proofs of can be easily modified accordingly.
### 6.2. Massey products for line bundles
Let us fix a quadruple of integers $`(d_0,d_1,d_2,d)`$ such that $`1d_0<\mathrm{min}(d_1,d_2)`$, $`d_0+d=d_1+d_2`$. Let $`_0`$, $`_1`$, $`_2`$ and $``$ be line bundles on elliptic curve $`E`$ of degrees $`d_0`$, $`d_1`$, $`d_2`$ and $`d`$ respectively, such that $`_0_1_2`$ (here and below we skip the sign of tensor product between line bundles for brevity). Then for any pair of sections $`s_1H^0(E,_1)`$, $`s_2H^0(E,_2)`$ and an element $`eH^1(^1)=H^1(E,_0_1^1_2^1)`$ such that the compositions $`s_1eH^1(E,_0_2^1)`$ and $`es_2H^1(E,_0_1^1)`$ are zero, the triple Massey product $`MP(s_1,e,s_2)`$ defined in 6.1 is an element of $`H^0(E,_0)`$ (since $`H^0(E,_0_1^1)=H^0(E,_0_2^1)=0`$). Thus, if we denote by $`K_{s_1,s_2}H^1(E,^1)`$ the kernel of the natural map
$$H^1(E,^1)H^1(E,_0_1^1)H^1(E,_0_2^1):e(s_1e,es_2)$$
then the Massey product defines a linear map
(6.2)
$$MP:K_{s_1,s_2}H^0(E,_0).$$
Assume that $`s_10`$, $`s_20`$. Let $`D`$ be the divisor of common zeroes of $`s_1`$ and $`s_2`$. Let us denote $`_i^{}=_i(D)`$ for $`i=0,1,2`$, $`^{}=(D)`$, so that we still have $`_0^{}^{}_1^{}_2^{}`$. Then we have the induced sections $`s_1^{}H^0(E,_1^{})`$ and $`s_2^{}H^0(E,_2^{})`$ and a canonical map
$$\phi :K_{s_1,s_2}K_{s_1^{},s_2^{}}$$
induced by the morphism $`^1^1(D)=(^{})^1`$. On the other hand, we can consider $`H^0(E,_0^{})=H^0(E,_0(D))`$ as a subspace in $`H^0(E,_0)`$.
###### Lemma 6.2.
For any $`eK_{s_1,s_2}`$ one has
$$MP(s_1,e,s_2)=MP(s_1^{},\phi (e),s_2^{}).$$
Proof. First let us notice that
$$H^0(E,_0^{}(_i^{})^1)=H^0(E,_0_i^1)=0$$
for $`i=1,2`$, so $`MP(s_1^{},\phi (e),s_2^{})`$ is defined. Let
$$0^1E𝒪_E0$$
be an extension representing $`e`$. Let $`f:_1^1E`$ be the lifting of $`s_1:_1^1𝒪_E`$. Then the image of $`f`$ belongs to the following subextension $`E^{}`$:
$$0^1E^{}𝒪_E(D)0.$$
Note that this extension represents the class $`\phi (e)`$. Let $`g^{}:E^{}_0_1^1(D)`$ be the lifting of the map $`s_2^{}:^1_0_1^1(D)`$. Then according to Proposition 6.1 we have
$$MP(s_1^{},\phi (e),s_2^{})=g^{}f.$$
On the other hand, the push-out of the extension $`E`$ by the morphism $`^1^1(D)`$ coincides with extension
$$0^1(D)E^{}(D)𝒪_E0.$$
Thus, we have an embedding $`i:EE^{}(D)`$, such that $`i|_E^{}:E^{}E^{}(D)`$ is the natural map. In particular, we can take $`g=g^{}i:E_0_1^1`$ as the lifting of the map $`s_2:^1_0_1^1`$. Applying Proposition 6.1 again we obtain
$$MP(s_1,e,s_2)=gf$$
which finishes the proof. ∎
Thus, it suffices to study the case when the sections $`s_1`$ and $`s_2`$ have no common zeroes. In this case one has an exact sequence
(6.3)
$$0_1^1_2^1\stackrel{\alpha }{}_1^1_2^1\stackrel{\beta }{}𝒪_E0$$
where $`\alpha (s)=(ss_2,ss_1)`$, $`\beta (t_1,t_2)=t_1s_1+t_2s_2`$. Tensoring by $`_0`$ and considering the corresponding sequence of cohomologies we get the following exact complex $`C`$:
$$0H^0(E,_0)\stackrel{\delta _1}{}H^1(E,^1)\stackrel{\delta _0}{}H^1(E,_0_1^1)H^1(E,_0_2^1)0$$
###### Proposition 6.3.
For any $`tH^0(E,_0)`$ one has
$$MP(s_1,\delta _1(t),s_2)=t.$$
Proof. The class $`\delta _1(t)\mathrm{Ext}^1(_0^1,_1^1_2^1)`$ is represented by the extension
$$0_1^1_2^1E_0^10$$
where $`E=\beta ^1(t(_0^1)).`$ The morphism $`_0^1_1^1\stackrel{s_1}{}_0^1`$ lifts to the morphism
$$f:_0^1_1^1E:s(st,0),$$
while the morphism $`_1^1_2^1\stackrel{s_2}{}_1^1`$ lifts to the morphism
$$g:E_1^1:(t_1,t_2)t_1.$$
According to Proposition 6.1 we have
$$MP(s_1,\delta _1(t),s_2)=gf,$$
so the result follows from the above formulas for $`f`$ and $`g`$. ∎
The spaces $`K_{s_1,s_2}`$ can be considered as stalks of the sheaf $`𝒦`$ over the appropriate moduli stack. The Massey product gives a morphism from $`𝒦`$ to the bundle with the fibre $`H^0(E,_0)`$. The previous proposition shows that over an open part this is an isomorphism inverse to $`\delta _1`$. In the section 7.1 we will compute the map $`MP`$ in terms of indefinite theta series using (2.3) and applying the above proposition we will get a proof of Theorem 4. In order to get modular (or Jacobi) forms from the coefficients of the map $`MP`$ we can try to map various standard line bundles (trivialized by theta functions) to $`𝒦`$. Below we will present two ways to do it. The first uses the determinants and gives Jacobi forms. The second approach (see section 7.2) is more direct and produces modular forms but it requires additional assumptions on the integers $`(d_0,d_1,d_2)`$.
### 6.3. Determinantal approach
Let us consider the $`1`$-dimensional vector space
$$M=detH^1(E,^1)detH^1(E,_0_1^1)^{}detH^1(E,_0_2^1)^{}$$
where for a vector space $`V`$ we denote by $`detV`$ its top-degree wedge power. We claim that for fixed $`s_1`$ and $`s_2`$ there is a canonical map
$$e_{s_1,s_2}:M^{d_01}H^1(E,^1)^{}K_{s_1,s_2}$$
Indeed, in general for a linear map $`f:VW`$ between vector spaces of dimensions $`n`$ and $`nk`$ one can construct a linear map
$$k_f:detVdetW^{}^{k1}V^{}V$$
such that its image belongs to $`\mathrm{ker}(\varphi )`$ as follows. Start with the morphism
$$^{nk}f^{}:detW^{}=^{nk}W^{}^{nk}V^{}$$
and then consider the following composition
$$\varphi :detW^{}^{k1}V^{}^{nk}V^{}^{k1}V^{}^{n1}V^{}$$
where the second arrow is given by the wedge product. It remains to use the isomorphism $`^{n1}V^{}VdetV^{}`$. It is easy to see that the image of $`k_f`$ belongs to $`\mathrm{ker}(f)`$. If $`f`$ is not surjective then $`k_f=0`$, otherwise, $`k_f`$ surjects onto $`\mathrm{ker}(f)`$. The map $`e_{s_1,s_2}`$ is defined by applying this construction to the map
$$\delta _0:H^1(E,^1)H^1(E,_0_2^1)H^1(E,_0_1^1)$$
induced by compositions with $`s_1`$ and $`s_2`$ (note that dimensions of the spaces are $`d_1+d+2d_0`$ and $`d_1+d_22d_0`$). Composing $`e_{s_1,s_2}`$ with the map (6.2) above we obtain a linear map
$$MP_{s_1,s_2}^{det}:M^{d_01}H^1(E,^1)^{}H^0(E,_0):\xi MP(s_1,e_{s_1,s_2}(\xi ),s_2).$$
It is easy to see that if the sections $`s_1`$ and $`s_2`$ have a common zero then $`\delta _0`$ is not surjective, hence, $`MP_{s_1,s_2}^{det}=0`$.
We refer to , Appendix A, for the definition of determinants of complexes.
###### Proposition 6.4.
Assume that the divisors of $`s_1`$ and $`s_2`$ do not intersect. Then the map $`\pm MP_{s_1,s_2}^{det}`$ is equal to the composition
$$M^{d_01}H^1(E,^1)^{}M^{d_01}H^0(E,_0)^{}H^0(E,_0)$$
where the first arrow is induced by $`\delta _1^{}`$ and the second arrow is induced by the isomorphism $`MdetH^0(E,_0)`$ given by the canonical trivialization of the determinant of the exact complex $`C`$.
Proof. This follows from Proposition 6.3 and from the following observation. For an exact complex of the form
$$0K\stackrel{\delta _1}{}V\stackrel{\delta _0}{}W0$$
the map
$$k_{\delta _0}:detVdetW^{}^{dimK1}V^{}K$$
constructed above coincides (up to a sign) with the composition
$$detVdetW^{}^{dimK1}V^{}\stackrel{\delta _0^{}}{}detVdetW^{}^{dimK1}K^{}K$$
where $`detVdetW^{}detK`$ by the canonical trivialization of the determinant of this complex. ∎
###### Corollary 6.5.
Assume that $`d_0=1`$. Fix some bases in the spaces $`H^1(E,^1)`$, $`H^1(E,_0_i^1)`$ and $`H^0(E,_0)`$. Then we have
$$MP_{s_1,s_2}^{det}=\pm det(C,B)^1$$
where $`B`$ is the corresponding basis of the complex $`C`$.
### 6.4. Serre duality in homological mirror symmetry
The identification of $`\mathrm{Ext}^1`$-spaces in the Fukaya category and the category of bundles on elliptic curve $`E=E_\tau `$ uses the Serre duality in the following form (see ):
$$\mathrm{Hom}(V_1,V_2)\mathrm{Ext}^1(V_2,V_1):f(gd\overline{z})_E𝑑z\mathrm{Tr}(fgd\overline{z}).$$
In other words, this is a composition of the natural map
$$\mathrm{Hom}(V_1,V_2)\mathrm{Ext}^1(V_2,V_1)H^1(E,𝒪_E)$$
and the map
$$\varphi :H^1(E,𝒪_E):\alpha _E𝑑z\alpha $$
where $`\alpha `$ is a $`(0,1)`$-form. The map $`\varphi `$ is in turn equal to the composition of the isomorphism
$$H^1(E,𝒪_E)H^1(E,\omega _E)$$
induced by the holomorphic $`1`$-form $`\omega _0=2\pi idz`$ and the functional
$$I:H^1(E,\omega _E):\eta \frac{1}{2\pi i}_E\eta $$
where $`\eta `$ is a $`(1,1)`$-form. The factor of $`\frac{1}{2\pi i}`$ in the definition of $`I`$ is important because then $`I`$ admits an algebraic definition (in particular, it is defined over the field of definition of $`E`$). In fact, we can define the functional $`I:H^1(C,\omega _C)`$ for any Riemann surface $`C`$ by the same formula. Then the following property of $`I`$ shows that it is algebraically defined.
###### Lemma 6.6.
For every point $`pC`$ let $`e_pH^1(C,\omega _C)`$ be the class defined by the boundary homomorphism $`H^0(C,𝒪_p)H^1(C,\omega _C)`$ coming from the exact sequence
$$0\omega _C\omega _C(p)\stackrel{\mathrm{Res}_p}{}𝒪_p0.$$
Then $`I(e_p)=1`$.
Proof. Let $`UC`$ be an open disk containing $`p`$. Consider the Cech complex of $`\omega _C`$ associated with the covering $`(U,Cp)`$:
$$𝒞^{}:\omega _C(U)\omega _C(Cp)\omega _C(Up)$$
where the differential sends $`(\alpha _U,\alpha _p)`$ to $`\alpha _U|_{Up}\alpha _p|_{Up}`$. Since $`H^1(Cp,\omega _C)=H^1(U,\omega _C)=0`$ the complex $`𝒞^{}`$ computes the cohomology of $`\omega _C`$. The residue map
$$\mathrm{Res}_p:𝒞^1=\omega _C(Up)$$
descends to the functional $`I^{}:H^1(C,\omega _C)`$. It suffices to prove that $`I=I^{}`$. To this end let us consider the Cech complex associated with the same covering and with the complex of sheaves $`\mathrm{\Omega }^{1,0}\overline{}\mathrm{\Omega }^{1,1}`$ on $`C`$:
$$𝒞𝒟^{}:\mathrm{\Omega }^{1,0}(U)\mathrm{\Omega }^{1,0}(Cp)\stackrel{d_1}{}\mathrm{\Omega }^{1,1}(U)\mathrm{\Omega }^{1,1}(Cp)\mathrm{\Omega }^{1,0}(Up)\stackrel{d_2}{}\mathrm{\Omega }^{1,1}(Up)$$
where
$$d_1(\alpha _U,\alpha _p)=(\overline{}\alpha _U,\overline{}\alpha _p,\alpha _U|_{Up}\alpha _p|_{Up}),$$
$$d_2(\eta _U,\eta _p,\beta )=\eta _U|_{Up}\eta _p|_{Up}\overline{}\beta .$$
The complex $`𝒞𝒟^{}`$ is concentrated in degrees $`[0,2]`$. We have natural morphisms of complexes $`𝒞^{}𝒞𝒟^{}`$ and $`\mathrm{\Omega }^{1,}(C)𝒟^{}`$ inducing isomorphisms on cohomologies. Now we define the functional $`\stackrel{~}{I}:𝒞𝒟^1`$ by the formula
$$\stackrel{~}{I}(\eta _U,\eta _p,\beta )=\frac{1}{2\pi i}_D\eta _U\frac{1}{2\pi i}_{CD}\eta _p+\frac{1}{2\pi i}_D\beta $$
where $`DU`$ is a smaller disk containing $`p`$. It is easy to check that $`\stackrel{~}{I}d_1=0`$, hence, $`\stackrel{~}{I}`$ descends to a functional
$$\stackrel{~}{I}:H^1(C,\omega _C)H^1(𝒞𝒟^{}).$$
If $`\eta `$ is a global $`(1,1)`$-form then
$$\stackrel{~}{I}(\eta |_U,\eta |_{Cp},0)=\frac{1}{2\pi i}_C\eta ,$$
hence $`\stackrel{~}{I}=I`$. On the other hand, if $`\alpha `$ is a holomorphic $`1`$-form on $`Up`$ then
$$\stackrel{~}{I}(0,0,\alpha )=\mathrm{Res}_p(\alpha ),$$
hence $`\stackrel{~}{I}=I^{}`$. ∎
### 6.5. Boundary homomorphism
Below we will need to calculate the matrices of $`\delta _1`$ and $`\delta _0`$ with respect to the standard bases of the terms of the complex $`C`$ (coming from Serre duality and the bases of theta functions in the spaces of global sections) for any pair of sections $`s_1,s_2`$ with no common zeroes. The components of $`\delta _0`$ are just the compositions with $`s_1`$ and $`s_2`$ so they are given by theta functions. On the other hand, the map
$$\delta _1:H^0(E,_0)H^1(E,^1)$$
is the composition with the class $`\delta (s_1,s_2)H^1(E,_1^1_2^1)`$ corresponding to the extension (6.3). Via Serre duality $`\delta _1`$ corresponds to a bilinear form
$$B_{s_1,s_2}:H^0(E,_0)H^0(E,)H^0(E,_1_2)$$
induced by the functional on $`H^0(E,_1_2)`$ dual to the class $`\delta (s_1,s_2)`$ (recall that we always use the trivialization of $`\omega _E`$ given by the form $`2\pi idz`$).
###### Lemma 6.7.
Assume that $`s_1`$ and $`s_2`$ have no common zeroes. Then one has
$$B_{s_1,s_2}(s,t)=\underset{xZ(s_1)}{}\mathrm{Res}_x(\frac{2\pi is(z)t(z)dz}{s_1(z)s_2(z)})$$
where $`Z(s_1)`$ is the divisor of zeroes of $`s_1`$.
Proof. Applying the octahedron axiom to the composition of arrows $`_1^1_1^1_2^1\stackrel{\beta }{}𝒪_E`$ one can easily derive that the class $`\delta (s_1,s_2)`$ can be represented as the following composition:
$$𝒪_E𝒪_{Z(s_1)}_1^1_2^1[1]$$
where the first arrow is the canonical one, the second arrow comes from the exact triangle
$$_1^1_2^1\stackrel{s_1}{}_2^1𝒪_{Z(s_1)}_1^1_2^1[1]$$
(here we use the natural trivialization of $`_2|_{Z(s_1)}`$ induced by $`s_2`$). In other words, this element in $`\mathrm{Ext}^1(𝒪_{Z(s_1)},_1^1_2^1)`$ corresponds to the functional on $`H^0(E,_1_2|_{Z(s_1)})`$ which at the point $`xZ(s_1)`$ is equal to
$$\mathrm{Res}_x(\frac{2\pi idz}{s_1(z)s_2(z)})$$
## 7. Calculations
In this section we will use homological mirror symmetry to relate the Massey products considered in section 6.2 to indefinite theta series. We keep the notations of the previous section.
### 7.1. Proof of Theorem 4
Let us fix a pair of complex numbers $`v_1\tau w_1`$ and $`v_2\tau w_2`$ where $`v_i,w_i`$. Now consider the following $`4`$ objects in the Fukaya category: $`(L_0=L(0,0),0)`$, $`(L_1=L(d_1,d_1v_1),d_1w_1)`$, $`(L_2=L(d_0d_2,d_2v_2),d_2w_2)`$, $`(L_3=L(d_0,0),0)`$. Let $`_\tau `$ be the basic line bundle on $`E=E_\tau `$ such that $`\theta (z)=\theta (z,\tau )`$ is a section of $`_\tau `$. Under the equivalence with the category of bundles on $`E`$ our $`4`$ objects correspond to $`𝒪_E`$, $`_1`$, $`_0_2^1`$ and $`_0`$ respectively, where $`_1=t_{v_1\tau w_1}^{}_\tau ^{d_1}`$, $`_2=t_{v_2\tau +w_2}^{}_\tau ^{d_2}`$, $`_0=_\tau ^{d_0}`$. Let $`\mathrm{\Lambda }^+=\mathrm{\Lambda }^+(0,d_1,d_0d_2,d_0)`$ be the corresponding rank $`2`$ lattice. Let us denote by $`𝐱=(x_0,x_1,x_2,x_3)`$ the element of $`\mathrm{\Lambda }^+`$ such that $`x_0=v_1\tau w_1`$ and $`x_3=v_2\tau w_2`$ (thus, $`x_1=\frac{(dd_1)x_0+d_2x_3}{d}`$, $`x_2=\frac{(dd_2)x_3+d_1x_0}{d}`$).
We start by fixing bases in all the relevant vector spaces. According to section 2.2 we have a natural isomorphism
$$\mathrm{Hom}(L_0,L_1)H^0(E,_1)$$
which identifies the basis $`(e_{01}(0,a),a/d_1)`$ with $`(\theta _{d_1,a}(z+x_0,\frac{\tau }{d_1}),a/d_1)`$. Similarly,
$$\mathrm{Hom}(L_2,L_3)H^0(E,_2)$$
such that $`e_{23}(0,a)`$, $`a/d_2`$, corresponds to $`\theta _{d_2,a}(zx_3,\frac{\tau }{d_2})`$. On the other hand,
$$\mathrm{Hom}^1(L_1,L_2)H^1(E,^1)H^0(E,)^{}$$
in such a way that the basis $`(e_{12}(0,a),a/d)`$ is dual to the basis $`(\theta _{d,a}(z+x_0+x_1,\frac{\tau }{d}),a/d)`$. Finally, the space $`\mathrm{Hom}^1(L_0,L_2)H^1(E,_0_2^1)`$ (resp. $`\mathrm{Hom}^1(L_1,L_3)H^1(E,_0_1^1)`$) has the basis $`(e_{02}(0,a),a/(d_2d_0))`$ (resp. $`(e_{13}(0,a),a/(d_1d_0))`$) and the space $`\mathrm{Hom}(L_0,L_3)H^0(E,_0)`$ has the basis $`(e_{03}(0,a),a/d_0)`$ identified with $`(\theta _{d_0,a}(z,\frac{\tau }{d_0}),a/d_0)`$.
We are going to compute explicitly the map (6.2) for $`s_1=\theta _{d_1}(z+x_0,\frac{\tau }{d_1})H^0(E,_1)`$ and $`s_2=\theta _{d_2}(zx_3,\frac{\tau }{d_2})H^0(E,_2)`$ (the corresponding morphisms in the Fukaya category are $`e_{01}=e_{01}(0,0)`$ and $`e_{23}=e_{23}(0,0)`$).
First, using the formula (2.2) one can easily compute the matrix of $`\delta _0`$. Namely, we have
$`m_2(e_{01},e_{12}(0,k))={\displaystyle \underset{n/I_1}{}}\theta _{I_1,\frac{k}{d}n}(p_1x_1,p_1\tau )e_{02}(0,k+dn),`$
$`m_2(e_{12}(0,k),e_{23})={\displaystyle \underset{n/I_2}{}}\theta _{I_2,\frac{k}{d}+n}(p_2x_2,p_2\tau )e_{13}(0,k+dn),`$
where $`I_i=\frac{dd_i}{d_i}`$, $`p_i=\frac{d_id}{dd_i}`$, $`i=1,2`$. In these formulas we fix a representative of $`k`$ in $``$. The change of a representative corresponds to a shift of the summation variable $`n`$. Here is a better way to write these formulas which doesn’t require a choice of a representative of $`k`$:
$`m_2(e_{01},e_{12}(0,k))={\displaystyle \underset{n/dI_1,nk(d)}{}}\theta _{I_1,\frac{n}{d}}(p_1x_1,p_1\tau )e_{02}(0,n),`$
(7.1) $`m_2(e_{12}(0,k),e_{23})={\displaystyle \underset{n/dI_2,nk(d)}{}}\theta _{I_2,\frac{n}{d}}(p_2x_2,p_2\tau )e_{13}(0,n),`$
Note that the coefficient of $`m_2(e_{01},e_{12}(0,k))`$ (resp. $`m_2(e_{12}(0,k),e_{23})`$) with $`e_{02}(0,i)`$ (resp. $`e_{13}(0,j)`$) is non-zero only if $`ki(d,dd_1)`$ (resp. $`kj(d,dd_2)`$). Now it is easy to deduce that an element $`e=_{k/d}c_ke_{12}(0,k)`$ (where $`c_k`$) belongs to the subspace $`K_{s_1,s_2}=\mathrm{ker}(\delta _0)`$ if and only if
$`{\displaystyle \underset{n/dI_1,ni(dd_1)}{}}c_n\theta _{I_1,\frac{n}{d}}(p_1x_1,p_1\tau )=0,`$
(7.2) $`{\displaystyle \underset{n/dI_2,nj(dd_2)}{}}c_n\theta _{I_2,\frac{n}{d}}(p_2x_2,p_2\tau )=0`$
for all $`i/(dd_1)`$ and $`j/(dd_2)`$.
Assume first that $`x_1`$ and $`x_2`$ are sufficiently generic, so that the corresponding circles in $`^2/^2`$ form a transversal configuration. Then we can use (2.3) to compute the relevant triple Fukaya products. Note that the projection $`(n_0,n_1,n_2,n_3)(n_0,n_3)`$ defines an isomorphism of the lattice $`\mathrm{\Lambda }^+`$ with the lattice
$$\mathrm{\Lambda }_{d_1,d_2,d}=\{(m,n)^2:d_1md_2n(d)\}.$$
On the other hand, the projection $`(n_0,n_1,n_2,n_3)(n_1,n_2)`$ maps $`\mathrm{\Lambda }^+`$ onto $`\mathrm{\Lambda }_{d_1,d_2,d_0}`$. Thus, we have
$`m_3(e_{01},e_{12}(0,k),e_{23})=`$
$`{\displaystyle \underset{(n_1,n_2)^2/\mathrm{\Lambda }_{d_1,d_2,d_0}}{}}\mathrm{\Theta }_{\mathrm{\Lambda }_{d_1,d_2,d},Q,mn>0;(\frac{k+d_2n_2d_1n_1}{d_0}n_1,\frac{k+d_2n_2d_1n_1}{d_0}n_2)}(x_0,x_3;\tau )`$
$`e_{03}(0,k+d_2n_2d_1n_1)`$
for $`k/d`$, where the form $`Q`$ is given by
$$Q(m,n)=\frac{1}{d}(d_1(dd_1)m^2+2d_1d_2mn+(dd_2)d_2n^2)$$
Let $`(e_{03}^{}(0,l),0l<d_0)`$ be the dual basis to $`(e_{03}(0,l))`$. Then for $`e=_kc_ke_{12}(0,k)`$ we have
$$\stackrel{~}{F}_l:=e_{03}^{}(0,l),MP(s_1,e,s_2)=\underset{k}{}c_ke_{03}^{}(0,l),m_3(e_{01},e_{12}(0,k),e_{23})=\underset{k}{}C_{l,k}c_k$$
where $`C_{l,k}`$ is zero unless $`kl(d_0,d_1,d_2)`$ in which case
$$C_{l,k}=\mathrm{\Theta }_{\mathrm{\Lambda }_{d_1,d_2,d},Q,mn>0;(\frac{l}{d_0}in_1,\frac{l}{d_0}+in_2)}(x_0,x_3;\tau )$$
where $`n_1`$, $`n_2`$, and $`i`$ are some integers satisfying
$$k+d_2n_2d_1n_1=l+d_0i.$$
Denoting $`m_1=in_1`$, $`m_2=in_2`$ we can rewrite this formula as follows:
(7.3)
$$C_{l,k}=\mathrm{\Theta }_{\mathrm{\Lambda }_{d_1,d_2,d},Q,mn>0;(\frac{l}{d_0}+m_1,\frac{l}{d_0}+m_2)}(x_0,x_3;\tau )$$
where $`(m_1,m_2)^2/\mathrm{\Lambda }_{d_1,d_2,d}`$ is defined by the congruence
$$d_2m_2d_1m_1kl(d).$$
Thus, we have
$`\stackrel{~}{F}_l={\displaystyle \underset{(m_1,m_2)^2/\mathrm{\Lambda }_{d_1,d_2,d}}{}}c_{d_2m_2d_1m_1+l}\mathrm{\Theta }_{\mathrm{\Lambda }_{d_1,d_2,d},Q,mn>0;(\frac{l}{d_0}+m_1,\frac{l}{d_0}+m_2)}(x_0,x_3;\tau )=`$
$`{\displaystyle \underset{(m,n)^2}{}}ϵ(m,n)c_{d_2nd_1m+l}\mathrm{exp}(\pi i\tau Q(m{\displaystyle \frac{l}{d_0}},n+{\displaystyle \frac{l}{d_0}})2\pi i[d_1x_1(m{\displaystyle \frac{l}{d_0}})+d_2x_2(n+{\displaystyle \frac{l}{d_0}})])=`$
$`{\displaystyle \underset{(m,n)^2}{}}ϵ(m,n)c_{d_1md_2n+l}\mathrm{exp}(\pi i\tau Q(m+{\displaystyle \frac{l}{d_0}},n{\displaystyle \frac{l}{d_0}})+2\pi i[d_1x_1(m+{\displaystyle \frac{l}{d_0}})+d_2x_2(n{\displaystyle \frac{l}{d_0}})])`$
where $`ϵ(m,n)=0`$ unless $`(m\frac{l}{d_0}+\alpha (x_0))(n+\frac{l}{d_0}+\alpha (x_3))>0`$ in which case $`ϵ(m,n)=\mathrm{sign}(m\frac{l}{d_0}+\alpha (x_0))`$. Now an easy computation shows that the conditions (7.1) are equivalent to the system of equations
$`{\displaystyle \underset{m}{}}c_{d_1md_2n_0+l}\mathrm{exp}(\pi i\tau Q(m+{\displaystyle \frac{l}{d_0}},n_0{\displaystyle \frac{l}{d_0}})+2\pi id_1x_1(m+{\displaystyle \frac{l}{d_0}}))=0,`$
$`{\displaystyle \underset{n}{}}c_{d_1m_0d_2n+l}\mathrm{exp}(\pi i\tau Q(m_0+{\displaystyle \frac{l}{d_0}},n{\displaystyle \frac{l}{d_0}})+2\pi id_2x_2(n{\displaystyle \frac{l}{d_0}}))=0,`$
where $`m_0,n_0`$, $`0l<d_0`$. Therefore, in the notation of Theorem 4 we obtain
$$\stackrel{~}{F}_l=\underset{(m,n)^2}{}ϵ(m,n)c_{d_1md_2n+l}a_{m,n,l},$$
and the above system is equivalent to the condition (0.4). This implies that one can replace the summation scheme defining $`\stackrel{~}{F}_l`$ to the summation over $`m,n0`$ and $`m,n<0`$ (with signs “minus” and “plus”, respectively). Hence, we obtain $`\stackrel{~}{F}_l=F_l`$ (the latter series is defined in the formulation of Theorem 4), i.e.,
(7.4)
$$e_{03}^{}(0,l),MP(s_1,e,s_2)=F_l.$$
Now assume that the sections $`s_1`$ and $`s_2`$ have no common zeroes. Then we claim that formula (7.4) holds without any further genericity assumption on $`x_1`$ and $`x_2`$. Indeed, as we have seen in section 6.2, when the data $`(_0,_1,_2,s_1,s_2)`$ vary in such a way that $`s_1`$ and $`s_2`$ have no common zeroes, the vector spaces $`K_{s_1,s_2}`$ can be viewed as fibers of a vector bundle on the space of parameters. Furthermore, the map $`MP:K_{s_1,s_2}H^0(E,_0)`$ varies continuously with parameters. Since $`F_l`$ is also a continuous function of $`x_1`$, $`x_2`$ and $`(c_k)`$ varying in the vector bundle defined by (0.4), we derive that equation (7.4) holds whenever $`s_1`$ and $`s_2`$ have no common zeroes.
Finally, we are going to combine the result of section 6.5 with Proposition 6.3 to derive the system of equations for $`F_l`$. We have
$$\delta _1(e_{03}(0,l))=\underset{k/d}{}D_{k,l}e_{12}(0,k),$$
where
$$D_{k,l}=B_{s_1,s_2}(\theta _{d_0,l}(z,\frac{\tau }{d_0}),\theta _{d,k}(z+x_0+x_1,\frac{\tau }{d})).$$
The divisor $`Z(s_1)`$ consists of the points $`(x_0+z(a),a/d_1)`$ where
$$z(a)=\frac{\tau }{2}+\frac{1}{2d_1}+\frac{a}{d_1}$$
Therefore, according to Lemma 6.7 one has
$`D_{k,l}={\displaystyle \underset{a/d_1}{}}\mathrm{Res}_{x_0+z(a)}({\displaystyle \frac{2\pi i\theta _{d_0,l}(z,\frac{\tau }{d_0})\theta _{d,k}(z+x_0+x_1,\frac{\tau }{d})dz}{\theta _{d_1}(z+x_0,\frac{\tau }{d_1})\theta _{d_2}(zx_3,\frac{\tau }{d_2})}})=`$
$`{\displaystyle \underset{a/d_1}{}}{\displaystyle \frac{2\pi i\theta _{d_0,l}(x_0+z(a),\frac{\tau }{d_0})\theta _{d,k}(x_1+z(a),\frac{\tau }{d})}{\theta _{d_1}^{}(z(a),\frac{\tau }{d_1})\theta _{d_2}(x_1+x_2+z(a),\frac{\tau }{d_2})}}=`$
$`{\displaystyle \frac{2\pi i}{d_1\theta _{}^{}(\frac{d_1\tau +1}{2},d_1\tau )}}{\displaystyle \underset{a/d_1}{}}{\displaystyle \frac{\theta _{d_0,l}(\frac{(d_2d)x_1+d_2x_2}{d_0}+z(a),\frac{\tau }{d_0})\theta _{d,k}(x_1+z(a),\frac{\tau }{d})}{\theta _{d_2}(x_1+x_2+z(a),\frac{\tau }{d_2})}},`$
where we denote $`\theta _{}^{}(z,\tau )=\frac{}{z}\theta _{}(z,\tau ).`$ It is easy to see that this formula is equivalent to the formula for $`D_{k,l}`$ in Theorem 4. Now Proposition 6.3 implies that
$$e=\delta _1(\underset{l}{}F_le_{03}(0,l))=\underset{k,l}{}D_{k,l}F_le_{12}(0,k).$$
Equating the coefficients with $`e_{12}(0,k)`$ we get the system of linear equations on $`F_l`$. It remains to notice that disjointness of the divisors of $`s_1`$ and $`s_2`$ is equivalent to the condition (0.3) of Theorem 4.
### 7.2. Examples of modular indefinite theta series
Let us assume that $`(dd_i)|d`$ for $`i=1,2`$. Choose an integer $`f`$ such that $`f|d`$ and $`(dd_i)|f`$ for $`i=1,2`$ and set
$$c_k=\{\begin{array}{cc}1,k0(f),\hfill & \\ 0,k0(f)\hfill & \end{array}$$
Then the conditions (7.1) boil down to the following pair of equations:
$`\theta _{\frac{f}{d}}(p_1x_1,p_1\tau )=0,`$
$`\theta _{\frac{f}{d}}(p_2x_2,p_2\tau )=0.`$
Hence, we can set
(7.5)
$$x_i=s_i\frac{f\tau }{2d}+r_i\frac{dd_i}{2fd_i},i=1,2,$$
where $`s_i`$, $`r_i`$ are odd integers, to satisfy these equations. Now Theorem 4 implies that the series
$$F_0=\underset{(m,n)^2,d_1md_2n(f)}{\overset{indef}{}}\mathrm{exp}(\pi i\tau Q(m,n)+2\pi i(d_1x_1m+d_2x_2n))$$
multiplied by the appropriate factor $`\mathrm{exp}(\pi ic\tau )`$ with $`c`$, is modular. More precisely, we can apply Theorem 4 directly only if the condition (0.3) is satisfied, i.e. the corresponding section $`s_1`$ and $`s_2`$ have no common zeroes. Otherwise, we first use Lemma 6.2 to reduce to this case.
Let us denote by $`f_0`$ the least common multiple of $`dd_1`$ and $`dd_2`$. Notice that the congruence $`d_1md_2n0(f)`$ implies that $`(dd_1)m`$ and $`(dd_2)n`$ belong to $`f_0`$. So we can take $`\frac{(dd_1)m}{f_0}`$ and $`\frac{(dd_2)n}{f_0}`$ as the new summation variables. Then we get
$`F_0={\displaystyle \underset{(m,n)^2,mn(\frac{f}{f_0})}{\overset{indef}{}}}\mathrm{exp}(\pi i\tau {\displaystyle \frac{f_0^2}{d}}({\displaystyle \frac{d_1}{dd_1}}m^2+{\displaystyle \frac{2d_1d_2}{(dd_1)(dd_2)}}mn+{\displaystyle \frac{d_2}{dd_2}}n^2)+`$
$`\pi i\tau {\displaystyle \frac{ff_0}{d}}({\displaystyle \frac{d_1}{dd_1}}s_1m+{\displaystyle \frac{d_2}{dd_2}}s_2n)+\pi i{\displaystyle \frac{f_0}{f}}(r_1m+r_2n))`$
To find the factor $`\mathrm{exp}(\pi ic\tau )`$ above recall that for $`N>0`$ the functions
$$\mathrm{exp}(\pi i\tau N\lambda ^2)\theta _{N,i}(\lambda \tau +\mu ,\frac{\tau }{N})$$
where $`i`$, $`\lambda ,\mu `$, $`\lambda >0`$, and
$$\mathrm{exp}(\pi i\tau N/4)\theta _{}^{}(\frac{N\tau +1}{2},N\tau )$$
are modular. Hence,
$$\mathrm{exp}(\pi i\tau (d_0(\frac{f}{2dd_0}((d_2d)s_1+d_2s_2)+\frac{1}{2})^2+d(\frac{f}{2d}s_1+\frac{1}{2})^2d_2(\frac{f}{2d}(s_1+s_2)+\frac{1}{2})^2\frac{d_1}{4}))D_{k,l}$$
are modular. Simplifying we conclude that
(7.6)
$$\mathrm{exp}(\pi i\tau \frac{f^2}{4d^2d_0}(d_1(d_2d)s_1^2+2d_1d_2s_1s_2+d_2(d_1d)s_2^2))F_0$$
is modular.
Proof of Theorem 2. Let $`a,b,c,p`$ be positive integers such that $`a|b`$, $`c|b`$, $`p|(b/a+1)`$, $`p|(c/a+1)`$, $`D=b^2ac>0`$. Then we set $`d_1=b(b+a)`$, $`d_2=b(b+c)`$, $`d=(b+a)(b+c)`$, $`d_0=D`$, $`f=(b+a)(b+c)p/h`$, where $`h`$ is the greatest common divisor of $`b/a+1`$ and $`b/c+1`$, so that we are in the situation considered above. We can rewrite the series (7.6) using the change of variables
$$\tau =\frac{ach^2}{b(b+a)(b+c)}\tau ^{}.$$
Then the above argument implies the modularity of the series
$$q^{\frac{p^2ac(2bs_1s_2as_1^2cs_2^2)}{8D}}\underset{(m,n)^2,mn(p)}{\overset{indef}{}}\zeta _{2p}^{r_1m+r_2n}q^{bmn+a\frac{m^2+mps_1}{2}+c\frac{n^2+nps_2}{2}},$$
where $`\zeta _{2p}`$ is a primitive root of $`1`$ of order $`2p`$ and $`r_1,r_2,s_1,s_2`$ are odd integers. Note that we have
$$\zeta _{2p}^{r_1m+r_2n}=\{\begin{array}{cc}\zeta _p^{\frac{r_1+r_2}{2}m},\text{ if }mn(2p),\hfill & \\ \zeta _p^{\frac{r_1+r_2}{2}m},\text{ if }mn+p(2p).\hfill & \end{array}$$
Thus, the above series is equal to
$$q^{\frac{p^2ac(2bs_1s_2as_1^2cs_2^2)}{8D}}\underset{(m,n)^2,mn(p)}{\overset{indef}{}}(1)^{\frac{nm}{p}}\zeta _p^{\frac{r_1+r_2}{2}m}q^{bmn+a\frac{m^2+mps_1}{2}+c\frac{n^2+nps_2}{2}}.$$
Since $`\frac{r_1+r_2}{2}`$ can be an arbitrary integer we derive that for any residue $`r/p`$ the series
$$q^{\frac{p^2ac(2bs_1s_2as_1^2cs_2^2)}{8D}}\underset{(m,n)^2,mnr(p)}{\overset{indef}{}}(1)^{\frac{nm}{p}}q^{bmn+a\frac{m^2+mps_1}{2}+c\frac{n^2+nps_2}{2}}$$
is modular. ∎
Proof of Theorem 6. Let us denote the series $`F_l`$ considered above by $`F_l(s_1,s_2)`$ to show their dependence on a pair of odd integers $`s_1,s_2`$. Note that using the change of variables from the proof of Theorem 2 we can identify $`f_{s_1,s_2}`$ with $`F_0(s_1,s_2)`$ multiplied by some power of $`q`$. By definition $`F_l(s_1,s_2)=0`$ unless $`l(dd_1)+(dd_2)`$. On the other hand, an easy computation shows that for all $`l_1,l_2`$ one has
$$F_0(s_1+2\frac{h}{p}l_1,s_2+2\frac{h}{p}l_2)=\zeta F_{(dd_1)l_2(dd_2)l_1}(s_1,s_2)$$
where $`\zeta `$ is a root of unity. The assumption in Theorem 6 is precisely the condition (0.3) for the pair $`x_1,x_2`$ defined by (7.5). By Theorem 4 at least one series among $`F_l`$ is non-zero which implies our statement. ∎
The simplest examples of identities that can be derived from the above computations are obtained in the case $`d_0=1`$. Then $`d=d_1+d_21`$ so the conditions $`(dd_1)|d`$, $`(dd_2)|d`$ are satisfied only in the following cases (assuming that $`d_1d_2`$): (i) $`d_1=d_2=2`$; (ii) $`d_1=2`$, $`d_2=3`$; and (iii) $`d_1=3`$, $`d_2=4`$. In case (i) we obtain formula (0.5) taking $`f=1`$. In case (ii) we get identities (0.6), (0.7) and (0.8). More precisely, the case $`f=4`$ leads to (0.6) while in the case $`f=2`$ we get identities (0.7) and (0.8) corresponding to the cases $`s_1s_2(4)`$ and $`s_1s_2(4)`$. In case (iii) above the assumption (0.3) of Theorem 4 is never satisfied, so we do not get any new identity.
One can also consider the degenerate case $`d_0=d_1<d_2=d`$ in the above picture. The coefficients of the Massey products in this case are given by the series (0.9). Application of the above analysis implies that this is a meromorphic Jacobi form of weight $`1`$ as claimed in Theorem 8. On the other hand, we obtain its expression in terms of theta functions. The case $`a=1`$ is well known (see , Section 486, or (5.26), or ). In the case $`a=2`$ we get for any odd $`s`$
(7.7)
$$\underset{n}{}\frac{q^{\frac{n^2+ns}{2}}}{1q^{2n}u}=\phi (q^2)^3\frac{_nq^{2n^2+n(s2)}u^n}{(_n(1)^nq^{n^2n}u^n)(_nq^{2n^2+n(s2)})},$$
where $`\phi (q)=_{n1}(1q^n)`$.
### 7.3. String functions for $`A_1^{(1)}`$
In the paper the authors discovered a relation between the string functions of irreducible highest weight representations of affine Lie algebra $`𝔤`$ of type $`A_1^{(1)}`$ and Hecke’s modular forms for certain indefinite quadratic modules. Namely, for every dominant weight $`\mathrm{\Lambda }`$ and every weight $`\lambda `$ they define the string function $`c_\lambda ^\mathrm{\Lambda }(\tau )`$ which describes a part of the character of the irreducible $`𝔤`$-module with highest weight $`\mathrm{\Lambda }`$ (see ). One of the formulas they get for this function (, bottom of page 258) is
$$\eta (\tau )^3c_\lambda ^\mathrm{\Lambda }(\tau )=\underset{k,n\frac{1}{2},kn(),k|n|\text{or}k>|n|}{}(1)^{2k}\mathrm{sign}(k+\frac{1}{4})q^{(m+2)(k+A)^2m(n+B)^2},$$
where $`m`$ is the level of $`\mathrm{\Lambda }`$ (the value of $`\mathrm{\Lambda }`$ on the central generator), $`A`$ and $`B`$ are rational numbers determined by $`\mathrm{\Lambda }`$ and $`\lambda `$, such that $`2(m+2)A\pm 2mB`$ are odd integers. Taking $`kn`$ and $`k+n`$ as new summation variables we can rewrite this series as follows:
$$q^{\frac{2bs_1s_2s_1^2s_2^2}{8(b^21)}}\underset{k,n}{\overset{indef}{}}(1)^{k+n}q^{(m+1)kn+\frac{k^2+s_1k}{2}+\frac{n^2+s_2n}{2}},$$
where $`s_1=2(m+2)A+2mB`$, $`s_2=2(m+2)A2mB`$. As shown in this series is equal to
$$\mathrm{\Theta }_{^2,Q;(A,B)}^H,$$
where the quadratic form $`Q`$ on $`^2`$ is given by $`Q(x,y)=2(m+2)x^22my^2`$. Below we generalize this observation to other series considered in Theorem 2.
### 7.4. Proof of Theorem 3
As we observed above one can replace the condition $`(m+\frac{1}{2})(n+\frac{1}{2})>0`$ in the definition of the series from Theorem 2 by any condition of the form $`(m+\alpha )(n+\beta )>0`$ where $`\alpha ,\beta `$ (due to the vanishing of the similar sum along the lines parallel to the generators on the cone). So we can write this series as follows:
$$F=\underset{(m,n)pp+𝐜,(m+\frac{1}{4D})(n+\frac{1}{4D})>0}{}\mathrm{sign}(m+\frac{1}{4D})(1)^{m_0+n_0}q^{\frac{am^2+2bmn+cn^2}{2}},$$
where $`𝐜=(r+\frac{pc(bs_2as_1)}{2D},r+\frac{pa(bs_1cs_2)}{2D})`$, $`(m_0,n_0)^2`$ is defined by $`(m,n)=p(m_0,n_0)+𝐜`$. Let us take $`m^{}=m/p`$ and $`n^{}=(n+\frac{b}{c}m)/p`$ as the new summation variables. Note that if $`(m,n)=p(m_0,n_0)+𝐜`$ then
$$(m^{},n^{})=(m_0,n_0+\frac{b}{c}m_0)+(\frac{r}{p}+\frac{acs}{2D},r\frac{b/c+1}{p}+\frac{s_2}{2}),$$
where $`s=s_1+\frac{b}{a}s_2`$. Thus, $`(m^{},n^{})`$ runs through the coset $`^2+(\frac{r}{p}+\frac{acs}{2D},\frac{1}{2})`$ intersected with the cone $`(m^{}+\epsilon )(n^{}\frac{b}{c}m^{}+\epsilon )>0,`$ where $`\epsilon >0`$ is sufficiently small. Therefore, we can write
$$\pm F=\underset{(m,n)^2+(\frac{r}{p}+\frac{acs}{2D},\frac{1}{2}),(m+\epsilon )(n\frac{b}{c}m+\epsilon )>0}{}\mathrm{sign}(m+\epsilon )(1)^{n_0+(\frac{b}{c}+1)m_0}q^{p^2(cn^2\frac{D}{c}m^2)},$$
where $`(m,n)=(m_0,n_0)+(\frac{r}{p}+\frac{acs}{2D},\frac{1}{2})`$ (the sign in front of $`F`$ depends on the parity of $`r(b/c+1)/p+(s_21)/2`$). Now the lattice $`\mathrm{\Lambda }`$ defined in the formulation of the theorem comes into play. Namely, splitting the above sum in two pieces according to the parity of $`n_0+(\frac{b}{c}+1)m_0`$ we obtain
$`\pm F={\displaystyle \underset{(m,n)(\mathrm{\Lambda }+(\frac{r}{p}+\frac{acs}{2D},\frac{1}{2}))C}{}}\mathrm{sign}(m+\epsilon )q^{\frac{p^2}{2}Q(m,n)}`$
$`{\displaystyle \underset{(m,n)(\mathrm{\Lambda }+(\frac{r}{p}+\frac{acs}{2D},\frac{1}{2}))C}{}}\mathrm{sign}(m+\epsilon )q^{\frac{p^2}{2}Q(m,n)},`$
where $`C`$ is the cone $`(m+\epsilon )(n\frac{b}{c}m+ϵ)>0`$. It is convenient to identify the lattice $`\mathrm{\Lambda }`$ with a $``$-submodule in the field $`K=(\sqrt{D})`$ by associating to $`(m,n)\mathrm{\Lambda }`$ the element $`cn+m\sqrt{D}`$. Then we have
$$\frac{1}{2}Q(m,n)=\frac{1}{c}\mathrm{Nm}(cn+m\sqrt{D}).$$
For two non-zero elements $`k_1,k_2K`$ let us denote by $`k_1,k_2=_{>0}k_1+_{>0}k_2`$, $`[k_1,k_2]=_0k_1+_0k_2`$, $`k_1,k_2]=_0k_1+_{>0}k_2`$. The intersection of a $`\mathrm{\Lambda }`$-coset with the cone $`C`$ is equal to its intersection with the set $`[1,b+\sqrt{D}]1,b\sqrt{D}`$. Making the change of variables $`(m,n)(m,n)`$ in the second sum of the above expression for $`\pm F`$ we can write:
$`\pm F={\displaystyle \underset{(m,n)(\mathrm{\Lambda }+𝐜)([1,b+\sqrt{D}]1,b\sqrt{D})}{}}\mathrm{sign}(m+\epsilon )q^{\frac{p^2}{2}Q(m,n)}`$
$`{\displaystyle \underset{(m,n)(\mathrm{\Lambda }+𝐜)([1,b+\sqrt{D}]1,b\sqrt{D})}{}}\mathrm{sign}(m+\epsilon )q^{\frac{p^2}{2}Q(m,n)},`$
where $`𝐜=(\frac{r}{p}+\frac{acs}{2D},\frac{1}{2})`$. Since $`\mathrm{\Lambda }+𝐜`$ doesn’t contain zero, the obtained series is equal to
$$\underset{(m,n)(\mathrm{\Lambda }+𝐜)(b\sqrt{D},b+\sqrt{D}]b\sqrt{D},b+\sqrt{D}])}{}\mathrm{sign}(cn+m\sqrt{D})q^{\frac{p^2}{2}Q(m,n)}.$$
Let us consider the totally positive element $`ϵ=\frac{b+\sqrt{D}}{b\sqrt{D}}K`$. Since $`\mathrm{Nm}(ϵ)=1`$ the multiplication by $`ϵ`$ preserves the quadratic form $`Q`$ on $`\mathrm{\Lambda }`$. The direct computation shows that the multiplication by $`ϵ`$ preserves also $`\mathrm{\Lambda }+𝐜`$, so we have
$$\pm F=\underset{(m,n)(\mathrm{\Lambda }+𝐜)(Q>0)/G_ϵ}{}\mathrm{sign}(cn+m\sqrt{D})q^{\frac{p^2}{2}Q(m,n)},$$
where $`G_ϵK^{}`$ is the subgroup generated by $`ϵ`$. Let $`G`$ be the subgroup of the group $`\mathrm{ker}(\mathrm{Nm}:K^{}^{})`$ consisting of elements $`k`$ such that $`k`$ is totally positive and $`k(+𝐜)=+𝐜`$. Then $`G_ϵ`$ is a subgroup of finite index in $`G`$, so we have
$$\pm F=|G/G_ϵ|\underset{(m,n)(\mathrm{\Lambda }+𝐜)(Q>0)/G}{}\mathrm{sign}(cn+m\sqrt{D})q^{\frac{p^2}{2}Q(m,n)}=|G/G_ϵ|\mathrm{\Theta }_{\mathrm{\Lambda },Q;𝐜}^H(p^2\tau ).$$
### 7.5. Determinantal Jacobi forms
We are going to compute explicitly the map $`MP_{s_1,s_2}^{det}`$ (see section 6.3) in the setup of section 7.1. Note that we can trivialize the $`1`$-dimensional vector space $`M`$ using the canonical bases in the relevant vector spaces. First, let us compute the map
$$e_{s_1,s_2}:^{d_01}\mathrm{Hom}^1(L_1,L_2)^{}\mathrm{Hom}^1(L_1,L_2).$$
For every subset $`S/d`$ of cardinality $`dd_0+1`$ let us denote by $`e_{12}^S`$ the element in $`^{d_01}\mathrm{Hom}^1(L_1,L_2)^{}`$ which is induced by the projection
$$\mathrm{Hom}^1(L_1,L_2)^{d_01}:\underset{k}{}y_ke_{12}(0,k)(y_k,kS).$$
Then
$$e_{s_1,s_2}(e_{12}^S)=\underset{kS}{}c_{k,S}e_{12}(0,k)$$
where $`(c_{k,S},kS)`$ is the sequence of $`(dd_0)\times (dd_0)`$-minors (with signs) of the $`(dd_0)\times (dd_0+1)`$-matrix $`R`$ obtained by putting together the $`(dd_1)\times (dd_0+1)`$-matrix $`R_1=(A_{ik};i/(dd_1),kS)`$ and the $`(dd_2)\times (dd_0+1)`$-matrix $`R_2=(B_{jk};j/(dd_2),kS)`$, where $`A_{ik}`$ is zero unless $`ki(d,dd_1)`$ in which case
$$A_{ik}=\theta _{I_1,\frac{k}{d}+n_1(k,i)}(p_1x_1,p_1\tau ),$$
where $`n_1(k,i)/I_1`$ is characterized by the congruence $`dn_1(k,i)ki(dd_1)`$; similarly $`B_{jk}`$ is zero unless $`kj(d,dd_2)`$ in which case
$$B_{jk}=\theta _{I_2,\frac{k}{d}+n_2(k,j)}(p_2x_2,p_2\tau ),$$
where $`n_2(k,j)/I_2`$ is characterized by the congruence $`dn_2(k,j)jk(dd_2)`$.
Now we have
$$F_{l,S}=e_{03}^{}(0,l),MP_{s_1,s_2}^{det}(e_{12}^S)=\underset{kS}{}C_{l,k}c_{k,S},$$
where $`C_{l,k}`$ are defined by (7.3). In other words, $`F_{l,S}`$ is equal to the determinant of the $`(dd_0+1)\times (dd_0+1)`$-matrix obtained by putting together $`R_1`$, $`R_2`$ and the row $`(C_{l,k},kS)`$ of length $`dd_0+1`$.
Let us recall the definition of Jacobi forms from . Let $`\mathrm{\Lambda }`$ be a lattice, $`𝐧𝐧^{}`$ be a symmetric bilinear form on $`\mathrm{\Lambda }`$ with values in $``$, $`\mathrm{\Gamma }\mathrm{SL}(2,)`$ be a congruenz-subgroup. Let us denote $`Q(𝐧)=𝐧𝐧`$ (this corresponds to $`2Q`$ in the notation of ). Then a meromorphic function $`f(𝐳,\tau )`$ on $`\mathrm{\Lambda }_{}\times `$ is called a (meromorphic) Jacobi form of weight $`k`$ with respect to $`(\mathrm{\Lambda },Q,\mathrm{\Gamma })`$ if the following equations hold:
$$f(𝐳+𝐯\tau +𝐰,\tau )=(1)^{𝐯𝐰}\mathrm{exp}(\pi i\tau Q(𝐯)2\pi i𝐯𝐳)f(𝐳,\tau ),$$
$$f(\frac{z}{c\tau +d},\frac{a\tau +b}{c\tau +d})=(c\tau +d)^k\mathrm{exp}(\pi i\frac{cQ(𝐳)}{c\tau +d})f(z,\tau ),$$
for every $`𝐯,𝐰\mathrm{\Lambda }`$, $`\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)\mathrm{\Gamma }`$. Similar to the case of modular forms one can extend the above definition to half-integer weights $`k`$.
###### Theorem 9.
The function $`F_{l,S}`$ is a Jacobi form of weight $`(dd_0)/2+1`$ with respect to some sublattice $`\mathrm{\Lambda }^{}\mathrm{\Lambda }^+`$, some congruenz-subgroup $`\mathrm{\Gamma }\mathrm{SL}(2,)`$, and the quadratic form $`Q+Q_0`$, where
$$\mathrm{\Lambda }^+=\{𝐧=(n_0,n_1,n_2,n_3)^4:\underset{i}{}n_i=0,d_1n_1+(d_0d_2)n_2+d_0n_3=0\},$$
$$Q(𝐧)=d_1n_0n_1d_2n_2n_3,$$
$$Q_0(𝐧)=d(d_1n_1^2+d_2n_2^2).$$
Proof. We have the system of linear equations
$$\underset{l}{}D_{k,l}F_{l,S}=c_{k,S}$$
determining $`F_{l,S}`$. Using the functional equation for theta function we derive that $`D_{k,l}`$ are Jacobi forms of weight $`1`$ with respect to some congruenz-subgroup of $`\mathrm{SL}(2,)`$, some sublattice of $`\mathrm{\Lambda }^+`$ and the quadratic form $`Q`$. On the other hand, $`c_{k,S}`$ are Jacobi forms of weight $`(dd_0)/2`$ with respect to the quadratic form $`Q_0`$. Hence, $`F_{l,S}`$ are Jacobi forms of weight $`(dd_0)/2+1`$ with respect to $`Q+Q_0`$. ∎
Remark. It is easy to see that the form $`Q+Q_0`$ can be written as follows:
$$(Q+Q_0)(𝐧)=\frac{d_1d_2}{d_0}(x_1+x_2)^2+\frac{d(d_01)}{d_0}(d_1x_1^2+d_2x_2^2).$$
In particular, it is always positive-definite for $`d_0>1`$. In the case $`d_0=1`$ this form is degenerate which means that the corresponding functions $`F_{l,S}`$ have form $`f(x_1+x_2,\tau )`$, where $`f(z,\tau )`$ is a Jacobi form.
For example, for $`d_1=d_2=2`$, $`d_0=1`$ we have the unique choice of $`l`$ and $`S`$. Let us take $`z_1=x_0,z_2=x_3`$ as coordinates in $`\mathrm{\Lambda }_{}^+`$. Then the corresponding function is equal (up to a sign) to the following deteminant:
$$F(x_1,x_2;\tau )=det\left(\begin{array}{ccc}\mathrm{\Theta }_0(x_1,x_2;\tau )& \mathrm{\Theta }_1(x_1,x_2;\tau )& \mathrm{\Theta }_1(x_1,x_2;\tau )\\ \theta _0(6x_1,6\tau )& \theta _{\frac{1}{3}}(6x_1,6\tau )& \theta _{\frac{1}{3}}(6x_1,6\tau )\\ \theta _0(6x_2,6\tau )& \theta _{\frac{1}{3}}(6x_2,6\tau )& \theta _{\frac{1}{3}}(6x_2,6\tau )\end{array}\right)$$
where we denoted $`\theta _r=\theta _{,r}`$,
$`\mathrm{\Theta }_c={\displaystyle \underset{mnc(3),(m+\alpha (2x_2x_1))(n+\alpha (2x_1x_2))>0}{}}\mathrm{sign}(m+\alpha (2x_2x_1))\times `$
$`\mathrm{exp}(\pi i\tau {\displaystyle \frac{2}{3}}(m^2+4mn+n^2)+4\pi i(mx_1+nx_2)).`$
The function $`F`$ is a Jacobi form in $`(x_1+x_2,\tau )`$ of weight $`2`$ and index $`2`$. Using the equation $`D_{0,0}F=c_0`$ and the addition formulas for theta functions we derive the following identity:
$$F(x_1,x_2;\tau )=\theta _{\frac{1}{2}}(2(x_1+x_2)+\frac{1}{2},2\tau )^2\frac{\eta ^3(2\tau )_n\chi _3(n)q^{\frac{(4n+3)^2}{24}}}{\theta (\frac{1}{2},2\tau )\theta _{\frac{1}{4}}(0,4\tau )}$$
where $`\chi _3(n)`$ is the non-trivial Dirichlet character modulo $`3`$.
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# I INTRODUCTION
## I INTRODUCTION
Numerous experiments using unstable nuclear beams have succeeded to reveal exotic properties of $`\beta `$-unstable nuclei, including a neutron halo structure. Especially, the shift of the closed-shell structure is a characteristic behavior of systems with weakly bound neutrons. Recently, the contributions of the $`sd`$ shell have been analyzed in $`N=8`$ nuclei based on the shell model. A calculation has shown that the slow $`\beta `$-decay of <sup>12</sup>Be to <sup>12</sup>B can be explained by an admixture of the $`sd`$ shell in <sup>12</sup>Be ($`N=8`$) in which the closed $`p`$-shell component must be less than 30$`\%`$. This shows that the concept of magic numbers is vague in <sup>12</sup>Be.
On the other hand, an $`\alpha `$-$`\alpha `$ structure is well established in the Be and B region. Typically in <sup>9</sup>Be, a microscopic $`\alpha `$+$`\alpha `$+$`n`$ model has reproduced the properties of low-lying states including the level inversion of $`1/2^{}`$ in the $`p`$-shell and $`1/2^+`$ in the $`sd`$ shell. Also in <sup>10</sup>Be, the microscopic $`\alpha `$-cluster models have been applied , and we have quantitively shown the mechanism for the lowering of the $`sd`$ orbits related with the clustering of the core.
According to our previous results for <sup>10</sup>Be, the main properties of this nucleus have been well described by the $`\alpha `$+$`\alpha `$+n+n model, where the orbits for the two valence neutrons are classified based on the molecular-orbit (MO) model. The dominant configuration of the two valence neutrons for the second $`0^+`$ state is $`(1/2^+)^2`$, and the level inversion between $`1/2^{}`$ and $`1/2^+`$ in <sup>9</sup>Be also holds in <sup>10</sup>Be. Here, the two valence neutrons stay along the $`\alpha `$-$`\alpha `$ axis, and reduce the kinetic energy by enhancing the $`\alpha `$-$`\alpha `$ distance (up to around 4 fm). Therefore, the $`1/2^+`$ orbit in this case is not a spherical $`s`$-wave, but a polarized orbit with the $`d`$-wave component. This feature of $`sd`$ mixing in the $`1/2^+`$ orbit can be qualitatively interpreted in terms of deformed models ( expression in the Nilsson diagram).
The main purpose of this paper is to show the mechanism for the appearance of this $`(1/2^+)^2`$ configuration in low-lying energy in detail. In our previous work, the calculated excitation energy of the second $`0^+`$ state has been higher by $``$5 MeV without the spin-orbit interaction than the experimental one. The spin-orbit interaction has been found to decrease the excitation energy by 3 MeV by adding a single optimal wave function with $`S=1`$. Hence, it is necessary to analyze the contribution of the spin-orbit interaction more precisely to clarify the lowering mechanism of the second $`0^+`$ state. The other reduction mechanism of the kinetic energy by more accurate description of the neutron-tail should also be taken into account.
Firstly, the spin-orbit interaction is able to contribute to the second $`0^+`$ state, when spin-triplet states for the valence neutrons are included among the basis states. If the $`1/2^+`$ state consists of the pure s-orbit, naturally there is no contribution of the spin-orbit interaction. In Be isotopes, the state contains the $`d`$-orbit component. However, the spin-orbit interaction again vanishes, when the two valence neutrons occupy along the $`\alpha `$-$`\alpha `$ axis, since two neutrons with the same spatial distribution construct only the spin-singlet state. This has been the situation in traditional MO models. It will be shown that when one of the valence neutrons deviates from the $`\alpha `$-$`\alpha `$ axis, the spin-triplet state can be constructed, and the spin-orbit interaction strongly acts between this state and the original $`(1/2^+)^2`$ configuration with the spin-singlet. The extension of the model space also enables us to improve the description of the neutron-tail. We estimate the reduction of the kinetic energy due to the improvement of the neutron-tail. These improvements will be shown to largely decrease the excitation energy of the second $`0^+`$ state, and that the calculated energy just corresponds to the experimental value.
Using the present extended MO model, negative-parity states in <sup>10</sup>Be are also investigated. The mixing of different spin structure $`(S=0,S=1)`$ is found to also be important in the negative-parity states. The $`K`$-mixing effect strongly acts on the $`2^{}`$ state, and the calculated small-level spacing between $`1^{}`$ and $`2^{}`$ ($``$300 keV) agrees with the experimental value.
Also, in <sup>12</sup>Be these two important mechanism for lowering the $`(1/2^+)^2`$ configuration are examined. In <sup>12</sup>Be, the configuration for four valence neutrons around the $`\alpha `$-$`\alpha `$ core is considered to correspond to the two pair configurations, $`(3/2^{})^2`$ and $`(1/2^+)^2`$. These two pair configurations appear in the ground $`0^+`$ state and the second $`0^+`$ state of <sup>10</sup>Be, respectively. The energy surface for the ground state of <sup>12</sup>Be will be shown, in which the spin-orbit interaction is seen to make the binding for the $`(3/2^{})^2(1/2^+)^2`$ configuration properly stronger in comparison with the $`(3/2^{})^2(1/2^{})^2`$ configuration. This is related to the breaking of the $`N=8`$ (closed $`p`$-shell) neutron-magic number.
This paper is organized as follows. In Sec. II, we summarize a description of the single-particle orbits around the two $`\alpha `$ clusters based on the MO model. In Sec. III, the contribution of the spin-orbit interaction to the $`(1/2^+)^2`$ configuration and the reduction of the kinetic energy due to an improvement of the neutron-tail are considered. These effects are discussed in detail for the second $`0^+`$ state of <sup>10</sup>Be (A), for the negative parity states of <sup>10</sup>Be (B), and for <sup>12</sup>Be (C). The conclusion is given in Sec. IV.
## II extended Molecular Orbital model
We introduce a microscopic $`\alpha `$+$`\alpha `$+$`2n`$ model for <sup>10</sup>Be and $`\alpha `$+$`\alpha `$+$`4n`$ model for <sup>12</sup>Be. The neutron configurations are classified according to the molecular-orbit (MO) picture. The details of the framework are given in Ref.; here, the essential part is presented.
The Hamiltonian is the same as in Ref., and consists of a kinetic-energy term, a central two-body interaction term, a two-body spin-orbit interaction term, and a Coulomb interaction term:
$$=\underset{i}{}T_iT_{cm}+\underset{i<j}{}V_{ij}+\underset{i<j}{}V_{ij}^{ls}+\underset{i<j}{}\frac{e^2}{4r_{ij}}(1\tau _z^i)(1\tau _z^j).$$
(1)
The effective nucleon-nucleon interaction is Volkov No.2 for the central part and the G3RS spin-orbit term for the spin-orbit part, as follows:
$$V_{ij}=\{V_1e^{a_1r_{ij}^2}V_2e^{a_2r_{ij}^2}\}\{WMP^\sigma P^\tau +BP^\sigma HP^\tau \},$$
(2)
$$V_{ij}^{ls}=V_0^{ls}\{e^{a_1r_{ij}^2}e^{a_2r_{ij}^2}\}\stackrel{}{L}\stackrel{}{S}P_{31},$$
(3)
where $`P_{31}`$ is a projection operator onto the triplet odd state. The parameters are $`V_1=60.650`$ MeV, $`V_2=61.140`$ MeV, $`a_1=0.980`$ fm<sup>-2</sup> and $`a_2=0.309`$ fm<sup>-2</sup> for the central interaction, and $`V_0^{ls}`$ = 2000 MeV, $`a_1=5.00`$ fm<sup>-2</sup>, and $`a_2=2.778`$ fm<sup>-2</sup> for the spin-orbit interaction. We employ the Majorana exchange parameter, $`M=0.6`$ ($`W=0.4`$), the Bartlett exchange parameter, $`B=0.125`$, and the Heisenberg exchange parameter, $`H=0.125`$, for the Volkov interaction (using $`B`$ and $`H`$, there is no neutron-neutron bound state). All of these parameters are determined from the $`\alpha +n`$ and $`\alpha +\alpha `$ scattering phase shifts and the binding energy of the deuteron.
The total wave function is fully antisymmetrized and expressed by a superposition of terms centered to different relative distances between the two $`\alpha `$ clusters ($`d`$) with various configurations of the valence neutrons. The projection to the eigen-states of angular momentum $`(J)`$ is performed numerically. All nucleons are described by Gaussians with the oscillator parameter $`(s)`$ equal to 1.46 fm. The $`\alpha `$ clusters located at $`d/2`$ and $`d/2`$ on the $`z`$-axis consist of four nucleons:
$$\varphi ^{(\alpha )}=G_{R_\alpha }^pG_{R_\alpha }^pG_{R_\alpha }^nG_{R_\alpha }^n\chi _p\chi _p\chi _n\chi _n.$$
(4)
$`G`$ represents Gaussians:
$$G_{R_\alpha }=\left(\frac{2\nu }{\pi }\right)^{\frac{3}{4}}\mathrm{exp}[\nu (\stackrel{}{r}\stackrel{}{R}_\alpha )^2],\nu =1/2s^2,\stackrel{}{R}_\alpha =\{d\stackrel{}{e}_z/2,d\stackrel{}{e}_z/2\}.$$
(5)
Each valence neutron ($`\varphi _{ci}\chi _{ci}`$) around the $`\alpha `$-$`\alpha `$ core is expressed by a linear combination of local Gaussians:
$$\varphi _{ci}\chi _{ci}=\underset{j}{}g_jG_{R_n^j}\chi _{ci},$$
(6)
$$G_{R_n^j}=\left(\frac{2\nu }{\pi }\right)^{\frac{3}{4}}\mathrm{exp}[\nu (\stackrel{}{r}\stackrel{}{R}_n^j)^2],\nu =1/2s^2.$$
(7)
The level structure is solved by superposing basis states with different $`\alpha `$-$`\alpha `$ distances and configurations after the angular momentum projection.
In the MO model, valence neutrons are expressed by a linear combination of orbits around two $`\alpha `$ clusters. The orbit around each $`\alpha `$ cluster is called the atomic orbit (AO). Here, we note that the valence neutrons and neutrons in the $`\alpha `$ clusters are identical particles, and antisymmetrization imposes the AO forbidden space. The lowest AO has one node and parity minus, which is the $`p`$-orbit. When a linear combination of two AOs ($`p`$-orbit) around the left-$`\alpha `$ cluster and the right-$`\alpha `$ cluster is summed up by the same sign, the resultant MO also has negative parity and one node. In this case, the MO is restricted to spread along an axis perpendicular to the $`\alpha `$-$`\alpha `$ $`(z)`$ axis, which is the so-called $`\pi `$-orbit. It cannot spread along the $`z`$-axis where the two $`\alpha `$ clusters are already located. If two $`p`$-orbits are summed up by different signs, the resultant MO has two nodes and positive parity. This MO can spread to all directions, and the optimal direction becomes the $`z`$-direction. This is the so-called $`\sigma `$-orbit, and the energy becomes lower as the distance between two $`\alpha `$ clusters increases.
In Be isotopes, the low-lying orbits with $`K^\pi =3/2^{}`$, $`1/2^+`$, and $`1/2^{}`$ are important; these are classified based on the MO mode1. The $`1/2^+`$ state is the $`\sigma `$-orbit, and the spin-orbit interaction splits the $`\pi `$-orbit to $`3/2^{}`$ and $`1/2^{}`$. In the present framework, each valence neutron is introduced to have a definite $`K^\pi `$ at the zero limit of centers of local Gaussians $`(\{R_n^j\})`$ describing the spatial distribution of the orbit. The precise positions of $`\{R_n^j\}`$ are determined variationally for each $`\alpha `$-$`\alpha `$ distance before the angular-momentum projection. Since the values of $`\{R_n^j\}`$ are optimized to be finite, the orbits are not exactly the eigen-state of $`K^\pi `$, and are labeled as $`\overline{K}^\pi `$.
In <sup>10</sup>Be, when two valence neutrons occupy these single particle orbits, three configurations with the total $`\overline{K}=\overline{K}_1+\overline{K}_2=0`$ are generated. For the ground state, valence neutrons occupy orbits with $`\overline{K}^\pi =3/2^{}`$ and $`\overline{K}^\pi =3/2^{}`$; these $`\pi `$-orbits are expressed as linear combinations of $`p_x`$ and $`p_y`$:
$$\mathrm{\Phi }(3/2^{},3/2^{})=𝒜[\varphi _1^{(\alpha )}\varphi _2^{(\alpha )}(\varphi _{c1}\chi _{c1})(\varphi _{c2}\chi _{c2})].$$
(8)
$$\varphi _{c1}\chi _{c1}=\{(p_x+ip_y)_{+a}+(p_x+ip_y)_a\}|n,$$
(9)
$$\varphi _{c2}\chi _{c2}=\{(p_xip_y)_{+a}+(p_xip_y)_a\}|n.$$
(10)
Each MO is constructed from two AOs around a left and right $`\alpha `$ cluster. We introduce a variational parameter $`(a)`$ on the $`\alpha `$-$`\alpha `$ axis ($`z`$-axis); and $`(p_x\pm ip_y)_{+a}`$ is $`p_x\pm ip_y`$, whose center is $`+a`$ on the $`z`$-axis; and $`(p_x\pm ip_y)_a`$ is $`p_x\pm ip_y`$, whose center is $`a`$ on the $`z`$-axis. Here, the spatial part and the spin part of $`\overline{K}`$ in Eqs. (9) and (10) are introduced to be parallel, for which the spin-orbit interaction acts attractively. The spin-up valence neutron has $`\overline{K}_1^\pi =3/2^{}`$ ($`rY_{11}|n`$), and the spin-down valence neutron has $`\overline{K}_1^\pi =3/2^{}`$ ($`rY_{11}|n`$). There $`(p_x)_{+a}`$, $`(p_x)_a`$, $`(p_y)_{+a}`$, and $`(p_y)_a`$ are expressed as the combination of two Gaussians, whose centers are shifted by a variational parameter $`(b)`$ perpendicular to the $`z`$-axis. For example,
$$(p_x)_{+a}=G_{a\stackrel{}{e}_z+b\stackrel{}{e}_x}G_{a\stackrel{}{e}_zb\stackrel{}{e}_x},(p_y)_{+a}=G_{a\stackrel{}{e}_z+b\stackrel{}{e}_y}G_{a\stackrel{}{e}_zb\stackrel{}{e}_y}.$$
(11)
These parameters $`a`$ and $`b`$, are optimized by using the Cooling Method in antisymmetrized molecular dynamics (AMD) for each $`\alpha `$-$`\alpha `$ distance.
The $`\sigma `$-orbit pair, $`(1/2^+)^2`$, is prepared, where each $`\sigma `$-orbit is expressed by subtracting two AOs at $`+a=d/2`$ and $`a=d/2`$:
$$\mathrm{\Phi }(1/2^+,1/2^+)=𝒜[\varphi _1^{(\alpha )}\varphi _2^{(\alpha )}(\varphi _{c1}\chi _{c1})(\varphi _{c2}\chi _{c2})],$$
(12)
$$\varphi _{c1}\chi _{c1}=\{(\stackrel{}{p})_{+a}(\stackrel{}{p})_a\}|na=d/2,$$
(13)
$$\varphi _{c2}\chi _{c2}=\{(\stackrel{}{p})_{+a}(\stackrel{}{p})_a\}|na=d/2,$$
(14)
where,
$$(\stackrel{}{p})_{\pm a}=(G_{+\stackrel{}{b}}G_\stackrel{}{b})_{\pm a}.$$
(15)
Since the $`\sigma `$-orbit has two nodes, the direction of the orbit is not limited due to the Pauli principle. It can thus take all directions. Therefore, the direction and value of parameter $`\stackrel{}{b}`$ for the center of the Gaussians are optimized. As a result, it is optimized to the direction of the $`\alpha `$-$`\alpha `$ axis. This basis state shares the dominant component of the second $`0^+`$ state.
## III RESULTS obtained using the extended MO model space
### A $`J^\pi =0^+`$ states in <sup>10</sup>Be
Before giving the results, we summarize the calculated energies and the optimal $`\alpha `$-$`\alpha `$ distance for three configurations ($`(3/2^{})^2`$, $`(1/2^{})^2`$, and $`(1/2^+)^2`$) discussed in Ref.. In Fig. 1, the $`0^+`$ energy curves of $`\mathrm{\Phi }(3/2^{},3/2^{})`$, $`\mathrm{\Phi }(1/2^+,1/2^+)`$, and $`\mathrm{\Phi }(1/2^{},1/2^{})`$ in <sup>10</sup>Be are presented as a function of the $`\alpha `$-$`\alpha `$ distance $`(d)`$. These configurations are the dominant components of the first, second, and third $`0^+`$ states in <sup>10</sup>Be, respectively.
————–
Fig. 1
————–
The energy of $`\mathrm{\Phi }(1/2^+,1/2^+`$) (dominant component of $`0_2^+`$) becomes lower as $`d`$ increases, and the energy becomes minimum at an $`\alpha `$-$`\alpha `$ distance of $`45`$ fm. However, it should be noted that the energy difference between $`\mathrm{\Phi }(1/2^+,1/2^+)`$ and $`\mathrm{\Phi }(3/2^{},3/2^{})`$ in Fig. 1 is about 11 MeV. This is much larger than the experimental excitation energy of the second $`0^+`$ state (6.26 MeV) by 5 MeV.
There are two reasons for this lack of binding energy by 5 MeV in the second $`0^+`$ state. One is that the contribution of the spin-orbit interaction is not taken into account; the other is that the tail of the valence neutron is not correctly described. As for the spin-orbit interaction, in the traditional MO model, the spin-orbit interaction has not contributed to the $`0^+`$ state with this configuration. This is because, the two valence neutrons occupy the same spatial configurations with opposite spin direction, and construct a spin-singlet state. Since it is impossible for two valence neutrons with the same spin direction to occupy the same $`\sigma `$-orbit along the $`\alpha `$-$`\alpha `$ axis, the spin-triplet state is prepared by making one of the valence neutrons deviate from the $`\alpha `$-$`\alpha `$ axis. As for the tail effect of the valence neutrons, in this analysis, the tail behavior is expressed by the superposition of local Gaussians.
These two improvements are shown for the $`\mathrm{\Phi }(1/2^+,1/2^+)`$ configuration with the optimal $`\alpha `$-$`\alpha `$ distance of 4 fm. The optimized parameter $`\stackrel{}{b}`$ is $`2.02\stackrel{}{e}_z`$ (fm). As listed in Table I, the spin-singlet $`\mathrm{\Phi }(1/2^+,1/2^+)`$ configuration has a $`0^+`$ energy of $`46.3`$ MeV. As shown in Eqs. (12)-(15), each valence neutron is described as a linear combination of four local Gaussians. When we optimize their linear combination, the energy slightly decreases to $`47.7`$ MeV (noted as $`S=0`$ $`(\sigma )^2`$ in Table I).
————–
Table I
————–
Next, the contribution of the spin-orbit interaction is taken into account to this spin-singlet state by preparing the spin-triplet state ($`S_z=1`$). One of the valence neutrons ($`\varphi _{c1}\chi _{c1}`$) occupies the $`\sigma `$-orbit, and another one ($`\varphi _{c2}\chi _{c2}`$) with the same spin-direction deviates from the $`\alpha `$-$`\alpha `$ axis. The later one is described by a local Gaussian centered at $`\stackrel{}{R}`$; $`\varphi _{c2}\chi _{c2}=G_\stackrel{}{R}|n`$. The calculated energy surface for the second $`0^+`$ state as a function of the parameter $`\stackrel{}{R}`$ is given in Fig. 2. The orthogonality of the ground $`0^+`$ state to the ground $`0^+`$ state is ensured in this calculation. The point on the $`x`$-$`z`$ plain in Fig. 2 shows the position of the parameter $`\stackrel{}{R}`$, and the contour map shows the total energy obtained by taking into account this coupling with the spin-triplet states.
————-
Fig. 2
————-
The lowest energy of $`50.9`$ MeV is seen at $`\stackrel{}{R}=(x,z)=(1.0,2.0)`$ (fm) on the energy surface. The contribution of the spin-orbit interaction decreases the energy by more than 3 MeV in comparison with the case of the spin-singlet state. Although the energy of such a spin-triplet state, itself, is much higher than that of $`\mathrm{\Phi }(1/2^+,1/2^+)`$ by about 20 MeV, the large coupling between the spin-singlet and spin-triplet states enables the spin-orbit interaction to attractively act. As a result, more than half of the under-binding by 5 MeV for the second $`0^+`$ state is explained by this spin-orbit contribution.
To see this effect more correctly, a lot of spin-triplet states with different $`\stackrel{}{R}`$ values are adopted. The employed $`\stackrel{}{R}`$ values are 0.0, 1.0, 2.0 (fm) for the $`x`$-direction (perpendicular to the $`\alpha `$-$`\alpha `$ direction) and 0.0, 1.0, 2.0, 3.0, 4.0 (fm) for the $`z`$-direction (the $`\alpha `$-$`\alpha `$ direction). In Table I, the column noted as spin-orbit gives the energy of the second $`0^+`$ state when these 15 spin-triplet states with different $`\stackrel{}{R}`$ values are superposed. The calculated energy is $`51.2`$ MeV; this is lower than the energy minimum point in Fig. 2 by only 300 keV. This is because, the spin-triplet basis states around the core have a large overlap with each other, and those apart from the core do not contribute to this coupling. This result shows that the contribution of the spin-orbit interaction for the $`(1/2^+)^2`$ state is approximately taken into account when we employ only one Slater determinant on the energy surface.
In this calculation, many spin-triplet states are employed; however, it should be noted that the spin-singlet state is only represented by the $`\mathrm{\Phi }(1/2^+,1/2^+)`$ configuration. Thus, to improve the tail of the valence neutron in the second $`0^+`$ state of <sup>10</sup>Be, it is necessary to superpose further the $`S_z=0`$ basis states with many different $`\stackrel{}{R}`$ values, which contain the spin-singlet component. The distance of $`\alpha `$-$`\alpha `$ is 4 fm, where one of the valence neutrons ($`\varphi _{c1}\chi _{c1}`$) occupies the $`\sigma `$-orbit, and the other valence neutrons ($`\varphi _{c2}\chi _{c2}`$) with the opposite spin-direction is described by local Gaussian centered at $`\stackrel{}{R}`$. The later one has the form of $`\varphi _{c2}\chi _{c2}=G_\stackrel{}{R}|n`$. We superpose 15 basis states with $`S_z=0`$, and the employed $`\stackrel{}{R}`$ values are 0.0, 1.0, 2.0 (fm) for the $`x`$-direction (perpendicular to the $`\alpha `$-$`\alpha `$ direction) and 0.0, 1.0, 2.0, 3.0, 4.0 (fm) for the $`z`$-direction (the $`\alpha `$-$`\alpha `$ direction). As listed in Table I, the energy becomes $`52.7`$ MeV. These basis states increase the binding of the second $`0^+`$ state by 1.5 MeV compared with $`51.2`$ MeV for the third row. Here, it is noted that these $`S_z=0`$ basis states contain both the spin-single and spin-triplet components. However, the increase of the binding energy by adding these $`S_z=0`$ basis states is due to the spin-singlet basis states, since it is considered that the functional space of the spin-triplet states are already exhausted by the $`S_z=1`$ states. These results show that by taking into account both the spin-orbit interaction (spin-triplet basis states) and the tail effect (spin-singlet basis states), the under-binding for the second $`0^+`$ state by 5 MeV is fully explained. The calculated excitation energy of the second $`0^+`$ state agrees with the experimental one.
To confirm the component of the spin-triplet configurations introduced, the expectation value of the spin-square $`(\widehat{S}^2`$) in the second $`0^+`$ state of <sup>10</sup>Be was calculated. The value is 0.26. This means that the spin-triplet configuration mixes in this state by about 13$`\%`$.
In this analysis, the $`\alpha `$-$`\alpha `$ distance is restricted to 4 fm, but we now examine the contribution of the spin-orbit interaction with respect to the $`\alpha `$-$`\alpha `$ distance.
————–
Table II
————–
Table II gives the contribution of spin-orbit interaction for the second $`0^+`$ state with the $`\alpha `$-$`\alpha `$ distance of 3 fm, 4 fm, and 5 fm. The column noted as $`S=0`$ $`(\sigma )^2`$ gives the energy for the $`\mathrm{\Phi }(1/2^+,1/2^+)`$ configuration, where the linear combination of Gaussians is optimized. The column noted as $`+S_z=1`$ gives the energy when we include the spin-orbit interaction by employing the spin-triplet basis states. The results show that the smaller is the $`\alpha `$-$`\alpha `$ distance, the larger is the contribution of the coupling with the spin-triplet state. When the $`\alpha `$-$`\alpha `$ distance is 5 fm, the coupling increases the binding energy by about 3 MeV; however, when the $`\alpha `$-$`\alpha `$ distance is 3 fm, the coupling increases the binding energy by about 4.5 MeV. Therefore, the coupling with the spin-triplet becomes stronger as the $`\alpha `$-$`\alpha `$ distance becomes small. The column noted as $`+S_z=0,1`$ gives the energy when we further include the spin-singlet basis state to express the tail of the valence neutron. This effect does not have a strong dependence on the $`\alpha `$-$`\alpha `$ distance, and almost constantly decreases the binding energy by 1.5 MeV.
Finally, the calculation where all of these basis states with different $`\alpha `$-$`\alpha `$ distance are superposed is performed based on the generator coordinate method (GCM). When we superpose all of the basis states in Table II, the calculated energy of the second $`0^+`$ state becomes $`54.3`$ MeV. Here, we superpose the basis states with $`\alpha `$-$`\alpha `$ distances of 3, 4, 5, and 6 fm, and the spin-triplet states are represented by one Slater determinant for each $`\alpha `$-$`\alpha `$ distance, which correspond to the minimum point on the energy surface (Fig. 2). This calculated energy of $`54.3`$ MeV is lower than a result calculated with only the $`\mathrm{\Phi }(1/2^+,1/2^+)`$ configuration (employed $`\alpha `$-$`\alpha `$ distances are the same) by 4.6 MeV.
### B Negative-parity states in <sup>10</sup>Be
We examine this spin-mixing effect for the negative-parity states in <sup>10</sup>Be. Here, one of the valence neutrons is introduced to occupy a $`\pi `$-orbit with negative parity; another one occupies a $`\sigma `$-orbit with positive parity. In this way, two configurations with $`S_z=0`$ and $`S_z=1`$ can be constructed. The $`S_z=0`$ basis state ($`\mathrm{\Phi }(3/2^{},1/2^+)`$) has $`K^\pi =1^{}`$.
$$\mathrm{\Phi }(3/2^{},1/2^+)=𝒜[\varphi _1^{(\alpha )}\varphi _2^{(\alpha )}(\varphi _{c1}\chi _{c1})(\varphi _{c2}\chi _{c2})].$$
(16)
$$\varphi _{c1}\chi _{c1}=\{(p_x+ip_y)_{+a}+(p_x+ip_y)_a\}|n,$$
(17)
$$\varphi _{c2}\chi _{c2}=\{(\stackrel{}{p})_{+a}(\stackrel{}{p})_a\}|n.$$
(18)
This basis state gives $`47.7`$ MeV for $`1^{}`$ and $`46.7`$ MeV for $`2^{}`$ at the optimal $`\alpha `$-$`\alpha `$ distance of 3 fm. This energy splitting between $`1^{}`$ and $`2^{}`$ by 1 MeV is much larger than the experimental value of 300 keV (experimentally, the excitation energies of the $`1^{}`$ state and the $`2^{}`$ state are 5.96 MeV and 6.26 MeV, respectively). At first, we only examine the tail effect of the valence neutrons with $`S_z=0`$. We introduce a parameter $`\stackrel{}{R}`$ for $`\varphi _{c2}\chi _{c2}`$, which originally occupies the $`\sigma `$-orbit:
$$\varphi _{c2}\chi _{c2}=G_\stackrel{}{R}|n.$$
(19)
However, even if we superpose states with many different positions ($`\stackrel{}{R}`$), the energy splitting between these states is still about 1 MeV ($`1^{}`$ state is $`49.9`$ MeV and the $`2^{}`$ state is $`48.9`$ MeV). The employed $`\stackrel{}{R}`$ values are 0.0, 1.0, 2.0 (fm) for the $`x`$-direction (perpendicular to the $`\alpha `$-$`\alpha `$ direction) and 0.0, 1.0, 2.0, 3.0, 4.0 (fm) for the $`z`$-direction ($`\alpha `$-$`\alpha `$ direction).
Therefore, to improve this discrepancy, we next include the basis states with $`S_z=1`$; $`\mathrm{\Phi }(3/2^{},1/2^+)`$. The $`S_z=1`$ basis state generates $`K^\pi =2^{}`$ band. As a result of the calculation, the energy of the $`2^{}`$ state with $`K^\pi =2^{}`$ (spin triplet) is found to be accidentally close to the energy of the $`2^{}`$ state with $`K^\pi =1^{}`$ projected from the $`S_z=0`$ basis states. Thus, the coupling effect between the $`S_z=0`$ basis states and the $`S_z=1`$ basis states is larger for the $`2^{}`$ state than for the $`1^{}`$ state, and the energy splitting becomes comparable to the experimental one. To see this effect, we superpose many $`S_z=1`$ basis states characterized by the Gaussian center $`\stackrel{}{R}`$; $`\varphi _{c2}\chi _{c2}=G_\stackrel{}{R}|n`$. When we include many states with different $`\stackrel{}{R}`$ values in our bases states, the energies become $`51.6`$ MeV ($`1^{}`$) and $`51.2`$ MeV ($`2^{}`$). The employed $`\stackrel{}{R}`$ values are 0.1, 1.0, 2.0 (fm) for the $`x`$-direction and 0.0, 1.0, 2.0, 3.0, 4.0 (fm) for the $`z`$-direction. The experimental small level splitting between $`1^{}`$ and $`2^{}`$ is found to be evidence of the spin vibration in the $`2^{}`$ state.
### C The structure of <sup>12</sup>Be
The large contribution of this spin-orbit interaction will be discussed concerning <sup>12</sup>Be. In the previous subsection, we have discussed for <sup>10</sup>Be that this effect is more important as the $`\alpha `$-$`\alpha `$ becomes smaller, and in <sup>12</sup>Be, we show that the distance is smaller than that of the second $`0^+`$ state of <sup>10</sup>Be. In <sup>12</sup>Be, four valence neutrons rotate around two $`\alpha `$ clusters and, mainly, two configurations are important for the $`0^+`$ ground state. One is $`(3/2^{})^2(1/2^{})^2`$ for the four valence neutrons, which corresponds to the closed $`p`$-shell configuration of the neutrons at the zero limit of the $`\alpha `$-$`\alpha `$ distance. The other configuration is $`(3/2^{})^2(1/2^+)^2`$, where two of the four valence neutrons occupy the $`\sigma `$-orbit. It is necessary to compare the energy of these two configurations as a function of the $`\alpha `$-$`\alpha `$ distance. The configurations of $`(3/2^{})^2(1/2^{})^2`$ and $`(3/2^{})^2(1/2^+)^2`$ are constructed as follows:
$$\mathrm{\Phi }(3/2^{},3/2^{},1/2^{},1/2^{})=𝒜[\varphi _1^{(\alpha )}\varphi _2^{(\alpha )}(\varphi _{c1}\chi _{c1})(\varphi _{c2}\chi _{c2})(\varphi _{c3}\chi _{c3})(\varphi _{c4}\chi _{c4})]$$
(20)
$$\varphi _{c1}\chi _{c1}=\{(p_x+ip_y)_{+a}+(p_x+ip_y)_a\}|n,$$
(21)
$$\varphi _{c2}\chi _{c2}=\{(p_xip_y)_{+a}+(p_xip_y)_a\}|n,$$
(22)
$$\varphi _{c3}\chi _{c3}=\{(p_xip_y)_{+a}+(p_xip_y)_a\}|n,$$
(23)
$$\varphi _{c4}\chi _{c4}=\{(p_x+ip_y)_{+a}+(p_x+ip_y)_a\}|n,$$
(24)
$$\mathrm{\Phi }(3/2^{},3/2^{},1/2^+,1/2^+)=𝒜[\varphi _1^{(\alpha )}\varphi _2^{(\alpha )}(\varphi _{c1}\chi _{c1})(\varphi _{c2}\chi _{c2})(\varphi _{c3}\chi _{c3})(\varphi _{c4}\chi _{c4})]$$
(25)
$$\varphi _{c1}\chi _{c1}=\{(p_x+ip_y)_{+a}+(p_x+ip_y)_a\}|n,$$
(26)
$$\varphi _{c2}\chi _{c2}=\{(p_xip_y)_{+a}+(p_xip_y)_a\}|n,$$
(27)
$$\varphi _{c3}\chi _{c3}=\{(\stackrel{}{p})_{+a}(\stackrel{}{p})_a\}|na=d/2,$$
(28)
$$\varphi _{c4}\chi _{c4}=\{(\stackrel{}{p})_{+a}(\stackrel{}{p})_a\}|na=d/2.$$
(29)
The dotted line in Fig. 3 shows $`0^+`$ energies of $`(3/2^{})^2(1/2^+)^2`$ and the dashed line shows $`(3/2^{})^2(1/2^{})^2`$ with respect to the $`\alpha `$-$`\alpha `$ distance.
————-
Fig. 3
————-
When the $`\alpha `$-$`\alpha `$ distance is small, for example 2 fm, the dominant configuration of the four valence neutrons is $`(3/2^{})^2(1/2^{})^2`$ for the ground state, which corresponds to the closed $`p`$-shell configuration at the limit of the $`\alpha `$-$`\alpha `$ distance, zero. On the other hand, the $`(3/2^{})^2(1/2^+)^2`$ configuration for the four valence neutrons becomes lower as the $`\alpha `$-$`\alpha `$ distance is increased. However, it is still higher than the $`(3/2^{})^2(1/2^{})^2`$ configuration by a few MeV. Here, we show the importance of the coupling between the $`(3/2^{})^2(1/2^+)^2`$ configuration and the spin-triplet state for the last two neutrons ($`(1/2^+)^2`$). The spin-triplet state for $`(1/2^+)^2`$ is introduced as follows:
$$\mathrm{\Phi }(3/2^{},3/2^{},1/2^+,G_\stackrel{}{R}|)=𝒜[\varphi _1^{(\alpha )}\varphi _2^{(\alpha )}(\varphi _{c1}\chi _{c1})(\varphi _{c2}\chi _{c2})(\varphi _{c3}\chi _{c3})(\varphi _{c4}\chi _{c4})]$$
(30)
$$\varphi _{c1}\chi _{c1}=\{(p_x+ip_y)_{+a}+(p_x+ip_y)_a\}|n,$$
(31)
$$\varphi _{c2}\chi _{c2}=\{(p_xip_y)_{+a}+(p_xip_y)_a\}|n,$$
(32)
$$\varphi _{c3}\chi _{c3}=\{(\stackrel{}{p})_{+a}(\stackrel{}{p})_a\}|na=d/2,$$
(33)
$$\varphi _{c4}\chi _{c4}=G_\stackrel{}{R}|n.$$
(34)
The coupling between the $`(3/2^{})^2(1/2^+)^2`$ configuration and the spin-triplet states is shown in Fig. 4 as a function of the parameter $`\stackrel{}{R}`$.
————-
Fig. 4
————-
The point on the $`x`$-$`z`$ plain in Fig. 4 shows the Gaussian-center $`\stackrel{}{R}`$ in the spin-triplet state. The contour map shows the energy calculated by taking into account this coupling with the spin-triplet state. This $`\alpha `$-$`\alpha `$ distance is 3 fm, and this optimal distance is smaller than that for the second $`0^+`$ state of <sup>10</sup>Be by about 1 fm, since two of the four valence neutrons occupy the $`(3/2^{})^2`$ configuration and increase the binding energy of the system. Because of this effect, the coupling effect of the spin-triplet states is more important than in the case of the second $`0^+`$ state of <sup>10</sup>Be. The minimal point of the surface shows that the energy gain due to the couplings is about 4 MeV in <sup>12</sup>Be.
The solid line in Fig. 3 shows the $`0^+`$ energy calculated by taking into account the coupling effect between the $`(3/2^{})^2(1/2^+)^2`$ configuration and the $`S_z=0,1`$ basis states for the last two valence neutrons. Due to the spin-orbit coupling, the energy is almost the same as that of $`(3/2^{})^2(1/2^{})^2`$ corresponding to the closed $`p`$-shell configuration. Furthermore, the energy of $`(3/2^{})^2(1/2^+)^2`$ is suggested to become even lower than $`(3/2^{})^2(1/2^{})^2`$ when the pairing effect between $`(3/2^{})^2`$ and $`(1/2^{})^2`$ is taken into account. These effects plays crucial roles in accounting for breaking of the $`N=8`$ magic number.
## IV CONCLUSION
The second $`0^+`$ state of <sup>10</sup>Be has been shown to be characterized by the $`\sigma `$-orbit of the two valence neutrons in terms of the molecular orbit (MO) model. The two valence neutrons stay along the $`\alpha `$-$`\alpha `$ axis (the $`1/2^+`$ orbit) and reduce the kinetic energy by enlarging the $`\alpha `$-$`\alpha `$ distance. This simple description for the second $`0^+`$ state gives a higher excitation energy by 5 MeV compared to the experimental one. To improve the description, spin-triplet basis states which have not been included in the traditional MO models are prepared by allowing a deviation of the neutron orbit from one just along the $`\alpha `$-$`\alpha `$ axis. Thus, the spin-orbit interaction can be taken into account, and the calculated second $`0^+`$ state becomes lower by 3.5 MeV. The precise description of the neutron-tail also decreases the excitation energy by 1.5 MeV. The discrepancy between the experimental excitation energy and the calculated one is compensated by these two effects.
The spin-mixing effects were studied for the negative parity state of <sup>10</sup>Be. If we restrict ourselves to the $`S_z=0`$ basis state, this energy splitting between the $`1^{}`$ state and $`2^{}`$ state is about 1 MeV, which is much larger than the experimental value of 303 keV. As a result of adding basis states with $`S_z=1`$, where the spin directions of the two valence neutrons are the same, the $`K`$-mixing occurs especially for the $`2^{}`$ state. The contribution of the $`S_z=1`$ state is larger for the $`2^{}`$ state. The resultant energy splitting becomes comparable to the experimental one. Therefore, the experimental small level splitting between $`1^{}`$ and $`2^{}`$ is considered to be a result of the spin vibration induced by the spin-orbit interaction in the $`2^{}`$ state.
The coupling with the spin-triplet basis states is also important in the case of <sup>12</sup>Be. Without the spin-triplet basis state, the energy of the configuration $`(3/2^{})^2(1/2^+)^2`$ is much higher than that of the closed $`p`$-shall configuration ($`(3/2^{})^2(1/2^{})^2`$) by 4 MeV. However, the energy of $`(3/2^{})^2(1/2^+)^2`$ is drastically decreased by coupling with the spin-triplet states. This is because the effect becomes stronger as the $`\alpha `$-$`\alpha `$ distance becomes shorter, and <sup>12</sup>Be has an optimal $`\alpha `$-$`\alpha `$ distance around 3 fm, which is smaller than the second $`0^+`$ state of <sup>10</sup>Be by 1 fm. The study shows that an energy of $`(3/2^{})^2(1/2^+)^2`$ is almost the same as $`(3/2^{})^2(1/2^{})^2`$, or even lower. This effect is suggested to play a crucial role in accounting for the dissipation of the $`N=8`$ magic number in <sup>12</sup>Be. It is an interesting subject to analyze the binding mechanism and properties of the ground state by taking into account the pairing mixing among states with configurations of $`(3/2^{})^2(1/2^{})^2`$, $`(3/2^{})^2(1/2^+)^2`$, and $`(1/2^{})^2(1/2^+)^2`$. A detail analysis is going to be performed not only for this state, but also for excited states where new states with cluster structure have been recently observed.
###### Acknowledgements.
The authors would like to thank Prof. I. Tanihata for various effective suggestions. They also thank other members of RI beam science laboratory in RIKEN for discussions and encouragements. One of the authors (N.I) thanks fruitful discussions with Prof. R. Lovas, Prof. H. Horiuchi, Prof. T. Otsuka, Prof. Y. Abe, Prof. K. Katō, Prof. K. Yabana, Dr. A. Ohnishi, and Dr. Y. Kanada-En’yo.
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# Macroscopic Quantum Self-trapping and Atomic Tunneling in Two-species Bose-Einstein Condensates
## I Introduction
Recently, much attention has been paid to the investigations of systems consisting of two weakly interacting Bose-Einstein condensates \[1-17\] due to the appearance of quantum interference \[18-34\] and new macroscopic quantum phenomena \[35-37\]. In principle, such condensate systems can be produced in a double trap with two condensates coupled by quantum tunneling and ground collisions, or in a system with two different magnetic sublevels of an atom, in which case the two species condensates correspond two electronic states involved. The coupling between two condensates could be realized by the near-resonant dipole-dipole interaction.
The first experiment involving interactions between two condensates was performed with atoms evaporatively cooled in the $`|F=2,M_f=2`$ and $`|1,1`$ spin states of $`{}_{}{}^{87}Rb`$. One of the latest experimental advances in this direction is the realization of measurements of relative phase in two-component Bose-Einstein condensates. In experiments at JILA , two condensates in two different internal atomic states are produced by using a single two-photon coupling pulse. The two condensates have a well defined relative phase. After a time during which the condensates evolve in the trapping potentials, the two condensates interfere through mixing coherently the two internal atomic states. Then, the relative phase of the two condensates is obtained from the spatial interference pattern. The realization of measurements of the relative phase between two condensates opened the fascinating possibility of experimentally examing the phase-related phenomena in Bose condensates, such as atomic Josephson effect and macroscopic quantum self-trapping (MQST) .
Theoretical studies of such systems began in Ho and Shenoy’s work which shown that binary mixtures of condensates of alkali atoms have a great variety of ground state and vortex structures. Then, the stability and collective excitations of two-species condensate systems have been extensively studied. More recently, Smerzi and coworkers have shown that the in a system of two Bose condensates the quantum coherent atomic tunneling between two condensates induces two types of interesting effects. One is an atomic Josephson effect in Bose condensates, which is a generalization of the sinusoidal Josephson effects familiar in superconductors. The other is macroscopic quantum self-trapping (MQST), which is a kind of a self-locked population imbalance between two Bose condensates. It arises because of the interatomic nonlinear self-interaction. The MQST has a quantum nature, involving the coherence of a macroscopic number of atoms in the two condensates. It has been known that the MQST depends upon the trap parameters, the total atoms and initial states of the system and is self-maintained in a closed conserved system without external drives. As pointed out in Ref., it is easier to observe the MQST in Bose condensates than self-trapping phenomena in other systems, such as the single-electron Coulomb Blockade effect arising from the Coulomb interaction between electrons, single polaron trapping in a medium which arises from single electrons, interacting with a polarizable lattice, and external gravitational effects on He II baths .
In Ref., Smerzi and coworkers only considered the MQST induced by interatomic nonlinear self-interaction in each condensate. However, in a system consisting of two Bose condensates, there are not only nonlinear self-interaction but also interspecies nonlinear interaction. Questions that naturally arise are, what is the effect of the interspecies nonlinear interaction on the MQST and quantum coherent atomic tunneling? Does interspecies nonlinear interaction strengthen or weaken the MQST and the atomic tunneling current between them?
In this paper, we present a theoretical treatment of the MQST and the quantum coherent atomic tunneling in a more general two-species Bose condensate system in terms of a two-mode approximate model and the rotating wave approximation. Our treatment involves not only the interatomic nonlinear self-interaction in each species but also the interspecies nonlinear interaction. We find that the presence of the interspecies nonlinear interaction gives rise to new insight to the MQST and the atomic tunneling between the two condensates. This paper is organized as follows. In Sec. II, we establish our model and present an approximate analytic solution. In Sec. III, we discuss the collapse and revival (CR) phenomenon on population imbalance between two condensates. In Sec. IV, we investigate the MQST in the two-condensate system, and discuss the dependence of the MQST upon the initial states and the tunneling interaction and the nonlinear interactions. In Sec. V, we study quantum dynamics of the atomic tunneling current and its $`dc`$ characteristics. We shall conclude our paper with discussions and remarks in the last section.
## II Model and Solution
We consider a zero-temperature two-species Bose condensate system in which the atoms interact via $`aa`$ and $`bb`$ and $`ab`$ elastic collisions, and there is a Josephson-like coupling term denoted by $`a^{}b`$ and $`ab^{}`$. In the formalism of the second quantization, Hamiltonian of such a system can be written as
$`\widehat{H}`$ $`=`$ $`\widehat{H}_1+\widehat{H}_2+\widehat{H}_{int}+\widehat{H}_{Jos},`$ (1)
$`\widehat{H}_i`$ $`=`$ $`{\displaystyle }d𝐱\widehat{\psi }_i^{}(𝐱)[{\displaystyle \frac{\mathrm{}^2}{2m}}^2+V_i(𝐱)`$ (3)
$`+U_i\widehat{\psi }_i^{}(𝐱)\widehat{\psi }_i(𝐱)]\widehat{\psi }_i(𝐱),(i=1,2),`$
$`\widehat{H}_{int}`$ $`=`$ $`U_{12}{\displaystyle 𝑑𝐱\widehat{\psi }_1^{}(𝐱)\widehat{\psi }_2^{}(𝐱)\widehat{\psi }_1(𝐱)\widehat{\psi }_2(𝐱)},`$ (4)
$`\widehat{H}_{Jos}`$ $`=`$ $`\mathrm{\Lambda }{\displaystyle 𝑑𝐱[\widehat{\psi }_1^{}(𝐱)\widehat{\psi }_2(𝐱)+\widehat{\psi }_1(𝐱)\widehat{\psi }_2^{}(𝐱)]},`$ (5)
where $`\widehat{\psi }_i(𝐱)`$ and $`\widehat{\psi }_i^{}(𝐱)`$ are the atomic field operators which annihilate and create atoms at position $`𝐱`$, respectively, they satisfy the commutation relation
$$[\widehat{\psi }_i(𝐱),\widehat{\psi }_j^{}(𝐱^{})]=\delta _{ij}\delta (𝐱𝐱^{}).$$
(6)
In Eq.(1), $`\widehat{H}_1`$ and $`\widehat{H}_2`$ describe the evolution of each condensate in the absence of interspecies interaction. $`\widehat{H}_{int}`$ describes interspecies collisions. $`\widehat{H}_{Jos}`$ is the Josephson-like tunneling coupling term. Atoms are confined in harmonic potentials $`V_i(𝐱)(i=1,2)`$ of frequencies $`\omega _i`$. Interactions between atoms are described by a nonlinear self-interaction term $`U_i=4\pi \mathrm{}^2a_i^{sc}/m`$ and a term that corresponds the nonlinear interaction between different condensates $`U_{12}=4\pi \mathrm{}^2a_{12}^{sc}/m`$, where $`a_i^{sc}`$ is $`s`$-wave scattering length of condensate $`i`$ and $`a_{12}^{sc}`$ that between condensate 1 and 2. For simplicity, throughout this paper we let $`\mathrm{}=1`$ and assume that $`a_1^{sc}=a_2^{sc}=a^{sc}`$, and $`V_1(𝐱)=V_2(𝐱)`$.
It is well known that the above Hamiltonian can be reduced to two-mode boson Hamiltonian \[30-32, 35,36\] through expanding the atomic field operators over single-particle states :
$$\widehat{\psi }_i(𝐱)=\widehat{a}_i\varphi _{iN}(𝐱)+\stackrel{~}{\psi }_i(𝐱),$$
(7)
where $`\widehat{a}_i^{}=𝑑𝐱\varphi _{iN}(𝐱)\widehat{\psi }_i^{}(𝐱)`$ create particles with distributions $`\varphi _{iN}(𝐱)`$ with $`[\widehat{a}_i,\widehat{a}_i^{}]=1`$. The first term in the mode expansion (6) acts only on the condensate state vector, whereas the second term $`\stackrel{~}{\psi }_i(𝐱)`$ accounts for noncondensed atoms. Substituting the mode expansions of the atomic field operators into the Hamiltonian (1), retaining only the first term representing the condensates, we arrive at the following two-mode approximate Hamiltonian
$`\widehat{H}`$ $`=`$ $`\omega _0(\widehat{a}_1^{}\widehat{a}_1+\widehat{a}_2^{}\widehat{a}_2)+q(\widehat{a}_1^2\widehat{a}_1^2+\widehat{a}_2^2\widehat{a}_2^2)`$ (9)
$`+g(\widehat{a}_1^{}\widehat{a}_2+\widehat{a}_2^{}\widehat{a}_1)+2\chi \widehat{a}_1^{}\widehat{a}_1\widehat{a}_2^{}\widehat{a}_2,`$
where the frequency and the coupling constants are defined by
$`\omega _0`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{2}{}}}{\displaystyle }d𝐱[{\displaystyle \frac{1}{2m}}|\varphi _{iN}(𝐱)|^2+V(𝐱)(|\varphi _{iN}(𝐱)|^2],`$ (10)
$`q`$ $`=`$ $`U_0{\displaystyle 𝑑𝐱(|\varphi _{1N}(𝐱)|^2+|\varphi _{2N}(𝐱)|^2)},`$ (11)
$`g`$ $`=`$ $`\mathrm{\Lambda }{\displaystyle 𝑑𝐱(\varphi _{1N}^{}(𝐱)\varphi _{2N}(𝐱)+\varphi _{1N}(𝐱)\varphi _{2N}^{}(𝐱))},`$ (12)
$`\chi `$ $`=`$ $`{\displaystyle \frac{1}{2}}U_{12}{\displaystyle 𝑑𝐱|\varphi _{1N}^{}(𝐱)|^2|\varphi _{2N}(𝐱)|^2}.`$ (13)
From Eqs.(6) and (7) we can see that the two-mode approximation essentially consists in neglecting all modes except the condensate modes. At zero temperature, this amounts to ignoring the atoms which have left the condensate mode due to the two-body interactions. In other words, what the two-mode approximation involves is only the first order effects of interactions. The mode expansion of the condensate function over single-particle states (6) makes the condensate shape not to be changed, this limits migration of condensed atoms from one condensate to the other. The constraint on the shapes of condensates implies that the two-mode approximation can be applied only for weak nonlinearity. The valid conditions of the two-mode approximation were demonstrated in Refs., which indicate that this approximation provides a reasonably accurate picture for weak many-body interactions , i.e., for small number of condensed atoms. For large condensates, the mode functions of condensates are altered due to the collisional interactions, and the two-mode approximation breaks down. As shown in Ref., a simple estimate shows that this happens when the number of atoms $`Na^{sc}r_0`$, where $`a^{sc}`$ is a typical scattering length and $`r_0`$ is a measure of the trap size. If we consider a large trap with the size $`r_0=10\mu `$m and the typical scattering length $`a^{sc}=5`$ nm, the two-mode approximation is applicable for $`N2000`$. This is the case which we consider here. We shall show that the MQST and atomic tunneling between the two condensates are strongly affected by the nonlinear many-body interactions.
We note that the two-mode approximate Hamiltonian has the same form with that of a two-mode nonlinear optical directional coupler . The two-mode Hamiltonian (7) can not be exactly solved, but for weak nonlinear interactions a closed analytical solution can be obtained under the rotating wave approximation suggested by Alodjanc et al. .
In order to obtain an approximate analytic solution of the Hamiltonian (7), we introduce a new pair of bosonic operators $`\widehat{A}_1`$ and $`\widehat{A}_2`$ by the following expressions::
$$\widehat{a}_1=\frac{1}{\sqrt{2}}(\widehat{A}_1e^{igt}i\widehat{A}_2e^{igt}),\widehat{a}_2=\frac{1}{\sqrt{2}}(\widehat{A}_1e^{igt}+i\widehat{A}_2e^{igt}),$$
(14)
where $`\widehat{A}_1`$ and $`\widehat{A}_2`$ are slowly varying operators, they satisfy the usual bosonic commutation relations: $`[\widehat{A}_i,\widehat{A}_j]=0`$, and $`[\widehat{A}_i,\widehat{A}_j^{}]=\delta _{ij}`$ with $`\widehat{A}_i^{}`$ being the hermitain conjugation of $`\widehat{A}_j`$. Then the Hamiltonian (7) reduces to the following form
$`\widehat{H}`$ $`=`$ $`\omega \widehat{N}+{\displaystyle \frac{1}{4}}q[3\widehat{N}^2(\widehat{A}_1^{}\widehat{A}_1\widehat{A}_2^{}\widehat{A}_2)^2]`$ (17)
$`+g(\widehat{A}_1^{}\widehat{A}_1\widehat{A}_2^{}\widehat{A}_2)+{\displaystyle \frac{\chi }{2}}\widehat{N}^2`$
$`\chi \widehat{A}_1^{}\widehat{A}_1\widehat{A}_2^{}\widehat{A}_2+\widehat{H}^{},`$
where the detuning is given by $`\omega =\omega _0(\chi +q)/2`$, the total number operator $`\widehat{N}`$ is a conserved constant which is given by
$$\widehat{N}=\widehat{a}_1^{}\widehat{a}_1+\widehat{a}_2^{}\widehat{a}_2=\widehat{A}_1^{}\widehat{A}_1+\widehat{A}_2^{}\widehat{A}_2,$$
(18)
and $`\widehat{H}^{}`$ is a nonresonant term which is given by
$$\widehat{H}_1^{}=\frac{1}{2}(\chi q)(\widehat{A}_1^2\widehat{A}_2^2e^{i4gt}+\widehat{A}_2^2\widehat{A}_1^2e^{i4gt}),$$
(19)
which oscillates at the frequency $`4g`$. The account of the fast oscillating term results only in some addtional oscillations which play no essential role in the evolution of the measurable quantities specifying the macroscopic quantum phenomena of the two-condensate system, so that it is fully negligible. This means the rotating wave approximation . After neglecting the nonresonant term $`H^{}`$, we get the following approximate Hamiltonian:
$`\widehat{H}_A`$ $`=`$ $`\omega \widehat{N}+g(\widehat{A}_1^{}\widehat{A}_1\widehat{A}_2^{}\widehat{A}_2)`$ (22)
$`+{\displaystyle \frac{1}{4}}q[3\widehat{N}^2(\widehat{A}_1^{}\widehat{A}_1\widehat{A}_2^{}\widehat{A}_2)^2]`$
$`+{\displaystyle \frac{1}{2}}\chi \widehat{N}^2\chi \widehat{A}_1^{}\widehat{A}_1\widehat{A}_2^{}\widehat{A}_2.`$
In order to solve the Hamiltonian (16) we introduce two Fock spaces of $`(\widehat{A}_1,\widehat{A}_2)`$ and $`(\widehat{a}_1,\widehat{a}_2)`$ in which the bases are defined by
$`|n,m)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{n!m!}}}\widehat{A}_1^n\widehat{A}_2^m|0,0),`$ (23)
$`|n,m`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{n!m!}}}\widehat{a}_1^n\widehat{a}_2^m|0,0.`$ (24)
where $`n`$ and $`m`$ take non-negative integers. Obviously, $`\widehat{H}_A`$ is diagonal in the Fock space of $`(\widehat{A}_1,\widehat{A}_2)`$, and we find that
$`\widehat{H}_A|n,m)`$ $`=`$ $`E(n,m)|n,m),`$ (25)
$`E(n,m)`$ $`=`$ $`\omega (n+m)+g(nm)`$ (28)
$`+{\displaystyle \frac{1}{4}}(3q+\chi )(n+m)^2{\displaystyle \frac{1}{4}}q(nm)^2`$
$`\chi nm.`$
Consider two coherent states defined in Fock spaces of $`(\widehat{A}_1,\widehat{A}_2)`$ and $`(\widehat{a}_1,\widehat{a}_2)`$, respectively,
$`|\alpha _1,\alpha _2`$ $`=`$ $`D_{\widehat{a}_1}(\alpha _1)D_{\widehat{a}_2}(\alpha _2)|0,0,`$ (29)
$`|u_1,u_2)`$ $`=`$ $`D_{\widehat{A}_1}(u_1)D_{\widehat{A}_2}(u_2)|0,0),`$ (30)
where $`D_{\widehat{a}_i}(\alpha _i)`$ and $`D_{\widehat{A}_i}(u_i)`$ are displacement operators defined by
$`D_{\widehat{a}_i}(\alpha _i)`$ $`=`$ $`\mathrm{exp}(\alpha _i\widehat{a}_i+\alpha _i^{}\widehat{a}_i^{}),`$ (31)
$`D_{\widehat{A}_i}(u_i)`$ $`=`$ $`\mathrm{exp}(u_i\widehat{A}_i+u_i^{}\widehat{A}_i^{}).`$ (32)
Note the fact that $`|0,0)=|0,0`$, we can find a useful relation to connect $`|\alpha _1,\alpha _2`$ and $`|u_1,u_2)`$ with each other
$`|\alpha _1,\alpha _2`$ $`=`$ $`|{\displaystyle \frac{\alpha _1+\alpha _2}{\sqrt{2}}},{\displaystyle \frac{i(\alpha _1\alpha _2)}{\sqrt{2}}}),`$ (33)
$`|\alpha _1,\alpha _2)`$ $`=`$ $`|{\displaystyle \frac{\alpha _1i\alpha _2}{\sqrt{2}}},{\displaystyle \frac{\alpha _1+i\alpha _2}{\sqrt{2}}}.`$ (34)
Following the arguments of Bose broken symmetry, we assume that the two condensates are initially in the coherent states $`|\alpha _1`$ and $`|\alpha _2`$, which are eigenstates of $`\widehat{a}_1`$ and $`\widehat{a}_2`$, respectively. Then the wave function of the two species condensate system at time $`t`$ can be explicitly expressed as
$`|\mathrm{\Phi }(t)`$ $`=`$ $`e^{\frac{1}{2}N}{\displaystyle \underset{n,m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{\sqrt{n!m!}}}u_1^n(iu_2)^m`$ (36)
$`\times e^{iE(n,m)t}|n,m),`$
where
$`u_1`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(\alpha _1+\alpha _2),u_2={\displaystyle \frac{1}{\sqrt{2}}}(\alpha _1\alpha _2)`$ (37)
$`N`$ $`=`$ $`|\alpha _1|^2+|\alpha _2|^2=|u_1|^2+|u_2|^2,`$ (38)
where we have used Eqs.(25), (26) in the derivation of Eq.(27).
## III collapse and revivals of Population imbalance
In this section we show that the two condensate system under our consideration exhibits a collapse and revival phenomenon of population imbalance between two condensates. Denote the number difference of atoms between the two condensates by
$$D(t)=N_1(t)N_2(t).$$
(39)
Then from Eq.(27) we can find that at time $`t`$, the number of atoms in each condensate $`N_i(t)=\widehat{a}_i^{}\widehat{a}_i`$ is given by
$`N_i(t)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\{N(1)^i2|u_1||u_2|\mathrm{cos}[4gt+\theta (t)]`$ (41)
$`\times e^{2N\mathrm{sin}^2\frac{1}{2}(q\chi )t}\},(i=1,2)`$
where we have used the following symbols:
$$u_i=|u_i|e^{i\phi _{u_i}},\theta (t)=(\phi _{u_2}\phi _{u_1})+u_{21}\mathrm{sin}(q\chi )t,$$
(42)
with $`u_{21}`$ and $`\phi _{u_i}`$ being defined by
$`u_{21}`$ $`=`$ $`|u_2|^2|u_1|^2,`$ (43)
$`\phi _{u_i}`$ $`=`$ $`\mathrm{tan}^1\{{\displaystyle \frac{|\alpha _1|\mathrm{sin}\phi _1(1)^i|\alpha _2|\mathrm{sin}\phi _2}{|\alpha _1|\mathrm{cos}\phi _1(1)^i|\alpha _2|\mathrm{cos}\phi _2}}\}.`$ (44)
Then, the population difference is given by
$`D(t)`$ $`=`$ $`2|u_1||u_2|\mathrm{cos}[4gt+\theta (t)]`$ (46)
$`\times e^{2N\mathrm{sin}^2(\frac{1}{2}(q\chi )t},`$
where $`N=|\alpha _1|^2+|\alpha _2|^2`$ is the total number of the atoms in the two condensates. Eq.(35) indicates that the population imbalance periodically oscillates with the time evolution. From Eq.(35) we can see that $`D(t)`$ exhibits collapse and revival phenomenon which is a kind of nonclassical effect well known in the Jaynes-Cummings model to describe interaction between a single-mode radiation field and a two-level atom. The CR is also found in a Bose condensate system . From Eq.(35) we see that the CR of the population imbalance in the two-condensate system depends on the tunneling interaction ($`g`$) and interatomic nonlinear interactions ($`q`$ and $`\chi `$). When $`g>|q\chi |/8`$, since the function $`\mathrm{cos}[4gt+\theta (t)]`$ is the rapidly varying part in (35), so that the shape of the CR is determined by the envelope function $`\mathrm{exp}[2N\mathrm{sin}^2(\frac{1}{2}(q\chi )t]`$. The maximal revivals take place at time $`t=2n\pi /|q\chi |`$, where $`n`$ is an integer. When $`g<|q\chi |/8`$, the function $`\mathrm{exp}[2N\mathrm{sin}^2(\frac{1}{2}(q\chi )t]`$ becomes the rapidly varying part in (35), the CR then is determined by the envelope function $`\mathrm{cos}[4gt+\theta (t)]`$. In Fig. 1 we plot the evolution of the population difference between the two condensates with respect to the time which is in units of $`|q\chi |`$, when the two condensates are in the initial state of $`\alpha _1=5`$ and $`\alpha _2=4`$, and the tunneling coupling is $`g=25|q\chi |`$. Fig. 1 clearly indicates the CR phenomenon of the population difference.
It is worthwhile to note that when $`q=\chi `$, the population difference (35) becomes
$$D(t)=2|u_1||u_2|\mathrm{cos}[4gt+(\phi _{u_2}\phi _{u_1})],$$
(47)
which is a simple sinusoidal oscillation, no CR occurs. The suppression of the CR can be explained by looking at the expression (35) which indicates that both the self-interaction ($`q`$) and the interspecies interaction ($`\chi `$) can induce the CR, but the CR produced by one can weaken that by another. It is the CR produced by the self-interaction completely counteracts the CR by the interspecies interaction that leads to the suppression of the CR.
When there is no nonlinear interactions, i.e., $`q=\chi =0`$, from Eq.(35) it is easy to find that the population imbalance has the same form with that of the case $`q=\chi 0`$. This means that the CR vanishes when the nonlinearity vanishes. Hence, the CR of the population imbalance is a consequence of nonlinear interactions in condensates. The CR of the oscillatory transfer of atoms between the two condensates constitutes a novel macroscopic quantum phenomenon induced by interatomic nonlinear interactions for the two species condensate system.
## IV Macroscopic Quantum Self-trapping
In this section we are concerned with the MQST. The MQST effect is characterized by the nonzero time mean value of the fractional population imbalance between the two condensates defined by
$$p(t)=\frac{N_1(t)N_2(t)}{N}.$$
(48)
From Eqs.(35) and (37) we get that
$`p(t)`$ $`=`$ $`{\displaystyle \frac{2|u_1||u_2|}{N}}\mathrm{cos}[4gt+\theta (t)]`$ (50)
$`\times e^{2N\mathrm{sin}^2(\frac{1}{2}(q\chi ))t}.`$
In order to investigate the MQST, we expand the above equation as
$`p(t)`$ $`=`$ $`{\displaystyle \frac{2|u_1||u_2|}{N}}e^N{\displaystyle \underset{n,m=\mathrm{}}{\overset{+\mathrm{}}{}}}J_n(u_{21})I_m(N)`$ (53)
$`\times \mathrm{cos}\{[(n+m)(q\chi )+4g]t`$
$`+(\phi _{u_2}\phi _{u_1})\},`$
where $`J_n(A)`$ and $`I_n(A)`$ are Bessel function and modified Bessel function. From Eq.(35) it is easy to find that when the tunneling interaction and nonlinear interactions satisfy the condition:
$$4g=K(\chi q),$$
(54)
where $`K`$ is an integer, we can get a nonzero time-averaged value of population imbalance
$`\overline{p}`$ $`=`$ $`{\displaystyle \frac{2|u_1||u_2|}{N}}e^N{\displaystyle \underset{n=\mathrm{}}{\overset{+\mathrm{}}{}}}J_n(u_{21})I_{Kn}(N)`$ (56)
$`\times \mathrm{cos}(\phi _{u_2}\phi _{u_1}),`$
which indicates the existence of the MQST. Eq.(40) is the condition under which the MQST happens.
From Eq.(39) we can see that when the tunneling interaction vanishes and nonlinear self-interaction equals nonlinear interspecies interaction, i.e., $`g=0`$ and $`q=\chi `$, we arrive at a constant population imbalance
$$p(t)=\frac{2|u_1||u_2|}{N}\mathrm{cos}(\phi _{u_2}\phi _{u_1}).$$
(57)
This is a time-independent state, called the self-trapping stationary state, which is the consequence of competing between nonlinear self-interaction and nonlinear interspecies interaction.
When there exists the tunneling coupling, i.e., $`g0`$, from Eq.(40) we can find the critical value of the tunneling coupling at which the MQST happens $`g_c=|q\chi |/4`$. This critical value $`g_c`$ depends upon only the difference between the nonlinear self-interaction and interspecies nonlinear interaction, not the nonlinear self-interaction and interspecies nonlinear interaction themselves. Therefore, it becomes possible that the MQST occurs only when the tunneling coupling equals or exceeds the critical value $`g_c`$.
In Fig.2 we plot the time evolution of the fractional population imbalance when the two condensates are in the initial state of $`\alpha _1=10`$ and $`\alpha _2=0`$ for (a) $`K=1`$, (b) $`K=20`$. Here the time is in units of $`|q\chi |`$. From Fig.2 (a) and (b) we can see that the weaker the tunneling coupling $`(g)`$, the more apparent the MQST becomes. This indicates that the MQST is an effect induced by interatomic nonlinear interactions $`q`$ and $`\chi `$ not the tunneling interaction $`(g)`$. In what follows we shall discuss the dependence of the MQST in detail upon the initial states and nonlinear interactions for specific cases.
A. The initial state dependence
We here discuss the dependence of the MQST on the initial states of the two condensates in the following four cases.
Case 1: $`|\alpha _1|=|\alpha _2|`$, $`\phi _1\phi _2`$. In this case, the two condensates initially have the same number of atoms but different phases. From Eq.(32) and (33) we get that $`|u_1|^2=N\mathrm{cos}^2[(\phi _1\phi _2)/2]`$, $`|u_2|^2=N\mathrm{sin}^2[(\phi _1\phi _2)/2]`$, and $`u_{21}=N\mathrm{cos}(\phi _1\phi _2)`$. Making use of Eq.(28), from Eqs.(35) and we can see that the fractional population periodically evolves with respect to time, $`p(t)0`$ except that $`\phi _1\phi _2=n\pi `$, where $`n`$ is an integer. If $`4g/(\chi q)=K`$ (an integer), we get the locked population imbalance
$`\overline{p}`$ $`=`$ $`|\mathrm{sin}(\phi _1\phi _2)|\mathrm{cos}(\phi _{u_1}\phi _{u_2})e^N`$ (59)
$`\times {\displaystyle \underset{n=\mathrm{}}{\overset{+\mathrm{}}{}}}J_n(N\mathrm{cos}(\phi _1\phi _2))I_{Kn}(N),`$
which implies that the MQST does exist, even if the two condensates initially have the same number of the atoms provided that they have different initial phases. And less the total atomic number is, stronger the MQST becomes.
Case 2: $`|\alpha _1||\alpha _2|`$, $`\phi _1=\phi _2`$. In this case the two condensates initially have the same phases but the different number of atoms. From Eq.(38) we see that when $`4g/(\chi q)=K`$ (an integer), the MQST occurs with the following locked population imbalance:
$$\overline{p}=\frac{|\alpha _1|^2+|\alpha _2|^2}{2N}e^N\underset{n=\mathrm{}}{\overset{+\mathrm{}}{}}J_n(2|\alpha _1|\alpha _2|)I_{Kn}(N).$$
(60)
Case 3: $`N=|\alpha _1|^2`$, $`|\alpha _2|=0`$. In this case, the system starts with all atoms being in one condensate. Making use of Eq.(28), we get that $`u_1=u_2=\alpha _1/\sqrt{2}`$, and $`\theta (t)=0`$. The fractional population evolution is given by
$$p(t)=e^N\underset{n=\mathrm{}}{\overset{+\mathrm{}}{}}I_n(N)\mathrm{cos}\{[n(q\chi )+4g]t,$$
(61)
So that when the tunneling interaction and the nonlinear interactions satisfy the relation $`4g/(\chi q)=K`$ (an integer), we can see the appearance of the MQST with the locked population imbalance
$$\overline{p}=e^NI_K(N).$$
(62)
Case 4: $`|\alpha _1|=|\alpha _2|`$ and $`\phi _1=\phi _2`$. In this case the two condensates initially have the same number of atoms and the same phases. From Eq.(28) we have $`u_2=0`$. Making use of Eq.(38) we can see that the oscillations of the population imbalance vanish, i.e., $`p(t)=0`$, and no MQST occurs. This is in agreement with the result in Ref..
B. The dependence on nonlinear interactions
Then, we turn to the dependence of the MQST upon the tunneling coupling ($`g`$) and the nonlinear interactions between atoms, which are described by the parameters $`q`$ and $`\chi `$ corresponding to self-interactions and interspecies interactions, respectively.
Case 1: $`g=0`$, $`q0`$, and $`\chi 0`$. In this case there is no tunneling interaction, but there exists interatomic nonlinear interactions. From Eq.(39) we find that the fractional population imbalance becomes
$`p(t)`$ $`=`$ $`{\displaystyle \frac{2|u_1||u_2|}{N}}e^N{\displaystyle \underset{n,m=\mathrm{}}{\overset{+\mathrm{}}{}}}J_n(u_{21})I_m(N)`$ (64)
$`\times \mathrm{cos}\{[(n+m)(q\chi )]t+(\phi _{u_2}\phi _{u_1})\},`$
which indicates that no MQST occurs if the coupling of interatomic self-interaction does not equals that of interspecies nonlinear interaction, i.e., $`q\chi `$. However, when the self-interaction equals the interspecies interaction, i.e., $`q=\chi `$, we can observe the self-trapping stationary state with a constant population imbalance $`p(t)=p(0)`$.
Case 2: $`g0`$, $`q=\chi =0`$. In this case, we consider only the effect of the tunneling coupling while the interatomic nonlinear interactions are not involved. Eq.(38) tells us that the MQST does not occur, although there exists oscillations of the population imbalance between the two condensates. This further confirms the validity of Smerzi et al.’s conclusion which the MQST arises from the interatomic nonlinear interaction.
Case 3: $`g0`$, and $`q=\chi 0`$. In this case there exist both the tunneling interaction and the nonlinear interactions, but self-interaction equals interspecies interaction. Frme Eq.(38) we can find that the population imbalance exhibits a simple oscillation with
$$p(t)=\frac{2|u_1||u_2|}{N}e^N\mathrm{cos}[4gt+(\phi _{u_2}\phi _{u_1})],$$
(65)
which means that the MQST vanishes.
Case 4: $`g0`$, $`q0`$, $`\chi =0`$, or $`g0`$, $`q=0`$, $`\chi 0`$. In this case, there exist the tunneling interaction and one of the self-interaction and the interspecies interaction. It is easy to see that the fractional population imbalance Eq.(38) reduces to
$`p(t)`$ $`=`$ $`{\displaystyle \frac{|u_1||u_2|}{N}}e^N{\displaystyle \underset{n,m=\mathrm{}}{\overset{+\mathrm{}}{}}}J_n(u_{21})I_m(N)`$ (67)
$`\times \mathrm{cos}\{[(n+m)\kappa +4g]t+(\phi _{u_2}\phi _{u_1})\},`$
where $`\kappa =q`$ or $`\chi `$. Eq.(49) reflects the fact that when $`4g/\kappa =n+m=K`$ (an integer), the MQST happens with the nonzero $`\overline{p}`$ given by Eq.(41). This implies that both the nonlinear self-interaction in each condensate and the interspecies nonlinear interactions contribute to the MQST. Since the values of $`K`$ to determine $`\overline{p}`$ have opposite signs for the self-interaction ($`\chi `$) and the interspecies interaction ($`q`$), the MQST produced by the self-interaction can weaken that by the interspecies interaction. It is the competition between the MQST induced by the nonlinear self-interactions of each condensate and that by the interspecies nonlinear interaction that leads to the quenching of the MQST in the above case 3.
## V The Coherent Atomic Tunneling Current
In this section, we study quantum dynamics of the coherent atomic tunneling current between two condensates and its $`dc`$ characteristics, and discuss the influence of the initial state of condensates and the tunneling interaction and the nonlinear interactions. The coherent atomic tunneling current between the two condensates is defined by $`I(t)=\dot{N}_1(t)\dot{N}_2(t)`$. Making use of Eq.(31), it is straightforward to get that
$`I(t)`$ $`=`$ $`2|u_1||u_2|\{4g\mathrm{sin}(\theta (t))`$ (70)
$`+(q\chi )[|u_1|^2\mathrm{sin}((q\chi )t\theta (t))`$
$`+|u_2|^2\mathrm{sin}((q\chi )t+\theta (t))]\}.`$
This indicates that the atomic tunneling current periodically changes, atoms periodically transfer between the two condensates with the time evolution.
In order to see $`dc`$ characteristic of the atomic tunneling current, we expand the atomic tunneling current (50) as the following expression
$`I(t)`$ $`=`$ $`2|u_1||u_2|e^N{\displaystyle \underset{n,m=\mathrm{}}{\overset{+\mathrm{}}{}}}\{8gJ_n(u_{21})`$ (74)
$`(q\chi )[|u_1|^2J_{n1}(u_{21})|u_2|^2J_{n+1}(u_{21})]\}`$
$`\times I_m(N)\mathrm{sin}\{[(m+n)(q\chi )+4g]t`$
$`+(\phi _{u_2}\phi _{u_1})\},`$
which implies that when the tunneling coupling, and nonlinear couplings satisfy the condition:
$$4g=K(\chi q),$$
(75)
we get the $`dc`$ component of the atomic tunneling current with the following form,
$`I_{dc}(K)`$ $`=`$ $`2|u_1||u_2|(q\chi )\mathrm{sin}(\phi _{u_2}\phi _{u_1})e^N`$ (78)
$`\times {\displaystyle \underset{n=\mathrm{}}{\overset{+\mathrm{}}{}}}\{2KJ_n(u_{21})[|u_1|^2J_{n1}(u_{21})`$
$`|u_2|^2J_{n+1}(u_{21})]\}I_{Kn}(N),`$
where $`K`$ is an integer. This indicates that the $`dc`$ component of the atomic tunneling current exhibits a step structure with respect to the integer $`K`$. This step structure is a resonant phenomenon among the tunneling interaction and nonlinear interactions with the resonant condition given by Eq.(51). It is the analogue of the Shapiro steps observed in the superconductor Josephson junction , so that we call the steps in the step structure of the $`dc`$ component of the atomic tunneling current the Shapiro-like steps. In what follows we discuss in detail the dependence of the atomic tunneling current and the Shapiro-like steps upon the initial states and nonlinear interactions for some specific cases.
A. The initial state dependence
In this subection we discuss the initial-state dependence of the atomic tunneling current and the Shapiro-like steps for the following four cases.
Case 1: $`|\alpha _1|=|\alpha _2|`$, $`\phi _1\phi _2`$. In this case, the two condensates initially have the same number of atoms but different phases. From Eq.(50) we can get the expression of the atomic tunneling current
$`I(t)`$ $`=`$ $`N|\mathrm{sin}(\phi _1\phi _2)|\{4g\mathrm{sin}(\theta (t))`$ (81)
$`+(q\chi )N[\mathrm{cos}^2{\displaystyle \frac{1}{2}}(\phi _1\phi _2)\mathrm{sin}((q\chi )t\theta (t))`$
$`+\mathrm{sin}^2{\displaystyle \frac{1}{2}}(\phi _1\phi _2)\mathrm{sin}((q\chi )t+\theta (t))]\},`$
where
$$\theta (t)=(\phi _{u_1}\phi _{u_2})N\mathrm{cos}(\phi _1\phi _2)\mathrm{sin}((q\chi )t.$$
(82)
Form Eq.(54) we can obtain the $`dc`$ component of the tunneling current with the following result,
$`I_{dc}(K)`$ $`=`$ $`N(q\chi )|\mathrm{sin}(\phi _1\phi _2)|\mathrm{sin}(\phi _{u_2}\phi _{u_1})`$ (87)
$`\times e^N{\displaystyle \underset{n=\mathrm{}}{\overset{+\mathrm{}}{}}}\{2KJ_n(N\mathrm{cos}(\phi _1\phi _2))`$
$`N[\mathrm{cos}^2{\displaystyle \frac{1}{2}}(\phi _1\phi _2)J_{n1}(N\mathrm{cos}(\phi _1\phi _2))`$
$`\mathrm{sin}^2{\displaystyle \frac{1}{2}}(\phi _1\phi _2)J_{n+1}(N\mathrm{cos}(\phi _1\phi _2))]\}`$
$`\times I_{Kn}(N),`$
where $`n`$ is an integer. Eq.(56) gives rise to the Shapiro-like steps of the atomic tunneling current.
It is interesting to note that when the initial phases of two condensates satisfy the condition: $`\phi _1\phi _2=n\pi `$, where $`n`$ is an integer, we can find zero atomic tunneling current. This means that the blockade of the atomic tunneling happens.
When the initial phases satisfy the relation: $`\phi _1\phi _2=(2n+1)\pi /2`$, $`n`$ being an integer, we find that
$`I(t)`$ $`=`$ $`4gN\mathrm{sin}(\phi _{u_2}\phi _{u_1})`$ (89)
$`+(q\chi )N\mathrm{cos}(\phi _{u_2}\phi _{u_1})\mathrm{sin}(2(q\chi )t),`$
which indicates that the atomic tunneling current is a simple superpositon of a alternating current with the sinusoidal oscillations and a $`dc`$ current. The $`dc`$ component is
$$I_{dc}=4gN\mathrm{sin}(\phi _{u_2}\phi _{u_1}),$$
(90)
which implies that the $`dc`$ component of the atomic tunneling current depends only upon the tunneling coupling, it is independent of the nonlinear interactions in two condensates, and increases linearly with the tunneling strength $`(g)`$ and the the total number of the atoms $`(N)`$. In this case no Shapiro-like steps appears.
Case 2: $`|\alpha _1||\alpha _2|`$, $`\phi _1=\phi _2`$. In this case, the two condensates initially have the same phases but the different number of atoms. The atomic tunneling current is given by
$`I(t)`$ $`=`$ $`|\alpha _1^2\alpha _2^2|\{4g\mathrm{sin}(\theta (t))`$ (93)
$`+{\displaystyle \frac{1}{2}}(q\chi )[(\alpha _1+\alpha _2)^2\mathrm{sin}((q\chi )t\theta (t))`$
$`+(\alpha _1\alpha _2)^2\mathrm{sin}(q\chi )t+\theta (t))]\},`$
where
$$\theta (t)=2\alpha _1\alpha _2\mathrm{sin}(q\chi )t.$$
(94)
And the $`dc`$ component of the atomic tunneling current has the following form,
$`I_{dc}(K)`$ $`=`$ $`|\alpha _1^2\alpha _2^2|(q\chi )\mathrm{sin}(\phi _{u_2}\phi _{u_1})`$ (98)
$`\times e^N{\displaystyle \underset{n=\mathrm{}}{\overset{+\mathrm{}}{}}}\{2KJ_n(2\alpha _1\alpha _2)`$
$`{\displaystyle \frac{1}{2}}[(\alpha _1\alpha _2)^2J_{n1}(2\alpha _1\alpha _2)`$
$`(\alpha _1\alpha _2)^2J_{n+1}(2\alpha _1\alpha _2)]\}I_{Kn}(N).`$
In FIG. 3, we plot the time evolution of the atomic tunneling current between the two condensates. Results are shown for the case of $`g=0.25`$ and $`q\chi =0.1`$ when the two condensates are in the initial state with $`\alpha _1=5`$ and $`\alpha _2=4`$. FIG. 3 indicates that the atomic tunneling current exhibits complicated oscillating behaviors.
Case 3: $`N=|\alpha _1|^2`$, $`|\alpha _2|=0`$. In this case, the system starts with all atoms being in one condensate. Taking into account Eqs.(32) and (33), from Eq.(50) we find the atomic tunneling current to be
$$I(t)=(q\chi )N^2\mathrm{sin}(q\chi )t,$$
(99)
which is a pure sinusoidal alternating atomic current with the period $`T=2\pi /|q\chi |`$ without $`dc`$ component. It is worthwhile to note that the atomic tunneling current (62) depends on only the difference between the self-coupling $`q`$ and the interspecies coupling $`\chi `$, it is independent of the tunneling coupling $`g`$ at all. Thus, we can conclude that the nonlinear interactions can induce the atomic tunneling, even if there is no tunneling coupling between two condensates.
Case 4: $`|\alpha _1|=|\alpha _2|`$ and $`\phi _1=\phi _2`$. In this case, the two condensates initially have the same number of atoms and phases. we see that the atomic tunneling current vanishes, and no Shapiro-like step occurs.
The above analyses indicate that the atomic tunneling current and the Shapiro-like steps strongly depend on the initial number of atoms in each condensate and the initial phase difference between the two condensates.
B. The dependence on nonlinear interactions
In what follows we show that the interatomic nonlinear interactions also significantly affect the atomic tunneling current and the Shapiro-like steps.
If we consider only the influence of the tunneling coupling while the interatomic nonlinear interactions are not involved, i.e., $`g0`$, $`q=\chi =0`$, we find that
$$I(t)=8g|u_1||u_2|\mathrm{sin}[4gt+(\phi _{u_2}\phi _{u_1})].$$
(100)
This atomic tunneling current is a pure sinusoidally alternating atomic current with the period $`T=\pi /2g`$. without $`dc`$ component. Hence no Shapiro-like steps appears.
It is interesting to note that when the interatomic nonlinear interactions are involved, but $`q=\chi 0`$, we get the same results with those of the case of $`q=\chi =0`$ for the atomic tunneling current and the Shapiro-like steps. This indicates that the contributions of the nonlinear self-interaction ($`q`$) in each condensate to the atomic tunneling can counteract that of the interspecies nonlinear interactions ($`\chi `$). It is the competition between the nonlinear self-interaction and the interspecies nonlinear interaction that leads to a simple form of the atomic tunneling current and the suppression of the Shapiro-like steps.
In particular, we note that when $`g=0`$, i.e., there is no tunneling coupling, and $`q\chi `$, we find that
$`I(t)`$ $`=`$ $`2|u_1||u_2|(q\chi )[|u_1|^2\mathrm{sin}((q\chi )t\theta (t))`$ (102)
$`+|u_2|^2\mathrm{sin}((q\chi )t+\theta (t))],`$
where $`\theta (t)`$ is given by Eq.(32). In FIG. 4, we display the time evolution of the atomic tunneling current (64) when the tunneling coupling vanishes. Results are shown for the case of $`q\chi =0.1`$ when the two condensates are in the initial state with $`\alpha _1=10`$ and $`\alpha _2=5`$. From FIG.4 we see that the evolution of the atomic tunneling current exhibits the CR phenomenon.
Eq.(64) implies that the atomic tunneling current is nonzero, but no the Shapiro-like steps appears. From Eq.(64) we can see that even though there is no the tunneling coupling between the two condensates $`(g=0)`$, the atomic tunneling between the two condensates may happen. This atomic tunneling is completely induced by the nonlinearity of interatomic interactions which are characterized by interatomic collisions ($`q`$ and $`\chi `$). Therefore, we may conclude that the nonlinearity of interatomic interactions in the two condensates can lead to the atomic tunneling between the two condensates.
## VI Concluding Remarks
We have studied the MQST and the quantum coherent atomic tunneling in a two species Bose condensate system in the presence of nonlinear self-interaction of each species, the interspecies nonlinear interaction, and the Josephson-like tunneling interaction, and have given new insight to the MQST and the atomic tunneling. We have shown that the interatomic nonlinear interactions in the two condensates induce not only the MQST but also the CR of the population difference between two condensates, The CR phenomenon can be considered as a novel macroscopic quantum effect. We have indicated that the nonlinear interactions significantly affect the atomic tunneling, and the Shapiro-like steps of the atomic tunneling current. Comparing with Smerzi and coworkers’ work , the present work involves the interspecies nonlinear interaction. The involvement of the interspecies nonlinear interaction gives rise to new characteristics on the MQST and the atomic tunneling. We have shown that when both the nonlinear self-interaction $`(q)`$ and the interspecies nonlinear interaction $`(\chi )`$ present at the same time, the atomic tunneling dynamics and the MQST and the Shapiro-like steps depends upon the difference $`(q\chi )`$, not $`q`$ and $`\chi `$ themselves. We have also found that the interspecies nonlinear interaction generates the MQST at the same level with nonlinear self-interaction. However, contribution from the interspecies interaction to the MQST and that from the self-interaction weaken themselves with each other. It is the competing effects between the nonlinear self-interaction in each species and the interspecies nonlinear interaction that leads to the quenching of the MQST and the suppression of the CR and the Shapiro-like steps of the atomic tunneling current. Especially, we have revealed that the nonlinearity of interatomic interactions in the two condensates can induce the coherent atomic tunneling between two condensates occurs, even though there does not exist the Josephson-like tunneling coupling. It should be mentioned that these results are obtained under the two-mode approximation and the rotating wave approximation, so they are valid for weak nonlinear interactions between atoms. Finally, It should be noted that in order to observe macroscopic quantum phenomena such as MQST and Shapiro-like steps in Bose condensates, one has to control various interactions between atoms. This controlling can be carried out through manipulating interatomic scattering lengthes. Several theoretical and experimental approaches \[46-51\] to alter the scattering length have been proposed. In particular, recent experiments on Feshbach resonances in a Bose condensate have indicated that the scattering length of ultracold atoms can be altered through Feshbach resonance. These experimental progresses provide the possibility to observe the MQST and Shapiro-like steps in Bose condensates.
Acknowledgments
L. M. Kuang would like to acknowledge the Abdus Salam International Center for Theoretical Physics, Trieste, for its hospitality where part of this work was done. This work was supported in part by the climbing project of China, NSF of China, Educational Committee Foundation and NSF of Hunan Province, special project of NSF of China via Institute of theoretical Physics, Academia Sinica.
Figure Captions
FIG. 1. Diagram of the time evolution of the population difference between the two condensates. The time is in units of $`|q\chi |`$. Results are shown for the case of $`g=25|q\chi |`$, when the two condensates are in the initial state with $`\alpha _1=5`$ and $`\alpha _2=4`$.
FIG.2. Diagram of the time evolution of the fractional population imbalance. The time is in units of $`|q\chi |`$. Results are shown for (a) $`K=1`$, and (b) $`K=20`$ when the two condensates are in the initial state with $`\alpha _1=10`$ and $`\alpha _2=0`$.
FIG.3. The atomic tunneling current between the two condensates as a function of time $`t`$ (in arbitrary units). Results are shown for the case of $`g=0.25`$ and $`q\chi =0.1`$ when the two condensates are in the initial state with $`\alpha _1=5`$ and $`\alpha _2=4`$.
FIG.4. The atomic tunneling current between the two condensates as a function of time (in arbitrary units) when there does not exist the tunneling coupling. Results are shown for the case of $`q\chi =0.1`$ when the two condensates are in the initial state with $`\alpha _1=10`$ and $`\alpha _2=5`$.
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# Untitled Document
Gherardelli linkage and complete intersections
MSC-class: 14M10 (Primary) 14M06, 1407 (Secondary).
Davide Franco
Dipartimento di Mathematica, Università di Ferrara
via Machiavelli 35, 44100 Ferrara, Italy
E-mail: frv@dns.unife.it
Steven L. Kleiman
Department of Mathematics, Room 2-278 MIT,
77 Mass Ave, Cambridge, MA 02139-4307, USA
E-mail: Kleiman@math.mit.edu and
Alexandru T. Lascu
Dipartimento di Mathematica, Università di Ferrara
via Machiavelli 35, 44100 Ferrara, Italy
E-mail: lsl@dns.unife.it
March 9, 2000
Abstract. Our main theorem characterizes the complete intersections of codimension 2 in a projective space of dimension 3 or more over an algebraically closed field of characteristic 0 as the subcanonical and self-linked subschemes. In order to prove this theorem, we’ll prove the Gherardelli linkage theorem, which asserts that a partial intersection of two hypersurfaces is subcanonical if and only if its residual intersection is, scheme-theoretically, the intersection of the two hypersurfaces with a third.
1. Introduction
Our main result is Theorem (3.2). It characterizes the complete intersections of codimension 2 in $`𝐏^n`$ where $`n3`$, over an algebraically closed field of characteristic 0, among the Cohen–Macaulay $`X`$ as those that are subcanonical and self-linked. This characterization was formulated by Ellia (pvt. comm.), who proved it in a joint work with Beorchia \[3, Thm. 5, p. 556\] assuming $`X`$ is smooth. In Remark 6.1 on p. 557, Beorchia and Ellia said they don’t know whether the smoothness “can be avoided.” It can! Furthermore, $`X`$ can be reducible and nonreduced.
More precisely, an $`X`$ is said to be $`a`$-subcanonical if its dualizing sheaf $`𝜔_X`$ is of the form $`𝜔_X=𝒪_X(a)`$. An $`X`$ is said to be self-linked by two hypersurfaces $`F_1`$ and $`F_2`$ if $`X`$ is equal to its own residual scheme in the complete intersection of $`F_1`$ and $`F_2`$. For example, suppose $`X`$ is the complete intersection of $`F_1`$ and $`F_3`$. Then $`X`$ is self-linked by $`F_1`$ and $`F_2`$ where $`F_2:=2F_3`$ or where $`F_2`$ is, more generally, any hypersurface such that $`F_1F_2=F_12F_3`$. Furthermore, $`X`$ is $`a`$-subcanonical where $`a`$ is the following integer: denote the degree of $`F_i`$ by $`m_i`$; then $`a:=m_1+m_3n1`$. Now, Theorem (3.2) says that this is, in fact, the only example!
This second formulation of Theorem (3.2) is more refined than the first. After all, the first says nothing much about the hypersurfaces $`F_i`$ involved. In particular, the first does not suggest anything like the equation $`m_2=2m_3`$. Indeed, in Corollary 4 on p. 557, Beorchia and Ellia offered an alternative proof of the first formulation in the case where $`X`$ is a curve and $`m_3m_2m_1`$. The proof is correct, but the case is vacuous!
Our proof of Theorem (3.2) follows, to a fair extent, the lines of Beorchia and Ellia’s proof of their Theorem 5. In both proofs, a key step is to split the normal bundle of $`X`$ in $`𝐏^n`$. At this stage, if $`n4`$ and $`X`$ is smooth, then we’re done simply because the normal bundle splits; indeed, Basili and Peskine \[2, p. 87\] proved that then $`X`$ is a complete intersection. However, in order to prove Theorem (3.2) in full generality, we must split the normal bundle with care. For example, consider the twisted cubic space curve $`X`$; its normal bundle is split because $`X`$ is rational, and it is known that $`X`$ is self-linked by a quadric cone and a cubic surface, but $`X`$ is, of course, not a complete intersection, nor even subcanonical.
To split the normal bundle, we’ll use the Gherardelli linkage theorem (2.5). It asserts that, when two hypersurfaces $`F_1`$ and $`F_2`$ of $`𝐏^n`$ intersect partially in an $`X`$, then $`X`$ is subcanonical if and only if its residual scheme $`Y`$ is, scheme-theoretically, of the form $`Y=F_1F_2F_3`$ where $`F_3`$ is a suitable hypersurface. (Such a $`Y`$ is called a quasi-complete intersection.) In particular, if $`X`$ is subcanonical, and is self-linked by $`F_1`$ and $`F_2`$, then $`X=F_1F_2F_3`$. In this case, we’ll form the conormal bundles of $`X`$ in $`F_3`$ and in $`𝐏^n`$, and we’ll split the natural map from the latter bundle onto the former.
We’ll then conclude that some multiple of $`X`$ is numerically equivalent to a hypersurface section of $`F_3`$, at least after we’ve replaced $`F_3`$ by an integral component; we’ll simply apply Braun’s main theorem \[4, p. 403\]. (Braun followed the lines of Ellingsrud, Gruson, Peskine, and Strømme’s remarkable proof of the theorem in the case of a curve on a smooth connected surface. This case had been treated earlier, in a very different fashion, by Griffiths, Harris, and Hulek. See Braun’s paper \[4, p. 411\] for all the references.) Finally, to conclude that $`X`$ is a complete intersection, Beorchia and Ellia used Gruson and Peskine’s work on space curves. Instead, we’ll make a direct geometric argument, and obtain our more refined statement of Theorem (3.2).
If $`n6`$ and $`X`$ is smooth, then, since $`X`$ is a quasi-complete intersection, it is, in fact, a complete intersection by Faltings’ Korollar of Satz 3 \[7, p. 398\]. This line of proof is significant because it is valid in any characteristic, whereas Basili and Peskine work in characteristic 0, and we must too although only to apply Braun’s theorem. Beorchia and Ellia \[3, p. 556\] suggested that there might be a problem in characteristic 2 by pointing out the following result, due in part to Rao \[15, p.272\] and in part to Migliore \[12, p. 185\]: a double line in $`𝐏^3`$ of arithmetic genus $`2`$ or less is self-linked if and only if the characteristic is 2. We’ll pursue this suggestion in Example (3.4). On the other hand, it would be nice to know whether Theorem (3.2) is valid except for certain $`X`$ of small dimension in characteristic 2.
The Gherardelli linkage theorem holds in greater generality than that stated above. In Theorem (2.5), we’ll replace $`𝐏^n`$ by any Gorenstein projective scheme $`P`$ having pure dimension 2 or more and satisfying this vanishing condition: $`H^q(𝒪_P(m))=0`$ for three specific values of the pair $`(q,m)`$. For example, $`P`$ can be a complete intersection in $`𝐏^n`$. Thus we’ll recover Theorem 2(i) of Fiorentini and Lascu \[10, p. 170\], where, in addition, $`X`$ and $`Y`$ are assumed to have no common components; in fact, our proof was inspired by theirs. Beorchia and Ellia \[3, p. 556\] proved the existence of the hypersurface $`F_3`$ directly in the case at hand by using the mapping cone. Earlier, Rao \[14, pp. 209–10\] proved (burying it among other things) a version of the Gherardelli linkage theorem, in which the condition that $`X`$ be subcanonical is replaced by the condition that $`X`$ be the zero scheme of a section of a rank-2 vector bundle on $`𝐏^n`$; these two conditions are equivalent by a famous theorem of Serre’s (see \[8, Prop. 3, p. 346\]). On the other hand, our Theorem (3.2) does not hold even if $`𝐏^n`$ is replaced by a smooth hypersurface $`P`$, as we’ll see in Example (3.3).
To prove the Gherardelli linkage theorem (2.5), we’ll use the Noether linkage sequence (2.3.1), which presents the dualizing sheaf of a partial intersection in any Gorenstein ambient scheme $`P`$ having pure dimension 2 or more. The case where $`P`$ is a complete intersection in $`𝐏^n`$ was treated in \[10, Lem. 1\] and in \[13, 1.6\] and was used in \[15, p. 253\]. The general case is, as we’ll see, no more difficult to prove.
In short, in Section 2, we’ll review some basic linkage theory, including the Peskine–Szpiro linkage theorem (compare with \[13, 1.3, p. 274\] and \[5, 21.23, p. 541\]), the Noether linkage sequence, and the Gherardelli linkage theorem. This theory is all more or less well known, but has not always been developed exactly as here, and it is all essential for our work in Section 3. In Section 3, we’ll prove our main theorem, our characterization of complete intersections of codimension 2 in $`𝐏^n`$. Finally, we’ll discuss two examples; the first shows that the ambient projective space cannot be replaced even by a smooth hypersurface, and the second shows that our characterization fails in characteristic 2.
2. Gherardelli linkage
Proposition (2.1) (Peskine–Szpiro linkage theorem) Let $`Z`$ be a Gorenstein scheme, $`XZ`$ be a proper closed subscheme, and $`Y`$ the residual scheme of $`X`$. If $`X`$ is Cohen–Macaulay of pure codimension $`0`$, then so is $`Y`$; furthermore, then $`X`$ is also the residual scheme of $`Y`$.
Proof. Let $`_{X/Z}`$ and $`_{Y/Z}`$ denote the ideals. Then we have
$$_{Y/Z}:=𝐴𝑛𝑛_{𝒪_Z}_{X/Z}\stackrel{}{}𝐻𝑜𝑚_{𝒪_Z}(𝒪_X,𝒪_Z),$$
$`(\mathrm{2.1.1})`$
where the equation holds by definition and the isomorphism is given by evaluation at 1.
It is a basic fact (see \[5, 21.21, p. 538\]) that, on the category of maximal (dimensional) Cohen–Macaulay $`𝒪_Z`$-modules $``$, the functor,
$$𝒟():=𝐻𝑜𝑚_{𝒪_Z}(,𝒪_Z),$$
is dualizing. Now, $`𝒟`$ interchanges the two basic exact sequences,
$$0_{X/Z}𝒪_Z𝒪_X0\text{ and }0_{Y/Z}𝒪_Z𝒪_Y0;$$
indeed, $`𝒟`$ carries the first sequence to the second thanks to (2.1.1), and so, as $`𝒟`$ is dualizing, $`𝒟`$ carries the second sequence back to the first. Thus, $`𝒪_Y=𝒟(_{X/Z})`$ and $`_{X/Z}=𝒟(𝒪_Y)`$. The latter equation implies that $`X`$ is the residual scheme of $`Y`$. The former equation implies that $`𝒪_Y`$ is maximal Cohen–Macaulay, because $`_{X/Z}`$ is so since, at any $`xX`$,
$$depth_{X/Z,x}\mathrm{min}(depth𝒪_{Z,x},\mathrm{\hspace{0.17em}1}+depth𝒪_{X,x})$$
(see \[5, 18.6b, p. 451\]). The proof is now complete.
Setup (2.2) Let $`P`$ be a complete scheme defined over an algebraically closed field of arbitrary characteristic. Assume that $`P`$ is Gorenstein of pure dimension at least 2, and equip $`P`$ with an invertible sheaf $`𝒪_P(1)`$, which is not necessarily ample. For $`i=1,\mathrm{\hspace{0.17em}2}`$ let $`f_iH^0(𝒪_P(m_i))`$ be a section, and $`F_i:f_i=0`$ its scheme of zeros. Set
$$Z:=F_1F_2,$$
and assume that $`Z`$ has pure codimension 2.
Let $`XZ`$ be a proper closed subscheme, and assume that $`X`$ is Cohen–Macaulay of pure codimension 2 in $`P`$. Let $`YZ`$ be the residual scheme of $`X`$. By the Peskine–Szpiro linkage theorem (2.1), also $`Y`$ is Cohen–Macaulay of pure codimension 2 in $`P`$, and $`X`$ is also the residual scheme of $`Y`$.
Proposition (2.3) (Noether linkage sequence) In the setup of (2.2), the dualizing sheaves and the ideals in $`P`$ are related by the following short exact sequence:
$$0_{Z/P}𝜔_P(m_1+m_2)_{Y/P}𝜔_P(m_1+m_2)𝜔_X0.$$
$`(\mathrm{2.3.1})`$
Proof. First, note the following two equations:
$$𝜔_Z=𝜔_P(m_1+m_2)|_Z\text{ and }𝜔_X=_{Y/Z}𝜔_Z.$$
$`(\mathrm{2.3.2})`$
The first equation is standard, and results from basic duality theory (see \[1, Ch. 1\] for example):
$$𝜔_Z=𝐸𝑥𝑡_P^2(𝒪_Z,𝜔_P)=𝐻𝑜𝑚_Z(det(_{Z/P}/_{Z/P}^2),𝜔_P|Z).$$
The second equation in (2.3.2) results from a series of three other equations:
$$𝜔_X=𝐻𝑜𝑚(𝒪_X,𝜔_Z)=𝐻𝑜𝑚(𝒪_X,𝒪_Z)𝜔_Z=_{Y/Z}𝜔_Z.$$
These hold by elementary duality theory, by the invertiblity of $`𝜔_Z`$, and by (2.1.1) above.
Finally, the Noether linkage sequence (2.3.1) results from the basic sequence,
$$0_{Z/P}_{Y/P}_{Y/Z}0,$$
by tensoring it with $`𝜔_P(m_1+m_2)`$ and then using the two equations in (2.3.2).
Remark (2.4) According to Enriques \[6, vol. 3, p. 534\], Noether obtained the preceding proposition in the special case where $`P`$ is the projective 3-space, Noether having stated it virtually as follows:
If the curve $`X`$ is the partial intersection of two surfaces $`F_1`$ and $`F_2`$ of degrees $`m_1`$ and $`m_2`$, meeting further in a curve $`Y`$, then the surfaces of degree $`m_1+m_24`$ passing through $`Y`$ cut on $`X`$ the complete canonical series.
To derive this statement, take (2.3.1), replace $`𝜔_P`$ by $`𝒪_P(4)`$, and extract cohomology, obtaining the following exact sequence:
$$H^0(_{Y/P}(m_1+m_24))H^0(𝜔_X)H^1(_{Z/P}(m_1+m_24)).$$
The third term vanishes as $`Z`$ is a complete intersection, and Noether’s statement follows.
Theorem (2.5) (Gherardelli linkage) Preserve the setup of (2.2). Let $`m_3>0`$. If an $`f_3H^0(𝒪_P(m_3))`$ exists such that $`Y=F_1F_2F_3`$ where $`F_3:f_3=0`$, then
$$𝜔_X=𝜔_P(m_1+m_2m_3)|_X.$$
The converse holds if, in addition,
$$H^1(𝒪_P(m_3m_1))=0,H^1(𝒪_P(m_3m_2))=0,\text{ and }H^2(𝒪_P(m_3m_1m_2))=0.$$
Proof. Assume an $`f_3`$ exists. Then $`Y=ZF_3`$. Hence, multiplication by $`f_3`$ gives a surjection $`\mu :𝒪_Z(m_3)_{Y/Z}`$. Its kernel $`𝐴𝑛𝑛_{Y/Z}(m_3)`$ is equal to $`_{X/Z}(m_3)`$ because $`X`$ is also the residual scheme of $`Y`$ thanks to (2.1). So $`\mu `$ induces an isomorphism $`𝒪_X(m_3)\stackrel{}{}_{Y/Z}`$. Hence $`𝜔_X`$ has the asserted form thanks to (2.3.2).
Conversely, assume $`𝜔_X=𝜔_P(m_1+m_2m_3)|_X`$. Then twisting the Noether linkage sequence (2.3.1) yields the following exact sequence:
$$0_{Z/P}(m_3)_{Y/P}(m_3)𝒪_X0.$$
$`(\mathrm{2.4.1})`$
Extracting cohomology yields the next exact sequence:
$$H^0(_{Y/P}(m_3))H^0(𝒪_X)H^1(_{Z/P}(m_3)).$$
Assume the additional vanishing conditions. Then $`H^1(_{Z/P}(m_3))=0`$ thanks to the twisted Koszul resolution,
$$0𝒪_P(m_3m_1m_2)𝒪_P(m_3m_1)𝒪_P(m_3m_2)_{Z/P}(m_3)0.$$
Hence, we may lift $`1H^0(𝒪_X)`$ to an $`f_3H^0(_{Y/P}(m_3))`$. Set $`F_3:f_3=0`$.
In (2.4.1), we may replace $`𝒪_X`$ by $`_{Y/Z}(m_3)`$. Hence $`_{Y/Z}(m_3)`$ is generated by the image of $`f_3`$ in $`H^0(_{Y/Z}(m_3))`$. Therefore, $`Y=ZF_3`$, and the proof is complete.
3. Complete intersections
Definition (3.1) Let $`P`$ be a Gorenstein scheme, $`X`$ a closed Cohen–Macaulay subscheme. We’ll say that $`X`$ is subcanonical in $`P`$ if $`P`$ is equipped with an invertible sheaf $`𝒪_X(1)`$, and if, for some integer $`\alpha `$, we have
$$𝜔_X=𝜔_P(\alpha )|_X.$$
Assume $`P`$ has pure dimension at least 3, and $`X`$ has pure codimension 2. We’ll say that $`X`$ is self-linked in $`P`$ by two effective Cartier divisors $`F_1`$ and $`F_2`$ if they meet properly in a subscheme $`Z`$ containing $`X`$, and if $`X`$ is equal to the residual scheme $`Y`$ of $`X`$ in $`Z`$.
Theorem (3.2)Let $`P`$ be a projective space of dimension $`n3`$ over an algebraically closed field of characteristic $`0`$. Let $`XP`$ be a closed subscheme that is Cohen–Macaulay of pure codimension $`2`$. Assume $`X`$ subcanonical and self-linked. Then $`X`$ is a complete intersection.
In fact, say $`X`$ is self-linked by hypersurfaces $`F_1`$ and $`F_2`$ of degrees $`m_1`$ and $`m_2`$. Then, after $`F_1`$ and $`F_2`$ are switched if need be, $`m_2`$ is even, and there is a hypersurface $`F_3`$ of degree $`m_2/2`$ such that $`X=F_1F_3\text{ and }Z=F_12F_3\text{ where }Z:=F_1F_2`$.
Proof. Since $`P`$ is smooth and $`X`$ is subcanonical, $`X`$ is Gorenstein. Hence, since $`X`$ has pure codimension $`2`$, it is locally a complete intersection in $`P`$ by one of Serre’s results \[5, 21.10, p. 537\]. Hence, on $`X`$, the conormal sheaf $`_{X/P}/_{X/P}^2`$ is locally free of rank 2.
By another celebrated theorem of Serre’s, $`H^i(𝒪_P(j))=0`$ for $`i=1,2`$ and for any $`j`$ since $`n3`$. Hence, by the Gherardelli linkage theorem (2.5), there is a hypersurface $`F_3`$ such that $`X=ZF_3`$.
Let $`xX`$. For $`i=1,2,3`$ let $`\phi _i𝒪_{P,x}`$ generate the ideal of $`F_i`$. Then $`_{X/P,x}`$ is generated by $`\phi _1`$, $`\phi _2`$, and $`\phi _3`$, but not by $`\phi _1`$ and $`\phi _2`$, since $`X=ZF_3`$, but $`XZ`$. Since $`_{X/P,x}`$ is generated by two elements, it must be generated either by $`\phi _1`$ and $`\phi _3`$ or by $`\phi _2`$ and $`\phi _3`$. Hence $`X`$ is a Cartier divisor on $`F_3`$.
For $`i=1,2`$ set $`Z_i:=F_iF_3`$. Let $`xX`$. Then, by the preceding paragraph, $`_{X/P,x}`$ is equal either to $`_{Z_1/P,x}`$ or to $`_{Z_2/P,x}`$. Put geometrically, $`X`$ is equal, in a neighborhood of $`x`$ in $`P`$, either to $`Z_1`$ or to $`Z_2`$.
For $`i=1,2,3`$ say $`F_i:f_i=0`$. For $`i=1,2`$ form the greatest common divisor $`g_i`$ of $`f_i`$ and $`f_3`$, and set $`G_i:g_i=0`$.
First, suppose both $`G_1`$ and $`G_2`$ are nonempty, and let $`x`$ be a common point. Since $`G_1`$ is a component of both $`F_1`$ and $`F_3`$, their intersection $`Z_1`$ is not equal to $`X`$ in a neighborhood of $`x`$. Similarly, $`Z_2`$ is not equal to $`X`$ in a neighborhood of $`x`$. This conclusion stands in contradiction to our conclusion above, that $`X`$ is equal, in a neighborhood of $`x`$ in $`P`$, either to $`Z_1`$ or to $`Z_2`$. Therefore, not both $`G_1`$ and $`G_2`$ are nonempty; say $`G_2`$ is empty.
Then $`Z_2`$ has pure codimension 2 in $`P`$, and $`Z_2X`$. If $`Z_2=X`$, then $`X=F_2F_3`$. So suppose not, and we’ll prove that $`X=F_1F_3`$. Form the residual scheme $`X_2`$ of $`X`$ in $`Z_2`$. By general principles, $`X_2`$ is a Cartier divisor on $`F_3`$ because $`X`$ and $`Z_2`$ are so; moreover, $`Z_2=X+X_2`$.
Suppose $`G_1`$ is nonempty. Set $`C:=G_1F_2`$. Then $`C`$ is a hypersurface section of $`F_2`$. So $`C`$ has a point $`x`$ in common with $`X_2`$, which also lies on $`F_2`$. Then $`xX`$, because $`CZ`$ and $`Z`$ has the same support as $`X`$. Since $`G_1`$ is a component of both $`F_1`$ and $`F_3`$, their intersection $`Z_1`$ is not equal to $`X`$ in a neighborhood of $`x`$. Since $`x`$ lies on both $`X_2`$ and $`X`$, also $`Z_2`$ is not equal to $`X`$ in a neighborhood of $`x`$. As before, there is a contradiction. Therefore, $`G_1`$ is empty.
Then $`Z_1`$ has pure codimension 2 in $`P`$, and $`Z_1X`$. If $`Z_1=X`$, then $`X=F_1F_3`$ as claimed. So suppose not, and form the residual scheme $`X_1`$ of $`X`$ in $`Z_1`$. By general principles, $`X_1`$ too is a Cartier divisor on $`F_3`$. After a bit of work, we’ll achieve a contradiction.
First, we’ll construct a natural splitting of the natural surjection,
$$_{X/P}/_{X/P}^2_{X/F_3}/_{X/F_3}^2.$$
$`(\mathrm{3.2.1})`$
To do so, in $`_{X/P}/_{X/P}^2`$ , form the image $``$ of $`_{Z/P}`$; we are going to show that $``$ maps isomorphically onto $`_{X/F_3}/_{X/F_3}^2`$. Since $``$ maps surjectively, and since $`_{X/F_3}/_{X/F_3}^2`$ is invertible as $`X`$ is a Cartier divisor on $`F_3`$, we need only show that $``$ is invertible.
Let $`xX`$. Say, as above, that $`_{X/P,x}=_{Z_1/P,x}`$. Set $`W:=F_12F_3`$. Then $`WZ`$; indeed, $`_{F_3/P}^2_{Z/P}`$ because $`_{X/Z}=𝐴𝑛𝑛_{X/Z}`$ since $`X`$ is self-linked. Since also $`WZ_1`$, there is a natural commutative diagram,
$`0_{Z_1/W}𝒪_W𝒪_{Z_1}0`$ $`.`$$`u`$$`v`$$`w`$ $`0_{X/Z}𝒪_Z𝒪_X0`$
Clearly, $`_{Z_1/W}=𝒪_P(F_3)|_{Z_1}`$. Moreover, $`_{X/Z}=𝜔_X𝜔_P(m_1+m_2)^1`$ thanks to (2.3.2) with $`Y:=X`$. Thus the source of $`u`$ is invertible on $`Z_1`$, and the target is invertible on $`X`$. Now, $`_{X/P,x}=_{Z_1/P,x}`$. Hence, $`w`$ is an isomorphism at $`x`$; in other words, $`X`$ and $`Z`$ are the same scheme in a neighborhood of $`x`$. Also, $`u`$ is surjective at $`x`$, and its source and target are invertible sheaves on the same scheme in a neighborhood of $`x`$; hence, $`u`$ is an isomorphism at $`x`$. So $`v`$ is an isomorphism at $`x`$. Hence, $`_{W/P,x}=_{Z/P,x}`$.
So, in $`_{X/P}/_{X/P}^2`$, the images of $`_{W/P}`$ and $`_{Z/P}`$ are equal at $`x`$. The image of $`_{W/P}`$ is equal to $`𝒪_P(F_1)|_X`$ at $`x`$; indeed, the latter sheaf maps naturally into the former, and this map is surjective since $`XF_3`$, and injective at $`x`$ since its natural image is a direct summand of $`_{X/P}/_{X/P}^2`$ at $`x`$, as $`_{X/P,x}=_{Z_1/P,x}`$. The image of $`_{Z/P}`$ is $``$ by definition. Thus $``$ is invertible at $`x`$. Since $`xX`$ is arbitrary, $``$ is invertible. Thus $`\stackrel{}{}_{X/F_3}/_{X/F_3}^2`$, and (3.2.1) splits.
Let $`F`$ be any irreducible component of $`F_3`$, and equip $`F`$ with its reduced structure. Since $`F`$ is a hypersurface, $`F`$ meets $`X`$. Set $`V:=XF`$. Then $`V`$ is a Cartier divisor on $`F`$, and hence $`V`$ is locally a complete intersection in $`P`$. Consider the natural commutative diagram of sheaves on $`V`$,
$`|V\left(_{X/F_3}/_{X/F_3}^2\right)|V`$ $``$$``$ $`_{V/P}/_{V/P}^2_{V/F}/_{V/F}^2`$
The top horizontal map is an isomorphism because it is the restriction of an isomorphism. The right vertical map is an isomorphism because it is surjective and its source and target are invertible. Therefore, the lower horizontal map splits.
Since the lower map splits, since $`V`$ is a Cartier divisor on $`F`$ and is locally a complete intersection in $`P`$, since $`F`$ is reduced, irreducible, and closed, and since $`P`$ is a projective space of dimension $`n3`$ over an algebraically closed field of characteristic $`0`$, Braun’s main theorem \[4, p. 26\] implies that some multiple of $`V`$ is numerically equivalent to a hypersurface section of $`F`$.
Since $`F`$ is a hypersurface, $`F`$ meets both $`X_1`$ and $`X_2`$, which are supposedly nonempty. For $`i=1,\mathrm{\hspace{0.17em}2}`$ set $`V_i:=X_iF`$. Then $`V_i`$ is a Cartier divisor on $`F`$, and $`V+V_i=F_iF`$. Hence, some multiple of $`V_i`$ too is numerically equivalent to a hypersurface section of $`F`$. Therefore, $`V_1`$ and $`V_2`$ have a common point $`x`$. Then $`x`$ lies on both $`Z_1`$ and $`Z_2`$, so on their intersection, which is $`X`$. However, there is no neighborhood of $`x`$ in which either $`Z_1`$ or $`Z_2`$ is equal to $`X`$ because $`x`$ lies on both $`X_1`$ and $`X_2`$. Thus, we’ve achieved the desired contradiction. Therefore, $`X=F_1F_3`$.
Then $`W=Z`$ everywhere, by the reasoning above; in other words, $`Z=F_12F_3`$. Finally, set $`m_3:=\mathrm{deg}F_3`$. Then $`\mathrm{deg}Z=2m_1m_3`$. Now, $`Z:=F_1F_2`$, so $`\mathrm{deg}Z=m_1m_2`$. Hence $`2m_3=m_2`$. The proof is now complete.
Example (3.3) Most of the proof of Theorem (3.2) works without change in the relative case where $`P`$ is a smooth projectively Cohen–Macaulay variety of pure dimension at least 3. However, to apply Braun’s theorem, we must know that the surjection (3.2.1) splits when $`P`$ is replaced by the ambient projective space; the proof shows that (3.2.1) itself splits, but this splitting is insufficient. The theorem does not hold even when $`P`$ is replaced by a smooth hypersurface, as the following example shows.
Let $`P`$ be a smooth quadric hypersurface in $`𝐏^4`$. Let $`F_1`$ be the section of $`P`$ by a hyperplane $`H_1`$ that is tangent to $`P`$ at a point $`x`$. Then $`F_1`$ is a cone in $`H_1`$ with vertex at $`x`$ and with base a smooth (plane) conic $`C`$. Fix $`yC`$. Then $`y`$ determines a generator $`X`$ of the cone $`F_1`$. Let $`H_2`$ be a hyperplane in $`𝐏^4`$ that cuts $`H_1`$ in the plane spanned by $`x`$ and by the tangent line to $`C`$ at $`y`$. Then $`X`$ is a line, so subcanonical in $`P`$. Moreover, $`X`$ is self-linked in $`P`$ by $`F_1`$ and $`F_2`$ with $`F_2:=H_2P`$. However, $`X`$ is not the complete intersection of two hypersurface sections of $`P`$ since any such complete intersection has even degree in $`𝐏^4`$.
Example (3.4) Theorem (3.2) is not valid in positive characteristic without some further restriction on $`X`$. Indeed, we are going to see that, in characteristic 2, there exists an example of an irreducible, but nonreduced, Cohen–Macaulay space curve $`X`$, which is subcanonical and self-linked, yet is not a complete intersection.
Ferrand \[8, p. 345\] explained how to put a subcanonical double structure on a line (indeed, on any complete curve that is locally a complete intersection) in $`𝐏^3`$ in any characteristic; moreover, the double curve can have arbitrarily negative arithmetic genus. Now, Migliore \[12, p. 185\] proved that, in characteristic 2, a double line $`X`$ is self-linked if its arithmetic genus is $`2`$ or less. Such an $`X`$ is not a complete intersection, because every complete intersection $`Z`$ has nonnegative arithmetic genus by (2.3.2).
References
AK A. Altman and S. Kleiman, “Introduction to Grothendieck duality theory,” Lecture Notes in Math. 146, Springer-Verlag, 1970.
BP B. Basili and C. Peskine, Décomposition du fibré normal des surfaces lisses de $`𝐏_4`$ et structures doubles sur les solides de $`𝐏_5`$, Duke Math. J. 69 (1993), 87–95.
BE V. Beorchia and Ph. Ellia, Normal bundle and complete intersections, Rend. Sem. Univ. Pol. Torino 48 (1990), 553–62.
B R. Braun, On the normal bundle of Cartier divisors on projective varieties, Arch. Math. 59 (1992), 403–11.
E D. Eisenbud, “Commutative Algebra,” GTM 150, Springer-Verlag 1994.
EC F. Enriques and O. Chisini, “Lezioni sulla teoria geometrica delle equazioni e delle funzioni algebriche,” Zanichelli, Bologna, 1915.
Fa G. Faltings, Ein Kriterium for vollständige Durchshnitte, Inventiones math. 62 (1980), 393–401.
F D. Ferrand, Courbes gauches et fibrés de rang $`2`$, C. R. Acad. Sc. Paris 281 (1975), 345–7.
FL81 M. Fiorentini and A. T. Lascu, Una formula di geometria numerativa, Ann. Univ. Ferrara 27 (1981), 201–27.
FL87 M. Fiorentini and A. T. Lascu, Projective embeddings and linkage, Rend. Sem. Mat. Fis. Milano LXVII (1987), 161–82.
Gh G. Gherardelli, Sulle curve sghembe algebriche intersezioni complete di tre superficie, Rend. Acc. d’Italia IV (1943), 460–62.
M J. Migliore, On linking double lines, Trans. Amer. Math. Soc. 294 (1986), 177–85.
PS C. Peskine and L. Szpiro, Liaison des variétés algébriques. I, Inventiones math. 26 (1974), 271–302.
R79 P. Rao, Liaison among curves in $`𝐏^3`$, Inventiones math. 50 (1979), 205–17.
R82 P. Rao, On self-linked curves, Duke Math. J. 49 (1982), 251–73.
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# L1551NE - Discovery of a Binary Companion
## 1 Introduction
Young T Tauri stars in Taurus have been found to have a high incidence of multiplicity, with the fraction of close ($``$100A.U.) companions found to be 0.40 $`\pm `$ 0.08 (Ghez et al. 1993; Leinert et al. 1993; Simon et al. 1995; Patience et al. 1998), using speckle and occultation techniques. Similar surveys of the Hyades cluster found a smaller fraction of 0.30 $`\pm `$ 0.06, and a still lower fraction of 0.14 $`\pm `$ 0.03 for G dwarfs in the solar neighborhood (Patience et al. 1998), suggesting an evolutionary effect and/or environmental effects during formation. Deeply embedded class I and class 0 protostars cannot be surveyed using similar techniques, since they are not visible at optical wavelengths and are generally too deeply embedded even at infrared wavelengths. Millimeter/sub-millimeter-wavelength interferometry at sub-arcsecond resolution can resolve close companions if they have circumstellar disks. For example, L1551 IRS5 was recently shown to be a binary based on 7mm interferometric observations at the VLA (Rodríguez et al. 1998). We present here evidence that L1551 NE also has a binary companion.
L1551NE (B1950 $`4^h28^m50.5^s`$ +$`18^{}02^{}10\mathrm{"}`$ (Draper, Warren-Smith & Scarrott 1985)) is a young stellar object in the L1551 molecular cloud, at a distance of 160 pc (Snell 1981). Discovered by Emerson et al. (1984) from IRAS data, it is the second brightest embedded source in the Taurus complex after L1551 IRS5, with L$`{}_{bol}{}^{}6L_{}`$. It has a molecular outflow (Moriarty-Schieven, Butner & Wannier 1995). The radial density distribution of L1551NE has been modeled by Barsony & Chandler (1993) from 800$`\mu m`$ images, and by Butner et al. (1995) from 100$`\mu m`$ and 200$`\mu m`$ observations. Both found that the density distribution implied by the radial intensity profile was much shallower than the $`n(r)r^{1.5}`$ predicted by the “inside-out” collapse model of Terebey, Shu & Cassen (1984). Moriarty-Schieven, Butner & Wannier (1995) have suggested that L1551 NE may be a class 0 source.
## 2 Observations
The observations were taken with the Owens Valley Radio Observatory millimeter-wave array 1994 November 12 and 1997 January 17, at a wavelength of 1.3mm (230.65 GHz). The 1994 observations were taken with the array (consisting at the time of six 10m antennae) in the compact “$`l`$” configuration (maximum baseline 115m, minimum baseline 30m), and in 1997 the data were taken using the extended “$`h`$” configuration (maximum baseline 241.7m, minimum baseline 35m). The effective resolution of the synthesized beam was $``$2” and $``$1” for the two configurations.
The observations were made in “snapshot” mode, with a series of ten minute integrations interspersed with observations of other sources and the phase calibrators (0528+134 and 3C84), to yield a total integration time of $``$1 hour in the “$`l`$” configuration and $``$2 hours in the “$`h`$” configuration. The system temperatures were $``$1200K and $``$1600K respectively. Neptune was used as the primary flux calibrator. The continuum bandwidth is 1GHz.
The data were reduced using the standard reduction package “mma”, and then exported to AIPS which was used to generate and clean images (using the ‘IMAGR’ task and uniform weighting). The clean box included only a small area covering the main source and the apparent binary companion. The angular resolution of the cleaned image was 1.29” $`\times `$ 1.07” at position angle -74.
## 3 Data
Figure 1 presents a contour image of L1551NE. Clearly seen are the primary source (henceforth source A) at (B1950) 04<sup>h</sup>28<sup>m</sup>50.552<sup>s</sup> 1802’09.85”, and another weaker source 1.43” to the southeast (source B). Emission also appears to extend to the north and to the east of the primary source, and surrounding source B.
To verify that we are seeing real sources and not phase instabilities, we generated images using the same techniques of another source, IRAS 04169+2702, which was observed during the same two days and interspersed with L1551NE. Phase errors would be manifested as “anomalous” sources or structures. No such anomalous sources are seen in the image of IRAS 04169+2702. In addition, we cleaned the “$`l`$” configuration data separately from the “$`h`$” configuration data. The high-resolution data clearly had two peaks, while the low-resolution data showed an extended disk-like structure with long axis through the line joining the two sources. Thus we believe that the structure seen here in L1551NE is real.
In Table 1 we present the positions, sizes and flux densities of the sources. The single-dish flux density at 1.35mm is also shown.
## 4 Discussion
There are three distinct components apparent in the image shown in Figure 1; a brighter, possibly extended source at the field center (source A), a weaker, probably point-like source (B) 1.43” south-east of A, and diffuse, low-level emission which appears to surround both sources and extend $``$2” to the northwest and east of A.
Sources A and B were fit with elliptical gaussians using the AIPS task JMFIT. The primary source A was found to have a size 1.53”x1.28” (i.e. it may have been slightly resolved with a size $`2\sigma `$ larger than the beamsize), peak intensity $``$ 0.33 Jy and integrated intensity $``$ 0.47 Jy (i.e. $``$40% ($`3\sigma `$) larger than the peak intensity). If it is slightly resolved (deconvolved size $``$ 0.82”x0.70”), then its size is $``$131x112 A.U. at a distance of 160pc. Source B is located 1.43” (229 A.U. at 160pc) south-east of A, and its size and intensity are consistent with it being unresolved, i.e. $`<`$100 A.U. We estimate the mass within each source using $`M_D=\frac{S_\nu D^2}{B_\nu (T)}\frac{4}{3}\frac{a\rho }{Q_\nu }`$, and display these in Table 1. We used the integrated flux densities, assumed a distance of 160 pc (Snell 1981), used dust temperature of 42K (Moriarty-Schieven et al. 1994), and dust emissivities from Hildebrand (1983) $`\frac{4}{3}\frac{a\rho }{Q_\nu }=0.1(\lambda /250)^\beta `$ g cm<sup>-2</sup>, assuming $`\beta `$ = 1 (Moriarty-Schieven et al. 1994), and have assumed a gas to dust mass ratio of 100. The derived masses are approximately 0.044 and 0.014 M for the emission from source A and B respectively, and 0.022M for the envelope. The total mass for the two sources plus envelope is $``$0.078 M.
Rodríguez, Anglada & Raga (1995) imaged the L1551NE region at $`\lambda `$3.5c m and found a continuum source located within 1.5$`\sigma `$ of our source A. Rodríguez et al attribute the $`\lambda `$3.5cm emission to shocks associated with the outflow (Moriarty-Schieven, Butner & Wannier 1995). Thus source A is the likely origin of the L1551NE outflow. Rodrǵuez et al. found a possible second source located $``$0.6” east of A, which is not coincident with our source B, for which no $`\lambda `$3.5cm emission was detected. If their source 0.6” east of A is real, this would suggest that L1551NE is at least a triple star system. However, another protostar capable of generating an outflow jet should have had a circumstellar disk large enough to be detected by our observations. Possibly the eastern $`\lambda `$3.5cm source is due to a jet from A but offset from the source, or is a background object.
The single-dish flux density (Butner et al. 2000) found using the 14m JCMT telescope (FWHM 20”)) is not significantly different from the total intensity found in our image of L1551NE (primary beam FWHM 28”). Thus only a small amount, if any, of the single-dish flux can have been resolved out by the interferometer. However, considerable extended emission was seen at $``$850 $`\mu `$m by Barsony & Chandler (1993) and Moriarty-Schieven et al. (1999), and at 200 $`\mu `$m by Butner et al. (1995). Barsony & Chandler and Butner et al. modeled this extended emission as an envelope whose radial density distribution decreases very slowly with distance from the protostar.
Weak, low-level emission can be seen extending $``$1-2” to the north and east of source A, and perhaps encompassing source B. This extended structure has a disk-like appearance, of dimension $`5^{\prime \prime }\times 2^{\prime \prime }`$ ($``$800x300pc with long axis at position angle $``$-12. This is roughly perpendicular to the axis of the conical reflection nebula emanating from L1551NE (Draper et al. 1985; Hodapp 1995), and hence of the molecular outflow (Moriarty-Schieven et al. 1995). This disk-like structure may represent a circum-binary disk.
###### Acknowledgements.
JAP was supported by a Hawai’i Space Grant College Fellowship which is funded by the NASA Undergraduate Space Grant Fellowship program. The Owens Valley millimeter-wave array is supported by NSF grant AST-96-13717.
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# Detection of new sources of methanol emission at 95 GHz with the Mopra telescope
## 1 Introduction
Methanol, OH and H<sub>2</sub>O masers are all frequently associated with massive star formation regions, however, methanol masers offer more possibilities for the study of star forming regions than either OH or H<sub>2</sub>O, because there are numerous transitions in the microwave region of spectrum. One of the most widespread methanol masers is the $`7_06_1A^+`$ transition at 44 GHz. About 50 masers from this transition have been detected in the northern hemisphere (Morimoto, Ohishi & Kanzawa 1985, Haschick, Menten & Baan 1990, Bachiller et al. 1990, Kalenskii et al. 1992, Kalenskii et al. 1994) and a similar number have been found in the southern hemisphere (Slysh et al. 1994). According to the empirical scheme of Menten (1991), the 44 GHz $`7_06_1A^+`$ transition is a class I methanol maser. Class I methanol masers differ from class II methanol masers in that they are not directly associated with compact Hii regions and OH masers, and the two classes also emit in different transitions.
The $`8_07_1A^+`$ transition at 95 GHz is an analogue of the 44-GHz masing transition (see Fig. 1 of Val’tts et al. (1999)). The upper ($`8_0`$) energy level of the 95-GHz transition is 18.5K higher than the upper ($`7_0`$) energy level of the 44-GHz transition and modelling (Cragg et al. 1992) shows that strong maser emission is also expected from this transition. Val’tts et al. (1995) carried out a search for 95-GHz methanol masers with the Onsala radiotelescope which detected a large number, mostly at the position of 44-GHz masers and with spectra similar to the spectra of the 44-GHz masers. The intensity of 95-GHz masers was found to correlate with intensity of 44-GHz masers, and was on average about 0.5 of the intensity 44-GHz masers. In this paper we present the results of the first search for 95-GHz methanol masers in the southern hemisphere, carried out with the ATNF Mopra telescope in Australia. This search completes a whole sky survey of class I methanol masers source in two transitions, at 44 and 95 GHz.
## 2 Observations
The observations were carried out in the period from July 1 to 17, 1997, using the Mopra 22-m millimeter-wave telescope of the ATNF. The assumed rest frequency of the $`8_07_1A^+`$ transition of methanol was 95.169489 GHz (De Lucia et al. 1989). At this frequency only the inner 15 metres of the Mopra antenna is illuminated and the aperture efficiency is 41%, which implies that one Kelvin of antenna temperature corresponds to 40 Jy. The half-power beamwidth of the Mopra antenna at 95 GHz is 52<sup>′′</sup>. The observations were performed in a position switching mode with reference positions offset 30 in declination. The antenna pointing was checked and corrected every 12 hours by making observations of 86-GHz SiO masers, the nominal pointing accuracy when this procedure is followed is 10<sup>′′</sup> rms. For the majority of sources observations were only made at the nominal position, as time limitations did not permit the observation of grids at offset positions (to accurately determine the position) to be made. A grid of observations were observed toward a number of interesting sources, and for some cases where the 95-GHz emission was anomalously low. A number of these grids found the strongest emission to be offset from the nominal position and in some cases this will be due to antenna pointing errors. However, many of the sources have only been observed in lower frequency class I transitions and the offset may be due to limited accuracy to which the source position has been previously determined.
A cryogenically cooled low-noise SIS mixer was used in the receiver. The single side-band receiver noise temperature was 110 K and the system temperature varied between 220 K and 320 K depending on weather conditions and the elevation of the telescope. An ambient temperature load (assumed to have a temperature of 290 K) was regularly placed in front of the receiver to enable calibration using the method of Kutner & Ulich (1981). This corrects the observed flux density for the effects of atmospheric absorption, ohmic losses and rear-ward spillover. Variations in the ambient temperature of a few percent occurred during the observations, and the estimated uncertainty of the absolute flux density scale is 10%.
For the majority of the observations the back-end was a 64 MHz wide 1024-channel autocorrelator with a frequency resolution of 62.5 kHz. This yields a velocity resolution at 95 GHz of 0.236 km s<sup>-1</sup>with uniform weighting and 0.394 km s<sup>-1</sup>with Hanning smoothing. Some sources were reobserved with the correlator configured to a 32 MHz bandwidth, also with 1024 spectral channels. This yields a velocity resolution of 0.118 km s<sup>-1</sup>with uniform weighting and 0.197 km s<sup>-1</sup>with Hanning smoothing. For each source a uniformly weighted spectrum was produced with a velocity width of approximately 80 km s<sup>-1</sup>centred on the velocity of the previously detected class I or II methanol maser emission. The spectrum was then Hanning smoothed to improve the signal-to-noise ratio of weak sources. For the Mopra observations, the spectra in all figures and Gaussian parameters (peak flux density, velocity and full width half maximum) in all tables are Hanning smoothed data collected with the 64-MHz correlator configuration, unless otherwise noted in Table 1.
The source list was compiled primarilty from 44-GHz class I methanol masers detected in the southern hemisphere by Slysh et al. (1994) and Haschick et al. (1990). The observing list also included 12 class I methanol masers detected at 36 GHz at Puschino (Kalenskii, private communication), 55 class II methanol sources which had not previously been searched for class I maser emission (Caswell et al. 1995, Ellingsen et al. 1996, Ellingsen et al. 1999, Walsh et al. 1997) and 12 Hii regions for which no 44-GHz emission was detected by Slysh et al. (1994).
## 3 Results
A total 153 sites were searched for 95-GHz methanol emission, with detections observed toward 85. Table 1 contains a list of all detected sources along with the Gaussian parameters of their spectral features. Their spectra are shown in Fig. 1. The sources toward which no 95-GHz emission was detected are listed in Table 5, along with the 3$`\sigma `$ limit (typically less than 5 Jy). The kinematic distance to each source was estimated using the rotation curve of Brand & Blitz (1993). For many of the sources the model yields two distance estimates, The majority are likely to be at the nearer distance, but we have not attempted to resolve the ambiguity and where it occurs both distances are listed in Table 1 & 5. 86 known 44- or 36-GHz class I methanol sources were observed, with 95-GHz emission being detected toward 71 (83%). A total of 54 class II methanol maser sites which had not previously been searched for class I methanol emission were also observed, yielding a total of 13 detections (24%).
Seventy nine of the 85 sources in Table 1 are new detections, the other 6 were discovered previously by Val’tts et al. (1995) at Onsala. In Table 7 we compare the Gaussian parameters of the strongest features of 6 sources observed at both Onsala and Mopra. The radial velocities agree within 0.4 km s<sup>-1</sup>in all cases, which is approximately half the spectral resolution of the Onsala observations. The flux densities agree to within a factor of 2, which is reasonable, taking into account differences in the spectral resolution, absolute calibration and possible pointing errors. Examination of Table 7 shows that some of the masers (e.g. OMC-2 and NGC2264) have very narrow lines of 0.2 km s<sup>-1</sup>, which were completely unresolved with the 0.7 km s<sup>-1</sup>spectral resolution of the Onsala observations. The Onsala results have not been corrected for spectral smoothing, and for this reason, very narrow features in OMC-2, NGC2264 and W51Met2 will have lower peak flux density than that observed at Mopra. Below we give comments for some of the more interesting sources.
### 3.1 Comments on individual sources
OMC-2. The 95-GHz methanol maser in this source consists of a very intense narrow line similar to that observed at 44 GHz (Haschick et al. 1990). The peak flux densities at both frequencies are almost equal, as are the line widths.
NGC2264. This source shows an intense narrow line at 95 GHz with a shoulder on the positive velocity side. The spectrum of the 95-GHz methanol maser is very similar to that at 44-GHz, but the peak flux density is a factor of two lower.
305.21+0.21. The spectra of the 95- and 44-GHz masers are very similar, showing a single intense narrow line with the 95-GHz flux density only 25% lower than the 44-GHz flux density.
333.13-0.43. The spectrum of this source in Fig.1 shows a weak broad line. Comparing this spectrum with the 44-GHz spectrum we suspected that 95-GHz flux density was too low, possibly because of a pointing error. To check this we made a five-point map consisting of an observations at the central position and offset by $`\pm 20^{\prime \prime }`$ in both right ascension and declination. Fig. 7 contains the spectra from this map, which shows that the strongest emission is offset to the south-west of the nominal position and that there are at least two rather strong narrow lines confirming that this is a maser.
333.23-0.05. This source (Fig.1) has an intense isolated spectral feature, which is best fitted by two Gaussians. A five-point map (Fig. 8) shows that the source is in fact offset from the nominal position by approximately $``$20<sup>′′</sup> in right ascension. The 95-GHz maser spectrum is similar to that at 44 GHz, where it is a factor of 3 stronger.
335.59-0.29. The 95-GHz spectrum of this source is best fitted with 4 Gaussian components, while at 44 GHz only one intense component is present at the same radial velocity as the strongest 95-GHz component.
337.91-0.47. The 95-GHz methanol maser spectrum is essentially identical to the 44-GHz spectrum, but the peak flux density at 44 GHz is a factor of four higher.
338.92+0.56. The 95-GHz spectrum of this source has two narrow features, at radial velocities $``$62.9 km s<sup>-1</sup>and $``$60.1 km s<sup>-1</sup>, the same as the 44-GHz spectrum (Slysh et al. 1994). However, a five-point map centred on the nominal position of the source (Fig. 9) shows that the strongest emission is to the north, where the flux density is at least a factor of two higher. From the spectrum of the offset (0, 20<sup>′′</sup>) one can see also a third detail, at the radial velocity $``$65.2 km s<sup>-1</sup>, which is not present in the 44-GHz spectrum. The five-point map also shows that the position of the spectral feature at $``$60.1 km s<sup>-1</sup>is shifted in right ascension relative to the spectral feature at $``$62.9 km s<sup>-1</sup>by about 20<sup>′′</sup>.
343.12-0.06. The 95-GHz spectrum of this source has been fitted with four Gaussian components, all of which have counterparts in the 44-GHz spectrum (Slysh et al. 1994), although the relative intensities are different at the two frequencies. A five-point map of the source (Fig. 10) shows that the emission in the source is spread over an area of at least 20<sup>′′</sup>, so the difference in relative intensities is most likely due to the different telescope beamwidths and pointing errors between the 44- and 95-GHz observations.
345.01+1.79 The 95-GHz spectrum of this source has been fitted by two Gaussians separated by 1.2 km s<sup>-1</sup>. A 9-point map of this source (Val’tts 1998) was used to determined position of both components. The stronger component at -13.1 km s<sup>-1</sup>has a position which coincides within the errors with the position of the southern 6.7-GHz methanol maser (Norris et al. 1993). The weaker component at $``$14.3 km s<sup>-1</sup>is displaced from the stronger one by $``$6.6 $`\pm `$7.9<sup>′′</sup> in right ascension and 10.2$`\pm `$5.3<sup>′′</sup> in declination, and coincides within the errors with the position of the northern 6.7-GHz methanol maser. Thus 95-GHz methanol masers are present in both regions of methanol maser emission in this source.
345.00-0.22. The 95-GHz methanol maser spectrum is similar to that at 44 GHz, except for the presence of an additional feature at $``$24 km s<sup>-1</sup>, which is absent from the 44-GHz spectrum.
351.16+0.70 (NGC6334B). The 95-GHz spectrum is similar to 44-GHz spectrum, although due to blending of spectral features the number of Gaussian components we are able to fit is different, 4 at 95 GHz compared to 7 at 44 GHz.
351.24+0.67 (NGC6334C). The 95-GHz emission from this source is anomalously low compared to 44-GHz emission (Slysh et al. 1994), possibly due to poor pointing, or variable weather conditions.
351.41+0.64 (NGC6334F) A broad line with some weaker narrow features is present in the 95-GHz spectrum; we have fitted the profile with two broad components and one narrow at the radial velocity $``$6.4 km s<sup>-1</sup>. It is probably the counterpart of the -6.32 km s<sup>-1</sup>component in the 44-GHz spectrum (Slysh et al. 1994), but other narrow 44-GHz features are not present in the 95-GHz spectrum.
NGC6334I(N). This is one of the strongest methanol masers at both 95- and 44-GHz (Haschick et al. 1990). The spectra are similar at the two frequencies, considering likely differences in pointing, since the maser is known to be spread over an area more than 30<sup>′′</sup> in extent as can be seen from our map (Fig. 11). The relative positions of the four strongest spectral features were determined from this map and are given in Table 9. A comparison with the map of the source at 44 GHz (Kogan & Slysh 1998) shows that there is a general agreement between the two maps if larger position errors of 95-GHz map are taken into account. All four spectral features are point-like within the Mopra beamwidth of 52<sup>′′</sup>.
351.64$``$1.26. A five-point map of this source shows that the single weak feature becomes stronger at offsets of ($``$20<sup>′′</sup>,0) and (0,$``$20<sup>′′</sup>) (Fig. 12), which implies that the true source position is to the south-east of the the nominal position.
351.78-0.54. The 95-GHz spectrum consists of four narrow components, all of which are present in the 44-GHz spectrum, and in addition there is a broad component about 8 km s<sup>-1</sup>wide, which is absent at 44 GHz.
354.61+0.47. At least three Gaussian components can be fitted to the 95-GHz spectrum of this source, and two of them are present in the 44-GHz spectrum. The component at $``$21.4 km s<sup>-1</sup>is relatively stronger at 95 GHz and the component at $``$20.5 km s<sup>-1</sup>is present only at 95 GHz.
Sgr A-A. The 95-GHz spectrum of this source consists of a very broad component and possibly a weak narrow component at a velocity of 47 km s<sup>-1</sup>. It is almost identical to the 44-GHz spectrum (Haschick et al. 1990), with the narrow component relatively stronger at 44 GHz. A 9-point map (Fig. 13) of the source does not show the narrow component, probably because of the shorter integration time of map spectra. The broad component is extended in declination with an angular size of 150<sup>′′</sup>, and is unresolved (less than 30<sup>′′</sup>) in right ascension. This is one of the rare strong thermal sources in our sample.
M8E. A single very strong narrow line is present in the 95-GHz spectrum, which is similar to that at 44 GHz. This is one of the strongest masers in both the 44- and 95-GHz transitions, and its angular extent is known to be less than 0.2<sup>′′</sup> (Kogan & Slysh 1998, Slysh et al. 1999). A thorough discussion of the methanol emission from this source is given in Val’tts (1999a).
W33MetC. A single feature near 32.5 km s<sup>-1</sup>is present in 95-GHz spectrum. At 44 GHz there is also a weaker component at 36.28 km s<sup>-1</sup>which is not visible in 95-GHz spectrum. The 95-GHz emission in this source has been mapped by Pratap & Menten (1992), and the corresponding 44-GHz maser has been shown to coincide with its position (Slysh et al. 1999). A full discussion of the methanol emission from this source is given by Val’tts (1999b).
14.33-0.64. This intense class I methanol maser was discovered at 44 GHz at Parkes (Slysh et al. 1994). At 95-GHz the maser is also strong and its spectrum consists of four peaks, similar to the 44-GHz spectrum. A 44-GHz map of this maser is discussed in Slysh et al. (1999).
GGD27. 44-GHz class I methanol maser emission in this source was discovered by Kalenskii et al. (1992), and a VLA map of it is presented in Slysh et al. (1999). At both 95 and 44 GHz the spectrum is dominated by a strong narrow feature. A 9-point map of the 95-GHz emission (Fig. 14) shows that it is strongest at a position offset in right ascension by $``$23<sup>′′</sup>$`\pm `$8<sup>′′</sup> and in declination by $``$8<sup>′′</sup>$`\pm `$5<sup>′′</sup> from the nominal position.
L379IRS3. The 44-GHz methanol maser associated with this source was also discovered by Kalenskii et al. (1992) and mapped with the VLA by Kogan & Slysh (1998) and Slysh et al. (1999). The 95-GHz spectrum is very similar to that at 44 GHz and a map (Fig. 15) of the 95-GHz emission shows that the four spectral features are spread over an area of approximately 30<sup>′′</sup>, similar to the 44-GHz emission.
23.43-0.19. Two Gaussian components have been fitted to the detection at 95 GHz of this source, both are present in the 44-GHz spectrum (Slysh et al. 1994). In addition there are several components near 103 km s<sup>-1</sup>in the 44-GHz spectrum that are not present at 95 GHz.
IRAS18537+0749 (S76E). This rather strong source was discovered to be a class I methanol maser during observations of the $`4_13_0E`$ transition at 36 GHz at Puschino (Val’tts, private communication). In both transitions it appears that a blend of several narrow lines is producing a spectrum resembling a wide band of emission.
W51e1/e2. Only a broad component is present in our 95-GHz spectrum. The narrow spectral feature at 48.88 km s<sup>-1</sup>detected by Haschick et al. (1990) at 44 GHz is not visible in the 95-GHz spectrum, because it is shifted from W51e1/e2 by about 70<sup>′′</sup> and was outside the main beam of Mopra telescope (Pratap & Menten 1992).
## 4 Discussion
The spectra of the 95-GHz $`8_07_1A^+`$ methanol emission sources found in this survey are in general similar to the spectra of the corresponding 44-GHz $`7_06_1A^+`$ sources. The emission in the two transitions typically covers the same velocity range, has approximately the same number of spectral features with very similar radial velocities, and in some cases even the same relative intensities of the components. In Table 8 we list the single strongest feature in each spectrum at 95 GHz and the corresponding spectral features at 44 GHz of the sources detected by Slysh et al. (1994). One can see that there is always a corresponding spectral feature at 44 GHz to every 95 GHz spectral feature from Table 9, and their radial velocities agree in general to within 0.1 km s<sup>-1</sup>. The line width of the 95-GHz components is in general somewhat larger than the line width of the corresponding 44-GHz features, partly due to a lower spectral resolution in the 95 GHz observations, but nevertheless there are many very narrow 95-GHz features with a line width less than 1 km s<sup>-1</sup>. The peak flux density of the 95-GHz components is generally lower than the flux density of the 44-GHz features. Fig. 16 shows a comparison between the flux densities of spectral features with the same radial velocities from the two transitions. In constructing this plot data on all available sources were used, including the results of this work and of the observations at Onsala (Val’tts et al. 1995). The straight line (with a correlation coefficient $`r`$=0.73) shows the best fit linear dependence which was found to be :
$$S(95)=(0.32\pm 0.08)\times S(44)(8.1\pm 2.7)$$
(1)
Although the scatter is quite large, on average the 95-GHz methanol masers are a factor of 3 weaker than the 44-GHz masers. This result is consistent with the findings of Val’tts et al. (1995) who found a linear dependence between integrated fluxes at two frequencies. We have used peak flux densities rather than luminosities since our survey is flux density limited by the sensitivity of the instrument, and luminosities would give a spurious correlation due to the multiplication of the flux densities at the two frequencies by the same distance squared. The slope of the dependence between integrated flux densities found by Val’tts et al. (1995) is $`0.52\pm 0.05`$, which is larger than the slope $`0.32\pm 0.08`$ found in this paper for the peak flux densities. This difference may be due to the larger average line width of 95-GHz masers mentioned above. The correlation between the peak flux density and the observed similarity in the spectra of the two transitions is strong evidence in favour of the suggestion that the emission from both transitions arises from the same spatial location. A comparison of published high resolution maps of the 44- and 95-GHz class I methanol masers in DR21(OH) and W33MetC shows that their images are very alike and consist of the same number of isolated maser spots (Plambeck & Menten 1990, Pratap & Menten 1992, Kogan & Slysh 1998, Slysh et al. 1999), consistent with this hypothesis.
The two transitions belong to the class I methanol masers (Menten 1991), which are thought to be pumped through collisional excitation. The difference between the two transitions is that the upper level of the 95-GHz $`8_07_1A^+`$ transition is 18.5 K above the upper level of the 44-GHz $`7_06_1A^+`$ transition. Therefore the population of the former is expected to be lower than the population of the latter, resulting in the lower intensity of 95-GHz emission compared to the intensity of the 44-GHz transition, although it is difficult to estimate the difference without any knowledge of the kinetic temperature and particle density in the source.
We used LVG code to calculated the intensity ratios of the $`7_06_1A^+`$ and $`8_07_1A^+`$ transitions in a collisional excitation model for four different parameter sets. The model parameters and the intensity ratios are presented in Table 10. The collisional selection rules are based on the paper by Lees & Haque (1974) and imply that $`\mathrm{\Delta }K=0`$ collisions are preferred by a factor of four. For model 1 with a gas temperature 20 K and density $`0.56\times 10^5`$ cm<sup>-3</sup>, the ratio of the 44- and 95-GHz intensities is 3.3, i.e., close to the mean observed ratio. The 95 GHz intensity is lower due to the lower population of the $`8_0A^+`$ level relative to that of the $`7_0A^+`$ level and due to a weaker inversion at 95 GHz. Increasing either the gas temperature or the density decreased the model ratio below the observed value. Thus, our results favour class I maser model with gas temperature about 20 K and density less than 10<sup>6</sup> cm<sup>-3</sup>.
## 5 Summary
1. As a result of a survey in the southern hemisphere 85 methanol emission sources were detected in the $`8_07_1A^+`$ transition at 95 GHz. This survey together with a similar Onsala survey (Val’tts et al. 1995) completes a whole sky survey of methanol emission at 95 GHz.
2. Most of the detected sources are class I methanol masers, and the majority of them have counterparts in other class I methanol transitions, such as the $`7_06_1A^+`$ at 44 GHz.
3. The previously found correlation between the methanol maser emission intensity at 44 and 95 GHz is confirmed here, using a larger sample of sources.
4. A maser model with collisional excitation based on LVG calculations can explain the observed intensity ratio at 44 and 95 GHz and gives constraints on the temperature and particle density.
## 6 Acknowledgements
I.E.V. is grateful to the ATNF for the hospitality, and to the staff of Mopra observatory for the help with the observations. The Australia Telescope is funded by the Commonwealth of Australia for operation as a National Facility managed by CSIRO. Travel to Australia for I.E.V. was aided by grant 96/1990 from the Australian Department of Industry, Science and Tourism. The work of I.E.V., V.I.S., S.V.K. and G.M.L. was partly supported by the grants 95-02-05826 and 98-02-16916 from the Russian Foundation for Basic Research and by the Federal Program ”Astronomiya” (Project N 1.3.4.2). S.P.E thanks the Queen’s trust for the computing system used to process the data from these observations. The authors would like to thank Ms V. Oakley and Mr J. Saab for their assistance during the observations and initial data processing.
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# 1 Introduction.
## 1 Introduction.
In the framework of charge transport in semiconductors, a technique widely used in order to find approximate solutions of the Boltzmann transport equation (BTE) is based on a spherical harmonics expansion (SHE) of the distribution function (Rahmat, White and Antoniadis, 1996; Vecchi and Rudan, 1998; Ventura, Gnudi and Baccarani, 1995; Liotta and Struchtrup, 2000). Recently an asymptotic solution of the SHE equations was found (Liotta and Majorana, 1999) in the case of a homogeneous (bulk) device with a simple parabolic band structure. Despite the very simple situation in which this solution was obtained, it has revealed to be very useful in order to develop new asymptotic hydrodynamical models describing the hot electron population in silicon devices (both in the homogeneous and non-homogeneous case). In particular see Anile and Mascali (2000) and Anile, Liotta and Mascali (2000), where this asymptotic solution was used in order to close the set of moment equations.
The aim of this work is to show the possibility of finding a new asymptotic solution generalizing that derived in Liotta and Majorana (1999) to the case of a non-parabolic band structure (Kane model). this solution reduces to the old one when the non-parabolicity parameter goes to zero. The importance is due to the fact that the Kane equation fits better the real band structure in the high field regime. Therefore this solution can be very useful in order to develop improved high field hydrodynamical models which could describe better the hot electron population. This solution is also interesting by itself in the framework of SHE models.
## 2 Basic equations.
We consider the case of unipolar semiconductor devices in which the current is essentially due to electrons (but the results can be generalized to holes). The semiclassical description of the electron transport is based on the BTE (Markowich et al., 1990; Ferry, 1991; Cercignani, 1987), which writes
$$\frac{f}{t}+𝐯(𝐤)_𝐱f\frac{\text{e}}{\mathrm{}}𝐄_𝐤f=Q(f),$$
(1)
here $`f(t,𝐱,𝐤)`$ is the electron distribution function, generally depending on time $`t`$, position $`𝐱`$ and wave vector $`𝐤`$ (belonging to the first Brillouin zone $`B`$). e is the absolute value of the electron charge, $`\mathrm{}`$ the reduced Planck constant, $`𝐄`$ the electric field. $`_𝐱`$ and $`_𝐤`$ denote the gradient with respect to $`𝐱`$ and $`𝐤`$ respectively. The group velocity $`𝐯(𝐤)`$ is determined by the conduction band structure: $`𝐯(𝐤)=\frac{1}{\mathrm{}}_𝐤\epsilon (𝐤)`$, where $`\epsilon (𝐤)`$ is the electron energy which depends on the wave vector. $`Q`$ is the collision operator which in the non-degenerate case has the form
$$Q(f)=W(𝐤,\stackrel{~}{𝐤})f(t,𝐱,\stackrel{~}{𝐤})𝑑\stackrel{~}{𝐤}f(t,𝐱,𝐤)W(\stackrel{~}{𝐤},𝐤)𝑑\stackrel{~}{𝐤},$$
(2)
$`W(𝐤,\stackrel{~}{𝐤})`$ representing the electron scattering rate from a state with wave vector $`\stackrel{~}{𝐤}`$ to one with wave vector $`𝐤`$.
We will consider a stationary and homogeneous situation (bulk device), neglecting the Poisson equation and taking into account only a constant externally applied electric field, then it will be $`f=f(𝐤)`$. Moreover, we will suppose that the electric field is directed along the $`𝐱`$-axis so to have a cylindrical symmetry around this axis and represent the two-dimensional momentum space by means of the polar coordinates $`k=|𝐤|\left(=\xi (\epsilon )\right)`$ and $`\theta =\mathrm{arccos}(𝐤𝐤_x/|𝐤||𝐤_x|)`$, where $`𝐤_x`$ is the projection of $`𝐤`$ along $`x`$. Therefore the distribution function can be expanded in Legendre polynomials of the angle $`\theta `$ (Rahamat et al., 1996; Liotta and Struchtrup, 2000)
$$f(𝐤)=\underset{n}{}f_n(\epsilon )P_n(\mathrm{cos}\theta ),$$
(3)
where $`P_n`$ is the $`n`$th order Legendre polynomial. This expansion will be computationally viable only if few spherical harmonics are enough to accurately represent the momentum space distribution. We will assume that the first two terms of the previous expression give a good approximation
$$f(𝐤)f_0(\epsilon )+f_1(\epsilon )\mathrm{cos}\theta .$$
(4)
The lowest order harmonic coefficient furnishes information about the isotropic part of the distribution function and $`f_0𝑑𝐤`$ yields the electron concentration. The first order harmonic coefficient describes the asymmetry of the distribution function in the direction of the applied electric field, and $`f_1\mathrm{cos}\theta 𝐯(𝐤)𝑑𝐤/f_0𝑑𝐤`$ gives the hydrodynamical velocity of the electron gas.
We will assume a spherically symmetric band structure of the Kane form (Ferry, 1991, Jacoboni and Lugli, 1989; Tomizawa, 1993)
$$\gamma (\epsilon ):=\epsilon (1+\alpha \epsilon )=\frac{\mathrm{}^2k^2}{2m^{}},$$
(5)
where $`m^{}`$ is the electron effective mass and $`\alpha `$ the constant non-parabolicity parameter. By putting $`\alpha =0`$ one obtains the usual parabolic band approximation. With this choice we can assume for the first Brillouin zone $`B\text{}^3`$ and we have $`𝐯(𝐤)=\frac{\mathrm{}𝐤}{m^{}(2\alpha \epsilon +1)}`$.
As regards collisions, we will take into account the interaction between electrons and non-polar optical phonons and that between electrons and acoustical phonons, the latter in the elastic approximation, valid when the thermal energy is much greater than that of the phonon involved in the scattering. We consider the electron scatterings with ionized impurities to be negligible, i.e. we assume the doping density to be low. Then the transition rate of the collision operator reads (Jacoboni and Lugli, 1989; Tomizawa, 1993)
$$W(𝐤,\stackrel{~}{𝐤})=\text{K}_{op}\left[n_{op}\delta (\epsilon \stackrel{~}{\epsilon }\mathrm{}\omega _{op})+(n_{op}+1)\delta (\epsilon \stackrel{~}{\epsilon }+\mathrm{}\omega _{op})\right]+\text{K}_{ac}\delta (\epsilon \stackrel{~}{\epsilon }),$$
(6)
where $`\epsilon =\epsilon (𝐤)`$, $`\stackrel{~}{\epsilon }=\epsilon (\stackrel{~}{𝐤})`$, $`n_{op}=\left(\mathrm{exp}\left(\frac{\mathrm{}\omega _{op}}{k_BT_L}\right)1\right)^1`$ is the thermal equilibrium optical phonon number and $`\text{K}_{op}`$ and $`\text{K}_{ac}`$ are respectively the non-polar optical and acoustical kernel coefficients (constant at a first approximation). $`\mathrm{}\omega _{op}`$ is the optical phonon energy, $`k_B`$ the Boltzmann constant and $`T_L`$ the lattice temperature. These choices are appropriate for silicon devices.
The SHE equations are easily obtained by inserting the expansion (3) into the BTE (1) and balancing the terms of the same order in $`P_n(\mathrm{cos}\theta )`$. To generate a closed set of equations, all coefficients of order higher than the first are set to be zero, see Rahmat et al. (1996) (a closure inspired by the Grad moment method, see Grad, 1958).
But for the aims of this paper it is preferable to perform a change of variables and write down a set of two coupled equations in the unknowns
$`N(\epsilon )=\sigma (\epsilon )f_0(\epsilon ),`$ (7)
$`P(\epsilon )={\displaystyle \frac{8}{3}}\pi {\displaystyle \frac{\sqrt{m^{}}}{\mathrm{}^3}}\gamma (\epsilon )f_1(\epsilon ),`$ (8)
where
$$\sigma (\stackrel{~}{\epsilon }):=_\text{}^3\delta (\epsilon (𝐤)\stackrel{~}{\epsilon })𝑑𝐤=4\sqrt{2}\pi \left(\frac{\sqrt{m^{}}}{\mathrm{}}\right)^3H(\stackrel{~}{\epsilon })(\gamma (\stackrel{~}{\epsilon }))^{\frac{1}{2}}\gamma ^{}(\stackrel{~}{\epsilon })$$
is the density of states . $`H(\epsilon )`$ is the Heaviside step function and $`\gamma ^{}(\epsilon )\frac{d\gamma }{d\epsilon }=(1+2\alpha \epsilon )`$. So doing, the expansion (4) writes
$$f(𝐤)\frac{N(\epsilon )}{\sigma (\epsilon )}+\left(\frac{8}{3}\pi \frac{\sqrt{m^{}}}{\mathrm{}^3}\gamma (\epsilon )\right)^1P(\epsilon )\mathrm{cos}\theta .$$
These new variables have also a direct physical interpretation:
$$_0^+\mathrm{}N(\epsilon )𝑑\epsilon =_\text{}^3f_0(\epsilon )𝑑𝐤,_0^+\mathrm{}P(\epsilon )𝑑\epsilon =_\text{}^3v(\epsilon )f_1(\epsilon )𝑑𝐤,$$
(9)
($`v(\epsilon )=|𝐯_x|`$) and are very suitable for our problem.
With these choices the equations of our SHE model, in the stationary homogeneous case, write:
$`\text{e}E{\displaystyle \frac{dP}{d\epsilon }}=G_1(N)`$ (10)
$`\text{e}E{\displaystyle \frac{d(g(\epsilon )N)}{d\epsilon }}+\text{e}Eh(\epsilon )N=G_2(P)`$ (11)
where
$`G_1(N)`$ $`=`$ $`\text{K}_{op}\sigma (\epsilon )\left[(n_{op}+1)N(\epsilon +\mathrm{}\omega _{op})+n_{op}N(\epsilon \mathrm{}\omega _{op})\right]`$
$`\text{K}_{op}\left[n_{op}\sigma (\epsilon +\mathrm{}\omega _{op})(n_{op}+1)\sigma (\epsilon \mathrm{}\omega _{op})\right]N(\epsilon ),`$
$`G_2(P)`$ $`=`$ $`\left[n_{op}\text{K}_{op}\sigma (\epsilon +\mathrm{}\omega _{op})+(n_{op}+1)\text{K}_{op}\sigma (\epsilon \mathrm{}\omega _{op})+\text{K}_{ac}\sigma (\epsilon )\right]P(\epsilon )`$
and
$$g(\epsilon ):=\frac{2}{3}\frac{\gamma (\epsilon )}{m^{}(\gamma ^{}(\epsilon ))^2},h(\epsilon ):=\frac{1}{m^{}\gamma ^{}(\epsilon )}\frac{4}{3}\frac{\alpha }{m^{}}\frac{\gamma (\epsilon )}{(\gamma ^{}(\epsilon ))^3}.$$
We would like to underline that equations (10)-(11) can be obtained directly from the BTE by using a new alternative procedure. It consists in multiplying both sides of equation (1) respectively by $`\delta (\epsilon (𝐤)\stackrel{~}{\epsilon })`$ and by $`𝐯(𝐤)\delta (\epsilon (𝐤)\stackrel{~}{\epsilon })`$ and then formally integrating with respect to $`𝐤`$ over the whole space $`\text{}^3`$. Some suitable closure relations are needed: in particular by assuming that $`f`$ depends on $`𝐤`$ only through the variable $`\epsilon `$ one obtains equations (10)-(11) in the general non-homogeneous, non-stationary case (Liotta and Majorana, 1999; Majorana, 1998). This method is similar to the method of frequency dependent moments of radiation hydrodynamics (Thorne, 1981).
## 3 Dimensionless equations and physical parameters.
It is useful to introduce dimensionless variables: let
$$t_{}:=\left[4\sqrt{2}\pi \left(\frac{\sqrt{m^{}}}{\mathrm{}}\right)^3\sqrt{k_BT_L}\text{K}_{op}n_{op}\right]^1,\mathrm{}_{}:=\left(\frac{k_BT_L}{m^{}}\right)^{\frac{1}{2}}t_{},\epsilon _{}:=k_BT_L,$$
$$w:=\frac{\epsilon }{\epsilon _{}},n(w):=u_{}\mathrm{}_{}^3N(\epsilon ),p(w):=u_{}\mathrm{}_{}^2t_{}P(\epsilon ),$$
$$\lambda :=\frac{\mathrm{}\omega _{op}}{k_BT_L},a:=\frac{n_{op}+1}{n_{op}}=e^\lambda ,\kappa :=\frac{\text{K}_{ac}}{n_{op}\text{K}_{op}},\zeta :=\text{e}E\frac{\mathrm{}_{}}{u_{}},\beta :=\alpha \epsilon _{}.$$
Moreover in the following we put $`\chi (w):=w(1+\beta w)`$ and $`\chi ^{}(w)\frac{d\chi }{dw}=(1+2\beta w)`$. By using these new variables, equations (10)-(11) become
$`\zeta {\displaystyle \frac{dp}{dw}}=\mu (w)\left[an(w+\lambda )+n(w\lambda )\right]\left[\mu (w+\lambda )+a\mu (w\lambda )\right]n(w)`$ (12)
$`\zeta \left[{\displaystyle \frac{d(r(w)n)}{dw}}+q(w)n\right]=\left[\mu (w+\lambda )+a\mu (w\lambda )+\kappa \mu (w)\right]p(w)`$ (13)
where
$`\mu (w)`$ $`:=`$ $`H(w)\left[\chi (w)\right]^{\frac{1}{2}}\chi ^{}(w)`$
$`r(w)`$ $`:=`$ $`{\displaystyle \frac{2}{3}}{\displaystyle \frac{\chi (w)}{\left[\chi ^{}(w)\right]^2}}`$
$`q(w)`$ $`:=`$ $`{\displaystyle \frac{1}{\chi ^{}(w)}}{\displaystyle \frac{4}{3}}\beta {\displaystyle \frac{\chi (w)}{\left[\chi ^{}(w)\right]^3}}.`$
We associate the following conditions to equations (12)-(13)
$`n(0)=0(p(0)=0),\underset{w+\mathrm{}}{lim}p(w)=0`$
$`n(w)0w0,{\displaystyle _0^+\mathrm{}}n(w)𝑑w>0and<+\mathrm{}.`$
Since equations (12)-(13) are linear and homogeneous, if a solution $`(n(w),p(w))`$, satisfying the above conditions exists, then, for every $`c>0`$, also $`(cn(w),cp(w))`$ is a solution.
The appropriate values of the physical parameters, in the case of a silicon device, are given in table I, were $`m_e`$ denotes the electron rest mass.
| $`m^{}=0.32m_e`$ | $`T_L=300`$ K | $`\mathrm{}\omega _{op}=0.063`$ eV |
| --- | --- | --- |
| $`\text{K}_{op}={\displaystyle \frac{\left(D_tK\right)^2}{8\pi ^2\rho \omega _{op}}}`$ | $`D_tK=11.4`$ eV $`\stackrel{}{\text{A}}`$<sup>-1</sup> | $`\rho =2330`$ Kg m<sup>-3</sup> |
| $`\text{K}_{ac}={\displaystyle \frac{k_BT_L}{4\pi ^2\mathrm{}v_0^2\rho }}\mathrm{\Xi }_d^2`$ | $`\mathrm{\Xi }_d=9`$ eV | $`v_0=9040`$ m sec<sup>-1</sup>. |
| $`\alpha =0.5eV^1`$ | | |
Table I. Values of the physical parameters used in this paper.
Using these parameters, we get $`\lambda 2.437`$, $`\kappa 5.986`$ and $`\beta 0.012926`$.
## 4 Asymptotic equations.
Now we want to show that it is possible to find an approximate solution of the equations (12)-(13) valid for high values of the electron energy.
It is useful to introduce a new variable $`\psi (w)`$ defined by
$$n(w)=\mu (w)\psi (w).$$
(14)
Equations (12)-(13) become
$`\zeta {\displaystyle \frac{dp}{dw}}`$ $`=`$ $`\mu (w)\left[a\mu (w+\lambda )\psi (w+\lambda )+\mu (w\lambda )\psi (w\lambda )\right]`$ (15)
$`\mu (w)\left[\mu (w+\lambda )+a\mu (w\lambda )\right]\psi (w)`$
$`{\displaystyle \frac{2}{3}}\zeta {\displaystyle \frac{\left[\chi (w)\right]^{\frac{3}{2}}}{\chi ^{}(w)}}{\displaystyle \frac{d\psi }{dw}}`$ $`=`$ $`\left[\mu (w+\lambda )+a\mu (w\lambda )+\kappa \mu (w)\right]p(w).`$ (16)
Because we look for an asymptotic form of the equations (15)-(16) for large values of the energy $`w`$, we expand the coefficients of the equations up to the zeroth order in $`\lambda `$: $`\mu (w\pm \lambda )\mu (w)`$. In this way we obtain a new set of equations
$`\zeta {\displaystyle \frac{dp_A}{dw}}`$ $`=`$ $`\mu ^2(w)\left[a\psi _A(w+\lambda )+\psi _A(w\lambda )(a+1)\psi _A(w)\right]`$ (17)
$`p_A(w)`$ $`=`$ $`{\displaystyle \frac{2}{3}}\zeta {\displaystyle \frac{\left[\chi (w)\right]^{\frac{3}{2}}}{\chi ^{}(w)\mu (w)\left[1+a+\kappa \right]}}{\displaystyle \frac{d\psi _A}{dw}},`$ (18)
where the subscript $`A`$ label the new unknowns. Substituting (18) into (17), it follows
$`{\displaystyle \frac{2}{3}}{\displaystyle \frac{\zeta ^2}{(1+a+\kappa )}}\left[{\displaystyle \frac{\chi (w)}{\left[\chi ^{}(w)\right]^2}}\psi _A^{\prime \prime }+\left({\displaystyle \frac{1}{\chi ^{}(w)}}{\displaystyle \frac{4\beta \chi (w)}{\left[\chi ^{}(w)\right]^3}}\right)\psi _A^{}\right]=`$
$`\chi (w)\left[\chi ^{}(w)\right]^2\left(a\psi _A\left(w+\lambda \right)+\psi _A\left(w\lambda \right)(a+1)\psi _A\left(w\right)\right),`$ (19)
where the primes denote derivatives with respect to $`w`$.
## 5 Approximate solution.
In order to find an approximate solution of equation (19) we expand the coefficients up to the first order in the non-parabolicity parameter $`\beta `$.
$`{\displaystyle \frac{2}{3}}{\displaystyle \frac{\zeta ^2}{(1+a+\kappa )}}\left[\left(w3\beta w^2\right)\psi _A^{\prime \prime }+\left(16\beta w\right)\psi _A^{}\right]=`$
$`\left(w+5\beta w^2\right)\left(a\psi _A\left(w+\lambda \right)+\psi _A\left(w\lambda \right)(a+1)\psi _A\left(w\right)\right).`$ (20)
This choice is justified by the smallness of $`\beta `$ and by the Kane equation itself, which is of the first order in the non-parabolicity parameter.
We will search for solutions of equation (20) having the form
$$\psi _A(w)=e^{f(w)},withf(w):=\eta _0w+\eta _1\beta w^2,$$
(21)
where $`\eta _0`$ and $`\eta _1`$ are functions of the applied electric force $`\zeta `$.
It is useful to observe that $`f(w\pm \lambda )=f(w)+f(\pm \lambda )\pm 2\eta _1\beta \lambda w`$. Expanding the following quantities up to the first order in $`\beta `$ :
$`e^{\pm 2\eta _1\beta \lambda w}`$ $``$ $`1\pm 2\eta _1\beta \lambda w,`$
$`e^{f(\pm \lambda )}`$ $``$ $`e^{\pm \eta _0\lambda }\left(1+\eta _1\lambda ^2\beta \right),`$
substituting (21) into (20) and retaining only terms up to first order in $`\beta `$, we obtain, after dropping the common factor $`e^{f(w)}`$, the equation
$`{\displaystyle \frac{2}{3}}{\displaystyle \frac{\zeta ^2}{(1+a+\kappa )}}\left[\eta _0+\left(\eta _0^2+4\beta \eta _16\beta \eta _0\right)w+\left(3\beta \eta _0^2+4\beta \eta _0\eta _1\right)w^2\right]=`$
$`\left[\left(ae^{\eta _0\lambda }+e^{\eta _0\lambda }\right)\left(1+\beta \eta _1\lambda ^2\right)(a+1)\right]w+`$
$`\left[ae^{\eta _0\lambda }\left(2\eta _1\beta \lambda +5\beta \right)+e^{\eta _0\lambda }\left(2\eta _1\beta \lambda +5\beta \right)5\beta (a+1)\right]w^2.`$ (22)
If we divide both sides of equation (22) by $`w^2`$, and neglect the terms in $`\frac{1}{w^2}`$, but not those in $`\frac{1}{w}`$ (in some sense we are serching for a ”weakly asymptotic” solution), we obtain the following system of two transcendent equations in the unknowns $`\eta _0`$ and $`\eta _1`$
$`{\displaystyle \frac{2}{3}}{\displaystyle \frac{\zeta ^2}{(1+a+\kappa )}}\left(\eta _0^26\beta \eta _0+4\beta \eta _1\right)=\left(ae^{\eta _0\lambda }+e^{\eta _0\lambda }\right)\left(1+\beta \eta _1\lambda ^2\right)(a+1)`$ (23)
$`{\displaystyle \frac{2}{3}}{\displaystyle \frac{\zeta ^2}{(1+a+\kappa )}}\left(4\eta _0\eta _13\eta _0^2\right)=ae^{\eta _0\lambda }\left(5+2\eta _1\lambda \right)+e^{\eta _0\lambda }\left(52\eta _1\lambda \right)5(a+1),`$ (24)
where in the second equation we have dropped the common factor $`\beta `$. If we are able to solve the previous system, it is possible to obtain $`\eta _0`$ and $`\eta _1`$ as functions of the applied electric force $`\zeta `$, and then $`\psi _A(w)`$. Moreover, by using equation (18), one can find $`p_A(w)`$.
Henceforth we will indicate as asymptotic solution of the SHE equations the expressions of $`n_A`$ anp $`p_A`$ which can be obtained by means of the approximate solution of (17)-(18) which have been found.
## 6 Discussion of the solution and comparison with numerical results.
An analytical solution of the equations (23)-(24) has turned out to be impossible. Therefore we have limited ourselves to a graphical and numerical analysis. Given a value of the applied electric force $`\zeta `$, the requirement that the functions be integrable in $`[0,+\mathrm{}[`$ tells us that both $`\eta _0`$ and $`\eta _1`$ have to assume negative values. In fact, we found a negative solution of the system (23)-(24) in the domain $`[1,0]\times [1,0]`$ of the $`(\eta _0,\eta _1)`$ plane, for all the values of the electric field in the explored range. We used a simple MapleV algorithm in order to find the solutions. In table II we give some of the values we found.
| $`E`$ (V/cm) | $`\eta _0`$ | $`\eta _1`$ |
| --- | --- | --- |
| $`0.0`$ | $`1.0`$ | $`0.0`$ |
| $`1.0\times 10^2`$ | $`0.9999636274`$ | $`0.0001493648`$ |
| $`1.0\times 10^3`$ | $`0.9963693066`$ | $`0.0148712049`$ |
| $`5.0\times 10^3`$ | $`0.9136708716`$ | $`0.3294754506`$ |
| $`1.0\times 10^4`$ | $`0.7162321385`$ | $`0.8333494338`$ |
| $`3.0\times 10^4`$ | $`0.2511364699`$ | $`0.5681222515`$ |
| $`5.0\times 10^4`$ | $`0.1270766899`$ | $`0.2772663271`$ |
| $`7.0\times 10^4`$ | $`0.0780599266`$ | $`0.1591756496`$ |
| $`1.0\times 10^5`$ | $`0.0452815087`$ | $`0.0850850173`$ |
Table II. Some values of $`\eta _0`$ and $`\eta _1`$ as functions of the electric field.
In figures 1, 2 and 3, we compare the asymptotic solution $`(n_A,p_A)`$ and the numerical solution $`(n_N,p_N)`$ of equations (12)-(13), the latter obtained by a suitable numerical technique (Liotta and Majorana, 1999), for some values of the applied electric field.
As is possible to see, the agreement between the numerical and asymptotic solutions is good in all the energy range, despite the obtained asymptotic solution should be good only for high enough energy values (the agreement is obviously very good in this region). We want to observe that for low values of the electric field ($`|𝐄|10^3`$ V/cm), instead of formula (18), we use the following expression for $`p_A`$:
$$p_A(w)=\frac{1}{\mu (w+\lambda )+a\mu (w\lambda )+\kappa \mu (w)}\frac{2}{3}\zeta \frac{\left[\chi (w)\right]^{\frac{3}{2}}}{\chi ^{}(w)}\frac{d\psi _A}{dw}.$$
(25)
This choice is not coherent with the expansion, but allows a better agreement with the numerical solution and puts in evidence a discontinuity in the derivatives at the point $`w=\lambda `$ (in dimensional variables $`\epsilon =\mathrm{}\omega _{op}`$) due to the term $`\mu (w\lambda )`$, that is zero for $`0w\lambda `$. This term is also present in the original set of equations (12)-(13). Then, such discontinuity is also expected in the solutions.
As a measure of the difference between the numerical and asymptotic solution we can also compute $`V_{as}=p_A𝑑\epsilon /n_A𝑑\epsilon `$ and $`V_{num}=p_N𝑑\epsilon /n_N𝑑\epsilon `$, which give, in dimensional units, the electron hydrodynamical velocity. In table III we compare the values we found.
| $`E`$ (V/cm) | $`V_{as}(m/sec)`$ | $`V_{num}(m/sec)`$ |
| --- | --- | --- |
| $`1.0\times 10^2`$ | $`1.4879\times 10^3`$ | $`1.4537\times 10^3`$ |
| $`1.0\times 10^3`$ | $`1.4831\times 10^4`$ | $`1.3858\times 10^4`$ |
| $`5.0\times 10^3`$ | $`7.0223\times 10^4`$ | $`5.1060\times 10^4`$ |
| $`1.0\times 10^4`$ | $`6.0652\times 10^4`$ | $`7.3611\times 10^4`$ |
| $`3.0\times 10^4`$ | $`9.8325\times 10^4`$ | $`9.7714\times 10^4`$ |
| $`5.0\times 10^4`$ | $`1.0106\times 10^5`$ | $`9.9872\times 10^4`$ |
| $`7.0\times 10^4`$ | $`9.5124\times 10^4`$ | $`9.8395\times 10^4`$ |
| $`1.0\times 10^5`$ | $`8.2528\times 10^4`$ | $`9.5101\times 10^4`$ |
Table III. Comparison between the electron hydrodynamical velocities calculated by using respectively the asymptotic and numerical solutions.
The behaviour of the hydrodynamical velocity is the typical one when the Kane equation is used (see Tomizawa, 1993, p. 100, fig. 3.11). We do not consider higher values of the electric field because in this case also the Kane equation becomes inadequate to describe the conduction band, and a full band structure should be used (Vecchi and Rudan, 1998).
## 7 Conclusions and acknowledgments.
We have shown that it is possible to find an ”asymptotic solution” of the SHE equation (12)-(13) in which the dependence on the applied electric field is given implicitly through the system of transcendent equations (23)-(24). The solution of the system (23)-(24) can be obtained by using standard numerical techniques. The agreement with the numerical solution is good in all the explored range of applied electric fields.
If we put $`\beta =0`$, we recover the parabolic band case and the asymptotic solution reduces to the asymptotic one already found by Liotta and Majorana (1999) (L-M solution). The L-M solution was used by Anile and Mascali (2000) to obtain a two fluid hydrodynamical model, where the electron population is splitted in two sub-populations: low-energy and high-energy electrons, separated by a threshold energy. In order to close the moment equations relatively to the hot electrons they used a distribution function whose form is directly suggested by the L-M solution, whereas relatively to cold electrons a a standard maximum-entropy distribution function is utilized. A detailed study of the resulting high-energy hydrodynamical model is given in Anile, Liotta and Mascali (2000). Some works about the extension of this model to the case of Kane band, using the asymptotic solution found in this paper, are in progress. A better description of the hot electron population is expected.
The author acknowledges support from Italian CNR (Prog. N.96.03855.CT01), from TMR (Progr. n. ERBFMRXCT970157), from Italian MURST (Prot. n. 9801169828-005) and from ”Convenzione quadro Università di Catania-ST Microelectronics”.
The author wish also to thank prof. A. M. Anile, prof. A. Majorana and dr. G. Mascali for useful discussions and precious suggestions.
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# Influence of new reaction rates on 18F production in novae
## 1 Introduction
Classical novae emit gamma-ray radiation at and below 511 keV during the early epochs after the explosion. This emission is produced by electron-positron annihilation in the expanding envelope, and the subsequent Comptonization of the resulting gamma-ray photons, and it shows a line, at 511 keV, and a continuum, between 20 and 511 keV (Gómez-Gomar et al. (1998)). The positrons responsible for this emission come mainly from the disintegration of <sup>18</sup>F (Leising & Clayton (1987); Gómez-Gomar et al. (1998)), because its lifetime ($`\tau `$=158 min) is such that positrons are emitted at the “right time”, i.e., when the expanding envelope starts to be transparent to gamma-ray radiation. Therefore, the amount of radiation emitted strongly depends on the <sup>18</sup>F content of the nova envelope.
The synthesis of <sup>18</sup>F in novae depends largely on some key nuclear reaction rates of <sup>18</sup>F destruction and production which are far from being well known. This is the case, in particular, of the <sup>18</sup>F(p,$`\gamma )^{19}`$Ne and <sup>18</sup>F(p,$`\alpha )^{15}`$O reactions. Recent experimental studies (Graulich et al. (1997); Utku et al. (1998)) drastically improved the knowledge of these reaction rates with respect to previous studies (Wiescher & Kettner (1982)). In a recent paper (Hernanz et al. (1999)), we have analyzed the influence of these rates (Utku et al. (1998)) of <sup>18</sup>F(p,$`\gamma )^{19}`$Ne and <sup>18</sup>F(p,$`\alpha )^{15}`$O on the final yields of <sup>18</sup>F for different models of CO and ONe novae. The effect of the new rates was important, since a factor of 10 reduction in the yields and in the resulting gamma-ray fluxes was obtained for all the models. Therefore, we concluded that a more detailed analysis of the reaction rates was necessary, in order to predict the gamma-ray emission of classical novae.
The rates proposed by Utku et al. (1998) were limited to relatively high resonance energies and temperature domains, more appropriate for temperatures typical of X–ray bursts than for those of classical novae. The reason is that for their determination neither the influence of uncertainties on low energy resonance strengths (based on assumed reduced width) nor the effect of all low energy resonance tails were considered. At higher temperatures the rates are more reliable since their main contributions come from two directly measured resonances. The purpose of this paper is to provide rates for the <sup>18</sup>F(p,$`\gamma )^{19}`$Ne and <sup>18</sup>F(p,$`\alpha )^{15}`$O reactions valid in the domain of temperature of nova nucleosynthesis, incorporating the latest experimental data. In addition to the nominal rate, we provide upper and lower limits. Since other nuclear reactions also affect <sup>18</sup>F synthesis in novae, a global analysis is done, including results from the recent NACRE compilation (Angulo et al. (1999)). With the new rates affecting <sup>18</sup>F synthesis, new nova models have been computed, in order to determine the mass of <sup>18</sup>F they eject and the impact of the new yields on their early gamma-ray emission.
The organisation of this paper is the following. In section 2 we describe the nucleosynthesis of <sup>18</sup>F. In section 3 we discuss in detail the <sup>18</sup>F+p rates and the corresponding uncertainties while in the following section (4), we briefly discuss other recently published rates. In section 5 we present new results that show the influence of these new rates on <sup>18</sup>F production. Following the conclusion, an appendix gives some analytical approximations to the rates.
## 2 Synthesis of <sup>18</sup>F in classical novae
In this Section, we will focus on the nuclear paths leading to <sup>18</sup>F synthesis, through a detailed analysis of a representative model of an ONe nova : a 1.25 M white dwarf, accreting solar-like matter at a rate of $`2\times 10^{10}`$ M.yr<sup>-1</sup> , assuming a 50% degree of mixing with the outermost core of composition taken from Ritossa, García–Berro & Iben (1996) (mainly <sup>16</sup>O and <sup>20</sup>Ne with traces of other Ne, Na and Mg isotopes). The evolutionary sequences leading to a nova outburst have been computed with an updated version of the code SHIVA (see José & Hernanz 1998), a one-dimensional, implicit, hydrodynamical code in Lagrangian formulation, which follows the course of the explosion from the onset of accretion up to the expansion and ejection stages. The nuclear reaction network is described in José et al. (1999), and in Hernanz et al. (1999). In particular, we use the reaction rates for proton captures on <sup>18</sup>F and <sup>17</sup>O based on Utku et al. (1998) and Landré et al. (1989), respectively.
The synthesis of <sup>18</sup>F in nova outbursts takes place within the hot-CNO cycle (see fig. 1). Regardless of the nova type (either CO or ONe, according to the composition of the underlying white dwarf), the initial <sup>16</sup>O abundance is quite large (up to 25% by mass). Hence, initial <sup>16</sup>O is the main source of the formation of <sup>18</sup>F, either through the chain <sup>16</sup>O(p,$`\gamma `$)<sup>17</sup>F(p,$`\gamma `$)<sup>18</sup>Ne($`\beta ^+`$)<sup>18</sup>F or through <sup>16</sup>O(p,$`\gamma `$)<sup>17</sup>F($`\beta ^+`$)<sup>17</sup>O(p,$`\gamma `$)<sup>18</sup>F, which reflects a competition between (p,$`\gamma `$) reactions and $`\beta ^+`$-decays.
Snapshots of the evolution of isotopes relevant to <sup>18</sup>F synthesis (i.e., <sup>16,17</sup>O, <sup>17,18,19</sup>F and <sup>18,19</sup>Ne) are shown in Fig. 2. At the early stages of the explosion, when the temperature achieved at the burning shell reaches $`T_{bs}5\times 10^7`$ K (Fig. 2, first panel), the main nuclear activity in this mass region is driven by <sup>16</sup>O(p,$`\gamma `$)<sup>17</sup>F, followed by its $`\beta ^+`$-decay to <sup>17</sup>O. However, the temperature is not high enough to burn a significant fraction of <sup>16</sup>O. On the contrary, the abundance of <sup>17</sup>O rises by several orders of magnitude with respect to its initial content. Both <sup>17,18</sup>F are being synthesized at the burning shell by means of proton captures on <sup>16,17</sup>O, respectively, whereas the amount of <sup>19</sup>F remains essentially unchanged. The minor activity driven by proton capture reactions onto <sup>17,18</sup>F results also in a moderate increase of both <sup>18,19</sup>Ne (below $`10^{10}`$ by mass).
A similar trend is found when $`T_{bs}7\times 10^7`$ K (Fig. 2, second panel). The temperature rise increases the number of proton captures onto <sup>16</sup>O, leading to <sup>17</sup>F, which rapidly decays into <sup>17</sup>O. The effect of convection, which already extends through most of the envelope (i.e., $`150`$ km), is shown in the smooth distribution of both <sup>17,18</sup>F, previously synthesized in the burning shell and efficiently carried away to the outer envelope. <sup>19</sup>F is in turn reduced down to 50% due to <sup>19</sup>F(p,$`\alpha `$)<sup>16</sup>O, which dominates <sup>19</sup>Ne($`\beta ^+`$)<sup>19</sup>F.
When $`T_{bs}`$ reaches $`10^8`$ K (Fig. 2, third panel), <sup>17</sup>O shows a flat profile along the envelope, with a mean mass fraction of $`3\times 10^3`$, due to <sup>16</sup>O(p,$`\gamma `$)<sup>17</sup>F($`\beta ^+`$)<sup>17</sup>O, which dominates destruction through <sup>17</sup>O(p,$`\gamma `$)<sup>18</sup>F. With respect to fluorine isotopes, we stress an important synthesis of <sup>17,18</sup>F at this stage, driven by <sup>16</sup>O(p,$`\gamma `$)<sup>17</sup>F and <sup>17</sup>O(p,$`\gamma `$)<sup>18</sup>F, whereas <sup>19</sup>F is almost fully destroyed. <sup>18,19</sup>Ne continue to rise due to proton captures onto <sup>17,18</sup>F.
A dramatic change is found as soon as the burning shell reaches $`2\times 10^8`$ K (Fig. 2, fourth panel). At this time, the mass fraction of <sup>16</sup>O has been reduced down to $`3\times 10^2`$. For the first time, <sup>17</sup>O decreases (by one order of magnitude with respect to the amount shown in previous panel), since destruction through <sup>17</sup>O(p,$`\gamma `$)<sup>18</sup>F dominates <sup>17</sup>F($`\beta ^+`$)<sup>17</sup>O in the vicinity of the burning shell. On the contrary, <sup>17</sup>F exhibits a significant rise, up to $`3\times 10^2`$ by mass (i.e., one order of magnitude higher than in previous panel), so that the amount of <sup>17</sup>F becomes larger than that of <sup>17</sup>O in almost all the envelope. Again, <sup>16</sup>O(p,$`\gamma `$)<sup>17</sup>F dominates both <sup>17</sup>O(p,$`\gamma `$)<sup>18</sup>Ne and <sup>17</sup>F($`\beta ^+`$)<sup>17</sup>O near the burning shell. The evolution of <sup>19</sup>F is still dominated by <sup>19</sup>F(p,$`\alpha `$)<sup>16</sup>O, which in fact has no significant influence on the <sup>16</sup>O content. Also noticeable is the dramatic rise of both <sup>18,19</sup>Ne, which increase by several orders of magnitude, mainly due to the fact that at such temperatures, proton captures onto <sup>17,18</sup>F become faster than the corresponding $`\beta ^+`$-decays. It is worth noticing that at this stage, a non-negligible path leading to <sup>18</sup>F synthesis is now driven by <sup>18</sup>Ne($`\beta ^+`$)<sup>18</sup>F, basically at the outer envelope shells, which compensates an efficient destruction through proton captures near the burning shell.
The burning shell achieves a maximum temperature of $`2.51\times 10^8`$ K (Fig. 2, fifth panel). <sup>16</sup>O is reduced to 0.096 by mass in the burning shell. Whereas <sup>17</sup>O is destroyed by (p,$`\alpha `$) and (p,$`\gamma `$) reactions at the burning shell, its mean abundance in the envelope increases due to <sup>17</sup>F($`\beta ^+`$)<sup>17</sup>O which dominates in almost all the envelope. On the contrary, <sup>17</sup>F is generated in the burning shell by means of <sup>16</sup>O(p,$`\gamma `$)<sup>17</sup>F, whereas it is being destroyed by $`\beta ^+`$-decays in the outer envelope. The evolution of <sup>18</sup>F reflects a competition between destruction and creation modes at different locations of the envelope: whereas its content is reduced by both <sup>18</sup>F(p,$`\alpha `$)<sup>15</sup>O reactions at the burning shell and by <sup>18</sup>F($`\beta ^+`$)<sup>18</sup>O at the outer envelope, a dominant source for <sup>18</sup>F synthesis through <sup>17</sup>O(p,$`\gamma `$)<sup>18</sup>F is found at the intermediate shells. Proton captures onto <sup>17,18</sup>F and convective transport continue to pump <sup>18,19</sup>Ne to the outer envelope shells, which is at the origin of the rise of <sup>19</sup>F at this stage (through <sup>19</sup>Ne($`\beta ^+`$)<sup>19</sup>F).
Shortly after, due to the sudden release of energy from the short-lived species <sup>13</sup>N, <sup>14,15</sup>O and <sup>17</sup>F, the envelope begins to expand. As a result of the drop in temperature, (p,$`\gamma `$) and (p,$`\alpha `$) reactions are basically restricted to the vicinity of the burning shell, whereas most of the envelope is dominated by $`\beta ^+`$-decays (Fig. 2, sixth and seventh panels). Hence, whereas <sup>17</sup>O is powered by <sup>17</sup>F($`\beta ^+`$)<sup>17</sup>O along the envelope, the amount of <sup>17</sup>F and <sup>18,19</sup>Ne decreases significantly. The <sup>18</sup>F abundance at these late-time stages of the outburst increases due to <sup>17</sup>O(p,$`\gamma `$)<sup>18</sup>F, which is dominant at the intermediate layers of the envelope (whereas in the vicinity of the burning shell <sup>18</sup>F is efficiently destroyed by <sup>18</sup>F(p,$`\alpha `$)<sup>15</sup>O, the outer envelope is dominated by <sup>18</sup>F($`\beta ^+`$)<sup>18</sup>O). Moreover, destruction of <sup>19</sup>F through <sup>19</sup>F(p,$`\alpha `$)<sup>16</sup>O at the burning shell is dominated by <sup>19</sup>Ne($`\beta ^+`$)<sup>19</sup>F elsewhere.
The final stages of the outburst (Fig. 2, eighth panel), as the envelope expands and cools down, are dominated by the release of nuclear energy by $`\beta ^+`$-decays such as <sup>18</sup>F($`\beta ^+`$)<sup>18</sup>O, or <sup>17</sup>F($`\beta ^+`$)<sup>17</sup>O. The resulting mean abundance of <sup>18</sup>F in the ejected shells in this Model is X(<sup>18</sup>F)= $`2.7\times 10^4`$. Most of the envelope is, however, dominated by the large abundances of <sup>16,17</sup>O (X(<sup>16</sup>O)= $`6.8\times 10^2`$, X(<sup>17</sup>O)= $`3.9\times 10^2`$). A residual <sup>17</sup>F, which is still decaying into <sup>17</sup>O, and a non-negligible amount of <sup>19</sup>F ($`3\times 10^6`$, by mass) are also present.
In summary, the synthesis of <sup>18</sup>F in classical novae is essentially controlled by five proton-capture reactions, <sup>16</sup>O(p,$`\gamma `$)<sup>17</sup>F, <sup>18</sup>F(p,$`\gamma `$)<sup>19</sup>Ne, <sup>18</sup>F(p,$`\alpha `$)<sup>15</sup>O, <sup>17</sup>O(p,$`\gamma `$)<sup>18</sup>F, <sup>17</sup>O(p,$`\alpha `$)<sup>14</sup>N, and several $`\beta ^+`$-decays. Accordingly, the corresponding reaction rates deserve further attention and their uncertainties are discussed in Sect. 3 for the <sup>18</sup>F+p rates and Sect. 4 for the other rates.
## 3 The <sup>18</sup>F(p,$`\alpha )^{15}`$O and <sup>18</sup>F(p,$`\gamma )^{19}`$Ne reaction rates
Uncertainties on the <sup>18</sup>F(p,$`\gamma )^{19}`$Ne and <sup>18</sup>F(p,$`\alpha )^{15}`$O reaction rates used to reach many orders of magnitude due to the very limited spectroscopic data available for <sup>19</sup>Ne in the domain of interest. The rate estimated by Wiescher & Kettner (1982) has become obsolete by the recent measurements of Rehm et al. (1995), Coszach et al. (1995), Rehm et al. (1996, 1997) and in particular by the work of Graulich et al. (1997), Utku et al. (1998) and Butt et al. (1998). We concentrate on the faster <sup>18</sup>F(p,$`\alpha )^{15}`$O reaction but most of the discussion also applies to the <sup>18</sup>F(p,$`\gamma )^{19}`$Ne one.
Before 1997 and below the excitation energy ($`E_x`$) corresponding to the resonance energy ($`E_r`$) of 0.4 MeV, only two levels were known $`E_x`$ = 6.437 and 6.742 MeV but with unknown or uncertain spins and parity ($`J^\pi `$)<sup>1</sup><sup>1</sup>1For more detailed explanations of the standard nuclear physics notations used, see e.g. the appendix in José et al. (1999).. On the contrary, various levels were known in the corresponding region of <sup>19</sup>F, the conjugate nucleus (see Fig. 3). For instance, in the $`E_r`$=0.–1. MeV region, only 6 levels were known and many spins unknown in <sup>19</sup>Ne, while 19 levels had been observed in the corresponding <sup>19</sup>F region. Hence, for the estimation of the rates, Wiescher & Kettner (1982) considered three known levels at $`E_x`$ = 6.437, 6.742 and 6.862 MeV and postulated three others. The corresponding rate was highly uncertain since the location of these three postulated levels was approximate and because other levels found in <sup>19</sup>F are expected to have counterparts in <sup>19</sup>Ne in the region of interest. This is the case, in particular, of the $`E_x,J^\pi `$ = 6.429, 1/2<sup>-</sup>, 6.497, 3/2<sup>+</sup> and 6.527, 3/2<sup>+</sup> levels of <sup>19</sup>F corresponding to strong expected resonances because of their low centrifugal barrier, i.e. low transfered orbital angular momentum ($`\mathrm{}_p`$) of 1, 0 and 0 respectively.
Even though the recent experiments have not been able to find all the counterparts of the <sup>19</sup>F levels (see Fig. 3), they provided direct measurement for the strengths of two resonances which likely dominate the rate in the domain of nova nucleosynthesis and the location of several new <sup>19</sup>Ne levels. The strength of the resonance corresponding to the $`E_x`$ = 6.742 MeV level ($`E_r`$=330 keV) has recently been measured directly by Graulich et al. (1997) with a <sup>18</sup>F beam provided by the Louvain–La–Neuve facility. This level is thought to be the analog of the 6.787, 3/2<sup>-</sup> <sup>19</sup>F level. The measured strength ($`\omega \gamma `$), $`3.5\pm 1.6`$ eV (Graulich et al. (1997)), is in good agreement with the Wiescher & Kettner (1982) estimate and accordingly does not induce a significant change in the rate. On the contrary, the resonance associated with the $`E_x,J^\pi `$ = 7.067 MeV, 3/2<sup>+</sup> level ($`E_r`$=659 keV) strongly modifies the rates. It has been studied by Rehm et al. (1995, 1996), Coszach et al. (1995), Graulich et al. (1997), Utku et al. (1998) and Butt et al. (1998) and its strength has been measured directly (Coszach et al. (1995); Graulich et al. (1997)). It is located well outside of the Gamow peak but due to its total width ($`\mathrm{\Gamma }`$30 keV), the contribution of its tail alone is greater than the Wiescher & Kettner (1982) rate in the domain of nova nucleosynthesis. However, no counterpart of this broad level was known in <sup>19</sup>F. The closest 3/2<sup>+</sup> <sup>19</sup>F level is located at $`E_x`$ = 7.262 but its width is much smaller ($`\mathrm{\Gamma }<`$6 keV). To remove this incompatibility, Rehm et al. (1996) claimed that the 7.067 MeV width could be smaller ($`\mathrm{\Gamma }`$14 keV). A 7.114, 7/2<sup>+</sup> level was known in <sup>19</sup>F with the possibility that it hides an unresolved $`J^\pi `$ = 3/2<sup>+</sup> level as suggested by Smotrich et al. (1961). A recent gamma ray spectroscopy study (Butt et al. (1998)) of <sup>19</sup>F produced by the <sup>15</sup>N($`\alpha ,\gamma )^{19}`$F reaction has confirmed the presence of a 7.101 MeV, 3/2<sup>+</sup>, $`\mathrm{\Gamma }`$ = $`28\pm 1`$ keV level. Hence the high width found by Coszach et al. (1995) and Utku et al. (1998) for the $`E_x,J^\pi `$ = 7.067 MeV, 3/2<sup>+</sup>, <sup>19</sup>Ne level can be understood. (Note however that this analog assignment has been questioned very recently by Fortune & Sherr (2000).)
Utku et al. (1998) have also found three new <sup>19</sup>Ne levels and have made tentative assignments of analog levels in <sup>19</sup>F (see Fig.3). They recognized the 6.437 MeV, 1/2<sup>-</sup> level ($`E_r`$=26 keV) as the analog of the 6.429, 1/2<sup>-</sup>, $`\mathrm{\Gamma }`$ = 280 keV one in <sup>19</sup>F from its large measured width ($`216\pm 19`$ keV). Even though it is located at a very low energy (in the context of novae) it is so broad that its tail may lead to a significant contribution to the rate (see Fig.4). The 6.449 MeV 3/2<sup>+</sup> ($`E_r`$=38 keV) level is one of the three new ones found by Utku et al. (1998). It is also of great importance since the contribution of its tail can dominate the astrophysical factor in the relevant energy range (see Fig.4). In comparison the two other new levels give a smaller (6.698 MeV 5/2<sup>+</sup>, $`E_r`$=287 keV) or negligible contribution (6.419 MeV 3/2<sup>+</sup>, $`E_r`$=8 keV).
A new rate has been provided by Utku et al. (1998), using their new data together with available spectroscopic information on <sup>19</sup>Ne and <sup>19</sup>F supplemented by estimates of radiative and proton widths when missing. This rate is already higher by a factor of up to $`3\times 10^2`$ when compared with Wiescher & Kettner (1982). Nevertheless, even though this rate is now set on a firmer basis than the Wiescher & Kettner (1982) one, it remains quite uncertain as it depends directly on the estimated proton widths. Except for the two resonances whose strengths have been measured directly by Graulich et al. (1997), the proton widths are assumed to be equal to 0.1 or 0.01 Wigner limit, i.e. $`\theta _p^2`$=0.1 or 0.01, according to their parity (Utku et al. (1998)). For the resonances considered here, one has $`\mathrm{\Gamma }_p\mathrm{\Gamma }_{Tot}\mathrm{\Gamma }_\alpha `$ so that their strengths are ($`\omega \gamma \omega \mathrm{\Gamma }_p`$) proportional to chosen $`\theta _p^2`$ values. Accordingly, to obtain lower and upper limits for the rate, we used the two extreme cases $`\theta _p^2`$=0 and $`\theta _p^2`$=1 respectively and regardless of parity but kept Utku et al. (1998) $`\theta _p^2`$ values for the nominal rate (see Fig.6). The lower limit is quite certain as it is given by the contribution of the two directly measured resonances whose parameters are now reliably determined. It remains higher than the Wiescher & Kettner (1982) rate above $``$50$`\times 10^6`$ K. The nominal rate is conditioned by the rather arbitrary choice of $`\theta _p^2`$=0.1 or 0.01 that we adopt following Utku et al. (1998). However, our nominal rate is significantly higher \[up to a factor of $``$10 and 30 around 10<sup>8</sup> K for (p,$`\alpha `$) and (p,$`\gamma `$) reactions\] than the one from Utku et al. (1998) for typical nova temperatures (see Fig.6). This is due to the inclusion of the contribution of the tail ($`\mathrm{\Gamma }_{Tot}=\mathrm{\Gamma }_\alpha `$=4.3 keV) of the 6.449 MeV level (see Fig. 4) not considered by Utku et al. (1998). The corresponding rates (p,$`\alpha `$) are presented in Fig.6 and have been calculated by numerical integration of the Breit-Wigner formula for all broad resonances. For the (p,$`\gamma `$) rates, we used the same procedure to obtain the low, nominal and high rates (see Figs. 5 and 7). Radiative widths were taken from <sup>19</sup>F analog levels as in Utku et al. (1998) except for the important resonance at 659 keV as Bardayan et al. (1999) provided an experimental value ($`\mathrm{\Gamma }_\gamma `$=0.39 eV) for the analog level. We also adopted the direct capture (DC) contribution from Utku et al. (1998) even though the spectroscopic factors used (<sup>18</sup>O$``$p instead of <sup>18</sup>F$``$p or even <sup>18</sup>F$``$n) may not be the more appropriates. These rates have been calculated using all the available spectroscopic data collected in Table II of Utku et al. (1998) except for the inclusion of a new gamma width for the 659 keV resonance and for a different proton width for the 26 keV resonance \[1% of the Wigner limits i.e. 2.5$`\times 10^{17}`$ eV, as for other negative parity resonances, instead of 6.6$`\times 10^{20}`$ eV of unknown origin in Utku et al. (1998)\]. These numbers are taken at face value to calculate the nominal rate only. For the low and high rates unknown proton widths are allowed to vary between 0 and the Wigner limit as discussed above. It is clear that this range ($`0<\theta _p^2<1`$) represents the most extreme values as are the corresponding rates. This should be kept in mind when interpreting the astrophysical implications.
However a few points have been neglected due to the lack of experimental or theoretical information. First of all, the identification of all $`E_x`$ $`<`$ 6.6 MeV analog levels is not complete yet (see Fig. 3) and is not certain. In particular, new contributions to the nominal and high rates corresponding to missing levels cannot be ruled out. The tail of the 38 keV resonance has been calculated using the large total width derived from the <sup>19</sup>F analog level proposed by Utku et al. (1998). Another 3/2<sup>+</sup> level, with smaller width, lies 30 keV below in <sup>19</sup>F. It has been assigned (Utku et al. (1998)) to be the analog level corresponding to the 8 keV (E<sub>x</sub>=6.419 MeV) resonance in <sup>18</sup>F+p (see Fig. 3). It is possible that the analog assignment for these two levels is inversed or that they are mixed but at least one or the other resonance (8 or 38 keV) should be broad enough to dominate the rates at low nova temperature. For the (p,$`\gamma `$) rates, we neglected any difference of radiative widths between analog levels and the uncertainty on the DC contribution in front of the uncertainty on the proton widths. These rates do not either include possible interference effects arising from the interference between the 3/2<sup>+</sup> resonances located at 38 and 659 keV. In a favourable case (when $`\theta _p^2`$0.01 for the 38 keV resonance), destructive interference could reduce the low rate by a factor of $``$10 in the sensitive $`T10^8`$ K region. However, at present the rates provided in the appendix are those we recommend.
## 4 Other reactions affecting <sup>18</sup>F production
According to the recent compilation of charged-particle induced thermonuclear reaction rates NACRE (Angulo et al. 1999), the rate for <sup>16</sup>O(p,$`\gamma `$)<sup>17</sup>F is rather well known and suffers from a little uncertainty (i.e., a factor of $``$ 2 between the high and low rates). Even though <sup>17</sup>F has a much shorter lifetime than <sup>18</sup>F, its destruction by proton capture is rather well known at nova temperatures and should not affect the analysis of Sect. 2. Below $``$4$`\times 10^8`$ K, the <sup>17</sup>F(p,$`\gamma )^{18}`$Ne reaction rate is dominated by the effect of direct capture on bound <sup>18</sup>Ne levels. This contribution is expected to suffer little uncertainty since it has been calculated (Wiescher & Kettner (1982); García et al. (1991)) using experimental neutron spectroscopic factors from the mirror nucleus <sup>18</sup>O. The lowest resonances are located at $`E_r`$ = 0.586, 0.655 and 0.600 MeV (Bardayan et al. (1999)) too high to contribute, even by their tails, to the rate below $``$4$`\times 10^8`$ K. Partial widths are reasonably known through experiments, shell model calculations or from the conjugate levels (Wiescher & Kettner (1982); García et al. (1991)) and are such that $`\mathrm{\Gamma }_p\mathrm{\Gamma }_\gamma `$ (i.e. $`\omega \gamma =\omega \mathrm{\Gamma }_\gamma `$). A long standing uncertainty (Wiescher & Kettner (1982); Wiescher, Görres & Thielemann (1988); García et al. (1991)) concerned the location of the $`\mathrm{}`$=0, 0.600 MeV resonance associated with the 4.524 MeV, 3<sup>+</sup> level. Fortunately, the resonance energy first measured by García et al. (1991), has been experimentally confirmed very recently by Bardayan et al. (1999). Hence the uncertainty associated with the <sup>17</sup>F(p,$`\gamma )^{18}`$Ne rate appears negligible within the domain of nova nucleosynthesis.
According to NACRE (Angulo et al. (1999)), the <sup>17</sup>O(p,$`\alpha )^{14}`$N (Fig. 8) and <sup>17</sup>O(p,$`\gamma )^{18}`$F (Fig. 9) reaction rates present large uncertainties at temperatures below a few 10<sup>8</sup> K. The increase in the rates around 50$`\times 10^6`$ (Angulo et al. (1999)) is due to the contribution of a $`E_r`$ = 66 keV resonance (Landré et al. (1989); Blackmon et al. (1995)), not included in Caughlan & Fowler (1988). The large uncertainty around 2$`\times 10^8`$ K (Angulo et al. (1999)), well within the range of temperature reached in nova outbursts, is due to the unknown contribution of a 179.5 keV resonance associated with the E<sub>x</sub> = 5786 keV (<sup>18</sup>F) level. The upper limits for the strengths \[(p,$`\gamma `$) and (p,$`\alpha `$)\] of this resonance are calculated (Angulo et al. (1999)) from the partial widths extracted by Rolfs, Berka & Azuma (1973) and the upper limit for the proton width from Landré et al. (1989). As usual, for the calculation of the contribution of this resonance to the NACRE low, recommended and high rates, its strength has been taken equal to 0, 10% and 100% of the experimental upper limit. Accordingly, the NACRE recommended rate is somewhat arbitrary in the temperature domain affected by the 179.5 keV resonance (Figs. 8 and Figs. 9)
## 5 Results: effect of new rates on <sup>18</sup>F synthesis
New models of a 1.25 M ONe nova have been computed with the SHIVA code, in order to analyze the effect of the new nuclear reaction rates. For the <sup>18</sup>F+p rates, we have adopted the low, high and nominal prescriptions described in section 3, whereas for the <sup>17</sup>O+p ones we have adopted the NACRE rates (see Table LABEL:t:yields). It is important to stress that the energetics of the explosion remains practically unchanged with the new <sup>18</sup>F+p and <sup>17</sup>O+p reaction rates, since these reactions are not the ones affecting most the evolution (see José & Hernanz (1998) and José et al. (1999) for a deep analysis of the reaction rates which have the largest influence in nova explosions). Therefore, the maximum temperature attained in the burning shell at the base of the accreted envelope, T<sub>max</sub>=2.51$`\times 10^8`$ K, the mean kinetic energy and the ejected mass, $`\mathrm{\Delta }`$$`M_{ejec}`$=1.8$`\times 10^5`$M , are unchanged with respect to previous models (i.e., José et al. (1999)).
Two previously computed evolutionary sequences, with the same white dwarf mass and composition, are adopted as old models for the purpose of comparison: a case with Wiescher & Kettner (1982) rates (see José & Hernanz (1998) for the yields and Gómez-Gomar et al. (1998) for the $`\gamma `$-ray spectra) and a case with Utku et al. (1998) rates (see Hernanz et al. (1999) for the yields and the spectra). In both of them, the Caughlan & Fowler (1988) updated by Landré et al. (1989) rates for <sup>17</sup>O+p were adopted. In Table LABEL:t:yields we show the <sup>18</sup>F yields obtained with the new rates, as well as those with the old models. We see that yields with the nominal <sup>18</sup>F+p rates are smaller by a factor of 30 than those with the older rates, as expected. In fact, the reduction in <sup>18</sup>F production is even larger (by a factor of 3) than the one obtained with Utku et al. (1998) rates (see Hernanz et al. (1999)), as a consequence of the inclusion of the tail of the resonance at 38 keV (see section 3).
Concerning the effect of the new <sup>17</sup>O+p rates, a reduction by an extra factor of 2 is obtained when the NACRE recommended rates, instead of the old Caughlan & Fowler (1988) and Landré et al. (1989) ones, are adopted (together with <sup>18</sup>F+p nominal rates; see Table LABEL:t:yields). In summary, our nominal yield of <sup>18</sup>F is 4.84$`\times 10^5`$, which is 60 times smaller than the yield with the old Wiescher & Kettner (1982) rates (José & Hernanz (1998); Gómez-Gomar et al. (1998)) and 6 times smaller than that with Utku et al. (1998) rates (Hernanz et al. (1999)). The consequences of the reduced <sup>18</sup>F yields are quite large for the early gamma-ray output, which is directly related to the amount of <sup>18</sup>F synthesized. The reduced content of <sup>18</sup>F in an expanding envelope with similar physical and dynamical properties, translates directly into a reduced positron annihilation gamma-ray flux, roughly by the same factor as the <sup>18</sup>F decrease. This was shown in Hernanz et al. (1999), where complete computations of the gamma-ray spectra were performed with the <sup>18</sup>F yields obtained with Utku et al. (1998) rates, for different masses and compositions of the white dwarf. With the rates presented here, the reduction of the predicted fluxes becomes larger by a factor of 6, for the combined nominal <sup>18</sup>F+p and recommended <sup>17</sup>O+p rates. In summary, if we compare with results in Gómez-Gomar et al. (1998), where <sup>18</sup>F yields were computed with Wiescher & Kettner (1982), the 511 keV and continuum fluxes for a 1.25 M ONe nova should be reduced by a factor of 60.
Another important point to stress is the impact of the uncertainty in the rates on the computed yields. The above nominal yields rely upon recommended or nominal rates which contain assumed values ($`\theta _p^2`$=0.01 or 0.1 for <sup>18</sup>F+p and 10% of an experimental upper limit for <sup>17</sup>O+p). Hence, the nominal yield should not be dissociated from the large remaining uncertainties. From our results (see Table LABEL:t:yields) we conclude that the range between high and low presciptions for the <sup>18</sup>F+p rates translates into an uncertainty (maximum versus minimum yield ratio) of a factor of 310, a quite large value, whereas for the <sup>17</sup>O+p rates the uncertainty is a factor of 10. (It is worth noticing that the results obtained with the Utku et al. (1998) rates fall within the limits obtained with the new rates; see Table LABEL:t:yields). The corresponding uncertainties in the gamma-ray fluxes are of the same order of magnitude. This points out the interest of more accurate determinations of the <sup>18</sup>F+p and <sup>17</sup>O+p rates.
## 6 Conclusions
We have investigated the <sup>18</sup>F formation and destruction in nova outbursts, identified the key reactions (proton capture on <sup>18</sup>F and <sup>17</sup>O) and analysed their rates. The proton capture rates on <sup>18</sup>F are higher than the Wiescher & Kettner (1982) ones at nova temperatures due to the contribution of the tail of the 659 keV resonance whose large measured width (Coszach et al. (1995); Graulich et al. (1997)) has been indirectly confirmed (Butt et al. (1998)). Another important contribution comes from the tail of the 38 keV resonance which was neglected in previous studies. Its nominal contribution is larger than the 659 keV one but is proportional to its assumed reduced proton width $`\theta _p^2`$. The strengths (proton widths) of the low lying resonances are unknown and induce large uncertainties (factors of 100 to 1000) in the rates at nova temperatures. We have provided updated nominal rates for the two capture reactions together with upper and lower limits. The <sup>17</sup>O+p rates also display some large uncertainties at nova temperatures according to the recent compilation of Angulo et al. (1999).
We have used these new nuclear physics results in a fully hydrodynamical nova code to calculate the <sup>18</sup>F yields in novae for different rates : our low, nominal and high <sup>18</sup>F+p rates, and the low, recommended and high <sup>17</sup>O+p rates from Angulo et al. (1999). These results have been compared with models computed with the old Wiescher & Kettner (1982) rates (José & Hernanz (1998)), and with more recent models with the Utku et al. (1998) rates (Hernanz et al. (1999)). Two important results have been obtained. First, there is always a reduction of the amount of <sup>18</sup>F synthesized in a nova explosion, with the nominal rates for the <sup>18</sup>F+p reactions both alone and combined with the recommended rates for the <sup>17</sup>O+p reactions. The nominal <sup>18</sup>F yield (nominal <sup>18</sup>F+p and recommended for <sup>17</sup>O+p) is 4.84$`\times 10^5`$ by mass, which is 60 times smaller than the one obtained with Wiescher & Kettner (1982) rates and 6 times smaller than the one with Utku et al. (1998) rates (and Caughlan & Fowler (1988) and Landré et al. (1989) for <sup>17</sup>O+p). The impact on the early gamma-ray spectrum of the nova is a reduction of the flux by the same amount (with respect to Gómez-Gomar et al. (1998) and Hernanz et al. (1999), respectively). Second, the yields are found to be very sensitive to the rates with resulting combined (<sup>18</sup>F+p and <sup>17</sup>O+p) uncertainties of more than three orders of magnitude. This supports the need of new experimental and theoretical studies to improve the knowledge of the <sup>18</sup>F+p and <sup>17</sup>O+p rates and, consequently, allow for a larger reliability of the predictions of annihilation gamma–ray fluxes from novae, to be observed by current and future instruments.
###### Acknowledgements.
This work was partially supported by PICS 319, PB98-1183-C02, PB98-1183-C03 and ESP98-1348.
Appendix: Updated rates
We give here the tabulated reaction rates for proton capture on <sup>18</sup>F resulting from numerical integrations (see § 3 and Figs. 6 and 7). For convenience, we also provide here formulas that approximate the nominal rates (As usual, T9 is the temperature in GK, T9XY=T9\**(X/Y) and T9LN=LOG(T9) in FORTRAN notations.)
The following <sup>18</sup>F(p,$`\alpha )^{15}`$O nominal rate formula includes the contribution of the two measured resonances (330 and 659 keV) with additional contributions from the possible resonances at 38 keV, (with high energy tail) and 287 keV.
```
SVDir = 9.13e10/T923*EXP(-18.052/T913-0.672*T92)*(1.+! Tail
& 0.0231*T913+6.12*T923+0.988*T9+9.92*T943+4.07*T953) ! 659 keV
& +5.78e5/T932*EXP(-3.830/T9)+9.91e8/T932*EXP(-7.648/T9)+ ! 330, 659 keV
& 2.81e-6/T932*EXP(-0.441/T9)+2706.*EXP(3.8319*T9LN-1.3450* ! 38 keV
& T9LN**2-.0001/T9**3)+4.46e4/T932*EXP(-3.331/T9) ! 287 keV
```
The following <sup>18</sup>F(p,$`\gamma )^{19}`$Ne nominal rate formula contains contributions from direct capture and from the 38, 287, 330 and 659 keV resonances.
```
SVDir = 3.98e7/T923*EXP(-18.052/T913)* ! DC
& (1.+0.0231*T913+0.0885*T923+0.0143*T9)+
& 1.34e3/T932*EXP(-3.830/T9)+1.51e4/T932*EXP(-7.648/T9)
& +EXP(-0.56781+3.6850*T9LN-1.3636*T9LN**2-.0001/T9**3)
& +8.26e-10/T932*EXP(-0.441/T9)+12.2/T932*EXP(-3.331/T9)
```
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# Phase changes in 38-atom Lennard-Jones clusters. II: A parallel tempering study of equilibrium and dynamic properties in the molecular dynamics and microcanonical ensembles
## I Introduction
The simulation of systems having complex potential energy surfaces (PES) is often difficult owing to the problem of quasi-ergodicity. Quasi-ergodicity can arise in systems having several energy minima separated by high energy barriers. When such situations occur, as for example in proteins, glasses or clusters, the system can become trapped in local basins of the energy landscape, and the ergodic hypothesis fails on the time scale of the simulation. In the canonical ensemble, the high energy regions of the PES are exponentially suppressed and barrier crossings become rare events. Calculations of equilibrium properties when phase space is thus partitioned require methods that overcome quasi-ergodicity by enhanced barrier crossing. Many techniques have been proposed to address this problem, including the use of generalized ensembles such as multicanonical or Tsallisian, simulated tempering, configurational or force bias Monte Carlo, or various versions of the jump-walking algorithm. Most of these techniques have been introduced for Monte Carlo (MC) simulations rather than molecular dynamics (MD) simulations. These techniques have been applied to a variety of sampling and optimization problems, and phase changes in clusters have often been considered as a benchmark to test these methods.
The double-funnel energy landscape of the 38-atom Lennard-Jones (LJ) cluster has been investigated in detail by Doye, Miller and Wales, who recently also estimated the inter-funnel rate constants using master equation dynamics. This landscape is challenging for simulation because of the high free-energy barrier separating the two funnels. In the preceding paper (hereafter referred to as I), we have shown how the parallel tempering algorithm can be used to deal with this particularly complex system for Monte Carlo simulations in the canonical ensemble. Achieving ergodicity in microcanonical simulations is much harder than in the canonical ensemble, because the system is unable to cross any energy barrier higher than the total energy available. The 38-atom Lennard-Jones cluster is fundamentally non-ergodic in a range of energies. This non-ergodicity may not be a serious problem when considering one particular cluster on a short time scale. However, in a statistical sample of such systems it is important to have ergodic results.
To allow MD simulations to cross the high energy barriers, one may think of heating the system (by increasing its kinetic energy), followed by a cooling back to its initial thermodynamic state. Although this process is straightforward, the dynamics becomes biased and non physical during the heating and cooling processes. Moreover, it is difficult to control accurately the heating and cooling rates that should be chosen for any system. This latter aspect is particularly critical for the 38-atom Lennard-Jones cluster where the narrow and deepest funnel is hard to reach even at high temperatures.
Because of the inherent difficulties of molecular dynamics, MC approaches can be especially useful for dealing with the problem of crossing high energy barriers. Monte Carlo methods have been developed in previous work in the microcanonical ensemble. In these approaches the microcanonical sampling is at fixed energy without any additional constraints. Such methods can be contrasted with isoenergetic molecular dynamics where the total, linear and angular momenta are also conserved. These additional constraints must be considered even at zero angular momentum. To differentiate microcanonical simulations, where only the energy is fixed, from molecular dynamics simulations, where additional constraints are imposed, we identify the former to be simulations in the microcanonical ensemble and identify the latter simulations to be in the molecular dynamics ensemble. The differences in the two ensembles can be particularly important when the angular momentum is large enough to induce structural (centrifugal) distortions. Because dynamical properties are calculated using molecular dynamics methods, in this work we find that a combination of Monte Carlo and molecular dynamics methods are most convenient for developing ergodic approaches to dynamics.
In this paper, we adapt the parallel tempering method to both the microcanonical and molecular dynamics ensembles. The application of parallel tempering in the molecular dynamics ensemble requires the incorporation of the conservation of the total linear and angular momenta into the probabilities. In order to extract ergodic dynamical properties, we combine Monte Carlo methods with molecular dynamics to develop a hybrid ergodic MC/MD algorithm. The efficiency of the simulation tools developed in this work is demonstrated by applications to the 38-atom Lennard-Jones cluster, which exhibits a solid-solid transition prior to melting. This transition to an equilibrium phase between truncated octahedral and icosahedral geometries makes this cluster an ideal candidate for investigating how the ergodic hypothesis can influence the dynamical behavior of a complex system.
The contents of the remainder of this paper are as follows. In the next section, we recall the basic principle of Monte Carlo sampling in the microcanonical ensemble, and we present the simple modifications needed to include parallel tempering. We test microcanonical parallel tempering methods on the 38-atom Lennard-Jones cluster, and compare the microcanonical results with those found in I using the canonical ensemble. We focus on equilibrium properties, including the caloric curve and the isomers distributions. In section III we review the characteristics of the molecular dynamics ensemble at fixed total linear and angular momenta and fixed total energy. We extend the parallel tempering Monte Carlo method to the MD ensemble, and we combine microcanonical parallel tempering with molecular dynamics to produce an ergodic MD method. We also apply these methods to several dynamical properties of LJ<sub>38</sub>; in particular the diffusion constant and the largest Lyapunov exponent. We summarize our findings and discuss our results in section IV.
## II Parallel tempering Monte Carlo in the microcanonical ensemble
The fundamental quantity in the microcanonical ensemble is the density of states $`\mathrm{\Omega }`$. For a system having $`N`$ identical particles, volume $`V`$ and total energy $`E`$, $`\mathrm{\Omega }`$ is defined by
$$\mathrm{\Omega }(N,V,E)=\frac{1}{N!h^{3N}}\delta [H(𝐫,𝐩)E]d^{3N}rd^{3N}p$$
(1)
where $`h`$ is Planck’s constant and where $`H(𝐫,𝐩)`$ denotes the classical Hamiltonian function of the coordinates $`𝐫`$ and momenta $`𝐩`$ of the $`N`$ particles. Knowing the microcanonical density of states $`\mathrm{\Omega }`$, one can calculate the canonical partition function $`Q(N,V,T)`$ by a Laplace transformation. The kinetic part of the Hamiltonian $`H`$ is quadratic in the momenta, and Eq. (1) can be partly integrated to give
$$\mathrm{\Omega }(N,V,E)=\left(\frac{2\pi m}{h^2}\right)^{3N/2}\frac{1}{N!\mathrm{\Gamma }(3N/2)}\mathrm{\Theta }[EU(𝐫)][EU(𝐫)]^{3N/21}d^{3N}r.$$
(2)
In Eq.(2), $`\mathrm{\Gamma }`$ is the Gamma function, $`m`$ is the mass of each particle, $`U(𝐫)=HK`$ is the potential energy and $`\mathrm{\Theta }`$ is the Heaviside step function: $`\mathrm{\Theta }(x)=1`$ if $`x0`$, 0 otherwise. Microcanonical averages of a coordinate-dependent variable $`A(𝐫)`$ can be expressed
$$A(N,V,E)=\frac{{\displaystyle \mathrm{\Theta }[EU(𝐫)][EU(𝐫)]^{3N/21}A(𝐫)d^{3N}r}}{{\displaystyle \mathrm{\Theta }[EU(𝐫)][EU(𝐫)]^{3N/21}d^{3N}r}}.$$
(3)
The microcanonical entropy $`S`$ can be defined by $`S(N,V,E)=k_B\mathrm{ln}\mathrm{\Omega }(N,V,E)`$ with $`k_B`$ the Boltzmann constant. The thermodynamic temperature $`T(N,V,E)`$ is given by the thermodynamic relation $`(S/E)_{N,V}=1/T`$, and can be obtained from a microcanonical average
$$T(N,V,E)=\frac{2}{3N2}\frac{1}{K^1}.$$
(4)
This expression is slightly different from the kinetic temperature $`2K/3N`$, which is a consequence of our choice in the definition of the entropy. As discussed by Pearson and co-workers, it is also possible to define the entropy using the phase space volume
$$\mathrm{\Phi }(N,V,E)=_0^E\mathrm{\Omega }(N,V,E^{})𝑑E^{}.$$
(5)
Definitions of the temperature based on $`\mathrm{\Omega }`$ differ from the temperature based on $`\mathrm{\Phi }`$ to order $`1/N`$, and the two definitions agree only in the thermodynamic limit.
Monte Carlo simulations can be used to explore the microcanonical ensemble by performing a random walk in configuration space. In the standard Metropolis scheme, a trial move from configuration $`𝐫_o`$ to configuration $`𝐫_n`$ is accepted with the probability
$$\mathrm{a}cc(𝐫_o𝐫_n)=\mathrm{min}(1,\frac{\rho _E(𝐫_n)T(𝐫_n𝐫_o)}{\rho _E(𝐫_o)T(𝐫_o𝐫_n)}),$$
(6)
where $`T(𝐫_o𝐫_n)`$ is a trial probability. The acceptance probability expressed in Eq.(6) insures detailed balance so that the random walk visits points in configuration space proportional to the equilibrium distribution $`\rho _E(𝐫)`$ defined by
$$\rho _E(𝐫)=\zeta ^1\mathrm{\Theta }[EU(𝐫)][EU(𝐫)]^{3N/21}$$
(7)
where $`\zeta `$ is the normalization. In practice, $`T(𝐫_o𝐫_n)`$ is a uniform distribution of points of width $`\mathrm{\Delta }`$ centered about $`𝐫_o`$, and $`\mathrm{\Delta }`$ is adjusted as a function of the energy so that not too many trial moves are either accepted or rejected.
Implementation of microcanonical Monte Carlo is as easy as its canonical version. Because Monte Carlo methods are based on random walks in configuration space, in principle the system can cross energy barriers higher than the available energy. However, in difficult cases like LJ<sub>38</sub>, large atomic displacements having poor acceptance ratios are needed to reach ergodicity.
Parallel tempering has proved to be an important approach to insure ergodicity in canonical Monte Carlo simulations, and parallel tempering can be easily adapted to the microcanonical ensemble by replacing the Boltzmann factors in the acceptance probability by the microcanonical weight $`\rho _E(𝐫)`$. In the parallel tempering scheme, several microcanonical MC simulations are performed simultaneously at different total energies $`\{E_i\}`$. With some predetermined probability, two simulations at energies $`E_i`$ and $`E_j`$ attempt to exchange their current configurations, respectively $`𝐫_i`$ and $`𝐫_j`$, and this exchange is accepted with probability
$`\mathrm{min}(1,{\displaystyle \frac{\rho _{E_i}(𝐫_j)\rho _{E_j}(𝐫_i)}{\rho _{E_i}(𝐫_i)\rho _{E_j}(𝐫_j)}}).`$
The acceptance ratio is analogous to the canonical expression given in I. In microcanonical simulations the potential energies must be smaller than $`\mathrm{min}(E_i,E_j)`$; otherwise the move is rejected. Parallel tempering microcanonical MC works in the same way as in standard canonical MC. As with canonical parallel tempering MC, the gaps between adjacent total energies must be chosen to be small enough so that exchanges between the corresponding trajectories are accepted with a reasonable probability.
By using a histogram reweighting technique, it is possible to extract from the MC simulations the density of states $`\mathrm{\Omega }`$, and then all the thermodynamic quantities in both the microcanonical and the canonical ensembles. The procedure is similar to that described in Ref. , and relies on the calculation of the distribution $`P(U,E)`$ of potential energy $`U`$ at the total energy $`E`$. $`P`$ is fitted to the microcanonical form $`P(U,E)=\mathrm{\Omega }_C(U)(EU)^{3N/21}/\mathrm{\Omega }(E)`$, where $`\mathrm{\Omega }_C`$ stands for the configurational density of states, and $`\mathrm{\Omega }(E)`$ is extracted by convolution of $`\mathrm{\Omega }_C(U)`$ and $`(EU)^{3N/21}`$.
We have tested the parallel tempering Monte Carlo algorithm in the microcanonical ensemble on the 38-atom Lennard-Jones cluster previously investigated. Forty different total energies ranging from $`172.4737\epsilon `$ to $`124\epsilon `$ have been used, and the same simulation conditions have been chosen as in I. In addition to a constraining sphere of radius $`2.25\sigma `$ to prevent evaporation, exchanges have been attempted every 10 passes, with the same method for choosing exchanging trajectories as described in the previous article. The simulations are begun with random configurations of the cluster geometry, and consist of $`1.3\times 10^{10}`$ points accumulated following equilibration moves consisting of $`95\times 10^6`$ Metropolis points(no exchanges) followed by $`190\times 10^6`$ points using parallel tempering. The microcanonical heat capacity calculated in this fashion and shown in Fig. 1, is qualitatively the same as the canonical heat capacity \[see I\]. The melting peak in the microcanonical heat capacity occurs at the same calcuated temperature as the temperature of the melting peak in the canonical heat capacity, and there are slope change regions at temperatures that correspond to equilibrium between the icosahedral basin and the truncated octahedral global minimum structure. The present simulations are also used to obtain structural insight about the cluster as a function of total energy. We have calculated the order parameter $`Q_4`$ as defined in I as a function of temperature, and compared the classification into the three categories of isomers (truncated octahedral, icosahedral or liquid-like) using the energy criterion also outlined in I.
In Fig. 2 we show the caloric curve $`T(E)`$ determined from our parallel tempering microcanonical MC simulations. We also present the canonical curve for comparison. The melting transition near $`T0.166\epsilon /k_B`$ is reflected in the change in slope of the temperature as a function of the energy. The microcanonical curve does not display a van der Waals loop, and remains very close to the canonical curve. The average value of the order parameter $`Q_4`$ is displayed in the lower panel of Fig. 2 as a function of the total energy. As has been discussed in I for the canonical simulation, the order parameter begins to drop at energies where there is the onset of isomerization transitions to the icosahedral basin (near $`E=160\epsilon `$), and the order parameter reaches its lowest value at the melting transition. The isomer distributions have been evaluated either using the parameter $`Q_4`$ or using the energy criterion (see the discussion in paper I). The results have been plotted in Fig. 3 as a function of the total energy. The behavior of the isomer distributions as a function of energy is similar to the canonical distributions as a function of temperature, and the cluster exhibits equilibrium between truncated octahedral and icosahedral geometries in the energy range $`160\epsilon E150\epsilon `$, prior to the solid-like to liquid-like phase change. As in the canonical case, the icosahedral distribution is a symmetric function of the energy when the energy criterion is used rather than the definition based on $`Q_4`$. This difference reflects the differences between the two definitions of icosahedral and liquid basins. The oscillatory structure observed at the peak of $`P_{Q_4}`$ for the icosahedral distribution in the upper panel of Fig. 3 is smaller than the calculated errors (two standard deviations of the mean are shown). Whether the observed structure would persist for a longer simulation is not known to us. Because the definition that assigns configurations to the icosahedral basin is arbitrary, we have chosen not to investigate this structure further.
It is useful to contrast the current results with previous constant energy studies of LJ<sub>38</sub>. Previous simulations have used molecular dynamics methods where no attempt has been made to insure ergodicity. To contrast these past studies with the molecular dynamics technique discussed in the next section of this paper, we define standard molecular dynamics to represent the usual molecular dynamics method where no special procedure is introduced to insure ergodicity. Simulations of LJ<sub>38</sub> using standard molecular dynamics invariably lead to a caloric curve with a clear van der Waals loop and a melting temperature higher than that inferred from Fig. 2. From the results of Ref. , the cluster is trapped in the octahedral basin, and the system does reflect the true dynamical coexistence state between the truncated octahedron and the icosahedral basin. This is the common situation encountered in MD simulations of the LJ<sub>38</sub> system; the cluster chooses either to remain trapped in the octahedral basin, or to escape and coexist between the icosahedral solid-like and liquid-like forms. Because the system is unable to return from the octahedral basin, the microcanonical temperature decreases. In the usual case, van der Waals loops arise when there are large energy gaps between the lowest-energy isomers. In the specific case of LJ<sub>38</sub>, it appears that the presence of extra (icosahedral) isomers only slightly higher in energy than the octahedral structure eliminates this loop in the ergodic microcanonical caloric curve.
In order to extract dynamical quantities, the Monte Carlo method we have presented must be modified to sample the MD ensemble. The modification is the subject of the next section.
## III Ergodic molecular dynamics
The molecular dynamics ensemble differs from the microcanonical ensemble in that two quantities are conserved in addition to the total energy $`E`$, volume $`V`$ and number of particles $`N`$. These two quantities are the total linear momentum $`𝐏`$ and total angular momentum $`𝐋`$. If their values are prescribed, the density of states remains the fundamental quantity of interest, and is now defined by
$`\mathrm{\Omega }(N,V,E,𝐏,𝐋)=`$ (8)
$`{\displaystyle \frac{1}{N!h^{3N}}}{\displaystyle \delta [H(𝐫,𝐩)E]\delta \left(𝐏\underset{i=1}{\overset{N}{}}𝐩_i\right)\delta \left(𝐋\underset{i=1}{\overset{N}{}}𝐫_i\times 𝐩_i\right)d^{3N}rd^{3N}p}.`$ (9)
As is the case in the microcanonical ensemble \[see Eq.(2)\], for atomic systems the momentum integrations in Eq.(9) can be evaluated explicitly. Because the thermodynamic properties are not affected by the translational motion of the center of mass, we can assume that $`𝐏=0`$. We then obtain
$`\mathrm{\Omega }(N,V,E,𝐏=0,𝐋)=`$ (10)
$`\left({\displaystyle \frac{2\pi m}{h^2}}\right)^{3N/23}{\displaystyle \frac{1}{N!\mathrm{\Gamma }(3N/23)}}{\displaystyle \mathrm{\Theta }[EU_𝐋(𝐫)][EU_𝐋(𝐫)]^{3N/24}\frac{d^{3N}r}{\sqrt{\mathrm{d}et𝐈}}},`$ (11)
where $`𝐈`$ is the inertia matrix and $`U_𝐋(𝐫)=U(𝐫)+𝐋^{}𝐈^1𝐋/2`$ is the effective rovibrational energy. This effective potential energy includes the kinetic energy contribution of the rotating cluster considered as a rigid body. The landscape of rotating clusters has been investigated by Miller and Wales in order to study cluster evaporation. Averages in the MD ensemble are now expressed as
$$A=\frac{{\displaystyle \mathrm{\Theta }[EU_𝐋(𝐫)][EU_𝐋(𝐫)]^{3N/24}A(𝐫)\frac{d^{3N}r}{\sqrt{\mathrm{d}et𝐈}}}}{{\displaystyle \mathrm{\Theta }[EU_𝐋(𝐫)][EU_𝐋(𝐫)]^{3N/24}\frac{d^{3N}r}{\sqrt{\mathrm{d}et𝐈}}}}.$$
(12)
As in the microcanonical ensemble, we define the entropy in the molecular dynamics ensemble by $`S=k_B\mathrm{ln}\mathrm{\Omega }`$. The differences between the microcanonical and molecular dynamics ensembles are the exponent $`3N/2`$ which is reduced by 3 owing to the geometrical constraints, the potential energy which now includes the contribution of the centrifugal energy, and the weight $`1/\sqrt{\mathrm{d}et𝐈}`$ which is a consequence of the conservation of the angular momentum. Monte Carlo simulations can sample the MD ensemble by performing a random walk in configuration space. The acceptance probability from configuration $`𝐫_o`$ to configuration $`𝐫_n`$ is
$$\mathrm{a}cc(𝐫_o𝐫_n)=\mathrm{min}(1,\frac{\rho _{E,𝐋}(𝐫_n)T(𝐫_n𝐫_o)}{\rho _{E,𝐋}(𝐫_o)T(𝐫_o𝐫_n)})$$
(13)
in the Metropolis scheme. The microcanonical weight $`\rho _E(𝐫)`$ is now replaced by the MD weight $`\rho _{E,𝐋}`$ given by
$$\rho _{E,𝐋}(𝐫)=\zeta ^1\frac{1}{\sqrt{\mathrm{d}et𝐈}}\mathrm{\Theta }[EU_𝐋(𝐫)][EU_𝐋(𝐫)]^{3N/24},$$
(14)
where $`\zeta `$ is a normalization. The expression for the acceptance probability is similar to Eq. (6), and a practical implementation of Monte Carlo in the MD ensemble is made in the same way as in the true microcanonical ensemble, given the vector $`𝐋`$. Parallel tempering can be also easily combined with the MC simulations. The acceptance probability of exchanging the configurations $`𝐫_i`$ and $`𝐫_j`$ initially at the total energies $`E_i`$ and $`E_j`$ respectively is then
$`\mathrm{min}(1,\left({\displaystyle \frac{[E_iU_𝐋(𝐫_j)][E_jU_𝐋(𝐫_i)]}{[E_iU_𝐋(𝐫_i)][E_jU_𝐋(𝐫_j)]}}\right)^{3N/24})`$
provided that all quantities inside brackets are positive (otherwise the move is rejected). It is remarkable that the geometrical weights have canceled in this expression.
The Monte Carlo method we have just described allows sampling of configuration space rigorously equivalent to the sampling we would obtain using molecular dynamics, but with the additional possibility of crossing the energy barriers higher than the available energy. The method can be used in its present form to extract equilibrium properties only dependent on the energy or geometry, as has been illustrated in the previous section. To compute dynamical quantities, the method can also provide a database of configurations representative of a given total energy. Instead of performing a few very long MD simulations that are in principle unable to reach other parts of the energy surface separated by barriers higher than the available energy, we choose to perform a statistical number of short simulations starting from configurations obtained by parallel tempering Monte Carlo in the MD ensemble with same total energy and angular momentum. By construction, if the MC method is correctly ergodic, then the hybrid MD method we have suggested can be expected to yield ergodic dynamical observables.
We now illustrate this ergodic molecular dynamics method on the LJ<sub>38</sub> problem. Two essentially dynamical parameters have been calculated. The first is the self diffusion constant $`D`$, obtained from the derivative of the average mean square atomic displacement
$$D=\frac{1}{6}\frac{d}{dt}[𝐫(t)𝐫(0)]^2,$$
(15)
where the average is taken over all particles of the system and over all short MD simulations. The other parameter is the largest Lyapunov exponent $`\lambda `$, that measures the exponential rate of divergence of the distance between two initially close trajectories in the phase space. If we write the equation describing the Hamiltonian dynamics in condensed form as $`\dot{\psi }(t)=F(\psi )`$ where $`F`$ is a nonlinear function and $`\psi =\{𝐫,𝐩\}`$ the phase space point, then a small perturbation $`\delta \psi `$ evolves according to the simple equation $`d\delta \psi /dt=(F/\psi )\delta \psi `$. The largest Lyapunov exponent $`\lambda `$ is obtained from the time evolution of the vector $`\delta \psi `$:
$$\lambda =\underset{t\mathrm{}}{lim}\underset{\delta \psi (0)0}{lim}\frac{1}{t}\mathrm{ln}\frac{\delta \psi (t)}{\delta \psi (0)}.$$
(16)
In Eq.(16), $``$ is a metric on the phase space. In principle, any metric can be used, and we choose the Euclidian metric including both the momenta and the coordinates. The numerical procedure involves a periodic renormalization of the vector $`\delta \psi `$ to prevent its exponential divergence. The successive lengths are accumulated and contribute to the average value of $`\lambda `$.
In I, the clusters have been defined using a hard sphere constraining potential. Because the angular momentum is not conserved after reflection from such hard wall boundaries, in the molecular dynamics simulations we have chosen a soft repulsive spherical wall $`U_c`$ defined with respect to the center of mass of the cluster for each particle by
$$U_c(𝐫)=\{\begin{array}{cc}0,\hfill & r<R_c\hfill \\ \kappa (rR_c)^4/4,\hfill & rR_c\text{.}\hfill \end{array}$$
(17)
In this equation, the atomic distances $`r`$ are measured with respect to the cluster center of mass. The simulations have been performed setting the angular momentum to zero for simplicity. We stress that even in this case (with $`𝐋=0`$), the weight $`1/\sqrt{\mathrm{d}et𝐈}`$ must be included in the Monte Carlo probabilities so that we effectively sample the MD ensemble. The actual thermodynamic behavior in the MD ensemble at zero angular momentum is nevertheless nearly identical to the microcanonical behavior.
The application to the LJ<sub>38</sub> cluster has been made by performing $`10^{10}`$ MC steps following $`10^7`$ equilibration steps in a parallel tempering simulation in the MD ensemble. The same 40 total energies have been chosen as in the previous section, and $`10^5`$ configurations have been stored every $`10^5`$ steps for each simulation. Short molecular dynamics runs of $`10^4`$ time steps following $`10^3`$ equilibration steps have been performed for each of these configurations, with the same total energy as the corresponding MC trajectory of origin, and with zero total linear and angular momenta as well. The parameters used for the constraining wall are respectively $`R_c=2.25\sigma `$ and $`\kappa =100\epsilon `$, for both the MC and MD runs. A simple Verlet algorithm has been used to propagate the MD trajectory with the time step $`\delta t=0.01`$ reduced LJ units. The propagation of the tangent trajectory to calculate the Lyapunov exponent has been determined with a fourth order Runge-Kutta scheme. The final values of $`D`$ and $`\lambda `$ are an average over the $`10^5`$ MD simulations. The variations of $`D`$ and $`\lambda `$ with total energy are depicted in Fig. 4. In both cases, two curves have been plotted, calculated either from standard molecular dynamics (with $`10^8`$ time steps following $`10^7`$ equilibration steps, and starting initially from the lowest-energy structure), or from our hybrid ergodic molecular dynamics method. For both quantities, the two MD schemes clearly yield distinct values in the energy range where equilibrium between truncated octahedral and icosahedral geometries occurs. The thermodynamic temperature, not plotted here, has the same variations as the caloric curve of Fig. 2 when calculated with ergodic MD. Standard molecular dynamics predicts a van der Waals loop centered at $`T0.18\epsilon /k_B`$. For standard MD, the cluster is trapped in the icosahedral basin and is, in practice, unable to reach the octahedral basin. Only the equilibrium between the icosahedral basin and liquid-like structures occurs. As can be seen from the upper panel of Fig. 4, this change in curvature of the temperature is also present for the diffusion constant, which exhibits strong variations at the energy where the octahedral structure vanishes when standard MD is used. In contrast, the variations in ergodic MD are smooth.
The melting temperature implied by the largest Lyapunov exponent is also higher in standard MD than in ergodic MD, even though the variations of the Lyapunov exponent are continuous in both MD schemes. Indeed, using ergodic molecular dynamics we observe a shift of the curve obtained by standard MD toward the lower energies. As shown by Hinde, Berry, and Wales, the Lyapunov exponent and the Kolmogorov entropy are quantities essentially dependent on the local properties of the energy landscapes. One contribution comes from the negative curvature of the landscape, and another contribution is the fluctuation of positive curvature. Both contributions are affected by the cluster being trapped either inside the truncated octahedral basin or inside the icosahedral basin. In this latter case in particular, the different isomers belonging to the icosahedral basin are connected through regions of negative curvature, while only one isomer defines the octahedral funnel.
Because ergodic molecular dynamics allows the cluster to be found in both basins prior to melting, the dynamical behavior is likely to be very different (and more chaotic) with respect to the dynamical behavior of the cluster confined to the octahedral funnel. This difference is precisely what we observe on the lower panel of Fig. 4.
## IV Conclusion
In this paper, we have explored the parallel tempering method in simulations in the microcanonical ensemble. The implementation of the parallel tempering algorithm in this ensemble is straightforward, the Boltzmann factor $`\mathrm{exp}(\beta U)`$ being replaced by the microcanonical weight $`(EU)^{3N/21}`$. Application to the LJ<sub>38</sub> cluster has shown the thermodynamic behavior in the microcanonical ensemble to be similar to the behavior in the canonical ensemble. The solid-liquid phase change is preceded by a solid-solid phase change where the cluster is in equilibrium between truncated octahedral and icosahedral geometries. This phase equilibrium is well reproduced in the simulations owing to the power of parallel tempering. The calculated microcanonical caloric curve, which does not display a van der Waals loop, is consistent with the single peaked heat capacity observed in I.
We have extended the parallel tempering microcanonical Monte Carlo algorithm to sample the molecular dynamics ensemble at constant total energy, linear momentum and angular momentum. Combined with standard molecular dynamics, this method circumvents the lack of connectivity between regions of the potential energy surface. The method can ensure ergodicity in microcanonical simulations, which is much more difficult to achieve than in the canonical ensemble. Ironically, this ergodic MD method can be viewed as the counterpart of the techniques developed by Chekmarev and Krivov to study the dynamics of systems confined to only one catchment basin in the energy surface.
We have performed a statistical number of short molecular dynamics runs starting from configurations stored periodically in parallel tempering Monte Carlo simulations. These simulations sample the MD ensemble at the same total energies, linear and angular momenta as the standard molecular dynamics runs. In fact, the length of the MD runs is mainly dictated by the large number of starting configurations. One may think of reducing drastically this number, to allow for the calculation of parameters varying on longer time scales. Unfortunately, if the energy landscape is not known in advance, then it is hard to guess how important are the contributions of the basins not selected as starting configurations. In the case of LJ<sub>38</sub> having only 3 main regions on the energy surface, one possibility is to compute a dynamical property as the average value over 3 different simulations starting either from the truncated octahedral geometry, one icosahedral geometry or a low-lying liquid geometry, all carried out at the same total energy. However, as we have seen in Fig. 3, it is not obvious how to choose properly the weights of each basin in this average because of the difficulty in distinguishing between icosahedral and liquid structures in many cases. For this reason, we believe that the first parallel tempering MC step of the hybrid ergodic method is essential in the vicinity of phase changes to capture many starting configurations that are used subsequently in standard molecular dynamics. The enhanced sampling offered by parallel tempering can also act as a statistical representation of the energy surface at a given total energy, and the long time dynamics may be further investigated by using master equations after searching the saddle points.
We have calculated two dynamical quantities with the present hybrid MD/MC method, the diffusion constant and the largest Lyapunov exponent in the 38-atom Lennard-Jones cluster. The variations of both quantities with the total energy are significantly different when evaluated with standard (non-ergodic) molecular dynamics or with our hybrid ergodic MD method. These results emphasize the different contributions of the two funnels of the energy landscape to the average value of the parameters estimated.
The algorithms developed in this investigation allow the calculation of thermodynamic, structural, or dynamical properties of systems such as LJ<sub>38</sub> that can be expressed as phase space or time averages. Parallel tempering works using a criterion based on the potential energy but not on the geometry. Consequently permutational isomers can be introduced in the course of the simulation. Quantities such as fluctuations of configuration-dependent properties are much more difficult to extract than actual averages. For instance, the Lindemann index $`\delta `$, which measures the root mean square bond length fluctuation, is often considered to be a reliable parameter for detecting melting in atomic and molecular systems. This quantity cannot be properly estimated with the ergodic MD scheme, and the same difficulty persists for other methods based on the use of different trajectories.
Although the idea of combining Monte Carlo sampling with standard molecular dynamics can be applied to other techniques such as jump-walking, we believe that parallel tempering is the key to the success in the case of LJ<sub>38</sub>. As in the canonical version, the equilibrium phase between truncated octahedral and icosahedral structures is correctly reproduced in an energy range preceding the melting region, because in this range configurations may be accessed either from higher energy trajectories containing mainly icosahedral geometries, or from lower energy trajectories acting as a reservoir for the octahedral geometry. As noticed by Falcioni and Deem, the parallel tempering algorithm is especially useful at low temperatures, or in our case, at low energy. The long relaxation times inherent in systems like clusters, proteins, critical or glassy liquids, are a serious difficulty for standard simulation methods. We expect the present ergodic method to be particularly useful to deal with the dynamics of such systems.
The method we have presented works at constant total energy. It is possible to improve ergodicity in constant-temperature MD either by using canonical parallel tempering as in the work of Sugita and Okamoto, or by coupling parallel tempering canonical Monte Carlo to short Nosé-Hoover trajectories. In the Nosé-Hoover approach such molecular dynamics simulations do conserve a zero angular momentum, so a rigorous MC sampling should include the geometrical weight $`1/\sqrt{\mathrm{d}et𝐈}`$ in the probabilities also in this case. The present microcanonical scheme can be easily used for rotating bodies, which makes the method suitable for investigating the strong influence of centrifugal effects on phase changes in atomic clusters.
## Acknowledgments
Some of this work has been motivated by the attendance of two of us (DLF and FC) at a recent CECAM meeting on ‘Overcoming broken ergodicity in simulations of condensed matter systems.’ We would like to thank CECAM, J.E. Straub and B. Smit who organized the meeting, and those who attended the workshop for stimulating discussions. This work has been supported in part by the National Science Foundation under grant numbers CHE-9714970 and CDA-9724347. This research has been supported in part by the Phillips Laboratory, Air Force Material Command, USAF, through the use of the MHPCC under cooperative agreement number F29601-93-0001. The views and conclusions contained in this document are those of the authors and should not be interpreted as necessarily representing the official policies or endorsements, either expressed or implied, of Phillips Laboratory or the US Government.
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# 1 Introduction
## 1 Introduction
By now a considerable amount of evidence has been accumulated for the Matrix theory for M-theory, originally proposed by Banks, Fischler, Shenker and Susskind and later re-interpreted by Susskind in the framework of discrete light-cone quantization. In particular, just to mention only the direct comparison with eleven dimensional supergravity, complete agreement for the multi-graviton scattering (including the recoil effects) at 2-loop and that for the two-body potential between arbitrary fermionic as well as bosonic objects at 1-loop can be cited as highly non-trivial and remarkable.
Despite such impressive pieces of evidence as well as general supportive arguments , the deep reason and the mechanism of the agreement are yet to be fully understood. Evidently, one of the keys should be the understanding of the structure and the role of the symmetries. The well-known symmetries present both in the Matrix theory and in the supergravity theory are the global $`Spin(9)`$ invariance, the CPT invariance and the invariance under 16 supersymmetries. Perhaps less familiar is the generalized conformal invariance, the generalization of the conformal symmetry that plays the major role in the AdS/CFT correspondence of Maldacena. Besides these symmetries, the supergravity possesses the general coordinate invariance while the Matrix theory has the Yang-Mills type gauge symmetry, which must be deeply connected. Finally, although not yet clearly identified, the agreement of the multi-body scattering amplitudes strongly suggests that the eleven dimensional Lorentz invariance is present in a highly non-trivial manner in the Matrix theory.
Except for the eleven dimensional Lorentz invariance, the above-mentioned symmetries of the Matrix theory are easily recognized in the original action. However, since the supergravity interactions between various objects arise only after introducing the corresponding backgrounds and integrating out the quantum fluctuations around them, the realizations of these symmetries are in general modified for the effective action of interest. Besides, being an off-shell quantity, the form of the effective action depends on the choice of the gauge as well as on the definition of the fields. Further, there is always a vast degree of ambiguities as one can add total derivatives. For these reasons, the study of how the symmetries govern the structure of the effective action becomes quite non-trivial.
In this article, we shall focus on the supersymmetry (SUSY), considered to be the most powerful among the ones listed above. In fact it has been claimed that a large degree of SUSY, $`𝒩=16`$, present in the theory imposes strong restrictions on the form of the effective action and leads to a number of “non-renormalization theorems”. However, upon close examinations one finds that such assertions in the existing literature can still be challenged. This is essentially due to the lack of completely off-shell consideration. As we shall elaborate in some detail in Sec. 3, to understand precisely to what extent SUSY is responsible in determining the effective action, one must allow the background fields to have arbitrary time-dependence. This in turn inevitably leads to the necessity of examining the gauge (BRST) symmetry, another important ingredient of Matrix theory, as SUSY and gauge symmetry are known to be intimately intertwined off-shell. Such an analysis has not been performed in the past.
What makes the off-shell analysis difficult is that we do not as yet have the off-shell unconstrained superfield formulation in the case of $`𝒩=16`$ supersymmetry. This prompts us to resort to the conventional means, namely the Ward identity for the symmetries in question. In order to be able to check against the available explicit result for the effective action, one would like to obtain the Ward identity in so-called the “background gauge”, in which all the calculations have been performed. Although the derivation of such a Ward identity is expected to be a text-book matter, this was not to be: One must carefully disentangle the dependence on the background field of the effective action by making use of BRST Ward identities. The result is a somewhat complicated Ward identity, in which the supersymmetry and the BRST symmetry are intertwined. A notable feature of our Ward identity is that one can read off the effective quantum-corrected SUSY transformations in closed form and this should serve as a starting point of various truly off-shell investigations. As a simple application, we compute the transformation laws to the lowest order in the derivative expansion at 1-loop and analyze the restriction imposed by SUSY on the effective action for a background with arbitrary time-dependence, to the corresponding order. We find that at this order the effective action is indeed fully determined by the requirement of supersymmetry, which agrees with the explicit calculation including the normalization.
The rest of the paper is organized as follows: In Sec. 2, we recall the symmetries of the original action of Matrix theory and the BRST symmetry associated with the background gauge fixing. Then in Sec. 3 we make some important remarks on the determination of the effective action and its symmetries, to emphasize the necessity of completely off-shell analysis. Having clarified the issue, we proceed in Sec. 4 to the derivation of the SUSY Ward identity for the effective action in the background gauge. By carefully separating the different origins of the dependence on the background field, we obtain the desired Ward identity together with the closed expressions for the effective SUSY transformation laws. This constitutes the main result of this work. As an application, we explicitly work out in Sec. 5 the SUSY transformation laws to the lowest order in the derivative expansion at 1-loop and show that the effective action to the corresponding order is completely determined by the Ward identity. Sec. 6 is devoted to a short summary and discussions of future problems.
## 2 Action and its Symmetries
Let us begin by recalling the action of the Matrix theory and its symmetries, which at the same time serves to set our notations.
The basic action of the Matrix theory can be written as
$`S_0`$ $`=`$ $`\text{Tr}{\displaystyle }dt\{{\displaystyle \frac{1}{2}}[D_t,X_m]^2+{\displaystyle \frac{g^2}{4}}[X_m,X_n]^2`$ (2.1)
$`+{\displaystyle \frac{i}{2}}\mathrm{\Theta }^T[D_t,\mathrm{\Theta }]+{\displaystyle \frac{g}{2}}\mathrm{\Theta }^T\gamma ^m[X_m,\mathrm{\Theta }]\},`$
$`D_t`$ $`=`$ $`_tigA,`$ (2.2)
where $`N\times N`$ hermitian matrices $`X_{ij}^m(t),A_{ij}(t),\mathrm{\Theta }_{\alpha ,ij}(t)`$ stand for the bosonic, the gauge, and the fermionic fields respectively. The middle Latin indices $`m,n,\mathrm{}`$, running from $`1`$ to $`9`$, denote the spatial directions, while the Greek letters such as $`\alpha ,\beta =116`$ are used for the $`SO(9)`$ spinor indices. The $`16\times 16`$ $`\gamma `$-matrices $`\gamma ^m`$ are real symmetric and satisfy $`\{\gamma ^m,\gamma ^n\}=\delta ^{mn}`$.
To facilitate the quantum computations, it is convenient to define the theory by going to the “Euclidean formulation”. Introduce the Euclidean time $`\tau `$, the gauge field $`\stackrel{~}{A}`$, and the action $`\stackrel{~}{S}_0`$ by<sup>1</sup><sup>1</sup>1Fermions are not transformed.
$`\tau `$ $``$ $`it,\stackrel{~}{A}iA,\stackrel{~}{S}_0iS_0.`$ (2.3)
Then the action and the covariant derivative become
$`\stackrel{~}{S}_0`$ $`=`$ $`\text{Tr}{\displaystyle }d\tau \{{\displaystyle \frac{1}{2}}[D_\tau ,X_m]^2{\displaystyle \frac{g^2}{4}}[X_m,X_n]^2`$ (2.4)
$`+{\displaystyle \frac{1}{2}}\mathrm{\Theta }^T[D_\tau ,\mathrm{\Theta }]{\displaystyle \frac{1}{2}}g\mathrm{\Theta }^T\gamma ^m[X_m,\mathrm{\Theta }]\},`$
$`D_\tau `$ $`=`$ $`_\tau ig\stackrel{~}{A}.`$ (2.5)
Besides the obvious $`Spin(9)`$ symmetry, this action is invariant under the following transformations:
1. Gauge transformations with a gauge parameter matrix $`\mathrm{\Lambda }`$:
$`\delta _\mathrm{\Lambda }\stackrel{~}{A}`$ $`=`$ $`[D_\tau ,\mathrm{\Lambda }],\delta _\mathrm{\Lambda }X_m=ig[\mathrm{\Lambda },X_m],\delta _\mathrm{\Lambda }\mathrm{\Theta }=ig[\mathrm{\Lambda },\mathrm{\Theta }].`$ (2.6)
2. Supersymmetry transformations with a spinor parameter $`ϵ_\alpha `$:
$`\delta _ϵ\stackrel{~}{A}`$ $`=`$ $`ϵ^T\mathrm{\Theta },\delta _ϵX^m=iϵ^T\gamma ^m\mathrm{\Theta },`$ (2.7)
$`\delta _ϵ\mathrm{\Theta }`$ $`=`$ $`i\left([D_\tau ,X_m]\gamma ^m+{\displaystyle \frac{g}{2}}[X_m,X_n]\gamma ^{mn}\right)ϵ`$ (2.8)
where $`\gamma ^{mn}\frac{1}{2}[\gamma ^m,\gamma ^n]`$ is real anti-symmetric. When the system possesses symmetries other than the supersymmetry, the supersymmetry algebra may close not only on the space-time translation but also on the generators of such additional symmetries. In fact in the present case it is well-known that it involves the gauge symmetry with a field-dependent gauge function. For example on $`\stackrel{~}{A}`$,
$`[\delta _ϵ,\delta _\lambda ]\stackrel{~}{A}`$ $`=`$ $`2iϵ^T\lambda _\tau \stackrel{~}{A}+\delta _\mathrm{\Lambda }\stackrel{~}{A},`$ (2.9)
$`\text{where}\mathrm{\Lambda }`$ $`=`$ $`2i(ϵ^T\gamma ^m\lambda X_mϵ^T\lambda \stackrel{~}{A}),`$ (2.10)
and similarly for the other fields. Moreover, for a system with $`𝒩=16`$ supersymmetry, such as the Matrix theory, formulation in terms of unconstrained superfields is not known and hence the algebra closes only up to the equations of motion in general. Thus it is expected that the proper understanding of the supersymmetry of the Matrix theory must include the analysis of these non-trivial features.
3. Generalized conformal transformations: If we rescale the fields $`X_m`$ and $`\stackrel{~}{A}`$ by a factor of $`g`$, such as $`X_mX_m/g`$, and allow $`g`$ to depend on $`\tau `$ to a linear order, $`\stackrel{~}{S}_0`$ is invariant under a generalization of the conformal transformations . In particular, the invariance under the special conformal transformation defined by
$`\delta _K\stackrel{~}{A}`$ $`=`$ $`2ϵ\tau \stackrel{~}{A},\delta _KX_m=2ϵ\tau X_m,\delta _K\mathrm{\Theta }_\alpha =0`$
$`\delta _K\tau `$ $`=`$ $`ϵ\tau ^2,\delta _Kg=3ϵ\tau g,`$ (2.11)
$`ϵ`$ $`=`$ $`\text{an infinitesimal bosonic parameter},`$
imposes a useful restriction on the form of the effective action.
Besides these well-established symmetries, the remarkable agreement between the 11-dimensional supergravity calculations and the 2-loop Matrix theory calculations for multi-body scattering processes strongly suggests that the Matrix theory actually possesses 11-dimensional Lorentz symmetry in a highly non-trivial manner.
In this article, we shall focus on how the first two of these symmetries, which are intimately intertwined, are implemented in the quantum effective action of the supergravitons. In the M-theory interpretation of the Matrix theory, the coordinates and the spin degrees of freedom of these supergravitons are represented by the diagonal backgrounds for $`X^m`$ and $`\mathrm{\Theta }_\alpha `$ respectively. We shall denote them by $`B_m`$ and $`\theta _\alpha `$ respectively and separate them from the quantum parts $`Y_m`$ and $`\mathrm{\Psi }_\alpha `$ as
$`X_m`$ $`=`$ $`{\displaystyle \frac{1}{g}}B_m+Y_m,`$ (2.12)
$`\mathrm{\Theta }_\alpha `$ $`=`$ $`{\displaystyle \frac{1}{g}}\theta _\alpha +\mathrm{\Psi }_\alpha .`$ (2.13)
As was already emphasized in the introduction and will be further elaborated in the next section, it is important to take $`B_m(\tau )`$ and $`\theta _\alpha (\tau )`$ as arbitrary backgrounds, not satisfying any equations of motion. Only in this way we can unambiguously determine how much restrictions are imposed by the supersymmetry on the effective action for these background fields.
To quantize the theory, we need to fix the gauge. Although our derivation of the Ward identity, to be presented in Sec. 4, can be readily adapted to any choice of gauge, the actual computations are extremely cumbersome except in the standard background gauge. It is specified by the gauge-fixing function of the form
$`G`$ $`=`$ $`_\tau \stackrel{~}{A}+i[B^m,X_m].`$ (2.14)
In fact essentially all the existing explicit calculations have been performed in this gauge. However, as it will become clear, the naive use of this gauge leads to a subtle but important complication in deriving the correct Ward identity. To avoid this problem, we will tentatively use a different function $`\stackrel{~}{B}_m`$ in place of $`B_m`$ and write the gauge-fixing function as
$`\stackrel{~}{G}`$ $`=`$ $`_\tau \stackrel{~}{A}+i[\stackrel{~}{B}^m,X_m].`$ (2.15)
Later at an appropriate stage, we will set $`\stackrel{~}{B}_m=B_m`$.
The corresponding ghost action can be readily obtained by the standard BRST method. The BRST transformations for the quantum part of the fields are given by
$`\delta _B\stackrel{~}{A}`$ $`=`$ $`[D_\tau ,C],\delta _BY_m=ig[X_m,C],`$
$`\delta _B\mathrm{\Psi }`$ $`=`$ $`ig\{C,\mathrm{\Theta }\},`$ (2.16)
$`\delta _BC`$ $`=`$ $`igC^2,\delta _B\overline{C}=ib,\delta _Bb=0.`$
$`C`$, $`\overline{C}`$ and $`b`$ are, respectively, the ghost, the anti-ghost and the Nakanishi-Lautrup auxiliary fields. The background fields are not transformed. Then the combined gauge-ghost action $`\stackrel{~}{S}_{gg}`$ is generated by
$`\stackrel{~}{S}_{gg}`$ $`=`$ $`\delta _B\text{Tr}{\displaystyle 𝑑\tau \left[\frac{1}{i}\overline{C}\left(\stackrel{~}{G}\frac{\alpha }{2}b\right)\right]}.`$ (2.17)
We will henceforth set the gauge parameter $`\alpha `$ to be 1. This leads, after integrating over the $`b`$ field, to the familiar gauge-ghost action
$`\stackrel{~}{S}_{gg}`$ $`=`$ $`\text{Tr}{\displaystyle 𝑑\tau \frac{1}{2}\left(_\tau \stackrel{~}{A}+i[\stackrel{~}{B}^m,X_m]\right)^2}`$ (2.18)
$`i\text{Tr}{\displaystyle 𝑑\tau \left(\overline{C}_\tau [D_\tau ,C]g\overline{C}[\stackrel{~}{B}^m,[X_m,C]]\right)}.`$
## 3 Remarks on the Determination of Effective Action and its Symmetries
Before starting the derivation of the off-shell supersymmetry Ward identity, we wish to make some important remarks on the determination of the effective action and its symmetries, which point to the necessity of off-shell analysis. Although many of the remarks will apply for general backgrounds, for clarity of discussions we shall consider the so-called “source-probe configuration”. In the case of $`U(N+1)`$ gauge group, it is defined as the situation where a probe supergraviton interacts with $`N`$ supergravitons situated at the origin that act as a heavy source. The background fields representing this situation are
$`B_m`$ $`=`$ $`\mathrm{diag}(r_m,0,0,\mathrm{},0),\theta _\alpha =\mathrm{diag}(\theta _\alpha ,0,0,\mathrm{},0).`$ (3.19)
Here, $`r_m`$ represents the coordinate of the probe and $`\theta _\alpha `$ its spin content. As usual, the spin of the source is neglected.
As already emphasized in the introduction, the primary feature of Matrix theory is that it generates the supergravity interactions among various objects only after (i) introducing the corresponding backgrounds and (ii) integrating over the quantum fluctuations around them. Because of this, realizations of the symmetries of the effective theory are in general modified non-trivially. Besides, being an off-shell quantity, the form of the effective action is affected by (A) the gauge choice, (B) the (re)definition of the fields, and (C) the freedom of adding total derivatives. Of course the on-shell S-matrix elements do not depend on these factors. However, the determination of the full ( i.e. quantum-corrected) on-shell condition itself requires the knowledge of the off-shell effective action<sup>2</sup><sup>2</sup>2This was clearly demonstrated in , where the agreement between the supergravity and the Matrix theory calculation was achieved with the recoil corrections. Thus, in order to understand the symmetry structure of the effective theory fully, it is necessary to perform an off-shell analysis with (A) $``$ (C) properly taken into account.
Now since an exact analysis is practically impossible, one often needs to make some approximations. In doing so, one must make sure that they are logically consistent for one’s aim. For the present purpose, some of the often used approximations are not appropriate. For example, the eikonal approximation, where one tries to reconstruct the effective action from the eikonal phase shift, can be dangerous and misleading. In fact the answer depends on the form of the effective Lagrangian assumed. As a simple illustration, consider the 1-loop eikonal phase shift for $`v1`$ given by
$`\stackrel{~}{\mathrm{\Gamma }}_1^e`$ $`=`$ $`{\displaystyle \frac{v^3}{b^6}}+0\times {\displaystyle \frac{v^5}{b^{10}}}{\displaystyle \frac{3}{2}}{\displaystyle \frac{v^7}{b^{14}}}+𝒪(v^9)`$
$`r_m`$ $`=`$ $`v_m\tau +b_m,vb=0,b_m=\text{impact parameter}.`$
If one assumes the effective Lagrangian to be of the form $`=(v,r)`$, then the effective action that reproduces this phase shift is uniquely determined to be
$`\stackrel{~}{S}_1`$ $`=`$ $`{\displaystyle 𝑑\tau \left(\frac{15}{16}\frac{v^4}{r^7}+0\times \frac{v^6}{r^{11}}\frac{9009}{4096}\frac{v^8}{r^{15}}+𝒪(v^{10})\right)}.`$
However, even restricting to $`𝒪(v^4)`$, the most general form allowed for the effective action contains 6 independent structures, after eliminating total-derivative ambiguities:
$`\stackrel{~}{S}_1`$ $`=`$ $`{\displaystyle }d\tau (A{\displaystyle \frac{v^4}{r^7}}+B{\displaystyle \frac{v^2(vr)^2}{r^9}}+C{\displaystyle \frac{(vr)^4}{r^{11}}}`$
$`+D{\displaystyle \frac{v^2(ar)}{r^7}}+E{\displaystyle \frac{(ar)^2}{r^7}}+F{\displaystyle \frac{a^2}{r^5}}).`$
On the other hand, there is only one condition, $`(15/16)=A+(1/7)B+(1/21)C`$, required for the correct phase shift, and hence 5 parameters remain undetermined. The situation at $`𝒪(v^6)`$ is even more striking. Although the eikonal phase shift vanishes, the explicit computation reveals that there are 5 non-vanishing independent structures present in the effective action.
It should be clear from these illustrations that the only logically consistent procedure, not affected by the total-derivative ambiguities, is to use off-shell backgrounds with arbitrary $`\tau `$-dependence<sup>3</sup><sup>3</sup>3This was emphasized in the context of generalized conformal symmetry in . Related discussion can also be found in . and to classify terms by derivative expansion according to the “order” defined by
order $`=`$ $`\text{\# of }_\tau +{\displaystyle \frac{1}{2}}\text{\# of fermions}.`$ (3.20)
If necessary, one may combine this with the usual loop expansion.
Having emphasized the importance of off-shell considerations, we now make some related comments on the general arguments on the restrictions imposed by supersymmetry, often referred to as SUSY non-renormalization theorems. They can be roughly classified into two categories.
The first type of argument, devised by Paban et al , relies on the closure property of SUSY transformations. For example, at $`𝒪(v^2)`$ they first make a choice of the definition of the fields so that the action takes the form $`{\displaystyle 𝑑\tau f(r)v^2}`$ and take the $`𝒪(v^0)`$ SUSY transformation laws in that basis to be the standard ones without any correction, $`\delta _ϵr^m=iϵ\gamma ^m\theta ,\delta _ϵ\theta _\alpha =i(v/ϵ)_\alpha `$. Then demanding that the closure is canonical, namely $`[\delta _ϵ,\delta _\lambda ]=2\lambda ^Tϵ_\tau `$, they show that there cannot be a correction to the transformation laws at $`𝒪(v^2)`$ and hence the $`𝒪(v^2)`$ effective action is tree-exact. Although the argument is quite simple and plausible, it is unclear why the closure should be canonical off-shell and further it is not obvious if the $`𝒪(v^0)`$ transformation laws must be of the standard form in a particular basis adopted. In general, field redefinitions affect the form of the SUSY transformations and hence they must be considered as a pair.
The other type of argument is known as the SUSY completion method, which makes use of the chain of relations produced among terms with different number of $`\theta `$’s by SUSY transformations. For example, at $`𝒪(v^4)`$ one expects relations of the form
$`v^4v^3\theta ^2v^2\theta ^4v\theta ^6\theta ^8.`$ (3.21)
By showing that the top form, $`\theta ^8`$ term in this case, is not renormalized beyond 1-loop, one wishes to infer the non-renormalization of all the other terms in the chain, in particular the $`v^4`$ term. This method appears efficient, but some care is needed in drawing firm conclusions. One problem is that sometimes the chain starting from the top form stops at an intermediate stage. Put differently, one may form a super-invariant not containing the top form. An example already occurs at $`𝒪(v^2)`$, where the tree-level expression $`{\displaystyle 𝑑\tau \left((v^2/2g^2)+(\theta \dot{\theta }/2g^2)\right)}`$ is SUSY-complete without a $`\theta ^4`$ term. More non-trivial example is seen at $`𝒪(v^6)`$: Although $`\theta ^{12}`$ term was shown to vanish at 1-loop, the bosonic contribution at $`𝒪(v^6)`$ nonetheless exists.
An attempt at filling this gap was made in . In this work all the connections in the chain (3.21) were examined and it was concluded that SUSY is indeed powerful enough to fix the effective action at this order up to an overall constant. Again one must be cautious in accepting this conclusion: In this analysis $`v`$ and $`\theta `$ were taken to be $`\tau `$-independent and hence the assumed form of the effective Lagrangian was not the most general one allowed in the proper derivative expansion with arbitrary backgrounds. Later it was recognized, however, that the higher derivative terms neglected in this analysis can actually be absorbed into the tree-level Lagrangian by a suitable field re-definition, which appeared to resurrect the validity of the analysis made in . Unfortunately, the problem still persists: By such a field re-definition the higher derivatives are simply shifted into the SUSY transformation laws and one must reanalyze the issue with such modifications.
Thus one sees that although the existing analyses are highly plausible they are not air-tight. In view of the importance of precise understanding of the role of supersymmetry and its connection with gauge symmetry, it is desirable to perform an unambiguous off-shell analysis with arbitrary backgrounds. This motivates us to the study of the Ward identity, to be described in the next two sections.
## 4 SUSY Ward Identity for the Effective Action in the Background Gauge
Having argued the importance of off-shell analysis for arbitrary trajectories, we shall now derive the SUSY Ward identity for the effective action $`\stackrel{~}{\mathrm{\Gamma }}`$ in the standard background gauge, (2.14), used exclusively in the actual computations.
To make use of the well-established method, let us further split the quantum fluctuations $`Y_m`$ and $`\mathrm{\Psi }_\alpha `$ into two parts, the diagonal and the off-diagonal, in the manner
$`Y_{m,ij}={\displaystyle \frac{\widehat{y}_{m,i}}{g}}\delta _{ij}+\widehat{Y}_{m,ij},\mathrm{\Psi }_{\alpha ,ij}={\displaystyle \frac{\widehat{\psi }_{\alpha ,i}}{g}}\delta _{ij}+\widehat{\mathrm{\Psi }}_{\alpha ,ij},`$ (4.22)
and introduce the sources only for the diagonal fields:
$`\stackrel{~}{S}_s={\displaystyle 𝑑\tau \left(J_{m,i}\widehat{y}_{m,i}+\eta _{\alpha ,i}\widehat{\psi }_{\alpha ,i}\right)}.`$ (4.23)
The Euclidean generating functionals are defined by
$`\stackrel{~}{Z}[J,\eta ]`$ $`=`$ $`{\displaystyle 𝒟\mu \mathrm{exp}\left(\stackrel{~}{S}_{tot}\right)}=\mathrm{exp}\left(\stackrel{~}{W}[J,\eta ]\right),`$ (4.24)
$`\stackrel{~}{S}_{tot}`$ $``$ $`\stackrel{~}{S}_0+\stackrel{~}{S}_{gg}+\stackrel{~}{S}_s,𝒟\mu 𝒟\stackrel{~}{A}𝒟Y𝒟\mathrm{\Psi }𝒟C𝒟\overline{C},`$
where $`\stackrel{~}{W}[J,\eta ]`$ is the one for the connected functions. By making the change of integration variables corresponding to the supersymmetry transformations, one obtains the primitive form of the Ward identity
$`0=\delta _ϵ\stackrel{~}{S}_{gg}+\delta _ϵ\stackrel{~}{S}_s.`$ (4.25)
Here and in what follows, $`𝒪`$ for an operator $`𝒪`$ means
$`𝒪={\displaystyle \frac{{\displaystyle 𝒟\mu 𝒪e^{\stackrel{~}{S}_{tot}}}}{{\displaystyle 𝒟\mu e^{\stackrel{~}{S}_{tot}}}}}.`$ (4.26)
We now rewrite this identity (4.25) in terms of the generating functional $`\stackrel{~}{\mathrm{\Gamma }}`$, which is 1PI (1-particle-irreducible) with respect to the diagonal fields. Define as usual the classical fields and $`\stackrel{~}{\mathrm{\Gamma }}`$ by
$`y_{m,i(\tau )}`$ $``$ $`{\displaystyle \frac{\delta \stackrel{~}{W}}{\delta J_{m,i}(\tau )}},\psi _{\alpha ,i}(\tau ){\displaystyle \frac{\delta \stackrel{~}{W}}{\delta \eta _{\alpha ,i}(\tau )}},`$ (4.27)
$`\stackrel{~}{\mathrm{\Gamma }}[y,\psi ]`$ $``$ $`W[J,\eta ]{\displaystyle 𝑑\tau (J_{m,i}y_{m,i}+\eta _{\alpha ,i}\psi _{\alpha ,i})}.`$ (4.28)
Then the sources are expressed in terms of $`\stackrel{~}{\mathrm{\Gamma }}`$ as
$`J_{m,i}(\tau )={\displaystyle \frac{\delta \stackrel{~}{\mathrm{\Gamma }}}{\delta y_{m,i}(\tau )}},\eta _{\alpha ,i}(\tau )={\displaystyle \frac{\delta \stackrel{~}{\mathrm{\Gamma }}}{\delta \psi _{\alpha ,i}(\tau )}}.`$ (4.29)
Therefore, the contribution to the Ward identity from the variation of the source action can be written as
$`\delta _ϵ\stackrel{~}{S}_s={\displaystyle 𝑑\tau \left(\frac{\delta \stackrel{~}{\mathrm{\Gamma }}}{\delta y_{m,i}(\tau )}\delta _ϵ\widehat{y}_{m,i}(\tau )+\frac{\delta \stackrel{~}{\mathrm{\Gamma }}}{\delta \psi _{\alpha ,i}(\tau )}\delta _ϵ\widehat{\psi }_{\alpha ,i}(\tau )\right)}.`$ (4.30)
As for the contribution $`\delta _ϵ\stackrel{~}{S}_{gg}`$ from the gauge-ghost part, a direct calculation yields a rather complicated expression, which constitutes an inhomogeneous term in the Ward identity regarded as a functional integro-differential equation for $`\stackrel{~}{\mathrm{\Gamma }}`$. This is undesirable since what we wish to understand is how the supersymmetry acts on the effective action $`\stackrel{~}{\mathrm{\Gamma }}`$. Fortunately, it was noted long ago in the context of four-dimensional super Yang-Mills theory that one can reexpress such a term in a form similar to (4.30). The first step is to note that the supersymmetry transformations (2.7) and (2.8) commute with the BRST transformation (2.16) on all the fields, as can be checked straightforwardly. Therefore, starting from (2.17), $`\delta _ϵ\stackrel{~}{S}_{gg}`$ can be written as
$`\delta _ϵ\stackrel{~}{S}_{gg}`$ $`=`$ $`\delta _ϵ\delta _B\text{Tr}{\displaystyle 𝑑\tau \left[\frac{1}{i}\overline{C}\left(\stackrel{~}{G}\frac{1}{2}b\right)\right]}`$ (4.31)
$`=`$ $`\delta _B\delta _ϵ\text{Tr}{\displaystyle 𝑑\tau \left[\frac{1}{i}\overline{C}\left(\stackrel{~}{G}\frac{1}{2}b\right)\right]}=\delta _B𝒪_ϵ,`$
where
$`𝒪_ϵ{\displaystyle \frac{1}{i}}\text{Tr}{\displaystyle 𝑑\tau \overline{C}\delta _ϵ\stackrel{~}{G}}`$ (4.32)
is a fermionic composite operator. This expression, being an expectation value of a BRST-exact form, vanishes in the ordinary vacuum. However, in the presence of external sources, it becomes proportional to the sources, and hence to the functional derivatives of $`\stackrel{~}{\mathrm{\Gamma }}`$. Let us collectively denote by $`\varphi `$ and $`J`$ the basic fields and the corresponding sources respectively and consider the generating functional with a source $`j`$ for the operator $`𝒪_ϵ`$:
$`Z[J,j]={\displaystyle 𝒟\varphi e^{(S_{tot}+J\varphi +j𝒪_ϵ)}}.`$ (4.33)
Now make a change of variables corresponding to the BRST transformation. We get
$`0={\displaystyle 𝒟\varphi ((1)^{|\varphi |}J\delta _B\varphi j\delta _B𝒪_ϵ)e^{(S_{tot}+J\varphi +j𝒪_ϵ)}},`$ (4.34)
where $`|\varphi |`$ is 0 (1) if $`\varphi `$ is bosonic (fermionic). By differentiating with respect to $`j`$ once, setting $`j=0`$, and then expressing the source $`J`$ in terms of $`\stackrel{~}{\mathrm{\Gamma }}`$, one easily obtains the following BRST Ward identity:
$`\delta _B𝒪_ϵ`$ $`=`$ $`{\displaystyle 𝑑\tau \frac{\delta \stackrel{~}{\mathrm{\Gamma }}}{\delta \varphi (\tau )}\delta _B\varphi (\tau )𝒪_ϵ}.`$ (4.35)
In this way, we get
$`\delta _ϵ\stackrel{~}{S}_{gg}={\displaystyle 𝑑\tau \left(\frac{\delta \stackrel{~}{\mathrm{\Gamma }}}{\delta y_{m,i}(\tau )}\delta _B\widehat{y}_{m,i}(\tau )𝒪_ϵ+\frac{\delta \stackrel{~}{\mathrm{\Gamma }}}{\delta \psi _{\alpha ,i}(\tau )}\delta _B\widehat{\psi }_{\alpha ,i}(\tau )𝒪_ϵ\right)}.`$ (4.36)
Putting all together, we arrive at the following SUSY Ward identity expressed solely in terms of the derivatives of $`\stackrel{~}{\mathrm{\Gamma }}`$:
$`0`$ $`=`$ $`{\displaystyle }d\tau ({\displaystyle \frac{\delta \stackrel{~}{\mathrm{\Gamma }}}{\delta y_{m,i}(\tau )}}(\delta _ϵ\widehat{y}_{m,i}(\tau )+\delta _B\widehat{y}_{m,i}(\tau )𝒪_ϵ)`$ (4.37)
$`{\displaystyle \frac{\delta \stackrel{~}{\mathrm{\Gamma }}}{\delta \psi _{\alpha ,i}(\tau )}}(\delta _ϵ\widehat{\psi }_{\alpha ,i}(\tau )\delta _B\widehat{\psi }_{\alpha ,i}(\tau )𝒪_ϵ)).`$
Normally it is now a simple matter to convert this into the desired Ward identity for the effective action as a functional of the backgrounds $`B^m`$ and $`\theta _\alpha `$: One would rewrite the derivatives with respect to $`y_m`$ and $`\psi _\alpha `$ into those with respect to $`B^m`$ and $`\theta _\alpha `$ and then set $`y_m=\psi _\alpha =0`$. This procedure is indeed valid for the fermions since $`\widehat{\psi }_\alpha `$ and $`\theta _\alpha `$ always appear in the combination $`\widehat{\psi }_\alpha +\theta _\alpha `$ in the original action and it is a simple matter to prove that this gets converted to $`\psi _\alpha +\theta _\alpha `$ in $`\stackrel{~}{\mathrm{\Gamma }}`$.
This is not so for the bosonic field. While most of the $`B^m`$ dependence comes from the splitting $`X_m=(1/g)(B_m+\widehat{y}_m)+\widehat{Y}_m`$, where $`B_m`$ and $`\widehat{y}_m`$ appear together, $`B^m`$ in the gauge-ghost sector (in the standard background gauge) is not accompanied with $`\widehat{y}_m`$. From the point of view of the Ward identity above, it is an independent extra parameter field which should be distinguished from the bona fide background field. This is why we chose to start out with a different symbol $`\stackrel{~}{B}^m`$ for this field.
Now the problem we face is the following. In order to be able to apply the Ward identity to the case of the standard background gauge, we need to express the dependence on $`\stackrel{~}{B}^m`$ again in the form of the functional derivative of $`\stackrel{~}{\mathrm{\Gamma }}`$ with respect to $`B^m`$. In other words, we must disentangle the two different types of $`B^m`$ dependence buried in the standard background gauge formulation and construct the Ward identity which takes both types of dependence into account. Although this is a technical rather than a conceptual problem, it is again a manifestation of the gauge theory nature of the Matrix theory, which has often been neglected.
The problem can be solved as follows. Since $`\stackrel{~}{B}^m`$ appears only in the gauge-ghost action and the variation with respect to it commutes with the BRST transformation, we get from (2.17)
$`{\displaystyle \frac{\delta \stackrel{~}{S}_{gg}}{\delta \stackrel{~}{B}^{m,i}(\tau )}}=\delta _B𝒪_{m,i}(\tau ),`$ (4.38)
where
$`𝒪_{m,i}={\displaystyle \underset{j}{}}\left(\overline{C}_{ji}Y_{m,ij}\overline{C}_{ij}Y_{m,ji}\right),`$ (4.39)
and $`\stackrel{~}{B}^{m,i}`$ stands for the diagonal elements $`\stackrel{~}{B}_{ii}^m`$. The expectation value of the left-hand side can be expressed in terms of the generating functional $`\stackrel{~}{W}`$ as $`\delta \stackrel{~}{W}/\delta \stackrel{~}{B}_{m,i}`$, which is equal to $`\delta \stackrel{~}{\mathrm{\Gamma }}/\delta \stackrel{~}{B}_{m,i}`$ since it is a parameter field. On the other hand the expectation value of the right-hand side can be treated in exactly the same way as we treated $`\delta _B𝒪_ϵ`$. In this way we get the identity
$`{\displaystyle \frac{\delta \stackrel{~}{\mathrm{\Gamma }}}{\delta \stackrel{~}{B}_{m,i}(\tau )}}`$ $`=`$ $`{\displaystyle }d\tau ^{}({\displaystyle \frac{\delta \stackrel{~}{\mathrm{\Gamma }}}{\delta y_{n,j}(\tau ^{})}}\delta _B\widehat{y}_{n,j}(\tau ^{})𝒪_{m,i}(\tau )`$ (4.40)
$`+{\displaystyle \frac{\delta \stackrel{~}{\mathrm{\Gamma }}}{\delta \psi _{\alpha ,j}(\tau ^{})}}\delta _B\widehat{\psi }_{\alpha ,j}(\tau ^{})𝒪_{m,i}(\tau )).`$
Now let us replace $`y_m`$ and $`\psi _\alpha `$ with $`B_m`$ and $`\theta _\alpha `$ respectively and then set $`y_m=\psi _\alpha =0`$ in the previous Ward identity (4.37) and in the relation (4.40) above. In the limit $`\stackrel{~}{B}_mB_m`$, the total variation of $`\stackrel{~}{\mathrm{\Gamma }}`$ with respect to $`B_m`$, which we denote by $`\mathrm{\Delta }\stackrel{~}{\mathrm{\Gamma }}/\mathrm{\Delta }B_m`$ to avoid confusion, is
$`{\displaystyle \frac{\mathrm{\Delta }\stackrel{~}{\mathrm{\Gamma }}}{\mathrm{\Delta }B_{m,i}(\tau )}}=\left({\displaystyle \frac{\delta \stackrel{~}{\mathrm{\Gamma }}}{\delta B_{m,i}(\tau )}}+{\displaystyle \frac{\delta \stackrel{~}{\mathrm{\Gamma }}}{\delta \stackrel{~}{B}_{m,i}(\tau )}}\right)_{\stackrel{~}{B}=B}.`$ (4.41)
Substituting (4.40) we then get
$`{\displaystyle \frac{\mathrm{\Delta }\stackrel{~}{\mathrm{\Gamma }}}{\mathrm{\Delta }B_{m,i}(\tau )}}`$ $`=`$ $`{\displaystyle 𝑑\tau ^{}\frac{\delta \stackrel{~}{\mathrm{\Gamma }}}{\delta B_{n,j}(\tau ^{})}T_{nj;mi}(\tau ,\tau ^{})}`$ (4.42)
$`{\displaystyle 𝑑\tau ^{}\frac{\delta \stackrel{~}{\mathrm{\Gamma }}}{\delta \theta _{\alpha ,j}(\tau ^{})}\delta _B\widehat{\psi }_{\alpha ,j}(\tau ^{})𝒪_{m,i}(\tau )},`$
with
$`T_{nj;mi}(\tau ,\tau ^{})`$ $``$ $`\delta _{mn}\delta _{ij}\delta (\tau \tau ^{})\delta _B\widehat{y}_{n,j}(\tau ^{})𝒪_{m,i}(\tau ).`$ (4.43)
By inverting this relation, we can express the partial variation $`\delta \stackrel{~}{\mathrm{\Gamma }}/\delta B_m`$ in terms of the total variation $`\mathrm{\Delta }\stackrel{~}{\mathrm{\Gamma }}/\mathrm{\Delta }B_m`$:
$`{\displaystyle \frac{\delta \stackrel{~}{\mathrm{\Gamma }}}{\delta B_{m,i}(\tau )}}`$ $`=`$ $`{\displaystyle 𝑑\tau ^{}T_{mi,nj}^1(\tau ,\tau ^{})\frac{\mathrm{\Delta }\stackrel{~}{\mathrm{\Gamma }}}{\mathrm{\Delta }B_{n,j}(\tau ^{})}}`$ (4.44)
$`+{\displaystyle 𝑑\tau ^{}T_{mi,nj}^1(\tau ,\tau ^{})𝑑\tau ^{\prime \prime }\frac{\delta \stackrel{~}{\mathrm{\Gamma }}}{\delta \theta _{\alpha ,k}(\tau ^{\prime \prime })}\delta _B\psi _{\alpha ,k}(\tau ^{\prime \prime })𝒪_{n,j}(\tau ^{})},`$
where $`T^1`$ is the inverse of $`T`$.
Finally, the correct Ward identity in the standard background gauge is obtained by substituting this expression into (4.37). Once the limit $`\stackrel{~}{B}_mB_m`$ is taken, the total variation $`\mathrm{\Delta }\stackrel{~}{\mathrm{\Gamma }}/\mathrm{\Delta }B_m`$ can be identified with $`\delta \stackrel{~}{\mathrm{\Gamma }}/\delta B_m`$, which denotes the usual functional derivative of the effective action computed in the standard background gauge ( i.e. with $`B_m`$ in the gauge-fixing term.) With this understood, the result can be put in the desired form
$`0={\displaystyle 𝑑\tau \left(\mathrm{\Delta }_ϵB_{m,i}(\tau )\frac{\delta \stackrel{~}{\mathrm{\Gamma }}}{\delta B_{n,j}(\tau )}+\mathrm{\Delta }_ϵ\theta _{\alpha ,i}(\tau )\frac{\delta \stackrel{~}{\mathrm{\Gamma }}}{\delta \theta _{\alpha ,i}(\tau )}\right)},`$ (4.45)
where the effective SUSY transformation laws are given by
$`\mathrm{\Delta }_ϵB_{m,i}(\tau )`$ $`=`$ $`{\displaystyle 𝑑\tau ^{}T_{mi,nj}^1(\tau ^{},\tau )(\delta _ϵ\widehat{y}_{n,j}(\tau ^{})+\delta _B\widehat{y}_{n,j}(\tau ^{})𝒪_ϵ)},`$ (4.46)
$`\mathrm{\Delta }_ϵ\theta _{\alpha ,i}(\tau )`$ $`=`$ $`\delta _ϵ\widehat{\psi }_{\alpha ,i}(\tau )\delta _B\widehat{\psi }_{\alpha ,i}(\tau )𝒪_ϵ`$
$`{\displaystyle 𝑑\tau ^{}𝑑\tau ^{\prime \prime }T_{mk,nj}^1(\tau ^{\prime \prime },\tau ^{})\delta _B\widehat{\psi }_{\alpha ,i}(\tau )𝒪_{n,j}(\tau ^{})\delta _ϵ\widehat{y}_{m,k}(\tau ^{\prime \prime })+\delta _B\widehat{y}_{m,k}(\tau ^{\prime \prime })𝒪_ϵ}.`$
(As said before $`y=\psi =0`$ is understood.) This is the main result of this section.
Note the following features:
* We have succeeded in putting the Ward identity in the form where the supersymmetry transformation laws for the effective action are cleanly identified in closed forms.
* As expected, the supersymmetry and the gauge (BRST) symmetry are non-trivially intertwined. Naively, one might expect that the effective transformation laws are obtained as the expectation values of the original transformation laws (2.7) and (2.8). In our notation, they are represented by $`\delta _ϵ\widehat{y}_{n,j}`$ and $`\delta _ϵ\widehat{\psi }_{\alpha ,i}`$ in $`\mathrm{\Delta }_ϵB_{m,i}`$ and $`\mathrm{\Delta }_ϵ\theta _{\alpha ,i}`$ respectively. The actual transformation laws, (4.46) and (LABEL:delth), are much more complicated. One can see, however, that the corrections to the naive laws all involve $`\delta _B`$, i.e. the BRST transformation. Since the quantization of the system inevitably requires a gauge fixing, this is a universal feature, not special to the standard background gauge adopted here.
* As has already been remarked, the transformation laws derived above are exact, albeit somewhat formal at this stage. In particular, there is no inherent distinction between the tree level contribution and the quantum corrections. Thus it is far from obvious that the anti-commutator of the effective transformations would close solely on the translation generator, as it does at the tree level: We have so far not been able to produce a proof.
In the next section, we will carefully examine the structure of this Ward identity at the 1-loop level and draw implications.
## 5 Explicit Calculations and Implications
We now compute the effective SUSY transformation laws explicitly and study the implications of the Ward identity.
### 5.1 Source-Probe Situation
Although the Ward identity derived in the previous section is valid for any background, we shall restrict ourselves to the source-probe situation, since the existing calculations of the effective action itself, to be compared later, are more complete for this configuration. The background fields representing this situation for the gauge group $`U(N+1)`$ were already described in (3.19), which we display again for convenience:
$`B_m=\mathrm{diag}(r_m,0,0,\mathrm{},0),\theta _\alpha =\mathrm{diag}(\theta _\alpha ,0,0,\mathrm{},0).`$ (5.48)
$`r_m`$ represents the coordinate of the probe and $`\theta _\alpha `$ its spin content. We shall denote the time derivative of the coordinate by $`v_m_\tau r_m`$.
To facilitate the computations, it is convenient to introduce the following notations: For any matrix $`U`$ we define $`U_IU_{1I},U_I^{}U_{I1}`$, where $`I=2N+1`$. (The symbol here does not stand for complex conjugation.) We shall call such index $`I`$ the off-diagonal matrix vector index. Then the “11” component of a product of two matrices becomes $`(UV)_{11}=_{I=2}^{N+1}U_IV_I^{}`$, which we abbreviate as $`UV^{}`$.
With this convention, the basic quantities appearing in the Ward identity take the following forms:
$`\delta _ϵ\widehat{y}_m`$ $`=`$ $`iϵ^T\gamma _m\theta ,`$ (5.49)
$`\delta _ϵ\widehat{\psi }_\alpha `$ $`=`$ $`i(v/_{\alpha \beta }ig^2\gamma _{\alpha \beta }^mAY_m^{}A^{}Y_m+{\displaystyle \frac{g^2}{2}}\gamma _{\alpha \beta }^{mn}Y_mY_n^{}Y_m^{}Y_n)ϵ_\beta ,`$ (5.50)
$`\delta _B\widehat{y}_m`$ $`=`$ $`ig^2(Y_mC^{}Y_m^{}C),`$ (5.51)
$`\delta _B\widehat{\psi }_\alpha `$ $`=`$ $`ig^2(C\mathrm{\Psi }_\alpha ^{}C^{}\mathrm{\Psi }_\alpha ),`$ (5.52)
$`𝒪_ϵ`$ $`=`$ $`iϵ_\alpha {\displaystyle }d\tau ({\displaystyle \frac{1}{g}}\overline{C}_{11}(\dot{\theta }_\alpha +\dot{\widehat{\psi }}_\alpha )+\overline{C}(_\tau +r/)_{\alpha \beta }\mathrm{\Psi }_\beta ^{}+\overline{C}^{}(_\tau r/)_{\alpha \beta }\mathrm{\Psi }_\beta `$ (5.53)
$`+\overline{C}_{IJ}\dot{\mathrm{\Theta }}_{JI}),`$
$`𝒪_m`$ $`=`$ $`\overline{C}^{}Y_m\overline{C}Y_m^{}.`$ (5.54)
### 5.2 One-Loop at Order Two
In this article, we shall perform the calculation of the simplest non-trivial contributions to the effective transformation laws, namely those which govern the order 2 part of the 1-loop effective action. Here the order is defined as in (3.20), i.e. the number of derivatives plus twice the number of fermions. Since the tree action is already of order 2, this means that we need to compute $`\mathrm{\Delta }_ϵr_m`$ and $`\mathrm{\Delta }_ϵ\theta _\alpha `$ to order 0 at 1-loop, where the transformation parameters $`ϵ_\alpha `$ are considered to be of order $`1/2`$.
Such contributions are further classified according to the number of $`\theta `$’s involved. Because $`r_m`$ and $`\theta _\alpha `$ are of order 0 and $`1/2`$ respectively, we need to consider $`\mathrm{\Delta }_ϵr_m`$ to linear order in $`\theta `$, while for $`\mathrm{\Delta }_ϵ\theta _\alpha `$ we need $`\theta ^2`$ contributions as well. In what follows, it is more convenient to refer only to the number of additional $`\theta `$’s relative to the tree level contribution. For instance, “a correction to $`𝒪(\theta ^0)`$” for $`\mathrm{\Delta }_ϵr_m`$ refers to a term of the form $`ϵ\theta f(r)`$.
At 1-loop order, since the terms involving the BRST variation $`\delta _B`$ start only at 1-loop, the expressions (4.46) and (LABEL:delth) for $`\mathrm{\Delta }_ϵr_m`$ and $`\mathrm{\Delta }_ϵ\theta `$ simplify to
$`\mathrm{\Delta }_ϵr^m(\tau )`$ $`=`$ $`{\displaystyle \frac{1}{i}}(ϵ\gamma ^m\theta )(\tau )+{\displaystyle 𝑑\tau ^{}\delta _B\widehat{y}_m(\tau )𝒪_n(\tau ^{})\frac{1}{i}(ϵ\gamma ^n\theta )(\tau ^{})}+\delta _B\widehat{y}_m(\tau )𝒪_ϵ,`$ (5.55)
$`\mathrm{\Delta }_ϵ\theta _\alpha (\tau )`$ $`=`$ $`\delta _ϵ\widehat{\psi }_\alpha (\tau )\delta _B\widehat{\psi }_\alpha (\tau )𝒪_ϵ{\displaystyle 𝑑\tau ^{}\delta _B\widehat{\psi }_\alpha (\tau )𝒪_m(\tau ^{})\frac{1}{i}(ϵ\gamma ^m\theta )(\tau ^{})}.`$ (5.56)
The expectation values of the composite operators themselves simplify considerably at this order. It is easy to check that at 1-loop only the 2-point functions contribute. Moreover, since there are no mixing between the fields with different off-diagonal matrix vector indices in such propagators, we always have the structure $`U_I(\tau )V_J^{}(\tau ^{})=\delta _{IJ}U(\tau )V^{}(\tau ^{})`$, where $`U`$ and $`V^{}`$ can be regarded as single-component fields. $`\delta _{IJ}`$’s then contract to produce a factor of $`N`$, which goes with the coupling $`g^2`$. With these remarks, the relevant multi-body expectation values become
$`\delta _ϵ\widehat{\psi }_\alpha (\tau )`$ $`=`$ $`i(v/ϵ)_\alpha (\tau )+g^2N(\gamma _{\alpha \beta }^mA(\tau )Y_m^{}(\tau )A^{}(\tau )Y_m(\tau )`$
$`+{\displaystyle \frac{i}{2}}\gamma _{\alpha \beta }^{mn}Y_m(\tau )Y_n^{}(\tau )Y_m^{}(\tau )Y_n(\tau ))ϵ_\beta (\tau ),`$
$`\delta _B\widehat{y}_n(\tau )𝒪_ϵ`$ $`=`$ $`ϵ_\alpha g^2N{\displaystyle }d\tau ^{}(C^{}(\tau )\overline{C}(\tau ^{})(+r/)_{\alpha \beta }\mathrm{\Psi }_\beta ^{}(\tau ^{})Y_m(\tau )`$ (5.58)
$`C(\tau )\overline{C}^{}(\tau ^{})(r/)_{\alpha \beta }\mathrm{\Psi }_\beta (\tau ^{})Y_m^{}(\tau )),`$
$`\delta _B\widehat{\psi }_\beta (\tau )𝒪_ϵ`$ $`=`$ $`ϵ_\alpha g^2N{\displaystyle }d\tau ^{}(C^{}(\tau )\overline{C}(\tau ^{})(+r/)_{\alpha \gamma }\mathrm{\Psi }_\gamma ^{}(\tau ^{})\mathrm{\Psi }_\beta (\tau )`$
$`C(\tau )\overline{C}^{}(\tau ^{})(r/)_{\alpha \gamma }\mathrm{\Psi }_\gamma (\tau ^{})\mathrm{\Psi }_\beta ^{}(\tau ))`$
$`gϵ_\beta {\displaystyle 𝑑\tau ^{}C\mathrm{\Psi }^{}(\tau )\overline{C}_{11}\dot{\widehat{\psi }}_\beta (\tau ^{})}+gϵ_\beta {\displaystyle 𝑑\tau ^{}C^{}\mathrm{\Psi }(\tau )\overline{C}_{11}\dot{\widehat{\psi }}_\beta (\tau ^{})},`$
$`\delta _B\widehat{y}_n(\tau ^{})𝒪_m(\tau )`$ $`=`$ $`ig^2N(Y_m(\tau ^{})Y_n^{}(\tau )(C^{}(\tau ^{})\overline{C}(\tau )`$ (5.60)
$`+Y_m^{}(\tau ^{})Y_n(\tau )(C(\tau ^{})\overline{C}^{}(\tau )),`$
where $``$ stands for the derivative with respect to $`\tau `$. We will now evaluate these expressions to the relevant order by perturbation theory.
#### 5.2.1 Corrections at $`𝒪(\theta ^0)`$
Let us begin with the corrections at $`𝒪(\theta ^0)`$. First we need to compute the tree-level propagators. The part of the Lagrangian quadratic in the fields without $`\theta `$ is of the form
$``$ $`=`$ $`Y_m(^2+r^2)Y_m^{}+\stackrel{~}{A}(^2+r^2)\stackrel{~}{A}^{}2iv_m(\stackrel{~}{A}Y_m^{}\stackrel{~}{A}^{}Y_m)`$ (5.61)
$`+{\displaystyle \frac{1}{2g^2}}\widehat{\psi }\widehat{\psi }+\mathrm{\Psi }_\alpha (+r/)_{\alpha \beta }\mathrm{\Psi }_\beta ^{}`$
$`+i\overline{C}(^2+r^2)C^{}+i\overline{C}^{}(^2+r^2)C+i\overline{C}_{11}(^2)C_{11}.`$
Following the remark already made on the trivial dependence on the off-diagonal matrix vector indices, we may treat the fields $`Y_m,\stackrel{~}{A},\mathrm{\Psi }_\alpha `$ and $`C,\overline{C}`$ in $``$ as if they were single-component. The propagators for the massless fields $`\widehat{\psi }_\alpha `$ and $`C_{11},\overline{C}_{11}`$ turn out to be needed only at $`𝒪(\theta ^2)`$, to be discussed in the next subsection.
The simplest is the ghost propagator, which can be directly read off as
$`\overline{C}(\tau )C^{}(\tau ^{})=\overline{C}^{}(\tau )C(\tau ^{})`$
$`=C^{}(\tau )\overline{C}(\tau ^{})=C(\tau )\overline{C}^{}(\tau ^{})=i\tau |\mathrm{\Delta }|\tau ^{},`$ (5.62)
where
$`\mathrm{\Delta }`$ $``$ $`(^2+r^2)^1.`$ (5.63)
The propagators for the $`\mathrm{\Psi }_\alpha `$ system are given by
$`\mathrm{\Psi }_\alpha (\tau )\mathrm{\Psi }_\beta ^{}(\tau ^{})`$ $`=`$ $`\tau |D_{\alpha \beta }^{}|\tau ^{},`$ (5.64)
$`\mathrm{\Psi }_\alpha ^{}(\tau )\mathrm{\Psi }_\beta (\tau ^{})`$ $`=`$ $`\tau |D_{\alpha \beta }^+|\tau ^{},`$ (5.65)
where
$`D^\pm `$ $``$ $`(\pm r/)^1.`$ (5.66)
¿From the relation $`(\pm r/)(r/)=(^2+r^2)(1\pm \mathrm{\Delta }v/)`$, we can expand $`D^\pm `$ in powers of the velocity $`v(\tau )`$ as
$`D^\pm `$ $`=`$ $`(r/)(1\mathrm{\Delta }v/)^1\mathrm{\Delta }=(r/)(\mathrm{\Delta }\pm \mathrm{\Delta }v/\mathrm{\Delta }+\mathrm{}).`$ (5.67)
To the order of interest, we will only need the $`v`$-independent part.
Now as for the $`Y_m`$-$`\stackrel{~}{A}`$ system, we have a mixing term with an arbitrary coefficient function $`v_m(\tau )`$ and it can only be resolved in the derivative expansion. Expanding in powers of the mixing term, we readily obtain
$`Y^m(\tau )^{}Y^n(\tau ^{})`$ $`=`$ $`\delta _{mn}\tau |\mathrm{\Delta }|\tau ^{}+𝒪(v^2),`$
$`\stackrel{~}{A}(\tau )\stackrel{~}{A}^{}(\tau ^{})`$ $`=`$ $`\tau |\mathrm{\Delta }|\tau ^{}+𝒪(v^2),`$ (5.68)
$`Y_m(\tau )\stackrel{~}{A}^{}(\tau ^{})`$ $`=`$ $`\stackrel{~}{A}(\tau )Y_m^{}(\tau ^{})=\tau |2i\mathrm{\Delta }v_m\mathrm{\Delta }|\tau ^{}+𝒪(v^3).`$
The only other 2-point functions appearing at this order are $`\mathrm{\Psi }_\alpha (\tau )Y_m^{}(\tau ^{})`$ and
$`\mathrm{\Psi }_\alpha ^{}(\tau )Y_m(\tau ^{})`$, which can be readily computed by inserting the vertices of the form
$`\theta _\alpha \gamma _{\alpha \beta }^m(Y_m\mathrm{\Psi }_\beta ^{}Y_m^{}\mathrm{\Psi }_\beta ).`$ (5.69)
The results are
$`\mathrm{\Psi }_\alpha (\tau )Y_m^{}(\tau ^{})`$ $`=`$ $`\tau |D_{\alpha \beta }^{}(\gamma ^m\theta )_\beta \mathrm{\Delta }|\tau ^{},`$ (5.70)
$`\mathrm{\Psi }_\alpha ^{}(\tau )Y_m(\tau ^{})`$ $`=`$ $`\tau |D_{\alpha \beta }^+(\gamma ^m\theta )_\beta \mathrm{\Delta }|\tau ^{}.`$ (5.71)
With this preparation, the calculations of the various expectation values in the Ward identity can be performed efficiently to the desired order with the use of the so-called “normal ordering method” developed in . The essence of this method is to first rearrange the order of a product of various operators and functions into the standard form $`f(\tau )^m\mathrm{\Delta }^n`$, using the commutation relations such as $`[\mathrm{\Delta },]=2\mathrm{\Delta }(rv)\mathrm{\Delta }`$ etc., and then use Baker-Campbell-Hausdorff and the Gaussian integration formulas to evaluate it. A useful list of formulas so obtained are collected in the appendix of . Below, we shall give a sample calculation of this sort and then simply list the results for the needed expectation values.
As an example, let us consider $`\delta _B\widehat{y}_n(\tau )𝒪_ϵ`$ given in (5.58). Substituting the expressions for the various propagators already computed, it becomes
$`\delta _B\widehat{y}^n(\tau )𝒪_ϵ`$ $`=`$ $`iϵ^Tg^2N\tau |\mathrm{\Delta }[(+r/)D^++(r/)D^{}]\gamma ^n\theta \mathrm{\Delta }|\tau `$ (5.72)
$`=`$ $`2iϵ^Tg^2N\tau |\mathrm{\Delta }(^2+r^2)\mathrm{\Delta }\gamma ^n\theta \mathrm{\Delta }|\tau .`$
To the order of interest, the normal-ordering is trivial and we get
$`\delta _B\widehat{y}_n(\tau )𝒪_ϵ=2iϵ^Tg^2N\gamma _n\theta \tau |\mathrm{\Delta }^2|\tau =iϵ^T\gamma _n\theta {\displaystyle \frac{1}{2}}{\displaystyle \frac{g^2N}{r^3}}.`$ (5.73)
In a similar manner, we obtain the following results:
$`\delta _B\widehat{y}_n(\tau ^{})𝒪_m(\tau )`$ $`=`$ $`2g^2N\delta _{nm}\tau ^{}|\mathrm{\Delta }|\tau ^2,`$ (5.74)
$`{\displaystyle }d\tau ^{}\delta _B\widehat{y}_m(\tau )𝒪_n(\tau ^{}){\displaystyle \frac{1}{i}}(ϵ\gamma ^n\theta )(\tau ^{})`$ $`=`$ $`iϵ\gamma _m\theta {\displaystyle \frac{1}{2}}{\displaystyle \frac{g^2N}{r^3}},`$ (5.75)
$`\delta _B\widehat{\psi }_\beta (\tau )𝒪_ϵ`$ $`=`$ $`0,`$ (5.76)
$`\delta _B\widehat{y}_m(\tau )𝒪_ϵ`$ $`=`$ $`iϵ\gamma _m\theta {\displaystyle \frac{1}{2}}{\displaystyle \frac{g^2N}{r^3}},`$ (5.77)
$`AY_m^{}A^{}Y_m`$ $`=`$ $`{\displaystyle \frac{N}{i}}{\displaystyle \frac{v_m}{r^3}},`$ (5.78)
$`Y_mY_n^{}Y_m^{}Y_n`$ $`=`$ $`0.`$ (5.79)
Making use of these formulas, we find the effective SUSY transformation laws at $`𝒪(\theta ^0)`$ to be
$`\mathrm{\Delta }_ϵr^m`$ $`=`$ $`iϵ^T\gamma ^m\theta \left(1+{\displaystyle \frac{g^2N}{r^3}}\right),`$ (5.80)
$`\mathrm{\Delta }_ϵ\theta `$ $`=`$ $`iv/ϵ\left(1{\displaystyle \frac{g^2N}{r^3}}\right).`$ (5.81)
Already at this order, there are non-trivial corrections to the tree level laws. For $`\mathrm{\Delta }_ϵ\theta `$ the entire correction came from the expectation value of the non-linear part of the original SUSY transformation law, namely from $`ig[\stackrel{~}{A},Y_m]\gamma ^m`$ contained in $`[D_\tau ,X_m]\gamma ^m`$. On the other hand, the same procedure is not applicable for $`\mathrm{\Delta }_ϵr_m`$: If one naively took the expectation value of the basic transformation law $`\delta _ϵX^m=iϵ\gamma ^m\mathrm{\Theta }`$, one would not get any corrections to all orders. What we want is the effective SUSY transformation operating on the effective action $`\stackrel{~}{\mathrm{\Gamma }}`$ and this can only be obtained through the analysis of the Ward identity, as we have done. One can see from the calculations outlined above that exactly half of the quantum correction for $`\mathrm{\Delta }_ϵr^m`$ came from $`\delta _ϵ\stackrel{~}{S}_{gg}`$ and the other half was produced from the procedure of taking into account the dependence on the extra parameter field we originally called $`\stackrel{~}{B}_m`$.
### 5.3 Corrections at $`𝒪(\theta ^2)`$
Now we move on to the corrections at $`𝒪(\theta ^2)`$. Since $`\theta ^2`$ is of order 1, we need to compute such contributions only for $`\mathrm{\Delta }_ϵ\theta _\alpha `$. The procedure is entirely similar to the $`𝒪(\theta ^0)`$ case but the calculations are more involved and we relegate the details to the Appendix A. As shown there, 7 types of diagrams contribute. Two of them, diagrams (B2-a,b), involve genuine 3 point vertices as well as massless propagators given by
$`\widehat{\psi }_\alpha (\tau )\widehat{\psi }_\beta (\tau ^{})=g^2\tau |{\displaystyle \frac{1}{}}|\tau ^{},C_{11}(\tau )\overline{C}_{11}(\tau ^{})=i\tau |{\displaystyle \frac{1}{^2}}|\tau ^{}.`$ (5.82)
They are singular in the infrared, but such singularities cancel in the end result. After some calculations, we find the $`𝒪(\theta ^2)`$ corrections to $`\mathrm{\Delta }_ϵ\theta _\alpha `$ to be
$`\mathrm{\Delta }_ϵ^{\theta ^2}\theta _\alpha `$ $`=`$ $`{\displaystyle \frac{3ig^2N}{16r^5}}(r_l(\theta \gamma ^{mnl}\theta )(\gamma ^{nm}ϵ)_\alpha +2r_l(\theta \gamma ^{ml}\theta )(\gamma ^mϵ)_\alpha 4r_l(ϵ\gamma ^m\theta )(\theta \gamma ^{ml})_\alpha `$ (5.83)
$`+4(ϵ\theta )(r/\theta )_\alpha 4(ϵr/\theta )\theta _\alpha ).`$
Although it appears quite complicated, it can be drastically simplified with the use of the $`SO(9)`$ Fierz identities described in the Appendix B. The relevant identity is
$`0`$ $`=`$ $`r_l(\theta \gamma ^{mnl}\theta )(\gamma ^{nm}ϵ)_\alpha +2r_l(\theta \gamma ^{nl}\theta )(\gamma ^nϵ)_\alpha +4r_l(ϵ\gamma ^n\theta )(\gamma ^{nl}\theta )_\alpha `$ (5.84)
$`+4r_l(ϵ\theta )(\gamma ^l\theta )_\alpha +12r_l(ϵ\gamma ^l\theta )\theta _\alpha .`$
Applying this to (5.83), we get
$`\mathrm{\Delta }_ϵ^{\theta ^2}\theta _\alpha ={\displaystyle \frac{3ig^2N}{r^5}}(ϵr/\theta )\theta _\alpha .`$ (5.85)
Let us summarize the results. To order 2 at 1-loop, the effective supersymmetry transformations laws are found to be
$`\mathrm{\Delta }_ϵr^m`$ $`=`$ $`iϵ^T\gamma ^m\theta \left(1+{\displaystyle \frac{g^2N}{r^3}}\right),`$ (5.86)
$`\mathrm{\Delta }_ϵ\theta _\alpha `$ $`=`$ $`i(v/ϵ)_\alpha \left(1{\displaystyle \frac{g^2N}{r^3}}\right){\displaystyle \frac{3ig^2N}{r^5}}(ϵr/\theta )\theta _\alpha .`$ (5.87)
#### 5.3.1 Closure of the Effective Algebra and Field Redefinition
It is of interest to examine the closure property of the transformation laws obtained above. By a simple calculation, one immediately finds
$`[\mathrm{\Delta }_{ϵ_1},\mathrm{\Delta }_{ϵ_2}]r_m`$ $`=`$ $`2\dot{r}_m(ϵ_2ϵ_1)+𝒪(g^4),`$ (5.88)
$`[\mathrm{\Delta }_{ϵ_1},\mathrm{\Delta }_{ϵ_2}]\theta _\alpha `$ $`=`$ $`2\dot{\theta }_\alpha (ϵ_2ϵ_1)+𝒪(g^4).`$ (5.89)
Thus to the 1-loop order, the closure turned out to be precisely canonical.
As we remarked at the end of Sec.4, this is not a feature guaranteed by the general analysis. One way to appreciate this is to consider a field redefinition which makes the form of $`\mathrm{\Delta }_ϵr^m`$ to be the same as the one at the tree level. On general grounds, the most general form of $`\mathrm{\Delta }_ϵr^m`$ and $`\mathrm{\Delta }_ϵ\theta _\alpha `$, to the order we are considering, are
$`\mathrm{\Delta }_ϵr^m`$ $`=`$ $`{\displaystyle \frac{1}{i}}(ϵ\gamma ^m\theta )\left(1+g^2F(r)\right),`$ (5.90)
$`\mathrm{\Delta }_ϵ\theta _\alpha `$ $`=`$ $`i(v/ϵ)_\alpha +g^2G_{\alpha \beta }(r,\theta )ϵ_\beta .`$ (5.91)
¿From (5.90), the desired field redefinition can be read off as
$`\stackrel{~}{\theta }_\alpha \left(1+g^2F(r)\right)\theta _\alpha .`$ (5.92)
The transformation law for this new field is then
$`\mathrm{\Delta }_ϵ\stackrel{~}{\theta }_\alpha =i(v/ϵ)_\alpha +g^2\left(i(v/ϵ)_\alpha F(r)+\mathrm{\Delta }_ϵF(r)\stackrel{~}{\theta }_\alpha +G_{\alpha \beta }(r,\stackrel{~}{\theta })ϵ_\beta \right)+𝒪(g^4).`$ (5.93)
It is quite non-trivial that the $`𝒪(g^2)`$ part of this expression vanishes exactly.
### 5.4 Implication of the Ward Identity
Having found the transformation laws, we are now ready to analyze the consequence of the Ward identity on the structure of $`\stackrel{~}{\mathrm{\Gamma }}`$.
Let us write the effective action up to 1-loop as $`\stackrel{~}{\mathrm{\Gamma }}=\stackrel{~}{\mathrm{\Gamma }}_0+\stackrel{~}{\mathrm{\Gamma }}_1`$, the subscript denoting the number of loops. The tree-level action is given by
$`\stackrel{~}{\mathrm{\Gamma }}_0={\displaystyle 𝑑\tau \left(\frac{v^2}{2g^2}+\frac{\theta \dot{\theta }}{2g^2}\right)}.`$ (5.94)
As for $`\stackrel{~}{\mathrm{\Gamma }}_1`$, it is easy to convince oneself that the most general structure at order 2, up to total derivatives, is
$`\stackrel{~}{\mathrm{\Gamma }}_1`$ $`=`$ $`{\displaystyle }d\tau (A{\displaystyle \frac{v^2}{r^3}}+B{\displaystyle \frac{(vr)^2}{r^5}}+C{\displaystyle \frac{\theta ^T\dot{\theta }}{r^3}}+D{\displaystyle \frac{\theta ^Tr/\dot{\theta }}{r^4}}`$ (5.95)
$`+E{\displaystyle \frac{\theta ^T\gamma ^{mn}r_mv_n\theta }{r^5}}+F{\displaystyle \frac{r^lr^k(\theta \gamma ^{lm}\theta )(\theta \gamma ^{mk}\theta )}{r^7}}),`$
where $`AF`$ are numerical constants<sup>4</sup><sup>4</sup>4The structure of the form $`r^lr^k(\theta \gamma ^{lmn}\theta )(\theta \gamma ^{mkn}\theta )`$ can be expressed in terms of the last term in (5.95) via a Fierz identity and hence can be omitted..
Now we demand that $`\mathrm{\Delta }_ϵ\stackrel{~}{\mathrm{\Gamma }}`$ vanish to the order of interest. The simplest way to proceed is as follows. First, look at the $`𝒪(\theta ^5)`$ terms, which can only be produced from the last term of (5.95) by the tree-level part of the transformation $`\mathrm{\Delta }_ϵr_m`$. They read
$`{\displaystyle 𝑑\tau F\frac{i}{r^9}r^k\left(2r^2(ϵ\gamma ^l\theta )(\theta \gamma ^{lm}\theta )+7r^lr^j(ϵ\gamma ^j\theta )(\theta \gamma ^{lm}\theta )\right)(\theta \gamma ^{mk}\theta )}.`$ (5.96)
It can be checked that the integrand does not vanish by any of the Fierz identities<sup>5</sup><sup>5</sup>5We have also checked this numerically.. Thus we find $`F=0`$. Next, demand that the $`𝒪(\theta )`$ part of $`\mathrm{\Delta }_ϵ\stackrel{~}{\mathrm{\Gamma }}`$ vanish. It is straightforward to show that this reduces the allowed form of $`\stackrel{~}{\mathrm{\Gamma }}_1`$ to be
$`\stackrel{~}{\mathrm{\Gamma }}_1={\displaystyle 𝑑\tau \left(A\frac{v^2}{r^3}+(A+N)\frac{\theta \dot{\theta }}{r^3}+\frac{3A}{2}\frac{\theta \gamma ^{mn}r_mv_n\theta }{r^5}\right)},`$ (5.97)
where $`A`$ remains undetermined. Finally, look at the $`𝒪(\theta ^3)`$ part of $`\mathrm{\Delta }_ϵ\stackrel{~}{\mathrm{\Gamma }}`$. The contributions arising from the tree level transformation of $`\stackrel{~}{\mathrm{\Gamma }}_1`$ are
$`{\displaystyle \frac{3iN}{r^5}}r_m(ϵ\gamma ^m\theta )(\theta \dot{\theta })`$
$`+{\displaystyle \frac{3iA}{2r^5}}\left(2r_m(ϵ\gamma ^m\theta )(\theta \dot{\theta })v_n(ϵ\gamma ^m\theta )(\theta \gamma ^{mn}\theta )r_m(ϵ\gamma ^n\dot{\theta })(\theta \gamma ^{mn}\theta )\right)`$
$`+{\displaystyle \frac{15iA}{4r^7}}r_mr_nv_l\left((ϵ\gamma ^n\theta )(\theta \gamma ^{ml}\theta )+(ϵ\gamma ^m\theta )(\theta \gamma ^{nl}\theta )\right).`$ (5.98)
On the other hand, the 1-loop level transformation applied to $`\stackrel{~}{\mathrm{\Gamma }}_0`$ produces
$`{\displaystyle \frac{3iN}{r^5}}r_m(ϵ\gamma ^m\theta )(\theta \dot{\theta }),`$ (5.99)
which cancels the first term of (5.98). The remaining terms in (5.98) are all proportional to $`A`$. Now note that while the four-fermion structures in the second line of (5.98) have only one “free index”, contracted with an arbitrary vector $`r_m(\tau )`$ or $`v_m(\tau )`$, the ones in the last line carry three free indices. Since the Fierz identities can only relate structures with the same number of free indices, expressions in these two lines cannot cancel each other. Moreover, it is easy to check, using the Fierz identities given in the Appendix A, that the second line does not vanish by itself. This then proves $`A=0`$.
In summary, we have found that the order 2 contribution at 1-loop for the effective action in the background gauge is completely determined by the requirement of supersymmetry and takes the form
$`\stackrel{~}{\mathrm{\Gamma }}_1={\displaystyle 𝑑\tau N\frac{\theta \dot{\theta }}{r^3}}.`$ (5.100)
This indeed agrees, including the overall normalization, with the direct calculation performed in . It can be easily checked that the normalization is directly linked to the magnitude of the $`𝒪(\theta ^0)`$ quantum corrections in (5.86) and (5.87), which cannot be determined by the closure property alone.
As we shall discuss in the concluding section, the power of our off-shell Ward identity can only be fully utilized at the next order in the derivative expansion, where the most general form of the action unavoidably contains many terms with higher derivatives, such as $`\ddot{r}_m`$ and $`\ddot{\theta }_\alpha `$. Nevertheless, it is gratifying that already at order 2 it has enabled us to see explicitly how the supersymmetry and the gauge symmetry intimately work together to dictate the form of the effective action.
## 6 Summary and Discussions
In this paper, we have derived the exact supersymmetric off-shell Ward identity for Matrix theory as a step toward answering the “old” yet important unsettled problem: “To what extent do the symmetries, in particular the supersymmetry, determine the low-energy effective action of Matrix theory?” Our work was motivated by the observation that the existing analyses are incomplete in that off-shell trajectories with arbitrary time-dependence have not been fully considered. An important aspect of our Ward identity is that it allows the quantum-corrected effective supersymmetry transformation laws to be directly identified in closed form. They exhibit an intricate interplay with the gauge (BRST) symmetry of the theory, a feature not properly appreciated previously.
As an application, we computed the explicit form of these transformation laws at 1-loop to the lowest order in the derivative expansion, and examined if the invariance under them determines the form of the effective action to the corresponding order. We found that the answer is affirmative, confirming the earlier result. This is as expected since at this order the higher derivatives, such as the acceleration etc., can be eliminated from the effective action by partial integration and the analysis is essentially the same as in the existing literature.
The full significance of our off-shell Ward identity should become apparent starting from the next order, i.e. from order 4, where complete elimination of higher derivatives will no longer be possible. There will be a considerable number of independent structures allowed in the most general effective action. Even the proper listing of them requires careful analysis due to the total derivative ambiguities and the existence of non-trivial Fierz identities. Nonetheless, a preliminary investigation indicates that, with an aid of computerized calculation, it appears feasible to determine whether SUSY alone is enough to fix the form of the effective action at order 4 for arbitrary trajectories.
Another important direction into which to extend our present work is to apply our Ward identity to genuinely multi-body configurations. To find out whether the remarkable agreement with supergravity in such a situation is due to supersymmetry alone would certainly deepen our understanding of the Matrix theory further.
At the more formal and structural level, we should mention that a further study should be made on the issue of the closure property of the effective SUSY transformations. As we have shown, at the lowest order the closure turned out to be canonical. So far, however, we have not been able to answer whether this persists at higher orders and loops. An analysis based on the general closed form expressions for the transformation laws should shed light on this intriguing question.
We hope to be able to report on these and other related issues in the near future.
Acknowledgment
We would like to express our special gratitude to Y. Okawa for a number of clarifying discussions and his interest in our work. Y. K acknowledges the warm hospitality extended to him by the organizers at the third international symposium on Frontiers of Fundamental Physics (Hyderabad, India), where a preliminary version of this work was presented. This work is supported in part by Grant-in-Aid for Scientific Research on Priority Area #707 “Supersymmetry and Unified Theory of Elementary Particles”, Grant-in-Aid for Scientific Research No. 09640337, and Grant-in-Aid for International Scientific Research (Joint Research) No. 10044061, from Japan Ministry of Education, Science and Culture.
Appendix A: Calculations of $`𝒪(\theta ^2)`$ terms in $`\mathrm{\Delta }_ϵ\theta _\alpha `$
In this appendix, we exhibit some details of the calculations of $`𝒪(\theta ^2)`$ terms in $`\mathrm{\Delta }_ϵ\theta _\alpha `$ at 1-loop order.
At this order, what we need to evaluate is (see (5.56))
$`\mathrm{\Delta }_ϵ\theta _\alpha (\tau )=\mathrm{\Delta }_ϵ\theta _\alpha ^A(\tau )+\mathrm{\Delta }_ϵ\theta _\alpha ^B(\tau )+\mathrm{\Delta }_ϵ\theta _\alpha ^C(\tau ),`$ (A.1)
where
$`\mathrm{\Delta }_ϵ\theta _\alpha ^A(\tau )`$ $``$ $`\delta _ϵ\widehat{\psi }_\alpha (\tau ),`$ (A.2)
$`\mathrm{\Delta }_ϵ\theta _\alpha ^B(\tau )`$ $``$ $`\delta _B\widehat{\psi }_\alpha (\tau )𝒪_ϵ,`$ (A.3)
$`\mathrm{\Delta }_ϵ\theta _\alpha ^C(\tau )`$ $``$ $`{\displaystyle 𝑑\tau ^{}\delta _B\widehat{\psi }_\alpha (\tau )𝒪_m(\tau ^{})\frac{1}{i}(ϵ\gamma ^m\theta )(\tau ^{})}.`$ (A.4)
Hereafter in this appendix, $`\mathrm{\Delta }_ϵ\theta _\alpha `$ will refer only to the $`𝒪(\theta ^2)`$ part. Also, we shall omit the overall factor of $`N`$, except in the final expression. The relevant Feynman diagrams are shown in Fig.1.
Calculation of $`\mathrm{\Delta }_ϵ\theta _\alpha ^A(\tau )`$: The explicit expression is
$`\mathrm{\Delta }\theta _\alpha ^A=g^2\stackrel{~}{A}Y_m^{}Y_m\stackrel{~}{A}^{}(\gamma ^mϵ)_\alpha +i{\displaystyle \frac{g^2}{2}}Y_mY_n^{}Y_nY_m^{}(\gamma ^{mn}ϵ)_\alpha .`$ (A.5)
To compute the expectation values above to $`𝒪(\theta ^2)`$, we need to perform the second order perturbation using the $`𝒪(\theta )`$ vertices
$`\stackrel{~}{S}_{\stackrel{~}{A}\mathrm{\Psi }}`$ $`=`$ $`i{\displaystyle 𝑑\tau \theta _\alpha (\stackrel{~}{A}\mathrm{\Psi }_\alpha ^{}\mathrm{\Psi }_\alpha \stackrel{~}{A}^{})},`$ (A.6)
$`\stackrel{~}{S}_{Y\mathrm{\Psi }}`$ $`=`$ $`{\displaystyle 𝑑\tau (\gamma ^m\theta )_\alpha (Y_m\mathrm{\Psi }_\alpha ^{}Y_m^{}\mathrm{\Psi }_\alpha )}.`$ (A.7)
This generates the diagrams (A1) and (A2) in Fig. 1, for the first and the second term in (A.5) respectively. Consider for example the term $`\stackrel{~}{A}(\tau )Y_m^{}(\tau )`$. Using the tree-level propagators given in (5.64) and (5.68), this can be computed as
$`\stackrel{~}{A}(\tau )Y_m^{}(\tau )`$ $`=`$ $`i{\displaystyle 𝑑\tau ^{}𝑑\tau ^{\prime \prime }\stackrel{~}{A}(\tau )\stackrel{~}{A}^{}(\tau ^{})\theta _\alpha (\tau ^{})(\gamma ^n\theta (\tau ^{\prime \prime }))_\beta \mathrm{\Psi }_\alpha (\tau ^{})\mathrm{\Psi }_\beta ^{}(\tau ^{\prime \prime })}`$ (A.8)
$`\times Y_n(\tau ^{\prime \prime })Y_m^{}(\tau )`$
$`=`$ $`i\tau |\mathrm{\Delta }\theta _\alpha (+r/)_{\alpha \beta }(\gamma ^m\theta )_\beta \mathrm{\Delta }^2|\tau .`$
Performing the normal-ordering, neglecting the terms which generate derivatives, this becomes
$`\stackrel{~}{A}(\tau )Y_m^{}(\tau )=i(\theta r/\gamma ^m\theta )\tau |\mathrm{\Delta }^3|\tau ={\displaystyle \frac{3i}{16r^5}}(\theta r/\gamma ^m\theta )={\displaystyle \frac{3i}{16r^5}}r_l(\theta \gamma ^{ml}\theta ).`$ (A.9)
Likewise, one easily finds that $`\stackrel{~}{A}^{}(\tau )Y_m(\tau )`$ gives exactly the same contribution. The evaluation of the diagram (A2) proceeds in an entirely similar manner. In this way one finds
$`\mathrm{\Delta }\theta _\alpha ^A={\displaystyle \frac{3ig^2}{16r^5}}\left(2r_l(\theta \gamma ^{ml}\theta )(\gamma ^mϵ)_\alpha r_l(\theta \gamma ^{mnl}\theta )(\gamma ^{nm}ϵ)_\alpha \right).`$ (A.10)
Calculation of $`\mathrm{\Delta }_ϵ\theta _\alpha ^B(\tau )`$: This receives contributions from two classes of diagrams, (B1) and (B2) in Fig.1.
For (B1), we have
$`\mathrm{\Delta }_ϵ\theta _\alpha ^{B1}(\tau )`$ $`=`$ $`ϵ_\beta g^2{\displaystyle }d\tau ^{}(C^{}(\tau )\overline{C}(\tau ^{})(+r/)_{\beta \gamma }\mathrm{\Psi }_\gamma ^{}(\tau ^{})\mathrm{\Psi }_\alpha (\tau )`$ (A.11)
$`C(\tau )\overline{C}^{}(\tau ^{})(r/)_{\beta \gamma }\mathrm{\Psi }_\gamma (\tau ^{})\mathrm{\Psi }_\alpha ^{}(\tau )).`$
Insertions of the vertices $`\stackrel{~}{S}_{\stackrel{~}{A}\mathrm{\Psi }}`$ (A.6) and $`\stackrel{~}{S}_{Y\mathrm{\Psi }}`$ (A.7) twice generate the diagrams (B1-a) and (B1-b). Again neglecting the derivatives produced in the process of normal ordering we obtain
$`\mathrm{\Delta }\theta _\alpha ^{B1}={\displaystyle \frac{3ig^2}{16r^5}}\left(2(ϵ\theta )(r/\theta )_\alpha 2(ϵr/\theta )\theta _\alpha 2r_l(ϵ\gamma ^m\theta )(\theta \gamma ^{ml})_\alpha \right).`$ (A.12)
Now, in distinction to all the other contributions, the one from (B2) involves propagation of massless fields $`\widehat{\psi }_\alpha `$, $`C_{11}`$ and $`\overline{C}_{11}`$ in the intermediate steps. The original expression is, to the order of interest,
$`\mathrm{\Delta }_ϵ\theta _\alpha ^{B2}=gϵ_\beta {\displaystyle 𝑑\tau ^{}C\mathrm{\Psi }^{}(\tau )\overline{C}_{11}\dot{\widehat{\psi }}_\beta (\tau ^{})}+gϵ_\beta {\displaystyle 𝑑\tau ^{}C^{}\mathrm{\Psi }(\tau )\overline{C}_{11}\dot{\widehat{\psi }}_\beta (\tau ^{})}.`$ (A.13)
To extract $`𝒪(\theta ^2)`$ contributions, we need to use the following five types of vertices:
$`V_1`$ $``$ $`ig{\displaystyle 𝑑\tau \left(r^m\overline{C}^{}C_{11}Y^m+r^m\overline{C}Y^mC_{11}\right)},`$ (A.14)
$`V_2`$ $``$ $`{\displaystyle 𝑑\tau \theta \gamma ^n(\mathrm{\Psi }Y^mY^m\mathrm{\Psi }^{})},`$ (A.15)
$`V_3`$ $``$ $`g{\displaystyle 𝑑\tau (\overline{C}^{}\dot{C}_{11}\stackrel{~}{A}+\overline{C}^{}C_{11}\dot{\stackrel{~}{A}}\overline{C}\dot{\stackrel{~}{A}^{}}C_{11}\overline{C}\stackrel{~}{A}^{}\dot{C}_{11})},`$ (A.16)
$`V_4`$ $``$ $`i{\displaystyle 𝑑\tau \theta (\mathrm{\Psi }\stackrel{~}{A}^{}\stackrel{~}{A}\mathrm{\Psi }^{})},`$ (A.17)
$`V_5`$ $``$ $`{\displaystyle \frac{1}{g^2}}{\displaystyle 𝑑\tau \dot{\widehat{\psi }}\theta }.`$ (A.18)
First, inserting $`V_1,V_2`$ and $`V_5`$, we get the diagrams of the type (B2-a). This gives the contribution
$`\mathrm{\Delta }_ϵ\theta _\alpha ^{B2a}=2ig^2{\displaystyle 𝑑\tau ^{}𝑑\tau ^{\prime \prime }\tau |\mathrm{\Delta }|\tau ^{\prime \prime }\tau ^{\prime \prime }|\frac{1}{}|\tau ^{}r^n(\tau ^{\prime \prime })\tau ^{\prime \prime }|\mathrm{\Delta }(\gamma ^n\theta )_\alpha \mathrm{\Delta }|\tau (ϵ\theta )(\tau ^{})}.`$
Similarly, use of $`V_3,V_4`$ and $`V_5`$ generates the diagrams of the type (B2-b), the contribution of which is worked out to be
$`\mathrm{\Delta }_ϵ\theta _\alpha ^{B2b}=2ig^2\tau |\mathrm{\Delta }(ϵ\theta )\mathrm{\Delta }r^m(\gamma ^m\theta )_\alpha \mathrm{\Delta }|\tau \mathrm{\Delta }_ϵ\theta _\alpha ^{B2a}.`$ (A.20)
Therefore,
$`\mathrm{\Delta }_ϵ\theta _\alpha ^{B2}=\mathrm{\Delta }_ϵ\theta _\alpha ^{B2a}+\mathrm{\Delta }_ϵ\theta _\alpha ^{B2b}=2ig^2(ϵ\theta )\tau |\mathrm{\Delta }^3|\tau (r/\theta )_\alpha ={\displaystyle \frac{3ig^2}{16r^5}}2(ϵ\theta )(r/\theta )_\alpha .`$ (A.21)
Calculation of $`\mathrm{\Delta }_ϵ\theta _\alpha ^C(\tau )`$: Finally, consider $`\mathrm{\Delta }_ϵ\theta _\alpha ^C`$. It takes the form
$`\mathrm{\Delta }_ϵ\theta _\alpha ^C(\tau )`$ $`=`$ $`g^2{\displaystyle }d\tau ^{}(C(\tau )\overline{C}^{}(\tau ^{})Y_m(\tau ^{})\mathrm{\Psi }_\alpha ^{}(\tau )`$ (A.22)
$`+C^{}(\tau )\overline{C}(\tau ^{})Y_m^{}(\tau ^{})\mathrm{\Psi }_\alpha (\tau ))ϵ^T\gamma ^m\theta (\tau ^{}),`$
which is represented by the diagram $`(C)`$. Using the vertex (A.7) and proceeding similarly to the previous calculations, we obtain
$`\mathrm{\Delta }\theta _\alpha ^C={\displaystyle \frac{3ig^2}{16r^5}}\left(2(ϵr/\theta )\theta _\alpha 2r_l(ϵ\gamma ^m\theta )(\theta \gamma ^{ml})_\alpha \right).`$ (A.23)
Summary: Adding up all the contributions and reinstating the factor of $`N`$, the final result is
$`\mathrm{\Delta }_ϵ\theta _\alpha `$ $`=`$ $`{\displaystyle \frac{3ig^2N}{16r^5}}(r_l(\theta \gamma ^{mnl}\theta )(\gamma ^{nm}ϵ)_\alpha +2r_l(\theta \gamma ^{ml}\theta )(\gamma ^mϵ)_\alpha 4r_l(ϵ\gamma ^m\theta )(\theta \gamma ^{ml})_\alpha `$ (A.24)
$`+4(ϵ\theta )(r/\theta )_\alpha 4(ϵr/\theta )\theta _\alpha ).`$
Appendix B: $`SO(9)`$ Fierz Identities
In this appendix, we record the $`SO(9)`$ Fierz identities which are crucial in simplifying the $`𝒪(\theta ^2)`$ part of $`\mathrm{\Delta }_ϵ\theta _\alpha `$ at 1-loop.
Adapting the notations of Taylor and Raamsdonk, let us introduce the following quantities for $`n=0,1,2,3`$ (repeated indices are summed):
$`E_n`$ $`=`$ $`r_m(ϵ\gamma ^{a_1\mathrm{}a_nm}\theta )(\theta \gamma ^{a_n\mathrm{}a_1}\lambda ),`$ (B.1)
$`\overline{E}_n`$ $`=`$ $`r_m(ϵ\gamma ^{a_1\mathrm{}a_n}\theta )(\theta \gamma ^{a_n\mathrm{}a_1m}\lambda ),`$ (B.2)
$`F_n`$ $`=`$ $`r_m(\theta \gamma ^{a_1\mathrm{}a_nm}\theta )(\lambda \gamma ^{a_n\mathrm{}a_1}ϵ),`$ (B.3)
$`\overline{F}_n`$ $`=`$ $`r_m(\theta \gamma ^{a_1\mathrm{}a_n}\theta )(\lambda \gamma ^{a_n\mathrm{}a_1m}ϵ),`$ (B.4)
$`P`$ $`=`$ $`{\displaystyle \frac{1}{4!}}r_mϵ^{a_1a_2a_3a_4b_1b_2b_3b_4m}(ϵ\gamma ^{a_1a_2a_3a_4}\theta )(\theta \gamma ^{b_1b_2b_3b_4m}\lambda ),`$ (B.5)
$`Q`$ $`=`$ $`{\displaystyle \frac{1}{4!}}r_mϵ^{a_1a_2a_3a_4b_1b_2b_3b_4m}(\theta \gamma ^{a_1a_2a_3a_4}\theta )(\lambda \gamma ^{b_1b_2b_3b_4m}ϵ),`$ (B.6)
where $`\lambda `$ is an arbitrary spinor. Because $`\theta \mathrm{\Gamma }\theta =0`$ for any symmetric matrix $`\mathrm{\Gamma }`$, five of them actually vanish, namely
$`F_0=\overline{F}_0=\overline{F}_1=F_3=Q=0.`$ (B.7)
Since there are nine independent Fierz identities, only four structures are independent, which we take to be $`F_1,F_2,\overline{F}_2`$ and $`\overline{F}_3`$. Then the remaining quantities can be expressed in terms of them as<sup>6</sup><sup>6</sup>6The sign in front of $`48E_0^{}`$ in Eq. (B.20) of should be $`+`$.
$`E_0`$ $`=`$ $`{\displaystyle \frac{1}{16}}F_1{\displaystyle \frac{1}{32}}(F_2+\overline{F}_2){\displaystyle \frac{1}{96}}\overline{F}_3,`$ (B.8)
$`\overline{E}_0`$ $`=`$ $`{\displaystyle \frac{1}{16}}F_1{\displaystyle \frac{1}{32}}(F_2+\overline{F}_2)+{\displaystyle \frac{1}{96}}\overline{F}_3,`$ (B.9)
$`E_1`$ $`=`$ $`{\displaystyle \frac{3}{8}}F_1{\displaystyle \frac{1}{8}}(F_2\overline{F}_2){\displaystyle \frac{1}{48}}\overline{F}_3,`$ (B.10)
$`\overline{E}_1`$ $`=`$ $`{\displaystyle \frac{3}{8}}F_1+{\displaystyle \frac{1}{8}}(F_2\overline{F}_2){\displaystyle \frac{1}{48}}\overline{F}_3,`$ (B.11)
$`E_2`$ $`=`$ $`{\displaystyle \frac{7}{4}}F_1{\displaystyle \frac{1}{4}}(F_2+\overline{F}_2)+{\displaystyle \frac{1}{24}}\overline{F}_3,`$ (B.12)
$`\overline{E}_2`$ $`=`$ $`{\displaystyle \frac{7}{4}}F_1{\displaystyle \frac{1}{4}}(F_2+\overline{F}_2){\displaystyle \frac{1}{24}}\overline{F}_3,`$ (B.13)
$`E_3`$ $`=`$ $`{\displaystyle \frac{21}{4}}F_1+{\displaystyle \frac{3}{4}}(F_2\overline{F}_2)+{\displaystyle \frac{3}{8}}\overline{F}_3,`$ (B.14)
$`\overline{E}_3`$ $`=`$ $`{\displaystyle \frac{21}{4}}F_1{\displaystyle \frac{3}{4}}(F_2\overline{F}_2)+{\displaystyle \frac{3}{8}}\overline{F}_3,`$ (B.15)
$`P`$ $`=`$ $`{\displaystyle \frac{15}{2}}(F_2+\overline{F}_2).`$ (B.16)
¿From these relations, one easily finds the identity
$`0=F_22F_14\overline{E}_1+4\overline{E}_0+12E_0.`$
Removing $`\lambda `$, we get
$`0`$ $`=`$ $`r_l(\theta \gamma ^{mnl}\theta )(\gamma ^{nm}ϵ)_\alpha +2r_l(\theta \gamma ^{nl}\theta )(\gamma ^nϵ)_\alpha +4r_l(ϵ\gamma ^n\theta )(\gamma ^{nl}\theta )_\alpha `$ (B.17)
$`+4r_l(ϵ\theta )(\gamma ^l\theta )_\alpha +12r_l(ϵ\gamma ^l\theta )\theta _\alpha ,`$
which was used in Sec. 5.
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# Automatic Classification of Text Databases Through Query Probing
## 1 Introduction
Text databases abound on the Internet. Sometimes users can browse through their documents by following hyperlinks. In many other cases, text databases are “hidden” behind search interfaces, and their documents are only available through querying. For those databases, web search engines cannot crawl inside, and they just index the “front pages”, ignoring the contents of possibly rich sources of information. One example of such a search-only text database is the archive of a newspaper. Many newspapers do not offer a browsing interface for past issues, but they do offer search capabilities to retrieve old articles. This is the case, for example, for the New York Times newspaper.
One way of facilitating the access to such searchable databases is to build metasearchers. A metasearcher sends user queries to many search engines, retrieves and merges the results and then returns the combined results back to the user (e.g., \[GGMT99, MLY<sup>+</sup>98, XC98, SE95, GCGMP97, LG98\]). Alternatively, users can browse Yahoo-like directories to locate databases of interest and then submit queries to these databases. Some sites have started in the last two years to provide such services. For example, InvisibleWeb<sup>1</sup><sup>1</sup>1http://www.invisibleweb.com/ and SearchEngine Guide<sup>2</sup><sup>2</sup>2http://www.searchengineguide.com/ classify various search engines into a hierarchical classification scheme. The archive of the New York Times is classified as: News $``$ Regional $``$ National (USA) $``$ New York $``$ News. A user can then locate relevant text databases and submit queries only to them to obtain more accurate and focused results than when searching a more general text database. Other services (e.g., Copernic<sup>3</sup><sup>3</sup>3http://www.copernic.com/) combine the metasearching approach with “browsing”. Users can select a specific category (e.g., Recipes, Newspapers, etc.) and the metasearcher then sends the user queries to the searchable databases previously classified in the given category.
Unfortunately, existing approaches for text database classification involve manual intervention of a human expert and do not scale. In this paper we will describe a way of automating this classification process by issuing query probes to the text databases. More specifically, in Section 2 we define what it means to classify a text database. Then, in Section 3 we focus on the design of our query probing classification strategy. Finally, in Section 4 we present some initial experiments over web databases.
#### Related Work
Query probing has been used recently for characterizing different properties of text databases. Manually constructed query probes have been used in \[GWG96\] for the classification of text databases. \[CCD99\] probe text databases with queries to determine an approximation of their vocabulary and associated statistics. This technique requires retrieving the documents in the query results for further analysis. Finally, guided query probing has been used in \[MYL99\] to determine sources of heterogeneity in the algorithms used to index and search locally at each text database.
## 2 Classification of Text Databases
In this section we will describe two basic approaches for classifying text databases. One approach classifies a database into one category when the database contains a substantial number of documents in this category. The other approach classifies a database into one category when the majority of its documents are in this category.
###### Example 1
: Consider two databases $`D_1,D_2`$ with 1,000 and 10,000,000 documents, respectively, and a topic category “Health.” Suppose that $`D_1`$ contains 900 documents about health while $`D_2`$ contains 200,000 such documents. Our decision whether to classify $`D_1`$ and $`D_2`$ in the “Health” category will ultimately depend on how users will take advantage of our classification and the databases. Some users might prefer a “focus-oriented” classification (i.e., might be looking for text databases having mostly documents about health and little else). Such users might not want to process documents outside of their topic of interest, and might then prefer that database $`D_1`$ be classified in the “Health” category (90% of its documents are on health). In contrast, $`D_2`$ should not be classified in that category. Although $`D_2`$ has a large number of document on health, these documents represent only a small fraction of the database (i.e., 2%). Hence, it is likely that our “focus-oriented” users would be exposed to non-health documents while exploring $`D_2`$. Alternatively, other users might be looking for text databases having a sufficiently large number of documents on health. It might be unimportant for such users what else is at each database. These users might then prefer $`D_2`$ to be classified in the “Health” category because of its large number of documents on health (i.e., 200,000). $`D_1`$ (with 900 documents on health) might or might not be classified in that category, depending on what we consider a “sufficient large” number of documents.
Consider a set of categories $`C_1,\mathrm{},C_k`$ and a text database $`D`$ that we want to classify in one or more of these categories. Each of $`D`$’s documents has been classified in one of the categories $`C_1,\mathrm{},C_k`$ that we use to classify $`D`$. Given this classification of the documents in $`D`$ we can compute a vector $`C=(n_1,\mathrm{},n_k)`$, which indicates the number of documents $`n_i`$ in category $`C_i`$, for $`i=1,\mathrm{},k`$. Vector $`C`$ is a good summary of the contents of database $`D`$ and we will use it to classify the database, as we describe next. As illustrated in Example 1 above, to categorize databases we need to capture how “focused” $`D`$ is and how many documents it contains for a given category. For this we define the following two metrics.
###### Definition 1
: Consider a text database $`D`$ and a category $`C_i`$. Then the coverage of $`D`$ for $`C_i`$ is the number of documents in $`D`$ in category $`C_i`$, $`n_i`$. The specificity of $`D`$ on $`C_i`$ is the fraction of documents in $`D`$ in category $`C_i`$:
$`\text{Coverage}(D,C_i)`$ $`=`$ $`n_i`$
$`\text{Specificity}(D,C_i)`$ $`=`$ $`{\displaystyle \frac{n_i}{|D|}}`$
*Specificity* defines how “focused” a database is on a given category. One problem with the definition above is that we do not always know the number of documents in a database. We will discuss how we can approximate this value in Section 3. *Coverage* defines the “absolute” amount of information that a database contains about a specific category. An alternative definition for coverage could divide $`n_i`$ by the total number of documents in all databases. This would capture what fraction of the existing documents in category $`C_i`$ are present in a given database. Although this definition is interesting, it has the undesirable property of depending on a universe of known databases. On the Internet, databases come and go constantly so this definition would make the resulting classification scheme that we describe quite unstable. Moreover, since the *Coverage* value would have the same normalizing constant for all databases, excluding this factor will have no bearing on the relative ranking of databases by their coverage of a certain topic.
Using the definitions above, each database $`D`$ has a specificity and a coverage value for each category. We can use these values to decide how to classify $`D`$ into one or more of the categories. As described above, we could classify a database into one category when the majority of the documents it contains are of a specific category. Our classification could alternatively be based on the number of documents of a specific category that a database contains.
###### Definition 2
: Consider a database $`D`$ and a category $`C_i`$ and let $`\tau _s,\tau _c0`$ to be two pre-specified thresholds. Then $`D`$ is in category $`C_i`$ according to a “coverage-oriented” classification if Coverage$`(D,C_i)\tau _c`$. Similarly, $`D`$ is in category $`C_i`$ according to a “specificity-oriented” classification if Specificity$`(D,C_i)\tau _s`$.
###### Example 1
: (cont.) Consider the two databases $`D_1,D_2`$ described above, and the category “Health.” Using Definition 2, Coverage$`(D_1,`$“Health”$`)=900`$, since $`D_1`$ has 900 documents on health. Similarly, Coverage$`(D_2,`$“Health”$`)=200,000`$. If threshold $`\tau _c`$ for our “coverage-oriented” classification is set to, say, 10,000, then $`D_2`$ will be classified in category “Health” while $`D_1`$ will not, since it does not have a sufficient large number of documents in this category. Analogously, Specificity$`(D_1,`$“Health”$`)=\frac{900}{1000}=0.9`$ while Specificity$`(D_2,`$“Health”$`)=\frac{200,000}{10,000,000}=0.02`$. If threshold $`\tau _s`$ for our “specificity-oriented” classification is set to, say, 0.3 then $`D_1`$ will be classified in category “Health” while $`D_2`$ will not, since it is not sufficiently focused on health and holds too many documents in other categories.
The two alternative database classification schemes above assume that we somehow know the number of documents that each database has in each category, which is clearly unrealistic in most Internet settings. In effect, as discussed in the Introduction, many times we do not have access to a database’s contents other than through a query interface. In the next section we introduce techniques for approximating the classification of text databases in this limited-access scenario.
## 3 Classifying Databases through Probing
The previous section described how to classify a database given the number of documents it contains in each of our categories. Unfortunately, text databases do not export such metadata. In this section we introduce a technique to classify text databases in the absence of any information about their contents. Our technique starts by training a rule-based document classifier over our categories (Section 3.1) and then uses the classifier’s rules to design a set of probing queries (Section 3.2). The database will be classified based on the number of matches returned for each of these queries, without accessing the documents per se (Section 3.3).
### 3.1 Training a Document Classifier
Our technique for classifying databases over a set of categories $`C_1,\mathrm{},C_k`$ starts by training a rule-based document classifier over those categories. We use RIPPER, an off-the-shelf tool developed at AT&T Research Laboratories\[Coh95, Coh96\]. Given a set of training, pre-classified documents, this tool returns a classifier that might consist of rules like the following:
| Computers | IF mac |
| --- | --- |
| Computers | IF graphics windows |
| Religion | IF god christian |
| Hobbies | IF baseball |
The first rule indicates that if a document contains the term mac it should be classified in the *“Computers”* category. A document should also be classified into that category if it has the words graphics and windows. Similarly, if a document has the words god and christian, it is a *“Religion”* document, whereas if it has the word baseball, it is a *“Hobbies”* document.
Once we have trained a document classifier using a tool like RIPPER, we could apply it to every document in a database $`D`$ that we want to classify. This procedure would produce a close approximation to the $`C=(n_1,\mathrm{},n_k)`$ vector of category frequencies for $`D`$ (Section 2), which we could use to classify $`D`$ according to Definition 2. Unfortunately, we often do not have access to all the documents in a database, other than indirectly through a query interface, as discussed above. Next, we define a query probing strategy to deal with such databases.
### 3.2 Probing a Database
Our goal is to create a set of queries for each category that will retrieve exactly the documents for that category from the database we are classifying. We will construct these queries based on the document classifier discussed above. To create our queries, we turn each rule into a query. The number of matches for each query will be the number of documents in the database that satisfy the corresponding rule. These numbers will then be used to approximate the distribution of documents in categories within a text database, as the following example illustrates.
###### Example 2
: Consider a database $`D`$ with 500 documents, all about *“Computers,”* and suppose that our categories of interest are *“Computers,” “Hobbies,”* and *“Religion.”* Then $`D`$ has associated with it a vector $`C=(500,0,0)`$ (Section 2), showing the distribution of documents over these three categories. Suppose also that we have trained a rule-based document classifier and obtained the four rules shown above for the three categories. If we do not have access to all the documents of $`D`$, we can still characterize its contents by issuing probing queries constructed from the document classifier as discussed above. Our first probe will be the query mac. The database will return a result of the form “92 documents found.” We send a second query graphics AND windows. Again, we get a result like “288 documents found.” Queries god AND christian and baseball return 0 and 2 matches respectively. From these results we conclude that $`D`$ has 288+92=380 “Computers” documents, 0 “Religion” documents, and 2 “Hobbies” documents. Thus we approximate the ideal vector $`C`$, with $`C^{}=(380,0,2)`$.
RIPPER can produce either an ordered set of rules or an unordered set of rules. When the rules are ordered, the first rule that is satisfied by a document fires and gives a classification for that document. No subsequent rules are matched against that document. We should formulate our queries properly in order to simulate the actions of the classifier as much as possible. For example, if the rules above were ordered rules, our second probing query would have been graphics AND windows AND NOT mac, to avoid retrieving any documents that would match the first, earlier rule.
If the query interface of a database does not support the kind of queries described above, we break these queries into smaller pieces that we can send separately. A detailed description of this technique is beyond the scope of this paper. For completeness, we mention that we submit the probing queries in such a way that we can use the inclusion-exclusion principle to calculate the number of results that would have been returned for the original queries.
A significant advantage of our probing approach is that we do not need to retrieve documents to analyze the contents of a database \[CCD99\]. Instead, we count only the number of matches for these queries. Thus, in our approach we only require a database to report the number of matches for a given query. It is common for a database to return something like “$`X`$ documents found” before returning the actual results.
### 3.3 Using the Probing Results for Classification
After the probing phase, we have calculated an approximation of the coverage of a database for our categories. To calculate the specificity values, we would need the size of the database $`|D|`$, and we approximate it by $`|D|_{i=1}^kn_i`$. This means that we will use only the documents that are classified into the given categories to calculate the size of the database. This approach can give poor results when there are many documents that do not belong to any of the given categories. In such a case, it is also difficult to categorize this text database into the given classification scheme, since no category will accurately reflect its contents.
An extra step that we applied to our method to improve the results is the following. For each of the rules, we know the accuracy from the training phase of the classifier. For example, the rule Computers IF mac may have correctly classified 90 documents and incorrectly classified 10 other documents during the training phase, resulting in an accuracy of 0.9. We adjust our results from the probing phase by multiplying the number of documents matched by each rule by the accuracy of that rule. Also, for the set of rules that classified documents into one category, we know their “recall”, i.e., how many documents they recalled over all the documents in this category. For example if category Computers in the training phase had 150 documents and the rules retrieved 100, then the recall is 0.67. This means that only this portion of all the documents of this category were retrieved. To adjust our results further, we divide each element of the $`C^{}`$ approximation vector with the recall for this category. This regularization of the values $`n_i`$ helps account for the fact that rules generally do not (and need not) have perfect recall on real document databases.
## 4 Initial Experiments
Using RIPPER, we created a classifier using a collection of 20,000 newsgroup articles from the UCI KDD archive<sup>4</sup><sup>4</sup>4http://kdd.ics.uci.edu/. This collection has been used in previous text categorization experiments \[Joa96, Mit97\], and is composed of 1,000 newsgroup articles from each of 20 newsgroups. We further grouped the articles into five large groups according to their originating newsgroups: Computers (comp.\*), Science (sci.\*), Hobbies (rec.\*), Society (alt.atheism, talk.\*, soc.\*) and Misc (misc.sale). We have removed all the headers (except for the “Subject:” line), the e-mail addresses from the body of the articles and all punctuation. Subsequently, we eliminated all words that appeared in fewer than 3 documents in the collection and the 100 most frequent words. Such feature reduction is in accordance with Zipf’s Law \[Zip49\], which shows that there are many infrequently used words in document collections. For purposes of classification, however, such infrequent terms generally provide little discriminating power between classes (due to their rarity), and can thus be safely eliminated with little, if any, reduction in subsequent classification accuracy. Similarly, very frequent words, that often tend to appear in virtually all articles, will also provide little ability to make classification distinctions, and can likewise be eliminated. After this step we applied an information theoretic feature selection algorithm \[KS96, KS97\] to reduce the terms from about 40,000 to 5,000. This algorithm eliminates features that have the least impact on the class distribution of documents (as measured by the relative entropy of the distribution of the document class labels conditioned on the appearance of a given feature). Features that have little impact on the class distribution are likely to also have little discriminating power between classes, and can thus be eliminated without much adverse impact on the final classification accuracy. For training set we used a random sample of 10,000 documents and the remaining 10,000 documents were used for testing.
The initial document classifier generated by RIPPER consisted of 534 ordered rules. Many of the rules were covering very few (one or two) examples from the training set. These rules did not contribute much to the overall accuracy of the document classifier, and would result in too many probing queries during the classification stage. Thus, we decided to restrict the classifier to produce only rules that covered at least 50 examples from the training set. This resulted in a classifier with 29 ordered rules that included a total of 32 words. We also tried to produce a rule set that would include rules with negations (NOT clauses). The resulting classifier had 31 rules with much better accuracy, but, in this case, a total of 92 words were used to form the rules. The queries for this classifier were much longer and we opted to use the simpler classifier (that had only 29 rules and 32 words) for the sake of query efficiency. The rules given in Section 3.1 are, in fact, examples of rules used by this classifier.
After constructing the classifier, we have selected four sites from InvisibleWeb<sup>5</sup><sup>5</sup>5http://www.invisibleweb.com/ to test our method. These four sites are topically cohesive, and should be classified in the same category by both the specificity- and the coverage- oriented classification alternatives of Definition 2:
* Cora<sup>6</sup><sup>6</sup>6http://www.cora.jprc.com/: A repository of technical papers about Computer Science. This database should be classified under the category “Computers.”
* American Scientist<sup>7</sup><sup>7</sup>7http://www.amsci.org/: An on-line version of a magazine on science and technology. This repository should be classified under category “Science.”
* AllOutdoors<sup>8</sup><sup>8</sup>8http://www.alloutdoors.com/: A site with articles about fishing, hunting, and other outdoor activities. This site should be classified under category “Hobbies.”
* ReligionToday<sup>9</sup><sup>9</sup>9http://www.religiontoday.com/: A site with news and discussions about religion. This site should be classified under category “Society.”
We probed these sites using the techniques described in Section 3.2. One problem that arose during the probing phase was a limitation on the length of the queries that we could submit to the “American Scientist” site. We truncated the long queries by eliminating terms that did not cover any documents (e.g., instead of issuing a query baseball AND NOT god, if the query god returned 0 results, we issue only the query baseball).
The results of our probing phase can be seen in Figure 1. Consider, for example, the results for Cora. After submitting the queries for the class Computer, the database reported 1450 matches for all the queries. For classes Science, Hobbies, Society, and Misc, it reported 151, 95, 215, and 45 respectively. Using these coverage values we estimated specificity as in Section 2. The specificity values are depicted using the bars, and it can be clearly seen that the results indicate that Cora is a site that is “focused” on Computers. Similarly for the other sites, we probed them using the same rules. The results clearly indicate the focus of each site. For example, if we had a threshold value for specificity of $`\tau _s=0.6`$, then each site would be classified correctly. Moreover, to measure the significance of our results, we performed a Chi-squared test comparing the distribution of the classes for each database given by the probes to the uniform distribution. This test gives us a measure for how likely the skew in the class distribution (toward the correct class) is likely to have been gotten by chance. The Chi-square test reveals that the skews in the class distributions for each database are significant at the 99.9% level. Thus, it appears that, in every case, the probes generated by the RIPPER rules have accurately captured the concept represented by each class of documents.
## 5 Conclusions and Future Work
In this paper, we have described a method that uses probing queries produced by a classifier to classify a text database. We have also shown some promising initial experiments. The method managed to identify the right category for each database, using only the number of matches for a small set of queries and without retrieving any documents. Our technique could also be used to characterize web sites that offer a browsable interface as well. The only requirement is the existence of a search interface for the local contents, which many sites offer. By using only a small set of probe queries, we can get a coarse idea about the contents of a web site.
Our future work includes the expansion of our strategy into a hierarchical classification scheme. We will also explore the efficiency of our algorithm for various indexing environments and for search interfaces that support different sets of boolean operators. We also plan to compare our approach against an adaptation, for the database classification problem, of the technique in \[CCD99\]. Finally, we will expand our adjustment technique (that currently uses only the precision of each rule and the recall of each category) to use the full set of statistics (i.e., confusion matrices) from the document classifier. This could produce better approximations of the contents of the search-only text databases.
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# 1 Introduction
## 1 Introduction
The violation of CP was first observed in the neutral kaon system , where a deviation from the superweak theory has recently been confirmed , and CP violation in $`B`$-meson decays is strongly suggested by recent experiments . In addition to its interest for particle physics, CP non-conservation provides a key ingredient for cosmological baryogenesis, namely for explaining the underlying mechanism which caused matter to dominate over anti-matter in our observable universe . Although a fundamental understanding of the origin of CP violation is still lacking, most of the scenarios proposed in the existing literature indicate that Higgs interactions play a key role in mediating CP violation. For instance, CP violation is broken explicitly in the Standard Model (SM) by complex Yukawa couplings of the Higgs boson to quarks. Another appealing scheme of CP violation occurs in models with an extended Higgs sector, in which the CP symmetry of the theory is broken spontaneously by the ground state of the Higgs potential . Supersymmetric (SUSY) theories, including the minimal supersymmetric extension of the Standard Model (MSSM), predict an extended Higgs sector, and therefore may realize either or both the above two schemes of explicit and spontaneous CP violation. However, the Higgs potential of the MSSM is invariant under CP at the tree level, and any explicit or spontaneous breakdown of CP symmetry can arise only via radiative corrections. The case of purely spontaneous CP violation in the MSSM leads to an unacceptably light CP-odd scalar , as a result of the Georgi-Pais theorem , and hence such a scenario is ruled out experimentally.
It has recently been shown that the tree-level CP invariance of the MSSM Higgs potential may be violated sizeably by loop effects involving soft CP-violating trilinear interactions of the Higgs bosons to top and bottom squarks. A detailed study has shown that significant CP-violating effects of level crossing in the Higgs sector can take place in such a minimal SUSY scenario of explicit CP violation, which may lead to drastic modifications of the tree-level couplings of the Higgs particles to fermions and to the $`W^\pm `$ and $`Z`$ bosons . The latter can have important phenomenological consequences on the production rates of the lightest Higgs particle, even though the upper bound on its mass was found to be very similar to that found previously in the CP-conserving case . The MSSM predicts an upper bound on the lightest Higgs boson mass of approximately 110 (130) GeV for small (large) values of the ratio of Higgs vacuum expectation values $`\mathrm{tan}\beta 2`$ $`(>15)`$. On the other hand, experiments at LEP 2, running at center-of-mass energies $`\sqrt{s}=`$ 196–202 GeV, have placed a severe lower bound of approximately 108 GeV on the mass of the SM Higgs boson . LEP 2 is expected to run at center-of-mass energies up to $`\sqrt{s}=206`$ GeV during the year 2000. Consequently, for stop masses smaller than 1 TeV, a significant portion of the parameter space spanned by $`\mathrm{tan}\beta `$ and the CP-odd scalar mass $`M_A`$ can be tested for the CP-conserving case in this next round of experiments at LEP 2 . However, the explorable region of parameters is smaller for larger amounts of stop mixing and/or larger CP-violating phases. A decisive test of such a scenario can only be provided by the upgraded Tevatron collider and the LHC.
The earlier study of the renormalization-group (RG) improved effective potential of the MSSM with explicit CP violation in was based on an expansion of the Higgs quartic couplings in inverse powers of the arithmetic average of stop and sbottom masses, under the assumption that the mass splittings of the left- and right-handed stop and sbottom masses are small. The mass expansion of the one-loop effective potential was truncated up to renormalizable operators of dimension four. Although the above approach captures the basic qualitative features of the underlying dynamics under study, it is known that such a mass expansion has limitations when the third-generation squark mixing is large. Since the dominant CP-violating loop contributions to the effective Higgs-boson masses and mixing angles occur for large values of the third-generation squark-mixing parameters, it is necessary to provide a more complete one-loop computation of the effective MSSM Higgs potential with explicit radiative breaking of CP invariance, including non-renormalizable operators below the heavy scalar-quark scale and without resorting to any other kinematic approximations.<sup>1</sup><sup>1</sup>1We recall that the diagrammatic computation of scalar-pseudoscalar transitions is the focus of , whilst sbottom contributions and relevant $`D`$-term effects on the effective potential are not considered in .
To this end, we consider here the two-loop leading logarithms induced by top- and bottom-quark Yukawa couplings as well as those associated with QCD corrections, by means of RG methods. In the calculation of the RG-improved effective potential, we also include the leading logarithms generated by one-loop gaugino and higgsino quantum effects . Finally, we implement the potentially large two-loop corrections that are induced by the one-loop stop and sbottom thresholds in the top- and bottom-quark Yukawa couplings, which may become particularly relevant in the large-$`\mathrm{tan}\beta `$ regime. On the basis of the RG-improved effective potential, we present predictions for the Higgs-boson mass spectrum and the effective Higgs couplings to fermions and gauge bosons. The results of the analysis are compared with those obtained in the CP-conserving case and also with the results obtained by truncating the one-loop RG-improved effective potential up to renormalizable operators .
In this analysis, it is important to consider the constraints on the low-energy CP-violating parameters of the MSSM that originate from experimental upper limits on the electron and neutron electric dipole moments (EDMs) . Most of the EDM constraints affect the CP-violating couplings of the first two generations . Thus, making the first two generations of squarks rather heavy, much above the TeV scale , is a possibility that can drastically reduce one-loop contributions to the neutron EDM, without suppressing the CP-violating phases of the theory. Another interesting possibility for avoiding any possible CP crisis in the MSSM is to arrange for cancellations among the different EDM terms, either at the level of short-distance diagrams or (for the neutron EDM) at the level of the strong-interaction matrix elements for operators with $`s`$, $`u`$ and $`d`$ quark flavours . Alternatively, one might require a specific form of non-universality in the soft trilinear Yukawa couplings .
However, third-generation squarks can also give rise by themselves to observable effects on the electron and neutron EDMs through the three-gluon operator , through the effective coupling of the ‘CP-odd’ Higgs boson to the gauge bosons , and through two-loop gaugino/higgsino-mediated EDM graphs . These different EDM contributions of the third generation can also have different signs and add destructively to the electron and neutron EDMs. In our phenomenological discussion, we take into account the relevant EDM constraints related to the CP-violating parameters of the stop and sbottom sectors.
This paper is organized as follows: in Section 2 we calculate the complete one-loop CP-violating effective potential, and derive the analytic expressions for the effective charged and neutral Higgs-boson mass matrices. Technical details are given in Appendices A and B. Section 3 describes our approach to determining the RG-improved Higgs-boson mass matrices, after including the leading two-loop logarithms associated with Yukawa and QCD corrections. Section 4 is devoted to the calculation of the effective top- and bottom-quark Yukawa couplings, in which one-loop threshold effects of the third-generation squarks are implemented. In Section 5, we discuss the phenomenological implications of representative CP-violating scenarios compatible with EDM constraints for direct Higgs searches at LEP 2 and the upgraded Tevatron collider. We also compare the results of our analysis with those obtained using the mass-expansion method . Finally, in Section 6 we summarize our conclusions.
## 2 CP-Violating One-Loop Effective Potential
In this Section, we first describe the basic low-energy structure of the MSSM that contains explicit CP-violating sources, such as soft CP-violating trilinear interactions. Then we calculate the general one-loop CP-violating effective potential. Finally, after implementing the minimization tadpole conditions related to the Higgs ground state, we derive the effective charged and neutral Higgs-boson mass matrices.
CP violation is introduced into the MSSM through the Higgs superpotential and the soft supersymmetry-breaking Lagrangian:
$`W`$ $`=`$ $`h_l\widehat{H}_1^Ti\tau _2\widehat{L}\widehat{E}+h_d\widehat{H}_1^Ti\tau _2\widehat{Q}\widehat{D}+h_u\widehat{Q}^Ti\tau _2\widehat{H}_2\widehat{U}\mu \widehat{H}_1^Ti\tau _2\widehat{H}_2,`$ (2.1)
$`_{\mathrm{soft}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(m_{\stackrel{~}{g}}\lambda _{\stackrel{~}{g}}^a\lambda _{\stackrel{~}{g}}^a+m_{\stackrel{~}{W}}\lambda _{\stackrel{~}{W}}^i\lambda _{\stackrel{~}{W}}^i+m_{\stackrel{~}{B}}\lambda _{\stackrel{~}{B}}\lambda _{\stackrel{~}{B}}+\mathrm{h}.\mathrm{c}.)+\stackrel{~}{M}_L^2\stackrel{~}{L}^{}\stackrel{~}{L}+\stackrel{~}{M}_Q^2\stackrel{~}{Q}^{}\stackrel{~}{Q}`$ (2.2)
$`+\stackrel{~}{M}_U^2\stackrel{~}{U}^{}\stackrel{~}{U}+\stackrel{~}{M}_D^2\stackrel{~}{D}^{}\stackrel{~}{D}+\stackrel{~}{M}_E^2\stackrel{~}{E}^{}\stackrel{~}{E}+m_1^2\stackrel{~}{\mathrm{\Phi }}_1^{}\stackrel{~}{\mathrm{\Phi }}_1+m_2^2\mathrm{\Phi }_2^{}\mathrm{\Phi }_2(B\mu \stackrel{~}{\mathrm{\Phi }}_1^Ti\tau _2\mathrm{\Phi }_2`$
$`+\mathrm{h}.\mathrm{c}.)+(h_lA_l\mathrm{\Phi }_1^{}\stackrel{~}{L}\stackrel{~}{E}+h_dA_d\mathrm{\Phi }_1^{}\stackrel{~}{Q}\stackrel{~}{D}h_uA_u\mathrm{\Phi }_2^Ti\tau _2\stackrel{~}{Q}\stackrel{~}{U}+\mathrm{h}.\mathrm{c}.),`$
where $`\stackrel{~}{\mathrm{\Phi }}_1=i\tau _2\mathrm{\Phi }_1^{}`$ is the scalar component of the Higgs chiral superfield $`\widehat{H}_1`$ and $`\tau _2`$ is the usual Pauli matrix. The conventions followed throughout this paper, including the quantum-number assignments of the fields under the SM gauge group, are displayed in Table 1.
As can be seen from (2.1) and (2.2), the MSSM includes additional complex parameters with new CP-odd phases that are absent in the SM. These new CP-odd phases may reside in the following parameters: (i) the mass parameter $`\mu `$ describing the bilinear mixing of the two Higgs chiral superfields in the superpotential; (ii) the soft supersymetry-breaking gaugino masses $`m_{\stackrel{~}{g}}`$, $`m_{\stackrel{~}{W}}`$ and $`m_{\stackrel{~}{B}}`$ of the gauge groups SU(3)<sub>c</sub>, SU(2)<sub>L</sub> and U(1)<sub>Y</sub>, respectively; (iii) the soft bilinear Higgs-mixing mass $`B\mu `$; and (iv) the soft supersymmetry-breaking trilinear couplings $`A_f`$ of the Higgs bosons to sfermions. In addition, there may exist other large CP-odd phases, associated with flavor off-diagonal soft supersymmetry-breaking masses of squarks and sleptons. We assume that these off-diagonal masses are small and therefore do not give sizeable contributions to the effective Higgs potential . The number of independent CP-odd phases may be reduced if one prescribes a universality condition for all gaugino masses at the unification scale $`M_X`$; the gaugino masses will then have a common phase. Correspondingly, the different trilinear couplings $`A_f`$ may be considered to be all equal at $`M_X`$, i.e., $`A_fA`$. In this case, however, because of the different RG running of the phases of the trilinear couplings, their values at low energies will be different.<sup>2</sup><sup>2</sup>2For a discussion of the RG effects, see . Two CP-odd phases may further be eliminated by employing the following two global symmetries that govern the dimension-four operators in the MSSM Lagrangian:
* The U(1)<sub>Q</sub> symmetry specified by $`Q(\widehat{H}_1)=1`$, $`Q(\widehat{H}_2)=2`$, $`Q(\widehat{Q})=Q(\widehat{L})=0`$, $`Q(\widehat{U})=2`$ and $`Q(\widehat{D})=Q(\widehat{E})=1`$. This U(1)<sub>Q</sub> symmetry is broken by the $`\mu `$ parameter and the respective soft supersymmetry-breaking one, $`B\mu `$.
* The U(1)<sub>R</sub> symmetry acting on the Grassmann-valued coordinates, i.e., the $`\theta `$ coordinate of superspace carries charge 1. Under the $`R`$ transformation, the matter superfields and gaugino fields carry charge 1, whilst the Higgs superfields are $`R`$-neutral. The $`R`$ symmetry is violated by the gaugino masses, the trilinear couplings $`A_f`$ and the parameter $`\mu `$.
We concentrate on the parameters which may have a dominant CP-violating effect on the MSSM Higgs potential, under the assumption of a common phase for the gauginos; the latter is made less important by the fact that the one-loop gaugino corrections are subdominant compared to the ones induced by the third-generation squarks. As has been mentioned above, two CP-odd phases of the complex parameters $`\{\mu ,m_{12}^2,m_\lambda ,A\}`$ may be removed by employing the global symmetries (i) and (ii). Specifically, one of the Higgs doublets and the common phase of the gaugino fields can be rephased in a way such that the gaugino masses and $`B\mu `$ become real numbers. As a consequence, arg($`\mu `$) and arg($`A_{t,b}`$) are the only physical CP-violating phases in the MSSM which affect the Higgs sector in a relevant way.
It is obvious from (2.1) and (2.2) that the Yukawa interactions of the third-generation quarks, $`Q^T=(t_L,b_L)`$ and $`t_R,b_R`$, as well as their SUSY bosonic counterparts, $`\stackrel{~}{Q}^T=(\stackrel{~}{t}_L,\stackrel{~}{b}_L)`$ and $`\stackrel{~}{t}_R,\stackrel{~}{b}_R`$, play the most significant role in radiative corrections to the Higgs sector. Therefore, it is useful to give the interaction Lagrangians related to the $`F`$ and $`D`$ terms of the third generation:
$`_F`$ $`=`$ $`|h_b|^2|\mathrm{\Phi }_1^+\stackrel{~}{Q}|^2+|h_t|^2|\mathrm{\Phi }_2^Ti\tau _2\stackrel{~}{Q}|^2`$ (2.3)
$`(\mu h_b^{}\stackrel{~}{Q}^{}\mathrm{\Phi }_2\stackrel{~}{b}_R+\mu h_t^{}\stackrel{~}{Q}^{}i\tau _2\mathrm{\Phi }_1^{}\stackrel{~}{t}_R+\mathrm{h}.\mathrm{c}.)`$
$`\left(h_b^{}\stackrel{~}{b}_R\mathrm{\Phi }_1^Ti\tau _2+h_t^{}\stackrel{~}{t}_R\mathrm{\Phi }_2^{}\right)\left(h_bi\tau _2\mathrm{\Phi }_1^{}\stackrel{~}{b}_R^{}h_t\mathrm{\Phi }_2\stackrel{~}{t}_R^{}\right),`$
$`_D`$ $`=`$ $`{\displaystyle \frac{g_w^2}{4}}\left[\mathrm{\hspace{0.17em}2}|\mathrm{\Phi }_1^Ti\tau _2\stackrel{~}{Q}|^2+\mathrm{\hspace{0.17em}2}|\mathrm{\Phi }_2^{}\stackrel{~}{Q}|^2\stackrel{~}{Q}^{}\stackrel{~}{Q}(\mathrm{\Phi }_1^{}\mathrm{\Phi }_1+\mathrm{\Phi }_2^{}\mathrm{\Phi }_2)\right]`$ (2.4)
$`+{\displaystyle \frac{g^2}{4}}(\mathrm{\Phi }_2^{}\mathrm{\Phi }_2\mathrm{\Phi }_1^{}\mathrm{\Phi }_1)\left[\frac{1}{3}(\stackrel{~}{Q}^{}\stackrel{~}{Q})\frac{4}{3}(\stackrel{~}{t}_R\stackrel{~}{t}_R^{})+\frac{2}{3}(\stackrel{~}{b}_R\stackrel{~}{b}_R^{})\right].`$
where $`g_w`$ and $`g^{}`$ are the usual SU(2)<sub>L</sub> and U(1)<sub>Y</sub> gauge couplings. Further, the interaction Lagrangian of the Higgs bosons to the top and bottom quarks is given by
$$_{\mathrm{fermions}}=(h_b\overline{b}_RQ^T\mathrm{\Phi }_1^{}+\mathrm{h}.\mathrm{c}.)+(h_t\overline{t}_RQ^Ti\tau _2\mathrm{\Phi }_2+\mathrm{h}.\mathrm{c}.),$$
(2.5)
where $`h_t`$ and $`h_b`$ are the top and bottom Yukawa couplings, respectively.
With the help of the Lagrangians (2.3)–(2.5), we now proceed with the calculation of the one-loop effective potential. More explicitly, in the $`\overline{\mathrm{MS}}`$ scheme, the one-loop CP-violating effective potential is determined by
$$_V=_V^0+\frac{3}{32\pi ^2}\underset{q=t,b}{}\left[\underset{i=1,2}{}\stackrel{~}{m}_{q_i}^4\left(\mathrm{ln}\frac{\stackrel{~}{m}_{q_i}^2}{Q^2}\frac{3}{2}\right)\mathrm{\hspace{0.17em}2}\overline{m}_q^4\left(\mathrm{ln}\frac{\overline{m}_q^2}{Q^2}\frac{3}{2}\right)\right].$$
(2.6)
In (2.6), $`_V^0`$ is the tree-level Lagrangian of the MSSM Higgs potential
$`_V^0`$ $`=`$ $`\mu _1^2(\mathrm{\Phi }_1^{}\mathrm{\Phi }_1)+\mu _2^2(\mathrm{\Phi }_2^{}\mathrm{\Phi }_2)+m_{12}^2(\mathrm{\Phi }_1^{}\mathrm{\Phi }_2)+m_{12}^2(\mathrm{\Phi }_2^{}\mathrm{\Phi }_1)+\lambda _1(\mathrm{\Phi }_1^{}\mathrm{\Phi }_1)^2`$ (2.7)
$`+\lambda _2(\mathrm{\Phi }_2^{}\mathrm{\Phi }_2)^2+\lambda _3(\mathrm{\Phi }_1^{}\mathrm{\Phi }_1)(\mathrm{\Phi }_2^{}\mathrm{\Phi }_2)+\lambda _4(\mathrm{\Phi }_1^{}\mathrm{\Phi }_2)(\mathrm{\Phi }_2^{}\mathrm{\Phi }_1),`$
with
$`\mu _1^2`$ $`=`$ $`m_1^2|\mu |^2,\mu _2^2=m_2^2|\mu |^2,m_{12}^2=B\mu ,`$
$`\lambda _1`$ $`=`$ $`\lambda _2={\displaystyle \frac{1}{8}}(g_w^2+g^2),\lambda _3={\displaystyle \frac{1}{4}}(g_w^2g^2),\lambda _4={\displaystyle \frac{1}{2}}g_w^2.`$ (2.8)
Further, in (2.6), $`\overline{m}_i^2`$ (with $`i=t,b`$) and $`\stackrel{~}{m}_{q_k}^2`$ (with $`q_k=t_1,b_1,t_2,b_2`$) denote the eigenvalues of the quark and squark mass matrices $`^{}`$ and $`\stackrel{~}{}^2`$, respectively, which depend on the Higgs background fields. Specifically, $`^{}`$ reads
$$^{}=\left(\begin{array}{cc}|h_t|^2|\mathrm{\Phi }_2|^2& h_th_b^{}\mathrm{\Phi }_1^Ti\tau _2\mathrm{\Phi }_2\\ h_t^{}h_b\mathrm{\Phi }_1^{}i\tau _2\mathrm{\Phi }_2^{}& |h_b|^2|\mathrm{\Phi }_1|^2\end{array}\right),$$
(2.9)
with eigenvalues
$$\overline{m}_{t(b)}^2=\frac{1}{2}\left[|h_b|^2|\mathrm{\Phi }_1|^2+|h_t|^2|\mathrm{\Phi }_2|^2+()\sqrt{\left(|h_b|^2|\mathrm{\Phi }_1|^2+|h_t|^2|\mathrm{\Phi }_2|^2\right)^2\mathrm{\hspace{0.17em}4}|h_t|^2|h_b|^2|\mathrm{\Phi }_1^{}\mathrm{\Phi }_2|^2}\right].$$
(2.10)
It is easy to see that, for $`\varphi _{1,2}^\pm =0`$, (2.10) simplifies to the known expressions: $`\overline{m}_t^2=|h_t|^2|\varphi _2^0|^2`$ and $`\overline{m}_b^2=|h_b|^2|\varphi _1^0|^2`$.
The $`(4\times 4)`$ squark mass matrix $`\stackrel{~}{}^2`$ is more complicated. In the weak basis $`\{\stackrel{~}{Q}^T=(\stackrel{~}{t}_L,\stackrel{~}{b}_L),\stackrel{~}{U}=\stackrel{~}{t}_R^{},\stackrel{~}{D}=\stackrel{~}{b}_R^{}\}`$, $`\stackrel{~}{}^2`$ may be cast in the form:
$$\stackrel{~}{}^2=\left(\begin{array}{ccc}(\stackrel{~}{}^2)_{\stackrel{~}{Q}^{}\stackrel{~}{Q}}& (\stackrel{~}{}^2)_{\stackrel{~}{Q}^{}\stackrel{~}{U}^{}}& (\stackrel{~}{}^2)_{\stackrel{~}{Q}^{}\stackrel{~}{D}^{}}\\ (\stackrel{~}{}^2)_{\stackrel{~}{U}\stackrel{~}{Q}}& (\stackrel{~}{}^2)_{\stackrel{~}{U}\stackrel{~}{U}^{}}& (\stackrel{~}{}^2)_{\stackrel{~}{U}\stackrel{~}{D}^{}}\\ (\stackrel{~}{}^2)_{\stackrel{~}{D}\stackrel{~}{Q}}& (\stackrel{~}{}^2)_{\stackrel{~}{D}\stackrel{~}{U}^{}}& (\stackrel{~}{}^2)_{\stackrel{~}{D}\stackrel{~}{D}^{}}\end{array}\right),$$
(2.11)
with
$`(\stackrel{~}{}^2)_{\stackrel{~}{Q}^{}\stackrel{~}{Q}}`$ $`=`$ $`\stackrel{~}{M}_Q^2\mathbf{\hspace{0.17em}1}_2+|h_b|^2\mathrm{\Phi }_1\mathrm{\Phi }_1^{}+|h_t|^2\left(\mathrm{\Phi }_2^{}\mathrm{\Phi }_2\mathbf{\hspace{0.17em}1}_2\mathrm{\Phi }_2\mathrm{\Phi }_2^{}\right){\displaystyle \frac{1}{2}}g_w^2\left(\mathrm{\Phi }_1\mathrm{\Phi }_1^{}\mathrm{\Phi }_2\mathrm{\Phi }_2^{}\right)`$
$`+\left({\displaystyle \frac{1}{4}}g_w^2{\displaystyle \frac{1}{12}}g^2\right)\left(\mathrm{\Phi }_1^{}\mathrm{\Phi }_1\mathrm{\Phi }_2^{}\mathrm{\Phi }_2\right)\mathbf{\hspace{0.17em}1}_2,`$
$`(\stackrel{~}{}^2)_{\stackrel{~}{U}\stackrel{~}{Q}}`$ $`=`$ $`(\stackrel{~}{}^2)_{\stackrel{~}{Q}^{}\stackrel{~}{U}^{}}^{}=h_tA_t\mathrm{\Phi }_2^Ti\tau _2+h_t\mu ^{}\mathrm{\Phi }_1^Ti\tau _2,`$
$`(\stackrel{~}{}^2)_{\stackrel{~}{D}\stackrel{~}{Q}}`$ $`=`$ $`(\stackrel{~}{}^2)_{\stackrel{~}{Q}^{}\stackrel{~}{D}^{}}^{}=h_bA_b\mathrm{\Phi }_1^{}h_b\mu ^{}\mathrm{\Phi }_2^{},`$
$`(\stackrel{~}{}^2)_{\stackrel{~}{U}\stackrel{~}{U}^{}}`$ $`=`$ $`\stackrel{~}{M}_t^2+|h_t|^2\mathrm{\Phi }_2^{}\mathrm{\Phi }_2+{\displaystyle \frac{1}{3}}g^2\left(\mathrm{\Phi }_1^{}\mathrm{\Phi }_1\mathrm{\Phi }_2^{}\mathrm{\Phi }_2\right),`$
$`(\stackrel{~}{}^2)_{\stackrel{~}{D}\stackrel{~}{D}^{}}`$ $`=`$ $`\stackrel{~}{M}_b^2+|h_b|^2\mathrm{\Phi }_1^{}\mathrm{\Phi }_1{\displaystyle \frac{1}{6}}g^2\left(\mathrm{\Phi }_1^{}\mathrm{\Phi }_1\mathrm{\Phi }_2^{}\mathrm{\Phi }_2\right),`$
$`(\stackrel{~}{}^2)_{\stackrel{~}{U}\stackrel{~}{D}^{}}`$ $`=`$ $`(\stackrel{~}{}^2)_{\stackrel{~}{D}\stackrel{~}{U}^{}}^{}=h_th_b^{}\mathrm{\Phi }_1^Ti\tau _2\mathrm{\Phi }_2.`$ (2.12)
Here and in the following, we denote by $`\stackrel{~}{M}_Q^2`$, $`\stackrel{~}{M}_t^2`$ and $`\stackrel{~}{M}_b^2`$ the soft supersymmetry-breaking masses of the third generation of squarks. It is rather difficult to express the four eigenvalues $`\stackrel{~}{m}_{q_k}^2`$ ($`q_k=t_1,b_1,t_2,b_2`$) of $`\stackrel{~}{}^2`$ in a simple form. However, as we detail in Appendix A.3, it is not necessary to know the analytic form of $`\stackrel{~}{m}_{q_k}^2`$ in order to evaluate the Higgs-boson masses and mixing angles . Of course, for $`\varphi _{1,2}^\pm =0`$, the field-dependent squark eigenvalues simplify to
$`\stackrel{~}{m}_{t_1(t_2)}^2`$ $`=`$ $`{\displaystyle \frac{1}{2}}[\stackrel{~}{M}_Q^2+\stackrel{~}{M}_t^2+\mathrm{\hspace{0.17em}2}|h_t|^2|\varphi _2^0|^2+{\displaystyle \frac{g_w^2+g^2}{4}}(|\varphi _1^0|^2|\varphi _2^0|^2)`$
$`+()\sqrt{\left[\stackrel{~}{M}_Q^2\stackrel{~}{M}_t^2+x_t\left(|\varphi _1^0|^2|\varphi _2^0|^2\right)\right]^2+4|h_t|^2\left|A_t\varphi _2^0\mu ^{}\varphi _1^0\right|^2}],`$
$`\stackrel{~}{m}_{b_1(b_2)}^2`$ $`=`$ $`{\displaystyle \frac{1}{2}}[\stackrel{~}{M}_Q^2+\stackrel{~}{M}_b^2+\mathrm{\hspace{0.17em}2}|h_b|^2|\varphi _1^0|^2{\displaystyle \frac{g_w^2+g^2}{4}}(|\varphi _1^0|^2|\varphi _2^0|^2)`$ (2.13)
$`+()\sqrt{\left[\stackrel{~}{M}_Q^2\stackrel{~}{M}_b^2x_b\left(|\varphi _1^0|^2|\varphi _2^0|^2\right)\right]^2+4|h_b|^2\left|A_b^{}\varphi _1^0\mu \varphi _2^0\right|^2}],`$
where $`x_t=\frac{1}{4}(g_w^2\frac{5}{3}g^2)`$ and $`x_b=\frac{1}{4}(g_w^2\frac{1}{3}g^2)`$.
After having set the stage, we now derive the minimization conditions governing the ground state of the MSSM one-loop effective potential and then determine the Higgs-boson mass matrices. As usual, we consider the following linear expansion of the Higgs doublets $`\mathrm{\Phi }_1`$ and $`\mathrm{\Phi }_2`$ around the ground state:
$$\mathrm{\Phi }_1=\left(\begin{array}{c}\varphi _1^+\\ \frac{1}{\sqrt{2}}(v_1+\varphi _1+ia_1)\end{array}\right),\mathrm{\Phi }_2=e^{i\xi }\left(\begin{array}{c}\varphi _2^+\\ \frac{1}{\sqrt{2}}(v_2+\varphi _2+ia_2)\end{array}\right),$$
(2.14)
where $`v_1`$ and $`v_2`$ are the moduli of the vacuum expectation values (VEVs) of the Higgs doublets and $`\xi `$ is their relative phase. Following , we require the vanishing of the total tadpole contributions
$`T_{\varphi _1(\varphi _2)}`$ $``$ $`<{\displaystyle \frac{_V}{\varphi _{1(2)}}}>=v_{1(2)}\left[\mu _{1(2)}^2+{\displaystyle \frac{v_1v_2}{v_{1(2)}^2}}\mathrm{Re}(m_{12}^2e^{i\xi })+\lambda _{1(2)}v_{1(2)}^2+{\displaystyle \frac{1}{2}}(\lambda _3+\lambda _4)v_{2(1)}^2\right]`$ (2.15)
$`{\displaystyle \frac{3}{16\pi ^2}}{\displaystyle \underset{q=t,b}{}}[{\displaystyle \underset{i=1,2}{}}<{\displaystyle \frac{\stackrel{~}{m}_{q_i}^2}{\varphi _{1(2)}}}>m_{\stackrel{~}{q}_i}^2(\mathrm{ln}{\displaystyle \frac{m_{\stackrel{~}{q}_i}^2}{Q^2}}1)\mathrm{\hspace{0.25em}2}<{\displaystyle \frac{\overline{m}_q^2}{\varphi _{1(2)}}}>m_q^2`$
$`\times (\mathrm{ln}{\displaystyle \frac{m_q^2}{Q^2}}1)],`$
$`T_{a_1(a_2)}`$ $``$ $`<{\displaystyle \frac{_V}{a_{1(2)}}}>=+()v_{2(1)}\mathrm{Im}(m_{12}^2e^{i\xi }){\displaystyle \frac{3}{16\pi ^2}}{\displaystyle \underset{q=t,b}{}}{\displaystyle \underset{i=1,2}{}}<{\displaystyle \frac{\stackrel{~}{m}_{q_i}^2}{a_{1(2)}}}>m_{\stackrel{~}{q}_i}^2`$ (2.16)
$`\times \left(\mathrm{ln}{\displaystyle \frac{m_{\stackrel{~}{q}_i}^2}{Q^2}}1\right),`$
where $`<\stackrel{~}{m}_{q_i}^2>=m_{\stackrel{~}{q}_i}^2`$, and the tadpole derivatives $`<\overline{m}_q^2/\varphi _{1(2)}>`$, $`<\stackrel{~}{m}_{q_k}^2/\varphi _{1(2)}>`$ and $`<\stackrel{~}{m}_{q_k}^2/a_{1(2)}>`$ are given in (A.1) and (A.2). Moreover, from (A.2), we readily see that $`T_{a_1}=\mathrm{tan}\beta T_{a_2}`$, with $`\mathrm{tan}\beta =v_2/v_1`$. This last fact allows us to perform an orthogonal rotation in the space spanned by the ‘CP-odd’ scalars $`a_1`$ and $`a_2`$,
$$\left(\begin{array}{c}a_1\\ a_2\end{array}\right)=\left(\begin{array}{cc}\mathrm{cos}\beta & \mathrm{sin}\beta \\ \mathrm{sin}\beta & \mathrm{cos}\beta \end{array}\right)\left(\begin{array}{c}G^0\\ a\end{array}\right).$$
(2.17)
The Higgs potential then has a flat direction with respect to the $`G^0`$ field, i.e., $`_V/G^0=0`$, and the $`G^0`$ field becomes the true would-be Goldstone boson eaten by the longitudinal component of the $`Z`$ boson.
We observe from (2.16) that a relative phase $`\xi `$ between the two Higgs vacuum expectation values is induced radiatively in the $`\overline{\mathrm{MS}}`$ scheme . However, we should stress that the phase $`\xi `$ is renormalization-scheme dependent. For example, one may adopt a renormalization scheme, slightly different from the $`\overline{\mathrm{MS}}`$ one, in which $`\xi `$ is set to zero order-by-order in perturbation theory . This can be achieved by requiring the bilinear Higgs-mixing mass $`m_{12}^2`$ to be real at the tree level, but to receive an imaginary counter-term (CT), $`\mathrm{Im}m_{12}^2`$, at higher orders, which is determined by the vanishing of the CP-odd tadpole parameters $`T_{a_1}`$ and $`T_{a_2}`$ for $`\xi =0`$. As we detail below, the scheme of renormalization of $`\mathrm{Im}(m_{12}^2e^{i\xi })`$ does not directly affect the renormalization scheme of other physical kinematic parameters of the theory to one loop, such as Higgs-boson masses and $`\mathrm{tan}\beta `$. In fact, it has been explicitly demonstrated in that physical CP-violating transition amplitudes, such as scalar-pseudoscalar transitions, are independent of the renormalization subtraction point $`Q^2`$ and the choice of phase $`\xi `$. In the following, we adopt the $`\xi =0`$ scheme of renormalization, as irrelevant $`\xi `$-dependent phases in the effective chargino and neutralino mass matrices can thereby be completely avoided.
In the remaining part of this section, we evaluate the one-loop effective Higgs-boson mass matrices. Employing the tadpole conditions $`T_{\varphi _1}=T_{\varphi _2}=0`$ and $`T_{a_1}=T_{a_2}=0`$ allows one to substitute the mass parameters $`\mu _1^2`$ and $`\mu _2^2`$, and $`\mathrm{Im}m_{12}^2`$ into the effective potential (2.6). After performing the above substitutions, we can express the charged-Higgs-boson mass matrix as follows:
$`(_\pm ^2)_{ij}`$ $`=`$ $`(1)^{i+j}{\displaystyle \frac{v_1v_2}{v_iv_j}}(\mathrm{Re}m_{12}^2+{\displaystyle \frac{1}{4}}g_w^2v_1v_2)+{\displaystyle \frac{3}{16\pi ^2}}{\displaystyle \underset{q=t,b}{}}[{\displaystyle \underset{k=1,2}{}}(<{\displaystyle \frac{^2\stackrel{~}{m}_{q_k}^2}{\varphi _i^+\varphi _j^{}}}>`$ (2.18)
$`{\displaystyle \frac{\delta _{ij}}{v_i}}<{\displaystyle \frac{\stackrel{~}{m}_{q_k}^2}{\varphi _j}}>{\displaystyle \frac{i(1\delta _{ij})}{v_i}}<{\displaystyle \frac{\stackrel{~}{m}_{q_k}^2}{a_j}}>\left)m^2_{\stackrel{~}{q}_k}\right(\mathrm{ln}{\displaystyle \frac{m_{\stackrel{~}{q}_k}^2}{Q^2}}1)`$
$`2(<{\displaystyle \frac{^2\overline{m}_q^2}{\varphi _i^+\varphi _j^{}}}>{\displaystyle \frac{\delta _{ij}}{v_i}}<{\displaystyle \frac{\overline{m}_q^2}{\varphi _j}}>)m_q^2(\mathrm{ln}{\displaystyle \frac{m_q^2}{Q^2}}1)].`$
Since $`\mathrm{det}_\pm ^2=0`$ in (2.18), the square of the charged Higgs-boson mass $`M_{H^+}^2`$ may be determined by the matrix element $`(_\pm ^2)_{12}`$:
$`M_{H^+}^2`$ $`=`$ $`{\displaystyle \frac{v^2}{v_1v_2}}\{(\mathrm{Re}m_{12}^2+{\displaystyle \frac{1}{4}}g_w^2v_1v_2){\displaystyle \frac{3}{16\pi ^2}}{\displaystyle \underset{q=t,b}{}}[{\displaystyle \underset{k=1,2}{}}(<{\displaystyle \frac{^2\stackrel{~}{m}_{q_k}^2}{\varphi _1^+\varphi _2^{}}}>{\displaystyle \frac{i}{v_1}}<{\displaystyle \frac{\stackrel{~}{m}_{q_k}^2}{a_2}}>)`$ (2.19)
$`\times m_{\stackrel{~}{q}_k}^2(\mathrm{ln}{\displaystyle \frac{m_{\stackrel{~}{q}_k}^2}{Q^2}}1)2<{\displaystyle \frac{^2\overline{m}_q^2}{\varphi _1^+\varphi _2^{}}}>m_q^2(\mathrm{ln}{\displaystyle \frac{m_q^2}{Q^2}}1)]\}.`$
The self-energy derivatives appearing in (2.18) and (2.19): $`<^2\overline{m}_q^2/\varphi _1^+\varphi _2^{}>`$ and $`<^2\stackrel{~}{m}_{q_k}^2/\varphi _1^+\varphi _2^{}>`$, as well as the tadpole terms $`<\overline{m}_{q_k}^2/\varphi _j>`$, $`<\stackrel{~}{m}_{q_k}^2/\varphi _j>`$ and $`<\stackrel{~}{m}_{q_k}^2/a_j>`$, are exhibited in Appendix A.
By analogy, in the weak basis $`\{\varphi _1,\varphi _2,a_1,a_2\}`$, the neutral-Higgs-boson mass matrix takes on the form<sup>3</sup><sup>3</sup>3Notice that our convention differs from that given in , as the neutral-Higgs-boson mass matrix $`_0^2`$ in that work is expressed in the weak basis $`\{a_1,a_2,\varphi _1,\varphi _2\}`$.
$$_0^2=\left(\begin{array}{cc}_S^2& _{SP}^2\\ (_{SP}^2)^T& _P^2\end{array}\right),$$
(2.20)
where $`_S^2`$, $`_P^2`$ and $`_{SP}^2`$ denote the two-by-two matrices of the scalar, pseudoscalar and scalar-pseudoscalar squared mass terms of the neutral Higgs bosons. Observe that the presence of CP-violating self-energy terms leads to mass eigenstates with no well-defined CP quantum numbers. Therefore the CP-odd Higgs-boson mass $`M_A`$ cannot be identified with any of the neutral Higgs-boson masses. The individual matrix elements of $`_0^2`$ are given by
$`(_S^2)_{ij}`$ $`=`$ $`(1)^{i+j}{\displaystyle \frac{v_1v_2}{v_iv_j}}\mathrm{Re}m_{12}^2+{\displaystyle \frac{1}{4}}(g_w^2+g^2)v_iv_j+{\displaystyle \frac{3}{16\pi ^2}}{\displaystyle \underset{q=t,b}{}}\{{\displaystyle \underset{k=1,2}{}}[(<{\displaystyle \frac{^2\stackrel{~}{m}_{q_k}^2}{\varphi _i\varphi _j}}>`$ (2.21)
$`{\displaystyle \frac{\delta _{ij}}{v_i}}<{\displaystyle \frac{\stackrel{~}{m}_{q_k}^2}{\varphi _j}}>)m_{\stackrel{~}{q}_k}^2(\mathrm{ln}{\displaystyle \frac{m_{\stackrel{~}{q}_k}^2}{Q^2}}1)+<{\displaystyle \frac{\stackrel{~}{m}_{q_k}^2}{\varphi _i}}><{\displaystyle \frac{\stackrel{~}{m}_{q_k}^2}{\varphi _j}}>\mathrm{ln}{\displaystyle \frac{m_{\stackrel{~}{q}_k}^2}{Q^2}}]`$
$`\mathrm{\hspace{0.17em}2}[(<{\displaystyle \frac{^2\overline{m}_q^2}{\varphi _i\varphi _j}}>{\displaystyle \frac{\delta _{ij}}{v_i}}<{\displaystyle \frac{\overline{m}_q^2}{\varphi _j}}>)m_q^2(\mathrm{ln}{\displaystyle \frac{m_q^2}{Q^2}}1)`$
$`+<{\displaystyle \frac{\overline{m}_q^2}{\varphi _i}}><{\displaystyle \frac{\overline{m}_q^2}{\varphi _j}}>\mathrm{ln}{\displaystyle \frac{m_q^2}{Q^2}}]\},`$
$`(_{SP}^2)_{ij}`$ $`=`$ $`{\displaystyle \frac{3}{16\pi ^2}}{\displaystyle \underset{q=t,b}{}}{\displaystyle \underset{k=1,2}{}}[(<{\displaystyle \frac{^2\stackrel{~}{m}_{q_k}^2}{\varphi _ia_j}}>{\displaystyle \frac{(1\delta _{ij})}{v_i}}<{\displaystyle \frac{\stackrel{~}{m}_{q_k}^2}{a_j}}>)m_{\stackrel{~}{q}_k}^2(\mathrm{ln}{\displaystyle \frac{m_{\stackrel{~}{q}_k}^2}{Q^2}}1)`$ (2.22)
$`+<{\displaystyle \frac{\stackrel{~}{m}_{q_k}^2}{\varphi _i}}><{\displaystyle \frac{\stackrel{~}{m}_{q_k}^2}{a_j}}>\mathrm{ln}{\displaystyle \frac{m_{\stackrel{~}{q}_k}^2}{Q^2}}],`$
$`(_P^2)_{ij}`$ $`=`$ $`(1)^{i+j}{\displaystyle \frac{v_1v_2}{v_iv_j}}\mathrm{Re}m_{12}^2+{\displaystyle \frac{3}{16\pi ^2}}{\displaystyle \underset{q=t,b}{}}\{{\displaystyle \underset{k=1,2}{}}[(<{\displaystyle \frac{^2\stackrel{~}{m}_{q_k}^2}{a_ia_j}}>{\displaystyle \frac{\delta _{ij}}{v_i}}<{\displaystyle \frac{\stackrel{~}{m}_{q_k}^2}{\varphi _j}}>)`$ (2.23)
$`\times m_{\stackrel{~}{q}_k}^2(\mathrm{ln}{\displaystyle \frac{m_{\stackrel{~}{q}_k}^2}{Q^2}}1)+<{\displaystyle \frac{\stackrel{~}{m}_{q_k}^2}{a_i}}><{\displaystyle \frac{\stackrel{~}{m}_{q_k}^2}{a_j}}>\mathrm{ln}{\displaystyle \frac{m_{\stackrel{~}{q}_k}^2}{Q^2}}]`$
$`2(<{\displaystyle \frac{^2\overline{m}_q^2}{a_ia_j}}>{\displaystyle \frac{\delta _{ij}}{v_i}}<{\displaystyle \frac{\overline{m}_q^2}{\varphi _j}}>)m_q^2(\mathrm{ln}{\displaystyle \frac{m_q^2}{Q^2}}1)\}.`$
Again, the analytic expressions for the self-energy and tadpole derivatives with respect to the background Higgs fields are given in Appendix A.
Since $`G^0`$ does not mix with the other neutral fields, the $`(4\times 4)`$ matrix $`_0^2`$ reduces to a $`(3\times 3)`$ matrix, which we denote by $`_N^2`$. In the weak basis $`\{\varphi _1,\varphi _2,a\}`$, the reduced neutral mass-squared matrix $`_N^2`$ may be expressed by
$$_N^2=\left(\begin{array}{ccc}(_S^2)_{11}& (_S^2)_{12}& \frac{1}{\mathrm{cos}\beta }(_{SP}^2)_{12}\\ (_S^2)_{21}& (_S^2)_{22}& \frac{1}{\mathrm{sin}\beta }(_{SP}^2)_{21}\\ \frac{1}{\mathrm{cos}\beta }(_{SP}^2)_{12}& \frac{1}{\mathrm{sin}\beta }(_{SP}^2)_{21}& \frac{1}{\mathrm{sin}\beta \mathrm{cos}\beta }(_P^2)_{12}\end{array}\right).$$
(2.24)
In writing $`_N^2`$ in (2.24), we have used the properties of the matrix elements of $`_{SP}^2`$: $`(_{SP}^2)_{11}=\mathrm{tan}\beta (_{SP}^2)_{12}`$ and $`(_{SP}^2)_{22}=\mathrm{cot}\beta (_{SP}^2)_{21}`$, and likewise for $`_P^2`$.
Using the expressions (2.19) and (2.24), we determine the analytic forms of the RG-improved charged and neutral Higgs-boson masses in the next Section.
## 3 RG-Improved Higgs-Boson Mass Matrices
In this Section, we perform a one-loop RG improvement of the squared charged Higgs-boson mass $`M_{H^+}^2`$ and of the squared neutral Higgs-boson mass matrix $`_N^2`$. The RG improvement incorporates all leading two-loop logarithmic corrections to the Higgs-boson mass-matrix elements, which were already found in the CP-conserving case to give rise to significant contributions to the Higgs-boson masses and couplings. In particular, the upper bound on the lightest CP-even Higgs mass was found to be strongly affected by the two-loop logarithmic corrections . In carrying out the RG improvement, we follow the procedure outlined in , in which the improvement of the Higgs-boson mass-matrix elements was performed by carefully applying the process of decoupling of the third-generation squarks.
Within the framework of the RG approach, the dominant contributions to the Higgs-boson mass matrix $`^2`$ may be written conceptually as a sum of two terms:
$$^2(m_t)=\overline{}^2(m_t)+^{2,\mathrm{th}}(m_t).$$
(3.1)
The first term, $`\overline{}^2(m_t)`$, contains the genuine logarithmic contributions which determine the whole scale dependence of the one-loop effective potential. These contributions would be present, even if the left-right mixing of the stop and sbottom states were absent. The second term, $`^{2,\mathrm{th}}(m_t)`$, describes the threshold effect of the decoupling of the heavier stop and sbottom squarks and their respective mixing with the lighter states. At the one-loop level, the second term is manifestly scale independent. In (3.1), we are interested in evaluating the effective potential at $`m_t`$, since it has been shown that this is the scale at which two-loop corrections are minimized. As we explain below, the renormalization of the above two contributions must proceed in different ways.
Let us denote by $`Q_{tb}`$ the scale of the heaviest third-generation squark, which we assume to be higher than the electroweak scale. In the language of the RG approach, we have first to consider that the aforementioned threshold contribution is ‘frozen’ at the scale $`Q_{tb}`$, $`^{2,\mathrm{th}}(Q_{tb})\mathrm{\Delta }^2(Q_{tb})`$, with all the involved kinematic parameters defined at this particular scale. Then, we have to rescale the threshold contribution with the anomalous dimension factors of the relevant Higgs fields:
$$_{ij}^{2,\mathrm{th}}(m_t)=\mathrm{\Delta }_{ij}^2(Q_{tb})\xi _i^1(m_t)\xi _j^1(m_t),$$
(3.2)
where $`\xi _i(m_t)`$ is the anomalous dimension factor of the $`H_i`$ state to be determined below. The one-loop matrix elements $`\mathrm{\Delta }_{ij}^2(Q_{tb})`$ depend on the running quark masses at the scale $`Q_{tb}`$, which have to be conveniently re-expressed as functions of the corresponding running masses at $`m_t`$. Thus, the anomalous-dimension factors in combination with the one-loop relation between the quark masses at scales $`Q_{tb}`$ and $`m_t`$ yield sizeable two-loop corrections to the mass-matrix elements originating from the one-loop threshold effects.
As was already mentioned, the contribution $`\overline{}^2(m_t)`$ of the third-generation squarks to the effective potential describes the genuine one- and two-loop leading-logarithmic running of the Higgs quartic couplings. In this context, there are two important technical details that should be mentioned. First, we notice that, in the MSSM, the tree-level Higgs quartic couplings $`\lambda _i`$, with $`i=1,2,3,4`$, are all proportional to the squared gauge couplings $`g_w^2`$ and $`g^2`$ (cf. (2)). However, the one-loop $`\beta `$ functions of $`\lambda _i`$ can generally have appreciable values, as they are proportional to the fourth power of the top- and bottom-quark Yukawa couplings. As a result, the low-energy values of $`\lambda _i`$ differ significantly from their tree-level ones. The RG-improved approach followed here is crucial for implementing properly the potentially large logarithmic corrections to the Higgs quartic couplings.
The second technical remark pertains to the RG evolution of the Higgs quartic couplings $`\lambda _i`$, with $`i=5,6,7`$ (for the notation, see ), which are absent in the Born approximation to the MSSM Higgs potential. On field-theoretic grounds, these quartic couplings must have vanishing one-loop $`\beta `$ functions, and cannot be generated by RG running. However, these quartic couplings are radiatively induced by threshold effects, and have already been taken into account in $`\mathrm{\Delta }^2(Q_{tb})`$, given by (3.2).
Following the above discussion, we now proceed with the RG improvement of the Higgs-boson mass-matrix elements. To this end, we first need to compute the one-loop values of the quartic couplings $`\lambda _i`$, with $`i=1,2,3,4`$, where the decoupling of the stop and sbottom contributions at their appropriate thresholds is properly taken into account. The best way to calculate the latter effects is to consider the logarithmic part of the effective potential (2.6) in the limit where the squark mixing parameters vanish, i.e., $`\mu =A_t=A_b=0`$. The pertinent one-loop running quartic couplings, denoted by $`\lambda _i^{(1)}`$, may then be obtained by
$`\lambda _1^{(1)}`$ $`=`$ $`{\displaystyle \frac{3}{32\pi ^2}}[({\displaystyle \frac{g_w^2}{4}}{\displaystyle \frac{g^2}{12}})^2\mathrm{ln}\left({\displaystyle \frac{\stackrel{~}{M}_Q^2+m_t^2}{Q^2}}\right)+(|h_b|^2{\displaystyle \frac{g_w^2}{4}}{\displaystyle \frac{g^2}{12}})^2\mathrm{ln}\left({\displaystyle \frac{\stackrel{~}{M}_Q^2+m_b^2}{Q^2}}\right)`$ (3.3)
$`+{\displaystyle \frac{g^4}{9}}\mathrm{ln}\left({\displaystyle \frac{\stackrel{~}{M}_t^2+m_t^2}{Q^2}}\right)+(|h_b|^2{\displaystyle \frac{g^2}{6}})^2\mathrm{ln}\left({\displaystyle \frac{\stackrel{~}{M}_b^2+m_b^2}{Q^2}}\right)],`$
$`\lambda _2^{(1)}`$ $`=`$ $`{\displaystyle \frac{3}{32\pi ^2}}[(|h_t|^2{\displaystyle \frac{g_w^2}{4}}+{\displaystyle \frac{g^2}{12}})^2\mathrm{ln}\left({\displaystyle \frac{\stackrel{~}{M}_Q^2+m_t^2}{Q^2}}\right)+({\displaystyle \frac{g_w^2}{4}}+{\displaystyle \frac{g^2}{12}})^2\mathrm{ln}\left({\displaystyle \frac{\stackrel{~}{M}_Q^2+m_b^2}{Q^2}}\right)`$ (3.4)
$`+(|h_t|^2{\displaystyle \frac{g^2}{3}})^2\mathrm{ln}\left({\displaystyle \frac{\stackrel{~}{M}_t^2+m_t^2}{Q^2}}\right)+{\displaystyle \frac{g^4}{36}}\mathrm{ln}\left({\displaystyle \frac{\stackrel{~}{M}_b^2+m_b^2}{Q^2}}\right)],`$
$`\lambda _3^{(1)}`$ $`=`$ $`{\displaystyle \frac{3}{16\pi ^2}}\{|h_t|^2|h_b|^2[\mathrm{ln}\left({\displaystyle \frac{\stackrel{~}{M}_Q^2+m_t^2}{Q^2}}\right)+\mathrm{ln}\left({\displaystyle \frac{\mathrm{max}(\stackrel{~}{M}_t^2+m_t^2,\stackrel{~}{M}_b^2+m_b^2)}{Q^2}}\right)]`$ (3.5)
$`\left[\left({\displaystyle \frac{g_w^2}{4}}+{\displaystyle \frac{g^2}{12}}\right)\left(|h_t|^2{\displaystyle \frac{g_w^2}{4}}\right){\displaystyle \frac{g^2}{12}}\left({\displaystyle \frac{g_w^2}{4}}{\displaystyle \frac{g^2}{12}}\right)\right]\mathrm{ln}\left({\displaystyle \frac{\stackrel{~}{M}_Q^2+m_t^2}{Q^2}}\right)`$
$`\left[\left({\displaystyle \frac{g_w^2}{4}}{\displaystyle \frac{g^2}{12}}\right)\left(|h_b|^2{\displaystyle \frac{g_w^2}{4}}\right)+{\displaystyle \frac{g^2}{12}}\left({\displaystyle \frac{g_w^2}{4}}+{\displaystyle \frac{g^2}{12}}\right)\right]\mathrm{ln}\left({\displaystyle \frac{\stackrel{~}{M}_Q^2+m_b^2}{Q^2}}\right)`$
$`+{\displaystyle \frac{g^2}{3}}(|h_t|^2{\displaystyle \frac{g^2}{3}})\mathrm{ln}\left({\displaystyle \frac{\stackrel{~}{M}_t^2+m_t^2}{Q^2}}\right)+{\displaystyle \frac{g^2}{6}}(|h_b|^2{\displaystyle \frac{g^2}{6}})\mathrm{ln}\left({\displaystyle \frac{\stackrel{~}{M}_b^2+m_b^2}{Q^2}}\right)\},`$
$`\lambda _4^{(1)}`$ $`=`$ $`{\displaystyle \frac{3}{16\pi ^2}}\{|h_t|^2|h_b|^2[\mathrm{ln}\left({\displaystyle \frac{\stackrel{~}{M}_Q^2+m_t^2}{Q^2}}\right)+\mathrm{ln}\left({\displaystyle \frac{\mathrm{max}(\stackrel{~}{M}_t^2+m_t^2,\stackrel{~}{M}_b^2+m_b^2)}{Q^2}}\right)]`$ (3.6)
$`{\displaystyle \frac{g_w^2}{2}}(|h_t|^2{\displaystyle \frac{g_w^2}{4}})\mathrm{ln}\left({\displaystyle \frac{\stackrel{~}{M}_Q^2+m_t^2}{Q^2}}\right){\displaystyle \frac{g_w^2}{2}}(|h_b|^2{\displaystyle \frac{g_w^2}{4}})\mathrm{ln}\left({\displaystyle \frac{\stackrel{~}{M}_Q^2+m_b^2}{Q^2}}\right)\}.`$
Moreover, we need to know the one-loop running of the soft supersymmetry-breaking parameter $`\mathrm{Re}m_{12}^2`$. Gathering the relevant logarithmic terms present in the effective potential (2.6), we find
$`\mathrm{Re}m_{12}^{2(1)}`$ $`=`$ $`{\displaystyle \frac{3}{16\pi ^2}}[|h_t|^2\mathrm{Re}(\mu A_t)\mathrm{ln}\left({\displaystyle \frac{\mathrm{max}(\stackrel{~}{M}_Q^2+m_t^2,\stackrel{~}{M}_t^2+m_t^2)}{Q^2}}\right)`$ (3.7)
$`+|h_b|^2\mathrm{Re}(\mu A_b)\mathrm{ln}\left({\displaystyle \frac{\mathrm{max}(\stackrel{~}{M}_Q^2+m_b^2,\stackrel{~}{M}_b^2+m_b^2)}{Q^2}}\right)].`$
The analytic form of $`\mathrm{\Delta }^2(Q_{tb})`$ in (3.2) can now be obtained by subtracting the one-loop Born-improved mass matrix $`^{2(0)}`$ from its total one-loop contribution $`^{2(1)}`$. Here, we have in mind the charged and neutral Higgs-boson mass matrices $`_\pm ^{2(1)}`$ and $`_N^{2(1)}`$ calculated in Section 2. More explicitly, $`\mathrm{\Delta }^2(Q_{tb})`$ is given by
$$\mathrm{\Delta }^2(Q_{tb})=^{2(1)}(Q_{tb})^{2(0)}[\mathrm{Re}m_{12}^{2(1)}(Q_{tb}),\lambda _i^{(1)}(Q_{tb})],$$
(3.8)
where $`^{2(0)}`$ represents the tree-level functional form of $`^2`$, expressed in terms of $`\lambda _i^{(1)}`$ and $`\mathrm{Re}m_{12}^{2(1)}`$. Furthermore, it is essential to stress again that the kinematic parameters involved in (3.8), such as masses and couplings, are evaluated at the scale $`Q_{tb}`$.
Another important ingredient in the RG improvement of the Higgs-boson mass matrices is the analytic two-loop result for the Higgs quartic couplings $`\lambda _1,\mathrm{},\lambda _4`$. As has been done in (3.3)–(3.6), we have to include two-loop leading logarithms, by appropriately considering the stop and sbottom thresholds. These two-loop leading logarithmic contributions to the Higgs quartic couplings, which we denote by $`\lambda _i^{(2)}`$, can be determined by solving iteratively the RG equations . In this way, we obtain
$`\lambda _1^{(2)}`$ $`=`$ $`{\displaystyle \frac{6|h_b|^4}{(32\pi ^2)^2}}\left({\displaystyle \frac{3}{2}}|h_b|^2+{\displaystyle \frac{1}{2}}|h_t|^2\mathrm{\hspace{0.25em}8}g_s^2\right)\left[\mathrm{ln}^2\left({\displaystyle \frac{\stackrel{~}{M}_Q^2+m_b^2}{Q^2}}\right)+\mathrm{ln}^2\left({\displaystyle \frac{\stackrel{~}{M}_b^2+m_b^2}{Q^2}}\right)\right],`$ (3.9)
$`\lambda _2^{(2)}`$ $`=`$ $`{\displaystyle \frac{6|h_t|^4}{(32\pi ^2)^2}}\left({\displaystyle \frac{3}{2}}|h_t|^2+{\displaystyle \frac{1}{2}}|h_b|^2\mathrm{\hspace{0.25em}8}g_s^2\right)\left[\mathrm{ln}^2\left({\displaystyle \frac{\stackrel{~}{M}_Q^2+m_t^2}{Q^2}}\right)+\mathrm{ln}^2\left({\displaystyle \frac{\stackrel{~}{M}_t^2+m_t^2}{Q^2}}\right)\right],`$ (3.10)
$`\lambda _3^{(2)}`$ $`=`$ $`{\displaystyle \frac{3|h_t|^2|h_b|^2}{(16\pi ^2)^2}}\left(|h_t|^2+|h_b|^2\mathrm{\hspace{0.25em}8}g_s^2\right)`$ (3.11)
$`\times \left[\mathrm{ln}^2\left({\displaystyle \frac{\stackrel{~}{M}_Q^2+m_t^2}{Q^2}}\right)+\mathrm{ln}^2\left({\displaystyle \frac{\mathrm{max}(\stackrel{~}{M}_t^2+m_t^2,\stackrel{~}{M}_b^2+m_b^2)}{Q^2}}\right)\right],`$
$`\lambda _4^{(2)}`$ $`=`$ $`\lambda _3^{(2)}.`$ (3.12)
For later convenience, we define collectively the sum of the tree, one-loop and two-loop quartic couplings as follows:
$$\overline{\lambda }_i=\lambda _i+\lambda _i^{(1)}+\lambda _i^{(2)},$$
(3.13)
with $`i=1,2,3,4`$. Similarly, the sum of the tree, one-loop and two-loop contributions to the soft-bilinear Higgs mixing may be defined as
$$\mathrm{Re}\overline{m}_{12}^2=\mathrm{Re}m_{12}^2+\mathrm{Re}m_{12}^{2(1)}+\mathrm{Re}m_{12}^{2(2)}.$$
(3.14)
As we see below, knowledge of the two-loop contribution $`\mathrm{Re}m_{12}^{2(2)}`$ is not required in the one-loop RG improvement of the MSSM Higgs potential.
Given the above definitions of the quartic couplings and the soft Higgs-mixing parameter in (3.13) and (3.14), we can express the one- and two-loop leading logarithmic contributions $`\overline{}^2(m_t)`$ to $`^2(m_t)`$ by means of the two-loop Born-improved mass matrix:
$$\overline{}^2(m_t)=^{2(0)}[\mathrm{Re}\overline{m}_{12}^2(m_t),\overline{\lambda }_i(m_t)].$$
(3.15)
Note that $`\overline{}^2(m_t)`$ also includes the tree-level terms. As has also been stated explicitly in (3.15), $`\overline{}^2(m_t)`$ is expressed in terms of mass and coupling parameters evaluated at the top-quark-mass scale.
The last ingredient for completing the programme of RG improvement of the Higgs-boson mass matrices is knowledge of the analytic expressions for the anomalous dimension factors that occur in (3.2). These analytic expressions are given below for the charged and neutral Higgs-boson cases separately.
Adopting the framework outlined above , it is not difficult to compute the RG-improved charged Higgs-boson mass $`M_{H^+}^2(m_t)`$ at the top-mass scale through the relation:
$$M_{H^+}^2(m_t)=\overline{M}_{H^+}^2(m_t)+\left[\xi _1^+(m_t)\xi _2^{}(m_t)\right]^1(\mathrm{\Delta }M_{H^+}^2)^{\stackrel{~}{t}\stackrel{~}{b}}(Q_{tb})+(M_{H^+}^{2(1)})^{tb}(m_t),$$
(3.16)
where $`Q_{tb}^2=\mathrm{max}(\stackrel{~}{M}_Q^2+m_t^2,\stackrel{~}{M}_t^2+m_t^2,\stackrel{~}{M}_b^2+m_b^2)`$, and $`\xi _1^+(m_t)`$ and $`\xi _2^{}(m_t)`$ are the anomalous dimension factors of the charged Higgs fields $`\varphi _1^+`$ and $`\varphi _2^{}`$, respectively:
$$\xi _1^+(m_t)=1+\frac{3|h_b|^2}{32\pi ^2}\mathrm{ln}\frac{Q_{tb}^2}{m_t^2},\xi _2^{}(m_t)=1+\frac{3|h_t|^2}{32\pi ^2}\mathrm{ln}\frac{Q_{tb}^2}{m_t^2}.$$
(3.17)
Further, $`\overline{M}_{H^+}^2(m_t)`$ is the squared two-loop Born-improved charged Higgs-boson mass given by
$$\overline{M}_{H^+}^2(m_t)=\frac{\mathrm{Re}\overline{m}_{12}^2(m_t)}{\mathrm{sin}\beta (m_t)\mathrm{cos}\beta (m_t)}+\frac{1}{2}\overline{\lambda }_4(m_t)v^2(m_t)$$
(3.18)
and $`(\mathrm{\Delta }M_{H^+}^2)^{\stackrel{~}{t}\stackrel{~}{b}}`$ is the one-loop scale-invariant part that contains the stop and sbottom contributions:
$`(\mathrm{\Delta }M_{H^+}^2)^{\stackrel{~}{t}\stackrel{~}{b}}(Q_{tb})`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{sin}\beta (m_t)\mathrm{cos}\beta (m_t)}}[(_\pm ^{2(1)})_{12}^{\stackrel{~}{t}\stackrel{~}{b}}(Q_{tb})+\mathrm{Re}m_{12}^{2(1)}(Q_{tb})`$ (3.19)
$`+{\displaystyle \frac{1}{2}}\lambda _4^{(1)}(Q_{tb})v_1(Q_{tb})v_2(Q_{tb})],`$
where the scale at which the kinematic parameters are to be evaluated has been indicated explicitly. In (3.19), the term between brackets is the threshold contribution to the off-diagonal matrix element of the charged-Higgs-boson mass matrix, where $`(_\pm ^{2(1)})_{12}^{\stackrel{~}{t}\stackrel{~}{b}}`$ denotes the one-loop contribution of the third-generation squarks to $`(_\pm ^2)_{12}`$. Finally, $`(M_{H^+}^{2(1)})^{tb}`$ describes the one-loop quark contribution to $`M_{H^+}^2`$ (see (3.16)).
We remark that the charged-Higgs-boson mass matrix receives the common anomalous dimension factor $`\xi _1^+(m_t)\xi _2^{}(m_t)`$, even though different matrix elements of $`_\pm ^2`$ are involved. This is because $`_\pm ^2`$ must possess a vanishing determinant at any RG scale $`Q^2`$, as one of its mass eigenstates must correspond to the massless would-be Goldstone boson $`G^+`$ that forms the longitudinal component of the $`W^+`$ boson. As a consequence, the following relations among the matrix elements of the RG-frozen part $`\mathrm{\Delta }_\pm ^2`$ of the charged-Higgs-boson mass matrix are obtained:
$`(\mathrm{\Delta }_\pm ^2)_{11}(Q_{tb})`$ $`=`$ $`\mathrm{tan}\beta (Q_{tb})(\mathrm{\Delta }_\pm ^2)_{12}(Q_{tb}),`$
$`(\mathrm{\Delta }_\pm ^2)_{22}(Q_{tb})`$ $`=`$ $`\mathrm{cot}\beta (Q_{tb})(\mathrm{\Delta }_\pm ^2)_{21}(Q_{tb}).`$ (3.20)
After including the RG running due to the Higgs-boson anomalous dimensions, we find
$`{\displaystyle \frac{(\mathrm{\Delta }_\pm ^2)_{11}(Q_{tb})}{[\xi _1^+(m_t)]^2}}`$ $`=`$ $`\mathrm{tan}\beta (m_t){\displaystyle \frac{(\mathrm{\Delta }_\pm ^2)_{12}(Q_{tb})}{\xi _1^+(m_t)\xi _2^{}(m_t)}},`$
$`{\displaystyle \frac{(\mathrm{\Delta }_\pm ^2)_{22}(Q_{tb})}{[\xi _2^{}(m_t)]^2}}`$ $`=`$ $`\mathrm{cot}\beta (m_t){\displaystyle \frac{(\mathrm{\Delta }_\pm ^2)_{12}(Q_{tb})}{\xi _1^+(m_t)\xi _2^{}(m_t)}},`$ (3.21)
where we have used the RG relation:
$$\mathrm{tan}\beta (Q_{tb})=\frac{\xi _1^+(m_t)}{\xi _2^{}(m_t)}\mathrm{tan}\beta (m_t).$$
(3.22)
As a consequence of this last relation, it is evident that the RG running of the different matrix elements of $`\mathrm{\Delta }_\pm ^2`$ may be expressed in terms of the running of $`(\mathrm{\Delta }_\pm ^2)_{12}`$ and the value of $`\mathrm{tan}\beta `$ at the scale $`m_t`$.
Correspondingly, the RG-improved neutral-Higgs-boson mass matrix $`_N^2`$ may be computed by
$`(_N^2)_{ij}(m_t)`$ $`=`$ $`(\overline{}_N^2)_{ij}(m_t)+\left[\xi _{ij}^{\stackrel{~}{t}}(m_t)\right]^1(\mathrm{\Delta }_N^2)_{ij}^{\stackrel{~}{t}}(Q_t)+\left[\xi _{ij}^{\stackrel{~}{b}}(m_t)\right]^1(\mathrm{\Delta }_N^2)_{ij}^{\stackrel{~}{b}}(Q_b)`$ (3.23)
$`+(_N^{2(1)})_{ij}^{tb}(m_t),`$
with $`i,j=1,2,3`$, $`Q_t^2=\mathrm{max}(\stackrel{~}{M}_Q^2+m_t^2,\stackrel{~}{M}_t^2+m_t^2)`$ and $`Q_b^2=\mathrm{max}(\stackrel{~}{M}_Q^2+m_b^2,\stackrel{~}{M}_b^2+m_b^2)`$. Notice that, unlike the charged-Higgs-boson case, one has to introduce here two decoupling scales $`Q_t`$ and $`Q_b`$, as the stop/top and sbottom/bottom loop effects occur separately in the threshold contributions. The parameters $`\xi _{ij}^{\stackrel{~}{q}}`$ (with $`q=t,b`$) are the anomalous-dimension factors related to the neutral Higgs-boson fields
$`\xi _{ij}^{\stackrel{~}{q}}(m_t)=\xi _{ji}^{\stackrel{~}{q}}(m_t)=\{\begin{array}{c}\xi _i^{\stackrel{~}{q}}(m_t)\xi _j^{\stackrel{~}{q}}(m_t),\mathrm{for}i,j=1,2\hfill \\ \xi _1^{\stackrel{~}{q}}(m_t)\xi _2^{\stackrel{~}{q}}(m_t),\mathrm{for}i=1,2,3\mathrm{and}j=3\hfill \end{array},`$ (3.26)
with
$`\xi _1^{\stackrel{~}{t}}(m_t)`$ $`=`$ $`1+{\displaystyle \frac{3|h_b|^2}{32\pi ^2}}\mathrm{ln}{\displaystyle \frac{Q_t^2}{m_t^2}},\xi _2^{\stackrel{~}{t}}(m_t)=1+{\displaystyle \frac{3|h_t|^2}{32\pi ^2}}\mathrm{ln}{\displaystyle \frac{Q_t^2}{m_t^2}},`$
$`\xi _1^{\stackrel{~}{b}}(m_t)`$ $`=`$ $`1+{\displaystyle \frac{3|h_b|^2}{32\pi ^2}}\mathrm{ln}{\displaystyle \frac{Q_b^2}{m_t^2}},\xi _2^{\stackrel{~}{b}}(m_t)=1+{\displaystyle \frac{3|h_t|^2}{32\pi ^2}}\mathrm{ln}{\displaystyle \frac{Q_b^2}{m_t^2}}.`$ (3.27)
We should observe that, for the very same reasons as in the charged-Higgs-boson case, the vanishing of the determinants of $`_P^2`$ and $`_{SP}^2`$ at any $`Q^2`$ scale leads to the common anomalous dimension factor $`\xi _{i3}(m_t)=\xi _{3i}(m_t)=\xi _1(m_t)\xi _2(m_t)`$ in the calculation of $`_N^2(m_t)`$ in (3.23). Since this last fact involves the matrix elements $`(_P^2)_{12}`$, $`(_{SP}^2)_{12}`$ and $`(_{SP}^2)_{21}`$, the corresponding matrix elements of $`(\mathrm{\Delta }_N^2)_{i3}^{\stackrel{~}{q}}`$ in (3.23) are given by
$`(\mathrm{\Delta }_N^2)_{13}^{\stackrel{~}{q}}(Q_q)`$ $`=`$ $`{\displaystyle \frac{(_{SP}^2)_{12}^{\stackrel{~}{q}}(Q_q)}{\mathrm{cos}\beta (m_t)}},(\mathrm{\Delta }_N^2)_{23}^{\stackrel{~}{q}}(Q_q)={\displaystyle \frac{(_{SP}^2)_{21}^{\stackrel{~}{q}}(Q_q)}{\mathrm{sin}\beta (m_t)}},`$
$`(\mathrm{\Delta }_N^2)_{33}^{\stackrel{~}{q}}(Q_q)`$ $`=`$ $`{\displaystyle \frac{(\mathrm{\Delta }_P^2)_{12}^{\stackrel{~}{q}}(Q_q)}{\mathrm{sin}\beta (m_t)\mathrm{cos}\beta (m_t)}},`$ (3.28)
where the RG-scale dependence of the involved quantities has been displayed explicitly.
In (3.23), the $`(3\times 3)`$ matrices $`\overline{}_N^2`$ and $`(\mathrm{\Delta }_N^2)^{\stackrel{~}{q}}`$ (with $`q=t,b`$) describe respectively the two-loop Born-improved effects and the one-loop threshold contributions associated with the decoupling of the heavy squark states:
$`\overline{}_N^2(m_t)`$ $`=`$ $`_N^{2(0)}[\mathrm{Re}\overline{m}_{12}^2(m_t),\overline{\lambda }_1(m_t),\overline{\lambda }_2(m_t),\overline{\lambda }_{34}(m_t)],`$ (3.29)
$`(\mathrm{\Delta }_N^2)^{\stackrel{~}{q}}(Q_q)`$ $`=`$ (3.30)
$`(_N^{2(1)})^{\stackrel{~}{q}}(Q_q)_N^{2(0)}[\mathrm{Re}m_{12}^{2(1),\stackrel{~}{q}}(Q_q),\lambda _1^{(1),\stackrel{~}{q}}(Q_q),\lambda _2^{(1),\stackrel{~}{q}}(Q_q),\lambda _{34}^{(1),\stackrel{~}{q}}(Q_q)],`$
where $`\lambda _{34}^{(1)}=\lambda _3^{(1)}+\lambda _4^{(1)}`$ (likewise $`\overline{\lambda }_{34}=\overline{\lambda }_3+\overline{\lambda }_4`$) and $`_N^{2(0)}`$ is the tree-level functional form of $`_N^2`$:
$`_N^{2(0)}`$ $`=`$ (3.38)
$`{\displaystyle \frac{\mathrm{Re}m_{12}^2}{\mathrm{sin}\beta \mathrm{cos}\beta }}\left(\begin{array}{ccc}\mathrm{sin}^2\beta & \mathrm{sin}\beta \mathrm{cos}\beta & 0\\ \mathrm{sin}\beta \mathrm{cos}\beta & \mathrm{cos}^2\beta & 0\\ 0& 0& 1\end{array}\right)v^2\left(\begin{array}{ccc}2\lambda _1\mathrm{cos}^2\beta & \lambda _{34}\mathrm{sin}\beta \mathrm{cos}\beta & 0\\ \lambda _{34}\mathrm{sin}\beta \mathrm{cos}\beta & 2\lambda _2\mathrm{sin}^2\beta & 0\\ 0& 0& 0\end{array}\right),`$
with $`\lambda _{34}=\lambda _3+\lambda _4`$. As in the charged-Higgs-boson case, we write $`(_N^{2(1)})^{\stackrel{~}{q}}`$ to denote the one-loop part of $`_N^2`$ containing the contributions of the third-generation squarks, and $`(_N^{2(1)})^{tb}`$ to denote its fermionic one-loop counterpart.
The resulting RG-improved Higgs-boson mass matrix $`_N^2(m_t)`$ is a symmetric, positive-definite $`(3\times 3)`$ matrix, and can therefore be diagonalized by an orthogonal transformation as follows:
$$O^T_N^2(m_t)O=\mathrm{diag}[M_{H_1}^2(m_t),M_{H_2}^2(m_t),M_{H_3}^2(m_t)],$$
(3.39)
where we have defined the Higgs fields such that their RG-improved masses satisfy the inequality:
$$M_{H_1}(m_t)M_{H_2}(m_t)M_{H_3}(m_t).$$
(3.40)
Notice that our convention in (3.39) differs from that chosen in , as we assign the Higgs fields in the reversed order. Analytic expressions for $`M_{H_i}(m_t)`$ and $`O`$ are presented in Appendix B.
Before closing this Section, two important remarks are in order. First, we observe that the free kinematic parameters of the MSSM Higgs sector are
$`M_{H^+}(m_t),\mathrm{tan}\beta (m_t),\mu (Q_{tb}),A_t(Q_{tb}),A_b(Q_{tb}),`$
$`\stackrel{~}{M}_Q^2(Q_{tb}),\stackrel{~}{M}_t^2(Q_{tb}),\stackrel{~}{M}_b^2(Q_{tb}).`$ (3.41)
In fact, the soft Higgs-mixing parameter $`\mathrm{Re}\overline{m}_{12}^2(m_t)`$ may be substituted by the squared RG-improved mass $`M_{H^+}^2(m_t)`$ of the charged Higgs boson (cf. (3.16) and (3.18)) in the neutral Higgs-boson mass matrix $`_N^2(m_t)`$ in (3.23).
Secondly, we reiterate the fact that $`\mathrm{Im}m_{12}^2`$ can be renormalized independently, without affecting the renormalization of the physical parameters of the theory . As was stressed in Section 3, the $`\xi =0`$ scheme of renormalization gives rise to a considerable simplification, since we can get rid of the radiatively-induced phase $`\xi `$ between the two Higgs vacuum expectation values in the analytic expressions of the Higgs-boson masses and mixing angles. For example, within the above $`\xi =0`$ scheme, the mass renormalization of $`H^+`$ may be entirely reabsorbed by a corresponding renormalization of $`\mathrm{Re}m_{12}^2`$ and $`\lambda _4`$. In other words, it can be shown that $`M_{H^+}`$ is $`Q^2`$-independent, after including the RG running of $`\mathrm{Re}m_{12}^2`$ and $`\lambda _4`$, denoted as $`\gamma _{\mathrm{Re}m_{12}^2}`$ and $`\beta _{\lambda _4}`$. For simplicity, we assume that only the third generation of squarks contributes to $`\gamma _{\mathrm{Re}m_{12}^2}`$, since fermions do not contribute to the RG running of $`\mathrm{Im}m_{12}^2`$. The analytic forms of $`\gamma _{\mathrm{Re}m_{12}^2}`$ and $`\beta _{\lambda _4}`$ are given by
$`\gamma _{\mathrm{Re}m_{12}^2}`$ $`=`$ $`{\displaystyle \frac{3}{16\pi ^2}}\left[|h_t|^2\mathrm{Re}(\mu A_t)+|h_b|^2\mathrm{Re}(\mu A_b)\right],`$ (3.42)
$`\beta _{\lambda _4}`$ $``$ $`{\displaystyle \frac{d\lambda _4^{(1)}}{d\mathrm{ln}Q^2}}={\displaystyle \frac{3}{16\pi ^2}}\left[\mathrm{\hspace{0.17em}2}|h_t|^2|h_b|^2{\displaystyle \frac{g_w^2}{2}}\left(|h_t|^2+|h_b|^2\right)+{\displaystyle \frac{g_w^4}{4}}\right].`$ (3.43)
Obviously, the RG running of $`\mathrm{Re}m_{12}^2`$ due to $`\stackrel{~}{t}`$ and $`\stackrel{~}{b}`$ is only relevant for non-zero values of $`\mu A_t`$ and $`\mu A_b`$. Employing (2.19) and examining only the $`\mathrm{ln}Q^2`$-dependent part, one can verify that
$$\frac{dM_{H^+}^2}{d\mathrm{ln}Q^2}\gamma _{\mathrm{Re}m_{12}^2}\frac{1}{2}\beta _{\lambda _4}v_1v_2+\frac{3}{32\pi ^2}\left(<\frac{^2\mathrm{Tr}\stackrel{~}{}^4}{\varphi _1^+\varphi _2^{}}>\frac{i}{v_1}<\frac{\mathrm{Tr}\stackrel{~}{}^4}{a_2}>\right)=0,$$
(3.44)
as it should be. As can also be seen from (3.44), an important role in this proof is played by the necessary CP-odd tadpole term $`<\mathrm{Tr}\stackrel{~}{}^4/a_2>`$ .
## 4 Effective Top and Bottom Yukawa Couplings
In addition to the RG improvement of the Higgs-boson mass matrices discussed in the previous Section, we consider here a further improvement related to the non-logarithmic threshold corrections to the top- and bottom-quark Yukawa couplings. Specifically, apart from the usual RG running, the effective top- and bottom-quark Yukawa couplings obtain additional non-logarithmic threshold contributions, which are induced by the decoupling of the heavy SUSY states at a high scale, e.g., $`Q_{tb}`$. For the bottom-quark Yukawa case, the one-loop RG relation between the bottom mass and the bottom-quark Yukawa coupling at the scale $`Q_{tb}`$ receives quantum corrections that also include terms proportional to $`\mathrm{tan}\beta `$ . Since these last terms can be significant for large values of $`\mathrm{tan}\beta `$,<sup>4</sup><sup>4</sup>4For the $`\tau `$-lepton Yukawa coupling, the corresponding enhanced $`\mathrm{tan}\beta `$ terms are much smaller, because they are proportional to the weak gauge couplings . we must resum them within the RG approach, so that the actual size of the radiative corrections to the Higgs-boson masses and couplings can properly be extracted . For the top-quark case, instead, although one-loop suppressed, the respective corrections can still give rise to an enhancement of up to 4 GeV in the prediction for the lightest Higgs-boson mass , and therefore should be included in the computation.
There may also be important CP-violating one-loop corrections to the bottom- and top-quark Yukawa couplings, in addition to the CP-violating effects induced by the radiative mixing of the Higgs states, which were considered in Sections 2 and 3 in detail. In the leptonic sector, these CP-violating vertex corrections are generally small . However, the CP-violating radiative corrections to the couplings of the Higgs bosons to $`b`$ quarks are significant , because of the large Yukawa and colour-enhanced QCD interactions . In particular, the radiative effects of the Higgs-boson couplings to the bottom quarks can be further enhanced, if the respective Higgs-mass eigenstate couples predominantly to the Higgs doublet $`\mathrm{\Phi }_2`$ , as the tree-level $`b`$-quark Yukawa coupling is suppressed in this case. For a general discussion of the form and the origin of these finite Yukawa corrections to the third-generation quark masses, the reader is referred to the original literature . In the following, we give a brief discussion of the non-logarithmic corrections to the top and bottom Yukawa couplings, and pay special attention to the CP-violating vertex effects.
We start our discussion by considering the effective Lagrangian of the $`b`$-quark Yukawa coupling :
$$_{\varphi ^0\overline{b}b}=(h_b+\delta h_b)\varphi _1^0\overline{b}_Rb_L+\mathrm{\Delta }h_b\varphi _2^0\overline{b}_Rb_L+\mathrm{h}.\mathrm{c}.,$$
(4.1)
with
$`{\displaystyle \frac{\delta h_b}{h_b}}`$ $`=`$ $`{\displaystyle \frac{2\alpha _s}{3\pi }}m_{\stackrel{~}{g}}^{}A_bI(m_{\stackrel{~}{b}_1}^2,m_{\stackrel{~}{b}_2}^2,|m_{\stackrel{~}{g}}|^2){\displaystyle \frac{|h_t|^2}{16\pi ^2}}|\mu |^2I(m_{\stackrel{~}{t}_1}^2,m_{\stackrel{~}{t}_2}^2,|\mu |^2),`$ (4.2)
$`{\displaystyle \frac{\mathrm{\Delta }h_b}{h_b}}`$ $`=`$ $`{\displaystyle \frac{2\alpha _s}{3\pi }}m_{\stackrel{~}{g}}^{}\mu ^{}I(m_{\stackrel{~}{b}_1}^2,m_{\stackrel{~}{b}_2}^2,|m_{\stackrel{~}{g}}|^2)+{\displaystyle \frac{|h_t|^2}{16\pi ^2}}A_t^{}\mu ^{}I(m_{\stackrel{~}{t}_1}^2,m_{\stackrel{~}{t}_2}^2,|\mu |^2),`$ (4.3)
where $`\alpha _s=g_s^2/(4\pi )`$ is the SU(3)<sub>c</sub> coupling strength, and $`I(a,b,c)`$ is the one-loop function
$$I(a,b,c)=\frac{ab\mathrm{ln}(a/b)+bc\mathrm{ln}(b/c)+ac\mathrm{ln}(c/a)}{(ab)(bc)(ac)}.$$
(4.4)
The $`b`$-quark Yukawa coupling $`h_b(Q_{tb})`$ is then related to the running $`b`$-quark mass $`m_b(Q_{tb})`$ by
$$h_b=\frac{g_wm_b}{\sqrt{2}M_W\mathrm{cos}\beta [\mathrm{\hspace{0.17em}1}+\delta h_b/h_b+(\mathrm{\Delta }h_b/h_b)\mathrm{tan}\beta ]},$$
(4.5)
where $`\delta h_b/h_b`$ and $`\mathrm{\Delta }h_b/h_b`$ are given in (4.2) and (4.3), respectively. The running $`b`$-quark mass $`m_b(Q_{tb})`$ is obtained by means of the RG running of the $`b`$-quark mass from the scale $`m_b`$. In (4.5), we have redefined the right-handed $`b`$-quark superfield, so that the physical $`b`$-quark mass is positive. Under such a field redefinition, only the Yukawa coupling $`h_b`$ becomes complex, while the phases of $`\delta h_b/h_b`$ and $`\mathrm{\Delta }h_b/h_b`$ as well as those of the supersymmetry-breaking parameters do not change. Moreover, since only the moduli of the Yukawa couplings $`h_b`$ and $`h_t`$ enter the field-dependent quark and squark masses in (2.10) and (2), the neutral Higgs-boson mass matrices remain unaffected by the above field redefinition. Also, we have checked that the very same property of invariance under rephasings of $`h_t`$ and $`h_b`$ persists for the charged Higgs-boson mass matrix as well. At this point, it is interesting to observe that the sum $`\delta h_b+\mathrm{\Delta }h_b\mathrm{tan}\beta `$ in (4.5) receives two sorts of quantum corrections, one originating from QCD effects and another from a chargino-mediated graph. The QCD correction is proportional to the hermitean conjugate of the sbottom-mixing parameter $`X_b=A_b\mu ^{}\mathrm{tan}\beta `$, whilst the chargino-induced diagram depends linearly on the stop-mixing parameter $`X_t=A_t\mu ^{}\mathrm{cot}\beta `$.
The effective Lagrangian describing the $`t`$-quark Yukawa coupling is given by
$$_{\varphi ^0\overline{t}t}=\mathrm{\Delta }h_t\varphi _1^0\overline{t}_Rt_L+(h_t+\delta h_t)\varphi _2^0\overline{t}_Rt_L+\mathrm{h}.\mathrm{c}..$$
(4.6)
The corresponding relation for $`h_t`$ as a function of $`m_t`$ may easily be determined analogously by the effective Lagrangian (4.6), and reads
$$h_t=\frac{g_wm_t}{\sqrt{2}M_W\mathrm{sin}\beta [\mathrm{\hspace{0.17em}1}+\delta h_t/h_t+(\mathrm{\Delta }h_t/h_t)\mathrm{cot}\beta ]},$$
(4.7)
with
$`{\displaystyle \frac{\mathrm{\Delta }h_t}{h_t}}`$ $`=`$ $`{\displaystyle \frac{2\alpha _s}{3\pi }}m_{\stackrel{~}{g}}^{}\mu ^{}I(m_{\stackrel{~}{t}_1}^2,m_{\stackrel{~}{t}_2}^2,|m_{\stackrel{~}{g}}|^2)+{\displaystyle \frac{|h_b|^2}{16\pi ^2}}A_b^{}\mu ^{}I(m_{\stackrel{~}{b}_1}^2,m_{\stackrel{~}{b}_2}^2,|\mu |^2),`$ (4.8)
$`{\displaystyle \frac{\delta h_t}{h_t}}`$ $`=`$ $`{\displaystyle \frac{2\alpha _s}{3\pi }}m_{\stackrel{~}{g}}^{}A_tI(m_{\stackrel{~}{t}_1}^2,m_{\stackrel{~}{t}_2}^2,|m_{\stackrel{~}{g}}|^2){\displaystyle \frac{|h_b|^2}{16\pi ^2}}|\mu |^2I(m_{\stackrel{~}{b}_1}^2,m_{\stackrel{~}{b}_2}^2,|\mu |^2).`$ (4.9)
As in the case of the $`b`$-quark Yukawa coupling, we have to make a judicious phase rotation of the right-handed $`t`$-quark superfield, such that the physical top-quark mass becomes positive. Again, one can show that such a field redefinition does not change the analytic results of the RG analysis.
At this stage, it is important to remark that, within the RG-resummation approach described in Section 3, the non-logarithmic corrections must be treated as threshold effects and hence they should only contribute to the RG-frozen part of the Higgs-boson mass matrices, generically denoted as $`\mathrm{\Delta }^2(Q_{tb})`$. Therefore, the decoupling procedure for the heavy squark states requires that the effective $`b`$\- and $`t`$-quark Yukawa couplings given by (4.5) and (4.7) are evaluated at the scale $`Q_{tb}`$. As we discuss in Section 5, these additional Yukawa corrections can lead to observable effects in Higgs-boson searches.
It is now straightforward to obtain the interaction Lagrangians of the Higgs-boson mass eigenstates $`H_i`$ to the up- and down-type quarks, collectively denoted as $`u`$ and $`d`$. Taking into account both CP-violating self-energy and vertex effects, we find
$$_{H\overline{f}f}=\underset{i=1}{\overset{3}{}}H_i\left[\frac{g_wm_d}{2M_W}\overline{d}\left(g_{H_idd}^S+ig_{H_idd}^P\gamma _5\right)d+\frac{g_wm_u}{2M_W}\overline{u}\left(g_{H_iuu}^S+ig_{H_iuu}^P\gamma _5\right)u\right],$$
(4.10)
with
$`g_{H_idd}^S`$ $`=`$ $`{\displaystyle \frac{1}{h_d+\delta h_d+\mathrm{\Delta }h_d\mathrm{tan}\beta }}\{\mathrm{Re}(h_d+\delta h_d){\displaystyle \frac{O_{1i}}{\mathrm{cos}\beta }}+\mathrm{Re}(\mathrm{\Delta }h_d){\displaystyle \frac{O_{2i}}{\mathrm{cos}\beta }}`$ (4.11)
$`[\mathrm{Im}(h_d+\delta h_d)\mathrm{tan}\beta \mathrm{Im}(\mathrm{\Delta }h_d)]O_{3i}\},`$
$`g_{H_idd}^P`$ $`=`$ $`{\displaystyle \frac{1}{h_d+\delta h_d+\mathrm{\Delta }h_d\mathrm{tan}\beta }}\{[\mathrm{Re}\left(\mathrm{\Delta }h_d\right)\mathrm{Re}(h_d+\delta h_d)\mathrm{tan}\beta ]O_{3i}`$ (4.12)
$`\mathrm{Im}(h_d+\delta h_d){\displaystyle \frac{O_{1i}}{\mathrm{cos}\beta }}\mathrm{Im}(\mathrm{\Delta }h_d){\displaystyle \frac{O_{2i}}{\mathrm{cos}\beta }}\},`$
$`g_{H_iuu}^S`$ $`=`$ $`{\displaystyle \frac{1}{h_u+\delta h_u+\mathrm{\Delta }h_u\mathrm{cot}\beta }}\{\mathrm{Re}(h_u+\delta h_u){\displaystyle \frac{O_{2i}}{\mathrm{sin}\beta }}+\mathrm{Re}(\mathrm{\Delta }h_u){\displaystyle \frac{O_{1i}}{\mathrm{sin}\beta }}`$ (4.13)
$`[\mathrm{Im}(h_u+\delta h_u)\mathrm{cot}\beta \mathrm{Im}(\mathrm{\Delta }h_u)]O_{3i}\},`$
$`g_{H_iuu}^P`$ $`=`$ $`{\displaystyle \frac{1}{h_u+\delta h_u+\mathrm{\Delta }h_u\mathrm{cot}\beta }}\{[\mathrm{Re}(\mathrm{\Delta }h_u)\mathrm{Re}(h_u+\delta h_u)\mathrm{cot}\beta ]O_{3i}`$ (4.14)
$`\mathrm{Im}(h_u+\delta h_u){\displaystyle \frac{O_{2i}}{\mathrm{sin}\beta }}\mathrm{Im}(\mathrm{\Delta }h_u){\displaystyle \frac{O_{1i}}{\mathrm{sin}\beta }}\},`$
where the Higgs scalar and pseudoscalar couplings are normalized with respect to their SM values.
Finally, it is interesting to investigate the behaviour of self-energy- and vertex-type CP violation in the decoupling limit of a heavy charged Higgs boson in the MSSM. Thus, for values of the charged Higgs mass $`M_{H^+}M_Z`$, one has $`O_{31}0`$, while $`O_{11}\mathrm{cos}\beta `$ and $`O_{21}\mathrm{sin}\beta `$. In this limit, the scalar components of the $`H_1dd`$ and $`H_1uu`$ couplings acquire the known SM form, given by $`g_wm_d/(2M_W)`$ and $`g_wm_u/(2M_W)`$, respectively, where
$`m_d`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(h_d+\delta h_d+\mathrm{\Delta }h_d\mathrm{tan}\beta \right)v_1,`$
$`m_u`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(h_u+\delta h_u+\mathrm{\Delta }h_u\mathrm{cot}\beta \right)v_2`$ (4.15)
have already been defined to be positive in (4.5) and (4.7). For similar reasons, the pseudoscalar parts of the $`H_1dd`$ and $`H_1uu`$ couplings vanish, as they are proportional to $`O_{31}`$ and $`\mathrm{Im}m_{d,u}=0`$. On the other hand, in the same large $`M_{H^+}`$ limit, the scalar and pseudoscalar couplings of both the two heaviest Higgs bosons $`H_2`$ and $`H_3`$ to the up and down fermions do not vanish. We can therefore conclude that CP-violating self-energy and vertex effects do not decouple in the heavy neutral Higgs sector. In the next section, we demonstrate explicitly the aforementioned (non-)decoupling features of CP violation by analyzing specific phenomenological examples.
## 5 Phenomenological Discussion
In this Section we discuss the phenomenological implications of radiative Higgs-sector CP violation in the MSSM for Higgs-boson searches at high-energy colliders. We focus our attention on the physics potential for discovering Higgs bosons with mixed CP parities at LEP 2 and the upgraded Tevatron collider, and also comment on the enhanced search capabilities offered by the LHC.
At the LEP 2 and Tevatron colliders, neutral Higgs bosons are predominantly produced via the Higgs-strahlung processes in $`e^+e^{}`$ and $`q\overline{q}`$ collisions, such as $`e^+e^{}Z^{}ZH_i`$ , $`q\overline{q}Z^{}ZH_i`$ and $`q\overline{q}W^\pm ^{}W^\pm H_i`$ , with $`i=1,2,3`$. If the next-to-lightest Higgs boson is not too heavy, Higgs bosons can also be produced copiously in pairs through the reactions: $`e^+e^{}Z^{}H_iH_j`$ and $`q\overline{q}Z^{}H_iH_j`$. In addition to the Higgs-boson masses, the Higgs-boson couplings to the gauge fields play an essential role in our forthcoming discussion. The effective Lagrangians governing the interactions of the Higgs bosons with the $`W^\pm `$ and $`Z`$ bosons are given by
$`_{HVV}`$ $`=`$ $`g_wM_W{\displaystyle \underset{i=1}{\overset{3}{}}}g_{H_iVV}\left(H_iW_\mu ^+W^{,\mu }+{\displaystyle \frac{1}{2\mathrm{cos}^2\theta _w}}H_iZ_\mu Z^\mu \right),`$ (5.1)
$`_{HH^{}W^\pm }`$ $`=`$ $`{\displaystyle \frac{g_w}{2}}{\displaystyle \underset{i=1}{\overset{3}{}}}g_{H_iH^{}W^+}(H_ii\underset{\mu }{\overset{}{}}H^{})W^{+,\mu }+\mathrm{h}.\mathrm{c}.,`$ (5.2)
$`_{HHZ}`$ $`=`$ $`{\displaystyle \frac{g_w}{2\mathrm{cos}\theta _w}}{\displaystyle \underset{j>i=1}{\overset{3}{}}}g_{H_iH_jZ}(H_i\underset{\mu }{\overset{}{}}H_j)Z^\mu ,`$ (5.3)
where $`\mathrm{cos}\theta _w=M_W/M_Z`$, $`\underset{\mu }{\overset{}{}}\underset{\mu }{\overset{}{}}\underset{\mu }{\overset{}{}}`$, and
$`g_{H_iVV}`$ $`=`$ $`\mathrm{cos}\beta O_{1i}+\mathrm{sin}\beta O_{2i},`$ (5.4)
$`g_{H_iH_jZ}`$ $`=`$ $`O_{3i}\left(\mathrm{cos}\beta O_{2j}\mathrm{sin}\beta O_{1j}\right)O_{3j}\left(\mathrm{cos}\beta O_{2i}\mathrm{sin}\beta O_{1i}\right),`$ (5.5)
$`g_{H_iH^{}W^+}`$ $`=`$ $`\mathrm{cos}\beta O_{2i}\mathrm{sin}\beta O_{1i}+iO_{3i}.`$ (5.6)
For completeness, we have included in (5.2) the interactions of the charged Higgs bosons $`H^\pm `$ with the neutral Higgs and $`W^{}`$ bosons. Note that the couplings $`H_iZZ`$ and $`H_iW^+W^{}`$ are related to the $`H_iH_jZ`$ couplings through
$$g_{H_kVV}=\epsilon _{ijk}g_{H_iH_jZ}.$$
(5.7)
Moreover, unitarity provides the constraint
$$\underset{i=1}{\overset{3}{}}g_{H_iVV}^2=1.$$
(5.8)
Evidently, if two Higgs-boson couplings to gauge bosons are known, this is sufficient to determine the complete set of the couplings $`g_{H_iVV}`$ and $`g_{H_iH_jZ}`$ .
In the above calculation of the effective Higgs-gauge-boson couplings, we have assumed that the dominant contributions arise from Higgs-mixing effects. As opposed to the $`b`$-quark Yukawa case discussed in Section 4, proper vertex corrections to the $`H_iZZ`$ and $`H_iH_jZ`$ couplings do not contain strong-coupling- or $`\mathrm{tan}\beta `$-enhanced diagrams. Therefore, naive dimensional analysis suggests that these corrections are suppressed relative to their tree-level values by loop factors of the kind: $`(3\alpha _w/4\pi )(|\mu A_t|/m_{\stackrel{~}{t}_1}^2),(3\alpha _w/4\pi )(|\mu A_t|^2v^2/m_{\stackrel{~}{t}_1}^6)\stackrel{<}{_{}}10^2`$, where $`\stackrel{~}{t}_1`$ is the heaviest stop squark (see (2)). In the following, we neglect proper vertex corrections to the $`H_iZZ`$ and $`H_iH_jZ`$ couplings.
For our phenomenological discussion of CP violation, we consider the following two representative values for $`\mathrm{tan}\beta `$: (i) $`\mathrm{tan}\beta =4`$ and (ii) $`\mathrm{tan}\beta =20`$. For definiteness, unless stated otherwise, the soft supersymmetry-breaking and $`\mu `$ parameters are set to the values
$`M_{\mathrm{SUSY}}`$ $`=`$ $`\stackrel{~}{M}_Q=\stackrel{~}{M}_t=\stackrel{~}{M}_b=0.5\mathrm{TeV},\mu =2\mathrm{TeV},|A_t|=|A_b|=1\mathrm{TeV},`$
$`|m_{\stackrel{~}{B}}|`$ $`=`$ $`|m_{\stackrel{~}{W}}|=0.3\mathrm{TeV},|m_{\stackrel{~}{g}}|=1\mathrm{TeV},`$ (5.9)
As can be seen in (5), we have chosen relatively large values for the stop and sbottom mixing parameters $`A_t`$, $`A_b`$ and $`\mu `$, as well as a common left- and right-handed squark mass $`M_{\mathrm{SUSY}}`$, which leads to enhanced CP-violating effects of the CP-odd phases arg $`(\mu A_{t,b})`$ on the Higgs sector.
As was mentioned in the Introduction, large CP-odd phases may lead to rather large EDM contributions, thereby violating the known upper bounds on the electron and neutron EDMs $`d_e`$ and $`d_n`$: $`d_e/e<\mathrm{0.5\; 10}^{26}`$ cm and $`d_n/e<\mathrm{0.6\; 10}^{25}`$ cm at the 2-$`\sigma `$ level. One phenomenologically interesting possibility for avoiding the possible CP crisis is to make the first two generations of squarks rather heavy with masses much above the TeV scale , keeping the third generation relatively light with masses of order 0.5 TeV. In such a scenario, CP violation may only reside in the third generation. For our illustrations, we shall take the $`\mu `$ parameter to be real and assume that the only CP-odd phases in the theory are $`\mathrm{arg}(A_t)=\mathrm{arg}(A_b)`$ and $`\mathrm{arg}(m_{\stackrel{~}{g}})`$. Again, as one way to avoid the one-loop EDM constraints , we have taken a gluino mass of 1 TeV in (5). However, in such a scheme, one has to worry about the fact that Higgs-boson two-loop contributions to the EDMs might still become sizeable. For the low-$`\mathrm{tan}\beta `$ scenario in (5), the two-loop EDM contributions are of the order of the experimental EDM upper bounds mentioned above. Since these two-loop EDM effects depend almost linearly on $`\mathrm{tan}\beta `$, one might then need to arrange for cancellations among the different one- and two-loop EDM terms at the level of 10% for the scenario with $`\mathrm{tan}\beta =20`$. We believe that this can be achieved without excessive fine-tuning of the CP-violating parameters of the theory.
As was already noticed in , the radiative corrections to the lightest Higgs boson $`H_1`$ depend crucially on the stop mixing parameter $`|X_t|=|A_t\mu ^{}\mathrm{cot}\beta |`$. Specifically, the radiatively-corrected $`H_1`$-boson mass increases as $`|X_t|`$ increases, reaching a maximum when $`|X_t|/M_{\mathrm{SUSY}}2.45`$. Then, as $`|X_t|`$ further increases, the radiative corrections to $`H_1`$-boson mass decrease and may even become negative, driving the latter to very small, experimentally excluded values. A distinctive feature of the CP-violating SUSY scenario compared to the CP-conserving one is that $`|X_t|`$ can be increased by varying only the phase $`\mathrm{arg}(A_t)`$ from zero to higher values, but holding fixed $`|A_t|`$ and $`|\mu |`$. For similar reasons, high values of $`|X_t|`$ induced by large values of $`\mathrm{arg}(A_t)`$ can make the mass of the lightest stop squark $`\stackrel{~}{t}_2`$ very low, so as to violate present experimental constraints, i.e., $`m_{\stackrel{~}{t}_2}\stackrel{>}{_{}}100`$ GeV. Furthermore, light stop quarks, with $`m_{\stackrel{~}{t}_2}\stackrel{<}{_{}}300`$ GeV and large $`|X_t|`$ values, can give rise to observably large contributions to the electroweak precision parameter $`\mathrm{\Delta }\rho `$ . For the scenarios under discussion, $`m_{\stackrel{~}{t}_2}`$ is always larger than about 300 GeV, for all the parameter space for which the $`H_1`$-boson mass acquires acceptable values. Therefore, apart from the bounds derived by EDM constraints, the requirement that the lightest Higgs-boson mass is positive can be used naively to set an upper bound on the phase of $`A_t`$.
In Fig. 1 we give numerical predictions for the two lightest Higgs-boson masses $`M_{H_1}`$ and $`M_{H_2}`$, and for the three relevant $`H_iZZ`$ couplings squared as a function of $`\mathrm{arg}(A_t)`$, for two different values of $`\mathrm{arg}(m_{\stackrel{~}{g}})`$: $`\mathrm{arg}(m_{\stackrel{~}{g}})=0`$ (solid lines) and $`\mathrm{arg}(m_{\stackrel{~}{g}})=\pi /2`$ (dashed lines). We first discuss the scenario with $`\mathrm{tan}\beta =4`$, for which the values of the remaining soft supersymmetry-breaking parameters and $`\mu `$ are given in (5). Since our interest is to analyze dominant CP-violating effects for a light Higgs sector, we present predictions for a relatively small charged-Higgs-boson mass, $`M_{H^+}=150`$ GeV. In the CP-conserving limit of the theory ($`\mathrm{arg}(A_t)=0`$), the mass $`M_{H_1}`$ of the lightest neutral Higgs boson is close to 85 GeV, while the square of the $`H_1ZZ`$ coupling, $`g_{H_1ZZ}^2`$, is approximately equal to 0.8. These values of masses and couplings are now excluded by Higgs-boson searches at LEP 2 . However, this situation changes crucially once CP-violating phases become relevant. As the phase of $`A_t`$ increases, two important effects take place. First, as was mentioned above, the stop mixing parameter $`|X_t|`$ becomes larger, giving rise to larger $`H_1`$-boson masses. Second, the mass-matrix terms describing the scalar-pseudoscalar mixing are enhanced, thereby effectively leading to large modifications in the couplings of the Higgs bosons to gauge bosons. This second effect can be attributed entirely to CP violation. In fact, as can be seen from Fig. 1(b), for $`\mathrm{arg}(A_t)80`$ degrees, $`g_{H_1ZZ}^2`$ gets very suppressed, implying that LEP 2 cannot detect the Higgs boson $`H_1`$ via $`e^+e^{}Z^{}ZH_1`$. On the other hand, for the same range of values of $`\mathrm{arg}(A_t)`$, i.e., $`\mathrm{arg}(A_t)=80^{}`$–95, the $`H_2`$ and $`H_3`$ bosons have significant couplings to the $`Z`$ bosons. Although the $`H_3`$ boson is too heavy to be detected at LEP 2 in this case, the $`H_2`$ boson has a mass of $`105`$$`110`$ GeV and $`g_{H_2ZZ}^20.8`$–0.6, which may be probed at LEP 2 in this year’s run. For larger values of the phase of $`A_t`$, $`95^{}<\mathrm{arg}(A_t)<110^{}`$, the discovery of a Higgs boson at LEP 2 is more challenging. The lightest Higgs boson $`H_1`$ acquires a mass below 90 GeV, but the $`H_1ZZ`$ coupling is too small to allow experimental detection through the reaction $`e^+e^{}Z^{}ZH_1`$. In addition, the $`H_2`$ boson becomes too heavy to allow for discovery via $`e^+e^{}Z^{}ZH_2`$, with $`g_{H_2ZZ}^2\stackrel{<}{_{}}0.7`$.
As was discussed in , the $`H_1`$ and $`H_2`$ bosons may also be searched for in the channel $`e^+e^{}Z^{}H_1H_2`$. Since the squared coupling $`g_{H_1H_2Z}^2=g_{H_3ZZ}^20.2`$ almost independently of $`\mathrm{arg}(A_t)`$, a careful experimental analysis of the parameter region of interest will be necessary to determine whether the two lightest Higgs bosons can be observed for such large mass differences ($`M_{H_1}M_{H_2}40`$ GeV) and such small $`g_{H_1H_2Z}`$ couplings.
In Fig. 1 we also present predictions for a gluino phase of 90 degrees (dashed lines). Since the vertex corrections are generally small for low or moderate values of $`\mathrm{tan}\beta `$, they are expected to induce only small corrections to the Higgs-boson masses and mixings. This last fact is reflected in Fig. 1, even though the coupling $`g_{H_1ZZ}^2`$ ($`g_{H_2ZZ}^2`$) gets slightly smaller (larger) for larger values of the phase of $`A_t`$.
Fig. 2 shows the changes in the predictions for the same choice of parameters as in Fig. 1, but with $`M_{\mathrm{SUSY}}`$, $`\mu `$ and $`|A_t|=|A_b|`$ rescaled by a factor of 2. This rescaling leads to a slight increase (decrease) of the Higgs-boson mass $`M_{H_1}`$ ($`M_{H_2}`$), while $`g_{H_1ZZ}^2`$ exhibits a slightly different quantitative dependence on the phase of $`A_t`$, especially for the region in which $`g_{H_1ZZ}^2`$ is very small. Thus, although $`M_{H_1}`$ becomes small for $`\mathrm{arg}(A_t)>115`$ degrees, the $`H_1ZZ`$ coupling gets sizeable again, well within the capabilities of LEP 2 to test.
We now investigate more quantitatively the predictions of a large-$`\mathrm{tan}\beta `$ scenario for the Higgs-boson masses and couplings. We adopt the scenario given in (5), with $`\mathrm{tan}\beta =20`$ and $`M_{H^+}=150`$ GeV. In this large-$`\mathrm{tan}\beta `$ scenario, one has $`|\mu |\mathrm{cot}\beta |A_t|`$, and the effective stop mixing parameter $`|X_t||A_t|`$ is almost independent of $`\mathrm{arg}(A_t)`$. Therefore, as is seen in Fig. 3(a), the Higgs-boson masses $`M_{H_1}`$ and $`M_{H_2}`$ do not exhibit any significant variation as a function of $`\mathrm{arg}(A_t)`$. In contrast to the Higgs-boson masses, Fig. 3(b) shows that there is a non-trivial dependence of the squared couplings $`g_{H_iZZ}^2`$ on $`\mathrm{arg}(A_t)`$. Furthermore, the next-to-lightest Higgs boson $`H_2`$ is heavy enough to render its search through the $`e^+e^{}Z^{}ZH_2`$ reaction kinematically inaccessible at LEP 2. For similar reasons, we find that, for all values of $`\mathrm{arg}(A_t)`$, the $`H_1`$ and $`H_2`$ bosons are rather too heavy to be produced via the $`e^+e^{}Z^{}H_1H_2`$ channel at LEP 2. As a result, Higgs-boson searches at LEP 2 tend to be more efficient for small values of the $`A_t`$ phases, for which the lightest neutral Higgs-boson mass is close to 100 GeV and $`g_{H_1ZZ}^2`$ is non-negligible ($`g_{H_1ZZ}^2`$ 0.3). In addition, for large $`A_t`$ phases, $`\mathrm{arg}(A_t)\stackrel{>}{_{}}80^{}`$, the $`H_1VV`$ coupling (with $`V=Z,W`$) is rather suppressed, so that the $`H_1`$ Higgs boson, although it becomes lighter with a mass in the range 90–95 GeV, will be elusive at LEP 2, and may also escape detection via the corresponding channel at the upgraded Tevatron. However, the next-to-lightest Higgs boson $`H_2`$ has couplings of order unity to the $`Z`$ and $`W`$ bosons. Present simulations show that a neutral Higgs boson, such as $`H_2`$, with $`M_{H_2}180`$ GeV and a SM-like coupling strength to vector gauge bosons can be tested at the Tevatron collider with a total integrated luminosity of 10 fb<sup>-1</sup>. However, discovery of such a Higgs boson at the 5-$`\sigma `$ level would demand a total integrated luminosity of 30 fb<sup>-1</sup>, and would have a reach up to $`M_{H_2}120`$ GeV . Finally, even though the $`H_1H_3Z`$ coupling is close to unity ($`g_{H_1H_3Z}^2=g_{H_2ZZ}^2`$) for $`\mathrm{arg}(A_t)>100^{}`$, the Higgs-pair production of $`H_1`$ and $`H_3`$ is not kinematically allowed at LEP 2, since $`M_{H_3}M_{H^+}`$. Further studies will be necessary to investigate the potential of this production mechanism at the Tevatron.
It is interesting to present predictions for the neutral Higgs-boson masses and their couplings to gauge bosons for lower values of the charged Higgs-boson mass $`M_{H^+}`$ in the above large-$`\mathrm{tan}\beta `$ scenario. In Fig. 4, we plot numerical estimates for the same kinematic parameters as in Fig. 3, but with $`M_{H^+}=135`$ GeV. In this case, the $`H_1`$-boson mass varies approximately between 80 and 65 GeV, and the $`H_1ZZ`$ coupling rapidly decreases as the phase of $`A_t`$ increases. The two heaviest neutral Higgs bosons $`H_2`$ and $`H_3`$ have masses in the range between 120 and 130 GeV. Hence, these two Higgs bosons cannot be produced via $`e^+e^{}Z^{}ZH_2`$ or $`e^+e^{}Z^{}ZH_3`$ at LEP 2. However, the $`H_2`$ and $`H_3`$ bosons may still be accessed via $`e^+e^{}Z^{}H_1H_2`$ or $`H_1H_3`$. Interestingly, the squared couplings $`g_{H_1H_2Z}^2=g_{H_3ZZ}^2`$ and $`g_{H_1H_3Z}^2=g_{H_2ZZ}^2`$ exhibit a cross-over as a function of $`\mathrm{arg}(A_t)`$. The crossing point of the two squared couplings is when $`\mathrm{arg}(A_t)90^{}`$. For $`\mathrm{arg}(A_t)=180^{}`$, one of the squared couplings goes to 0 and the other to 1, depending on the phase of the gluino mass. In this case, the two heaviest neutral Higgs bosons become almost degenerate in mass. For the whole range of values of $`\mathrm{arg}(A_t)`$, either the $`H_2`$ or $`H_3`$ Higgs boson can be tested at the upgraded Tevatron collider provided a total integrated luminosity of 10 fb<sup>-1</sup> per detector is available
Figs. 5(a) and (b) show the degree of mass splitting between the two heaviest neutral Higgs bosons $`H_2`$ and $`H_3`$, for the same choice of parameters as in Figs. 1 and 2, but for $`M_{H^\pm }=`$ 200, 300, 400, 500 GeV. As was already observed in , even though the $`H_2`$ and $`H_3`$ bosons are almost degenerate in the CP-conserving limit of the theory, they can have a degree of splitting up to 30% for a maximal CP-violating phase $`\mathrm{arg}(A_t)90^{}`$. The comparison of the Fig. 5(a) with (b) reveals that this last result is almost independent of the common scale factor of $`M_{\mathrm{SUSY}}`$, $`\mu `$ and $`|A_t|`$. Also, the degree of mass splitting is not much affected by the value of the gluino phase $`\mathrm{arg}(m_{\stackrel{~}{g}})`$: the results for $`\mathrm{arg}(m_{\stackrel{~}{g}})=90^{}`$ are slightly higher than those of $`\mathrm{arg}(m_{\stackrel{~}{g}})=0`$. In this vein, it is interesting to mention that large CP-violating scalar-pseudoscalar mixings can lead to observable phenomena of resonant CP violation at high-energy colliders .
In Figs. 6 and 7, we examine the behaviours of the scalar and pseudoscalar parts of the $`H_1bb`$ coupling as functions of the CP-odd phase $`\mathrm{arg}(A_t)`$, for two different charged Higgs-boson masses, $`M_{H^+}=150`$ and 300 GeV. As was done in , we find that the best way of analyzing such a behaviour is in terms of the CP-even and CP-odd quantities: $`[(g_{H_1bb}^S)^2+(g_{H_1bb}^P)^2]`$ and $`2g_{H_1bb}^Sg_{H_1bb}^P/[(g_{H_1bb}^S)^2+(g_{H_1bb}^P)^2]`$. For example, Higgs-boson branching ratios are proportional to the first quantity, while the second one will only occur in CP-violating observables. In other words, $`2g_{H_1bb}^Sg_{H_1bb}^P/[(g_{H_1bb}^S)^2+(g_{H_1bb}^P)^2]`$ gives a measure of the CP-violating component in the $`H_1bb`$ coupling. If we compare the predictions for $`M_{H^+}=150`$ GeV with those for $`M_{H^+}=300`$ GeV in Figs. 6 and 7, we find that the CP-violating component of the $`H_1bb`$ coupling reduces in magnitude, for large values of the charged Higgs-boson mass. Such a decoupling behaviour of the CP-violating $`H_1bb`$ component is in agreement with our observation, which we already made at the end of Section 4. From Figs. 6(b) and 7(b), we see that the impact of the gluino phase on the CP-violating component of the $`H_1bb`$ coupling is more important for the large-$`\mathrm{tan}\beta `$ scenario. This may be attributed to the fact that the radiatively-induced term $`\mathrm{\Delta }h_b\mathrm{tan}\beta `$, which crucially depends on the gluino phase and $`\mathrm{tan}\beta `$, has a dominant contribution to the $`H_1bb`$ coupling.
There can be a cancellation or a strong suppression of the coupling of the lightest Higgs boson $`H_1`$ to the bottom quarks, depending on the magnitude of the CP-violating phases and of the products $`A_t\mu `$, $`A_b\mu `$ and $`m_{\stackrel{~}{g}}\mu `$. This cancellation usually takes place for moderate values of the charged Higgs mass and large values of $`\mathrm{tan}\beta `$. Such an effect is also present in the CP-conserving case, for specific signs and magnitudes of the above products involving the trilinear terms $`A_{t,b}`$ and the gluino mass, and has been discussed in detail in . Figure 8(b) illustrates such a cancellation for the CP-violating SUSY model under discussion. For example, we observe that for the set of SUSY parameters considered in Fig. 8, the $`H_1bb`$ coupling can be strongly suppressed for $`\mathrm{arg}(A_t)=\mathrm{arg}(A_b)15^{}`$. Moreover, we have checked that for this same set of parameters the $`H_1ZZ`$ coupling is almost SM-like. In addition, Fig. 8(a) shows that the mass of the lightest Higgs boson is of order 105 GeV, practically independent of the CP-violating phase. This is therefore an extremely interesting example, since the lightest Higgs mass is in the mass range that may be within the reach of LEP 2, and its production cross section will be SM-like. However, the main decay channel, $`H_1b\overline{b}`$, can be strongly suppressed if the CP-violating phases $`\mathrm{arg}(A_t)`$ and $`\mathrm{arg}(m_{\stackrel{~}{g}})`$ lie in a specific range. Therefore, in such a scenario, the detection of the $`H_1`$ boson may in principle be impossible, even in the final run of LEP 2. To make a conclusive statement on this possibility, one should study in detail the capability of LEP 2 to detect such a light $`H_1`$ boson via its decays into $`\tau `$ pairs, or into other hadronic modes. Most intriguingly, the set of parameters considered in this example also allow for a light right-handed stop squark and moderate mixing parameter, $`|X_t|/\stackrel{~}{M}_Q`$, of the type necessary to allow the possibility of electroweak baryogenesis . Hence, if such a Higgs boson cannot be discovered at LEP 2 via other decay channels, the final phase of LEP 2 will leave an open window for electroweak baryogenesis. A careful study of the CP-violating phases required for electroweak baryogenesis and the detection capabilities of LEP 2 for alternative decay modes becomes essential for testing this exciting scenario.
We shall briefly comment on the enhanced LHC capabilities for Higgs-boson searches . The LHC has a considerably higher reach than LEP 2 and the upgraded Tevatron collider in the search for heavier Higgs bosons, and hence has more chances to unravel the complete Higgs-boson spectrum of the MSSM with explicit CP violation. At the LHC, Higgs bosons may be copiously produced via a wide variety of processes which depend in many different ways on the couplings of the neutral and charged Higgs bosons both to gauge bosons and fermions . In the case of the CP-violating version of the MSSM under study, we have shown that mixing between states with different CP parities can dramatically modify those couplings and, hence, importantly affect the associated production and decay mechanisms. Studies including CP-violating effects on the gluon-fusion production of Higgs bosons at the LHC have already been considered in the literature . Our work provide the basic tool to improve further those studies, and to perform a complete analysis of the CP-violating effects on the many other Higgs-boson search mechanisms available at LHC. The LHC, together with the information gathered from experiments at LEP 2 and the upgraded Tevatron collider, will be capable of providing a thorough test of the MSSM Higgs sector and shed light on the possibility of explicit radiative breaking of CP invariance in supersymmetry.
Finally, it is interesting to make a comparative analysis between our results and those obtained previously in . In the latter work, the effective RG-improved potential was expanded up to renormalizable operators of dimension 4. The expansion was performed in powers of the stop-mass splitting, $`m_{\stackrel{~}{t}_1}^2m_{\stackrel{~}{t}_2}^2`$, relative to the arithmetic average of the squared stop masses, $`M_{\mathrm{SUSY}}^2=\frac{1}{2}(m_{\stackrel{~}{t}_1}^2+m_{\stackrel{~}{t}_2}^2)`$. Moreover, the two-loop effect originating from the one-loop radiative corrections to the Yukawa couplings of the top and bottom quarks was not taken into account in the computation of the effective potential in . Nevertheless, for moderate values of $`\mathrm{tan}\beta `$ and for all soft supersymmetry-breaking masses equal to $`M_{\mathrm{SUSY}}`$, the deviations of the results presented in with our results are expected to be small, for small values of the stop-mixing parameter $`X_t=|A_t\mu ^{}\mathrm{cot}\beta |`$; the deviations will only grow for increasing values of $`X_t`$. For larger values of $`\mathrm{tan}\beta `$, instead, the impact of the bottom-mass quantum corrections, which were omitted in , is significant. In fact, only in the limit of small values of $`|\mu |`$, in which case the bottom-mass corrections are small, are both approaches guaranteed to give comparable numerical estimates.
In Fig. 9, we show the predictions for the mass of the lightest neutral Higgs boson $`H_1`$ and its coupling to $`Z`$ bosons, as obtained obtained by our RG approach (solid lines) and the operator-expansion method of (dashed lines). For the sake of comparison, we consider the same input parameters as those chosen in Fig. 1, for vanishing gluino phase and three different values of the charged Higgs-boson mass: $`M_{H^+}=150`$, 200 and 500 GeV. As was discussed above, we find that the predictions of the two works are in excellent agreement with one another for small values of $`X_t`$. For large values of $`X_t`$, instead, we observe larger quantitative differences in the results obtained by the two approaches, even though the qualitative behaviour of the two predictions exhibits a quite analogous functional dependence.
Our one-loop RG-improved approach overcomes the limitations present in earlier analyses. In particular, our RG approach allows for a rather precise determination of the radiative effects on a generic Higgs-boson mass spectrum, even in cases of large stop mixings and/or large hierarchies between the left- and right-handed stop masses. Also, within our RG approach, the important effects of the one-loop corrections to the quark Yukawa couplings are incorporated in the computation of the Higgs-boson masses and in their respective Higgs-boson couplings to gauge and fermion fields. A Fortran code that computes the Higgs-boson masses and couplings, including all the CP-violating effects as presented in this work, may be found in .
## 6 Conclusions
We have performed a complete one-loop RG improvement of the effective Higgs potential in the MSSM, in which CP violation is induced radiatively by soft CP-violating trilinear interactions that involve the Higgs fields and the stop and sbottom squarks. Earlier studies of the neutral Higgs-boson mass spectrum were based on a number of particular assumptions and/or kinematic approximations. The present work goes well beyond those studies, and extends the most detailed analysis , in which the one-loop RG-improved effective potential was expanded up to renormalizable operators of dimension 4, assuming a moderate mass splitting among the stop squarks. This assumption seemed to impose a serious limitation, given that CP-violating effects exhibit an enhanced behaviour for large values of the stop-mixing parameter $`X_t`$. The results obtained using the present RG approach confirm, however, the qualitative phenomenological features found in , for the Higgs-boson masses and their couplings to fermions and to the $`W^\pm `$ and $`Z`$ bosons. It offers very accurate predictions, at the same level as the most accurate calculations in the CP-conserving case (implying an uncertainty of order 3 GeV) , for the Higgs-boson masses and for the whole range of the MSSM parameter space in the presence of non-trivial CP-violating phases. More specifically, the present study also includes two-loop leading logarithms associated with QCD effects and $`t`$\- and $`b`$-quark Yukawa couplings. It also contains all dominant two-loop non-logarithmic contributions to the one-loop effective potential, which are induced by one-loop threshold effects on the $`t`$\- and $`b`$-quark Yukawa couplings due to the decoupling of the third-generation squarks (as considered in in the CP-conserving limit). These one-loop threshold terms, $`\delta h_{u,d}`$ and $`\mathrm{\Delta }h_{u,d}`$, strongly depend on the phase of the gluino mass and so introduce new CP violation into the MSSM effective potential at the two-loop level.
Large radiative effects of CP violation in the Higgs sector of the MSSM can have important phenomenological consequences on Higgs-boson searches at LEP 2, the Tevatron and the LHC. We have explicitly demonstrated that the radiatively-induced CP violation in the MSSM Higgs potential can lead to important effects of mass and coupling level crossing among the three neutral Higgs particles. These CP-violating effects of level crossing in the Higgs sector modify drastically the Higgs-boson couplings to the up- and down-type quarks, and to the $`W^\pm `$ and $`Z`$ bosons. In particular, CP violation in the lightest Higgs sector becomes relevant for a relatively light charged Higgs boson, with $`M_{H^+}\stackrel{<}{_{}}160`$ GeV. For instance, for $`M_{H^+}=150`$ GeV and $`\mathrm{tan}\beta =4`$, even a neutral Higgs boson as light as 60 GeV may escape detection at LEP 2. However, the upgraded Tevatron may have the physics potential to explore such CP-violating scenarios at low $`\mathrm{tan}\beta `$ values, which may remain at the edge of accessibility even during the final LEP 2 run.
We have also studied the effects induced by a non-trivial CP-odd gluino phase, which enters the effective potential at the two-loop level. The presence of a gluino phase gives rise to small but non-negligible changes in the Higgs-boson mass spectrum and in the couplings of the Higgs fields to the $`W^\pm `$ and $`Z`$ bosons. In this context, we have also found that the product of the scalar times the pseudoscalar coupling of the lightest Higgs boson $`H_1`$ to the bottom quarks has a very strong dependence on the gluino phase. This product of couplings gives a measure of CP violation in the $`H_1bb`$ coupling, which can even be of order unity for relatively small charged Higgs-boson masses.
Finally, it is worth stressing that CP violation decouples from the lightest Higgs sector in the large-mass limit of a heavy charged Higgs boson. This decoupling property, which was known to hold for the CP-violating self-energy effect , has now been shown to be valid for the CP-violating vertex effects as well. As a result, the predictions for the lightest Higgs-boson mass and its couplings to gauge bosons in the above decoupling regime of the theory will be practically identical to the corresponding predictions in the CP-conserving case. However, unlike the lightest Higgs sector, CP violation does not decouple from the heaviest Higgs sector in the MSSM, opening up new possibilities for studying enhanced effects in CP-violating Higgs scalar-pseudoscalar transitions at the LHC or future muon colliders , where the heavy MSSM Higgs bosons can be resonantly produced.
In conclusion: the present analysis has shown that the MSSM with explicit radiative breaking of CP invariance constitutes a very rich theoretical framework, introducing new challenges in the search for fundamental Higgs scalars at LEP 2, the Tevatron and the LHC.
### Acknowledgements
We thank Manuel Drees for useful discussions. A.P. thanks the theory group of Fermilab for the kind hospitality extended to him while part of this work was done. Work supported in part by the U.S. Department of Energy, High Energy Physics Division, under Contract W-31-109-Eng-38.
### Note Added
After completion of the work described here, we saw , which also computes the one-loop RG-improved effective potential of the MSSM with explicit CP violation. Here we further i) perform a more complete RG improvement of the effective potential, ii) develop a self-consistent treatment of the whole MSSM Higgs sector, taking into account the crucial one-loop relation between the charged Higgs-boson and neutral Higgs-boson mass matrices, and iii) include the dominant two-loop non-logarithmic corrections to the effective potential, which may also have an important impact on the $`H_1bb`$ coupling. In the limits where a comparison between the results was possible, we find reasonable agreement with .
## Appendix A Derivatives of Background-Field-Dependent <br>Masses
In this Appendix, we list analytic expressions pertaining to derivatives of quark and squark masses with respect to their background Higgs fields. These derivative expressions are very useful, as they constitute the building blocks of the general one-loop effective potential presented in Section 2. We have divided the Appendix into three subsections. In the first subsection, we list the derivatives of quark masses with respect to Higgs fields, while the next two subsections contain the corresponding expressions for derivatives of squark masses with respect to neutral and charged Higgs fields, respectively.
### A.1 Quark Derivatives
First, we give the derivatives related to non-vanishing tadpole contributions:
$$\frac{1}{v_2}<\frac{\overline{m}_t^2}{\varphi _2}>=|h_t|^2,\frac{1}{v_1}<\frac{\overline{m}_b^2}{\varphi _1}>=|h_b|^2,$$
(A.1)
where the operation $`<\mathrm{}>`$ denotes that the above expressions should be evaluated in the ground state of the Higgs potential.
Then, the self-energy-type derivatives involving neutral and charged Higgs bosons may be listed as follows:
$`<{\displaystyle \frac{^2\overline{m}_t^2}{\varphi _2^2}}>`$ $`=`$ $`<{\displaystyle \frac{^2\overline{m}_t^2}{a_2^2}}>=|h_t|^2,<{\displaystyle \frac{^2\overline{m}_b^2}{\varphi _1^2}}>=<{\displaystyle \frac{^2\overline{m}_b^2}{a_1^2}}>=|h_b|^2,`$
$`<{\displaystyle \frac{^2\overline{m}_t^2}{\varphi _1^+\varphi _1^{}}}>`$ $`=`$ $`{\displaystyle \frac{|h_b|^2m_t^2}{m_t^2m_b^2}},<{\displaystyle \frac{^2\overline{m}_b^2}{\varphi _1^+\varphi _1^{}}}>={\displaystyle \frac{|h_b|^2m_b^2}{m_b^2m_t^2}},`$
$`<{\displaystyle \frac{^2\overline{m}_t^2}{\varphi _2^+\varphi _2^{}}}>`$ $`=`$ $`{\displaystyle \frac{|h_t|^2m_t^2}{m_t^2m_b^2}},<{\displaystyle \frac{^2\overline{m}_b^2}{\varphi _2^+\varphi _2^{}}}>={\displaystyle \frac{|h_t|^2m_b^2}{m_b^2m_t^2}},`$ (A.2)
$`<{\displaystyle \frac{^2\overline{m}_t^2}{\varphi _1^+\varphi _2^{}}}>`$ $`=`$ $`<{\displaystyle \frac{^2\overline{m}_t^2}{\varphi _2^+\varphi _1^{}}}>=<{\displaystyle \frac{^2\overline{m}_b^2}{\varphi _1^+\varphi _2^{}}}>=<{\displaystyle \frac{^2\overline{m}_b^2}{\varphi _2^+\varphi _1^{}}}>={\displaystyle \frac{|h_th_b|m_tm_b}{m_b^2m_t^2}}.`$
Note that we have not listed derivatives that vanish.
### A.2 Derivatives of Squark Masses with Respect to <br>Neutral Higgs Fields
This section contains the derivatives of the field-dependent squark masses $`\stackrel{~}{m}_{t_{1,2}}^2`$ and $`\stackrel{~}{m}_{b_{1,2}}^2`$ with respect to neutral Higgs fields $`\varphi _{1,2}`$ and $`a_{1,2}`$. We first give the tadpole terms:
$`{\displaystyle \frac{1}{v_1}}<{\displaystyle \frac{\stackrel{~}{m}_{t_1(t_2)}^2}{\varphi _1}}>`$ $`=`$ $`{\displaystyle \frac{g_w^2+g^2}{8}}+(){\displaystyle \frac{1}{2\left(m_{\stackrel{~}{t}_1}^2m_{\stackrel{~}{t}_2}^2\right)}}[x_t(\stackrel{~}{M}_Q^2\stackrel{~}{M}_t^2+\frac{1}{2}x_tv^2\mathrm{cos}2\beta )`$
$`\mathrm{\hspace{0.17em}2}|h_t|^2(\mathrm{Re}(\mu A_t)\mathrm{tan}\beta |\mu |^2)],`$
$`{\displaystyle \frac{1}{v_1}}<{\displaystyle \frac{\stackrel{~}{m}_{b_1(b_2)}^2}{\varphi _1}}>`$ $`=`$ $`|h_b|^2{\displaystyle \frac{g_w^2+g^2}{8}}(+){\displaystyle \frac{1}{2\left(m_{\stackrel{~}{b}_1}^2m_{\stackrel{~}{b}_2}^2\right)}}[x_b(\stackrel{~}{M}_Q^2\stackrel{~}{M}_b^2\frac{1}{2}x_bv^2\mathrm{cos}2\beta )`$
$`+\mathrm{\hspace{0.17em}2}|h_b|^2(\mathrm{Re}(\mu A_b)\mathrm{tan}\beta |A_b|^2)],`$
$`{\displaystyle \frac{1}{v_2}}<{\displaystyle \frac{\stackrel{~}{m}_{t_1(t_2)}^2}{\varphi _2}}>`$ $`=`$ $`|h_t|^2{\displaystyle \frac{g_w^2+g^2}{8}}(+){\displaystyle \frac{1}{2\left(m_{\stackrel{~}{t}_1}^2m_{\stackrel{~}{t}_2}^2\right)}}[x_t(\stackrel{~}{M}_Q^2\stackrel{~}{M}_t^2+\frac{1}{2}x_tv^2\mathrm{cos}2\beta )`$
$`+\mathrm{\hspace{0.17em}2}|h_t|^2(\mathrm{Re}(\mu A_t)\mathrm{cot}\beta |A_t|^2)],`$
$`{\displaystyle \frac{1}{v_2}}<{\displaystyle \frac{\stackrel{~}{m}_{b_1(b_2)}^2}{\varphi _2}}>`$ $`=`$ $`{\displaystyle \frac{g_w^2+g^2}{8}}+(){\displaystyle \frac{1}{2\left(m_{\stackrel{~}{b}_1}^2m_{\stackrel{~}{b}_2}^2\right)}}[x_b(\stackrel{~}{M}_Q^2\stackrel{~}{M}_b^2\frac{1}{2}x_bv^2\mathrm{cos}2\beta )`$
$`\mathrm{\hspace{0.17em}2}|h_b|^2(\mathrm{Re}(\mu A_b)\mathrm{cot}\beta |\mu |^2)],`$
$`{\displaystyle \frac{1}{v_2}}<{\displaystyle \frac{\stackrel{~}{m}_{t_1(t_2)}^2}{a_1}}>`$ $`=`$ $`{\displaystyle \frac{1}{v_1}}<{\displaystyle \frac{\stackrel{~}{m}_{t_1(t_2)}^2}{a_2}}>=(+){\displaystyle \frac{|h_t|^2\mathrm{Im}(\mu A_t)}{m_{\stackrel{~}{t}_1}^2m_{\stackrel{~}{t}_2}^2}},`$
$`{\displaystyle \frac{1}{v_2}}<{\displaystyle \frac{\stackrel{~}{m}_{b_1(b_2)}^2}{a_1}}>`$ $`=`$ $`{\displaystyle \frac{1}{v_1}}<{\displaystyle \frac{\stackrel{~}{m}_{b_1(b_2)}^2}{a_2}}>=(+){\displaystyle \frac{|h_b|^2\mathrm{Im}(\mu A_b)}{m_{\stackrel{~}{b}_1}^2m_{\stackrel{~}{b}_2}^2}},`$ (A.3)
where the coupling parameters $`x_t`$ and $`x_b`$ are defined after (2). Then, the self-energy-type terms $`<^2\stackrel{~}{m}_{q_k}^2/a_ia_j>`$, with $`q_k=t_1,b_1,t_2,b_2`$ and $`i,j=1,2`$, are found to be
$`<{\displaystyle \frac{^2\stackrel{~}{m}_{t_1(t_2)}^2}{a_1^2}}>`$ $`=`$ $`{\displaystyle \frac{g_w^2+g^2}{8}}+(){\displaystyle \frac{1}{2\left(m_{\stackrel{~}{t}_1}^2m_{\stackrel{~}{t}_2}^2\right)}}\left[x_t\left(\stackrel{~}{M}_Q^2\stackrel{~}{M}_t^2+\frac{1}{2}x_tv^2\mathrm{cos}2\beta \right)+\mathrm{\hspace{0.17em}2}|h_t|^2|\mu |^2\right],`$
$`<{\displaystyle \frac{^2\stackrel{~}{m}_{b_1(b_2)}^2}{a_1^2}}>`$ $`=`$ $`|h_b|^2{\displaystyle \frac{g_w^2+g^2}{8}}(+){\displaystyle \frac{1}{2\left(m_{\stackrel{~}{b}_1}^2m_{\stackrel{~}{b}_2}^2\right)}}[x_b(\stackrel{~}{M}_Q^2\stackrel{~}{M}_b^2\frac{1}{2}x_bv^2\mathrm{cos}2\beta )`$
$`\mathrm{\hspace{0.17em}2}|h_b|^2|A_b|^2],`$
$`<{\displaystyle \frac{^2\stackrel{~}{m}_{t_1(t_2)}^2}{a_2^2}}>`$ $`=`$ $`|h_t|^2{\displaystyle \frac{g_w^2+g^2}{8}}(+){\displaystyle \frac{1}{2\left(m_{\stackrel{~}{t}_1}^2m_{\stackrel{~}{t}_2}^2\right)}}[x_t(\stackrel{~}{M}_Q^2\stackrel{~}{M}_t^2+\frac{1}{2}x_tv^2\mathrm{cos}2\beta )`$
$`\mathrm{\hspace{0.17em}2}|h_t|^2|A_t|^2],`$
$`<{\displaystyle \frac{^2\stackrel{~}{m}_{b_1(b_2)}^2}{a_2^2}}>`$ $`=`$ $`{\displaystyle \frac{g_w^2+g^2}{8}}+(){\displaystyle \frac{1}{2\left(m_{\stackrel{~}{b}_1}^2m_{\stackrel{~}{b}_2}^2\right)}}\left[x_b\left(\stackrel{~}{M}_Q^2\stackrel{~}{M}_b^2\frac{1}{2}x_bv^2\mathrm{cos}2\beta \right)+\mathrm{\hspace{0.17em}2}|h_b|^2|\mu |^2\right],`$
$`<{\displaystyle \frac{^2\stackrel{~}{m}_{t_1(t_2)}^2}{a_1a_2}}>`$ $`=`$ $`(+){\displaystyle \frac{|h_t|^2\mathrm{Re}(\mu A_t)}{m_{\stackrel{~}{t}_1}^2m_{\stackrel{~}{t}_2}^2}}+(){\displaystyle \frac{2v_1v_2|h_t|^4\mathrm{Im}^2(\mu A_t)}{(m_{\stackrel{~}{t}_1}^2m_{\stackrel{~}{t}_2}^2)^3}},`$
$`<{\displaystyle \frac{^2\stackrel{~}{m}_{b_1(b_2)}^2}{a_1a_2}}>`$ $`=`$ $`(+){\displaystyle \frac{|h_b|^2\mathrm{Re}(\mu A_b)}{m_{\stackrel{~}{b}_1}^2m_{\stackrel{~}{b}_2}^2}}+(){\displaystyle \frac{2v_1v_2|h_b|^4\mathrm{Im}^2(\mu A_b)}{(m_{\stackrel{~}{b}_1}^2m_{\stackrel{~}{b}_2}^2)^3}}.`$ (A.4)
In addition, the non-vanishing CP-violating self-energy terms $`<^2\stackrel{~}{m}_{q_k}^2/\varphi _ia_j>`$ are given by
$`<{\displaystyle \frac{^2\stackrel{~}{m}_{t_1(t_2)}^2}{\varphi _1a_2}}>`$ $`=`$ $`+(){\displaystyle \frac{|h_t|^2\mathrm{Im}(\mu A_t)}{m_{\stackrel{~}{t}_1}^2m_{\stackrel{~}{t}_2}^2}}\{\mathrm{\hspace{0.17em}1}{\displaystyle \frac{v_1^2}{\left(m_{\stackrel{~}{t}_1}^2m_{\stackrel{~}{t}_2}^2\right)^2}}[x_t(\stackrel{~}{M}_Q^2\stackrel{~}{M}_t^2+\frac{1}{2}x_tv^2\mathrm{cos}2\beta )`$
$`\mathrm{\hspace{0.17em}2}|h_t|^2(\mathrm{Re}(\mu A_t)\mathrm{tan}\beta |\mu |^2)]\},`$
$`<{\displaystyle \frac{^2\stackrel{~}{m}_{b_1(b_2)}^2}{\varphi _1a_2}}>`$ $`=`$ $`+(){\displaystyle \frac{|h_b|^2\mathrm{Im}(\mu A_b)}{m_{\stackrel{~}{b}_1}^2m_{\stackrel{~}{b}_2}^2}}\{\mathrm{\hspace{0.17em}1}+{\displaystyle \frac{v_1^2}{\left(m_{\stackrel{~}{b}_1}^2m_{\stackrel{~}{b}_2}^2\right)^2}}[x_b(\stackrel{~}{M}_Q^2\stackrel{~}{M}_b^2\frac{1}{2}x_bv^2\mathrm{cos}2\beta )`$
$`+\mathrm{\hspace{0.17em}2}|h_b|^2(\mathrm{Re}(\mu A_b)\mathrm{tan}\beta |A_b|^2)]\},`$
$`<{\displaystyle \frac{^2\stackrel{~}{m}_{t_1(t_2)}^2}{\varphi _2a_1}}>`$ $`=`$ $`(+){\displaystyle \frac{|h_t|^2\mathrm{Im}(\mu A_t)}{m_{\stackrel{~}{t}_1}^2m_{\stackrel{~}{t}_2}^2}}\{\mathrm{\hspace{0.17em}1}+{\displaystyle \frac{v_2^2}{\left(m_{\stackrel{~}{t}_1}^2m_{\stackrel{~}{t}_2}^2\right)^2}}[x_t(\stackrel{~}{M}_Q^2\stackrel{~}{M}_t^2+\frac{1}{2}x_tv^2\mathrm{cos}2\beta )`$
$`+\mathrm{\hspace{0.17em}2}|h_t|^2(\mathrm{Re}(\mu A_t)\mathrm{cot}\beta |A_t|^2)]\},`$
$`<{\displaystyle \frac{^2\stackrel{~}{m}_{b_1(b_2)}^2}{\varphi _2a_1}}>`$ $`=`$ $`(+){\displaystyle \frac{|h_b|^2\mathrm{Im}(\mu A_b)}{m_{\stackrel{~}{b}_1}^2m_{\stackrel{~}{b}_2}^2}}\{\mathrm{\hspace{0.17em}1}{\displaystyle \frac{v_2^2}{\left(m_{\stackrel{~}{b}_1}^2m_{\stackrel{~}{b}_2}^2\right)^2}}[x_b(\stackrel{~}{M}_Q^2\stackrel{~}{M}_b^2\frac{1}{2}x_bv^2\mathrm{cos}2\beta )`$ (A.5)
$`\mathrm{\hspace{0.17em}2}|h_b|^2(\mathrm{Re}(\mu A_b)\mathrm{cot}\beta |\mu |^2)]\}.`$
Finally, the CP-conserving self-energy-type derivatives $`<^2\stackrel{~}{m}_{q_k}^2/\varphi _i\varphi _j>`$ have been calculated to be
$`<{\displaystyle \frac{^2\stackrel{~}{m}_{t_1(t_2)}^2}{\varphi _1^2}}>`$ $`=`$ $`{\displaystyle \frac{g_w^2+g^2}{8}}+(){\displaystyle \frac{1}{2\left(m_{\stackrel{~}{t}_1}^2m_{\stackrel{~}{t}_2}^2\right)}}[x_t(\stackrel{~}{M}_Q^2\stackrel{~}{M}_t^2+\frac{1}{2}x_tv^2(1+2\mathrm{cos}2\beta ))`$
$`+\mathrm{\hspace{0.17em}2}|h_t|^2|\mu |^2](+){\displaystyle \frac{v_1^2}{2\left(m_{\stackrel{~}{t}_1}^2m_{\stackrel{~}{t}_2}^2\right)^3}}[x_t(\stackrel{~}{M}_Q^2\stackrel{~}{M}_t^2+\frac{1}{2}x_tv^2\mathrm{cos}2\beta )`$
$`\mathrm{\hspace{0.17em}2}|h_t|^2(\mathrm{Re}(\mu A_t)\mathrm{tan}\beta |\mu |^2)]^2,`$
$`<{\displaystyle \frac{^2\stackrel{~}{m}_{b_1(b_2)}^2}{\varphi _1^2}}>`$ $`=`$ $`|h_b|^2{\displaystyle \frac{g_w^2+g^2}{8}}(+){\displaystyle \frac{1}{2\left(m_{\stackrel{~}{b}_1}^2m_{\stackrel{~}{b}_2}^2\right)}}[x_b(\stackrel{~}{M}_Q^2\stackrel{~}{M}_b^2\frac{1}{2}x_bv^2(1+2\mathrm{cos}2\beta ))`$
$`\mathrm{\hspace{0.17em}2}|h_b|^2|A_b|^2](+){\displaystyle \frac{v_1^2}{2\left(m_{\stackrel{~}{b}_1}^2m_{\stackrel{~}{b}_2}^2\right)^3}}[x_b(\stackrel{~}{M}_Q^2\stackrel{~}{M}_b^2\frac{1}{2}x_bv^2\mathrm{cos}2\beta )`$
$`+\mathrm{\hspace{0.17em}2}|h_b|^2(\mathrm{Re}(\mu A_b)\mathrm{tan}\beta |A_b|^2)]^2,`$
$`<{\displaystyle \frac{^2\stackrel{~}{m}_{t_1(t_2)}^2}{\varphi _2^2}}>`$ $`=`$ $`|h_t|^2{\displaystyle \frac{g_w^2+g^2}{8}}(+){\displaystyle \frac{1}{2\left(m_{\stackrel{~}{t}_1}^2m_{\stackrel{~}{t}_2}^2\right)}}[x_t(\stackrel{~}{M}_Q^2\stackrel{~}{M}_t^2+\frac{1}{2}x_tv^2(2\mathrm{cos}2\beta 1))`$
$`\mathrm{\hspace{0.17em}2}|h_t|^2|A_t|^2](+){\displaystyle \frac{v_2^2}{2\left(m_{\stackrel{~}{t}_1}^2m_{\stackrel{~}{t}_2}^2\right)^3}}[x_t(\stackrel{~}{M}_Q^2\stackrel{~}{M}_t^2+\frac{1}{2}x_tv^2\mathrm{cos}2\beta )`$
$`+\mathrm{\hspace{0.17em}2}|h_t|^2(\mathrm{Re}(\mu A_t)\mathrm{cot}\beta |A_t|^2)]^2,`$
$`<{\displaystyle \frac{^2\stackrel{~}{m}_{b_1(b_2)}^2}{\varphi _2^2}}>`$ $`=`$ $`{\displaystyle \frac{g_w^2+g^2}{8}}+(){\displaystyle \frac{1}{2\left(m_{\stackrel{~}{b}_1}^2m_{\stackrel{~}{b}_2}^2\right)}}[x_b(\stackrel{~}{M}_Q^2\stackrel{~}{M}_b^2\frac{1}{2}x_bv^2(2\mathrm{cos}2\beta 1))`$
$`+\mathrm{\hspace{0.17em}2}|h_b|^2|\mu |^2](+){\displaystyle \frac{v_2^2}{2\left(m_{\stackrel{~}{b}_1}^2m_{\stackrel{~}{b}_2}^2\right)^3}}[x_b(\stackrel{~}{M}_Q^2\stackrel{~}{M}_b^2\frac{1}{2}x_bv^2\mathrm{cos}2\beta )`$
$`\mathrm{\hspace{0.17em}2}|h_b|^2(\mathrm{Re}(\mu A_b)\mathrm{cot}\beta |\mu |^2)]^2,`$
$`<{\displaystyle \frac{^2\stackrel{~}{m}_{t_1(t_2)}^2}{\varphi _1\varphi _2}}>`$ $`=`$ $`(+){\displaystyle \frac{1}{2\left(m_{\stackrel{~}{t}_1}^2m_{\stackrel{~}{t}_2}^2\right)}}\left[\frac{1}{2}x_t^2v^2\mathrm{sin}2\beta +\mathrm{\hspace{0.17em}2}|h_t|^2\mathrm{Re}(\mu A_t)\right]+(){\displaystyle \frac{v_1v_2}{2\left(m_{\stackrel{~}{t}_1}^2m_{\stackrel{~}{t}_2}^2\right)^3}}`$
$`\times \left[x_t\left(\stackrel{~}{M}_Q^2\stackrel{~}{M}_t^2+\frac{1}{2}x_tv^2\mathrm{cos}2\beta \right)\mathrm{\hspace{0.17em}2}|h_t|^2\left(\mathrm{Re}(\mu A_t)\mathrm{tan}\beta |\mu |^2\right)\right]`$
$`\times \left[x_t\left(\stackrel{~}{M}_Q^2\stackrel{~}{M}_t^2+\frac{1}{2}x_tv^2\mathrm{cos}2\beta \right)+\mathrm{\hspace{0.17em}2}|h_t|^2\left(\mathrm{Re}(\mu A_t)\mathrm{cot}\beta |A_t|^2\right)\right],`$
$`<{\displaystyle \frac{^2\stackrel{~}{m}_{b_1(b_2)}^2}{\varphi _1\varphi _2}}>`$ $`=`$ $`(+){\displaystyle \frac{1}{2\left(m_{\stackrel{~}{b}_1}^2m_{\stackrel{~}{b}_2}^2\right)}}\left[\frac{1}{2}x_b^2v^2\mathrm{sin}2\beta +\mathrm{\hspace{0.17em}2}|h_b|^2\mathrm{Re}(\mu A_b)\right]+(){\displaystyle \frac{v_1v_2}{2\left(m_{\stackrel{~}{b}_1}^2m_{\stackrel{~}{b}_2}^2\right)^3}}`$ (A.6)
$`\times \left[x_b\left(\stackrel{~}{M}_Q^2\stackrel{~}{M}_b^2\frac{1}{2}x_bv^2\mathrm{cos}2\beta \right)+\mathrm{\hspace{0.17em}2}|h_b|^2\left(\mathrm{Re}(\mu A_b)\mathrm{tan}\beta |A_b|^2\right)\right]`$
$`\times \left[x_b\left(\stackrel{~}{M}_Q^2\stackrel{~}{M}_b^2\frac{1}{2}x_bv^2\mathrm{cos}2\beta \right)\mathrm{\hspace{0.17em}2}|h_b|^2\left(\mathrm{Re}(\mu A_b)\mathrm{cot}\beta |\mu |^2\right)\right].`$
### A.3 Derivatives of Squark Masses with Respect to <br>Charged Higgs Fields
Here we evaluate the derivatives of the field-dependent squark masses with respect to charged Higgs fields. To calculate the expressions $`<^2\stackrel{~}{m}_{q_k}^2/\varphi _i^+\varphi _j^{}>`$ directly turns out to be a formidable task. The reason is that $`\stackrel{~}{m}_{q_k}^2`$ are the eigenvalues of a non-trivial $`(4\times 4)`$ squark mass matrix $`\stackrel{~}{}^2`$ (cf. (2.11)), and their analytic form is very complicated. Therefore, we proceed differently, using a mathematical trick which was first applied in .
First, we notice that $`<\stackrel{~}{m}_{q_k}^2/\varphi _i^\pm >=0`$, as a consequence of the fact that the true ground state of the effective potential should conserve charge. Then, one may make use of the eigenvalue equation:
$$\mathrm{det}(\stackrel{~}{}^2\stackrel{~}{m}_{q_k}^2\mathrm{𝟏}_4)=\stackrel{~}{m}_{q_k}^8+A\stackrel{~}{m}_{q_k}^6+B\stackrel{~}{m}_{q_k}^4+C\stackrel{~}{m}_{q_k}^2+D=0,$$
(A.7)
with
$`A`$ $`=`$ $`\mathrm{Tr}\stackrel{~}{}^2,`$
$`B`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(\mathrm{Tr}^2\stackrel{~}{}^2\mathrm{Tr}\stackrel{~}{}^4\right),`$
$`C`$ $`=`$ $`{\displaystyle \frac{1}{3}}\left(\mathrm{Tr}^3\stackrel{~}{}^2\mathrm{Tr}\stackrel{~}{}^6\right){\displaystyle \frac{1}{2}}\mathrm{Tr}\stackrel{~}{}^2\left(\mathrm{Tr}^2\stackrel{~}{}^2\mathrm{Tr}\stackrel{~}{}^4\right),`$
$`D`$ $`=`$ $`\mathrm{det}\stackrel{~}{}^2={\displaystyle \frac{1}{4}}\left(\mathrm{Tr}\stackrel{~}{}^8+A\mathrm{Tr}\stackrel{~}{}^6+B\mathrm{Tr}\stackrel{~}{}^4+C\mathrm{Tr}\stackrel{~}{}^2\right),`$ (A.8)
to obtain
$$<\frac{^2\stackrel{~}{m}_{q_k}^2}{\varphi _i^+\varphi _j^{}}>=<\frac{A_{ij}\stackrel{~}{m}_{q_k}^6+B_{ij}\stackrel{~}{m}_{q_k}^4+C_{ij}\stackrel{~}{m}_{q_k}^2+D_{ij}}{\underset{q_lq_k}{}(\stackrel{~}{m}_{q_k}^2\stackrel{~}{m}_{q_l}^2)}>,$$
(A.9)
with $`q_l,q_k=t_1,b_1,t_2,b_2`$, and $`A_{ij}=^2A/\varphi _i^+\varphi _j^{}`$, $`B_{ij}=^2B/\varphi _i^+\varphi _j^{}`$, etc.. For our purposes, it is sufficient to calculate the derivatives with respect to $`\varphi _1^+`$ and $`\varphi _2^{}`$. In particular, it proves convenient to use a representation in which the columns and rows ‘2’ and ‘3’ of the $`(4\times 4)`$ matrix $`\stackrel{~}{}^2`$ have been interchanged. With such a reordering, we find
$$<\stackrel{~}{}^2>\left(\begin{array}{cc}\stackrel{~}{}_t^2& 0\\ 0& \stackrel{~}{}_b^2\end{array}\right),<\frac{\stackrel{~}{}^2}{\varphi _1^+}>\left(\begin{array}{cc}0& \stackrel{~}{}_+\\ 0& 0\end{array}\right),<\frac{\stackrel{~}{}^2}{\varphi _2^{}}>\left(\begin{array}{cc}0& 0\\ \stackrel{~}{}_{}& 0\end{array}\right),$$
(A.10)
where $`\stackrel{~}{}_t^2`$ and $`\stackrel{~}{}_b^2`$ are the usual $`\stackrel{~}{t}`$\- and $`\stackrel{~}{b}`$\- $`(2\times 2)`$-mass matrices, respectively, and
$$\stackrel{~}{}_+=\left(\begin{array}{cc}\frac{1}{\sqrt{2}}(|h_b|^2\frac{1}{2}g_w^2)v_1& h_b^{}A_b^{}\\ h_t\mu ^{}& \frac{1}{\sqrt{2}}h_th_b^{}v_2\end{array}\right),\stackrel{~}{}_{}=\left(\begin{array}{cc}\frac{1}{\sqrt{2}}(|h_t|^2\frac{1}{2}g_w^2)v_2& h_t^{}A_t^{}\\ h_b\mu ^{}& \frac{1}{\sqrt{2}}h_t^{}h_bv_1\end{array}\right).$$
(A.11)
Then, the relevant coefficients $`A_{12}`$, $`B_{12}`$, $`C_{12}`$ and $`D_{12}`$ may be expressed in a compact form as follows:
$`<A_{12}>`$ $`=`$ $`0,`$
$`<B_{12}>`$ $`=`$ $`\mathrm{Tr}(\stackrel{~}{}_+\stackrel{~}{}_{}),`$
$`<C_{12}>`$ $`=`$ $`\left(\mathrm{Tr}\stackrel{~}{}_t^2+\mathrm{Tr}\stackrel{~}{}_b^2\right)\mathrm{Tr}(\stackrel{~}{}_+\stackrel{~}{}_{})\mathrm{Tr}(\stackrel{~}{}_t^2\stackrel{~}{}_+\stackrel{~}{}_{})\mathrm{Tr}(\stackrel{~}{}_b^2\stackrel{~}{}_{}\stackrel{~}{}_+),`$
$`<D_{12}>`$ $`=`$ $`\mathrm{Tr}(\stackrel{~}{}_t^2\stackrel{~}{}_+\stackrel{~}{}_b^2\stackrel{~}{}_{})\mathrm{Tr}(\stackrel{~}{}_t^4\stackrel{~}{}_+\stackrel{~}{}_{})\mathrm{Tr}(\stackrel{~}{}_b^4\stackrel{~}{}_{}\stackrel{~}{}_+)`$ (A.12)
$`+\left(\mathrm{Tr}\stackrel{~}{}_t^2+\mathrm{Tr}\stackrel{~}{}_b^2\right)\left[\mathrm{Tr}(\stackrel{~}{}_t^2\stackrel{~}{}_+\stackrel{~}{}_{})+\mathrm{Tr}(\stackrel{~}{}_b^2\stackrel{~}{}_{}\stackrel{~}{}_+)\right]`$
$`+{\displaystyle \frac{1}{2}}\left[\mathrm{Tr}\stackrel{~}{}_t^4+\mathrm{Tr}\stackrel{~}{}_b^4\left(\mathrm{Tr}\stackrel{~}{}_t^2+\mathrm{Tr}\stackrel{~}{}_b^2\right)^2\right]\mathrm{Tr}(\stackrel{~}{}_+\stackrel{~}{}_{}).`$
With the help of (A.3), it is straightforward to obtain the derivatives $`^2\stackrel{~}{m}_{q_k}^2/\varphi _1^+\varphi _2^{}`$. More explicitly, we have
$`<{\displaystyle \frac{^2\stackrel{~}{m}_{t_1}^2}{\varphi _1^+\varphi _2^{}}}>`$ $`=`$ $`{\displaystyle \frac{<B_{12}>m_{\stackrel{~}{t}_1}^4+<C_{12}>m_{\stackrel{~}{t}_1}^2+<D_{12}>}{\left(m_{\stackrel{~}{t}_1}^2m_{\stackrel{~}{b}_1}^2\right)\left(m_{\stackrel{~}{t}_1}^2m_{\stackrel{~}{t}_2}^2\right)\left(m_{\stackrel{~}{t}_1}^2m_{\stackrel{~}{b}_2}^2\right)}},`$
$`<{\displaystyle \frac{^2\stackrel{~}{m}_{t_2}^2}{\varphi _1^+\varphi _2^{}}}>`$ $`=`$ $`{\displaystyle \frac{<B_{12}>m_{\stackrel{~}{t}_2}^4+<C_{12}>m_{\stackrel{~}{t}_2}^2+<D_{12}>}{\left(m_{\stackrel{~}{t}_2}^2m_{\stackrel{~}{b}_1}^2\right)\left(m_{\stackrel{~}{t}_2}^2m_{\stackrel{~}{t}_1}^2\right)\left(m_{\stackrel{~}{t}_2}^2m_{\stackrel{~}{b}_2}^2\right)}},`$
$`<{\displaystyle \frac{^2\stackrel{~}{m}_{b_1}^2}{\varphi _1^+\varphi _2^{}}}>`$ $`=`$ $`{\displaystyle \frac{<B_{12}>m_{\stackrel{~}{b}_1}^4+<C_{12}>m_{\stackrel{~}{b}_1}^2+<D_{12}>}{\left(m_{\stackrel{~}{b}_1}^2m_{\stackrel{~}{t}_1}^2\right)\left(m_{\stackrel{~}{b}_1}^2m_{\stackrel{~}{t}_2}^2\right)\left(m_{\stackrel{~}{b}_1}^2m_{\stackrel{~}{b}_2}^2\right)}},`$
$`<{\displaystyle \frac{^2\stackrel{~}{m}_{b_2}^2}{\varphi _1^+\varphi _2^{}}}>`$ $`=`$ $`{\displaystyle \frac{<B_{12}>m_{\stackrel{~}{b}_2}^4+<C_{12}>m_{\stackrel{~}{b}_2}^2+<D_{12}>}{\left(m_{\stackrel{~}{b}_2}^2m_{\stackrel{~}{b}_1}^2\right)\left(m_{\stackrel{~}{b}_2}^2m_{\stackrel{~}{t}_2}^2\right)\left(m_{\stackrel{~}{b}_2}^2m_{\stackrel{~}{t}_1}^2\right)}}.`$ (A.13)
## Appendix B Higgs-Boson Masses and Mixing Angles
Here we present analytic expressions for the Higgs-boson masses $`M_{H_i}(m_t)`$ ($`i=1,2,3`$) and the corresponding $`(3\times 3)`$ orthogonal matrix $`O`$, after diagonalizing the RG-improved Higgs-boson mass matrix $`(_N^2)(m_t)`$.
For notational simplicity, we do not display explicitly the functional dependence of $`_N^2`$ on $`m_t`$. The mass eigenvalues of the $`(3\times 3)`$ matrix $`_N^2`$ are then obtained by solving the characteristic equation of cubic order:
$$x^3+rx^2+sx+t=0,$$
(B.1)
with
$`r`$ $`=`$ $`\mathrm{Tr}(_N^2),`$
$`s`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[\mathrm{Tr}^2(_N^2)\mathrm{Tr}(_N^4)\right],`$
$`t`$ $`=`$ $`\mathrm{det}(_N^2).`$ (B.2)
To this end, it proves useful to define the following auxiliary parameters:
$`p`$ $`=`$ $`{\displaystyle \frac{3sr^2}{3}},`$
$`q`$ $`=`$ $`{\displaystyle \frac{2r^3}{27}}{\displaystyle \frac{rs}{3}}+t,`$
$`D`$ $`=`$ $`{\displaystyle \frac{p^3}{27}}+{\displaystyle \frac{q^2}{4}}.`$ (B.3)
To ensure that the three eigenvalues are positive, it is necessary and sufficient to require that
$$D<0,r<0,s>0,t<0.$$
(B.4)
Imposing these inequalities on the kinematic parameters of the theory, we may express the three mass eigenvalues of $`_N^2`$ as
$`M_{H_1}^2`$ $`=`$ $`{\displaystyle \frac{1}{3}}r+2\sqrt{p/3}\mathrm{cos}\left({\displaystyle \frac{\phi }{3}}+{\displaystyle \frac{2\pi }{3}}\right),`$
$`M_{H_2}^2`$ $`=`$ $`{\displaystyle \frac{1}{3}}r+2\sqrt{p/3}\mathrm{cos}\left({\displaystyle \frac{\phi }{3}}{\displaystyle \frac{2\pi }{3}}\right),`$
$`M_{H_3}^2`$ $`=`$ $`{\displaystyle \frac{1}{3}}r+2\sqrt{p/3}\mathrm{cos}\left({\displaystyle \frac{\phi }{3}}\right),`$ (B.5)
with
$$\phi =\mathrm{arccos}\left(\frac{q}{2\sqrt{p^3/27}}\right)\mathrm{and}0\phi \pi .$$
(B.6)
Since the Higgs-boson mass matrix $`_N^2`$ is symmetric, we can diagonalize it by means of an orthogonal rotation $`O`$ as stated in (3.39). Furthermore, one can show that the Higgs-boson mass eigenvalues in (B) satisfy the desired mass hierarchy in accordance with the inequality of (3.40).
If $`M_{ij}^2`$, with $`i,j=1,2,3`$, denote the matrix elements of $`_N^2`$, the elements $`O_{ij}`$ can then be obtained by appropriately solving the underdetermined coupled system of equations, $`_kM_{ik}^2O_{kj}=M_{H_j}^2O_{ij}`$:
$`(M_{11}^2M_{H_i}^2)O_{1i}+M_{12}^2O_{2i}+M_{13}^2O_{3i}`$ $`=`$ $`0,`$
$`M_{21}^2O_{1i}+(M_{22}^2M_{H_i}^2)O_{2i}+M_{23}^2O_{3i}`$ $`=`$ $`0,`$
$`M_{31}^2O_{1i}+M_{32}^2O_{2i}+(M_{33}^2M_{H_i}^2)O_{3i}`$ $`=`$ $`0.`$ (B.7)
More explicitly, we have
$$O=\left(\begin{array}{ccc}|x_1|/\mathrm{\Delta }_1& x_2/\mathrm{\Delta }_2& x_3/\mathrm{\Delta }_3\\ y_1/\mathrm{\Delta }_1& |y_2|/\mathrm{\Delta }_2& y_3/\mathrm{\Delta }_3\\ z_1/\mathrm{\Delta }_1& z_2/\mathrm{\Delta }_2& |z_3|/\mathrm{\Delta }_3\end{array}\right),$$
(B.8)
where
$$\mathrm{\Delta }_i=\sqrt{x_i^2+y_i^2+z_i^2}$$
(B.9)
and
$`\left|x_1\right|`$ $`=`$ $`\begin{array}{cc}M_{22}^2M_{H_1}^2\hfill & \hfill M_{23}^2\\ M_{32}^2\hfill & \hfill M_{33}^2M_{H_1}^2\end{array},y_1=\mathrm{s}_{x_1}\left|\begin{array}{cc}M_{23}^2& M_{21}^2\\ M_{33}^2M_{H_1}^2& M_{31}^2\end{array}\right|,z_1=\mathrm{s}_{x_1}\left|\begin{array}{cc}M_{21}^2& M_{22}^2M_{H_1}^2\\ M_{31}^2& M_{32}^2\end{array}\right|,`$ (B.16)
$`x_2`$ $`=`$ $`\mathrm{s}_{y_2}\left|\begin{array}{cc}M_{13}^2& M_{12}^2\\ M_{33}^2M_{H_2}^2& M_{32}^2\end{array}\right|,\left|y_2\right|=\begin{array}{cc}M_{11}^2M_{H_2}^2\hfill & \hfill M_{13}^2\\ M_{31}^2\hfill & \hfill M_{33}^2M_{H_2}^2\end{array},z_2=\mathrm{s}_{y_2}\left|\begin{array}{cc}M_{12}^2& M_{11}^2M_{H_2}^2\\ M_{32}^2& M_{31}^2\end{array}\right|,`$ (B.23)
$`x_3`$ $`=`$ $`\mathrm{s}_{z_3}\left|\begin{array}{cc}M_{12}^2& M_{13}^2\\ M_{22}^2M_{H_3}^2& M_{23}^2\end{array}\right|,y_3=\mathrm{s}_{z_3}\left|\begin{array}{cc}M_{13}^2& M_{11}^2M_{H_3}^2\\ M_{23}^2& M_{21}^2\end{array}\right|,\left|z_3\right|=\begin{array}{cc}M_{11}^2M_{H_3}^2\hfill & \hfill M_{12}^2\\ M_{21}^2\hfill & \hfill M_{22}^2M_{H_3}^2\end{array}.`$ (B.30)
In (B.16), the abbreviation $`\mathrm{s}_x\mathrm{sign}(x)`$ is an operation that simply gives the sign of a real expression $`x`$.
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# Color Dynamics On Phase Space
## Abstract
We describe the properties of quark matter at zero temperature and finite baryon densities within microscopic Vlasov/molecular dynamics approaches. We use an inter-quark Richardson’s potential consistent with the indications of Lattice QCD calculations. The color degrees of freedom are explicitly taken into account. We explicitly demonstrate that the Vlasov approach alone is insufficient in the hadronization region. In order to overcome this problem we prepare the initial condition for many events using molecular dynamics with frictional cooling and a Thomas-Fermi approximation to the Fermi motion. These events are averaged and propagated in time using the Vlasov approach. We find some evidence for a second order phase transition from nuclear to quark matter at high baryon densities. An order parameter suitable to describe the phase transition is discussed. At low densities the quark condensate into approximately color white clusters (nucleon).
PACS : 12.39.Pn 24.85.1p
One of the open problems in theoretical nuclear and particle physics is how to obtain the well known nuclear properties starting from the quark degrees of freedom . This also includes the possibility of understanding the basic free nucleon-nucleon interaction from quark and gluons dynamics. Some kind of solution to this problem is becoming more and more needed with the new experiments done or planned using ultra-relativistic heavy ions at CERN and soon at RHIC. The search for a quark-gluon plasma (QGP) in such collisions is in fact one of the new and most exciting directions in physics at the border between nuclear and particle physics . Quantum ChromoDynamics (QCD) because of its difficulties (numerical and conceptual), has been applied so far to some limited cases such as quark matter at zero baryon ($`\rho _B`$) density and high temperature (T) . Furthermore in relativistic heavy ion collisions (RHIC) dynamics plays surely an important role and accordingly the theory should be dynamical.
Recently , we have proposed a dynamical approach based on the Vlasov equation to reproduce hadron masses and the properties of nuclear matter at finite $`\rho _B`$. Some works in the same spirit are discussed in . Our approach needs as inputs the interaction potential among quarks, which was borrowed from phenomenology i.e. the Richardson’s potential, and the quark masses which were fitted to reproduce known meson masses such as the $`\pi `$, the $`\varphi `$, the $`\eta _c`$ etc. When the particles are embedded in a dense medium such as in nuclear matter (NM) the potential becomes screened in a similar fashion as ions and electrons in condensed matter do, i.e. Debye screening (DS).
It is the purpose of this paper to refine that approach in one important aspect which is the treatment of the color degrees of freedom. In the previous works color degrees of freedom were implicitly taken into account through the use of a Debye radius that effectively screens the $`qq`$ interaction potential. In the present paper we give to the quarks explicitly a color (using the Gell-Mann matrices) and follow their dynamics in phase space solving the Vlasov equation (VE). Thus screening will be dynamically obtained. In general the self screening obtained in Vlasov dynamics is inadequate, which is the reason why it was not adopted in the earlier attempts . In fact we will show explicitly that the Vlasov approach alone gives a good description of the system at large densities only,i.e. in the QGP region. In order to overcome this problem we adopted the following strategy. We first prepare the initial conditions using molecular dynamics (MD) with frictional cooling for many events. The events are averaged and care is taken of antisymmetrization. These are the initial conditions for the Vlasov evolution. Since the VE fulfills the Liouville theorem, the initial phase space density remains constant in time. Thus the initial antisymmetrization and eventual clustering obtained in the cooling process are maintained.
We outline our approach on purely classical grounds, however the same results can be obtained within the Wigner transform formalism of the quantum BBGKY-hierarchy in the limit $`\mathrm{}0`$.
The exact (classical) one-body distribution function $`f_1(r,p,t)`$ satisfies the equation (BBGKY hierarchy):
$$_tf_1+\frac{\stackrel{}{\text{p}}}{E}_rf_1=d(2)_rV(𝐫,𝐫_\mathrm{𝟐})_pf_2(𝐫,𝐫_\mathrm{𝟐},𝐩,𝐩_\mathrm{𝟐},t)$$
(1)
$`E=\sqrt{p^2+m_i^2}`$ is the energy and $`m_i=10MeV`$ is the (u,d) quark mass. Here we assume the potential to be dependent on the relative coordinates only. A generalization to include a momentum dependent part is straightforward. $`f_2`$ is the two-body distribution function, which in the classical limit reads:
$$f_2(r,r_2,p,p_2,t)=\mathrm{\Sigma }_{\alpha \beta }^Q\delta (𝐫𝐫_\alpha )\delta (𝐩𝐩_\alpha )\times \delta (𝐫_\mathrm{𝟐}𝐫_\beta )\delta (𝐩_\mathrm{𝟐}𝐩_\beta )$$
(2)
where $`Q=q+\overline{q}`$ is the total number of quarks and anti quarks ($`\overline{q}=0`$ in this work). Inserting this equation into Eq.(1) gives:
$$_tf_1+\frac{\stackrel{}{\text{p}}}{E}_rf_1_\text{r}U_\text{p}f_1=0$$
(3)
Where $`U=\mathrm{\Sigma }_jV(𝐫,𝐫_𝐣)`$ is the exact potential. Let us now define $`f_1`$ and $`U`$ as sums of an ensemble averaged quantity plus the deviation from this average:
$$f_1=\overline{f}_1+\delta f_1;U=\overline{U}+\delta U$$
(4)
Substituting these equations in Eq.(3) and ensemble averaging gives:
$$_t\overline{f}_1+\frac{\stackrel{}{\text{p}}}{E}_r\overline{f}_1_\text{r}\overline{U}_\text{p}\overline{f}_1=<_r\delta U_p\delta f_1>$$
(5)
where one recognizes in the lhs the Vlasov term and in the rhs the Balescu-Lennard collision term . The mean-field is given by:
$$\overline{U}(𝐫)=\frac{\mathrm{𝟏}}{𝐍_{\mathrm{𝐞𝐯}}}𝚺_{\mathrm{𝐞𝐯}}𝚺_𝐣𝐕(𝐫,𝐫_𝐣)$$
(6)
For the purpose of this work we neglect the collision term in Eq.(5) and note that such term will be essential when dealing with RHIC. In agreement to LQCD calculations the interacting potential $`V(r)`$ for quarks is ($`\mathrm{}=1`$):
$`V(r_{i,j})=3\mathrm{\Sigma }_{a=1}^8{\displaystyle \frac{\lambda _i^a}{2}}{\displaystyle \frac{\lambda _j^a}{2}}[{\displaystyle \frac{8\pi }{332n_f}}\mathrm{\Lambda }(\mathrm{\Lambda }r_{ij}{\displaystyle \frac{f(\mathrm{\Lambda }r_{ij})}{\mathrm{\Lambda }r_{ij}}})+{\displaystyle \frac{8\pi }{9}}\overline{\alpha }{\displaystyle \frac{<\sigma _q\sigma _{\overline{q}}>}{m_qm_{\overline{q}}}}\delta (𝐫_{\mathrm{𝐢𝐣}})]`$ (7)
and
$`f(t)=14{\displaystyle \frac{dq}{q}\frac{e^{qt}}{[ln(q^21)]^2+\pi ^2}}`$ (8)
We fix the number of flavors $`n_f=2`$ and the parameter $`\mathrm{\Lambda }=0.25GeV`$ . In Eq.(7) we have added to the Richardson’s potential the chromomagnetic term (ct), very important to reproduce the masses of the hadrons in vacuum. Since in this work we will be dealing with finite nuclei, the ct can be neglected, we only notice that with the parameters choice discussed here, the hadron masses can be reproduced by suitably tuning the ct term .
The $`\lambda ^a`$ are the Gell-Mann matrices. From lattice calculations we expect that there is no color transport for distances of the order of $`0.20.3fm`$, which are distances much shorter than the ones we will be dealing with in this paper. Thus we will use the $`\lambda _{3,8}`$ only commuting diagonal Gell-Mann matrices (Abelian approximation).
Numerically the VE equation(5) is solved by writing the one body distribution function as:
$$\overline{f}_1(r,p,t)=\frac{1}{n_{tp}}\underset{i}{\overset{N}{}}\delta (rr_i(t))\delta (pp_i(t))$$
(9)
$`N=Qn_{tp}`$ is the number of such terms. Actually, N is much larger than the total quark number Q, so that we can say that each quark is represented by $`n_{tp}`$ terms called test particles(tp). Notice how well this fits in the previous discussion if we put $`n_{tp}=N_{ev}`$. Inserting Eq.(9) in the Vlasov equation gives the Hamilton equations of motion (eom) for the tp . The total number of tp (or corresponding $`N/Q`$ events)used in this work ranges from 5000 to 50000, and Q=150-300.
Initially, we distribute randomly the tp in a sphere of radius $`R=r_{0B}A^{1/3}`$ (the radius of the nucleus) in coordinate space and $`p_f`$ in momentum space. $`r_{0B}=(\frac{3}{4\pi \rho _B})^{1/3}`$, $`A=Q/3`$ and $`\rho _B`$ is the baryon density. $`p_f`$ is the Fermi momentum estimated in a simple Fermi gas model by imposing that a cell in phase space of size $`h=2\pi `$ can accommodate at most two identical quarks of different spins, flavors and colors. A simple estimate gives the following relation between the quark density , $`n_q`$, and the Fermi momentum:
$`n_q={\displaystyle \frac{g_q}{6\pi ^2}}p_f^3`$ (10)
The degeneracy number $`g_q=n_f\times n_c\times n_s`$, where $`n_c`$ is the number of colors and $`n_s`$ is the number of spins. For quarks and anti quarks 3 different colors are used red, yellow and blue (r,y,b) .
In figure(1), we plot the total energy per nucleon (top) and energy density (in units of the Fermi gas energy density) vs. baryon density. The full triangles give the results obtained by randomly distributing the tp as described above. We notice that a minimum at about $`\rho _c=2.08fm^3`$ is found with $`E_t/A=2GeV/A`$. Such a minimum is at much higher density and energy than expected for the ground state (gs) of a nucleus ($`\rho _0=0.16fm^3`$ and $`E_t/A=0.9380.016GeV/A`$). An important property of the system that we have described above is the following. If we rotate the quarks in color space, regardless of their position in r-space, the total energy will remain the same. This is indeed a ”pure” Vlasov solution and demonstrates explicitly the in-capability of the Vlasov approach to give clustering of quarks into nucleons. However this result is already instructive since it gives us an hint on where the quark and nuclear matter are located, i.e. above and below $`\rho _c`$ respectively. This result is qualitatively in agreement with the Hartree-Fock(HF) calculations of refs. (compare to fig.1 in ). For a discussion on why the $`E_t/A`$ increases at low densities in the HF/VE approaches we also refer to .
Of course distributing randomly the quarks in a sphere in r and p-space is not the most economical way to find the real gs of the system. In MD one searches for a minimum energy by introducing a friction term. The friction acts in such a way to lead the particles to a configuration for which the potential energy is a minimum. We cannot use this technique for our system since we are dealing with fermions and the friction term will destroy the initial antisymmetrization. In order to overcome this difficulty we adopt the following strategy. First we prepare $`N_{ev}`$ events as above, and we evolve them numerically solving the eom but with friction included. Because of the large number of particles interacting with attractive and repulsive forces, the quarks will slowly evolve to new positions where the potential is lower while keeping the initial root mean square radius approximately constant. When the averaged potential (over events) reaches a given value $`V_{min}`$, we look for the two closest particles $`(j,k)`$ to a quark $`i`$ in the same event. For these three quarks we know what the local density and the number of colors are. For instance if in a certain region we find two red and one blue quark, we use $`n_c=2`$ in Eq.(10) and calculate the local density from the knowledge of the distances of the 3 quarks. In this way the Fermi momentum is defined locally similar to the procedure used in nuclear or atomic physics (Thomas-Fermi approximation). We repeat this for all quarks $`i`$ in all events and calculate the total energy for this state . We let the system evolve again with friction included to a lower potential energy $`V_{min}=V_{min}\delta V`$, where $`\delta V`$ is a constant. We calculate the local density and local color numbers again and apply the Thomas-Fermi approximation to obtain the new total energy. We stop the procedure when the total energy is a minimum. The initial conditions so obtained are then propagated in time using the VE in order to maintain the initial antisymmetrization. We will show below, see fig.3, that the initial clustering obtained in MD is preserved as well during the Vlasov evolution.
The open triangles in fig.(1) are the result of the minimization. Notice especially at low densities the large decrease of the total energy of the system as compared to the Vlasov result. Now the calculated total energy at the nuclear gs is very close to the experimental value indicated by the full circle. However, we find slightly lower energies at lower densities, i.e. the gs of the nucleus is shifted in our calculations at about 1/10 of the experimental one. This should not be surprising in view of the simple potential that we have used. Also we have not tried to look for a best set of fitting parameters in these exploratory studies. At low $`\rho _B`$, the global invariance for rotations in color space is lost, i.e. if we exchange the colors of two $`\mathrm{𝑑𝑖𝑠𝑡𝑎𝑛𝑡}`$ quarks, the total energy of the system will change. At high densities (larger than $`2fm^3`$) the Vlasov and MD solutions are the same. This can be also seen in the bottom part of the figure where energy densities are given. We would like to stress again the qualitative agreement to ref. where a stochastic method had been used to calculate the g.s. energy of the system. This is quite evident if one compares our fig.1 to fig.1 of .
The energy density displayed in fig.1 (bottom) is a smooth function apart some fluctuations around $`2fm^3`$ density. From this result we can exclude a first order phase transition but a second order phase transition might be possible.
In order to check if a second order phase transition occurs we define an order parameter. As we discussed above the color degrees of freedom play an important role for our system. When the quarks of different colors are in suitable positions in r-space the potential energy largely decreases. It is also very important that the system is locally color white because in this way $`n_c=3`$ and the kinetic part will also decrease. Thus for extremely small densities the quarks should condensate in clusters of 3 and zero net color. When the density changes this picture gradually modifies and at very high densities it does not matter where the quarks are located and which color they have. Looking at Eq.(7) we see that this fact is a consequence of the $`\frac{3}{4}_{a=3,8}\lambda _i^a\lambda _j^a`$ term (equal to $`1`$ for identical color quarks and $`1/2`$ otherwise) . Thus we define an order parameter $`M_c`$ as:
$$M_c=\frac{1}{N}\underset{i=1}{\overset{N}{}}\underset{a=3,8}{}(\lambda _j^a\lambda _k^a+\lambda _i^a\lambda _j^a+\lambda _i^a\lambda _k^a)=M_{c_r}+\frac{1}{N}\underset{i=1}{\overset{N}{}}\underset{a=3,8}{}(\lambda _j^a\lambda _k^a++\lambda _i^a\lambda _k^a)$$
(11)
Where $`j,k`$ are the two quarks closest in r-space to quark $`i`$ as before. In Eq.(11) we have also defined a ”reduced” order parameter $`M_{c_r}`$ which tells us the color of the closest particle $`j`$ to quark $`i`$. From the properties of the (3,8) Gell-Mann matrices it is easy to derive the following results for $`M_c`$. If the 3 closest quarks have the same color, $`M_c=3`$. We stress that this case is practically impossible to be obtained because the corresponding potential energy would be very large and repulsive. If the 3 closest quark states have two colors $`M_c=0`$, this case is also recovered if the colors are randomly distributed such as in the ”pure” Vlasov solution. The case of 3 different color quarks gives $`M_c=3/2`$. The last is the ideal case of well isolated white nucleons. If this last case is recovered in the calculations at small densities then the system is $`\mathrm{𝑙𝑜𝑐𝑎𝑙𝑙𝑦}`$ invariant for rotation in color space. i.e. if we rotate the color states inside the nucleon, the total energy of the system will remain constant. Using similar arguments it is simple to show that if the closest particle to quark $`i`$ has always a different color $`M_{c_r}=1/2`$, if the two closest quarks have the same color $`M_{c_r}=1`$. If the closest quark color is randomly chosen $`M_{c_r}=0`$.
In figure (2) top , we plot the order parameter (opportunely normalized) vs. density (divided by a critical density-see below) for the MD case. The displayed $`M_c`$ are always positive i.e. it never happens that 3 equal quark color states are on average in the same region in r-space. In the top part of the figure we have distinguished two cases. The first one indicated by the full triangles corresponds to calculations where the average potential energy is larger than zero i.e. the linear term in Eq.(7) is dominant, small $`\rho _B`$. The full squares correspond to the case where the Coulomb term is dominant, large $`\rho _B`$ and negative mean field. The first case can be rather well fitted by the relation:
$$M_c|1\frac{\rho _B}{\rho _c}|^\beta $$
(12)
Where $`\rho _c=2.08fm^3`$ is the critical density and the two curves correspond to the critical exponent values $`\beta =1/3`$(full line) and $`1/2`$ (dashed line). The latter is the expected value of the critical exponent $`\beta `$ in a mean field approach . We notice that in ref. a phase transition of first order was found. However the potential and the kinetic term (non-relativistic) used there are quite different from ours. In LQCD calculations for Fermions at zero temperature the order of the phase transition depends on the quark masses. For small quark masses such as ours, the transition found is second order .
In order to better understand the behavior of $`M_c`$ we have repeated the calculations by turning off the Coulomb term (open triangles) or the linear part (open squares). At low densities the linear term is larger while the Coulomb term is dominant at high densities. The two terms are equally important around the critical density. However, there is an important difference between the two cases. In fact at low densities the kinetic part is rather small compared to the potential one, while it is rather large at high densities, Eq.(10). Furthermore, at high densities the strong coupling constant entering the Coulomb potential, Eq.(8), vanishes logarithmically. Thus it is clear that the order parameter increases again at high densities because of the friction used to lower the potential energy. But as soon as we ”turn on” the kinetic term we expect the bonds to be quickly broken. This is shown in figure (3) where the time evolution of $`M_c`$ is given. At high density we have a large value of $`M_c`$ (full circles) at time t=0fm/c. But such a value quickly decreases to a minimum value of about 0.3. The $`M_c`$ obtained at low densities is rather large and constant (open circles), while the one obtained at the critical density slightly increases (open squares). This result also shows that the VE can keep the initial clustering obtained from the MD initial conditions for a time sufficient to perform calculations for heavy ion collisions at relativistic energies. It also implies that the $`M_c`$ values obtained at high $`\rho _B`$ are an artifact of the calculation and we expect $`M_c`$ to saturate at about 0.3-0.4 at high densities. This is exactly the behavior expected for a order parameter, i.e. a power law dependence for densities below the critical one and constant otherwise. It is also instructive to notice that the two closest quarks have always a different color for low densities as it is suggested by $`M_{c_r}`$ in the bottom part of figure(2). The $`M_{c_r}`$ starts to be smaller than 1/2 for densities larger than $`\rho _0`$ and approaches zero very slowly with increasing density.
From figures (2) and (3) some important consequences can be derived.
i) The order parameter is never equal to 3/2 i.e. isolated white nucleons. As it is shown in figure 2, $`M_c`$ is still increasing with decreasing densities. Thus in the limit of very small densities we should get color white objects. However such a limit is hard to reach because of numerical fluctuations due to the confining potential. The maximum calculated $`M_c`$ value is about 0.9,a value in between 2 and 3 color states. Of course it is not always the same cluster (nucleon) to have $`n_c=2`$ or 3, but rather the number of colors in a cluster changes dynamically between 2 and 3. In order to understand this behavior imagine to have two clusters with one color exchanged, $`(r_1,y_1,y_2)`$ and $`(r_2,b_1,b_2)`$, located at very large distances. The contribution to the potential of cluster 1 due to cluster 2 is zero for the red quark (-1+1/2+1/2) and 3/2 for each of the yellow quarks. Thus the two cluster will be $`\mathrm{𝑎𝑡𝑡𝑟𝑎𝑐𝑡𝑒𝑑}`$ towards each other at low densities. On the other hand there is a repulsion between the two equal color quarks in the two clusters. These quarks will be pushed away from their original cluster and eventually the white color will be established. In other words the ”color migration” binds the clusters.
ii) The system is not locally ($`M_c=3/2`$) nor globally ($`M_c=0`$) color invariant. This is true at high densities as well, but there the potential energy is negligible as compared to the kinetic one and color invariance is (approximately) restored.
In conclusion in this work we have discussed microscopic Vlasov/MD approaches to finite nuclei starting from quark degrees of freedom with colors. In order to obtain the correct initial conditions we have introduced a method based on MD with frictional cooling plus a Thomas-Fermi approximation for the Fermi motion. We have shown that the method is able to describe at least qualitatively the well known features of nuclei near the ground state. At high densities a second order phase transition from nuclear to quark matter is predicted. Such a transition is due to the restoration of global color invariance at high densities and we have defined an order parameter accordingly. The approach can be refined in order to obtain a better description of the ground state of the nucleus. This can be used to simulate heavy ion collisions at ultra-relativistic energies after the introduction of a suitable collision term. Our approach can be very useful for the understanding of the quark gluon plasma formation and its signatures.
ACKNOWLEDGMENTS
We thank prof. J.B. Natowitz and the colleagues at the Cyclotron Institute-Texas $`A\&M`$ University for warm hospitality and financial support.
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# Composite fermions close to the one-half filling of the lowest Landau level revisited
\[
## Abstract
By strictly adhering to the microscopic theory of composite fermions (CFs) for the Landau-level filling fractions $`\nu _\mathrm{e}=p/(2p+1)`$, we reproduce, with remarkable accuracy, the surface-acoustic-wave (SAW)-based experimental results by Willett and co-workers concerning two-dimensional electron systems with $`\nu _\mathrm{e}`$ close to $`1/2`$. Our results imply that the electron band mass $`m_b`$, as distinct from the CF mass $`m_{}`$, must undergo a substantial increase under the conditions corresponding to $`\nu _\mathrm{e}1/2`$. We further establish that a finite mean-free path $`\mathrm{}_0`$ is essential for the observed linearity of the longitudinal conductivity $`\sigma _{xx}(q)`$ as deduced from the SAW velocity shifts.
\]
The fractional quantum Hall effect (FQHE) owes its existence to the electron-electron (e-e) interaction. The fermionic Chern-Simons field theory in $`2+1`$ dimensions unifies the FQHE with the integer QHE (IQHE) whose existence does not depend on e-e interaction . This is effected through the binding, brought about by the mediation of the Chern-Simons action, of $`2n`$, $`n=1,2,\mathrm{}`$, magnetic flux quanta to electrons, whereby the composite particles, that is CFs , are exposed to an effective magnetic flux density $`\mathrm{\Delta }B`$ whose corresponding IQH state determines the FQH state for the electrons; the FQH state associated with $`\nu _\mathrm{e}=p/(2np+1)`$ is the IQH state of CFs in which $`p`$ lowest Landau levels (LLs) are fully occupied. The sequence $`\nu _\mathrm{e}=p/(2np+1)`$ and its particle-hole conjugate $`\nu _\mathrm{e}^{}:=1p/(2np+1)`$ approach $`\nu _\mathrm{e}=1/(2n)`$ for $`p\mathrm{}`$. In this work we deal with the case where $`n=1`$ and states whose $`\nu _\mathrm{e}`$ are close to $`1/2`$. With $`\mathrm{\Delta }Bhn_e/(ep)`$, where $`e<0`$ stands for the charge of an electron and $`n_e`$ the planar number density of electrons, these states thus correspond to small $`\mathrm{\Delta }B`$. The state corresponding to $`\nu _\mathrm{e}=1/2`$ was proposed by Halperin, Lee and Read to be a compressible state of degenerate fermions which, in the case of full polarisation of the electron spins, is characterised by a cylindrical Fermi surface of radius $`k_F^{\mathrm{cf}}=\sqrt{4\pi n_e}`$. This has been borne out by several experiments .
In this work we particularly concentrate on a series of experimental results by Willett and co-workers , concerning 2DESs with $`\nu _\mathrm{e}`$ at and close to $`1/2`$ and establish that these can be remarkably accurately reproduced within the framework of the microscopic Chern-Simons field theory. To achieve this, it turns out to be essential that the band-electron mass $`m_b`$ be by one order of magnitude larger than the customarily-assumed value: for GaAs heterostructures, in which the 2DESs under consideration were realised, $`m_b`$ is customarily taken to be $`0.067\times m_e`$, where $`m_e`$ is the electron mass in vacuum. In what follows, we use the notation $`m_b^0=0.067\times m_e`$ and denote the electron band mass, as required for reproducing the indicated experimental results, by $`m_b`$. Since quantum fluctuations, with respect to the mean-field approximation for CFs give rise to mass renormalisation , we reserve $`m_{}`$ to denote the renormalised CF mass. Our numerical results imply $`m_b{}_{}{}^{>}16\times m_b^0`$ and $`m_{}0.5\times m_b{}_{}{}^{>}8\times m_b^00.54\times m_e`$. This value is in good accord with the CF mass as deduced both from the values of the energy gaps $`\mathrm{\Delta }_{\nu _e}`$ separating the CF LLs (determined from the activated temperature dependence of the longitudinal resistivities, $`\varrho _{xx}`$, centred around $`\nu _e`$’s close to $`1/2`$) and the amplitude of the oscillations in $`\varrho _{xx}`$, for varying $`\mathrm{\Delta }B`$, viewing these as the Shubnikov-de Haas oscillations corresponding to CFs . For completeness, according to Park and Jain , the scale of the low-energy excitations of a 2DES may be determined by either the ‘activation mass’ $`m_a`$ or the ‘polarisation mass’ $`m_p`$, the latter being often the case for $`\nu _\mathrm{e}>1`$. For $`\nu _\mathrm{e}=1/2`$ ($`n_e=1.26\times 10^{11}`$ cm<sup>-2</sup>) Kukushkin, et al. have reported a CF mass equal to $`2.27\times m_e`$ which the authors suggest to be the $`m_p`$ of CFs. In view of the fact that the experiments in Ref. are based on optical excitations, involving $`q0`$, we suggest the possibility that the mass measured by these authors may be not $`m_p`$, but $`m_b`$, consistent with the requirement of a Kohn’s theorem (note the aspect $`q0`$).
In the experiments by Willett, et al. , the relative change $`\mathrm{\Delta }v_s/v_s`$ in the velocity $`v_s`$ of SAWs as well as the damping $`\kappa `$ in their amplitudes, while propagating on the surface of samples at distance $`d`$ from the 2DES, were measured; here the change $`\mathrm{\Delta }v_s`$ is relative to the $`v_s`$ corresponding to the case where the conductivity of the 2DES is infinitely large (see further on). Theoretically, $`\mathrm{\Delta }v_s/v_s`$ is determined by the ‘on-the-mass-shell’ value of the longitudinal conductivity $`\sigma _{xx}(q,\omega )`$ of the 2DES as follows
$$\frac{\mathrm{\Delta }v_s}{v_s}=\frac{\alpha ^2/2}{1+[\sigma _{xx}(q)/\sigma _m]^2},$$
(1)
where $`\sigma _{xx}(q)\sigma _{xx}(q,\omega =v_sq)`$, $`\sigma _m`$ is a constant to be specified below, and $`\alpha ^2/2=`$ $`(e_{14}^2\stackrel{~}{ϵ}_r/\{Hϵ_0ϵ_r^2\})`$ $`|F(qd)|^2`$; here $`e_{14}0.145`$ Cm<sup>-2</sup> is the piezoelectric constant for Al<sub>x</sub>Ga<sub>1-x</sub>As, $`\stackrel{~}{ϵ}_r=ϵ_r\left(1+\{(ϵ_r1)/(ϵ_r+1)\}\mathrm{exp}(2qd)\right)^1`$ (c.f. Eq. (A4) in Ref. ), $`H28.8\times 10^{10}`$ Nm<sup>-2</sup>, $`ϵ_0=8.854\mathrm{}\times 10^{12}`$ Fm<sup>-1</sup> denotes the vacuum permittivity, $`ϵ_r`$ is the relative dielectric constant of the bulk of the host material which we take to be equal to $`12.4`$, and $`F(x):=A_1\mathrm{exp}(ax)\mathrm{sin}(bx+c)+A_2\mathrm{exp}(x)`$ (c.f. Eq. (51) in Ref. ), where $`A_13.18`$, $`a0.501`$, $`b0.472`$, $`c2.41`$ and $`A_23.10`$; further,
$$\sigma _m:=\frac{\omega e^2}{q^2v(q)}|_{\omega =v_sq}2ϵ_0\stackrel{~}{ϵ}_rv_s,$$
(2)
where in the second expression on the right-hand side we have replaced the e-e interaction function $`v(q)`$ by the Coulomb function $`v(q)=e^2/(2ϵ_0\stackrel{~}{ϵ}_rq)`$; we have further made use of the dispersion of acoustic phonons and employed $`\omega =v_sq`$, where $`v_s`$ stands for the sound velocity which in GaAs amounts to $`3010`$ ms<sup>-1</sup>.
Our calculations are based on (see also Eq. (208) in Ref. ) $`\mathrm{\Delta }v_s/v_s+i\kappa /q=(\alpha ^2/2)\left(\gamma v(q)K_{00}(q,\omega =v_sq)\right)`$, where $`K_{00}(q,\omega )`$ describes the change in $`n_e`$ to linear order in the external potential ; the minus sign here has its origin in our convention with regard to the sign of $`e`$. In the literature, the constant $`\gamma `$ is invariably identified with unity, this on account of the fact that $`\mathrm{\Delta }v_s`$ is the deviation of the measured $`v_s`$ with respect to the $`v_s`$ pertaining to the case where $`\sigma _{xx}(q)\mathrm{}`$. We maintain $`\gamma `$ in our considerations, following the fact that, experimentally, the reference $`v_s`$ does not correspond to an infinitely large $`\sigma _{xx}(q)`$ (owing to impurities, and $`q0`$).
In the present work we employ the ‘modified random phase approximation’ (MRPA) for $`K_{00}(q,\omega )`$ due to Simon and Halperin which takes account of the renormalisation of the mass of CFs and which approaches the RPA for $`K_{00}(q,\omega )`$ as $`q0`$, the latter coinciding to leading order (proportional to $`q^2`$) with the exact $`K_{00}(q,\omega )`$ for $`q0`$ and any $`\omega `$ at which both $`K_{00}^{\mathrm{rpa}}(q,\omega )`$ and $`K_{00}(q,\omega )`$ are bounded, which is the case for $`|\omega |<\mathrm{\Delta }\omega _c:=e\mathrm{\Delta }B/m_b`$ (see in particular Eq. (5.6) in Ref. ). The $`K_{00}^{\mathrm{mrpa}}`$ has thus the property that it conforms with the requirement of a Kohn’s theorem according to which $`K_{00}(q0,\omega )`$ must be determined by $`m_b`$ rather than $`m_{}`$. For the explicit expression concerning $`K_{00}^{\mathrm{mrpa}}(q,\omega )`$ corresponding to $`\nu _\mathrm{e}=p/(2np+1)`$ in terms of elementary functions, we refer the reader to Ref. (we have presented and employed $`K_{00}^{\mathrm{mrpa}}(q,\omega =0)`$ in Ref. ). We mention that for the purpose of calculating $`K_{00}^{\mathrm{mrpa}}(q,\omega )`$, which depends on $`p`$ and $`n`$, for a continuous range of $`\mathrm{\Delta }B`$ around zero, we first determine $`\nu _\mathrm{e}`$ from $`\nu _\mathrm{e}=hn_e/(eB)`$, with $`B`$ the total applied magnetic flux density, and subsequently obtain the required $`p`$ from $`p=[\nu _\mathrm{e}/(12\nu _\mathrm{e})]`$ if $`\nu _\mathrm{e}<1/2`$, and $`p=[(1\nu _e)/(12(1\nu _\mathrm{e}))]`$ if $`\nu _\mathrm{e}>1/2`$, where $`[x]`$ denotes the greatest integer less than or equal to $`x`$. The $`\nu _\mathrm{e}`$ corresponding to the thus-obtained $`p`$ remains constant for certain ranges of values of $`B`$, which causes artificial stepwise-constant behaviour in the functions of $`\mathrm{\Delta }B`$ that depend on $`K_{00}^{\mathrm{mrpa}}(q,\omega )`$.
We model the effects of the impurity scattering through substituting $`\omega +i/\tau _0`$ for $`\omega `$ in $`K_{00}^{\mathrm{mrpa}}(q,\omega )`$. Here $`\tau _0`$ stands for the scattering time which is related to the mean-free path $`\mathrm{}_0=v_F\tau _0`$, where $`v_F`$ stands for the Fermi velocity. In general, the substitution $`\omega \omega +i/\tau _0`$ amounts to solving the Boltzmann equation within the framework of the relaxation-time approximation, neglecting the so-called current-conservation correction which has been found to have no significant consequences in contexts similar to that of our present considerations (see Sec. 4.1.4 in Ref. ). With (from here onwards we suppress ‘mrpa’) $`K_{00}^{}(q,\omega +i/\tau _0)\mathrm{Re}\left\{K_{00}(q,\omega +i/\tau _0)\right\}`$ and $`\gamma ^{}:=\mathrm{Re}\{\gamma \}`$, from the above-presented expression for $`\mathrm{\Delta }v_s/v_s+i\kappa /q`$ we obtain
$$\frac{\mathrm{\Delta }v_s}{v_s}=\frac{\alpha ^2}{2}\left(\gamma ^{}v(q)K_{00}^{}(q,v_sq+i/\tau _0)\right).$$
(3)
We eliminate $`\gamma ^{}`$, whose value has no influence on the form of $`\mathrm{\Delta }v_s/v_s`$, by requiring that $`\mathrm{\Delta }v_s/v_s`$, as a function of $`\mathrm{\Delta }B`$, coincide with the experimental $`\mathrm{\Delta }v_s/v_s`$ for $`\mathrm{\Delta }B0`$.
In Fig. 1 we present $`\mathrm{\Delta }v_s/v_s`$ as a function of $`\mathrm{\Delta }B`$ for the cases where $`n_e=6.92\times 10^{10}`$ cm<sup>-2</sup>, $`f:=\omega /(2\pi )=8.5`$ GHz (left panel), to be compared with Fig. 4 in Ref. , and $`n_e=1.6\times 10^{11}`$ cm<sup>-2</sup>, $`f=10.7`$ GHz (right panel), to be compared with Fig. 1 in Ref. . The excellent agreement between the theoretical results corresponding to an enhanced $`m_b`$ (with respect to $`m_b^0`$) and experimental results, in particular when these are compared with those obtained within the same theoretical framework in which $`m_{}`$ retains the same enhanced value as compared with $`m_b^0`$ but $`m_b=m_b^0`$ (curves (c)), strongly support the viewpoint that under the conditions where $`\nu _\mathrm{e}1/2`$, the bare band mass $`m_b^0`$ should be enhanced. This observation is compatible with the experimental finding with regard to the stronger than the theoretically-predicted divergence of the CF mass for $`\nu _\mathrm{e}1/2`$ . In this connection we should emphasise that the closer one approaches $`\nu _\mathrm{e}=1/2`$, the less sensitive $`\mathrm{\Delta }v_s/v_s`$ becomes with respect to the further increase of $`m_{}`$ (or $`m_b`$ for that matter); our choices $`m_b=16\times m_b^0`$ and $`m_{}=0.5\times m_b`$ are based on the consideration that the experimental features corresponding to $`\mathrm{\Delta }B`$ in the range $`0.11`$ T be reproduced. The results in Fig. 1 obtained through employing the semi-classical $`\sigma _{\mu \nu }`$, due to Cohen, Harrison and Harrison (CHH) (see also Appendix B in Ref. as well as Eqs. (2) and (3) in Ref. ), bring out the inadequacy of the semi-classical approach; curves marked by (d) unequivocally demonstrate the shortcoming of strictly adhering to the viewpoint that CFs behaved like non-interacting electrons exposed to a reduced magnetic field — curves marked by (b), which are similarly based on the CHH $`\sigma _{\mu \nu }`$, owe their resemblance to the experimental results to the fact that in their calculation explicit account has been taken of the conditions which are specific to the regime corresponding to $`\nu _\mathrm{e}1/2`$. With reference to our earlier work , it is appropriate to compare curve (a) in the right panel of Fig. 1 with the curves in Fig. 1 of the work by Mirlin and Wölfle and compare both with the experimental trace in Fig. 1 of Ref. . One observes that our present result, in contrast with those in Ref. , precisely reproduces almost all features of the experimental trace, such as the values of $`\mathrm{\Delta }v_s/v_s`$ at $`\mathrm{\Delta }B0,0.38,0.54,1.0`$ T.
From Eqs. (1) and (3) we obtain
$$\frac{\sigma _{xx}(q)}{\sigma _m}=\left(\left(\gamma ^{}v(q)K_{00}^{}(q,v_sq+i/\tau _0)\right)^11\right)^{1/2},$$
(4)
according to which $`\sigma _{xx}(q)`$ interestingly does not explicitly depend upon $`\alpha ^2/2`$. Unless we set $`\gamma ^{}=1`$, we eliminate $`\gamma ^{}`$ in Eq. (4) by requiring that for given values of $`\sigma _m`$ and $`q`$, $`\sigma _{xx}(q)`$ according to Eq. (4) yield the corresponding experimental SAW velocity shift.
In Fig. 2 we present our theoretical $`\sigma _{xx}(q)`$ in comparison with its SAW-derived experimental $`\sigma _{xx}(q)`$ by Willett, et al. (see the middle panel of Fig. 2 herein). The details in Fig. 2 again support our above finding with regard to $`m_b`$ and $`m_{}`$, that a mere enhancement of $`m_{}`$ with respect to $`m_b^0`$ is not sufficient (see inset \[B\]). We also observe that a finite $`\mathrm{}_0`$ is most crucial to the experimentally-observed linear behaviour of the SAW-deduced $`\sigma _{xx}(q)`$ for $`q`$ in the range $`(0.015/a_0,0.075/a_0`$), with $`a_0`$ the Bohr radius (see curve (c)); the original observation with regard to $`\sigma _{xx}(q)q`$ for $`q2/\mathrm{}_0`$ thus turns out to be relevant for values of $`q`$ far outside the experimental range. We note that $`\mathrm{}_00.5`$ $`\mu `$m coincides with that reported in the pertinent experimental articles. A further aspect that our present results in Fig. 2 clarify is that, in contrast to earlier observations (see the paragraph following Eq. (7.6) in Ref. and that following Eq. (211) in Ref. ), the available experimental results by no means are in conflict with the predictions of Eq. (2): our theoretical results for $`\sigma _{xx}(q)`$ in Fig. 2 have been obtained through multiplying $`\sigma _{xx}(q)/\sigma _m`$ by the same $`\sigma _m`$ that has been employed to determine $`\sigma _{xx}(q)`$ from the SAW-deduced $`\sigma _{xx}(q)/\sigma _m`$. Thus, rather than Eq. (2) being inadequate, the empirical method of determining $`\sigma _m`$ (which employs the dc conductivity ) should be considered as inappropriate.
We now briefly focus on the physical significance of the expression for $`\mathrm{\Delta }v_s/v_s`$ in Eq. (1). To this end, let $`\delta v_{\mathrm{ext}}(𝐫,t)`$ denote an applied time-dependent external potential, representing that corresponding to the SAWs. The change in the Hamiltonian of the system, following the application of $`\delta v_{\mathrm{ext}}`$, has the form $`\delta \widehat{H}(t)=\mathrm{d}^2r\delta v_{\mathrm{ext}}(𝐫,t)\widehat{\psi }^{}(𝐫t)\widehat{\psi }(𝐫t)`$, where $`\widehat{\psi }^{}(𝐫t)`$, $`\widehat{\psi }(𝐫t)`$ denote creation and annihilation field operators. Denoting the change in the energy of the system corresponding to $`\delta \widehat{H}(t)`$ by $`\delta E(t)`$, we have $`\delta E(t)=`$ $`\mathrm{d}^2r`$ $`\delta v_{\mathrm{ext}}(𝐫,t)\overline{n}_e(𝐫,t)`$, where $`\overline{n}_e(𝐫,t):=_0^1d\lambda n_e^{(\lambda )}(𝐫,t)`$. Here $`n_e^{(\lambda )}(𝐫,t)`$ stands for the instantaneous number density of the system corresponding to $`\delta v_{\mathrm{ext}}^{(\lambda )}(𝐫,t):=\lambda \delta v_{\mathrm{ext}}(𝐫,t)`$. By assuming $`\delta v_{\mathrm{ext}}(𝐫,t)\delta v_{\mathrm{ext}}(𝐫)\mathrm{exp}(i\omega _0t)`$, one obtains for $`\delta (\omega ):=dt`$ $`\delta E(t)\mathrm{exp}(i\omega t)`$, $`\delta (\omega )=\mathrm{d}^2r\mathrm{d}^2r^{}\delta v_{\mathrm{ext}}(𝐫)\overline{K}_{00}(𝐫,𝐫^{};\omega \omega _0)\delta v_{\mathrm{ext}}(𝐫^{})`$, where $`\overline{K}_{00}(𝐫,𝐫^{};\omega ):=_0^1d\lambda (1\lambda )K^{(\lambda )}(𝐫,𝐫^{};\omega )`$ . Under the assumption that $`\delta v_{\mathrm{ext}}(𝐫)`$ be weak, the dependence upon $`\lambda `$ of $`K_{00}^{(\lambda )}`$ can be neglected so that $`\overline{K}_{00}2^1K_{00}^{(0)}`$ where $`K_{00}^{(0)}`$ stands for the density-density response function of the uniform, unperturbed, system. To second order in $`\delta v_{\mathrm{ext}}(q)`$, one thus obtains $`\delta (\omega )=2^1K_{00}(q,\omega \omega _0)|\delta v_{\mathrm{ext}}(q)|^2`$. Let now $`\overline{\delta E}_t:=_0^td\tau \delta E(\tau )`$, from which one readily obtains $`\overline{\delta E}_t2^1|\delta v_{\mathrm{ext}}(q)|^2d\omega K_{00}(q,\omega \omega _0)\{1\mathrm{exp}(i\omega t)\}/(2\pi i\omega )`$. For large $`t`$, the integrand of the $`\omega `$ integral becomes highly oscillatory so that to the leading order in $`1/t`$, $`d\omega K_{00}(q,\omega \omega _0)\{1\mathrm{exp}(i\omega t)\}/(2\pi i\omega )K_{00}(q,\omega _0)d\omega \{1\mathrm{exp}(i\omega t)\}/(2\pi i\omega )=K_{00}(q,\omega _0)`$. Since $`K_{00}(q,\omega )`$ is an even function of $`\omega `$, we eventually obtain $`\overline{\delta E}_t2^1K_{00}(q,\omega _0)|\delta v_{\mathrm{ext}}(q)|^2`$, for $`t1/\omega _0`$. This expression, which coincides with Eq. (14) in Ref. , is the fundamental link between $`K_{00}(q,\omega )`$ and $`\mathrm{\Delta }v_s/v_s+i\kappa /q`$ presented above. These details make explicit first, that only small-amplitude perturbations are correctly accounted for by Eq. (1) (and similarly, Eq. (3)), and second, that the observation of “geometric resonance” and the “cyclotron frequency deduced from dc transport” though, as suggested by Willett, et al. , inconsistent with “a non-interacting, semi-classical quasiparticle model \[for CFs\]”, are in fact not inconsistent with the physical picture that the above derivation brings out: that the mechanism underlying the SAW-based experiments does not involve any resonance phenomenon in the usual sense and that the SAW experiments, which involve a long-time integration of the fluctuations in the total energy of 2DESs, unveil $`K_{00}(q,\omega )`$ at $`\omega =\omega _0v_sq`$ and $`q=\omega _0/v_s`$, independent of the magnitude of the CF cyclotron frequency $`\mathrm{\Delta }\omega _c`$ and consequently of that of the CF mass $`m_{}`$.
In conclusion, by strictly adhering to the microscopic theory of CFs, we have established that the SAW-based experimental results by Willett and co-workers close to $`\nu _\mathrm{e}=1/2`$ can be remarkably accurately reproduced provided the electron band mass $`m_b`$ be substantially enhanced with respect to $`m_b^0`$; in this picture, the (observed) large value of the CF mass $`m_{}`$ follows from a subsequent reduction of $`m_b`$ (owing to quantum fluctuations) rather than a direct enhancement of $`m_b^0`$. We have further established that a finite mean-free path $`\mathrm{}_0`$ is essential to the experimentally-observed linearity in the SAW-deduced $`\sigma _{xx}(q)`$ in the range $`1.56{}_{}{}^{<}q{}_{}{}^{<}7.68`$ $`\mu `$m<sup>-1</sup>, and that there exists no discrepancy between the theoretical and experimental values for $`\sigma _m`$.
I thank Professors D. E. Khmel’nitskiǐ and P. B. Littlewood for a discussion and Professor B. I. Halperin and Drs S. H. Simon and R. L. Willett for kindly clarifying some aspects concerning the empirical $`\sigma _m`$. With appreciation I acknowledge hospitality of Cavendish Laboratory.
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# The Verlinde formula for parabolic bundles
## 1 Introduction
Let $`\mathrm{\Sigma }^g`$ be a compact Riemann surface of genus $`g`$, and $`G`$ $`=SU(n)`$. We introduce the central element $`c`$ $`=\mathrm{diag}(e^{2\pi id/n},\mathrm{},e^{2\pi id/n})`$ for $`d`$ coprime to $`n.`$ In this paper we prove the Verlinde formula for the Riemann-Roch number of a line bundle over the moduli space $`_{g,1}(c,\mathrm{\Lambda })`$ of representations of the fundamental group of a Riemann surface of genus $`g`$ with one boundary component, for which the loop around the boundary is constrained to lie in the conjugacy class of $`c\mathrm{exp}(\mathrm{\Lambda })`$ (for $`\mathrm{\Lambda }𝐭_+)`$. (Here, $`𝐭`$ denotes the Lie algebra of the maximal torus $`T`$ of $`G`$, and $`𝐭_+`$ the fundamental Weyl chamber.) We also prove the Verlinde formula for the moduli space $`_{g,b}(c,𝚲)`$ of representations of the fundamental group of a Riemann surface of genus $`g`$ with $`s+1`$ boundary components for which the loop around the 0-th boundary component is sent to the central element $`c`$ and the loop around the $`j`$-th boundary component is constrained to lie in the conjugacy class of $`\mathrm{exp}(\mathrm{\Lambda }^{(j)})`$ for $`\mathrm{\Lambda }^{(j)}𝐭_+`$, where we have introduced the notation $`𝚲=(\mathrm{\Lambda }^{(1)},\mathrm{\Lambda }^{(2)},\mathrm{},\mathrm{\Lambda }^{(b)}).`$ The methods extend the proof we gave in Section 11 of for the Verlinde formula for the moduli space $`M(n,d)`$ of holomorphic vector bundles of coprime rank $`n`$ and degree $`d`$ and fixed determinant, which can alternatively be described as the space of representations of the fundamental group of a Riemann surface of genus $`g`$ with one boundary component into $`G`$ which send the loop around the boundary to the central element $`c`$. Our proofs are valid for $`\mathrm{\Lambda }`$ in a suitable neighbourhood of $`0`$ and $`𝚲`$ in a suitable neighbourhood of $`\mathrm{𝟎}`$.
Our earlier work used the Riemann-Roch formula to prove a formula for the dimension of the space of holomorphic sections of powers of a certain line bundle over $`M(n,d)`$; in the case of $`M(n,d)`$ the higher cohomology vanishes, so the Riemann-Roch formula gives the dimension of the zeroth cohomology group. In this paper we exploit the fact that the related moduli spaces $`_{g,1}(c,\mathrm{\Lambda })`$ and more generally $`_{g,b}(c,𝚲)`$ defined above (which appear in algebraic geometry as moduli spaces of bundles with parabolic structure, where the parameters $`\mathrm{\Lambda }^{(s)}`$ $`𝐭_+`$ are equivalent to the specification of weights; see for example ) fibre over $`M(n,d)`$ provided the weights are sufficiently small. When these spaces admit a prequantum line bundle, we may push forward along the fibre to obtain a formula for the Riemann-Roch number of the prequantum line bundle on the total space in terms of the evaluation of appropriate cohomology classes on the fundamental class of $`M(n,d)`$. Formulas for the intersection numbers in $`M(n,d)`$ were proved in ; we apply these formulas together with the fibration to recover the Verlinde formula in this more general situation.
The layout of this paper is as follows. In Section 2 we review results from on the cohomology ring of $`M(n,d)`$, while in Section 3 we summarize results on symplectic fibrations. Section 4 contains our proofs of the Verlinde formula: in Section 4.1 we first review the proof for $`M(n,d)`$ from , while Section 4.2 gives the proof for $`_{g,1}(c,\mathrm{\Lambda })`$ and Section 4.3 gives the proof for $`_{g,b}(c,𝚲)`$.
Acknowledgements: This paper relies heavily on the author’s earlier joint work with F. Kirwan, notably on the paper where the Verlinde formula for $`M(n,d)`$ is proved. We would like to acknowledge the hospitality of Université Paris-Sud (Orsay), where part of the work was completed.
## 2 Review of results on the cohomology of moduli spaces
Theorem (Atiyah-Bott 1982) The cohomology ring of $`M(n,d)`$ over $`𝐐`$ is generated by elements $`\{a_r,b_r^j,f_r\}`$ (for $`r=2,\mathrm{},n`$ and $`j=1,\mathrm{},2g`$), where
$$a_rH^{2r}(M(n,d)),$$
$$b_r^jH^{2r1}(M(n,d)),$$
$$f_rH^{2r2}(M(n,d)).$$
If $`U`$ is the universal vector bundle over $`M(n,d)\times \mathrm{\Sigma }^g`$, the elements $`a_r,b_r^j,f_r`$ are obtained by decomposing $`c_r(U)`$ in terms of
$$_{s=0}^2H^{2rs}(M(n,d))H^s(\mathrm{\Sigma }^g)$$
using the Künneth formula.
###### Remark 2.1
The Verlinde formula (for $`M(n,d)`$ as well as for $`_{g,1}(c,\mathrm{\Lambda })`$ and $`_{g,b}(c,𝚲)`$) may be deduced from pairings involving only $`f_2`$ and the $`a_j`$.
We review the results of , where formulas were proved for intersection pairings in the cohomology of $`M(n,d)`$.
###### Remark 2.2
If $`M`$ is a symplectic manifold equipped with the Hamiltonian action of a Lie group $`G`$, and $`0`$ is a regular value of the moment map $`\mu :M𝐠^{}`$, then the Kirwan map is the map
$$\kappa :H_G^{}(M)H_G^{}(\mu ^1(0))H^{}(M_{\mathrm{red}}).$$
In this paper we are using the notation $`\kappa `$ to refer to the restriction to the image of $`H_G^{}(\mathrm{pt})S(𝐠^{})^G`$ in the domain $`H_G^{}(M)`$. In the case where $`M_{\mathrm{red}}=M(n,d)`$, this restricted map is simply the Chern-Weil map for the universal bundle restricted to $`M(n,d)\times \{\mathrm{pt}\}.`$
###### Theorem 2.3
Let $`c=\mathrm{diag}(e^{2\pi id/n},\mathrm{},e^{2\pi id/n})`$ where $`d\{1,\mathrm{},n1\}`$ is coprime to $`n`$, and suppose that $`\eta H_G^{}(\{\mathrm{pt}\})`$ is a polynomial $`Q(\tau _2,\mathrm{},\tau _n)`$ in the equivariant cohomology classes $`\tau _rH_G^{2r}(\mathrm{pt})`$ $`S(𝐠^{})^G`$ for $`2rn`$ (where $`\tau _r`$ is the $`r`$-th elementary symmetric polynomial), which generate the $`G`$-equivariant cohomology of a point, and map under the Kirwan map $`\kappa `$ to classes $`a_rH^{2r}(M(n,d))`$. Let $`f_2`$ be the cohomology class of the symplectic form. Then the pairing $`\kappa (\eta )\mathrm{exp}(f_2)[M(n,d)]`$ is given by
$$_{M(n,d)}\kappa (\eta )\mathrm{exp}(f_2)=\frac{(1)^{n_+(g1)}}{n!}\mathrm{Res}_{Y_1=0}\mathrm{}\mathrm{Res}_{Y_{n1}=0}\left(\frac{_{wW_{n1}}e^{[[w\stackrel{~}{c}]],X}_{T^{2g}}\eta e^\omega }{𝒟^{2g2}_{1jn1}(\mathrm{exp}(Y_j)1)}\right),$$
where $`n_+=\frac{1}{2}n(n1)`$ is the number of positive roots of $`G=SU(n)`$ and $`X=(X_1.\mathrm{},X_n)𝐭𝐂`$ has coordinates $`Y_1=X_1X_2,\mathrm{},Y_{n1}=X_{n1}X_n`$ defined by the simple roots, while $`W_{n1}S_{n1}`$ is the Weyl group of $`SU(n1)`$ embedded in $`SU(n)`$ in the standard way using the first $`n1`$ coordinates. The quantity $`𝒟(X)=_{\gamma >0}\gamma (X)`$ is the product of the positive roots. The element $`\stackrel{~}{c}`$ is the unique element of $`𝐭`$ which satisfies $`\mathrm{exp}\stackrel{~}{c}=c`$ and belongs to the fundamental domain defined by the simple roots for the translation action on $`𝐭`$ of the integer lattice $`\mathrm{\Lambda }^I`$ (in other words, the fundamental alcove).
Also, the notation $`[[\gamma ]]`$ means the unique element which is in the fundamental domain defined by the simple roots for the translation action on $`𝐭`$ of the integer lattice and for which $`[[\gamma ]]`$ is equal to $`\gamma `$ plus some element of the integer lattice.
Theorem 2.3 is a special case of Theorem 8.1 of : this special case suffices to prove the Verlinde formulas.
## 3 Symplectic fibrations
Let $`\stackrel{~}{c}𝐭`$ be an element of the closed fundamental alcove $`D_+`$ satisfying $`\mathrm{exp}\stackrel{~}{c}=c`$. The interior of the fundamental alcove will be denoted $`D_+^o`$.
Let $`\mathrm{\Sigma }_n^g`$ denote an oriented two-manifold of genus $`g`$ with $`b`$ oriented boundary components $`S_1,\mathrm{},S_b`$.
###### Definition 3.1
Let $`𝚲=(\mathrm{\Lambda }^{(1)},\mathrm{},\mathrm{\Lambda }^{(b)})`$ be a collection of $`b`$ values in $`D_+`$. The moduli space of representations is defined by
$$_{g,b}(𝚲)=_{g,b}(𝚲)/G,$$
where
$$_{g,b}(𝚲)=\{\rho \mathrm{Hom}(\pi _1(\mathrm{\Sigma }_b^g),G):\rho ([S_a])\mathrm{Cl}(\mathrm{exp}\mathrm{\Lambda }^{(a)}),a=1,\mathrm{},b\}$$
(3.1)
and $`G`$ acts on $`_{g,b}(𝚲)`$ by conjugation. Here, $`\mathrm{Cl}(\mathrm{exp}\mathrm{\Lambda }^{[a]})`$ denotes the conjugacy class of $`\mathrm{exp}\mathrm{\Lambda }^{[a]}`$ in $`G`$.
The fundamental group of $`\mathrm{\Sigma }_b^g`$ is the free group on $`2g+b`$ generators with one relation:
$$\pi _1(\mathrm{\Sigma }_n^g)=<x_1,\mathrm{},x_{2g},y_1,\mathrm{},y_b:\underset{j=1}{\overset{g}{}}[x_j,x_{j+g}]=\underset{r=1}{\overset{b}{}}y_r>.$$
Thus we have
$$_{g,n}(𝚲)=\{(h_1,\mathrm{},h_{2g},\beta _1,\mathrm{},\beta _b)G^{2g+b}:\underset{j=1}{\overset{g}{}}[h_j,h_{j+g}]=\underset{r=1}{\overset{b}{}}\beta _r,\beta _r\mathrm{Cl}(\mathrm{exp}\mathrm{\Lambda }^{(r)})\}/G.$$
(3.2)
For later convenience we make a slight modification of this definition: We put
$$_{g,b}(c,𝚲)=\{(h_1,\mathrm{},h_{2g},\beta _1,\mathrm{},\beta _b)G^{2g+b}:\underset{j=1}{\overset{g}{}}[h_j,h_{j+g}]=c\underset{r=1}{\overset{b}{}}\beta _r,\beta _r\mathrm{Cl}(\mathrm{exp}\mathrm{\Lambda }^{(r)})\}/G,$$
(3.3)
where $`c=\mathrm{diag}(e^{2\pi id/n},\mathrm{},e^{2\pi id/n}).`$ In particular we have
$$_{g,1}(c,\mathrm{\Lambda })=\{(h_1,\mathrm{},h_{2g},\beta )G^{2g+1}:\underset{j=1}{\overset{g}{}}[h_j,h_{j+g}]=c\beta ,\beta \mathrm{Cl}(\mathrm{exp}\mathrm{\Lambda })\}/G.$$
Results on symplectic fibrations of moduli spaces were developed in , and used there for purposes distinct from the objectives of the present paper. We have
###### Theorem 3.2
There is a neighbourhood $`U`$ of $`0`$ in $`𝐭`$ such that if $`\mathrm{\Lambda }U`$ then there is a fibration
$$\pi :_{g,1}(c,\mathrm{\Lambda })M(n,d)$$
(3.4)
with fibre $`𝒪_\mathrm{\Lambda }`$ (the orbit of the adjoint action of $`G`$ on $`𝐠`$). Further, the symplectic form $`\omega _\mathrm{\Lambda }`$ on $`_{g,1}(\mathrm{\Lambda })`$ satisfies
$$\omega _\mathrm{\Lambda }=\pi ^{}\omega _{n,d}+\stackrel{~}{\mathrm{\Omega }}_\mathrm{\Lambda }$$
(3.5)
where $`\omega _{n,d}`$ is the symplectic form on $`M(n,d)`$ and $`\stackrel{~}{\mathrm{\Omega }}_\mathrm{\Lambda }`$ restricts on each fibre of $`\pi `$ to the standard Kirillov-Kostant symplectic form $`\mathrm{\Omega }_\mathrm{\Lambda }`$ on the coadjoint orbit $`𝒪_\mathrm{\Lambda }`$.
Proof: This follows by general results regarding symplectic fibrations associated to symplectic reduction at a regular value: see , Theorem 6.1 for a proof. Note that $`M(n,d)`$ is obtained by reducing an appropriate extended moduli space at a regular value of the moment map and that the 2-form $`\omega _\mathrm{\Lambda }`$ is nondegenerate in a neighbourhood of the preimage of this regular value under the moment map: see Proposition 5.5 of .
###### Remark 3.3
The region $`U`$ is characterized as the component of the complement of the union of the walls $`H_{v,n}`$ in $`𝐭_+`$ whose closure contains $`0`$, where
$$H_{v,n}=\{\mathrm{\Lambda }𝐭:v(\mathrm{\Lambda })=n\}$$
for one of the fundamental weights $`v`$ and $`n𝐙`$. Clearly the $`H_{v,n}`$ are the hyperplanes where the function encoding the volume of the $`_{g,1}(\mathrm{\Lambda })`$ is not smooth : in other words they bound the region of regular values of the moment map, which is the region where the symplectic fibration of Theorem 3.2 exists. We thank A. Szenes for this observation, which is explained in .
###### Proposition 3.4
If $`\mathrm{\Lambda }`$ is a regular element of $`𝐭^{}`$, the coadjoint orbit $`𝒪_\mathrm{\Lambda }`$ is diffeomorphic to the homogeneous space $`G/T`$ so its cohomology is given by
$$H^{}(𝒪_\mathrm{\Lambda })\frac{S(𝐭^{})}{S(𝐭^{})^W},$$
(3.6)
in other words the quotient of ring of polynomials on $`𝐭`$ by the subring of symmetric polynomials.
We have the following proposition:
###### Proposition 3.5
The space $`_{g,1}(c,\mathrm{\Lambda })`$ is a splitting manifold for the universal bundle $`U_{M(n,d)\times \{\mathrm{pt}\}}`$ over $`M(n,d)\times \{\mathrm{pt}\}M(n,d)\times \mathrm{\Sigma }_0^g:`$ in other words
$$\pi ^{}\left(U_{M(n,d)\times \{\mathrm{pt}\}}\right)=L_1\mathrm{}L_n$$
where $`c_1(L_j)=e_j`$ for a collection of classes $`e_j`$ in $`H^2\left(_{g,1}(c,\mathrm{\Lambda })\right)`$ (for $`j=1,\mathrm{},n`$). Here, when $`j=1,\mathrm{},n1`$, $`e_j`$ restricts on the fibres of $`\pi `$ to the generator $`\alpha _j`$ $`(j=1,\mathrm{},n1)`$ of $`H^2(G/T,𝐙)H^1(T,𝐙)`$ corresponding to the $`j`$-th fundamental weight of $`SU(n)`$ (an element of $`\mathrm{Hom}(T,U(1))`$, which is isomorphic to $`H^1(T,𝐙)`$) and $`e_n=(e_1+\mathrm{}+e_{n1}).`$
Proof: This follows from the algebro-geometric description of the moduli space of parabolic bundles (see for instance ): it is the moduli space parametrizing holomorphic bundles over $`\mathrm{\Sigma }_0^g`$ together with a flag in the fibre of each bundle over a basepoint $`(\{pt\})\mathrm{\Sigma }_0^g`$. The flag structure enables us naturally to split the universal bundle into a sum of holomorphic line bundles.
###### Proposition 3.6
If $`\tau _r`$ is the $`r`$-th elementary symmetric polynomial (for $`r=2,\mathrm{},n`$) then $`\tau _r(e_1,\mathrm{},e_n)=\pi ^{}a_r`$, where $`a_r=c_r(U_{M(n,d)\times \{\mathrm{pt}\}}).`$
Proof: See p. 284 (Section 21) of for results on the properties of splitting manifolds and flag bundles. There, it is proved that for a complex vector bundle $`E`$ over a complex manifold $`M`$ with splitting manifold $`\mathrm{Fl}(E)`$, we have
$$H^{}\left(\mathrm{Fl}(E)\right)=\frac{H^{}(M)[e_1,\mathrm{},e_n]}{_{i=1}^n(1+e_i)=c(E)},$$
(3.7)
where the $`e_jH^2(\mathrm{Fl}(E))`$ restrict (for $`j=1,\mathrm{},n`$) on the fibre $`U(n)/U(1)^n`$ $`G/T`$ (where $`G=SU(n)`$ and $`T`$ is its maximal torus) to the images under the coboundary map (in the Leray-Serre spectral sequence) of the elements $`H_{U(1)^n}^1(\{pt\},𝐙)`$ $`=\mathrm{Hom}(U(1)^n,U(1))`$ given by a basis for the weight lattice of $`U(1)^n`$.
The following is a standard result (see for instance , Lemma 7.22):
###### Proposition 3.7
Let $`\alpha _1,\mathrm{},\alpha _n`$ (subject to $`_{j=1}^n\alpha _j=0`$) be the basis for $`H_T^2(\{pt\})`$ (the second equivariant cohomology group of a point for the maximal torus $`T`$ of $`SU(n)`$) which was introduced in Proposition 3.5. Let $`\mathrm{\Lambda }=_{i=1}^{n1}\mathrm{\Lambda }_i\widehat{u}_i`$ where $`\widehat{u}_i`$ is the $`i`$-th simple root: note that the simple roots are the basis of $`𝐭`$ dual to the fundamental weights, which were introduced in Proposition 3.5 to define the generators $`\alpha _j`$. Then the standard Kirillov-Kostant symplectic form $`\mathrm{\Omega }_\mathrm{\Lambda }`$ on $`𝒪_\mathrm{\Lambda }`$ is given by
$$\mathrm{\Omega }_\mathrm{\Lambda }=\underset{j=1}{\overset{n}{}}\mathrm{\Lambda }_j\alpha _j.$$
## 4 Application to the Verlinde formula
### 4.1 The Verlinde formula for $`M(n,d)`$
###### Definition 4.1
The highest root $`\gamma _{\mathrm{max}}`$ is given by $`\gamma _{\mathrm{max}}(X)=X_nX_1`$ or $`\gamma _{\mathrm{max}}(X)=Y_1+\mathrm{}+Y_{n1}`$.
###### Definition 4.2
The Verlinde function $`V_{n,d}(g,k)`$ is given by
$$V_{n,d}(g,k)=\underset{\mu \mathrm{\Delta }(r)}{}\frac{e^{2\pi i\mu \rho ,\stackrel{~}{c}}}{(S_{0\mu }(k))^{2g2}}$$
where $`\rho `$ is half the sum of the positive roots and
$$S_{0\mu }(k)=\frac{1}{\sqrt{n}r^{(n1)/2}}\underset{\gamma >0}{}2\mathrm{sin}\pi \gamma ,\mu /r.$$
(See (A.44) and (3.16).) Here, $`\mathrm{\Lambda }^W`$ is the weight lattice, identified with points in $`𝐭`$. We have introduced the quantity
$$r=k+n;$$
(4.1)
we have also introduced
$$\mathrm{\Delta }(r)=\{\mu \mathrm{\Lambda }_{\mathrm{reg}}^W𝐭_+:\mu ,\gamma _{\mathrm{max}}<r\}.$$
The Verlinde formula is a formula for the dimension $`D_{n,d}(g,k)`$ of the space of holomorphic sections of powers of $``$ , where $``$ is a particular line bundle over $`M(n,d)`$: it has been proved by Beauville and Laszlo , Faltings , Kumar, Narasimhan and Ramanathan and Tsuchiya, Ueno and Yamada . In Bismut and Labourie have given a proof of the Verlinde formula using techniques from symplectic geometry.
In this section we review the results from Section 11 of , showing how the Verlinde formula follows from the formula (Theorem 2.3) for intersection pairings in $`M(n,d)`$.
A line bundle $``$ over $`M(n,d)`$ may be defined for which $`c_1()=nf_2`$, since $`nf_2H^2(M(n,d),𝐙)`$ (see ). Whenever $`k`$ is a positive integer divisible by $`n`$, we then define
$$D_{n,d}(g,k)=dimH^0(M(n,d),^{k/n}).$$
(4.2)
Verlinde’s conjecture says that the Verlinde function specifies the dimension of the space of holomorphic sections of $`^{k/n}`$:
###### Theorem 4.3
(Verlinde’s conjecture)
$$D_{n,d}(g,k)=V_{n,d}(g,k).$$
We review the method of Section 11 of , where we gave a proof of Verlinde’s conjecture for $`M(n,d)`$: an outline of the method we use was given by Szenes (Section 4.2).
In fact $`H^i(M(n,d),^m)=0`$ for all $`i>0`$ and $`m>0`$ (see Section 11 of for references and an outline of the proof). So $`D_{n,d}(g,k)`$ is given for $`k>0`$ by the Riemann-Roch formula:
$$D_{n,d}(g,k)=_{M(n,d)}\mathrm{ch}^{k/n}\mathrm{Td}M(n,d).$$
(4.3)
We use the following results to convert (4.3) into a form to which we may apply our previous results.
###### Lemma 4.4
For any complex manifold $`M`$ the Todd class of $`M`$ is given by
$$\mathrm{Td}(M)=e^{c_1(M)/2}\widehat{A}(M)$$
where $`c_1(M)`$ is the first Chern class of the holomorphic tangent bundle of $`M`$, and $`\widehat{A}(M)`$ is the $`A`$-hat genus of $`M`$.
Proof: See for example , pages 97-99.
###### Proposition 4.5
We have
$$\widehat{A}(M(n,d))=\kappa \left(\underset{\gamma >0}{}\frac{\gamma /2}{\mathrm{sinh}\gamma /2}\right)^{2g2}.$$
Proof: This is proved by Newstead in .
###### Lemma 4.6
We have
$$c_1(M(n,d))=2nf_2.$$
Proof: This is proved in , Théorème F .
Of course the Chern character of $`^{k/n}`$ is given by $`\mathrm{ch}^{k/n}=e^{kf_2}.`$ Thus we obtain
###### Proposition 4.7
The quantity $`D_{n,d}(g,k)`$ is given by
$$D_{n,d}(g,k)=_{M(n,d)}e^{(k+n)f_2}\kappa \left(\underset{\gamma >0}{}\frac{\gamma }{e^{\gamma /2}e^{\gamma /2}}\right)^{2g2}.$$
Proof: This follows immediately from (4.3), Lemmas 4.4 and 4.6 and Proposition 4.5.
###### Theorem 4.8
We have
$$D_{n,d}(g,k)=\frac{(1)^{n_+(g1)}}{n!}\underset{wW_{n1}}{}\mathrm{Res}_{Y_1=0}\mathrm{}\mathrm{Res}_{Y_{n1}=0}(e^{r[[w\stackrel{~}{c}]],X}_{T^{2g}}e^{r\omega }\times $$
$$\underset{\gamma >0}{}\left(\frac{\gamma (X)}{e^{\gamma (X)/2}e^{\gamma (X)/2}}\right)^{2g2}\frac{1}{_{j=1}^l(e^{rY_j}1)𝒟(X)^{2g2}}).$$
(4.4)
Proof: This is a direct consequence of Corollary 4.7 and Theorem 2.3. Note that because the factor $`e^{f_2}`$ in the statement of Theorem 2.3 has been replaced by $`e^{rf_2}`$, it is necessary to replace $`e^{[[w\stackrel{~}{c}]],X}`$ by $`e^{r[[w\stackrel{~}{c}]],X}`$, and $`e^{Y_j}1`$ by $`e^{rY_j}1`$.
We introduce $`Z_j=\mathrm{exp}Y_j`$. Since for any $`wW_{n1}`$ we have that
$$[[w\stackrel{~}{c}]]=[[w\stackrel{~}{c}]]_1\widehat{u}_1+[[w\stackrel{~}{c}]]_2\widehat{u}_2+\mathrm{}+[[w\stackrel{~}{c}]]_{n1}\widehat{u}_{n1}$$
in terms of the standard basis vectors
$$\widehat{u}_j=(0,\mathrm{},0,1,1,0,\mathrm{},0)$$
for the integer lattice $`\mathrm{\Lambda }^I`$ of $`𝐭`$, with $`n[[w\stackrel{~}{c}]]_j𝐙`$ for all $`j`$, and $`0[[w\stackrel{~}{c}]]_j<1`$ for all $`j`$. We obtain
$$e^{r[[w\stackrel{~}{c}]]_,X}=Z_1^{[[w\stackrel{~}{c}]]_1r}Z_2^{[[w\stackrel{~}{c}]]_2r}\mathrm{}Z_{n1}^{[[w\stackrel{~}{c}]]_{n1}r}.$$
(Recall that $`k`$ and $`r`$ are divisible by $`n`$ so $`e^{r\stackrel{~}{c},X}`$ is a well defined single valued function of $`Z_1`$, $`\mathrm{},Z_{n1}`$.) Thus we can equate $`D_{n,d}(g,k)`$ with
$$\frac{(1)^{n_+(g1)}}{n!}\underset{wW_{n1}}{}\mathrm{Res}_{Z_1=1}\mathrm{}\mathrm{Res}_{Z_{n1}=1}(\left(\underset{j=1}{\overset{n1}{}}\frac{1}{Z_j}\right)_{T^{2g}}e^{r\omega }\times $$
$$\frac{Z_1^{[[w\stackrel{~}{c}]]_1r}Z_2^{[[w\stackrel{~}{c}]]_2r}\mathrm{}Z_{n1}^{[[w\stackrel{~}{c}]]_{n1}r}}{_{\gamma >0}(\stackrel{~}{\gamma }^{1/2}\stackrel{~}{\gamma }^{1/2})^{2g2}(Z_1^r1)\mathrm{}(Z_{n1}^r1)})$$
$$=\frac{(1)^{n1+n_+(g1)}}{n!}\underset{wW_{n1}}{}\mathrm{Res}_{Z_1=1}\mathrm{}\mathrm{Res}_{Z_{n1}=1}(\left(\underset{j=1}{\overset{n1}{}}\frac{1}{Z_j}\right)\times $$
(4.5)
$$_{T^{2g}}e^{r\omega }\frac{Z_1^{[[w\stackrel{~}{c}]]_1r}Z_2^{[[w\stackrel{~}{c}]]_2r}\mathrm{}Z_{n1}^{[[w\stackrel{~}{c}]]_{n1}r}}{_{\gamma >0}(\stackrel{~}{\gamma }^{1/2}\stackrel{~}{\gamma }^{1/2})^{2g2}(Z_1^r1)\mathrm{}(Z_{n1}^r1)}).$$
Here, we have introduced $`\stackrel{~}{\gamma }`$ defined (for the root $`\gamma =u_r+u_{r+1}+\mathrm{}+u_{s1}`$, where the $`u_j𝐭^{}`$ are identified via the usual inner product with the $`\widehat{u}_j𝐭`$) by
$$\stackrel{~}{\gamma }(Z_1,\mathrm{},Z_{n1})=Z_r\mathrm{}Z_{s1}.$$
(4.6)
We also have
###### Lemma 4.9
$$_{T^2}e^\omega =n$$
and hence
$$_{T^{2g}}e^{r\omega }=r^{(n1)g}n^g.$$
Proof: This follows from the calculation given in Lemma 10.10 of .
The following may be proved by the same method as in Section 2 of (see also ):
###### Proposition 4.10
Suppose $`\nu `$ $`𝐭`$ is of the form $`\nu =_{j=1}^{n1}\nu _j\widehat{u}_j`$ with $`0\nu _j<1`$ for all $`j`$ (in other words $`\nu `$ is in the interior of the fundamental alcove). Define the meromorphic function $`f`$ on the complexification $`T^𝐂`$ of $`T`$ as follows:
$$f(Z)=(1)^{n1}(1)^{n_+(g1)}r^{(n1)(g1)}n^{g1}\frac{Z_1^{\nu _1r}\mathrm{}Z_{n1}^{\nu _{n1}r}}{_{\gamma >0}(\stackrel{~}{\gamma }^{1/2}\stackrel{~}{\gamma }^{1/2})^{2g2}}.$$
(4.7)
Then we have that
$$\frac{1}{(n1)!}\mathrm{Res}_{Z_1=1}\mathrm{}\mathrm{Res}_{Z_{n1}=1}\underset{wW_{n1}}{}\underset{j=1}{\overset{n1}{}}\left(\frac{r}{Z_j}\right)\frac{[[wf]](Z)}{_{j=1}^{n1}(Z_j^r1)}$$
(4.8)
$$=\underset{\mu \mathrm{\Delta }(r)}{}f(\mathrm{exp}2\pi i\mu /r).$$
Here, $`W_{n1}`$ is the permutation group on $`\{1,\mathrm{},n1\}`$ which is (isomorphic to) the Weyl group of $`SU(n1)`$, and $`[[wf]]`$ is the function
$$[[wf]](Z)=(1)^{n1}(1)^{n_+(g1)}r^{(n1)(g1)}n^{g1}\frac{Z_1^{[[w\nu ]]_1r}\mathrm{}Z_{n1}^{[[w\nu ]]_{n1}r}}{_{\gamma >0}(\stackrel{~}{\gamma }^{1/2}\stackrel{~}{\gamma }^{1/2})^{2g2}}.$$
(4.9)
For a root $`\gamma `$, the quantity $`\stackrel{~}{\gamma }`$ was defined by (4.6).
Remark: Notice that we have
$$\underset{\lambda \mathrm{\Delta }(r)}{}f(\mathrm{exp}2\pi i\lambda /r)=\frac{1}{n1}\underset{m_j=1}{\overset{r1}{}}f\left(e^{2\pi i(_jm_jw_j)/r}\right).$$
(Here, the $`w_j`$ are the fundamental weights, which are dual to the simple roots $`\{\widehat{u}_j\}`$.) The set $`\{X𝐭:`$ $`X=_j\lambda _j\widehat{u}_j,0\lambda _j<1,j=1,\mathrm{},n1\}`$ is a fundamental domain for the action of the integer lattice $`\mathrm{\Lambda }^I`$ on $`𝐭`$, while the set $`\{X𝐭_+𝐭:\gamma _{\mathrm{max}}(X)<1\}`$ is a fundamental domain for the affine Weyl group $`W_{\mathrm{aff}}`$ (the semidirect product of the Weyl group and the integer lattice), and $`\mathrm{\Lambda }^I`$ has index $`(n1)!`$ (rather than $`n!`$) in $`W_{\mathrm{aff}}`$ (in other words a fundamental domain for $`\mathrm{\Lambda }^I`$ contains $`(n1)!`$ fundamental domains for $`W_{\mathrm{aff}}`$).
Applying Proposition 4.10 we find (noting that $`(1)^{n1}=c^\rho `$ when $`n`$ and $`d`$ are coprime) that
$$D_{n,d}(g,k)=(1)^{n_+(g1)}r^{(n1)(g1)}n^{g1}c^\rho \underset{\lambda \mathrm{\Delta }(r)}{}\frac{e^{2\pi i\stackrel{~}{c},\lambda }}{_{\gamma >0}(e^{2\pi i\frac{\gamma }{2r},\lambda }e^{2\pi i\frac{\gamma }{2r},\lambda })^{2g2}}.$$
(4.10)
This gives
$$D_{n,d}(g,k)=(1)^{n_+(g1)}r^{(n1)(g1)}n^{g1}\underset{\lambda \mathrm{\Delta }(r)}{}\frac{e^{2\pi i\stackrel{~}{c},\lambda +\rho }}{_{\gamma >0}(2i\mathrm{sin}\pi \gamma ,\lambda /r)^{2g2}}$$
(4.11)
$$=r^{(n1)(g1)}n^{g1}\underset{\lambda \mathrm{\Delta }(r)}{}\frac{e^{2\pi i\stackrel{~}{c},\lambda +\rho }}{_{\gamma >0}\left(2\mathrm{sin}\frac{\pi \gamma ,\lambda }{r}\right)^{2g2}}.$$
(4.12)
Comparing with Definition 4.2, we see that $`D_{n,d}(g,k)=V_{n,d}(g,k)`$. This completes the proof of Theorem 4.3.
### 4.2 The Verlinde formula for $`_{g,1}(c,\mathrm{\Lambda })`$
There is a more general version of Verlinde’s conjecture which applies to the case of $`_{g,1}(c,\mathrm{\Lambda })`$. Let $`\mathrm{\Lambda }𝐭`$, and let $`k`$ be a positive integer divisible by $`n`$. If $`k\mathrm{\Lambda }`$ is in the weight lattice $`\mathrm{\Lambda }^W`$, then the cohomology class of $`k\omega `$ is a class in integral cohomology, and hence is the first Chern class of a line bundle over $`_{g,1}(c,\mathrm{\Lambda })`$ (denoted $`^{k/n}`$). Notice that $`kf_2`$ is automatically in integral cohomology, since $`nf_2`$ is in integral cohomology.
###### Definition 4.11
$$D_{n,d}(g,k,\mathrm{\Lambda })=\underset{j0}{}(1)^j\mathrm{dim}H^j(_{g,1}(c,\mathrm{\Lambda }),^{k/n})$$
Note that the argument sketched in Section 11 of does not generalize to show that $`H^j(_{g,1}(c,\mathrm{\Lambda }),^{k/n})=0`$ when $`j>0`$, unlike the situation for $`M(n,d)`$. (However, in Teleman has constructed an alternative argument showing the vanishing of these higher cohomology groups.)
We introduce $`\lambda =k\mathrm{\Lambda }`$.
###### Definition 4.12
The Verlinde function is
$$V_{n,d}(g,k,\mathrm{\Lambda })=\underset{\mu \mathrm{\Delta }(r)}{}\frac{e^{2\pi i\mu \rho ,\stackrel{~}{c}}S_{\lambda \mu }(r)}{(S_{0\mu }(r))^{2g1}}$$
where
$$S_{\lambda \mu }(r)=\frac{(\sqrt{1})^{n(n1)/2}}{\sqrt{n}r^{(n1)/2}}\underset{wW}{}(1)^we^{2\pi iw(\lambda +\rho ),\mu +\rho /r}.$$
(4.13)
###### Theorem 4.13
*\[*Verlinde’s conjecture for parabolic bundles\] There exists a neighbourhood $`U`$ of $`0`$ in $`𝐭`$ such that, for $`\mathrm{\Lambda }U`$ for which $`k\mathrm{\Lambda }\mathrm{\Lambda }^W`$, we have
$$D_{n,d}(g,k,\mathrm{\Lambda })=V_{n,d}(g,k,\mathrm{\Lambda }).$$
The proof of Theorem 4.13 proceeds by a sequence of lemmas. By the Riemann-Roch formula, we have
$$D_{n,d}(g,k,\mathrm{\Lambda })=_{_{g,1}(c,\mathrm{\Lambda })}\mathrm{ch}(^{k/n})\mathrm{Td}(_{g,1}(c,\mathrm{\Lambda }))$$
(4.14)
We now use the fibration from Theorem 3.2 to integrate over the fibre in order to obtain an integral over the base $`M(n,d)`$: the latter is then evaluated using Theorem 2.3. First we observe that because the Todd class is multiplicative, it decomposes as the product of Todd classes corresponding to the fibre and the base:
###### Lemma 4.14
$$\mathrm{Td}(_{g,1}(c,\mathrm{\Lambda }))=\mathrm{Td}(T_{\mathrm{vert}}_{g,1}(c,\mathrm{\Lambda }))\pi ^{}\mathrm{Td}(M(n,d)).$$
Here we have introduced $`T_{\mathrm{vert}}_{g,1}(c,\mathrm{\Lambda })`$, the vertical tangent bundle of the fibre of (3.4), so
$$\pi _{}\mathrm{Td}(_{g,1}(c,\mathrm{\Lambda }))=\mathrm{Td}(M(n,d))\pi _{}\mathrm{Td}(T_{\mathrm{vert}}_{g,1}(c,\mathrm{\Lambda })).$$
Next we recall that since the fibre of (3.4) is just a homogeneous space $`G/T`$, we can express its Todd class in terms of the generators introduced in Proposition 3.5.
###### Lemma 4.15
$$\mathrm{Td}(𝒪_\mathrm{\Lambda })=\left(\underset{\gamma >0}{}\frac{\gamma ()}{1e^{\gamma ()}}\right).$$
Here, $`=(e_1,\mathrm{},e_n)`$ is regarded as a member of $`H^2(𝒪_\mathrm{\Lambda })𝐭`$ so one can naturally pair it with the root $`\gamma 𝐭^{}`$ to obtain an element $`(\gamma ,)=\gamma ((e_1,\mathrm{},e_n))`$ in $`H^2(𝒪_\mathrm{\Lambda },𝐑).`$
Proof: This is proved (for example) in Section 14 of . Here, the $`e_i`$ were introduced in Proposition 3.5.
To obtain the Riemann-Roch number we must compute
$$_{_{g,1}(c,\mathrm{\Lambda })}\mathrm{ch}(^{k/n})\mathrm{Td}(_{g,1}(c,\mathrm{\Lambda })).$$
(4.15)
Using the fact that
$$\mathrm{ch}(^{k/n})=\mathrm{exp}(k\omega _\mathrm{\Lambda })$$
and the decomposition of $`\omega _\mathrm{\Lambda }`$ given by (3.5), we find that we must compute
$$_{_{g,1}(c,\mathrm{\Lambda })}\left(e^{k\pi ^{}\omega _{n,d}}\pi ^{}\mathrm{Td}(M(n,d))\right)\left(e^{k\stackrel{~}{\mathrm{\Omega }}_\mathrm{\Lambda }}\mathrm{Td}(T_{\mathrm{vert}}_{g,1}(c,\mathrm{\Lambda }))\right)$$
(4.16)
where $`T_{\mathrm{vert}}_{g,1}(c,\mathrm{\Lambda })`$ is the vertical tangent bundle of the fibres.
We shall integrate over the fibre appearing in (3.4): on this fibre we see that
$$e^{k\stackrel{~}{\mathrm{\Omega }}_\mathrm{\Lambda }}\mathrm{Td}(T_{\mathrm{vert}}_{g,1}(c,\mathrm{\Lambda }))$$
becomes
$$\mathrm{exp}\left(\underset{j=1}{\overset{n}{}}\lambda _je_j\right)\underset{\gamma >0}{}\frac{(\gamma ,)e^{(\gamma ,)/2}}{e^{(\gamma ,)/2}e^{(\gamma ,)/2}}$$
(4.17)
Recall that we have introduced
$$=(e_1,\mathrm{},e_n).$$
To do the integral we reorganize it in terms of Weyl invariant and Weyl anti-invariant cohomology classes. We replace the expression in (4.17) by
$$𝒯()\underset{\gamma >0}{}(\gamma ,),$$
(4.18)
where
$$𝒯()=\frac{1}{|W|}\underset{vW}{}(1)^v\frac{e^{v(\lambda +\rho ),}}{_{\gamma >0}e^{(\gamma ,)/2}e^{(\gamma ,)/2}.}$$
Recall that $`\rho `$ is half the sum of the positive roots.
Note that the quantity $`𝒯()`$ is invariant under the transformation $`w`$. Hence by Propositions 3.4 and 3.6, $`𝒯()`$ $`=\pi ^{}𝒮_\lambda (a_2,\mathrm{},a_n)`$ for an appropriate polynomial $`𝒮_\lambda (a_2,\mathrm{},a_n)`$ in the $`a_2,\mathrm{},a_n`$. Here, $`𝒮_\lambda `$ is defined by
$$𝒮_\lambda (\kappa (\tau _2),\mathrm{},\kappa (\tau _n))=\kappa \left(\frac{_{vW}(1)^ve^{(v(\lambda +\rho ),X)}}{_{\gamma >0}e^{(\gamma ,X)/2}e^{(\gamma ,X)/2}}\right).$$
(4.19)
The $`a_j`$ satisfy $`a_j=\kappa (\tau _j)`$ where $`\tau _j`$ is the $`j`$-th elementary symmetric polynomial ($`\tau _jS(𝐭^{})^W`$) and $`\kappa `$ is the Kirwan map. The integral over the fibre reduces to
$$_{M(n,d)}𝒮_\lambda (\kappa (\tau _2),\mathrm{},\kappa (\tau _n))e^{k\omega _{n,d}}\mathrm{Td}(M(n,d))_{𝒪_\mathrm{\Lambda }}\underset{\gamma >0}{}(\gamma ,)$$
(4.20)
where $`𝒮_\lambda `$ is defined in (4.19). Now the integral
$$_{𝒪_\mathrm{\Lambda }}\underset{\gamma >0}{}(\gamma ,)$$
gives the Euler characteristic of the orbit, which is simply $`|W|`$ $`=n!`$ (, Chap. 14). Thus the integral over the total space $`_{g,1}(c,\mathrm{\Lambda })`$ is reduced to an integral over the base space $`M(n,d)`$: it becomes
$$_{M(n,d)}𝒮_\lambda (\kappa (\tau _2),\mathrm{}\kappa (\tau _n))e^{k\omega _{n,d}}\mathrm{Td}(M(n,d)).$$
(4.21)
(Notice that the argument of $`\kappa `$ in (4.19) is the quotient of two Weyl anti-invariant functions of the variable $`X`$ and hence is Weyl invariant.)
A generalization of Theorem 2.3 (proved by an extension of the proof of Theorem 2.3 given in ) is
###### Theorem 4.16
In the notation from the statement of Theorem 2.3, let
$$\alpha =\kappa \left(\underset{vW}{}(1)^v\frac{e^{v\nu ,X}}{_{\gamma >0}(e^{\gamma ,X}e^{\gamma ,X})}\right)$$
be a formal cohomology class in $`H^{}(M(n,d))`$, where $`\nu 𝐭^{}`$. (Note that the argument of $`\kappa `$ is a Weyl invariant function of $`X`$.) Then $`\kappa (\eta )\alpha \mathrm{exp}(f_2)[M(n,d)]`$ is given by
$$_{M(n,d)}\kappa (\eta )\alpha \mathrm{exp}(f_2)=\frac{(1)^{n_+(g1)}}{n!}\mathrm{Res}_{Y_1=0}\mathrm{}\mathrm{Res}_{Y_{n1}=0}$$
$$\left(\frac{_{wW_{n1}}_{vW}(1)^ve^{[[w(\stackrel{~}{c}+v\nu )]],X}_{T^{2g}}\eta e^\omega }{𝒟_n^{2g2}_{\gamma >0}(e^{\gamma ,X/2}e^{\gamma ,X/2})_{1jn1}(\mathrm{exp}(Y_j)1)}\right),$$
Recall that the notation $`[[\gamma ]]`$ means the unique element which is in the fundamental domain defined by the simple roots for the translation action on $`𝐭`$ of the integer lattice and for which $`[[\gamma ]]`$ is equal to $`\gamma `$ plus some element of the integer lattice.
The expression in (4.21) is the type of integral computed by Theorem 4.16. A straightforward modification of the proof of Proposition 4.7 gives that the Riemann-Roch number is
$$V_{n,d}(g,k,\mathrm{\Lambda })=_{M(n,d)}e^{(k+n)f_2}\kappa \left(\underset{\gamma >0}{}\frac{\gamma (X)^{2g2}}{\left(e^{\gamma (X)/2}e^{\gamma (X)/2}\right)^{2g1}}\underset{vW}{}(1)^ve^{(v(\lambda +\rho ),X)}\right).$$
(4.22)
Using Theorem 4.16 this leads to (in a manner similar to the proof of Theorem 4.8)
$$V_{n,d}(g,k,\mathrm{\Lambda })=\frac{(1)^{n_+(g1)}}{n!}\underset{wW_{n1}}{}\mathrm{Res}_{Y_1=0}\mathrm{}\mathrm{Res}_{Y_{n1}=0}_{T^{2g}}e^{r\omega }\times $$
(4.23)
$$\underset{\gamma >0}{}(\frac{(\gamma (X))^{2g2}}{(e^{\gamma (X)/2}e^{\gamma (X)/2})^{2g1}}\frac{_{vW}(1)^ve^{r[[w(\stackrel{~}{c}+v(\lambda +\rho )/r)]],X}}{_{j=1}^l(e^{rY_j}1)𝒟(X)^{2g2}}.$$
This yields in turn (just as in (4.5))
$$V_{n,d}(g,k,\mathrm{\Lambda })=\frac{(1)^{n1+n_+(g1)}}{n!}\underset{wW_{n1}}{}\underset{vW}{}(1)^v\mathrm{Res}_{Z_1=1}\mathrm{}\mathrm{Res}_{Z_{n1}=1}\left(\underset{j=1}{\overset{n1}{}}\frac{1}{Z_j}\right)\times $$
(4.24)
$$_{T^{2g}}e^{r\omega }\frac{Z_1^{r[[w(\stackrel{~}{c}+v(\lambda +\rho )/r)]]_1}Z_2^{r[[w(\stackrel{~}{c}+v(\lambda +\rho )/r)]]_2}\mathrm{}Z_{n1}^{r[[w(\stackrel{~}{c}+v(\lambda +\rho )/r)]]_{n1}}}{_{\gamma >0}(\stackrel{~}{\gamma }^{1/2}\stackrel{~}{\gamma }^{1/2})^{2g1}(Z_1^r1)\mathrm{}(Z_{n1}^r1)}$$
Finally we may use Proposition 4.10 to obtain an expression equivalent to (4.24):
###### Theorem 4.17
We have
$$V_{n,d}(g,k,\mathrm{\Lambda })=\underset{\mu \mathrm{\Delta }(r)}{}\frac{e^{2\pi i\mu +\rho ,\stackrel{~}{c}}S_{\lambda \mu }(r)}{S_{0\mu }^{2g1}}$$
where
$$S_{\lambda \mu }(r)=\frac{(\sqrt{1})^{n(n1)/2}}{\sqrt{n}r^{(n1)/2}}\underset{vW}{}(1)^ve^{2\pi iv(\lambda +\rho ),\mu +\rho /r}.$$
(4.25)
Since the formula given in Theorem 4.17 is the formula for $`D_{n,d}(g,k,\mathrm{\Lambda })`$, we have proved Theorem 4.13. The proof of Theorem 4.17 is valid when $`\mathrm{\Lambda }U`$.
### 4.3 The Verlinde formula for $`_{g,b}(c,𝚲)`$
More generally, there are spaces $`_{g,b}(c,𝚲)`$ which fibre over $`M(n,d)`$ with fibre $`𝒪_{\mathrm{\Lambda }^{(1)}}\times \mathrm{}\times 𝒪_{\mathrm{\Lambda }^{(b)}}`$ for $`\mathrm{\Lambda }^{(j)}`$ sufficiently close to $`0`$. We introduce the notation $`𝚲=(\mathrm{\Lambda }^{(1)},\mathrm{},\mathrm{\Lambda }^{(b)})`$. We restrict to those $`\mathrm{\Lambda }^{(j)}`$ for which $`k\mathrm{\Lambda }^{(j)}\mathrm{\Lambda }^W`$; we denote $`k\mathrm{\Lambda }^{(j)}`$ by $`\lambda ^{(j)}`$.
###### Proposition 4.18
There is a neighbourhood $`𝐔`$ of the element $`(0,\mathrm{},0)`$ $`𝐭^b`$ such that if $`𝚲𝐔`$ then there is a fibration
$$\pi :_{g,b}(c,𝚲)M(n,d)$$
with fibre $`𝒪_{\mathrm{\Lambda }^{(1)}}\times \mathrm{}\times 𝒪_{\mathrm{\Lambda }^{(b)}}`$. Further, the symplectic form on $`_{g,b}(c,𝚲)`$ satisfies
$$\omega _𝚲=\pi ^{}\omega _{n,d}+\stackrel{~}{\mathrm{\Omega }}_𝚲$$
(4.26)
where $`\stackrel{~}{\mathrm{\Omega }}_𝚲`$ restricts on each fibre to the sum of the Kirillov-Kostant symplectic forms
$$\underset{j=1}{\overset{b}{}}\mathrm{\Omega }_{\mathrm{\Lambda }^{(j)}}$$
on the product of the coadjoint orbits $`𝒪_{\mathrm{\Lambda }^{(j)}}.`$
Proof: The proof is as in Theorem 3.2. We use the fact that $`M(n,d)`$ is the reduced space at the regular value $`\mathrm{𝟎}`$ of a symplectic manifold equipped with a Hamiltonian action of $`G^b`$, and that at orbits of the $`G^b`$ action close to $`0`$, the corresponding reduced spaces are the $`_{g,b}(c,𝚲)`$. This extends the argument of Proposition 5.5 of , using the fact that the action of $`G^b`$ is free on the zero locus of the moment map: this action is given in (5.9) and (5.10) of , and it is straightforward to verify that the action is free when the moment map takes the value $`\{\mathrm{𝟎}\}.`$
###### Remark 4.19
The region $`𝐔`$ is characterized as the component whose closure contains $`(0,\mathrm{},0)`$ in the complement in $`𝐭_{+}^{}{}_{}{}^{b}`$ of the union of the walls
$$\{(\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_b):\underset{(w_1,\mathrm{},w_b)W^b}{}w_j\mathrm{\Lambda }_jH_{v,n}\}.$$
Here the walls $`H_{v,n}`$ were introduced in (3.3). This observation is due to A. Szenes .
###### Proposition 4.20
We have
$$\pi ^{}\left[U|_{M(n,d)\times \{\mathrm{pt}\}}\right]=_{j=1}^n_{s=1}^bL_j^{(s)},$$
where $`c_1(L_j^{(s)})=e_j^{(s)}`$ restricts on the $`s`$-th orbit $`𝒪_{\mathrm{\Lambda }^{(s)}}`$ in the fibre of $`\pi `$ to the generator $`\alpha _j^{(s)}`$ of $`H^2(𝒪_{\mathrm{\Lambda }^{(s)}})`$ corresponding to the $`j`$-th fundamental weight of $`SU(n)`$.
As in Lemma 4.14, we have
$$\mathrm{Td}(_{g,b}(c,𝚲))=\mathrm{Td}(T_{\mathrm{vert}}_{g,b}(c,𝚲))\pi ^{}\mathrm{Td}(M(n,d)),$$
where $`\mathrm{Td}(T_{\mathrm{vert}}_{g,b}(c,𝚲))`$ is the vertical tangent bundle of $`_{g,b}(c,𝚲)`$. The integral over the fibres decomposes as a product (as in (4.17)) of factors of the form
$$_{𝒪_{\mathrm{\Lambda }^{(s)}}}\mathrm{exp}\left(\underset{j=1}{\overset{b}{}}\lambda _j^{(s)}e_j^{(s)}\right)\underset{\gamma >0}{}\frac{(\gamma ,^{(s)})e^{(\gamma ,^{(s)})/2}}{e^{(\gamma ,^{(s)})/2}e^{(\gamma ,^{(s)})/2}}$$
(4.27)
where $`^{(s)}=(e_1^{(s)},\mathrm{},e_n^{(s)}).`$ (Here we have introduced a basis $`\{e_j^{(s)}\}`$ for $`H^2(𝒪_{\mathrm{\Lambda }^{(s)}})`$, following Propositions 3.4 and 3.5.) Each of the integrals (4.27) is of the form
$$𝒮_{\lambda ^{(s)}}(\kappa (\tau _2),\mathrm{},\kappa (\tau _n))$$
where $`𝒮_\lambda `$ was defined in (4.19). From the product of these integrals over the $`𝒪_{\mathrm{\Lambda }^{(s)}}`$, we obtain
$$V_{n,d}(g,k,𝚲)=_{M(n,d)}e^{rf_2}\times $$
(4.28)
$$\times \kappa \left(\underset{\gamma >0}{}\frac{\gamma (X)^{2g2}}{\left(e^{\gamma (X)/2}e^{\gamma (X)/2}\right)^{2g2+b}}\underset{s=1}{\overset{b}{}}\underset{v_sW}{}(1)^{v_s}e^{(v_s(\lambda ^{(s)}+\rho ),X)}\right).$$
Using Theorem 4.16, (4.28) leads to
$$V_{n,d}(g,k,𝚲)=\frac{(1)^{n_+(g1)}}{n!}\underset{wW_{n1}}{}\mathrm{Res}_{Y_1=0}\mathrm{}\mathrm{Res}_{Y_{n1}=0}_{T^{2g}}e^{r\omega }\times $$
(4.29)
$$\times \underset{\gamma >0}{}\left(\frac{\gamma (X)^{2g2}}{(e^{\gamma (X)/2}e^{\gamma (X)/2})^{2g2+b}}\right)\underset{s=1}{\overset{b}{}}\underset{v_sW}{}(1)^{v_s}\frac{\mathrm{exp}(r[[w\stackrel{~}{c}+_s(v_s(\lambda ^{(s)}+\rho )/r]],X)}{_{j=1}^l(e^{rY_j}1)𝒟(X)^{2g2}}.$$
This leads in turn (just as in (4.5) and (4.24)) to
$$V_{n,d}(g,k,𝚲)=\frac{(1)^{n1+n_+(g1)}}{n!}\underset{wW_{n1}}{}\underset{v_1W}{}\mathrm{}\underset{v_bW}{}(1)^{v_1}\mathrm{}(1)^{v_b}\times $$
(4.30)
$$\times \mathrm{Res}_{Z_1=1}\mathrm{}\mathrm{Res}_{Z_{n1}=1}\left(\underset{j=1}{\overset{n1}{}}\frac{1}{Z_1}\times \mathrm{}\times \frac{1}{Z_{n1}}\right)_{T^{2g}}e^{r\omega }$$
$$\times \frac{Z_1^{r[[w(\stackrel{~}{c}+_{s=1}^bv_s(\lambda ^{(s)}+\rho )/r)]]_1}Z_2^{r[[w(\stackrel{~}{c}+_{s=1}^bv_s(\lambda ^{(s)}+\rho )/r)]]_2}\mathrm{}Z_{n1}^{r[[w(\stackrel{~}{c}+_{s=1}^bv_s(\lambda ^{(s)}+\rho )/r)]]_{n1}}}{_{\gamma >0}(\stackrel{~}{\gamma }^{1/2}\stackrel{~}{\gamma }^{1/2})^{2g2+b}(Z_1^r1)\mathrm{}(Z_{n1}^r1)}$$
Finally we may use Proposition 4.10 to recover
###### Theorem 4.21
The Riemann-Roch number of $`_{g,b}(c,𝚲)`$ is given by
$$_{_{g,b}(c,𝚲)}\mathrm{ch}(^{k/n})\mathrm{Td}(_{g,b}(c,𝚲))=$$
$$=\underset{\mu }{}\frac{e^{2\pi i\mu \rho ,\stackrel{~}{c}}_{s=1}^bS_{\mu \lambda ^{(s)}}(r)}{(S_{0\mu }(r))^{2g2+b}}.$$
The proof of Theorem 4.21 is valid when $`𝚲𝐔`$.
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# Lectures on Strings, D-branes and Gauge TheoriesLectures delivered by the first author in the Third Workshop on Gravitation and Mathematical Physics, Nov. 28-Dec. 3 1999, León Gto. México.
## I Introduction
String theory is by now, beyond the standard model of particle physics, the best and the most sensible understanding of all the matter and their interactions in an unified scheme. There are well known the ‘esthetical’ problems arising in the heart of the standard model of particles, such as: the abundance of free parameters, the origin of flavor and the gauge group, etc. It is also generally accepted that these problems require to be answered. Thus the standard model can be seen as the low energy effective theory of a more fundamental theory which can solve the mentioned problems. It is also clear that the quantum mechanics and general relativity cannot be reconciliated in the context of a perturbative quantum field theory of point particles. Hence the nonrenormalizability of the general relativity can be seen as a genuine evidence that it is just an effective field theory and new physics associated to the fast degrees of freedom should exist at higher energies. String theory propose that these fast degrees of freedom are precisely the strings at the perturbative level and at the non-perturbative level the relevant degrees of freedom are, in addition to the strings, higher-dimensional extended objects called D-branes (dual degrees of freedom). The study of theories involving D-branes is just in the starting stage and many surprises surely are coming up. Thus we are still at an exploratory stage of the whole structure of the string theory. Therefore the theory is far to be completed and we cannot give yet concrete physical predictions to take contact with collider experiments and/or astrophysical observations. However many aspects of theoretic character, necessary in order to make of string theory a physical theory, are quickly in progress. The purpose of these lectures are to overview the basic ideas to understand these progresses. This paper is an extended version of the lectures presented at the Third Workshop on Gravitation and Mathematical Physics at León Gto. México. We don’t pretend to be exhaustive and we will limit ourselves to describe some basic elements of string theory and some particular new developments as: non-BPS branes and noncommutative gauge theories. We apologize for omiting numerous original references and we prefer to cite review articles and some few seminal papers.
In Sec. II we overview the string and the superstring theories from the perturbative point of view. T-duality and $`D`$-branes is considered in Sec. III. The Sec. IV is devoted to describe the string dualities and the web of string theories connected by duality. M and F theories are also briefly described. Sec. V is devoted to review the non-BPS branes and their description in terms of topological K-Theory. Some recent results by Witten and Moore-Witten, about the classification of Ramond-Ramond fields is also described. The relation of string theory to noncommutative Yang-Mills theory and deformation quantization theory is the theme of the Sec. VI.
## II Perturbative String and Superstring Theories
In this section we overview some basic aspects of bosonics and fermionic strings. We focus mainly in the description of the spectrum of the theory in the light-cone gauge and the brief description of spectrum of the five consistent superstring theories (for details and further developments see for instance ).
First of all consider, as usual, the action of a relativistic point particle. It is given by $`S=m𝑑\tau \sqrt{\dot{X}^\mu \dot{X}_\mu }`$, where $`X^\mu `$ are $`D`$ functions representing the coordinates of the $`(D1,1)`$-dimensional Minkowski spacetime (the target space), $`\dot{X}^\mu \frac{dX^\mu }{d\tau }`$ and $`m`$ can be identified with the mass of the point particle. This action is proportional to the length of the world-line of the relativistic particle. In analogy with the relativistic point particle, the action describing the dynamics of a string (one-dimensional object) moving in a $`(D1,1)`$-dimensional Minkowski spacetime (the target space) is proportional to the area $`𝐀`$ of the worldsheet. We know from the theory of surfaces that such an area is given by $`𝐀=\sqrt{det(g)}`$, where $`g`$ is the induced metric (with signature $`(,+)`$) on the worldsheet. The background metric will be denoted by $`\eta _{\mu \nu }`$ and $`\sigma ^a=(\tau ,\sigma )`$ with $`a=0,1`$ are the local coordinates on the worldsheet. $`\eta _{\mu \nu }`$ and $`g_{ab}`$ are related by $`g_{ab}=\eta _{\mu \nu }_aX^\mu _bX^\nu `$ with $`\mu =0,1,\mathrm{},D1`$. Thus the classical action of a relativistic string is given by the Nambu-Goto action
$$S_{NG}[X^\mu ]=\frac{1}{2\pi \alpha ^{}}𝑑\tau 𝑑\sigma \sqrt{det(_aX^\mu _bX^\nu )},$$
(1)
where $`X^\mu `$ are $`D`$ embedding functions of the worldsheet into the target space $`X`$. Now introduce a metric $`h`$ describing the worldsheet geometry, we get a classically equivalent action to the Nambu-Goto action. This is the Polyakov action
$$S_P[X^\mu ,h_{ab}]=\frac{1}{4\pi \alpha ^{}}d^2\sigma \sqrt{h}h^{ab}_aX^\mu _bX^\nu \eta _{\mu \nu },$$
(2)
where the $`X^\mu `$’s are $`D`$ scalar fields on the worldsheet. Such a fields can be interpreted as the coordinates of spacetime $`X`$ (target space), $`h`$ = det$`(h^{ab})`$ and $`h_{ab}=_aX^\mu _bX^\nu \eta _{\mu \nu }`$.
Polyakov action has the following symmetries: $`(i)`$ Poincaré invariance, $`(ii)`$ Worldsheet diffeomorphism invariance, and $`(iii)`$ Weyl invariance (rescaling invariance). The energy-momentum tensor of the two-dimensional theory is given by
$$T^{ab}:=\frac{1}{\sqrt{h}}\frac{\delta S_P}{\delta h_{ab}}=\frac{1}{4\pi \alpha ^{}}\left(^aX^\mu ^bX_\mu \frac{1}{2}h^{ab}h^{cd}_cX^\mu _dX_\mu \right)=0.$$
(3)
Invariance under worldsheet diffeomorphisms implies that it should be conserved i.e. $`_aT^{ab}=0`$, while the Weyl invariance gives the traceless condition, $`T_a^a=0`$. The equation of motion associated with Polyakov action is given by
$$_a\left(\sqrt{h}h^{ab}_bX^\mu \right)=0.$$
(4)
Whose solutions should satisfy the boundary conditions for the open string: $`_\sigma X^\mu _0^{\mathrm{}=\pi }=0`$ (Neumann) and for the closed string: $`X^\mu (\tau ,\sigma )=X^\mu (\tau ,\sigma +2\pi )`$ (Dirichlet). Here $`\mathrm{}=\pi `$ is the characteristic length of the open string. The variation of $`S_P`$ with respect to $`h^{ab}`$ leads to the constraint equations $`T_{ab}=0`$. From now on we will work in the conformal gauge. In this gauge: $`h_{ab}=\eta _{ab}`$ and equations of motion (4) reduces to the Laplace equation whose solutions can be written as linear superposition of plane waves.
2.1 The Closed String
For the closed string the boundary condition $`X^\mu (\tau ,\sigma )=X^\mu (\tau ,\sigma +2\pi )`$, leads to the general solution of Eq. (4)
$$X^\mu =X_0^\mu +\frac{1}{\pi T}P^\mu \tau +\frac{i}{2\sqrt{\pi T}}\underset{n0}{}\frac{1}{n}\left\{\alpha _n^\mu exp\left(i2n(\tau \sigma )\right)+\stackrel{~}{\alpha }_n^\mu exp\left(i2n(\tau +\sigma )\right)\right\}$$
(5)
where $`X_0^\mu `$ and $`P^\mu `$ are the position and momentum of the center-of-mass of the string and $`\alpha _n^\mu `$ and $`\stackrel{~}{\alpha }_n^\mu `$ satisfy the conditions $`\alpha _n^\mu =\alpha _n^\mu `$ (left-movers) and $`\stackrel{~}{\alpha }_n^\mu =\stackrel{~}{\alpha }_n^\mu `$ (right-movers).
2.2 The Open String
For the open string the respective boundary condition is $`_\sigma X^\mu _0^{\mathrm{}=\pi }=0`$ (this is the only boundary condition which is Lorentz invariant) and the solution is given by
$$X^\mu (\tau ,\sigma )=X_0^\mu +\frac{1}{\pi T}P^\mu \tau +\frac{i}{\sqrt{\pi T}}\underset{n0}{}\frac{1}{n}\alpha _n^\mu exp\left(in\tau \right)\mathrm{cos}(n\sigma )$$
(6)
with the condition $`\alpha _n^\mu =\stackrel{~}{\alpha }_n^\mu .`$
2.3 Quantization
The quantization of the closed bosonic string can be carried over, as usual, by using the Dirac prescription to the center-of-mass and oscillator variables in the form
$$[X_0^\mu ,P^\nu ]=i\eta ^{\mu \nu },$$
$$[\alpha _m^\mu ,\alpha _n^\nu ]=[\stackrel{~}{\alpha }_m^\mu ,\stackrel{~}{\alpha }_n^\nu ]=m\delta _{m+n,0}\eta ^{\mu \nu },$$
$$[\alpha _m^\mu ,\stackrel{~}{\alpha }_n^\nu ]=0.$$
(7)
One can identify $`(\alpha _n^\mu ,\stackrel{~}{\alpha }_n^\mu )`$ with the creation operators and the corresponding operators $`(\alpha _n^\mu ,\stackrel{~}{\alpha }_n^\mu )`$ with the annihilation ones. In order to specify the physical states we first denote the center of mass state given by $`|0;P^\mu `$. The vacuum state is defined by $`\alpha _m^\mu |0,P^\mu =0`$ and $`P_m^\mu |0,P^\mu =p^\mu 0,P^\mu `$ and similar for the right movings. For the zero modes these states have negative norm (ghosts). However one can choice a suitable gauge where ghosts decouple from the Hilbert space when $`D=26`$. This is the subject of the following subsection.
2.4 Light-cone Quantization
Now we turn out to work in the so called light-cone gauge. In this gauge it is possible to solve explicitly the Virasoro constraints (3). This is done by removing the light-cone coordinates $`X^\pm =\frac{1}{\sqrt{2}}(X^0\pm X^D)`$ leaving only the transverse coordinates $`X^i`$ representing the physical degrees of freedom (with $`i=1,2,\mathrm{},D2`$). In this gauge the Virasoro constraints (3) are explicitly solved. Thus the independent variables are $`(X_0^{},P^+,X_0^j,P^j,\alpha _n^j,\stackrel{~}{\alpha }_n^j)`$. Operators $`\alpha _n^{}`$ and $`\stackrel{~}{\alpha }_n^{}`$ can be written in terms of $`\alpha _n^j`$ and $`\stackrel{~}{\alpha }_n^j`$ respectively as follows: $`\alpha _n^{}=\frac{1}{\sqrt{2\alpha ^{}}P^+}(_{m=\mathrm{}}^{\mathrm{}}:\alpha _{nm}^i\alpha _m^i:2A\delta _n)`$ and $`\stackrel{~}{\alpha }_n^{}=\frac{1}{\sqrt{2\alpha ^{}}P^+}(_{m=\mathrm{}}^{\mathrm{}}:\stackrel{~}{\alpha }_{nm}^i\stackrel{~}{\alpha }_m^i:2A\delta _n`$). For the open string we get $`\alpha _n^{}=\frac{1}{2\sqrt{2\alpha ^{}}P^+}(_{m=\mathrm{}}^{\mathrm{}}:\alpha _{nm}^i\alpha _m^i:2A\delta _n)`$. Here $`::`$ stands for the normal ordering.
In this gauge the Hamiltonian is given by
$$H=\frac{1}{2}(P^i)^2+NA(\mathrm{open}\mathrm{string}),H=(P^i)^2+N_L+N_R2A(\mathrm{closed}\mathrm{string})$$
(8)
where $`N`$ is the operator number, $`N_L=_{m=\mathrm{}}^{\mathrm{}}:\alpha _m\alpha _m:`$, and $`N_R=_{m=\mathrm{}}^{\mathrm{}}:\stackrel{~}{\alpha }_m\stackrel{~}{\alpha }_m:.`$ The mass-shell condition is given by $`\alpha ^{}M^2=(NA)`$ (open string) and $`\alpha ^{}M^2=2(N_L+N_R2A)`$ (closed string). For the open string, Lorentz invariance implies that the first excited state is massless and therefore $`A=1`$. In the light-cone gauge $`A`$ takes the form $`A=\frac{D2}{2}_{n=1}^{\mathrm{}}n`$. From the fact $`_{n=1}^{\mathrm{}}n^s=\zeta (s),`$ where $`\zeta `$ is the Riemann’s zeta function (which converges for $`s>1`$ and has a unique analytic continuation at $`s=1`$, where it takes the value $`\frac{1}{12}`$) then $`A=\frac{D2}{24}`$ and therefore $`D=26`$.
2.5 Spectrum of the Bosonic String
Closed Strings
The spectrum of the closed string can be obtained from the combination of the left-moving states and the the right-moving ones. The ground state ($`N_L=N_R=0`$) is given by $`\alpha ^{}M^2=4`$. That means that the ground state includes a tachyon. The first excited state ($`N_L=1=N_R`$) is massless and it is given by $`\alpha _1^i\stackrel{~}{\alpha }_1^j|0,P`$. This state can be naturally decomposed into irreducible representations of the little group $`SO(24)`$ as follows
$$\alpha _1^i\stackrel{~}{\alpha }_1^j0,P=\alpha _1^{[i}\stackrel{~}{\alpha }_1^{j]}0,P+\left(\alpha _1^{(i}\stackrel{~}{\alpha }_1^{j)}\frac{1}{D2}\delta ^{ij}\alpha _1^k\stackrel{~}{\alpha }_1^k\right)0,P$$
$$+\frac{1}{D2}\delta ^{ij}\alpha _1^k\stackrel{~}{\alpha }_1^k0,P.$$
(9)
The first term of the rhs is interpreted as a spin 2 massless particle $`g_{ij}`$ (graviton). The second term is a range 2 anti-symmetric tensor $`B_{ij}`$. While the last term is an scalar field $`\varphi `$ (dilaton). Higher excited massive states are combinations of representations of the little group SO$`(25).`$
Open Strings
For the open string, the ground state includes once again a tachyon since $`\alpha ^{}M^2=1`$. The first exited state $`N=1`$ is given by a massless vector field in 26 dimensions. The second excitation level is given by the massive states $`\alpha _2^i0,P`$ and $`\alpha _1^i\alpha _1^j0,P`$ which are in irreducible representations of the little group SO$`(25)`$.
2.6 Superstrings
In bosonic string theory there are two big problems. The first one is the presence of tachyons in the spectrum. The second one is that there are no spacetime fermions. Here is where superstrings come to the rescue. A superstring is described, despite of the usual bosonic fields $`X^\mu `$, by fermionic fields $`\psi _{L,R}^\mu `$ on the worldsheet. Which satisfy anticommutation rules and where the $`L`$ and $`R`$ denote the left and right worldsheet chirality respectively. The action for the superstring is given by
$$L_{SS}=\frac{1}{8\pi }d^2\sigma \sqrt{h}\left(h^{ab}_aX^\mu _bX_\mu +2i\overline{\psi }^\mu \gamma ^a_a\psi _\mu i\overline{\chi }_a\gamma ^b\gamma ^a\psi ^\mu \left(_bX_\mu \frac{i}{4}\overline{\chi }_b\psi _\mu \right)\right),$$
(10)
where $`\psi ^\mu `$ and $`\chi _a`$ are the superpartners of $`X^\mu `$ and the tetrad field $`e^a`$, respectively. In the superconformal gauge and in light-cone coordinates it can be reduced to
$$L_{SS}=\frac{1}{2\pi }\left(_LX^\mu _RX_\mu +i\psi _R^\mu _L\psi _{\mu R}+i\psi _L^\mu _R\psi _{\mu R}\right).$$
(11)
In analogy to the bosonic case, the local dynamics of the worldsheet metric is manifestly independent of quantum corrections if the critical spacetime dimension $`D`$ is 10. Thus the string oscillates in the 8 transverse dimensions. The action (10) is invariant under: $`(i)`$ worldsheet supersymmetry, $`(ii)`$ Weyl transformations, $`(iii)`$ super-Weyl transformations, $`(iv)`$ Poincaré transformations and $`(v)`$ Worldsheet reparametrizations. The equation of motion for the $`X^{}s`$ fields is the same that in the bosonic case (Laplace equation) and whose general solution is given by Eqs. (5) or (6). Equation of motion for the fermionic field is the Dirac equation in two dimensions. Constraints here are more involved and they are called the super-Virasoro constraints. However in the light-cone gauge, everything simplifies and the transverse coordinates (eight coordinates) become the bosonic physical degrees of freedom together with their corresponding supersymmetric partners. Analogously to the bosonic case, massless states of the spectrum come into representations of the little group SO(8) of SO$`(9,1)`$, while that the massive states lie into representations of the little group SO$`(9)`$.
For the closed string there are two possibilities for the boundary conditions of fermions: $`(i)`$ periodic boundary conditions (Ramond (R) sector) $`\psi _{L,R}^\mu (\sigma )=+\psi _{L,R}^\mu (\sigma +2\pi )`$ and $`(ii)`$ anti-periodic boundary conditions (Neveu-Schwarz (NS) sector) $`\psi _{L,R}^\mu (\sigma )=\psi _{L,R}^\mu (\sigma +2\pi )`$. Solutions of Dirac equation satisfying these boundary conditions are
$$\psi _L^\mu (\sigma ,\tau )=\underset{n}{}\overline{\psi }_n^\mu exp\left(in(\tau +\sigma )\right),\psi _R^\mu (\sigma ,\tau )=\underset{n}{}\psi _n^\mu exp\left(in(\tau \sigma )\right).$$
(12)
In the case of the fermions in the R sector $`n`$ is integer and it is semi-integer in the NS sector.
The quantization of the superstring come from the promotion of the fields $`X^\mu `$ and $`\psi ^\mu `$ to operators whose oscillator variables are operators satisfying the relations $`[\alpha _n^\mu ,\alpha _m^\nu ]_{}=n\delta _{m+n,0}\eta ^{\mu \nu }`$ and $`[\psi _n^\mu ,\psi _m^\nu ]_+=\eta ^{\mu \nu }\delta _{m+n,0},`$ where $`[,]_{}`$ and $`[,]_+`$ stand for commutator and anti-commutator respectively.
The zero modes of $`\alpha `$ are diagonal in the Fock space and its eigenvalue can be identified with its momentum. For the NS sector there is no fermionic zero modes but they can exist for the R sector and they satisfy a Clifford algebra $`[\psi _0^\mu ,\psi _0^\nu ]_+=\eta ^{\mu \nu }`$. The Hamiltonian for the closed superstring is given by $`H_{L,R}=N_{L,R}+\frac{1}{2}P_{L,R}^2A_{L,R}`$. For the NS sector $`A=\frac{1}{2}`$, while for the R sector $`A=0`$. The mass is given by $`M^2=M_L^2+M_R^2`$ with $`\frac{1}{2}M_{L,R}^2=N_{L,R}A_{L,R}`$.
There are five consistent superstring theories: Type IIA, IIB, Type I, SO(32) and $`E_8\times E_8`$ heterotic strings. In what follows of this section we briefly describe the spectrum in each one of them.
2.7 Type II Superstring Theories
In this case the theory consist of closed strings only. They are theories with $`𝒩=2`$ spacetime supersymmetry. For this reason, there are 8 scalar fields (representing the 8 transverse coordinates to the string) and one Weyl-Majorana spinor. There are 8 left-moving and 8 right-moving fermions.
In the NS sector there is still a tachyon in the ground state. But in the supersymmetric case this problem can be solved through the introduction of the called GSO projection. This projection eliminates the tachyon in the NS sector and it acts in the R sector as a ten-dimensional spacetime chirality operator. That means that the application of the GSO projection operator defines the chirality of a massless spinor in the R sector. Thus from the left and right moving sectors, one can construct states in four different sectors: $`(i)`$ NS-NS, $`(ii)`$ NS-R, $`(iii)`$ R-NS and $`(iv)`$ R-R. Taking account the two types of chirality $`L`$ and $`R`$ one has two possibilities:
$`a)`$ The GSO projections on the left and right fermions produce different chirality in the ground state of the R sector (Type IIA).
$`b)`$GSO projection are equal in left and right sectors and the ground states in the R sector, have the same chirality (Type IIB). Thus the spectrum for the Type IIA and IIB superstring theories is:
Type IIA
The NS-NS sector has a symmetric tensor field $`g_{\mu \nu }`$ (spacetime metric), an antisymmetric tensor field $`B_{\mu \nu }`$ and a scalar field $`\varphi `$ (dilaton). In the R-R sector there is a vector field $`A_\mu `$ associated with a 1-form $`C_1`$ ($`A_\mu C_1`$) and a rank 3 totally antisymmetric tensor $`C_{\mu \nu \rho }C_3`$. In general the R-R sector consist of $`p`$-forms $`G_p=dC_{p1}`$ (where $`C_p`$ are called RR fields) on the ten-dimensional spacetime $`X`$ with $`p`$ even i.e. $`G_0`$, $`G_2`$, $`G_4`$, …. In the NS-R and R-NS sectors we have two gravitinos with opposite chirality and the supersymmetric partners of the mentioned bosonic fields.
Type IIB
In the NS-NS sector Type IIB theory has exactly the same spectrum that of Type IIA theory. On the R-R sector it has a scalar field $`\chi C_0`$, an antisymmetric tensor field $`B_{\mu \nu }^{}C_2`$ and a rank 4 totally antisymmetric tensor $`D_{\mu \nu \rho \sigma }C_4`$ whose field strength is self-dual i.e., $`G_5=dC_4`$ with $`G_5=+G_5`$. Similar than for the case of Type IIA theory one has, in general, RR fields given by $`p`$-forms $`G_p=dC_{p1}`$ on the spacetime $`X`$ with $`p`$ odd i.e. $`G_1`$, $`G_3`$, $`G_5`$, …. The NS-R and R-NS sectors do contain gravitinos with the same chirality and the corresponding fermionic matter.
2.8 Type I Superstrings
In this case the $`L`$ and $`R`$ degrees of freedom are the same. Type I and Type IIB theories have the same spectrum, except that in the former one the states which are not invariant under the change of orientation of the worldsheet, are projected out. This worldsheet parity $`\mathrm{\Omega }`$ interchanges the left and right modes. Type I superstring theory is a theory of breakable closed strings, thus it incorporates also open strings. The $`\mathrm{\Omega }`$ operation leave invariant only one half of the spacetime supersymmetry, thus the theory is $`𝒩=1`$.
The spectrum of bosonic massless states in the NS-NS sector is: $`g_{\mu \nu }`$ (spacetime metric) and $`\varphi `$ (dilaton) from the closed sector and $`B_{\mu \nu }`$ is projected out. On the R-R sector there is an antisymmetric field $`B_{\mu \nu }`$ of the closed sector. The open string sector is necessary in order to cancel tadpole diagrams. A contribution to the spectrum come from this sector. Chan-Paton factors can be added at the boundaries of open strings. Hence the cancellation of the tadpole are needed 32 labels at each end. Therefore in the NS-NS sector there are 496 gauge fields in the adjoint representation of SO(32).
2.9 Heterotic Superstrings
This kind of theory involves only closed strings. Thus there are left and right sectors. The left-moving sector contains a bosonic string theory and the right-moving sector contains superstrings. This theory is supersymmetric on the right sector only, thus the theory contains $`𝒩=1`$ spacetime supersymmetry. The momentum at the left sector $`P_L`$ lives in 26 dimensions, while $`P_R`$ lives in 10 dimensions. It is natural to identify the first ten components of $`P_L`$ with $`P_R`$. Consistency of the theory tell us that the extra 16 dimensions should belong to the root lattice $`E_8\times E_8`$ or a $`𝐙_2`$-sublattice of the SO(32) weight lattice.
The spectrum consists of a tachyon in the ground state of the left-moving sector. In both sectors we have the spacetime metric $`g_{\mu \nu }`$, the antisymmetric tensor $`B_{\mu \nu }`$, the dilaton $`\varphi `$ and finally there are 496 gauge fields $`A_\mu `$ in the adjoint representation of the gauge group $`E_8\times E_8`$ or SO$`(32)`$.
## III Toroidal Compactification, $`T`$-duality and D-branes
D-branes are, despite of the dual fundamental degrees of freedom in string theory, extremely interesting and useful tools to study nonperturbative properties of string and field theories (for a classic review see ). Non-perturbative properties of supersymmetric gauge theories can be better understanding as the world-volume effective theory of some configurations of intersecting D-branes (for a review see ). D-branes also are very important to connect gauge theories with gravity. This is the starting point of the AdS/CFT correspondence or Maldacena’s conjecture. We don’t review this interesting subject in this paper, however the reader can consult the excellent review . Roughly speaking D-branes are static solutions of string equations which satisfy Dirichlet boundary conditions. That means that open strings can end on them. To explain these objects we follow the traditional way, by using T-duality on open strings we will see that Neumann conditions are turned out into the Dirichlet ones. To motivate the subject we first consider T-duality in closed bosonic string theory.
3.1 T-duality in Closed Strings
The general solution of Eq. (4) in the conformal gauge can be written as $`X^\mu (\sigma ,\tau )=X_R^\mu (\sigma ^{})+X_L^\mu (\sigma ^+)`$, where $`\sigma ^\pm =\sigma \pm \tau `$. Now, take one coordinate, say $`X^{25}`$ and compactify it on a circle of radius $`R`$. Thus we have that $`X^{25}`$ can be identified with $`X^{25}+2\pi Rm`$ where $`m`$ is called the winding number. The general solution for $`X^{25}`$ with the above compactification condition is
$$X_R^{25}(\sigma ^{})=X_{0R}^{25}+\sqrt{\frac{\alpha ^{}}{2}}P_R^{25}(\tau \sigma )+i\sqrt{\frac{\alpha ^{}}{2}}\underset{l0}{}\frac{1}{l}\alpha _{R,l}^{25}exp\left(il(\tau \sigma )\right)$$
$$X_L^{25}(\sigma ^+)=X_{0L}^{25}+\sqrt{\frac{\alpha ^{}}{2}}P_L^{25}(\tau +\sigma )+i\sqrt{\frac{\alpha ^{}}{2}}\underset{n0}{}\frac{1}{l}\alpha _{L,l}^{25}exp\left(il(\tau +\sigma )\right),$$
(13)
where
$$P_{R,L}^{25}=\frac{1}{\sqrt{2}}\left(\frac{\sqrt{\alpha ^{}}}{R}n\frac{R}{\sqrt{\alpha ^{}}}m\right).$$
(14)
Here $`n`$ and $`m`$ are integers representing the discrete momentum and the winding number, respectively. The latter has not analogous in field theory. While the canonical momentum is given by $`P^{25}=\frac{1}{\sqrt{2\alpha ^{}}}(P_L^{25}+P_R^{25})`$. Now, by the mass shell condition, the mass of the perturbative states is given by $`M^2=M_L^2+M_R^2`$, with
$$M_{L,R}^2=\frac{1}{2}P^\mu P_\mu =\frac{1}{2}(P_{L,R}^{25})^2+\frac{2}{\alpha ^{}}(N_{L,R}1).$$
(15)
We can see that for all states with $`m0`$, as $`R\mathrm{}`$ the mass become infinity, while $`m=0`$ implies that the states take all values for $`n`$ and form a continuum. At the case when $`R0`$, for states with $`n0`$, mass become infinity. However in the limit $`R0`$ for $`n=0`$ states with all $`m`$ values produce a continuum in the spectrum. So, in this limit the compactified dimension disappears. For this reason, we can say that the mass spectrum of the theories at radius $`R`$ and $`\frac{\alpha ^{}}{R}`$ are identical when we interchange $`nm`$. This symmetry is known as T-duality.
The importance of T-duality lies in the fact that the T-duality transformation is a parity transformation acting on the left and right moving degrees of freedom. It leaves invariant the left movers and changes the sign of the right movers (see Eq. (14))
$$P_L^{25}P_L^{25},P_R^{25}P_R^{25}.$$
(16)
The action of T-duality transformation must leave invariant the whole theory (at all order in perturbation theory). Thus, all kind of interacting states in certain theory should correspond to those states belonging to the dual theory. In this context, also the vertex operators are invariant. For instance the tachyonic vertex operators are
$$V_L=exp(iP_L^{25}X_L^{25}),V_R=exp(iP_R^{25}X_R^{25}).$$
(17)
Under T-duality, $`X_L^{25}X_L^{25}`$ and $`X_R^{25}X_R^{25}`$; and from the general solution Eq. (13), $`\alpha _{R,i}^{25}\alpha _{R,i}^{25}`$, $`X_{0R}^{25}X_{0R}^{25}`$. Thus, T-duality interchanges $`nm`$ (Kaluza-Klein modes $``$ winding number) and $`R\frac{\alpha ^{}}{R}`$ in closed string theory.
3.2 T-duality in Open Strings
Now, consider open strings with Neumann boundary conditions. Take again the $`25^{th}`$ coordinate and compactify it on a circle of radius $`R`$, but keeping Neumann conditions. As in the case of closed string, center of mass momentum takes only discrete values $`P^{25}=\frac{n}{R}`$. While there is not analogous for the winding number. So, when $`R0`$ all states with nonzero momentum go to infinity mass, and do not form a continuum. This behavior is similar as in field theory, but now there is something new. The general solutions are
$$X_R^{25}=\frac{X_0^{25}}{2}\frac{a}{2}+\alpha ^{}P^{25}(\tau \sigma )+i\sqrt{\frac{\alpha ^{}}{2}}\underset{l0}{}\frac{1}{l}\alpha _l^{25}exp\left(i2l(\tau \sigma )\right),$$
$$X_L^{25}=\frac{X_0^{25}}{2}+\frac{a}{2}+\alpha ^{}P^{25}(\tau +\sigma )+i\sqrt{\frac{\alpha ^{}}{2}}\underset{l0}{}\frac{1}{l}\alpha _l^{25}exp\left(i2l(\tau +\sigma )\right)$$
(18)
where $`a`$ is a constant. Thus, $`X^{25}(\sigma ,\tau )=X_R^{25}(\sigma ^{})+X_L^{25}(\sigma ^+)=X_0^{25}+\frac{2\alpha ^{}n}{R}\tau +oscillatorterms.`$ Taking the limit $`R0`$, only the $`n=0`$ mode survives. Because of this, the string seems to move in 25 spacetime dimensions. In other words, the strings vibrate in 24 transversal directions. T-duality provides a new T-dual coordinate defined by $`\stackrel{~}{X}^{25}(\sigma ,\tau )=X_L^{25}(\sigma ,\tau )X_R^{25}(\sigma ,\tau )`$. Now, taking $`\stackrel{~}{R}=\frac{\alpha ^{}}{R}`$ we have $`\stackrel{~}{X}^{25}(\sigma ,\tau )=a+2\stackrel{~}{R}\sigma n+oscillatorterms.`$ Using the boundary conditions at $`\sigma =0,\pi `$ one has $`\stackrel{~}{X}^{25}(\sigma ,\tau )_{\sigma =0}=a`$ and $`\stackrel{~}{X}^{25}(\sigma ,\tau )_{\sigma =\pi }=a+2\pi \stackrel{~}{R}n.`$ Thus, we started with an open bosonic string theory with Neumann boundary conditions, and T-duality and a compactification on a circle in the $`25^{th}`$ dimension, give us Dirichlet boundary conditions in such a coordinate. We can visualize this saying that an open string has its endpoints fixed at a hyperplane with 24 dimensions.
Strings with $`n=0`$ lie on a 24 dimensional plane space (D24-brane). Strings with $`n=1`$ has one endpoint at a hyperplane and the other at a different hyperplane which is separated from the first one by a factor equal to $`2\pi \stackrel{~}{R}`$, and so on. But if we compactify $`p`$ of the $`X^i`$ directions over a $`T^p`$ torus ($`i=1,\mathrm{},p`$). Thus, after T-dualizing them we have strings with endpoints fixed at hyperplane with $`25p`$ dimensions, the D$`(25p)`$-brane.
Summarizing: the system of open strings moving freely in spacetime with $`p`$ compactified dimensions on $`T^p`$ is equivalent, under T-duality, to strings whose enpoints are fixed at a D(25-p)-brane i.e. obeying Neumann boundary conditions in the $`X^i`$ longitudinal directions ($`i=1,\mathrm{},p`$) and Dirichlet ones in the transverse coordinates $`X^m`$ ($`m=p+1,\mathrm{},25`$).
The effect of T-dualizing a coordinate is to change the nature of the boundary conditions, from Neumann to Dirichlet and viceverse. If one dualize a longitudinal coordinate this coordinate will satisfies the Dirichlet condition and the D$`p`$-brane becomes a D$`(p+1)`$-brane. But if the dualized coordinate is one of the transverse coordinates the D$`p`$-brane becomes a D$`(p1)`$-brane.
T-duality also acts conversely. We can think to begin with a closed string theory, and compactify it on to a circle in the $`25^{th}`$ coordinate, and then by imposing Dirichlet conditions, obtain a D-brane. This is precisely what occurs in Type II theory, a theory of closed strings.
Spectrum and Wilson Lines
Now, we will see how emerges a gauge field on the D$`p`$-brane world-volume. Again, for the mass shell condition for open bosonic strings and because T-duality $`M^2=(\frac{n}{\alpha ^{}}\stackrel{~}{R})^2+\frac{1}{\alpha ^{}}(N1)`$. The massless state ($`N=1`$, $`n=0`$) implies that the gauge boson $`\alpha _1^\mu 0`$ ($`U(1)`$ gauge boson) lies on to the D24-brane world-volume. On the other hand, $`\alpha _1^{25}0`$ has a vev (vacuum expectation value) which describes the position $`\stackrel{~}{X}^{25}`$ of the D-brane after T-dualizing. Thus, we can say in general, there is a gauge theory $`U(1)`$ over the world volume of the D$`p`$-brane.
Consider now an orientable open string. The endpoints of the string carry charge under a non-Abelian gauge group. For Type II theories the gauge group is $`U(N)`$. One endpoint transforms under the fundamental representation N of $`U(N)`$ and the other one, under its complex conjugate representation (the anti-fundamental one) .
The ground state wave function is specified by the center of mass momentum and by the charges of the endpoints. Thus implies the existence of a basis $`k;ij`$ called Chan-Paton basis. States $`k;ij`$ of the Chan-Paton basis are those states which carry charge 1 under the $`i^{th}`$ $`U(1)`$ generator and $`1`$ under the $`j^{th}`$ $`U(1)`$ generator. So, we can decompose the wave function for ground state as $`k;a=_{i,j=1}^Nk;ij\lambda _{ij}^a`$ where $`\lambda _{ij}^a`$ are called Chan-Paton factors. From this, we see that it is possible to add degrees of freedom to endpoints of the string, that are precisely the Chan-Paton factors.
This is consistent with the theory, because the Chan-Paton factors have a Hamiltonian which do not posses dynamical structure. So, if one endpoint to the string is prepared in a certain state, it always will remains the same. It can be deduced from this, that $`\lambda ^aU\lambda ^aU^1`$ with $`U`$ $`U(N)`$. Thus, the worldsheet theory is symmetric under $`U(N)`$, and this global symmetry is a gauge symmetry in spacetime. So the vector state at massless level $`\alpha _1^\mu k,a`$ is a $`U(N)`$ gauge boson.
When we have a gauge configuration with non trivial line integral around a compactified dimension (i.e a circle), we said there is a Wilson line. In case of open strings with gauge group $`U(N)`$, a toroidal compactification of the $`25^{th}`$ dimension on a circle of radius $`R`$. If we choice a background field $`A^{25}`$ given by $`A^{25}=\frac{1}{2\pi R}diag(\theta _1,\mathrm{},\theta _N)`$ a Wilson line appears. Moreover, if $`\theta _i=0`$, $`i=1,\mathrm{},l`$ and $`\theta _j0`$, $`j=l+1,\mathrm{},N`$ then gauge group is broken: $`U(N)U(l)\times U(1)^{Nl}`$. It is possible to deduce that $`\theta _i`$ plays the role of a Higgs field. Because string states with Chan-Paton quantum numbers $`ij`$ have charges $`1`$ under $`i^{th}`$ $`U(1)`$ factor (and $`1`$ under $`j^{th}`$ $`U(1)`$ factor) and neutral with all others; canonical momentum is given now by $`P_{(ij)}^{25}\frac{n}{R}+\frac{(\theta _j\theta _i)}{2\pi R}.`$ Returning to the mass shell condition it results,
$$M_{ij}^2=\left(\frac{n}{R}+\frac{\theta _j\theta _i}{2\pi R}\right)^2+\frac{1}{\alpha ^{}}(N1).$$
(19)
Massless states ($`N=1,n=0`$) are those in where $`i=j`$ (diagonal terms) or for which $`\theta _j=\theta _i`$ $`(ij)`$. Now, T-dualizing we have $`\stackrel{~}{X}_{ij}^{25}(\sigma ,\tau )=a+(2n+\frac{\theta _j\theta _i}{\pi })\stackrel{~}{R}\sigma +oscillatorterms.`$ Taking $`a=\theta _i\stackrel{~}{R}`$, $`\stackrel{~}{X}_{ij}^{25}(0,\tau )=\theta _i\stackrel{~}{R}`$ and $`\stackrel{~}{X}_{ij}^{25}(\pi ,\tau )=2\pi n\stackrel{~}{R}+\theta _j\stackrel{~}{R}.`$ This give us a set of $`N`$ D-branes whose positions are given by $`\theta _j\stackrel{~}{R}`$, and each set is separated from its initial positions ($`\theta _j=0`$) by a factor equal to $`2\pi \stackrel{~}{R}`$. Open strings with both endpoints on the same D-brane gives massless gauge bosons. The set of $`N`$ D-branes give us $`U(1)^N`$ gauge group. An open string with one endpoint in one D-brane, and the other endpoint in a different D-brane, yields a massive state with $`M(\theta _j\theta _i)\stackrel{~}{R}`$. Mass decreases when two different D-branes approximate to each other, and are null when become the same. When all D-branes take up the same position, the gauge group is enhanced from $`U(1)^N`$ to $`U(N).`$ On the D-brane world-volume there are also scalar fields in the adjoint representation of the gauge group $`U(N)`$. The scalars parametrize the transverse positions of the D-brane in the target space $`X`$.
3.3 D-brane Actions and Ramond-Ramond Charges
D-Brane Action
With the massless spectrum on the D-brane world-volume it is possible to construct a low energy effective action. For open strings massless fields are interacting with the closed strings massless spectrum from the NS-NS sector. Let $`\xi ^a`$ with $`a=0,\mathrm{},p`$ the wold-volume coordinates. The effective action is the gauge invariant action well known as the Born-Infeld action
$$S_D=T_p_Wd^{p+1}\xi e^\mathrm{\Phi }\sqrt{det\left(G_{ab}+B_{ab}+2\pi \alpha ^{}F_{ab}\right)},$$
(20)
where $`T_p`$ is the tension of the D-brane, $`G_{ab}`$ is the world-volume induced metric, $`B_{ab}`$ is the induced antisymmetric field, $`F_{ab}`$ is the Abelian field strength on $`W`$ and $`\mathrm{\Phi }`$ is the dilaton field.
For $`N`$ D-branes the massless fields turns out to be $`N\times N`$ matrices and the action turns out to be non-Abelian Bon-Infeld action (for a nice review about the Born-Infeld action in string theory see )
$$S_D=T_p_Wd^{p+1}\xi e^\mathrm{\Phi }Tr\left(\sqrt{det\left(G_{ab}+B_{ab}+2\pi \alpha ^{}F_{ab}\right)}+O\left([X^m,X^n]^2\right)\right)$$
(21)
where $`m,n=p+1,\mathrm{},10`$. The scalar fields $`X^m`$ representing the transverse positions become $`N\times N`$ matrices and so, the spacetime become a noncommutative spacetime. We will come back later to this interesting point.
Ramond-Ramond Charges
D-branes are coupled to Ramond-Ramond (RR) fields $`G_p`$. The complete effective action on the D-brane world-volume $`W`$ which take into account this coupling is
$$S_D=T_p_Wd^{p+1}\xi \left\{e^\mathrm{\Phi }\sqrt{det\left(G_{ab}+B_{ab}+2\pi \alpha ^{}F_{ab}\right)}+i\mu _p_W\underset{p}{}C_{(p+1)}Tr\left(e^{2\pi \alpha ^{}(F+B)}\right)\right\}$$
(22)
where $`\mu _p`$ us the RR charge. RR charges can be computed by considering the anomalous behavior of the action at intersections of D-branes . Thus RR charge is given by
$$Q_{RR}=ch(j!E)\sqrt{\widehat{A}(TX)}$$
(23)
where $`j:WX`$. Here $`E`$ is the Chan-Paton bundle over $`X`$, $`\widehat{A}(TX)`$ is the genus of the spacetime manifold $`X`$. This gives an ample evidence that the RR charges take values not in a cohomology theory, but in fact, in a K-Theory . This result was further developed by Witten in the context of non-BPS brane configurations worked out by A. Sen. This subject will be reviewed below in Sec. V.
## IV String Duality
4.1 Duality in Field Theory
Duality is a notion which in the last years has led to remarkable advances in nonperturbative quantum field theory. It is an old known type of symmetry which by interchanging the electric and magnetic fields leaves invariant the vacuum Maxwell equations. It was extended by Dirac to include sources, with the well known price of the prediction of monopoles, which appear as the dual particles to the electrically charged ones and whose existence could not be confirmed up to now. Dirac obtained that the couplings (charges) of the electrical and magnetical charged particles are the inverse of each other, $`i.e.`$ as the electrical force is ‘weak’, and it can be treated perturbatively, then the magnetic force among monopoles will be ‘strong’ (for some reviews see ).
This duality, called ‘$`S`$-duality’, has inspired a great deal of research in the last years. By this means, many non-perturbative exact results have been established. In particular, the exact Wilson effective action of $`𝒩=2`$ supersymmetric gauge theories has been computed by Seiberg and Witten, showing the duality symmetries of these effective theories. It turns out that the dual description is quite adequate to address standard non-perturbative problems of Yang-Mills theory, such as confinement, chiral symmetry breaking, etc.
$`𝒩=4`$ supersymmetric gauge theories in four dimensions have vanishing renormalization group $`\beta `$-function. Montonen and Olive conjectured that (at the quantum level) these theories would possess an SL$`(2,𝐙)`$ exact dual symmetry. Many evidences of this fact have been found, although a rigorous proof does not exist at present. For $`𝒩=2`$ supersymmetric gauge theories in four dimensions, the $`\beta `$-function in general does not vanish. So, Montonen-Olive conjecture cannot be longer valid in the same sense as for $`𝒩=4`$ theories. However, Seiberg and Witten found that a strong-weak coupling ’effective duality’ can be defined on its low energy effective theory for the cases pure and with matter. The quantum moduli space of the pure theory is identified with a complex plane, the $`u`$-plane, with singularities located at the points $`u=\pm 1,\mathrm{}`$. It turns out that at $`u=\pm 1`$ the original Yang-Mills theory is strongly coupled, but effective duality permits the weak coupling description at these points in terms of monopoles or dyons (dual variables). $`𝒩=1`$ gauge theories are also in the class of theories with non-vanishing $`\beta `$-function. More general, for a gauge group SU$`(N_c)`$, an effective non-Abelian duality is implemented even when the gauge symmetry is unbroken. It has a non-Abelian Coulomb phase. Seiberg has shown that this non-Abelian Coulomb branch is dual to another non-Abelian Coulomb branch of a theory with gauge group SU$`(N_fN_c)`$, where $`N_f`$ is the number of flavors. $`𝒩=1`$ theories have a rich phase structure. Thus, it seems that in supersymmetric gauge theories strong-weak coupling duality can only be defined for some particular phases.
For non-supersymmetric gauge theories in four dimensions, the subject of duality has been explored recently in the Abelian as well as in the non-Abelian cases. In the Abelian case (on a curved compact four-manifold $`X`$) the $`CP`$ violating Maxwell theory partition function $`Z(\tau )`$, transforms as a modular form under a finite index subgroup $`\mathrm{\Gamma }_0(2)`$ of SL$`(2,𝐙)`$. The dependence parameter of the partition function is given by $`\tau =\frac{\theta }{2\pi }+\frac{4\pi i}{e^2}`$, where $`e`$ is the Abelian coupling constant and $`\theta `$ is the usual theta angle. In the case of non-Abelian non-supersymmetric gauge theories, strong-coupling dual theories can be constructed which results in a kind of dual “massive” non-linear sigma models. The starting Yang-Mills theory contains a $`CP`$-violating $`\theta `$-term and it turns out to be equivalent to the linear combination of the actions corresponding to the self-dual and anti-self-dual field strengths.
4.2 String Duality
In Sec. I we have described the massless spectrum of the five consistent superstring theories in ten dimensions. Additional theories can be constructed in lower dimensions by compactification of some of the ten dimensions. Thus the ten-dimensional spacetime $`X`$ looks like the product $`X=K^d\times 𝐑^{1,9d}`$, with $`K`$ a suitable compact manifold or orbifold. Depending on which compact space is taken, it will be the quantity of preserved supersymmetry.
All five theories and their compactifications are parametrized by: the string coupling constant $`g_S`$, the geometry of the compact manifold $`K`$, the topology of $`K`$ and the spectrum of bosonic fields in the NS-NS and the R-R sectors. Thus one can define the string moduli space of each one of the theories as the space of all associated parameters. Moreover, it can be defined a map between two of these moduli spaces. The dual map is defined as the map $`𝒮:^{}`$ between the moduli spaces $``$ and $`^{}`$ such that the strong coupling region of $``$ is interchanged with the weak-coupling region of $`^{}`$ and viceverse. One can define another map $`𝒯:^{}`$ which interchanges the volume $`V`$ of $`K`$ for $`\frac{1}{V}`$. One example of the map $`𝒯`$ is the equivalence, by T-duality, between the theories Type IIA compactified on $`𝐒^1`$ at radius $`R`$ and the Type IIB theory on $`𝐒^1`$ at raduis $`\frac{1}{R}`$. The theories Het($`E_8\times E_8`$) and Het(SO$`(32)`$) is another example. In this section we will follows the Sen’s review . Another useful reviews are .
4.3 Type I-SO(32)-Heterotic Duality
In order to analyze the duality between Type I and SO(32) heterotic string theories we recall from Sec. II the spectrum of both theories. These fields are the dynamical fields of a supergravity Lagrangian in ten dimensions. Type I string theory has in the NS-NS sector the fields: the metric $`g_{\mu \nu }^I`$, the dilaton $`\mathrm{\Phi }^I`$ and in the R-R sector: the antisymmetric tensor $`B_{\mu \nu }^I`$. Also there are 496 gauge bosons $`A_\mu ^{aI}`$ in the adjoint representation of the gauge group SO(32). For the SO(32) heterotic string theory the spectrum consist of: the spacetime metric $`g_{\mu \nu }^H`$, the dilaton field $`\mathrm{\Phi }^H`$, the antisymmetric tensor $`B_{\mu \nu }^H`$ and 496 gauge fields $`A_\mu ^{aH}`$ in the adjoint representation of SO(32). Both theories have spacetime supersymmetry $`𝒩=1`$. The effective action for the massless fields of the Type I supergravity effective action $`S_I`$ is defined at tree-level on the disk. Thus the string coupling constant $`g_s^I`$ arises in the Einstein frame as $`exp(\mathrm{\Phi }^I/4)`$. While the heterotic action $`S_H`$ is defined on the sphere and $`g_S^H`$ is given by $`exp(\mathrm{\Phi }^H/4)`$. The comparison of these two actions in the Einstein frame leads to the following identification of the fields
$$g_{\mu \nu }^I=g_{\mu \nu }^H,B_{\mu \nu }^I=B_{\mu \nu }^H$$
$$A_\mu ^{aI}=A_\mu ^{aH},\mathrm{\Phi }^I=\mathrm{\Phi }^H.$$
(24)
This give us many information, the first relation tell us that the metrics of both theories are the same. The second relation interchanges the $`B`$ fields in the NS-NS and the R-R sectors. That interchanges heterotic strings by Type I D1-branes. The third relation identifies the gauge fields coming from the Chan-Paton factors from the Type I side with the gauge fields coming from the 16 compactified internal dimensions of the heterotic string. Finally, the opposite sign for the dilaton relation means that the string coupling constant $`g_S^I`$ is inverted $`g_S^H=1/g_S^I`$ within this identification, and interchanges the strong and weak couplings of both theories leading to the explicit realization of the $`𝒮`$ map.
4.4 Type II-Heterotic Duality
Lower dimensional theories constructed up on compactification can have different spacetime supersymmetry. Thus it can be very useful to find dual pairs by compactifying two string theories with different spacetime supersymmetry on different spaces $`K`$ in such a way that they become to have the same spacetime supersymmetry.
Perhaps the most famous example is the $`S`$-dual pair between the Type II theory on $`K3`$ and the heterotic theory on $`T^4`$. To describe more generally these kind of dualities we first give some preliminaries. Let $`𝒜`$ and $``$ two different theories of the family of string theories. $`𝒜`$ and $``$ are compactified on $`K_𝒜`$ and $`K_{}`$ respectively. Consider the dual pair
$$𝒜/K_𝒜/K_{}$$
(25)
then we can construct the more general dual pair
$$𝒜/Q_𝒜/Q_{},$$
(26)
where $`K_𝒜Q_𝒜D`$ and $`K_{}Q_{}D`$ are fibrations and $`D`$ is an auxiliary finite dimensional manifold.
These insights are very useful to construct dual pairs for theories with eight supercharges. An example of this is the pair in six dimensions with $`𝒜=IIA`$, $`K_𝒜=K3`$ and $`=Het`$, $`K_{}=T^4`$ i.e.
$$IIA/K3Het/T^4.$$
(27)
From this a dual pair can be constructed in four dimensions with the auxiliary space $`D=\mathrm{𝐂𝐏}^1`$ being the complex projective space, thus we have
$$IIA/CYHet/K3\times T^2,$$
(28)
where $`T^4Q_{IIA}\mathrm{𝐂𝐏}^1`$ and $`K3Q_{Het}\mathrm{𝐂𝐏}^1`$ are fibrations. As can be observed the four-dimensional theories have $`𝒩=2`$ supersymmetry and the duality uses K3-fibrations.
4.5 M-Theory
We have described how to construct dual pairs of string theories. By the uses of the $`𝒮`$ and the $`𝒯`$ maps a network of theories can be constructed in various dimensions all of them related by dualities. However new theories can emerge from this picture, this is the case of M-theory. M-theory (the name come from ‘mystery’, ‘magic’, ‘matrix’, ‘membrane’, etc.) was originally defined as the strong coupling limit for Type IIA string theory . At the effective low energy action level, Type IIA theory is described by the Type IIA supergravity theory and it is known that this theory can be obtained from the dimensional reduction of the eleven dimensional supergravity theory (a theory known from the 70’s years). Let $`Y`$ be the eleven dimensional manifold, taking $`Y=X^{10}\times 𝐒_R^1`$ the compactification radius $`R`$ is proportional to $`g_{10,10}\mathrm{\Phi }`$. Thus the limit $`\mathrm{\Phi }\mathrm{}`$ corresponds to the limit $`R\mathrm{}`$ and thus the strong coupling limit of the Type IIA theory corresponds to the 11 dimensional supergravity. It is conjectured that there exist an eleven dimensional fundamental theory whose low energy limit is the 11 dimensional supergravity theory. At the present time the degrees of freedom are still unknown, through at the macroscopic level they should be membranes and fivebranes (also called M-two-branes and M-fivebranes). There is a proposal to describe dof of M-theory in terms of a gas of D0-branes. This is the called ‘Matrix Theory’. This proposal as been quite successful (for some reviews see and references therein).
4.6 Horava-Witten Theory
Just as the M-theory compactification on $`𝐒_R^1`$ leads to the Type IIA theory, Horava and Witten realized that orbifold compactifications leads to the $`E_8\times E_8`$ heterotic theory in ten dimensions (see for instance ). More precisely
$$\mathrm{M}/𝐒^1/𝐙_2E_8\times E_8Het$$
(29)
where $`𝐒^1/𝐙_2`$ is homeomorphic to the finite interval $`I`$ and the $`M`$-theory is thus defined on $`Y=X^{10}\times I`$. From the ten-dimensional point of view, this configuration is seen as two parallel planes placed at the two boundaries $`I`$ of $`I`$. Dimensional reduction and anomalies cancellation conditions imply that the gauge degrees of freedom should be trapped on the ten-dimensional planes $`X`$ with the gauge group being $`E_8`$ in each plane. While that the gravity is propagating in the bulk and thus both copies of $`X`$’s are only connected gravitationally.
4.7 F-Theory
$`F`$-Theory was formulated by C. Vafa, looking for an analog theory to M-Theory for describing non-perturbative compactifications of Type IIB theory (for a review see ). Usually in perturbative compactifications the parameter $`\lambda =a+iexp(\mathrm{\Phi }/2)`$ is taken to be constant. $`F`$-theory generalizes this fact by considering variable $`\lambda `$. Thus $`F`$-theory is defined as a twelve-dimensional theory whose compactification on the elliptic fibration $`T^2D`$, gives the Type IIB theory compactified on $`D`$ (for a suitable space $`D`$) with the identification of $`\lambda (\stackrel{}{z})`$ with the modulus $`\tau (\stackrel{}{z})`$ of the torus $`T^2`$. These compactifications can be related to the $`M`$-theory compactifications through the known $`S`$ mapping $`𝒮:IIAM/𝐒^1`$ and the $`𝒯`$ map between Type IIA and IIB theories. This gives
$$F/\times 𝐒^1M/.$$
(30)
Thus the spectrum of massless states of $`F`$-theory compactifications can be described in terms of $`M`$-theory. Other interesting $`F`$-theory compactifications are the Calabi-Yau compactifications
$$F/CYHet/K3.$$
(31)
4.8 Gravitational Duality
As a matter of fact, string theory constitutes nowadays the only consistent and phenomenologically acceptable way to quantize gravity. It contains in its low energy limit Einstein gravity. Thus, a legitimate question is the one of which is the ‘dual’ theory of gravity or, more precisely, how gravity behaves under duality transformations.
Gravitational analogs of non-perturbative gauge theories were studied several years ago, particularly in the context of gravitational Bogomolny bound. As recently was shown , there are additional non-standard $`p`$-branes in $`D=10`$ type II superstring theory and $`D=11`$ M-theory, and which are required by U-duality. These branes were termed ‘gravitational branes’ (‘G-branes’), because they carry global charges which correspond to the $`ADM`$ momentum $`P_M`$ and to its ‘dual’, a $`(D5)`$-form $`K_{M_1\mathrm{}M_{D5}},`$ which is related to the NUT charge. These charges are ‘dual’ in the same sense that the electric and magnetic charges are dual in Maxwell theory, but they appear in the purely gravitational sector of the theory. Last year, Hull has shown in that these global charges $`P`$ and $`K`$ arise as central charges of the supersymmetric algebra of type II superstring theory and M-theory. Thus the complete spectrum of BPS states should include the gravitational sector.
Finally, a different approach to the ‘gravitational duality’ was worked out by using some techniques of strong-weak coupling duality for non-supersymmetric Yang-Mills theories were applied to the MacDowell-Mansouri dynamical gravity (for a review see ). One would suspect that both approaches might be related in some sense. One could expect that the gauge theory of gravity would be realized as the effective low energy theory on the ‘G-branes’.
## V Non-BPS Branes and K-theory
5.1 Non-BPS Branes
The notion of D-branes as BPS states implies the existence of certain supersymmetric theory on the world-volume of the D-brane. However it is extremaly relevant the consideration of non-supersymmetric theories (in order to describe our non-supersymmetric world) and here is where it is important the construction of brane configurations without remanent supersymmetry. A. Sen proposed the construction of such non-supersymmetric configurations by considering pairs of D-branes and anti-D-branes (for a nice review see see also ). These configurations break all supersymmetry and the spectrum on the world-volume has a tachyon which cannot be cancelled by GSO projection. The presence of this tachyon leads to unstable brane configuration and the configuration decay into an stable BPS configuration . The classification of these stable D-branes was given by Witten in terms of topological K-Theory in the beautiful seminal paper (for a review of this exciting subject see ).
In order to fix some notation let $`X`$ be the ten-dimensional spacetime manifold and let $`W`$ be a $`(p+1)`$-dimensional submanifold of $`X`$. Branes or antibranes or both together can be wrapped on $`W`$. When configurations of $`N`$ coincident branes or antibranes only are wrapped on $`W`$, the world-volume spectra on $`W`$ consists of a vector multiplet and scalars in some representation of the gauge group. These configurations can be described through Chan-Paton bundles which are U$`(N)`$ gauge bundles $`E`$ over $`W`$ for Type II superstring theory and by SO$`(N)`$ or Sp$`(N)`$ bundles in Type I theory. Gauge fields from the vector multiplet define a U$`(N)`$ gauge connection for Type II theory (or SO$`(N)`$ or Sp$`(N)`$ gauge connection for Type I theory) on the (corresponding) Chan-Paton bundle. GSO projection cancels the usual tachyonic degrees of freedom. Something similar occurs for the anti-brane sector.
The description of coincident $`N_1`$ coincident $`p`$-branes and $`N_2`$ $`p`$-anti-branes wrapped on $`W`$ leads to the consideration pairs of gauge bundles $`(E,F)`$ (over $`W`$) with their respective gauge connections $`A`$ and $`A^{}`$. In the mixed configurations GSO projection fails to cancel the tachyon. Thus the system is unstable and may flow toward the annihilation of the brane-antibrane pairs with RR charge for these brane configurations being conserved in the process.
On the open string sector Chan-Paton factors are $`2\times 2`$ matrices constructed from the possible open strings stretched among the different types of branes. Brane-brane and antibrane-antibrane sectors correspond to the diagonal elements of this matrix. Off-diagonal elements correspond with the Chan-Paton labels of an oriented open string starting at a brane and ending at an antibrane and the other one to be the open string with opposite orientation.
The physical mechanism of brane-antibrane creation or annihilation without violation of conservation of the total RR charge, leads to consider physically equivalent configurations of $`N_1`$ branes and $`N_2`$ antibranes and the same configuration but with additional created or annihilated brane-antibrane pairs.
5.2 D-branes and K-Theory
The relevant mathematical structure describing the brane-antibrane pairs in general type I and II superstring theories is as follows:
1. $`G_1`$ and $`G_2`$ gauge connections $`A`$ and $`A^{}`$ on the Chan-Paton bundles $`E`$ and $`F`$ over $`W`$, respectively. Bundles $`E`$ and $`F`$ corresponding to branes and antibranes are topologically equivalent. The groups $`G_1`$ and $`G_2`$ are restricted to be unitary groups for Type II theories and symplectic or orthogonal groups for Type I theories.
2. Tachyon field $`T`$ can be seen as a section of the tensor product of bundles $`EF^{}`$ and its conjugate $`\overline{T}`$ as a section of $`E^{}F`$ (where $``$ denotes the dual of the corresponding bundle.)
3. Brane-antibrane configurations are described by pairs of gauge bundles $`(E,F)`$.
4. The physical mechanism of brane-antibrane creation or annihilation of a set of $`m`$ $`9`$-branes and $`9`$-antibranes is described by the same U$`(m)`$ (for Type II theories) or SO$`(m)`$ (for type I theories) gauge bundle $`H`$. This mechanism is described by the identification of pairs of gauge bundles $`(E,F)`$ and $`(EH,FH)`$. Actually instead of pairs of gauge bundles one should consider classes of pairs of gauge bundles $`[(E,F)]=[E][F]`$ identified as above. Thus the brane-antibrane pairs really determine an element of the K-theory group K$`(X)`$ of gauge bundles over $`X`$ and the brane-antibrane creation or annihilation of pairs is underlying the $`K`$-theory concept of stable equivalence of bundles. For 9-branes, the embedded submanifold $`W`$ coincides with $`X`$ and the thus brane charges take values in K-theory group of $`X`$.
Consistency conditions for 9-branes ($`p=9`$) in Type IIB superstring theory such as tadpole cancellation implies the equality of the ranks of the structure groups of the bundles $`E`$ and $`F`$. Thus $`rk(G_1)=rk(G_2)`$. The ‘virtual dimension’ $`d`$ of an element $`(E,F)`$ is defined by $`d=rk(G_1)rk(G_2)`$. Thus tadpole cancellation leads to a description of the theory in terms of pairs of bundles with virtual dimension vanishing, $`d=0`$. This is precisely the definition of reduced K-theory $`\stackrel{~}{\mathrm{K}}(X)`$. Thus consistency conditions implies to project the description to reduced K-theory.
In Type I string theory $`9\overline{9}`$ pairs are described by a class of pairs $`(E,F)`$ of SO$`(N_1)`$ and SO$`(N_2)`$ gauge bundles over $`X`$. Creation-annihilation is now described through the SO$`(k)`$ bundle $`H`$ over $`X`$. In Type I theories tadpole cancellation condition is $`N_1N_2=32`$. In this case equivalence class of pair bundles $`(E,F)`$ determines an element in the real K-theory group KO$`(X)`$. Tadpole cancellation $`N_1N_2=32`$, newly turns out into reduced real K-theory group $`\stackrel{~}{\mathrm{KO}}(X)`$.
Type IIA theory involves more subtle. It was argued by Witten in that configurations of brane-antibrane pairs are classified by the K-theory group of spacetime with an additional circle space $`𝐒^1\times X`$. K-theory group for type IIA configurations is K$`(𝐒^1\times X)`$.
5.3 Ramond-Ramond Fields and K-Theory
Ramond-Ramond charges are classified according to the K-theory groups. In this subsection we will review that the proper RR fields follows a similar classification. For details see the recent papers by Witten and by Moore and Witten .
It can be showed that RR fields do not satisfy the Dirac quantization condition. Thus for example,
$$_{W_p}\frac{G_p}{2\pi }𝐙.$$
(32)
The reason of this is the presence of chiral fermions on the brane. The phase of the fermions contribute with a gravitational term $`\lambda =_W\frac{1}{16\pi ^2}tr\left(RR\right)`$. This gives a correction to the Dirac quantization. In trying to extended it for the all RR fields $`G_p`$ in string theory it is necessary introduce new ideas as the notion of quantum self-duality and K-theory. Thus RR fields should be generalized in the context of K-Theory and we will see that in fact, they find an appropriate description within this context. Similar as the RR charge, the RR fields find a natural classification in terms of K-theory.
For self-dual RR fields it is a very difficult to find the quantum partition function. For the scalar field in two-dimensions it can be obtained by summing over only one of the periods of the 2-torus. It is not possible to sum simultaneously over both periods. This description can be generalized to any higher degree $`p`$-forms $`G_p`$. It can be done by defining a function $`\mathrm{\Omega }(x)`$ for $`x`$ in the lattice $`\{H^1(\mathrm{\Sigma },𝐙)\}`$ of periods such that
$$\mathrm{\Omega }(x+y)=\mathrm{\Omega }(x)\mathrm{\Omega }(y)(1)^{(x,y)},$$
(33)
where $`(x,y)xy=xy`$.
The partition function can be constructed easily from these data. One first has to identify the corresponding period lattice $`\lambda `$. After that, find the $`\mathrm{\Omega }`$ function as a $`𝐙`$-valued function on $`\mathrm{\Lambda }`$ such that it satisfies Eq. (33). Finally one has to construct the partition function.
Period Lattice for Ramond-Ramond Fields
Let $`X`$ be the spacetime manifold. One could suppose that that period lattice are: $`_{peven}H^p(X;𝐙)`$ for Type IIA theory and $`_{podd}H^p(X;𝐙)`$ for Type IIB theory. However this are not the right choice since the RR charges and fields actually take values in K-Theory, just as has been described in the last subsection. Thus, one can see that the period lattice for Type IIA theory is K$`(X)`$ and for Type IIB it is K$`{}_{}{}^{1}(X)`$. This is more obvious from the anomalous brane couplings. If $`X=𝐑\times Y`$ we have
$$\frac{dG}{2\pi }=\delta (Y)\sqrt{\widehat{A}(Y)}ch(TX).$$
(34)
Hence the period lattice is constructed from $`_{peven}H^p(X;𝐑)`$ generated by $`\sqrt{\widehat{A}}ch(TX)`$ for $`x=(E,F)K(X)`$. Still it is necessary to quantize the lattice by finding the $`\mathrm{\Omega }`$ function and its corresponding quantum partition function.
The $`\mathrm{\Omega }`$ Function
In K-Theory there exist a natural definition of the $`\mathrm{\Omega }`$ function given by the index theory
$$(x,y)=\mathrm{Index}\mathrm{of}\mathrm{Dirac}\mathrm{Operator}\mathrm{on}\mathrm{X}\mathrm{with}\mathrm{values}\mathrm{in}\mathrm{x}\overline{\mathrm{y}}=_X\widehat{A}(X)ch(x)ch(\overline{y}).$$
(35)
Thus the $`\mathrm{\Omega }`$ function can be defined as
$$\mathrm{\Omega }(x)=(1)^{j(x)},$$
(36)
where $`j(x)`$ is given by the mod 2 index of the Dirac operator with values in the real bundle $`x\overline{x}`$. It can be shown that this definition of $`\mathrm{\Omega }(x)`$ satisfies the relation (33). From this one can construct a quantum partition function which is compatible with $`(i)`$ T-duality, $`(ii)`$ Self-duality of RR fields, $`(iii)`$ the interpretation of RR fields in K-theory and $`(iv)`$ description of the brane anomalies.
## VI String Theory and Noncommutative Gauge Theory
6.1 Noncommutative D-branes From String Interactions
Finally in this section we describe briefly some new developments on the relation between string theory and Connes’s noncommutative Yang-Mills theory (for a survey on noncommutative geometry see the classic Connes book ). We do not pretend to be exhaustive but only to remark the key points of the recent exciting developments (for a nice review see ).
The roughly idea consists from the description of a string propagating in a flat background (spacetime) of metric $`g_{ij}`$ and a NS constant $`B`$-field $`B_{ij}`$. The action is given by
$$=\frac{1}{4\pi \alpha ^{}}_Dd^2\sigma \left(g_{ij}_aX^i^aX^j2\pi i\alpha ^{}B_{ij}\epsilon ^{ab}_aX^i_bX^j\right)$$
(37)
where $`D`$ is the disc. Or equivalently
$$=\frac{1}{4\pi \alpha ^{}}_Dd^2\sigma g_{ij}_aX^i^aX^j\frac{i}{2}_D𝑑\tau B_{ij}X^i_\tau X^j$$
Equations of motion from this action are subjected to the boundary condition
$$g_{ij}_nX^j+2\pi i\alpha ^{}B_{ij}_tX^j|_\mathrm{\Sigma }=0.$$
(38)
The propagator of open string vertex operators is given by
$$X^i(\tau )X^j(\tau ^{})=\alpha ^{}G^{ij}\mathrm{log}(\tau \tau ^{})^2+\frac{i}{2}\mathrm{\Theta }^{ij}\epsilon (\tau \tau ^{})$$
(39)
where
$$G^{ij}=\left(\frac{1}{g+2\pi \alpha ^{}B}\right)_S^{ij},\mathrm{\Theta }^{ij}=2\pi \alpha ^{}\left(\frac{1}{g+2\pi \alpha ^{}B}\right)_A^{ij}.$$
(40)
Here $`S`$ and $`A`$ stands for the symmetric and antisymmetric part of the involved matrix, and the logarithmic term determines the anomalous dimensions as usual. Thus $`G_{ij}`$ is the effective metric seen by the open strings. While, as was suggested by Schomerus , the antisymmetric part $`\mathrm{\Theta }^{ij}`$ determines the noncommutativity.
The product of tachyon vertex operators $`exp(ipX)`$ and $`exp(iqX)`$ for $`\tau >\tau ^{}`$ in the short distance singularity is written as
$$exp\left(ipX\right)(\tau )exp\left(iqX\right)(\tau ^{})(\tau \tau ^{})^{2\alpha ^{}G^{ij}p_iq_j}exp\left(\frac{1}{2}\mathrm{\Theta }^{ij}p_iq_j\right)exp\left(i(p+q)X\right)(\tau ^{})+\mathrm{}$$
(41)
or
$$exp\left(ipX\right)exp\left(iqX\right)exp\left(ipX\right)exp\left(iqX\right)exp\left(\frac{i}{2}\mathrm{\Theta }^{ij}p_iq_j\right)exp\left(i(p+q)X\right)$$
(42)
where $``$ is defined for any smooth functions $`F`$ and $`G`$ over $`X`$ and it is given by
$$FG=exp\left(\frac{i\mathrm{}}{2}\mathrm{\Theta }^{ij}\frac{}{u_i}\frac{}{v_j}\right)F(x+u)G(x+y).$$
(43)
Here the operation $``$ is associative $`F(GH)=(FG)H`$ and noncommutative $`FGGF`$. The above product can be written as $`FG=FG+i\{F,G\}+\mathrm{}`$ where $`\{F,G\}`$ is the Poisson bracket given by $`\mathrm{\Theta }^{ij}_iF_jG.`$ $`\mathrm{\Theta }`$ is determined in terms of $`B`$. Its give an associative and noncommutative algebra. In the limit $`\alpha ^{}0`$ (ignoring the anomalous dimensions of open string sector) the product of vertex operators turns out to be the Moyal product of functions on the spacetime $`X`$.
Now one can consider scattering amplitude (parametrized by $`G`$ and $`\mathrm{\Theta }`$) of $`k`$ gauge bosons of momenta $`p_i`$, polarizations $`\epsilon _i`$ and Chan-Paton wave functions $`\lambda _i,`$ $`ı=1,\mathrm{},k`$
$$A(\lambda _i,\epsilon _i,p_i)_{G,\mathrm{\Theta }}=Tr\left(\lambda _1\lambda _2\mathrm{}\lambda _k\right)𝑑\tau _i^{}\underset{i=1}{\overset{k}{}}\epsilon _i\frac{dX}{d\tau }exp\left(ip_iX\right)(\tau _i^{})_{G,\mathrm{\Theta }}.$$
(44)
The $`\mathrm{\Theta }`$ dependence come from the factor $`exp\left(\frac{i}{2}_{s>r}p_i^{(s)}p_j^{(r)}\mathrm{\Theta }^{ij}\right)`$. Thus amplitude factorizes as $`A(\lambda _i,\epsilon _i,p_i)_{G,\mathrm{\Theta }=0}exp\left(\frac{i}{2}_{s>r}p_i^{(s)}\mathrm{\Theta }^{ij}p_j^{(r)}\epsilon (\tau _r\tau _s)\right)`$ which depends only on the cycle ordering of the points $`\tau _1,\mathrm{},\tau _k`$ on the boundary of the disc $`D.`$
For $`B=0`$ the effective action is obtained under the assumption that the divergences are regularized through the Pauli-Villars procedure and it is given by
$$S_G=\frac{1}{g_{st}}d^nx\sqrt{G}\left(TrF_{ij}F^{ij}+\alpha ^{}\mathrm{corrections}\right).$$
(45)
The important case of the effective theory when $`\mathrm{\Theta }0`$ is incorporated through the phase factor and thus one have to replace the ordinary multiplication of wave functions by the $``$ product (effective action is computed by using point splitting regularization)
$$\widehat{S}_G=\frac{1}{g_{st}}d^nx\sqrt{G}G^{ii^{}}G^{jj^{}}\left(Tr\widehat{F}_{ii^{}}\widehat{F}^{jj^{}}+\alpha ^{}\mathrm{corrections}\right),$$
(46)
where $`\widehat{F}_{ij}=_i\widehat{A}_j_j\widehat{A}_ii\{\widehat{A}_i,\widehat{A}_j\}_M`$ is the noncommutative field strength. Here $`\{F,G\}_MFGGF`$. Thus we get a noncommutative Yang-Mills theory as the $`\mathrm{\Theta }`$ (or $`B`$) dependence of the effective action to all orders in $`\alpha ^{}`$. Gauge field transformation $`(\widehat{\lambda }\widehat{A})_{ij}=\widehat{\lambda }_{ik}A_j^k`$ and $`\delta \widehat{A}_i=_i\widehat{\lambda }+i\widehat{\lambda }\widehat{A}_ii\widehat{A}_i\widehat{\lambda }.`$
For the low varying fields the effective action is given by the Born-Infeld-Dirac action
$$S=\frac{1}{g_{st}(\alpha ^{})^2}d^nx\sqrt{\mathrm{det}\left(g+\alpha ^{}(F+B)\right)}.$$
(47)
The same effective action is described by noncommutative Yang-Mills theory but also by the standard Yang-Mills theory. They differ only in the regularization prescription. For the standard commutative case it is the Pauli-Villars one, while for the noncommutative case it is the point splitting prescription. The two frameworks are equivalent and thus there is a redefinition of the variable fileds and it can be seen ‘as a transformation connecting standard and noncommutative descriptions. The change of variables known as the Seiberg-Witten map is as follows
$$\widehat{A}_i=A_i\frac{1}{4}\mathrm{\Theta }^{kl}\{A_k,_lA_i+F_{li}\}+O(\mathrm{\Theta }^2)$$
$$\widehat{\lambda }=\lambda +\frac{1}{4}\mathrm{\Theta }^{kl}\{_l\lambda ,+A_j\}+O(\mathrm{\Theta }^2).$$
(48)
6.2 String Theory and Deformation Quantization
Very recently a renewed deal of excitation has been taken place in deformation quantization theory , since the Kontsevich’s seminal paper . In this paper Kontsevich proved by construction the existence of a star-product for any finite dimensional Poisson manifold. His construction is based on his more general statement known as the “formality conjecture”. The existence of such a star-product determines the existence of a deformation quantization for any Poisson manifold. Kontsevich’s proof was strongly motivated by some perturbative issues of string theory and topological gravity in two-dimensions, such as, matrix models, the triangulation of the moduli space of Riemann surfaces and mirror symmetry.
One of the main lessons of the stringly and D-brane descriptions of Kontsevich’s formula in that of the deformation quantization for any Poisson manifold requires necessarily of string theory. In addition this was confirmed in . The deformation parameter of this quantization is precisely the string scale $`\alpha ^{}`$ (or the string coupling constant) which in the limit $`\alpha ^{}0`$ it reproduces the field theory limit but in this limit the deformation quantization does not exist. The deformation arising precisely when $`\alpha ^{}0`$ is an indication that deformation quantization is an stringly phenomenon. Actually it was already suspected since the origin of the formality conjecture where several mathematical ingredients of string theory were present.
String action in a background NS constant $`B`$ field is
$$S=\frac{1}{4\pi \alpha ^{}}_Dd^2z_aX^i_aX^jG_{ij}+\frac{1}{4\pi \alpha ^{}}_D𝑑z𝑑\overline{z}J^i(z)\overline{J}^j(\overline{z})B_{ij},$$
(49)
where $`J^i(x)=2iX^i(z,\overline{z})`$ and $`\overline{J}^i(x)=2i\overline{}X^i(z,\overline{z})`$.
Define the function
$$F(X(x))=V[F](x):=\frac{1}{(2\pi )^{d/2}}d^dk\widehat{F}(k)V_k(x)$$
(50)
where $`V_k(x)=:exp\left(ik_iX^i(x)\right):`$ is the vertex operator. OPE between $`J^{}s`$ and $`V^{}s`$ operators leads to $`V[F](1)V[G](0)V[FG](0)+\mathrm{}`$. The introduction of a NS constant $`B`$ field in the action ‘deforms’ the OPE leading to
$$(V[F](1)V[G](0)^BV[FG](0)+\mathrm{}$$
(51)
where the $``$ product will be determined. It can be obtained by computing the $`N`$-point correlations functions for the complete action (49) (including the $`B`$-term)
$$\mathrm{\Phi }_1\mathrm{\Phi }_2\mathrm{}\mathrm{\Phi }_N^B=\frac{1}{Z}\mathrm{\Phi }_1\mathrm{\Phi }_2\mathrm{}\mathrm{\Phi }_Nexp\left(\frac{1}{4\pi \alpha ^{}}_{}𝑑z𝑑\overline{z}J^i(z)\overline{J}^j(\overline{z})B_{ij}\right)^B$$
$$=\frac{1}{Z}\underset{n=0}{\overset{\mathrm{}}{}}\left(\frac{1}{4\pi \alpha ^{}}\right)\frac{1}{n!}_{_n^\epsilon }d^dzd^d\overline{z}\mathrm{\Phi }_1\mathrm{\Phi }_2\mathrm{}\mathrm{\Phi }_N\underset{a=1}{\overset{n}{}}B_{i_aj_a}J^{i_a}(z_a)\overline{J}^{i_a}(\overline{z}_a)$$
(52)
where $`Z:=exp\left(\frac{1}{4\pi \alpha ^{}}_{}𝑑z𝑑\overline{z}J^i(z)\overline{J}^j(\overline{z})B_{ij}\right)^B`$ and $`_n^\epsilon :=\{(z_1,z_2,\mathrm{},z_n)|Im(z_a)>\epsilon ,|z_az_b|>\epsilon \mathrm{for}a0\}`$. We choice $`\mathrm{\Phi }_1=V[F](1)`$ and $`\mathrm{\Phi }_2[G](0)`$. Now using the usual OPE of the $`J^{}s`$ and $`V_k^{}s`$ operators and substituting all this at Eq. (52) we get deduce the explicit form fo the $``$ product and it is given by
$$FG=\underset{n}{}(4\pi \alpha ^{})^nW_nB_n(F,G),$$
(53)
where $`W_n`$ are the weight function
$$W_n:=\frac{1}{(2\pi )^{2n}}\frac{1}{n!}d^nzd^n\overline{z}\underset{a=1}{\overset{n}{}}\left(\frac{1}{z_a1}\frac{1}{\overline{z}}\frac{1}{\overline{z}1}\frac{1}{z_a}\right)$$
(54)
and $`B_n(F,G)`$ are bi-differential operators
$$B_n(F,G):=\mathrm{\Theta }^{i_1j_1}\mathrm{\Theta }^{i_2j_2}\mathrm{}\mathrm{\Theta }^{i_nj_n}_{i_1}_{i_2}\mathrm{}_{i_n}F_{j_1}_{j_2}\mathrm{}_{j_n}G,$$
(55)
where $`\mathrm{\Theta }^{ij}`$ is that given in Eq. (40).
This is precisely the formula given by Kontsevich in for the $``$ product on any Poisson manifold. In this case the Poisson manifold is the spacetime $`X`$ and $`\mathrm{\Theta }_{ij}`$ is the Poisson-like structure.
Acknowledgements
We are very grateful to the organizing committee of the Third Workshop on Gravitation and Mathematical Physics for the hospitality. One of us ‘O. L-B. is supported by a CONACyT graduate fellowship.
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# Metal-insulator transition in 2D: resistance in the critical region
## I Introduction
After several years of intensive experimental and theoretical efforts (see, e.g., Ref. \[\] for an extensive bibliography), even the basic features of the phenomenon known as the “metal-insulator transition in two dimensions”(2D MIT) remain to be the subjects of ongoing discussions. Is this phenomenon a true “quantum phase transition”(QPT) \[\] or can it be understood in terms of conventional physics of disordered conductors \[\] ? This question is at the heart of the whole discussion. Recently, we wrote a paper \[\] in favor of the second possibility. In particular, we argued that
1. in the “metallic phase”, 2D systems seem to behave as a quite conventional disordered metal rather than a distinctly new state of matter and
2. it is possible to explain the anomalous behavior of the resistivity, $`\rho (T)`$, by an interplay of two temperature dependences: the one given by a metallic-like, quasiclassical (Drude) resistivity \[\], $`\rho _\mathrm{d}(T)`$, and the other one arising from quantum localization.
We do realize that ii) represents a rather naive approach, at least because it does not fully take into account the Coulomb interaction between electrons, whereas the most pronounced effects have been observed in systems with a priori strong electron-electron correlations. Nevertheless, this approach allows one to describe qualitatively, and even semi-quantitatively, most of the results on electron transport both in metallic phase and near the transition point. This suggests that it is a good starting point for developing an adequate theoretical understanding of the MIT in two dimensions.
Conclusion i) is based on the analysis of existing experimental data, as summarized in Refs. \[\]. It has been further supported by recent experimental papers \[\], which showed that the transport properties of Si/Ge, Si MOSFET and p-GaAs structures in the metallic regime can be successfully interpreted in conventional terms. Authors of Refs. \[\] have also identified the contribution of electron-electron interactions to the resistance (via measuring the temperature-dependent part of the Hall resistance) and found that it remains quite small even when parameter $`r_s`$ is rather large. There is also a number of recent theoretical papers \[\] which, although differing in a particular mechanism for the $`T`$-dependence of the Drude resistivity, share the general spirit of propositions i) and ii).
An alternative point of view, expressed, for example, in Refs. \[\], is that the observed MIT-like behavior indicates a true quantum phase transition between an insulator and a novel metallic phase. Motivated by the importance of this controversy, i.e., QPT versus propositions i) and ii), for the field of quantum transport in 2D and its possible relevance for other realizations of quantum phase transitions, we decided to analyze in detail the arguments on both sides of the QPT question. In this paper, we discuss recent and some of the previously published experimental results on the resistivity of Si MOSFETs in the vicinity of the MIT, along with some of the theoretical interpretations of these results. We conclude that that there is no convincing experimental evidence for the QPT in the existing data.
As one of the main questions in the field is the behavior of a 2D system in the limit $`T0`$, we begin our analysis with a short summary of recent results by Kravchenko and Klapwijk (KK) \[\], who measured the temperature dependence of the resistivity down to a bath temperature as low as $`T_{\mathrm{bath}}=35`$ mK. This paper reports measurements on a single Si MOSFET sample with peak mobility $`\mu _{\mathrm{peak}}=27,000`$ cm<sup>2</sup>/Vs in the range from 35 mK to 1.2 K at five different electron concentrations. Digitized data from Ref. \[\] is shown in Fig. 1 (curves 1-5).
In Ref. \[\], the following two points are emphasized:
1. For $`n=n_1=6.85\times 10^{10}`$cm<sup>-2</sup> and $`n=n_2=7.17\times 10^{10}`$cm<sup>-2</sup> (curves 1 and 2), the resistivity decreases as temperature increases $`(d\rho /dT<0)`$; for $`n=n_4=7.57\times 10^{10}`$cm<sup>-2</sup> and $`n=n_5=7.85\times 10^{10}`$cm<sup>-2</sup> (curves 4 and 5), the resistivity increases with temperature $`(d\rho /dT>0)`$. Based on these observations, the authors argue that curves 1 and 2 correspond to an insulating phase, whereas curves 4 and 5 demonstrate a metallic behavior. At $`n=n_3=7.25\times 10^{10}`$cm<sup>-2</sup> (curve 3), only a weak ($`\pm 5\%`$) temperature dependence was observed within the interval of bath temperature from 35 mK to 1 K. The authors conclude that $`n_3=7.25\times 10^{10}`$cm<sup>-2</sup> is the critical electron density for their sample, i.e., it corresponds to the metal-insulator transition point. Note that this conclusion is based entirely on the temperature independence of curve 3 over a relatively narrow temperature range.
2. Another observation, emphasized in Ref. \[\], is that in the range 0.1 K$`<T<`$ 0.4 K the metallic $`\rho (T)`$ \- dependences (curves 4, 5) appear to be nearly linear and thus different from the exponential temperature dependence observed in other experiments. KK stress that in this region the resistivity shows no insulating up-turn at the lowest temperature achieved.
KK consider the fact that they observed no temperature dependence of the resistivity of their sample at a certain concentration $`n=n_3`$ as a strong argument against our model \[\]. They point out that such a precise cancellation would require a special relation between $`\rho _d(T)`$ and the scaling $`\beta (\rho )`$-function. KK argue that this relation between two objects of entirely different origin “would be a remarkable coincidence”, and is therefore unlikely.
Based on the arguments listed above, authors of Ref. \[ \] conclude that their experiments are consistent with the existence of the zero-temperature quantum phase transition .
Note that a rather narrow range of densities was explored in Ref. \[\]. To show the development of the MIT-phenomenon over a wider range of densities and temperatures, we present in Fig. 2 the results of Ref. \[\] for a sample with a close value of $`\mu _{\mathrm{peak}}`$. The encircled region of the $`\{T,n\}`$ plane (region KK in this figure) indicates roughly the domain of parameters explored in Ref. \[\]. The temperature in Fig. 2 is normalized to the Fermi energy $`E_F`$, in order to demonstrate that resistivities corresponding to quite different Fermi energies, and therefore $`r_s`$-values, exhibit similar $`T`$-dependences. A strong (50-100%) drop in $`\rho `$ is still present for $`n20\times 10^{11}`$ cm<sup>-2</sup>, which corresponds to $`r_s2`$ and $`\rho 0.01h/e^2`$. Thus, contrary to the popular opinion, the resistance drop is not intrinsic only to the low density range (high $`r_s`$-values).
We would like to stress once again that in the region of high densities and small resistivities quantum interference corrections of any kind (weak localization, exchange and Hartree interactions, spin-orbit, etc.) would amount only to a few percent variation (of both signs) in $`\rho `$ with $`T`$. The phonon contribution to the resistivity of Si-MOSFETs can be shown to be negligible at least for $`n20\times 10^{11}`$ cm<sup>-2</sup> in the relevant $`T`$-range. Therefore, there must be some other mechanism of the $`T`$-dependence at work, neither of quantum interference nor of phonon nature. The resistance drop of non-phonon nature persists up to the highest temperature of these measurements, which is $`3540`$ K for each curve in Fig. 2 (for conversion of densities into the Fermi-energies, use the formula: $`E_F`$\[K\]$`=7.31\times n[10^{11}`$ cm$`{}_{}{}^{2}]`$).
## II Metal, Insulator, and Critical Point. General Discussion
Contrary to a common belief, a metallic phase and thus a metal-insulator transition can in principle occur in two dimensions. For example, in a hypothetical situation of non-interacting electrons spin-orbit coupling results in resistivity vanishing with temperature, provided that disorder is sufficiently weak (“weak antilocalization” \[\]). However, if disorder is strong enough, it leads to localization. The metallic state is stable with respect to weak and short-ranged electron-electron interactions. The Coulomb interaction between electrons for small $`r_s`$ overwhelms weak antilocalization and thus destroys the metallic state with zero residual resistivity. Nevertheless, the Coulomb interaction itself brings in an antilocalizing contribution to the resistivity (triplet channel). For small $`r_s`$, this contribution is smaller than the one from the singlet channel, and thus the net effect of interactions is the enhancement of localization. It may be the case that the triplet channel wins for larger $`r_s`$ and a metallic state becomes possible. (In fact, this scenario follows from the RG treatment \[\].) However, if such a state exists, it must be distinct from the conventional “disordered Fermi liquid”.
So far, the experimental evidence is in contrast with the existence of a distinctly new state (see Refs.\[\]). In fact, at higher densities the resistivity exhibits an insulating up-turn in agreement with the weak localization theory. We believe that in the existing experiments the true low-temperature asymptotic behavior has not yet been achieved for lower densities (in the vicinity of the transition). Therefore, one cannot interpret the crossover between $`d\rho /dT>0`$ and $`d\rho /dT<0`$ at finite temperatures as a proof for the existence of two distinct phases. Nevertheless, it is instructive to adopt such an interpretation for a moment and to compare the experimentally observed $`\rho (T)`$ with the commonly accepted picture of a quantum phase transition. Let us first recall what is known about the transition in 3D.
We start our analysis by formulating the definitions of metal and insulator. Arguing about definitions is not a too rational thing to do – a definition cannot be right or wrong. We simply point out that for the well-studied case of MIT in 3D, the commonly held definitions of both phases differ from those adopted by KK. In 3D systems, $`\rho _{\mathrm{insulator}}(T0)\mathrm{}`$, whereas $`\rho _{\mathrm{metal}}(T0)`$ is finite (see \[\] and references therein). In fact, $`d\rho /dT`$ is negative in both the metallic and insulating phases, provided that the system is close to the transition point and the temperature is low enough.
For non-interacting particles, the conventional scenario of the 3D MIT is well supported by the perturbation theory and renormalization group arguments \[\]. According to this scenario, in the metallic phase close to the transition $`\rho (T)`$ increases as $`T^\alpha `$ as temperature decreases. The increase of $`\rho (T)`$ saturates at $`TT_{\mathrm{sat}}`$. Upon approaching the critical point, $`T_{\mathrm{sat}}`$ tends to zero. Therefore, exactly at the critical density the $`\rho (T)`$-dependence is a power law: $`\rho (T)T^\alpha `$. In the insulating phase, the resistivity behaves as $`\mathrm{exp}(T_0/T)`$ (nearest neighbor hopping) or as $`\mathrm{exp}[(T_0/T)^\beta ]`$ (variable range hopping). The exponent $`\alpha `$ is determined by mechanisms of dephasing. As a result, it is, strictly speaking, non-universal. Let the dephasing rate scale with temperature as $`1/\tau _\phi T^p`$. Then assuming that (a) one-parameter scaling holds, (b) this parameter is the conductance, and (c) $`p=1`$, one arrives at $`\alpha =1/3`$. Note that the electron-electron interactions can change the value of $`\alpha `$.
Numerous experiments on 3D doped semiconductors confirmed the power-law behavior of $`\rho `$ in the critical region. However, there is still no consensus regarding the value of $`\alpha `$: both $`\alpha =1/2`$ \[\] and $`\alpha =1/3`$ \[\] have been reported. It is also possible that $`\alpha `$ depends on whether a semiconductor is compensated or not. (For the review of an MIT in 3D see Ref. \[\].)
Returning to the MIT in 2D, we note that the assumption of a temperature-independent resistivity at the critical density is as doubtful as the statement that $`\alpha =1/3`$ in 3D. Indeed, both of these points can be justified only within the one- parameter scaling. This scaling does not seem to apply universally even in the 3D case. For a 2D system of noninteracting electrons, this scaling predicts no metallic state at all. To get a chance to describe a metallic state, one needs to add more ingredients, e.g., electron-electron interactions, to the theory. Once we deal with a problem which is richer than localization in quenched disorder, there are no reasons to assume that the one-parameter scaling still holds (see, e.g., Ref. \[\]). Therefore, a temperature-independent critical resistivity is an assumption rather than a law of nature.
As an example, consider the following density- and temperature-dependences of the resistivity (which do not follow from any of the existing theories, but do not contradict to any of the commonly respected general principles either):
$$\rho (n,T)=\rho _0(n)+\rho _1(n)\left[\frac{T}{T_1}\right]^\alpha \mathrm{exp}\left[\frac{T_0(n)}{T}\right].$$
(1)
It is natural to define the critical concentration, $`n_c`$, as the concentration at which $`T_0(n)`$ changes sign: $`T_0(n_c)=0`$; for $`n>n_c`$ (metal), $`T_0(n)`$ is negative, whereas for $`n>n_c`$ (insulator), $`T_0(n)`$ is positive. It is also assumed that $`\rho _0(n)`$$`\rho _1(n)`$ and $`T_1(n)`$ are smooth functions of the density in the vicinity of the transition point $`n=n_c`$. According to Eq.(1), exactly at the critical point ($`n=n_c`$)
$$\rho (n_c,T)=\rho _0+\rho _1\left(\frac{T}{T_1}\right)^\alpha .$$
(2)
As we see, the critical resistivity is temperature-dependent and can even diverge as $`T0`$. At the same time, the resistivity saturates at $`\rho =\rho _0(n)`$ for $`T<T_0`$ in the metallic regime. (In our example, $`\rho (T)`$ has a maximum at $`T=T_0/\alpha `$. We do not think that such a maximum is a mandatory feature of the MIT in two dimensions.) We conclude that, at least in this example, a temperature-independent resistivity is a signature of a metal rather than of a critical point.
Regardless of this particular and rather artificial example, we notice that it is neither easy nor a straightforward task to determine the critical point of an MIT. Quantum phase transitions are zero temperature phenomena, whereas experiments are performed at finite $`T`$. Therefore, to determine the critical concentration, one does not have another choice but to analyze the data taken at the lowest accessible temperatures. This analysis should prove that $`\rho (T)`$ indeed acquires insulating exponential behavior as soon as the concentration gets lower than the apparent value of $`n_c`$.
We find it more appropriate to use the onset of the insulating exponential behavior in $`\rho (T)`$ rather than the vanishing of the derivative $`d\rho /dT`$ for an experimental identification of the critical point. We are not aware of any reasons to assign the meaning of the critical point to a density $`n_c^{}`$ at which the temperature dependence of the resistivity is least pronounced, as it was done by KK. Indeed, direct measurements of the two quantities, $`n_c^{}`$ (from the temperature independent “separatrix”) and $`n_c`$ (from the disappearance of the non-Ohmic hopping behavior) confirm their systematic difference for high-$`\mu `$ samples: $`n_c^{}`$ is larger than $`n_c`$ by 1-5% (see, e.g., Fig. 3 in Ref. \[\]). A similar difference arises also from the comparison of critical behaviors in the thermoelectric power and conductance \[\]. Note that four out of five electron densities in Ref. \[\], $`n=n_1n_4`$, fall into this ambiguous interval.
## III Experimental determination of the critical point.
We now turn to the experimental data of Ref. \[\]. The authors assume that $`n_c`$ equals to $`7.25\times 10^{10}`$cm<sup>-2</sup> (curve 3 in Fig. 1). According to the definition of the critical point, proposed in Sec. II, this assumption implies that at lower densities ($`n=n_2=7.17\times 10^{10}`$cm<sup>-2</sup> and $`n=n_1=6.85\times 10^{10}`$cm<sup>-2</sup>), the resistivity increases exponentially as temperature decreases. It turns out that one can fit neither $`\rho (T,n=n_2)`$ nor $`\rho (T,n=n_1)`$ with a simple exponential dependence. However, $`\rho (T,n=n_1)`$ can be satisfactorily approximated by
$$\rho _{\mathrm{exp}}(T,n)=\rho _m^{}\mathrm{exp}\left[\left(T_{0m}(n)/T\right)^m\right]$$
(3)
with either $`m=1/2`$ or $`m=1/3`$ (dashed line $`A_1`$ and solid line $`B_1`$ in Fig. 1, respectively). The quality of the fit can be seen from the deviations of the experimental data from the straight line in the inset to Fig. 1. Although curve 2 in Fiq. 1 does not seem to behave exponentially, one can still try to fit it by function (3) in two different ways:
1. one can assume that $`T_{0m}`$ is proportional to $`|nn_c|`$. (This dependence was observed in previous studies of the MIT in Si MOSFETs \[\]). Then, using the value of $`T_{0m}`$ for $`n=n_1`$, one determines $`T_{0p}`$ for $`n=n_2`$ by linear extrapolation, whereas $`\rho _m^{}`$ is found by fitting the experimental data. The results of this procedure for $`m=1/2`$ are shown in Fig. 1 (dashed-and-dotted line C<sub>2</sub>). It is quite clear that such a fit does not work (the same is true for $`m=1/3`$, not shown).
2. Alternatively, one can simply make the best fit of curve 2 by function (3) for $`m=1/2,1/3`$, treating $`T_{0m}`$ and $`\rho _m^{}`$ as independent parameters. Results of this fit are depicted in Fig. 1 by curves $`A_2`$ and $`B_2`$, respectively. It is customary to assume that in the vicinity of the critical point, $`T_{0m}`$ scales as $`|nn_c|^{z\nu }`$. Given the values of $`T_{0m}`$ for $`n=n_{1,2}`$, one can follow the assumption of KK that $`n_c=n_3`$ and find
$`z\nu ={\displaystyle \frac{\left[\mathrm{ln}T_{0m}(n_1)/T_{0m}(n_2)\right]}{\mathrm{ln}\left[(n_3n_1)/(n_3n_2)\right]}}=`$ (4)
$`=\{\begin{array}{cc}2.23\mathrm{for}m=1/2;& \\ 3.57\mathrm{for}m=1/3.& \end{array}`$ (7)
It is quite alarming that these values of the critical exponent are substantially different from $`z\nu =1.21.6`$ found in previous studies \[\].
It is also often assumed that in the vicinity of the critical point $`\rho (T,n)`$ should be symmetric with respect to reflection about the separatrix \[\]. To check if this relation works, we reflect curve 1 ($`|\delta n|=n_3n_2=0.4\times 10^{10}`$ cm<sup>-2</sup>) about the would-be separatrix (curve 3). The reflected curve is shown by semi-open circles. Had reflection symmetry worked, the reflected curve must have been located in between curves 4 ($`\delta n=0.32`$ cm<sup>-2</sup>) and 5 ($`\delta n=0.6`$ cm<sup>-2</sup>), closer to curve 4. In fact, the reflected curve crosses both curves 4 and 5. Thus by this criterion $`n_3`$ is not a critical density.
Unfortunately, for $`n=n_2`$ KK present the data only down to $`80`$ mK, whereas for the rest of the densities, both smaller $`(n_1)`$ and larger $`(n_3,n_4,n_5)`$ than $`n_2`$, is shown down to 35 mK. If observed, a substantial increase in $`\rho (T,n_2)`$ with cooling from $`80`$ mK to 35 mK would provide KK with a strong argument in favor of $`n_3`$ being close to the critical density.
Within the assumption that the data presented in Fig. 1 can at all be interpreted within a paradigm of a quantum phase transition, the contradictions demonstrated above suggest at least one of the following three conclusions:
1. the separatrix corresponds not to density $`n_3`$ but to some smaller density;
2. the critical exponent is not universal for a given material but depends on the sample preparation, geometry, etc.;
3. the exponential increase of the resistivity at density $`n_2`$ is suppressed due to overheating (we will discuss this option in Sec. V in more details).
Alternatively, one may conclude that the data of Ref. \[ \] do not provide evidence for a quantum phase transition between two distinct states of matter, but rather indicate a crossover between two different types of (relatively) high temperature behavior.
## IV Resistivity of Si MOSFETs in the critical region
As it has already been mentioned, any scenario of a quantum phase transition implies that for the densities close to the critical point $`n_c`$ the resistivity $`\rho (T,n)`$ at sufficiently high temperatures demonstrates a critical behavior. This means that the deviation of $`\rho (T,n)`$ from the separatrix $`\rho (T,n_c)`$ is small for small values of the critical parameter $`|nn_c|^{z\nu }/T`$. Keeping this in mind, we now recall the main results of earlier experiments on the resistivity of high mobility Si MOSFETs in the critical region and compare them with KK’s results.
Consider, for example, the resistivity of a sample with $`\mu _{\mathrm{peak}}=55,000`$ cm<sup>2</sup>/V s, reproduced from Ref. \[\] in Fig. 3. Assume for a moment that metallic and insulating behaviors can indeed be distinguished by the sign of the derivative $`d\rho /dT`$, as suggested by KK. At first glance, the application of this criterion to the region $`T<1`$ K in Fig. 3 is rather unambiguous. Indeed, $`n_c^{}`$ determined from the condition of $`T`$-independent $`\rho `$ corresponds to the dashed separatrix.
On the other hand, using the onset of activated and non-Ohmic conduction as a criterion for the transition, we find that the critical density $`n_c`$ is lower, $`n0.719\times 10^{11}`$ cm<sup>-2</sup> and corresponds to the “tilted” (6th) curve from the top. At higher temperatures, some of the “metallic” curves change their slopes and follow quite closely those “insulating” curves, which correspond to lower densities (see solid curves in Fig. 3.) The region of the $`\{T,n\}`$ plane, where such a behavior occurs, is defined as the critical region. Certainly, the density range of the critical region is narrower for lower temperatures.
As a rule, the resistivity depends on the temperature in this region. This feature, which can be found in almost any published data, is demonstrated in Figs. 2 and 3. For example, $`\rho (T,n)`$ in Fig. 3 increases by approximately a factor of 2, as $`T`$ changes from 3 K to 1 K, both for $`n=0.719\times 10^{11}`$ cm<sup>-2</sup> (the 6th curve from the top) and $`n=0.773\times 10^{11}`$ cm<sup>-2</sup> (the 7th curve).
In fact, the $`\rho (T)`$-dependence at higher temperatures is usually more pronounced in higher-mobility samples. This is illustrated by Fig. 4, which is reproduced from Ref. \[\].
The slope of the dashed line for the sample with $`\mu _{\mathrm{peak}}=71,000`$ cm<sup>2</sup>/Vs (Fig. 4a) is negative and its absolute value is about 10 times bigger than in Fig. 4b ($`\mu _{\mathrm{peak}}=33,000`$ cm<sup>2</sup>/Vs). The reduction of the slope with $`\mu _{\mathrm{peak}}`$ persists down to relatively low mobilities until eventually the slope changes sign. This is seen from Fig. 5a, reproduced from Ref. \[\], where the authors presented $`\rho (T)`$ for a sample with $`\mu _{\mathrm{peak}}`$ as small as $`5,000`$ cm<sup>2</sup>/Vs.
Even for this low-mobility sample, the $`\rho (T)`$-dependence at the highest density measured ($`n=54.4\times 10^{11}`$cm<sup>-2</sup>) shows no signs of saturation in the accessible temperature region (above 100 mK), see Fig. 5b. There is no consistent theoretical understanding of this $`\mu _{\mathrm{peak}}`$-dependence of the slope yet. Such an understanding will probably come together with the description of the “critical resistance”, which apparently changes from about $`0.3e^2/h`$ at low mobilities to several units of $`e^2/h`$ at high $`\mu _{\mathrm{peak}}`$\[\], (see also discussion in Refs. \[\]). On the other hand, it is natural that for the intermediate mobilities, the slope is small, because it should vanish somewhere in the range from 5,000 to 30,000 cm<sup>2</sup>/Vs. The mobility of KK’s sample is 27,000 cm<sup>2</sup>/Vs, i.e., right in this range. Therefore one should not be too surprised by the fact that in some range of the concentrations the temperature dependence of the resistivity is rather flat, as it was observed by KK.
Last but not least: a variety of the $`\rho (T)`$\- dependences in the critical region indicates by itself the absence of universality in observed MITs \[\] and raises doubts that Si MOSFETs undergo a genuine quantum phase transition.
## V “How Cold are the Electrons?”
Heating of electrons by the applied source-drain field as well as by external noise is a common problem in low-temperature transport measurements. This problem turns out to be especially serious in Si MOSFETS in the vicinity of the “metal-insulator transition” and at temperatures $`100`$ mK due to \*) weak temperature dependence of the resistivity in this range of densities and \**) weak electron-phonon coupling. As a result, the interpretation of resistivity measurements becomes rather ambiguous. For example, the temperature interval, in which the resistivity appears to be temperature-independent, may look to be much broader than it actually is. The goal of this section is to demonstrate the seriousness of this problem.
In macroscopic Si MOSFETs, as well as in bulk Si, electrons couple to phonons only via the deformation potential. As this coupling is rather weak, the electron temperature, $`T_e`$, may exceed substantially the bath temperature, $`T_{\mathrm{bath}}`$; it is also not easy to control $`T_e`$. To provide convincing data, one has to use a “thermometer”capable of measuring $`T_e`$ directly. Any observable which depends on $`T_e`$, other than the (zero-magnetic-field) resistivity itself, can be used for this purpose. In earlier studies on high mobility Si-MOSFETs, Refs. \[\], which reported results of the resistivity measurements down to $`T_{\mathrm{bath}}=1825`$ mK, the electron temperature was attempted to be measured via
i) the amplitude of Shubnikov - de Haas (ShdH) oscillations,
ii) temperature dependence of the hopping resistivity,
iii) linearity of $`IV`$-characteristics, and
iv) $`T_e`$-dependence of the phase relaxation time.
Unfortunately, results on $`T_e`$ in the up-to-date measurements of 2D MIT are neither fully consistent with each other nor convincing enough below $`T_{\mathrm{bath}}=250`$ mK.
No independent measurement of $`T_e`$ is reported in Ref. \[\]. It is only mentioned that the source-to-drain bias, $`U`$, was less than 200 $`\mu `$V “to ensure that the total power dissipated in the sample was less than $`10^{13}`$ W”. As it is not clear whether this is enough to prevent electrons from being overheated, it is worth discussing electron heating in Si MOSFETs.
Electrons are driven out of the equilibrium by the applied voltage and/or by external noise. Let $`P`$ be the power deposited in the electron system. In the stationary state, all this power leaves the system either with electrons (through contacts) or with phonons. The phonon mechanism dominates at higher temperatures. If $`T`$ is low enough, this mechanism can be neglected compared to electron out-diffusion, in which thermal conductivity of the electrons determines the heat balance. As our task is to estimate heating at $`T`$100 mK, we first discuss what happens without phonons.
Given the temperature increment $`\mathrm{\Delta }T=T_eT_{\mathrm{bath}}`$, the power, which is carried out by electrons through the leads, can be estimated as (see Refs. \[\] and also \[\] for recent discussion):
$$P=\left(\frac{2\pi }{e}\right)^2\frac{T_{\mathrm{bath}}\mathrm{\Delta }T}{R},$$
(8)
where $`R`$ is the resistance which differs from the resistivity by the aspect ratio. To obtain a lower bound estimate of the bath temperature $`T_\mathrm{h}`$ at which heating becomes strong, i.e., $`\mathrm{\Delta }TT_\mathrm{h}`$, we neglect heating due to external noise and assume that the main reason for heating is the source-to-drain bias, $`U`$. Accordingly, $`P=U^2/R`$. Expressing the ratio $`\mathrm{\Delta }T/T_{\mathrm{bath}}`$ through $`P`$ via Eq. (8), we obtain
$$\frac{\mathrm{\Delta }T}{T_{\mathrm{bath}}}=PR\left(\frac{e}{2\pi T_{\mathrm{bath}}}\right)^2=\left(\frac{eU}{2\pi T_{\mathrm{bath}}}\right)^2.$$
(9)
Strictly speaking, Eqs.(8,9) are valid only when $`\mathrm{\Delta }TT_{\mathrm{bath}}`$. Nevertheless, one can use them to estimate $`T_\mathrm{h}`$ as
$$T_\mathrm{h}[\mathrm{mK}]=1.8U[\mu \mathrm{V}].$$
(10)
For $`U200\mu `$V, Eq. (10) gives $`T_\mathrm{h}=360`$ mK. For lower bias, e.g., for $`U=50\mu `$V this estimate gives $`T_\mathrm{h}=90`$ mK, whereas the experimental dependence $`T_e(P)`$ \[\] taken at $`T_{\mathrm{bath}}=140`$ mK shows that at this bias electrons are already overheated: $`T_e=1.4T_{\mathrm{bath}}`$. Despite its simple form, Eq. (10) is quite general and universal. For example, it is also valid in the presence of Fermi-liquid interactions, which do not affect the Wiedemann-Franz law \[ \].
Taking into account the phonon mechanism of energy relaxation does not change our estimate of characteristic temperature $`T_\mathrm{h}`$ substantially. This is so because $`T_\mathrm{h}`$ is low enough for out-diffusion to dominate over phonon emission. Indeed, the electron-phonon energy loss rate for weak overheating can be written as
$$P_{\mathrm{ph}}=\frac{\pi ^2}{3}\nu _FLW\frac{T_{\mathrm{bath}}\mathrm{\Delta }T}{\tau _{\mathrm{e}\mathrm{ph}}},$$
(11)
where $`\nu _F`$ is the density of states at the Fermi level, $`L`$ and $`W`$ are the sample length and width, respectively, and $`\tau _{\mathrm{e}\mathrm{ph}}`$ is the electron-phonon relaxation time. We assume that electrons are coupled to phonons via the screened deformation potential. The corresponding relaxation time for $`T_{\mathrm{bath}}\mathrm{}k_Fs=3.65\sqrt{n}`$ \[K\] (where $`s`$ is the speed of sound, and $`n`$ in $`10^{11}`$cm<sup>-2</sup>) can be written as \[\]
$$\frac{1}{\tau _{\mathrm{e}\mathrm{ph}}}=\frac{120}{\pi }\frac{\mathrm{\Xi }^2m^{}k_F^2a_0^3}{\mathrm{}^2Ms}\left(\frac{T_{\mathrm{bath}}}{k_Fs}\right)^3\left(\frac{T_{\mathrm{bath}}}{\kappa s}\right)^2,$$
(12)
where $`\mathrm{\Xi }`$ is the deformation potential constant, $`M`$ is the atomic mass, $`a_0`$ is the effective lattice constant defined in such a way that $`M/a_0^3`$ is equal to the mass density, and $`\kappa `$ is the inverse screening length. Comparing (8) to (11), one finds that out-diffusion dominates, provided that $`T_{\mathrm{bath}}`$ is smaller than
$$T^{}=\frac{\pi }{5}\mathrm{}k_Fs\frac{h}{\rho e^2}𝒫r_s^2\frac{\mathrm{\Xi }}{\epsilon _F}\left(\frac{a_0}{L}\right)^2,$$
(13)
where $`Ms^2\left(\mathrm{}^2/m^{}a_0^2\right)^2/\mathrm{\Xi }^3`$ and $`𝒫\mathrm{}/m^{}sa_0`$. Substituting material constants for a Si MOSFET into (13), one obtains
$$T^{}[\mathrm{K}]=0.3\left(\rho n^2\right)^{1/5}L^{2/5},$$
(14)
where $`\rho `$, $`n`$, and $`L`$ are measured in $`h/e^2`$, $`10^{11}`$ cm<sup>-2</sup>, and mm, respectively. Using experimental results on electron heating in Si MOSFETs \[\], one estimates the observed crossover temperature as 0.7 K, whereas Eq. (14) gives $`T^{}=0.6`$ K for the experimental conditions of Ref. \[\]. For the conditions of another experiment on heating \[\], Eq. (14) gives $`T^{}=0.3`$ K, which is just the lowest temperature of this measurement. No clear crossover from the phonon to out-diffusion mechanisms was observed in Ref. \[\] for $`T0.3`$ K (although the $`T`$-dependence of the energy-loss rate does start to slow down at $`T1`$ K). This is again consistent with our estimate. We thus conclude that Eq. (14) is consistent with the experimental data and can serve as at least a lower bound for $`T^{}`$.
For typical values of $`\rho e^2/h`$, $`n10^{11}`$ cm<sup>-2</sup>, and $`L=0.11`$ mm, Eq. (14) gives $`T^{}0.30.75`$ K, which is of the same order as $`T_\mathrm{h}`$ estimated above as $`0.36`$ K. Note that in deriving (14) we assumed that $`\rho `$ was $`T`$-independent. Taking the metallic-like $`\rho (T)`$-dependence into account enhances out-diffusion at low temperatures and thus shifts $`T^{}`$ towards even higher values. Also, taking into account possible external noise would further increase the value of $`T_\mathrm{h}`$. The lowest bath temperature in Ref. \[\] is $`T_\mathrm{b}=35`$ mK. Using Eq. (10), we see that in order to prevent electron heating, the bias voltage has to be much smaller than 20 $`\mu `$V, which is ten times smaller than the upper boundary for $`U`$ indicated in Ref. \[\].
We thus see that there are serious reasons to doubt that the electron temperature in KK’s measurements was below $`100`$ mK.
## VI Temperature-Independent Resistivity
Now it is time to analyze the first of the two main arguments that KK brought in favor of the “true quantum phase transition”, namely, the temperature-independent resistivity at what they believe to be the critical point. More specifically, KK claim that they observed no $`T`$-dependence (within a 10% margin) in the interval $`35\mathrm{mK}<T<1`$ K. They analyzed the data in terms of our recent theory of Anderson localization by temperature-dependent disorder \[\] and came to the conclusion that within this theory a nearly constant $`\rho (T)`$ would imply a “remarkable coincidence”.
We begin our discussion of this issue with summarizing briefly the argument of KK. Afterwards, we present our theoretical counterarguments.
The theory of Ref. \[\] describes the $`T`$-dependence of the observable resistivity, $`\rho `$, in the presence of two factors: i) the $`T`$-dependence of the Drude resistivity, $`\rho _\mathrm{d}`$, and ii) Anderson localization. Each of these two factors brings a $`T`$-dependence of its own and the resultant $`\rho (T)`$-dependence is described by the following scaling equation
$$\left[\frac{1}{\beta (\rho )}\gamma \right]\frac{d\mathrm{ln}\rho }{d\mathrm{ln}T}=\frac{p}{2}+\left[\frac{1}{\beta (\rho _\mathrm{d})}1\right]\frac{d\mathrm{ln}\rho _\mathrm{d}}{d\mathrm{ln}T},$$
(15)
where $`\rho `$ is measured in units of $`h/e^2`$, $`\beta (x)`$ is the Gell-Mann–Low function, whereas $`p`$ and $`\gamma `$ parameterize the phase-breaking time as
$$\tau _\phi T^p\rho ^{12\gamma }.$$
(16)
The $`T`$-dependence of $`\rho _\mathrm{d}`$, entering the RHS of Eq. (15), leads to a variety of behaviors in $`\rho (T)`$. In Ref. \[\] we discussed in particular the situation when $`\rho (T)`$ exhibits a maximum at a rather high temperature $`T_{\mathrm{max}}`$ (close to the starting temperature of the RG flow).
Is a very slow variation of $`\rho (T)`$ over a wide $`T`$-interval permitted in this model? KK answer this question negatively. Their argument goes as follows: to have $`d\rho /dT=0`$ within some interval of temperature, one has to require that the RHS of Eq. (15) is equal to zero within this interval. This is possible only if the Drude conductivity $`\rho _\mathrm{d}`$ is related in a specific way both to the $`\beta `$-function and to $`\tau _\phi `$. Such an exact relation is unlikely, given the different origin of these quantities.
This argument sounds to be formally correct. Nevertheless, as we will demonstrate shortly, it is possible to achieve a very small variation of $`\rho (T)`$ within a wide temperature interval without imposing any constraints of this kind on $`\rho _\mathrm{d}`$. It turns out that the variation of $`\rho (T)`$ is small enough and the $`T`$-interval is wide enough to agree with the experiment. This possibility results from the fact that there is a whole family of the $`\rho _\mathrm{d}(T)`$-functions, which are parameterized by the electron density and some other parameters, such as the peak mobility. It is the freedom in fine-tuning these parameters that allows one to suppress the variation of the observable resistivity dramatically.
For the purposes of this section, we adopt the following working definitions, consistent with those used in Ref. \[\]: by “metal” (“insulator”) we understand the region in which $`d\rho /dT>0(<0)`$ or $`d\rho _\mathrm{d}/dT>0`$; by “critical region” we understand the region in which $`\rho 1`$.
As was discussed in Ref. \[\], a maximum in $`\rho (T)`$ results from a competition between metallic-like $`\rho _\mathrm{d}(T)`$ and localization effects, the latter being controlled both by $`\rho _\mathrm{d}(T)`$ and $`\tau _\phi `$. At higher $`T`$, when $`\rho _\mathrm{d}`$ is sufficiently large, localization can already be strong enough to ensure the negative sign of $`d\rho /dT`$ despite $`\rho _\mathrm{d}(T)`$ decreasing with temperature. However, as $`T`$ goes down, $`\rho _\mathrm{d}`$ decreases, localization weakens and $`d\rho /dT`$ changes its sign from negative to positive, thus a maximum in $`\rho (T)`$ occurs. This can happen provided that certain conditions are satisfied. In particular, one has to require that
$$\rho _\mathrm{d}(T_0)>\rho _c,$$
(17)
where $`T_0`$ is the temperature at which $`\tau _\phi `$ is comparable to the transport mean free time (the flow of Eq. (15) starts at $`T=T_0`$) and $`\rho _c`$ is some critical resistivity determined by a particular form of $`\rho _\mathrm{d}(T)`$. For example, if $`\rho _\mathrm{d}T^q`$, critical resistivity $`\rho _c`$ is a certain function of $`p/q`$ (see Eq. (9) of Ref. \[\]). In this case, there is another condition for the maximum in $`\rho (T)`$ to occur, namely, $`p>2q`$.
As the resistivity $`\rho _\mathrm{d}(T_0)`$ depends on the electron density $`n`$, it can be tuned by the gate voltage. This tuning can drive the system between the domains of parameters corresponding either to a maximum or no maximum in $`\rho (T)`$. Let us start at $`\rho _\mathrm{d}(T_0)>\rho _c`$, so that the maximum in $`\rho `$ is at $`T=T_{\mathrm{max}}`$. As $`\rho _\mathrm{d}`$ approaches $`\rho _c`$ from above, $`T_{\mathrm{max}}`$ remains finite, while the maximum flattens out and disappears as soon as $`\rho _\mathrm{d}(T_0)`$ approaches $`\rho _c`$. At this moment, both $`d\rho /dT`$ and $`d^2\rho /dT^2`$ vanish, and the variation of $`\rho `$ around $`T_{\mathrm{max}}`$ takes place solely due to higher ($`n3`$) terms in the Taylor expansion of $`\rho `$.
How large is the temperature interval, in which this variation does not exceed some given value? Denoting,
$`t`$ $``$ $`\mathrm{ln}T,`$ (18)
$`t_{\mathrm{max}}`$ $``$ $`\mathrm{ln}T_{\mathrm{max}},`$ (19)
$`\beta _\mathrm{d}`$ $``$ $`\beta [\rho _\mathrm{d}(t_{\mathrm{max}})],`$ (20)
$`y(t)`$ $``$ $`\mathrm{ln}\rho _\mathrm{d}(t),`$ (21)
$`\beta _\mathrm{d}^{}`$ $``$ $`d\beta (\rho )/d\mathrm{ln}\rho |_{\rho =\rho _\mathrm{d}},`$ (22)
$`\beta _\mathrm{d}^{\prime \prime }`$ $``$ $`d^2\beta (\rho )/d\mathrm{ln}\rho ^2|_{\rho =\rho _\mathrm{d}},`$ (23)
we obtain from Eq. (15)
$`{\displaystyle \frac{d\rho }{dT}}|_{T=T_{\mathrm{max}}}=0`$ $``$ $`\dot{y}(t_{\mathrm{max}})={\displaystyle \frac{p}{2}}{\displaystyle \frac{\beta _\mathrm{d}}{\beta _\mathrm{d}1}}`$ (25)
$`{\displaystyle \frac{d^2\rho }{dT^2}}|_{T=T_{\mathrm{max}}}=0`$ $``$ $`\ddot{y}(t_{\mathrm{max}})={\displaystyle \frac{\beta _\mathrm{d}^{}}{\beta _\mathrm{d}(\beta _\mathrm{d}1)}}\left[\dot{y}(t_{\mathrm{max}})\right]^2.`$ (26)
The third (logarithmic) derivative of $`\rho `$ is then found from Eq. (15) to be
$$\frac{d^3\mathrm{ln}\rho }{dt^3}|_{t=t_{\mathrm{max}}}=(1/A)\left[S\frac{\beta _\mathrm{d}1}{\beta _\mathrm{d}}\frac{d^3y}{dt^3}|_{t=t_{\mathrm{max}}}\right],$$
(27)
where
$$A1/\beta (\rho )\gamma $$
(28)
and
$`S`$ $``$ $`{\displaystyle \frac{2\beta _\mathrm{d}^2\beta _\mathrm{d}\beta _\mathrm{d}^{\prime \prime }}{\beta _\mathrm{d}}}^3\left[\dot{y}(t_{\mathrm{max}})\right]^3{\displaystyle \frac{3\beta _\mathrm{d}^{}}{\beta _\mathrm{d}^2}}\ddot{y}(t_{\mathrm{max}})\dot{y}(t_{\mathrm{max}})`$ (30)
$`=`$ $`{\displaystyle \frac{p^3}{8(\beta _\mathrm{d}1)^4}}\left[\left(2\beta _\mathrm{d}^{}_{}{}^{}2\beta _\mathrm{d}\beta _\mathrm{d}^{^{\prime \prime }}\right)(\beta _\mathrm{d}1)+3\beta _\mathrm{d}^{^{}}\right].`$ (31)
(In going from (30) to (31), we used (25,26)).
For a sufficiently narrow interval of $`t`$ around $`t_{\mathrm{max}}`$, it suffices to keep only the lowest (cubic) term in the Taylor expansion of $`\mathrm{ln}\rho (T)`$ as a function of $`\mathrm{ln}(T/T_{\mathrm{max}})`$:
$`\mathrm{ln}\left({\displaystyle \frac{\rho (t)}{\rho (t_{\mathrm{max}})}}\right)=`$ $`(1/6A)[S{\displaystyle \frac{\beta _\mathrm{d}1}{\beta _\mathrm{d}}}{\displaystyle \frac{d^3\mathrm{ln}\rho _\mathrm{d}}{dt^3}}|_{t=t_{\mathrm{max}}}]\times `$ (32)
$`\times \left(tt_{\mathrm{max}}\right)^3.`$ (33)
As $`y`$ is supposed to be a smooth function of $`t`$, one can estimate $`d^ny/dt^n`$ as $`\overline{y}/\overline{t}^n`$, where $`\overline{y}`$ is the typical value of $`y`$ in the interval $`\overline{t}`$. In the “critical region”, $`\rho 1`$, which means that $`\overline{y}1`$ as well, and hence $`\beta _\mathrm{d}\beta _\mathrm{d}^{}\beta _\mathrm{d}^{\prime \prime }1`$. Thus the two terms in the square brackets in (33) are of the same order and (33) reduces to
$$\mathrm{ln}\rho (t)\mathrm{ln}\rho (t_{\mathrm{max}})\left(S/6A\right)(tt_{\mathrm{max}})^3.$$
(34)
Let us now estimate the number $`D`$ of decades in temperature
$$T_{\mathrm{high}}/T_{\mathrm{low}}=10^D,$$
(35)
for which
$$|\mathrm{ln}\rho (t)\mathrm{ln}\rho (t_{\mathrm{max}})|10\%.$$
(36)
(KK found that their $`\rho `$ is $`T`$-independent with this accuracy.) Using an interpolation formula $`\beta (\rho )=\mathrm{ln}\left(1+2\rho /\pi \right)`$ \[\] and choosing $`\rho _\mathrm{d}(T_{\mathrm{max}})=\pi /2`$, we have
$$D2.3|A|^{1/3}/p.$$
(37)
For $`p=1`$ and $`|A|=1`$, we get $`D2.3`$, i.e., the resistivity stays within a 10%-margin for more than two decades in $`T`$. (Note that a numerical coefficient of the order of unity, which we neglected in deriving (34), would enter (37) only under the cubic root, and is thus unimportant.) Fig. 6 demonstrates $`\rho (T)`$ calculated from Eq. (15) for $`\rho _\mathrm{d}`$ taken from the model of Ref. \[\]:
$`\rho _\mathrm{d}`$ $`=`$ $`\rho _1+\overline{\rho }\left(T/T_0\right)^q`$ (41)
$`\times \{\begin{array}{cc}\left(1+|\delta |T_0/cT\right)^q,\mathrm{for}\delta 0,\mathrm{`}\mathrm{`}\mathrm{insulator}^{\prime \prime };\hfill & \\ e^{\delta T_0/T},\mathrm{for}\delta >0,\mathrm{`}\mathrm{`}\mathrm{metal}^{\prime \prime },\hfill & \end{array}`$
where $`\delta =\delta _0(nn_c)/n_c`$ characterizes the distance from the critical point and $`c1`$. The dimensionless coefficient $`\delta _0`$ depends on details of the model (cf. Eq.(11) of Ref. \[\]). For the purposes of the present paper, we view Eq. (41) as simply a phenomenological form for $`\rho _\mathrm{d}(T)`$, regardless of the model \[\] it was originally derived from. This form is consistent with the experiment, if $`\delta _01`$.
The value of $`\delta =0`$ corresponds to the central line in Fig. 6, which on this scale is essentially constant for more than two decades. Zooming in (cf. the inset in Fig. 6), one sees that in fact $`\rho (\delta =0,T)`$ stays within a 10%-margin over exactly three decades in $`T`$.
KK claim that the observed $`\rho (T)`$ stays within this margin over only 1.5 decades, which is smaller than even the conservative estimate (37). (As it was discussed in Sec. V, a more realistic estimate of KK’s $`T`$-interval is one decade). A decrease in $`\delta `$ by 1% ($`\delta =0.01`$) leads to a weakly insulating behavior: $`\rho (T)`$ increases by 30% as $`T`$ decreases by 1.5 decades. An increase in $`\delta `$ by 4% ($`\delta =0.04`$) leads to a well-pronounced metallic behavior: $`\rho (T)`$ drops by about a factor of two over 1.5 decades in $`T`$. These features are in quantitative agreement with the experiment \[\]. (For the sake of simplicity, we assumed that parameters $`\rho _1`$ and $`\overline{\rho }`$ do not depend on the density, i.e., on $`\delta `$. Taking these dependences into account should lead to even better agreement with the experiment.)
We emphasize that no fine tuning between the $`\beta `$-function and $`\rho _\mathrm{d}`$ was used. If it is the competition between the decrease in Drude resistivity with temperature and the localization effects that leads to the maximum in the observable resistivity, then the $`\rho (T)`$-dependence flattens out, as the density approaches the threshold for a maximum. It turns out that if the exponent $`p`$ is not too big (which is usually the case), this flattening suppresses the variation of $`\rho (T)`$ below the experimentally observable level over several decades in temperature. Even an oversimplified model \[\] which does not take into account, e.g., electron-electron interactions (except for as a possible phase-breaking mechanism), can easily produce almost constant resistivity in a temperature interval, which is two orders of magnitude wider than the experimental one. Therefore, we do not find it too much of a “remarkable coincidence ” that at some density the system demonstrates a fairly constant $`\rho `$.
We now turn to the data presented in Ref. \[\]. The authors emphasize strongly that at $`n=n_3=7.25\times 10^{11}`$ cm<sup>-2</sup> the resistivity is almost temperature independent in the temperature range 35 mK - 1 K. They also add that $`\rho (T)`$ decreases with $`T`$ for $`T1.8`$ K. It would be interesting to know what happens at intermediate temperatures: $`1\mathrm{K}<T<1.8`$ K. This point being unclear, KK fill in the missing temperature interval by combining their result with the one obtained on another system. Indeed, they write:
“In combination with the results of Ref. (Ref. \[\] of this paper–AMP) where the temperature-independent curve (with essentially the same value of resistivity) was observed in the temperature range 250 mk-1.8 K in another 2D system in silicon, we further allege that there is no observable $`T`$-dependence at $`n_s=n_c`$ in the temperature range 35 mk-1.8 K…”. (Ref. \[\], p. 3, second paragraph.)
Parenthetically, we note that the resistivities of combined curves differ by 35%. As far as we understand, temperature intervals of experiments on different systems are not additive parameters, and thus the procedure described above is not justified.
Summarizing this part of our discussion, we take the liberty to describe the experimental situation in the following way:
In a sample from the intermediate peak mobility range (where temperature dependence of $`\rho `$ is known to be quite weak at high $`T`$), one can find a density such that $`\rho `$ does not change for more than 10% within about one order of magnitude in $`T`$. This statement is a result of measurements on a single sample in the interval 100 mK$`<T<1`$ K (when a realistic estimation of the electron heating is taken into account). (It will be two samples if results of Ref. \[\] for 250mK$`<T<`$1.8K are taken into consideration.)
In our opinion, there are neither theoretical nor experimental reasons to believe that the density corresponding to the weakest $`\rho (T)`$-dependence is indeed the critical one. i.e., that $`\rho (T0)\mathrm{}`$ and $`\rho (T0)`$ is finite for lower and higher densities, respectively. The lack of a pronounced temperature dependence of $`\rho `$ in a limited range of $`T`$ does not signal any remarkable phenomenon and is well-described by a simple theory of Anderson localization in a temperature-dependent disorder (Ref. \[\]).
## VII On the apparent linear temperature dependence of $`\rho (T)`$
We now turn to point (2) of Ref. \[)\] regarding the functional form of the observed $`\rho (T)`$-dependence. According to KK, this form is a) linear and b) qualitatively distinct from those observed by other authors. Fig. 7 reproduces Fig. 3b of Ref. \[\], where the “almost linear dependence” is demonstrated. The data in Fig. 7 is data 5 of Fig. 1, displayed in truncated ($`T<400`$ K) and full ($`T<1200`$ mK) intervals. The truncated interval, in which $`\rho (T)`$ is supposed to be “almost linear”, corresponds to Fig. 3b of Ref. \[\].
We begin with an attempt to fit untruncated data 5 of Figs. 1, 7 in the whole interval $`T<1.2`$ K by an empirical expression \[\]
$$\rho (T)=\rho _0+\rho _1\mathrm{exp}\left(\frac{T_0}{T}\right),$$
(42)
which is a simplified version of Eq. (1) with $`\alpha =0`$. This expression is known to work reasonably well for the metallic phase not only in Si MOSFETs \[\] but also in other 2D systems exhibiting the MIT-like behavior \[ \].
As is seen from Fig. 7, this attempt is quite successful. We achieved more than just acceptable fit in a much broader temperature interval than the one in which $`\rho (T)`$ looks “almost linear”. Of course, it is not a much of an achievement to fit a smooth curve by a function with three free parameters (to demonstrate a flexibility in choosing the fitting parameters, Fig. 7 presents two sets of parameters that can be used). On the other hand, we see no reason to argue that the $`\rho (T)`$-dependence of Ref. \[ \] at $`n=n_4`$ and $`n=n_5`$ (see Fig. 1) is qualitatively different from those observed in other experiments.
Moreover, such “almost linear” $`\rho (T)`$-dependences are typical for low-mobility samples in the “metallic” regime. For instance, compare Fig. 7 with Fig. 5b, in which $`\rho (T)`$ for a low-mobility sample is plotted over a much wider temperature range (90 mK - 4 K), and for $`n`$ about 50% higher than what can be called a critical concentration. The similarity of these two plots is quite clear.
We find it quite unlikely that the quasi-linear $`\rho (T)`$-dependences of Figs. 1 and 7 (as well as Fig. 5) manifest an exotic non-fermi-liquid state. Taking into account the whole set of features of the metallic state in Si-MOSFETs (which we listed in Ref. \[\]), we conclude that these dependences must have a classical explanation (see the next paragraph for our definition of this term). For example, one cannot exclude that such behavior is caused by the temperature dependence of the screening length \[\] or by the temperature dependence of the charge traps’ population in the oxide \[\].
We conclude this discussion by remarking that there seems to be some confusion in recent literature on MIT in 2D (see, e.g., Ref. \[\]) regarding the distinction between “quantum” and “classical” behaviors. We believe that this confusion is mainly semantic. We call a regime “classical” as long as electron transport can be described in the framework of the Boltzmann theory. This does not necessarily mean that electrons obey Boltzmann (non-degenerate) statistics -in fact, the Fermi-liquid theory was originally formulated in terms of the Boltzmann equation \[\]. It also does not exclude that a scattering cross-section is described in terms of quantum mechanics. (A more common name for this regime is “semi-classical” \[\], although the semi-classical approximation may not necessarily be employed for the calculation of a scattering cross-section.) However, quantum-mechanical effects arising due to interference of electron waves, e.g., weak localization, cannot be described by the Boltzmann equation. We call the regime “quantum”, if these effects are important. Although we agree with the authors of Ref. \[ \] in that “the formation of a Fermi surface” is an “intrinsic quantum-mechanical effect”, this does not imply that the metallic-like temperature dependence of the resistivity, whose onset is correlated with the Fermi-surface formation, is of the quantum nature, in a sense of the definition given above.
## VIII At what temperatures weak localization may be observed
There is one more point to discuss regarding the metallic behavior of Si MOSFETs. In addition to the almost linear dependence with $`d\rho /dT>0`$, KK emphasize that no signs of weak localization (WL) was observed, although they went to temperatures as low as $`T0.008T_F`$, where $`T_F`$ is the Fermi temperature. Giving the authors the benefit of the doubt, we assume the most favorable scenario in which electrons are not overheated. Nevertheless, under this assumption the apparent absence of WL is not as surprising as KK present it. In fact, we believe that even the bath temperature was too high to observe the localization upturn in the $`T`$-dependence of the resistivity. Indeed, experimental studies of high-mobility Si MOSFETs in the wide range of electron densities \[ \] show that WL effects are well-pronounced in the (perpendicular) magnetic field dependence of the resistivity. The observed magnetoresistance is in a quantitative agreement with the conventional theory \[\], even for densities rather close to the critical one. On the other hand, even for $`nn_c`$ observation of the WL effect in the zero-field $`\rho (T)`$-dependence \[\] (as well as the effects of electron-electron interactions) requires a precision of the order of $`\delta \rho (T)/\rho (T)(1/\pi )(e^2/h)\rho 0.1\%`$.
It turns out that the observed dependence in the range of “high temperatures”, $`T=(0.10.5)E_F`$, is too strong to be explained by WL and quantum interaction effects, and can be attributed to the $`T`$-dependence of the Drude resistivity. However, the slope $`|d\rho /dT|`$ decreases with $`T`$. For relatively large densities, one can reach the temperature region, $`T<T_Q`$, in which the slope becomes comparable to that of the WL correction, $`|d\rho /dT|_{WL}(e^2/\pi )\rho ^2T`$. In this region, quantum interference should contribute significantly to the $`\rho (T)`$-dependence. Authors of Ref. \[\] found an empirical relation: $`T_Q0.007E_F`$ (see Fig.(2)). As $`E_F`$ is proportional to $`n`$, it becomes progressively more challenging to cool electrons below $`T_Q`$ as $`n_c`$ is approached. $`E_F5.5`$ K for the highest electron density reported by KK and thus $`T_Q=40`$ mK. Therefore, in order just to reach the high temperature edge of the quantum transport region, where WL effects are seen, electrons should be cooled below 40 mK. The localization upturn may be expected to occur only at temperatures substantially smaller than even this one. In other words, the lowest temperature reported by KK (35 mK) is still too high for localization effects to be observed.
Why WL manifests itself so differently in the magnetoresistance and in the $`\rho (T)`$? This can be explained naturally by assuming that the Drude resistivity has a pronounced $`T`$-dependence, which masks quantum corrections. On the other hand, $`\rho _\mathrm{d}`$ is not expected to vary substantially in weak magnetic fields, thus WL can be seen in magnetoresistance.
## IX Conclusions
Some concluding remarks:
1. the apparent metal-insulator transition in 2D is a very interesting and unexpected phenomenon. Although we have limited our discussion to Si MOSFETs, which exhibit the strongest, among other 2D systems, “anomalous” metallic behavior, there is a whole variety of interesting effects not only in transport but also thermodynamic properties, e.g., compressibility \[\]), observed in various 2D heterostructures. Whether all these observations have a universal explanation or there is a number of different, system-specific mechanisms at work, remains to be seen.
2. we do not believe that the experimental results accumulated up-to-date provide a convincing evidence for this phenomenon being a true quantum phase transition between two distinct states of matter. More precisely, there is no evidence that the “metallic phase” (as defined by the sign of $`d\rho /dT`$) is a new state of matter. On the contrary, a whole set of features–see Refs.\[\]–is consistent with a conventional behavior of a disordered Fermi liquid. This statement does not negate i). There is still no consensus on the origin of the effect, and this problem requires the most serious and intensive investigations.
3. For a successful resolution of the problem, it is probably not sufficient to concentrate on the temperature dependence of the resistivity in the vicinity of the transition. Such features as a weak temperature dependence of the resistivity at a certain electron density in a limited (though large) temperature interval allow for different interpretations and, unfortunately, can not provide an unambiguous information about the zero-$`T`$ state of the system. In particular, it is impossible to determine even the critical density without making rather arbitrary assumptions.
4. There is a serious experimental difficulty that up to now prevented a substantial increase of the temperature interval, in which the resistivity can reliably be measured in Si-MOSFETs. It turns out that in existing samples phonons can hardly cool the electron gas below a fraction of a Kelvin. On the other hand, the systems are quite noisy, and thus the applied voltage cannot be reduced much below the currently used level. As a result, even in the absence of non-equilibrium noise, one cannot neglect electron heating at $`T100`$ mK.
###### Acknowledgements.
The work at Princeton University was supported by ARO MURI DAAG55-98-1-0270. D. L. M. acknowledges the financial support from NSF DMR-970338 and from the Research Corporation Innovation Award (RI0082). V. M. P. acknowledges the support from RFBR, INTAS, NWO, NATO (PST.CLG.976208) and Programs “Physics of solid-state nanostructures”, “Statistical Physics”, and “Integration”. We would like to thank the Center of Higher Studies (Oslo, Norway) and Institut für Halbleiterphysik Johannes Kepler Universität (Linz, Austria), where parts of this work were done. We are grateful to M. Reizer for numerous illuminating discussions and to A. F. Hebard, R. Fletcher, C. Marcus, G. W. Martin, and A. V. Varlamov for critical reading of the manuscript and valuable comments. We also appreciate the assistance of G. W. Martin in the manuscript preparation.
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# 1 Introduction
## 1 Introduction
Dark Matter (DM) is one of the intriguing problems in particle physics and cosmology. Several types of stable (or quasi stable) particles have been proposed as condidates for the cold DM. Among these condidates one has the nuclearites (strange quark matter), which are agglomerates of quarks u,d and s . Another example of cold DM candidate, coming from theories beyond the standard Model of particle physics, is the lightest supersymmetric particle (LSP) .
In theories where scalar fields carry a conserved global quantum number, $`Q`$, there may exist non-topological solitons which are stabilized by global charge conservation. They act like homogenous balls of matter, with $`Q`$ playing the role of the quantum number; Coleman called this type of matter Q-balls .
The conditions for the existence of absolutely stable Q-balls are satisfied in supersymmetric theories with low energy supersymmetry breaking<sup>1</sup><sup>1</sup>1According to ref. , abelian non-topological solitons with baryon and/or lepton quantum numbers naturally appear in the spectrum of the Minimal Supersymmetric Standard Model. . The role of conserved quantum number is played by the baryon number (The same reasoning applies to the lepton number for the sleptonic Q-balls). Q-balls can be considered like coherent states of squarks, sleptons and Higgs fields. Under certain assumptions about the internal self interaction of these particles and fields the Q-balls could be absolutely stable .
In this note we summarize the main physical and astrophysical properties of Q-balls, their interactions with matter, their energy losses in detectors, and the possibility of traversing the earth to reach any detector at the underground Gran Sasso Laboratory (Italy), for example MACRO . We neglect the possibility of : $`(i)`$ electromagnetic radiation emitted by Q-balls of high $`\beta `$, $`(ii)`$ strong interactions of Q-balls in the upper atmosphere capable of destroying the Q-ball, and $`(iii)`$ the possibility that neutral Q-balls (SENS) may become charged (SECS) and viceversa .
## 2 Properties of Q-balls
Q-balls could have been produced in the Early Universe, and may contribute to the cold DM. Several mechanisms could have lead to the formation of Q-balls in the Early Universe. They may have originated in the course of a phase transition, which is sometimes called “solitogenesis”, or they could have been produced via fusion processes, reminiscent of the Big Bang nucleosynthesis, which have been called “ solitosynthesis”; small Q-balls could be pair-produced in very high energy collisions (at high temperatures) .
The astrophysical consequenses of Q-balls in many ways resemble those of strange quark matter, nuclearites.
One of the peculiarities of Q-balls is that their mass grows as $`Q^{3/4}`$, while for nuclearites the mass grows linearly with the baryon number .
In the squark bag there is an almost uniform potential $`U(\varphi )`$, which may be taken as $`U(\varphi )M_S^4=`$constant in SUSY theories with low energy supersymmetry breaking . For large scalar $`\varphi `$, the mass $`M_Q`$ and radius $`R_Q`$ of Q-balls with baryon number $`Q`$ are given by :
$$M_Q=\frac{4\pi \sqrt{2}}{3}M_SQ^{3/4}5924\left(\frac{M_S}{1TeV}\right)Q^{3/4}(GeV)$$
$$1.2\times 10^{20}\left(\frac{M_S}{1TeV}\right)Q^{3/4}(g)$$
(1)
$$R_Q=\frac{1}{\sqrt{2}}M_S^1Q^{1/4}1.4\times 10^{17}\left(\frac{1TeV}{M_S}\right)Q^{1/4}(cm)$$
(2)
The parameter $`M_S`$ is the energy scale of the SUSY breaking symmetry. The first parts of formulae (1), (2) are in $`\mathrm{}=c=1`$ units; the last parts are in cgs units. A stability condition of the stable Q-ball mass $`M_Q`$ is related to the proton mass $`M_p`$ by
$$M_Q<QM_p$$
(3)
From Eq. 1 and Eq. 3 one has the stability constraint
$$Q>\frac{M_Q}{M_p}=1.6\times 10^{15}\left(\frac{M_S}{1TeV}\right)^4$$
(4)
In Fig. 1 the allowed region for stable Q-balls is indicated as a shaded region in the plane (Q, $`M_S`$). The dashed lines indicate the Q-ball number as function of $`M_S`$ assuming different Q-ball masses $`M_Q`$; the dotted lines represent Q($`M_S`$) for different Q-balls radii (R).
Relic Q-balls are expected to concentrate in galactic halos and to move at the typical galactic velocity $`v=\beta c10^3c`$. Assuming that Q-balls constitute the cold galactic dark matter with $`\rho _{DM}0.3GeV/cm^3`$, their number density is
$$N_Q\frac{\rho _{DM}}{M_Q}\rho _{DM}\left(\frac{3}{4\pi \sqrt{2}}\right)Q^{3/4}M_S^1$$
(5)
$$5\times 10^5Q^{3/4}\left(\frac{1TeV}{M_S}\right)cm^3$$
The corresponding flux is
$$\varphi \frac{1}{4\pi }N_Qv=\frac{\rho _{DM}v}{4\pi M_Q}$$
$$\frac{1}{4\pi }\rho _{DM}\left(\frac{4\pi \sqrt{2}}{3}\right)^1Q^{3/4}M_S^1v$$
(6)
$$10^2Q^{3/4}\left(\frac{M_S}{1TeV}\right)^1v(cm^2s^1sr^1)$$
In Fig. 2 we present the dependence of the Q-ball radius $`R_Q`$ from its mass $`M_Q`$ for two values of the $`M_S`$ parameter. Q-balls with $`M_Q10^8`$ GeV are unstable; Q-balls with $`M_Q10^{25}`$ GeV should be very rare.
In Fig. 3 we present the Q-ball flux versus Q-ball mass assuming that the density of the DM is $`\rho =0.3GeV/cm^3`$, that all of it is in the form of Q-balls and that $`\beta 10^3`$. If Q-balls are part of the DM, their flux cannot be higher than that of the top dotted-dashed line in Fig. 3.
When traversing normal matter, Q-balls could form some sort of bound states with quarks (which remain in their outer layers), acquiring a positive electric charge. Electrons around the Q-ball can be captured through the reaction $`ued\nu _e`$ or could leave the Q-ball charged if the rate of this reaction is small, otherwise they would neutralise it .
Q-balls can be classified in two classes: SECS (Supersymmetric Electrically Charged Solitons) and SENS (Supersymmetric Electrically Neutral Solitons). The interactions of Q-balls with matter and their detection differ drastically for SECS or SENS .
## 3 Interaction with matter
Q-balls could have velocities typical of objects bound to the galaxy, $`10^4<\beta <10^2`$. Thus we do not expect relativistic Q-balls. For reasons of completeness we estimate the energy losses even for relativistic Q-balls, which we indicate as dotted lines in Figs. 3, 4.
### 3.1 Interactions with matter of Q-balls type SECS
SECS are Q-balls with a net positive electric charge which tends to be mainly in the outer layer. The charge of SECS originates from the unequal rates of absorption in the condensate. The positive electric charge could be of one unit up to several tens. This positive electric charge may be neutralized by a surrounding cloud of electrons.
For small size Q-balls the positive charge interacts with matter (electrons and nuclei) via elastic or quasi elastic collisions. The cross section is similar to the Bohr cross section of hydrogenoid atoms :
$$\sigma =\pi a_0^210^{16}cm^2$$
(7)
where $`a_0`$ is the Bohr radius. The formula is valid for $`R_Qa_0`$, which happens for $`Q2.7\times \left(\frac{M_S}{1TeV}\right)^4`$. For $`Q2.7\times \left(\frac{M_S}{1TeV}\right)^4`$ the Q-ball radius is larger than $`a_0`$, and it increases with the Q-ball mass. Thus electrons could find themselves inside the Q-ball, for large Q-ball masses. One could assume that in such conditions the electronic capture is favoured and thus the SECS tend to become SENS.
The main energy losses of SECS passing through matter with velocities in the range $`10^4<\beta <10^2`$ are due to the interaction of the SECS positive charge: $`(i)`$ with the nuclei (nuclear contribution), and $`(ii)`$ with the electrons of the traversed medium (electronic contribution). The total energy loss is the sum of the two contributions.
SECS could cause the catalysis of proton decay, but only in the case for
$$R_Q\frac{2Z_Qe^2}{m_ev^2}$$
(8)
where $`v`$ and $`Z_Q`$ are respectively the velocity and the positive electric charge of the SECS, $`m_e`$ is the electron mass (see Fig. 2). For $`\beta =10^3,Z_Q=1`$ this corresponds to radii larger than $`10^3`$ fm and masses larger than 1 mg, which are very large values. The probability of catalysis by large mass SECS is reduced by the presence of their electric charge.
Electronic losses of SECS: The electronic contribution to the energy loss of SECS may be computed with the following formula
$$\frac{dE}{dx}=\frac{8\pi a_0e^2\beta }{\alpha }\frac{Z_Q^{7/6}N_e}{(Z_Q^{2/3}+Z^{2/3})^{3/2}}forZ_Q1$$
(9)
where $`\alpha `$ is the fine structure constant, $`\beta =v/c`$, $`Z_Q`$ is the positive charge of SECS, $`Z`$ is the atomic number of the medium and $`N_e`$ is the density of electrons in the medium. Electronic losses dominate for $`\beta >10^4`$ (see Figs 3 and 4 for the case of the dE/dx in the Earth).
Nuclear losses of SECS: The nuclear contribution to the energy loss of SECS is due to the interaction of the SECS positive charge with the nuclei of the medium and it is given by
$$\frac{dE}{dx}=\frac{\pi a^2\gamma NE}{ϵ}S_n(ϵ)$$
(10)
where
$$S_n(ϵ)\frac{0.56ln(1.2ϵ)}{1.2ϵ(1.2ϵ)^{0.63}},ϵ=\frac{aME}{Z_QZe^2M_Q}$$
(11)
and
$$a=\frac{0.885a_0}{(\sqrt{Z_Q}+\sqrt{Z})^{2/3}},\gamma =\frac{4M}{M_Q}$$
(12)
$`M_Q`$ is the mass of the incident Q-ball; $`M`$ is the mass of the target nucleus; $`Z_Qe`$ and $`Ze`$ are their electric charges; $`a`$ is the screening radius and $`a_0`$ is the Bohr radius; we assume that $`M_Q>>M`$. Nuclear losses dominates for $`\beta 10^4`$ (see Figs. 3 and 4).
The energy losses of SECS in the Earth mantle and Earth core: The energy losses of SECS in the Earth mantle and Earth core have been computed for different $`\beta `$-regions and for different charges of the Q-ball core, employing the same general procedures used for computing the energy losses in the Earth of magnetic monopoles and of nuclearites . The results of the calculations are presented in Figs. 3 and 4; the dashed lines indicate interpolations.
Considering the energy losses discussed above, we have computed for a specific velocity ($`v=250km/s`$) the angular acceptance ($`\mathrm{\Omega }/4\pi `$) of the MACRO detector (located at an average depth of 3700 m.w.e.) for SECS (with $`Z_Q=1`$), see Fig. 5. Notice that for $`M_Q10^{12}GeV`$ the angular acceptance is $`2\pi `$ (corresponding to SECS coming only from above), while for $`M_Q10^{22}GeV`$ the geometrical acceptance is $`4\pi `$.
Fig. 6 shows the accessible regions in the plane ($`M,\beta `$) for Q-balls type SECS with $`Z_Q=1`$, coming from above (to the right of the dashed curve) and from below (to the right of the solid line), respectively. Notice that SECS coming from above (below) may reach MACRO only il they have masses larger than $`10^{13}`$ GeV ($`10^{18}`$ GeV).
### 3.2 Interactions of Q-balls type SENS
The Q-ball interior of SENS is characterized by a large Vacuum Expectation Value (VEV) of squarks, and may be of sleptons and Higgs fields <sup>2</sup><sup>2</sup>2Also SECS have in their interior a large VEV of squarks, sleptons and Higgs fields; in addition they have a positive electric charge, and thus a coulomb potential which reduces the possibility of the catalysis of proton decay. . The $`SU(3)_c`$ symmetry is broken and deconfinement takes place inside the Q-ball.
If a nucleon enters the region of deconfinement, it dissociates into three quarks, some of which may then become absorbed in the supersymmetric condensate \[7-8\]. The nucleon enters the $`\stackrel{~}{q}`$ condensate and giving rise to processes like
$$qq\stackrel{~}{q}\stackrel{~}{q}$$
(13)
In practice one has
$$(Q)+Nucleon(Q+1)+pions$$
(14)
and less probably,
$$(Q)+Nucleon(Q+1)+kaons$$
(15)
The pions carry the electric charge of absorbed nucleon (for example $`\pi ^+\pi ^0`$).
If it is assumed that the energy released in (14) and (15) is the same as in typical hadronic processes (about 1 GeV per nucleon), this energy is carried by 2 or 3 pions (or kaons). The cross section for reactions (13) and (14) are determined by the Q-ball radius $`R_Q`$
$$\sigma =\pi R_Q^2=\frac{16\pi ^2}{9}M_Q^2Q^26\times 10^{34}Q^{1/2}\left(\frac{1TeV}{M_S}\right)^2cm^2$$
(16)
The corresponding mean free path $`\lambda `$ is
$$\lambda =\frac{1}{\sigma n}$$
(17)
where $`n`$ is the number of atoms per $`cm^3`$ in the traversed material. According to refs. \[6-8\] the energy loss of SENS moving with velocities in the range $`10^4<\beta <10^2`$ is constant and is given by
$$\frac{dE}{dx}\frac{\zeta }{\lambda }=\sigma n\zeta 6\times 10^{34}Q^{1/2}n\zeta $$
(18)
where $`\zeta =1(\rho /1gcm^3)GeV`$ is the energy released in one reaction. The energy loss of SENS is due to the energy released from the absorbed nucleon mass. Large mass SENS lose a small fraction of their kinetic energy and are able to traverse the earth without attenuation for all masses of our interest. Fig. 7 shows the energy loss of SENS with $`10^4\beta 10^2`$ plotted versus radius for $`M_S=1TeV`$.
Fig. 8 shows the mean free path $`\lambda `$ of Q-balls type SENS versus Q-ball number $`Q`$ for $`M_S=1TeV`$.
## 4 Sensitivity to SECS of scintillators, streamer tubes and nuclear track detectors
In this section we discuss the response of liquid scintillators, of streamer tubes and of nuclear track detectors (in particular the MACRO CR39) to Q-balls, assuming that for $`10^4<\beta <10^2`$ the electric charge of SECS is constant.
### 4.1 Light yield of Q-balls type SECS
The MACRO liquid sintillator has a density of 0.86 g/cm<sup>3</sup> and it is made of 96.4 $`\%`$ of mineral oil, 3.6 $`\%`$ of pseudocumene, 1.44 g/l of PPO, 1.44 mg/l of bis-MSB and 40 mg/l of antioxydant .
For SECS we distinguish two contributions to the light yield in scintillators: the primary yight yield and the secondary light yield.
The primary light yield: it is due to the direct excitation (and ionization that occurs only for $`\beta >10^3`$) produced by the SECS in the medium. The energy loss in the MACRO liquid scintillator is computed from the energy loss of protons in hydrogen and carbon
$$\left(\frac{dE}{dx}\right)_{SECS}=\frac{1}{14}\left[2\left(\frac{dE}{dx}\right)_H+12\left(\frac{dE}{dx}\right)_C\right]=SP=\frac{SL\times SH}{SL+SH}$$
(19)
where $`SP`$ is the stopping power of SECS, which reduces to the lowest stopping power (SL) at low $`\beta `$ and to the highest stopping power (SH) at high $`\beta `$. The stopping power coincides with the Bethe-Bloch formula for electric energy losses at relatively high velocities.
1. For electric charge $`q=1e`$ the energy loss of SECS in hydrogen and carbon is computed from adding an exponential factor coming from the experimental data (see ref. ) on slow protons, see ref. .
i) For $`10^5<\beta <5\times 10^3`$ we have the following formula
$$\left(\frac{dE}{dx}\right)_{SECS}=1.3\times 10^5\beta [1\mathrm{exp}\left(\frac{\beta }{7\times 10^4}\right)^2]\frac{MeV}{cm}$$
(20)
ii) For $`5\times 10^3<\beta <10^2`$ we used the following formula
$$SP=SP_H+SP_C=\left(\frac{dE}{dx}\right)_{SECS}$$
(21)
where
$$SP_H=\frac{SL_H\times SH_H}{SL_H+SH_H}$$
(22)
$$SP_C=\frac{SL_C\times SH_C}{SL_C+SH_C}$$
(23)
and
$$SL=A_1E^{0.45},SH=A_2Ln\left(1+\frac{A_3}{E}+A_4E\right)$$
(24)
where ($`A_{i=1,4}`$) are estimated in ref. and $`E`$ is the energy of a proton with velocity $`\beta `$.
2. For SECS with electric charge $`q=Z_1e`$ the energy losses for $`10^5<\beta <10^2`$ are given by
$$\left(\frac{dE}{dx}\right)_{SECS}=\frac{8\pi e^2a_0\beta }{\alpha }\frac{Z_Q^{7/6}n_e}{(Z_Q^{2/3}+Z^{2/3})^{3/2}}[1\mathrm{exp}(\frac{\beta }{7\times 10^4})^2]$$
(25)
where $`Z`$ is the atomic number of the target atom, $`n_e`$ the density of electrons and $`\alpha `$ is the fine structure constant.
The primary light yield of SECS is given by
$$\left(\frac{dL}{dx}\right)_{SECS}=A\left[\frac{1}{1+AB\frac{dE}{dx}}\right]\frac{dE}{dx}$$
(26)
where $`dE/dx`$ is the total energy loss of SECS; $`A`$ is a conversion constant of the energy losses in photons (light yield) and $`B`$ is the parameter describing the saturation of the light yield; both parameters depend only on the velocity of SECS . For example for the $`\beta `$-range, $`5\times 10^5<\beta <10^3`$, $`A=0.067`$ and $`B=0.66`$ cm/MeV.
The secondary light yield arises from recoiling particles: we consider the elastic or quasi-elastic recoil of hydrogen and carbon nuclei. The light yield $`L_p`$ from a hydrogen or carbon nucleus of given initial energy $`E`$ is computed as
$$L_p(E)=_0^E\frac{dL}{dx}(ϵ)S_{tot}^1𝑑ϵ$$
(27)
where $`S_{tot}`$ is the sum of electronic and nuclear energy losses, $`ϵ`$ is the excitation energy of the outer shell electrons. The nuclear energy losses are given in ref. . The secondary light yield is then
$$\left(\frac{dL}{dx}\right)_{\text{secondary}}=N_0^{T_m}L_p(T)\frac{d\sigma }{dT}𝑑T$$
(28)
where $`N`$ is the number density of nuclei in the medium $`T_m`$ is the maximum energy transferred and $`\frac{d\sigma }{dT}`$ is the differential scattering cross section, given in ref. .
Fig. 9 shows the light yield of SECS in the MACRO liquid scintillator; for reference are also indicated the light yields of MMs with $`g=g_D`$ and fast muons.
### 4.2 Energy losses of SECS in streamer tubes
The composition of the gas in the MACRO limited streamer tubes is 73% helium and 27% n-pentane, in volume . The pressure is about one atmosphere and the resulting density is low (in comparison with the density of the other detectors): $`\rho _{gas}=0.856\text{mg/cm}^3`$. The energy losses of MMs in the streamer tubes have been discussed in ref. .
The ionization energy losses of SECS in the MACRO streamer tubes for $`10^3<\beta <10^2`$ are computed with the same general procedure used for scintillators, using the density and the chemical composition of streamer tubes.
For $`q=13e`$ the ionization energy losses are calculated as in ref. (we have omitted the exponential factor which takes into account the energy gap in organic scintillators).
The threshold for ionizing n-pentane occurs for $`\beta 2\times 10^3`$.
Fig. 10 shows the ionization energy losses of SECS with electric charges $`q=1`$ and $`q=13`$ plotted versus $`\beta `$.
### 4.3 Restricted energy losses of SECS in the nuclear track detector CR39
The quantity of interest for nuclear track detectors is the Restricted Energy Loss (REL), that is, the energy deposited within $``$ 100 Å diameter from the track.
The REL in CR39 has already been computed for MMs of $`g=g_D`$ and $`g=3g_D`$ and for dyons with $`q=e`$, $`g=g_D`$ in ref. . We have checked these calculations and extended them to other cases of interest, see ref. .
The chemical composition of CR39 is $`(\text{C}_{12}\text{H}_{18}\text{O}_7)_n`$, and the density is 1.31 $`\text{g/cm}^3`$. For the computation of the REL, only energy transfers to atoms larger than $`12`$ eV are considered, because it is estimated that $`12`$ eV are necessary to break the molecular bonds .
At low velocities ($`3\times 10^5<\beta <10^2`$) there are two contributions to REL: the ionization and the atomic recoil contributions.
The ionization contribution, which becomes important only for $`\beta >2\times 10^3`$, was computed with Ziegler’s fit to the experimental data .
The atomic recoil contribution, which is important to REL for low $`\beta `$ values was calculated using the interaction potential between an atom and a SECS
$$V(r)=\frac{Z_QZe^2}{r}\varphi (r)$$
(29)
where $`r`$ is the distance between the core of SECS and the target atom, $`Z_Qe`$ is the electric charge of the SECS core, $`Z`$ is the atomic number of the target atom. The function $`\varphi (r)`$ is the screening function given by
$$\varphi (r)=\underset{1}{\overset{3}{}}C_iexp[\frac{b_ir}{a}]$$
(30)
where $`b_i`$ and $`c_i`$ are semiempirical constants given in Table 1 of ref. and the coefficients $`C_i`$ are restricted such that $`_1^3C_i=1`$, $`a`$ is the screening length of Eq. .
Assuming the validity of the potential of Eq. 29, we calculate the relation between the scattering angle $`\theta `$ and the impact parameter b. From this relation, the differential cross section $`\sigma (\theta )`$ is obtained as
$$\sigma (\theta )=(db/d\theta )b/\mathrm{sin}\theta $$
(31)
The relation between the transferred kinetic energy $`K`$ and the scattering angle $`\theta `$ is given by the relation
$$K=4E_{inc}\mathrm{sin}^2(\theta /2)$$
(32)
where $`E_{inc}`$ is the energy of the atom relative to the SECS in the center of mass system. The restricted energy losses are finally obtained by integrating the transferred energies as
$$\frac{dE}{dx}=N\sigma (K)𝑑K$$
(33)
where $`N`$ is the number density of atoms in the medium, $`\sigma (K)`$ is the differential cross section as function of the transferred kinetic energy $`K`$.
Fig. 11 shows the computed REL in CR39 for SECS with $`q=1e`$ and $`q=13e`$ plotted versus $`\beta `$.
## 5 Conclusions
Supersymmetric generalizations of the Standard Model allow for stable non-topological solitons, called Q-balls, which may be considered as bags of squarks and sleptons and thus have non-zero baryon and lepton numbers; they have positive electric charge (SECS) or neutral (SENS) small masses \[1-3\]. The Q-ball could be produced in the Early Universe, could affect the nucleosynthesis of light elements, and could lead to a variety of other astrophysical and cosmological consequences.
In this paper, we computed the energy losses of Q-balls of type SENS and SECS. Using these energy losses and a rough model of the Earth’s composition and density profiles, we have computed the angular acceptance of the an underground detector at the Gran Sasso Laboratory for Q-balls type SECS with $`v=250km/s`$ as function of the Q-ball mass $`M_Q`$. We have calculated the accessible region in the plane (mass, velocity) of SECS reaching the MACRO detector from above and from below.
We also presented an analysis of the energy deposited in the MACRO subdetectors: scintillators, streamer tubes and CR39 nuclear track detectors by SECS, in forms useful for their detection. In particular we computed the light yield in scintillators, the ionization in streamer tubes and the REL in nuclear track detectors.
The three MACRO subdetectors are sensitive to SECS with $`\beta 10^3`$ and masses larger than $`10^{13}GeV`$ ($`10^{19}GeV`$) for SECS coming from above (below). A flux upper limit may be obtained from MACRO at the level of few times $`10^{16}cm^2s^1sr^1`$.
SENS are more difficult to detect. MACRO scintillators could detect them, the streamer tubes have a limited efficiency and the CR39 detectors cannot see them.
Acknowledgements: We gratefully acknowledge the cooperation of many members of the MACRO collaboration, in particular of all the members of Bologna group. We thank A. Kusenko for stimulating discussions.
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# BOTTOM PRODUCTION
## 1 INTRODUCTION
In the context of the LHC experiments, the physics of bottom flavoured hadrons enters in different contexts. It can be used for QCD tests, it affects the possibilities of $`B`$ decays studies, and it is an important source of background for several processes of interest.
The physics of $`b`$ production at hadron colliders has a rather long story, dating back to its first observation in the UA1 experiment. Subsequently, $`b`$ production has been studied at the Tevatron. Besides the transverse momentum spectrum of a single $`b`$, it has also become possible, in recent time, to study correlations in the production characteristics of the $`b`$ and the $`\overline{b}`$.
At the LHC new opportunities will be offered by the high statistics and the high energy reach. One expects to be able to study the transverse momentum spectrum at higher transverse momenta, and also to exploit the large statistics to perform more accurate studies of correlations.
This chapter is organized as follows.
Section 2 is mostly theoretical. Its goal is to provide benchmark cross sections and distributions for the LHC, including rates relevant for the trigger requirements of the experiments. Furthermore, a discussion of the present status of $`b`$ production phenomenology at hadron colliders is given. In this context, one cannot forget that the theoretical status is a mixed success. On one side, the shape of distributions and correlations are reasonably well explained by perturbative QCD. On the other side, however, the observed cross sections at the Tevatron are larger than QCD predictions. It is hoped that further studies may help to understand the nature of the discrepancy. As of now, we see two possible explanations: either the absolute normalization of the cross section is not correctly predicted due to the presence of large higher order terms, or the shape of the distributions is distorted by some perturbative or non-perturbative effects (like, for example, fragmentation effects). With the wide $`p_T`$ range covered by the LHC experiments, and perhaps also with the possibility of performing more accurate studies of correlations, these two possibilities may be distinguished. The problem of fragmentation effects has been studied in this workshop also from the point of view of hadronization models in Monte Carlo programs, in Section 3. This study deals with the hadronization model in the HERWIG Monte Carlo program. Its aim was to understand whether, in simple realistic models of hadronization, the usual assumption of QCD factorization is really at work. In general, the problem of studying how realistic is the heavy flavour production mechanism implemented in shower Monte Carlo’s is quite important, and probably will require a considerable effort. Along this line, in Section 4, a problem in the heavy flavour production mechanisms implemented in Pythia is examined.
Further self-contained theoretical topics are dealt with in Sections 5 and 6. Section 5 deals with the charge asymmetry in $`b`$ production in $`pp`$ collisions. In this context, QCD is not of great help, since in perturbative QCD charge asymmetries turn out to be extremely small. Instead, studies are made within specific hadronization models, that are parametrized in such a way that they fit charm asymmetry data. This topic, besides being interesting in its own, since it deals with a phenomenon which is dominated by non-perturbative physics, has also an impact on $`CP`$ violation studies in $`B`$ decays.
In Section 6 quarkonium production is discussed. This subject has been intensively studied in recent years, following an initial CDF observation of a $`J/\mathrm{\Psi }`$ production rate much higher than theoretical predictions. This has triggered, from the theoretical side, the understanding that the fragmentation process is the dominant mechanism in quarkonium production. Besides this, a novel branch of applications of perturbative QCD, the NRQCD approach, has emerged, that may be useful to explain the production process.
In Section 7, the prospects for $`b`$ detection are discussed. It is shown that there is a complementarity between ATLAS/CMS and the LHCb experiment, with a certain region of overlap. In particular, the LHCb experiment can detect very low momentum heavy quarks, while the other experiments can reach the very high transverse momentum region. Some results on correlations measurements are also given, exploring the possibility of looking at one $`b`$ decaying into a $`J/\mathrm{\Psi }`$, and the other decaying semileptonically. Double heavy flavour production, charge asymmetry, polarization effects, and doubly-heavy meson production are also discussed.
In Section 8 the tuning of the multiple interaction parameters in Pythia is illustrated. The correct treatment of multiple interactions is important to model the multiplicity observables in both minimum-bias and heavy flavour events.
## 2 BENCHMARK CROSS SECTIONS<sup>1</sup><sup>1</sup>1Section coordinator: P. Nason and G. Ridolfi
### 2.1 Total cross sections
It is assumed that heavy flavour production in hadronic collisions can be described in the usual improved parton model approach, where light partons in the incoming hadrons collide and produce a heavy quark-antiquark pair via elementary strong interaction vertices, like, for example, in the diagram of fig. 1.
This description is appropriate for all hard processes in hadronic collisions, and thus, in the case of heavy flavours, is applicable as long as the mass of the heavy flavour can be considered sufficiently large. The perturbative QCD cross section for heavy flavour production has been computed to next-to-leading order accuracy (i.e. $`𝒪(\alpha _\mathrm{S}^3)`$) a long time ago , and a large amount of experimental and theoretical work has been done in this field. A relatively recent account of the status of this field can be found in ref. . It can be said that qualitatively the QCD description of heavy flavour production seems to be adequate also for charm production, while quantitatively large uncertainties are present in the calculation of the charm and bottom cross section. Only for a quark as heavy as the top quark the perturbative calculation seems (up to now) to predict the cross section with a good accuracy.
Large uncertainties are also found in the calculation of the bottom production cross section at the LHC. The largest uncertainty is due to unknown higher order effects, and it is traditionally quantified by estimating the scale dependence of the cross section when the renormalization and factorization scales are varied by a factor of 2 above and below their central value, which is usually taken equal to the heavy quark mass. Since this uncertainty is due to a limitation in our current theoretical knowledge, it is hard to overcome. Other sources of uncertainty are related to theoretical and experimental errors in the parameters that enter the perturbative calculation: the value of the strong coupling constant, the heavy quark mass, and the parton density functions.
We present here a benchmark study of $`b`$ total cross sections at the LHC, using the FMNR package for heavy flavour cross sections (the code for this package is available upon request to the authors). In the study we consider
* The dependence of the total cross section on the choice of the factorization and renormalization scales. We will use the values $`\mu =m_b,2m_b,m_b/2`$.
* The dependence on the parton density parametrization. We will use the sets MRST , MRST$`(g)`$, MRST$`(g)`$, MRST$`(\alpha _\mathrm{S})`$ and MRST$`(\alpha _\mathrm{S})`$. The first set is used as reference set. MRST$`(g)`$ and MRST$`(g)`$ have extreme gluon densities, MRST$`(\alpha _\mathrm{S})`$-MRST$`(\alpha _\mathrm{S})`$ have extreme values of the strong coupling constant: $`\mathrm{\Lambda }_5=220`$MeV for MRST, 164 MeV for MRST$`(\alpha _\mathrm{S})`$, 280 MeV for MRST$`(\alpha _\mathrm{S})`$. Cross section values obtained with the CTEQ4 set are very similar to the MRST set. We have preferred to use the MRST sets because they gave us the possibility to perform a study of sensitivity to $`\mathrm{\Lambda }`$ and to variations in the gluon density.
* The dependence on the $`b`$ quark mass: $`m_b=4.75\pm 0.25`$GeV.
Factorization and renormalization scale dependence of the total cross section at $`\sqrt{S}=14`$TeV is reported in table 1, where we have used the MRST parton densities, with $`\mathrm{\Lambda }_5=220`$MeV, and we have fixed the $`b`$ mass at the value $`m_b=4.75`$GeV.
Notice that:
* If we keep $`\mu _\mathrm{F}=\mu _\mathrm{R}`$, the full cross section variation is small (467 to 512 $`\mu b`$).
* The largest cross section corresponds to large $`\mu _\mathrm{F}`$ and small $`\mu _\mathrm{R}`$
* The smallest cross section corresponds to small $`\mu _\mathrm{F}`$ and large $`\mu _\mathrm{R}`$
This is understood since, at small $`x`$, the gluon density $`g(x)`$ grows with the scale, and $`\alpha _\mathrm{S}`$ decreases with the scale.
The dependence on the choice of parton density parametrization is shown in table 2.
As one can see, the sensitivity to the variation of the gluon density is small. Apparently, the constraints from HERA data are strong enough in the $`x`$ region where most of the $`b`$ production takes place. The dependence upon the strong coupling constant is instead larger, and can increase the upper limit of the cross section by about 40%. The last two sets have $`\mathrm{\Lambda }_5=164`$ and 288 MeV respectively, corresponding to $`\alpha _\mathrm{S}(M_Z)=0.1125`$ and 0.1225, which is a reasonably large range.
Mass uncertainties are quite important, especially if $`m_b`$ is allowed to take very small values. This can be seen from table 3.
We see that lowering the $`b`$ mass from 4.5 down to 4 GeV raises the upper limit of the cross section by about 50%. It is however unlikely that such small values are viable. A rough view of the status of the bottom mass determination is given in fig. 2,
which we obtained by taking the various determinations of the $`\overline{\mathrm{MS}}`$ bottom mass from the Particle Data Book, and rescale them by a factor of (1+0.09+0.06), to account for the two-loop correction needed to translate the $`\overline{\mathrm{MS}}`$ mass into the pole mass. As one can see, not all determinations are consistent among each other. A critical review of all determinations is beyond the scope of this workshop. We should however point out that recent progress has been made in the bottom mass determination. The reader can find a summary of these issues and further references in ref. . It is argued there that the bottom mass is determined with higher precision in processes where it is probed at short distances, like in the $`\mathrm{{\rm Y}}`$ mass, or via sum rules applied to the bottom vector current spectral function in $`e^+e^{}`$ annihilation. The mass extracted in this way can be reliably related to the so called $`\overline{\mathrm{MS}}`$ mass. The relation of the $`\overline{\mathrm{MS}}`$ mass to the pole mass is instead not so precise, because the perturbative expansion that relates the two quantities is not convergent. In ref. a preferred value of $`\overline{m}_b(m_b)=4.23\pm 0.08`$ is given, where $`\overline{m}_b(m_b)`$ is the $`\overline{\mathrm{MS}}`$ bottom mass at the scale of the bottom mass itself. The corresponding pole mass, obtained using the newly computed 3-loop relation between the $`\overline{\mathrm{MS}}`$ and pole mass , is $`4.98\pm 0.09`$ GeV. If one wanted to account for the uncertainties due to the lack of convergence of the perturbative expansion, the range obtained in this way should be enlarged by some amount, of the order of 100 MeV. The question arises whether it would be possible to eliminate this uncertainty by expressing the hadroproduction cross section in terms of the $`\overline{\mathrm{MS}}`$ mass. In our view, the answer is most likely no, since the bottom hadroproduction cross section does not have the same inclusive character as the sum rules applied to the $`e^+e^{}`$ bottom spectral function.
In the present work we thus used the traditional range $`4.5`$ GeV $`<m_b<5`$ GeV for the bottom pole mass in the hadroproduction process, keeping in mind that recent determinations seem to favour the upper region of this range. The sensitivity of the cross section to the bottom mass in this range is at most of $`\pm 10`$%, and it becomes much smaller if transverse momentum cuts are applied. Thus, as far as the LHC is concerned, this uncertainty is not very important.
The largest uncertainy in the cross section comes from the scale uncertainty, which is a (rather arbitrary) method to assess the possible impact of unknown higher order corrections. In the following we report a brief discussion of the origin of these large corrections. Radiative corrections for the total cross section are usually parametrized as follows. The total cross sections $`\sigma _{ij}`$ for the various parton subprocesses ($`\overline{q}q,qg,gg`$) have a perturbative expansions given by
$$\sigma _{ij}=\frac{\alpha _\mathrm{S}^2(\mu )}{m^2}\left[f_{ij}^{(0)}(\rho )+4\pi \alpha _\mathrm{S}\left(f_{ij}^{(1)}(\rho )+\overline{f}^{(1)}(\rho )\mathrm{log}\frac{\mu ^2}{m^2}\right)\right],$$
(1)
where $`\rho =4m_b^2/\widehat{s}`$ and $`\widehat{s}`$ is the squared partonic center-of-mass energy. The functions $`f_{ij}^{(0,1)}`$ for the $`\overline{q}q`$ and $`gg`$ subprocesses are displayed in fig. 3.
Notice the behaviour near threshold
$`f_{q\overline{q}}^{(1)}{\displaystyle \frac{f_{q\overline{q}}^{(0)}(\rho )}{8\pi ^2}}\left[{\displaystyle \frac{\pi ^2}{6\beta }}+{\displaystyle \frac{16}{3}}\mathrm{ln}^2\left(8\beta ^2\right){\displaystyle \frac{82}{3}}\mathrm{ln}(8\beta ^2)\right]`$
$`f_{gg}^{(1)}{\displaystyle \frac{f_{gg}^{(0)}(\rho )}{8\pi ^2}}\left[{\displaystyle \frac{11\pi ^2}{42\beta }}+12\mathrm{ln}^2\left(8\beta ^2\right){\displaystyle \frac{366}{7}}\mathrm{ln}(8\beta ^2)\right]`$
due to Coulomb $`1/\beta `$ singularities and to Sudakov double logarithms. Near threshold, these terms may require special treatement, such as resummation to all orders. Notice also the constant asymptotic behaviour of $`f_{gg}^{(1)}`$, which may cause problems far above threshold.
Plotting the cross section as a function of the partonic energy $`\widehat{s}`$ may help to understand the origin of large corrections.
We find that radiative corrections are large near the production threshold. This problem becomes more and more severe as we approach the production threshold. Thus, it is more important for production of $`b`$ at fixed target energies, or for production of $`t\overline{t}`$ pairs at colliders. Techniques to resum these large corrections to all orders of perturbation theory, at the NLO level are available , but it is found that little improvement is achieved for the bottom cross section at collider energies. Large corrections are also found far above threshold. This effect is bound to become more and more pronounced in the high energy limit. In order to reduce the scale uncertainties coming from these corrections, one should resum them at the next-to-leading order level. This problem has been discussed in the literature, so far, only at the leading logarithmic level . At the time of the completion of this workshop, no further progress has been achieved in this field.
In fig. 5 we present a study of the scale dependence of the total cross section as a function of $`\rho `$.
We find a large scale dependence near threshold, due to both renormalization and factorization scale variation, and a large scale dependence far from threshold. Here, the renormalization scale dependence plays a dominant role. Renormalization scale variations are mainly due to the large variation of the coupling constant in the $`𝒪(\alpha _\mathrm{S}^3)`$ terms. Where radiative corrections are small, a reasonable scale compensation takes place. Thus, both the threshold and the high energy regions, where corrections are large, are strongly affected. Factorization scale variation has a strong impact on threshold corrections, while in the high energy region we observe some compensation. In fact, the cross section near threshold increases with $`\mu _\mathrm{F}`$ near threshold, while above threshold the $`\mu _\mathrm{F}=m`$ value is above both the $`\mu _\mathrm{F}=m/2`$ and the $`\mu _\mathrm{F}=2m`$ curves, indicating the presence of some sort of compensation. As of now, it appears therefore that a better understanding of the high energy region will not strongly reduce the scale uncertainty, although it might, of course, improve our confidence in the error band we quote.
The study given here deals with total cross sections. It should be repeated with appropriate rapidity cuts, since this may reduce large effects due to the high energy limit. In general, we may expect that the cross section with rapidity and transverse momentum cuts may have smaller error bars than the total. It is particularly interesting to investigate directly cross sections for muons originating from $`B`$ decays, since muons are often used as trigger objects for $`B`$ physics. We have performed this study using a simple implementation of the $`B`$ semileptonic decay in the FMNR program, that will be described in more details in the following subsections. The results are shown in table 4.
The same results are also reported in fig. 6, since several features become more apparent there.
First of all, we point out that, as expected, there is a considerable reduction in the scale dependence in these muon rates. This is mostly due to the presence of cuts in the transverse momentum of the muon, that increases the total transverse energy that characterizes the cross section. Thus, while the ratio of the upper to the lower limit of the cross section is above a factor of 5 in the total rate, it is between a factor of 2 and 4 in the muon rates. The smallest values are achieved for the highest momentum cuts. A non perturbative fragmentation function of the Peterson form was also included in the calculation, with $`ϵ`$ parameter taking the values $`0`$ (i.e., no fragmentation), $`0.002`$ and $`0.006`$. More details on its implementation are given in the following subsections. Observe that for softer fragmentation functions (i.e. larger $`ϵ`$ parameter) the uncertainty is reduced, since they imply higher quark momenta. The reduction in the scale uncertainty is obtained at the price of introducing a sensitivity to the fragmentation function parameter. We considered as realistic values of $`ϵ`$ between 0.002 and 0.006. The corresponding variation of the cross section is not large. The impact of an intrinsic transverse momentum of the incoming partons (see the following subsections) is also studied. We have chosen the unrealistically large value $`k_T=4`$ GeV just to show that its effect is in all cases not a dramatic one.
### 2.2 Transverse momentum spectrum
#### 2.2.1 Benchmark single-inclusive distributions
The fixed-order, NLO result for single-inclusive $`b`$ production has several limitations in different regions of the phase space. In particular, one should be aware of the high-energy limit problem when $`p_T`$ is small compared to the incoming energy, of the logarithms of $`m_b/p_T`$ for high transverse momenta, and of further problems when approaching the threshold region. All these issues will be discussed in some detail in the next Sections. However, the fixed-order calculation at NLO provides a useful starting point for estimating the differential cross section. At this time, it is probably not useful to perform a cross section study with different sets of parton densities, and for different values of the $`b`$ mass. We limit ourselves to the MRST set, and we only study the scale dependence of the cross section. We do not include, at this stage, fragmentation effects, which, as shown in the following Sections, can be easily accounted for. In tables 5-8 we collect the results of this study.
The central values we obtained are also plotted in figs. 7 and 8, so that the wide kinematic range of heavy flavour production can be appreciated by a glance.
More detailed rapidity distributions at low momenta are shown in fig. 9.
First of all, we see that the differential cross section spans many orders of magnitude. At a luminosity of $`10^{34}\mathrm{cm}^2\mathrm{sec}^1`$ each $`\mu b`$ of cross section corresponds to $`10^4`$ events per second, or (roughly) $`10^{11}`$ events per year. Thus, at the level of $`10^{11}`$ in the plot there should be one event per year per bin of $`p_T`$ and $`y`$. The $`p_T`$ spectrum starts to drop fast for $`p_T`$ larger than the heavy quark mass, dropping even faster as the threshold region is approached. The rapidity distributions have the typical shape of a wide plateau, dropping at the edge of the phase space, and becoming narrower for larger transverse momenta. At the LHC the gluon fusion production mechanism is dominant, as can bee seen in fig. 10.
There one can see that the quark-antiquark annihilation component is below the gluon fusion component by more than one order of magnitude in the $`p_T`$ range considered, while the quark-gluon term becomes more important at larger $`p_T`$. We remind the reader that the cross section for $`qqb+X`$ is not included in the NLO calculation. One may thus worry about a loss of accuracy in the result, since the quark-quark luminosity at the LHC are by far the largest for high transverse momenta $`b`$ production. This problem, however, is dealt with appropriately in the resummation formalism for high $`p_T`$ heavy flavour production, where a quark-quark fusion contribution does indeed appear.
#### 2.2.2 Understanding Tevatron data
It is well known that Tevatron data for the integrated transverse momentum spectrum in $`b`$ production are systematically larger than QCD predictions. This problem has been around for a long time, although it has become less severe with time. The present status of this issue is summarized in fig. 11.
A similar discrepancy is also observed in UA1 data (see ref. for details).
The theoretical prediction has a considerable uncertainty, which is mainly due to neglected higher-order terms in the perturbative expansion. In our opinion, it is not unlikely that we may have to live with this discrepancy, which is certainly disturbing, but not strong enough to question the validity of perturbative QCD calculations. In other words, the QCD $`𝒪(\alpha _\mathrm{S}^3)`$ corrections for this process are above 100% of the Born term, and thus it is not impossible that higher order terms may give contributions of the same size. Nevertheless, it is useful to look for higher-order perturbative effects and non-perturbative effects that may enhance the cross section.
For values of $`p_T`$ much larger than the $`b`$ quark mass, large logarithms of the ratio $`p_T/m_b`$ arise in the coefficients of the perturbative expansion. Techniques are available to resum this class of logarithms to all orders. In ref. the NLO cross section for the production of a massless parton $`i`$ (a gluon or a massless quark) has been folded with the NLO fragmentation function for the transition $`ib`$ . The evolution of the fragmentation functions resums all terms of order $`\alpha _\mathrm{S}^n\mathrm{log}^n(p_T/m_b)`$ and $`\alpha _\mathrm{S}^{n+1}\mathrm{log}^n(p_T/m_b)`$. All the dependence on the $`b`$-quark mass lies in the boundary conditions for the fragmentation function. The result is then matched with the full NLO cross section, which contains the exact dependence on $`m_b`$ up to order $`\alpha _\mathrm{S}^3`$, in a way that avoids double counting. Corrections to the result of ref. are either of order $`\alpha _\mathrm{S}^4\mathrm{log}^i(p_T/m_b)`$, with $`i2`$, or of order $`\alpha _\mathrm{S}^4`$ times positive powers of $`m_b/\sqrt{p_T^2+m_b^2}`$.
Figures 12-13 show the differential and integrated $`b`$-quark $`p_T`$ distribution obtained in the fragmentation function approach of ref. , compared to the standard fixed-order NLO result. It should be noted that for high transverse momenta the scale dependence is significantly reduced with respect to the fixed-order calculation. Furthermore, it can be seen from fig. 13 that, for 10
>
pT
>
3010
>
subscript𝑝𝑇
>
3010\mathop{\mbox{\vbox{\hbox{$>$} \vskip-9.0pt\hbox{$\sim$}
\vskip-3.0pt}}}p_{T}\mathop{\mbox{\vbox{\hbox{$>$} \vskip-9.0pt\hbox{$\sim$}
\vskip-3.0pt}}}30 GeV, the result of the fragmentation-function approach lies slightly above the fixed-order NLO calculation. This has been interpreted in ref. as an evidence for large, positive higher order corrections. Unfortunately, their effect is not easy to quantify. These higher order terms are in fact computed in a massless approximation, and thus fail at low transverse momenta. In figs. 12-13 these terms are suppressed by a factor that becomes smaller and smaller at low $`p_T`$.
A more detailed discussion of this point can be found in the original reference. Here, we simply conclude that some evidence for large higher order terms in the intermediate transverse momentum region is present, although difficult to quantify.
Finally, notice that the overall effect of the inclusion of higher-order logarithms is a steepening of the $`p_T`$ spectrum. This is quite natural, since multiple radiation is accounted for in the resummation procedure.
It has been argued that an intrinsic transverse momentum for the incoming partons may explain the discrepancy observed at the Tevatron. In fact, large values (up to 4 GeV) of the average transverse momentum of the incoming partons have been invoked to explain direct photon production data . Such large values, much larger than typical QCD scales, are clearly incompatible with the usual application of perturbative QCD. Thus, evidence for such a large intrinsic transverse momentum cannot be claimed on the basis of a single observable. In other words, we would need evidence from several observables, all leading to a similar value of the intrinsic $`k_T`$, before we accept such a flaw in the usual perturbative QCD description. Nevertheless, in the following we will perform the exercise of applying very large intrinsic transverse momenta to the heavy flavour production process. This procedure will lead to an increase in the $`b`$ transverse momentum spectrum. We will also show, however, that other variables, very sensitive to an intrinsic transverse momentum, that should be strongly affected, do not show any evidence of that.
There are several possible ways to implement the presence of a non-zero transverse momentum of the colliding partons, and the choice is, to a large extent, arbitrary. We implemented it in the FMNR code in the following way. We call $`\stackrel{}{p}_\mathrm{T}(Q\overline{Q})`$ the total transverse momentum of the pair. For each event, in the longitudinal centre-of-mass frame of the heavy-quark pair, we boost the $`Q\overline{Q}`$ system to rest. We then perform a transverse boost, which gives the pair a transverse momentum equal to $`\stackrel{}{p}_\mathrm{T}(Q\overline{Q})+\stackrel{}{k}_\mathrm{T}(1)+\stackrel{}{k}_\mathrm{T}(2)`$; $`\stackrel{}{k}_\mathrm{T}(1)`$ and $`\stackrel{}{k}_\mathrm{T}(2)`$ are the transverse momenta of the incoming partons, which are chosen randomly, with their moduli distributed according to
$$\frac{1}{N}\frac{dN}{dk_\mathrm{T}^2}=\frac{1}{k_\mathrm{T}^2}\mathrm{exp}(k_\mathrm{T}^2/k_\mathrm{T}^2).$$
(2)
The reader can find more details in ref. .
In fig. 14 we show the effect of an intrinsic $`k_T`$ generated in this way, with the (unphysically large) choice $`k_\mathrm{T}=4`$ GeV (in fig. 14, the sensitivity to the $`ϵ_b`$ parameter in the fragmentation function is also shown; fragmentation will be discused in more detail in the next subsection.)
We see that, for $`p_T^{min}<20`$ GeV, the $`k_T`$ effect is sizeable, even in the presence of fragmentation, provided we allow for unphysically large intrinsic $`k_T`$.
It is fair to ask whether such large values are compatible with other observables. There is a particular class of observables that are particularly sensitive to the intrinsic transverse momentum. One example is the azimuthal distance $`\mathrm{\Delta }\varphi `$ between the directions of the produced $`b`$ and $`\overline{b}`$. The $`\mathrm{\Delta }\varphi `$ distribution is trivial at leading order: $`b`$ and $`\overline{b}`$ are emitted back-to-back, and therefore
$$\frac{d\sigma }{d\mathrm{\Delta }\varphi }\delta (\varphi \pi ).$$
(3)
An intrinsic $`k_T`$ of the colliding partons has the effect of smearing the $`\delta `$ function. For $`k_T=4`$GeV the effect is quite dramatic, as can be seen in fig. 15.
Is such an important effect consistent with the observed azimuth correlations? The CDF and D0 collaborations have measured the azimuthal correlation of muon pairs produced in $`b`$ decays. In order to compare with these data sets, we have implemented in the FMNR code the semileptonic decay of $`b`$ quarks. We have assumed that the muon energy is distributed according to the prediction of the spectator model with massless leptons. We have also checked that the muon energy distribution given by Pythia leads to similar results. Our results are shown in figs. 16 and 17, where CDF and D0 data are superimposed to the perturbative QCD prediction, with and without an intrinsic $`k_T`$ with $`k_T=4`$GeV. Tevatron data do not seem to favour such a large intrinsic transverse momentum. The measured distributions are more peaked at $`\mathrm{\Delta }\varphi =\pi `$ than the theoretical curve with $`k_T=4`$GeV. The effect of Peterson fragmentation is also shown in both cases.
We thus conclude that the data does not seem to favour large $`k_T`$ effects.
#### 2.2.3 Single-incusive distributions and correlations at the LHC
In this subsection, we will follow the pragmatic assumption that the discrepancy observed at the Tevatron may either be attributed to a problem in the overall normalization of the cross section, or to the presence of effects, either perturbative or not, that distort the spectrum. We will continue to model these effects as fragmentation effects and intrinsic transverse momentum effects, and see if the LHC can distinguish among the two. In fig. 18 we plot the $`b`$ cross section with a transverse momentum cut.
From the figure it is quite clear that the effects of fragmentation and the effects of an intrinsic transverse momentum kick manifest themselves in quite a different way. In particular, at $`p_T>20`$ GeV even the effect of a very large transverse momentum kick is small, while fragmentation has a strong impact. On the other hand, the transverse momentum kick increases the cross section in the intermediate $`p_T`$ region, with a maximum around 7 GeV. The $`p_T`$ coverage offered by the combined LHC experiments will allow an effective discrimination of the two kinds of effects. For completeness, we also show in fig. 19 a comparison of the fixed-order calculation of the single-inclusive spectrum, using the fixed-order calculation in two different schemes for the light flavour, and the matched-resummed result.
As in the Tevatron case, the band obtained with the resummation procedure is much narrower at large transverse momentum. The corresponding uncertainty does not include fragmentation function uncertainties, that will be discussed in more detail further on.
As an example of what could be discriminated at the LHC using correlations, we present in fig. 20 the azimuthal correlation of the muons coming from semileptonic $`B`$ decays, using typical cuts that are implemented in the LHC experiments triggers for $`B`$ studies.
The curves are obtained with different values of the $`ϵ_b`$ parameter for the fragmentation function, and with or without a very large intrinsic transverse momentum for the incoming partons. As one can expect, the $`ϵ_b`$ parameter affects only the total rate in this case, while the primordial transverse momentum has a considerable effect on the shape of the distribution. This example shows that, even with very simple experimental setup, at the LHC it will be possible to test important features of the differential distributions.
### 2.3 Fragmentation function formalism
In analogy with the case of charm production, the agreement between theory and data improves if one does not include any fragmentation effects. It is then natural to ask whether the fragmentation functions commonly used in these calculations are appropriate. Following the LEP measurements, fragmentation functions have appeared to be harder then previously thought. It will be interesting to see whether SLD new data will help in clarifying this issue.
The effect of a non-perturbative fragmentation function on the $`p_T`$ spectrum is easily quantified if one assumes a steeply-falling transverse momentum distribution for the produced $`b`$ quark
$$\frac{d\sigma }{dp_T}=Ap_T^M,$$
(4)
The corresponding distribution for the hadron is
$$\frac{d\sigma _{\mathrm{had}}}{dp_T}=A\widehat{p}_T^M\delta (p_Tz\widehat{p}_T)D(z)𝑑z𝑑\widehat{p}_T=Ap_T^M_0^1𝑑zz^{M1}D(z).$$
(5)
We can see that the hadron spectrum is proportional to the quark spectrum times the $`M^{\mathrm{th}}`$ moment of the fragmentation function $`D(z)`$. Thus, the larger the moment, the larger the enhancement of the spectrum.
In practice, the value of $`M`$ will be slightly dependent upon $`p_T`$. We thus define a $`p_T`$ dependent $`M`$ value
$$\frac{d\mathrm{log}\sigma (p_T>p_T^{\mathrm{cut}})}{d\mathrm{log}p_T^{\mathrm{cut}}}=M(p_T^{\mathrm{cut}})+1$$
(6)
and
$$\sigma _{\mathrm{had}}(p_T>p_T^{\mathrm{cut}})=\sigma (p_T>p_T^{\mathrm{cut}})\times _0^1𝑑zz^{M(p_T^{\mathrm{cut}})1}D(z).$$
(7)
This gives an excellent approximation to the effect of the fragmentation function, as can be seen from fig. 21.
Since the second moment of the fragmentation function is well constrained by $`e^+e^{}`$ data, it is sensible to ask for what shapes of the fragmentation function, for fixed $`z`$, one gets the highest value for $`z^{M1}`$. We convinced ourselves that the maximum is achieved by the functional form
$$D(z)=A\delta (z)+B\delta (1z)$$
(8)
which gives
$$z=\frac{B}{A+B};z^{M1}=\frac{B}{A+B}.$$
(9)
This is however not very realistic: somehow, we expect a fragmentation function which is concentrated at high values of $`z`$, and has a tail at small $`z`$. We convinced ourselves that, if we impose the further constraint that $`D(z)`$ should be monotonically increasing, one gets instead the functional form
$$D(z)=A+B\delta (1z),$$
(10)
which gives
$$z=\frac{A/2+B}{A+B};z^{M1}=\frac{A/M+B}{A+B}.$$
(11)
We computed numerically the $`M^{th}`$ moments of the Peterson form,
$$D(z)\frac{1}{z\left(1\frac{1}{z}\frac{ϵ}{1z}\right)^2}$$
(12)
of the form
$$D(z)z^\alpha (1z)^\beta $$
(13)
for $`\beta =1`$ (Kartvelishvili), for which
$$z^{M1}=\frac{\mathrm{\Gamma }(\alpha +M)\mathrm{\Gamma }(\alpha +\beta +2)}{\mathrm{\Gamma }(\alpha +1)\mathrm{\Gamma }(\alpha +\beta +M+1)},$$
(14)
of the form of Collins and Spiller
$$D(z)\frac{\left(\frac{1z}{z}+\frac{(2z)ϵ}{1z}\right)\left(1+z^2\right)}{\left(1\frac{1}{z}\frac{ϵ}{1z}\right)^2}$$
(15)
and of the form in eq. (10), at fixed values of $`z`$ corresponding to the choices $`ϵ_b=0.002`$ and $`0.006`$ in the Peterson form. We found that the $`p_T`$ distribution at the Tevatron, for $`p_T`$ in the range $`10`$ to $`100`$ GeV, behaves like $`p_T^M`$, with $`M`$ around 5. Therefore, we present in tables 9 and 10 values of the $`4^{\mathrm{th}}`$, $`5^{\mathrm{th}}`$ and $`6^{\mathrm{th}}`$ moments of the above-mentioned fragmentation functions.
We thus find that keeping the second moment fixed the variation of the hadronic $`p_T`$ distribution obtained by varying the shape of the fragmentation function among commonly used models is between 5% and 13% for both values of $`ϵ_b`$. It thus seems difficult to enhance the transverse momentum distribution by suitable choices of the form of the fragmentation function. With the extreme choice of eq. (10), one gets at most a variation of 50% for the largest value of $`ϵ_b`$ and $`M`$. It would be interesting to see if such an extreme choice is compatible with $`e^+e^{}`$ fragmentation function measurements.
## 3 A STUDY OF HEAVY QUARK NON-PERTURBATIVE FRAGMENTATION IN HERWIG<sup>3</sup><sup>3</sup>3Section coordinators: S. Frixione and M.L. Mangano
In this Section we present the results of a phenomenological study of the non-perturbative hadronization of $`b`$-quarks. According to the standard QCD picture, distributions for an observable hadron $`H`$ can be computed by convoluting the short-distance cross section $`\widehat{\sigma }(p)`$ with a fragmentation function $`D_H^{(h)}(z)`$ that describes the way in which the heavy quark $`h`$ hadronizes into $`H`$:
$$\mathrm{d}\sigma _H(p)=𝑑zD_H^{(h)}(z)d\widehat{\sigma }(p/z).$$
(16)
The precise definition of $`D_H^{(h)}(z)`$ depends on how much of the heavy quark evolution after its production is absorbed into the perturbative part $`\widehat{\sigma }(p)`$, and how much is assigned to the non-perturbative component parameterised by $`D_H^{(h)}(z)`$. Since perturbation theory (PT) is well defined for a massive quark, the standard prescription is to absorb into $`\widehat{\sigma }(p)`$ not only the hard matrix elements, but also the perturbative part of the fragmentation function, defined by the evolution in $`Q^2`$ down to a scale equal to the heavy quark mass $`m_h`$. $`D_H^{(h)}(z)`$ will therefore account for the transition of an “on-shell” quark $`h`$ into the hadron $`H`$. The assumptions built into eq. (16) are that $`D_H^{(h)}(z)`$ depends neither on the type of hard process, nor on the scale at which $`h`$ was produced. Under these assumptions, $`D_H^{(h)}(z)`$ can be extracted from data in one given reaction (typically, $`e^+e^{}`$), and eventually used to predict the cross section in some other reaction ($`p\overline{p}`$, DIS and so on).
QCD factorization theorems indicate that this universality of $`D_H^{(h)}(z)`$ holds in the asymptotic limit, and up to corrections of order $`m_h/Q`$, $`Q`$ being the scale of the hard process. The size of these corrections cannot be calculated, today, in any rigorous way. A possible approach to this problem is to turn to the phenomenological models of hadronization implemented in QCD-based parton-shower Monte Carlo (PSMC) codes. In PSMC the full final-state kinematical configuration is available at both the parton and hadron levels. Therefore, it is possible to “measure” $`D_H^{(h)}(z)`$ using eq. (16), both $`\mathrm{d}\sigma _H`$ and $`\mathrm{d}\widehat{\sigma }`$ being known. In the present section, we carry out this program using the PSMC HERWIG . HERWIG evolves quarks according to perturbative QCD down to small scales. The quarks are paired up at the end of the evolution into colour singlet clusters, which are then decayed to the physical hadrons using phenomenological distributions. The study of the heavy quark hadronisation process in HERWIG will allow us to test the universality assumption, and to measure the size of possible deviations.
We should stress that, at this moment, our conclusions are only relevant for the hadronization model implemented in HERWIG; other PSMC’s, which treat the hadronization process differently (for example, by adopting a string model), may well lead to different conclusions.
In order to precisely define our procedure for extracting $`D_H^{(h)}(z)`$, we need to consider in more details the way in which HERWIG generates events. Regardless of the type of initial-state particles, we can distinguish the following steps.
* Hard subprocess: at this stage, the PSMC generates the kinematics for the basic $`22`$ hard reaction. We denote the momentum of the $`b`$-quark (or antiquark) as $`p_b^{hard}`$.
* Parton shower: the partons resulting from the hard subprocess undergo successive branchings, until their virtuality is smaller than a fixed cutoff value. We denote the momentum of the $`b`$-quark at the end of this phase as $`p_b^{ps}`$.
* Gluon splitting and cluster formation: the gluons present at the end of the shower are decayed into light-quark pairs. Colour-singlet, two-body clusters are formed, according to colour parenthood and closeness in the phase-space. If there exist one or more cluster whose mass is too large (relative to a given threshold), part of the cluster rest energy is transformed into new $`q\overline{q}`$ pairs, and new clusters are defined. In this process, energy-momentum is redistributed among the cluster elements, and the momentum of the $`b`$-quark can therefore be modified with respect to $`p_b^{ps}`$. The momentum of the $`b`$ quark after completion of the clustering process will be denoted by $`p_b^{glsp}`$.
* Cluster decay and hadron formation: the clusters decay into observable hadrons, according to the flavour and to tabulated mass spectra. We therefore obtain $`b`$-flavoured hadrons, whose momentum we denote as $`p_B`$.
The hard subprocess and parton shower stages are based on perturbative QCD. Thus, we identify the predictions given by HERWIG at the end of the parton shower with the cross section $`\widehat{\sigma }`$ that appears in eq. (16). On the other hand, the gluon splitting and cluster decay stages do not contain QCD information, as they are performed according to a phenomenological model. The $`gq\overline{q}`$ splitting and the decay kinematics are induced by simple phase-space considerations. We thus identify these stages as the long-distance, non perturbative part of the process, which gives rise to $`D_H^{(h)}(z)`$. We therefore determine the fragmentation function by comparing the results for $`p_b^{ps}`$ and $`p_B`$, defining, on an event-by-event basis, the following variables:
$$z_1=\frac{\stackrel{}{p}_B\stackrel{}{p}_b^{ps}}{\left|\stackrel{}{p}_b^{ps}\right|^2},z_2=\frac{E_B+\stackrel{}{p}_B\widehat{p}_b^{ps}}{E_b^{ps}+\left|\stackrel{}{p}_b^{ps}\right|},$$
(17)
where $`\widehat{p}_b^{ps}=\stackrel{}{p}_b^{ps}/\left|\stackrel{}{p}_b^{ps}\right|`$. Our conclusions will apply to both $`z_1`$ and $`z_2`$; thus, we will collectively denote them by $`z`$. In hadronic collisions, the momenta and energies have to be substituted by transverse momenta and transverse energies respectively. Our results are summarized in table 11; we considered $`e^+e^{}`$ collisions at $`\sqrt{S}=91.2`$ GeV and $`p\overline{p}`$ collisions at $`\sqrt{S}=1.8`$ TeV. In the table, we present four of the (normalized) Mellin moments of the $`z`$ distribution, defined as follows:
$$\mu _n=𝑑zz^nD_H^{(h)}(z)/𝑑zD_H^{(h)}(z).$$
(18)
Usually, $`0z1`$. In the present case, as we will see, we can also have $`z>1`$; thus, in eq. (18) the range of integration coincide with the support of $`D_H^{(h)}(z)`$.
The Mellin moments appearing in table 11 have been evaluated by considering bins in $`\left|\stackrel{}{p}_b^{ps}\right|`$ (in the case of hadronic collisions, the momentum is the transverse one). In $`e^+e^{}`$ collisions larger (smaller) values of $`\left|\stackrel{}{p}_b^{ps}\right|`$ correspond to less (more) energy lost to gluons. In hadronic collisions larger (smaller) values of $`\left|\stackrel{}{p}_b^{ps}\right|`$ are more likely to correspond to larger (smaller) values of the hard process momentum before evolution. In either case, dependence of $`D_H^{(h)}(z)`$ on $`\left|\stackrel{}{p}_b^{ps}\right|`$ signals therefore a departure from universality.
By inspection of the table, we see that $`D_H^{(h)}(z)`$ is scale-independent to a very good extent (the situation appears to be slightly better in the case of $`e^+e^{}`$ collisions), except for the very low $`p_b^{ps}`$ region; this is what we should expect, since in that region the factorization theorem on which eq. (16) is based is bound to fail. On the other hand, there seems to emerge a clear difference between the fragmentation functions extracted from $`e^+e^{}`$ and $`p\overline{p}`$ “data”, the latter being substantially softer than the former. The first moment, which is the average value of the fragmentation variable, changes by about 10%. This variation can change the rate of predicted $`b`$-hadrons in hadronic collisions by almost 50%.
This suggests that transporting to hadronic collisions the non-perturbative fragmentation functions obtained by fitting $`e^+e^{}`$ data may not be correct. Of course, a much more detailed investigation on the subject is required before reaching a firm conclusion; however, this simple exercise of ours shows that universality should not be taken for granted.
We now concentrate on the separate role played in the fragmentation process by the purely perturbative evolution and by the non-perturbative gluon-splitting phase, before the cluster formation and decay take place. We shall confine ourselves to the case of $`e^+e^{}`$ collisions. The variables relevant to our study are the following:
1. Energy fraction retained during the perturbative evolution:
$$z_{ps}=\frac{2\left|\stackrel{}{p}_b^{ps}\right|}{\sqrt{S}}=\frac{\left|\stackrel{}{p}_b^{ps}\right|}{\left|\stackrel{}{p}_b^{hard}\right|},$$
(19)
where $`\sqrt{S}`$ is the $`e^+e^{}`$ CM energy.
2. Energy fraction retained during the gluon-splitting:
$$z_{glsp}=\frac{\left|\stackrel{}{p}_b^{glsp}\right|}{\left|\stackrel{}{p}_b^{ps}\right|}.$$
(20)
3. Energy fraction left after the perturbative evolution and the gluon-splitting:
$$z=z_{ps}\times z_{glsp}=\frac{2\left|\stackrel{}{p}_b^{glsp}\right|}{\sqrt{S}}=\frac{\left|\stackrel{}{p}_b^{glsp}\right|}{\left|\stackrel{}{p}_b^{hard}\right|},$$
(21)
The left panel in fig. 22 shows the three distributions for $`b`$ quarks at $`\sqrt{S}=91.2`$ GeV. The solid histogram represents the distribution of $`z_{ps}`$. The distribution has the shape of a Gribov-Lipatov, with no indication of a Sudakov turn-over at large $`z_{ps}`$. The dotted line is the distribution in $`z`$. A strong deformation of the purely perturbative curve is clearly seen. The dashed line corresponds to the $`z_{glsp}`$ distribution. This is part of what the MC treats as a non-perturbative component of the fragmentation process. The peak of the dashed histogram at $`z_{glsp}=1`$ corresponds to events where the cluster containing the heavy quark does not need to be further split, while the tail corresponds to events where the invariant mass of the heavy-quark cluster is too large, and additional light-quark pairs have to be produced by hand. Notice that almost as much energy is lost during this non-perturbative phase, as is lost during the perturbative evolution.
For comparison, we show the same set of curves for the evolution of the charm quark (right panel of fig. 22). Notice that while the effect of the perturbative evolution is to soften the quark spectrum relative to the $`b`$-quark case, the amount of energy lost due to gluon splitting is similar ($`z_{glsp}^c=0.82`$, as opposed to $`z_{glsp}^b=0.85`$). This is bizarre, since one expects the non-perturbative part to scale with $`1/m_h`$. The same result is found for the fragmentation of the $`s`$ quark (left panel of fig. 23). Here $`z_{glsp}^s`$ is $`0.81`$. Again, a violation of the expected $`1/m_h`$ scaling.
Things improve for the top quark, whose distributions for $`\sqrt{S}=2`$ TeV are shown on the right panel of fig. 23. The gluon-splitting part has only a minor impact on the overall spectrum of the top quark.
We are a bit bothered by the dominant role played by the gluon-splitting phase. By comparison, the next step in the evolution, namely the cluster formation and decay, plays only a minor role, as will be shown next. We would have anticipated that the cluster formation and decay should be the place where most of the non-perturbative physics should show up. This suggests that the thresholds for the perturbative evolution in the MC shold be lowered, so that the impact of the non-perturbative gluon splitting phase is reduced, and purely perturbative Sudakov effects can manifest themselves.
We now turn again to the non-perturbative part of the fragmentation function. The most striking feature, that cannot be inferred from the simple study of Mellin moments as done in table 11, is the presence of a double peak in the high-$`z`$ region (see the left panel of fig. 24). A first peak (which we will call peak A) is seen at $`z`$ values around 0.97. A second peak (peak B) is at $`z=1.01`$ (we have a $`z`$-bin size of $`0.02`$. We verified that the events contributing to the peak B do not have $`z=1+ϵ`$, i.e., the peak is not due to a roundoff error). The latter peak is higher than the former.
The origin of this double peak can be traced back to the following facts. First, the momentum of the emitted $`B`$ meson is very strongly correlated with the momentum of the $`b`$ quark which enters the cluster. Therefore, the $`z`$ distribution closely reflects the mass spectrum of the light hadron emitted, together with the $`B`$ meson, in the cluster decay. Second, the peak B is almost entirely due to events where the cluster decay into $`B+\pi `$: at this peak, $`z>1`$ because the mass of the pion is lighter than the mass of the lightest quarks in HERWIG. More in detail, we have observed the following facts.
* The structure of the double peak is strongly influenced by the value taken by the two input parameters CLDIR and CLSMR. If the default is used (CLDIR=1, CLSMR=0), the double peak is observed (see the plot on the left of fig.24). On the other hand, by setting CLSMR$``$0, the $`z`$ distributions display a single peak (broader that the previous ones) at about $`z=0.97`$. For small CLSMR values and large $`b`$ momenta, a second peak at $`z=1.01`$ tends to re-appear, although smaller than observed before.
* The double peak disappears also if one chooses CLDIR=0, as in older HERWIG versions. In this case, the $`z`$ distributions peak at about $`z=0.9`$, this peak being much broader that those obtained with CLDIR=1, regardless of the value of CLSMR.
* We then set CLDIR=1 and CLSMR=0. For any given $`B`$ meson, we looked for the parent cluster $`𝒞=\{b_Cq\}`$, and for the parent bottom quark, $`b_P`$ (the parent quark is defined as in the HERWIG routine HWCHAD). We observed what follows.
+ Plotting the $`z`$ distributions for the events with $`\stackrel{}{p}_{b_C}\stackrel{}{p}_{b_P}`$ (i.e. events where the original cluster was split), we see a single peak, at the same $`z`$ value as for the peak A (solid line, plot on the right of fig.24).
+ The $`z`$ distributions for events such that $`\stackrel{}{p}_{b_C}=\stackrel{}{p}_{b_P}`$ display again a double peak. The two peaks are at the same $`z`$ values as peaks A and B, the latter one being by far dominant.
+ Selecting only events with $`\stackrel{}{p}_{b_C}=\stackrel{}{p}_{b_P}`$, we found that the peak at the position of B corresponds to those clusters decaying into a $`B`$ meson and a $`\pi `$, while the peak at the position of A is relevant for all the other two-body decays (dotted and dashed lines respectively, plot on the right of fig.24).
Overall, notice also that the amount of energy retained after the gluon-splitting phase is of the same size as that retained at the end of the full hadronization process, indicating that cluster formation and decay have a minor impact on the total amount of energy lost during the non-perturbative part of the evolution.
We were also able to reproduce the previous findings with a very simple model. Given a momentum for a quark $`b`$, we generate randomly the momentum for a light quark $`q`$, to be combined with $`b`$ into a cluster, which eventually decays into a $`B`$ meson and a particle of given mass $`m_P`$. The momentum of the quark $`q`$ is allowed to have a (small) transverse momentum with respect to the direction of the quark $`b`$. After evaluating the cluster mass, we performed the decay in the rest frame of the cluster, either in a isotropic manner (thus mimicking the choice CLDIR=0), or by letting the momentum of the meson $`B`$ to be parallel to that of the quark $`b`$ (which corresponds to CLDIR=1 and CLSMR=0). In the latter case, depending upon the value of $`m_P`$, we got a peak for $`z<1`$ (if $`m_P>m_q`$) or $`z>1`$ (if $`m_P<m_q`$).
In conclusion, the $`z`$ distributions we find when using HERWIG seem not to contain a lot of dynamical information, the most important features being those implemented in the cluster-decay routine. If the decay is not smeared out (CLDIR=0), we get a structure which is very difficult to reconcile with the idea of fragmentation we have from QCD. After smearing, the distribution still has a $`z>1`$ tail which will be extremely difficult to fit with a function vanishing for $`z1`$. This problem is related to the fact that the mass of the lightest quarks in the MC is 320 MeV, that is much larger than the pion mass. We performed a test by reducing the light quark masses to 20 MeV, and increasing the shower cutoff VQCUT in such a way as to maintain the default value of the effective infrared threshold. The double peak structure, as expected, disappeared. It remains to be seen, however, whether such a small value of the quark masses is, more generally, acceptable.
## 4 A STUDY OF THE $`b\overline{b}`$ PRODUCTION MECHANISM IN PHYTIA<sup>5</sup><sup>5</sup>5Section coordinators: S. Gennai, A. Starodumov, F. Palla, R. Dell’Orso
### 4.1 Introduction
In this section, we present a study on $`b\overline{b}`$ production performed within the CMS collaboration using the Monte Carlo package Pythia 5.75 as an event generator. In particular, we investigate the influence of the cut-off on the hard interaction transverse momentum $`\widehat{P_t}`$ on the production of $`b\overline{b}`$ events.
In Monte Carlo programs, $`b\overline{b}`$ pairs in hadron collisions are produced by the mechanisms of gluon fusion, gluon splitting and flavour excitation. All of them give contributions of the same order to the total cross section, but they give rise to different kinematical configurations of the final state.
There are two ways to generate $`b\overline{b}`$ events in Pythia:
* Using a steering card MSEL=5, a gluon fusion mechanism ($`ggb\overline{b}`$) is mainly simulated. Each event contains at least one $`b\overline{b}`$ pair.
* Using MSEL=1, all QCD $`22`$ processes are simulated. In this case, all production mechanisms contribute to the $`b\overline{b}`$ production, but the probability to find a $`b\overline{b}`$ pair in the event is less than 1$`\%`$.
About one million events have been simulated in CMS with MSEL=1, in order to have a sample with all $`b\overline{b}`$ production mechanisms and default Pythiacut-off, not to introduce any bias in the kinematics. The selection efficiency of triggered events out of this sample is quite low. In order to have higher signal statistics, in some cases one can use kinematical cuts which are different from the Pythiadefault. .
### 4.2 $`b\overline{b}`$ production
Two samples have been prepared to investigate the influence of the $`\widehat{P_t}`$ cut on the production of $`b\overline{b}`$ events. Both of them have been generated using MSEL=1 and contain events with only one $`b\overline{b}`$ pair. Only events with $`\widehat{P_t}10`$ GeV have been selected in the samples. The first sample (SAMPLE A) has been generated with the default $`\widehat{P_t}`$ cut of 1 GeV and only events with $`\widehat{P_t}10`$ GeV were selected. The second sample (SAMPLE B) has been generated with $`\widehat{P_t}10`$ GeV. In both samples the following processes contribute:
$`ggq\overline{q}`$ (22)
$`gq_igq_i`$ (23)
$`gggg.`$ (24)
$`b\overline{b}`$ pair is produced by gluon splitting $`gb\overline{b}`$ in initial or final state shower evolution (processes (22) to (24)) or in the hard interaction (process (22)). For both samples A and B, we have computed the $`b\overline{b}`$ production cross section
$$\sigma _{b\overline{b}}^{tot}=\frac{N_{b\overline{b}}}{N^{tot}}\sigma ^{tot},$$
(25)
where $`N_{b\overline{b}}`$ is the number of $`b\overline{b}`$ events with $`\widehat{P_t}10`$ GeV, $`N^{tot}`$ is the total number of generated events, and $`\sigma ^{tot}`$ is the total cross section (given by Pythia). We find that
* for sample A, $`\sigma _{b\overline{b}}^{tot}=150`$ $`\mu `$b. The gluon fusion contribution is about 20 $`\mu `$b, while the gluon splitting contributions are $`30`$ $`\mu `$b and $`100`$ $`\mu `$b for processes (23) and (24) respectively;
* for sample B, $`\sigma _{b\overline{b}}^{tot}`$=257 $`\mu `$b. Gluon fusion and gluon splitting contributions are at the same level as in sample A. In this case, however, there are also contributions from the processes $`bgbg`$ and $`bqbq`$ of about 110 $`\mu `$b. In the following we will call these contributions flavour excitation.
Figure 25 illustrates the difference in the $`b\overline{b}`$ production cross sections due to the additional contribution of the flavour excitation mechanism in sample B.
The effect has the following explanation. When the default $`\widehat{P_t}`$ cut-off is used, Pythia generates processes in the low energy approximation, i.e. there are no heavy quarks inside the parton distribution. This approach changes if one uses a different $`\widehat{P_t}`$ cut-off: the parton distributions in this case include also $`b`$ and $`c`$ quarks. As a consequence, samples A and B are different in two respects: values of the cross sections, and set of production mechanisms. The difference in the cross section is not very important, because the results are usually normalized to the total $`b\overline{b}`$ cross section of 500 $`\mu `$b. On the other hand, the different production mechanisms could be more dangerous, as they can lead to different kinematical distributions, and therefore affect the efficiencies of physical selection.
### 4.3 Kinematics
The main kinematical parameters which define the signature of $`b\overline{b}`$ event are the transverse momenta and pseudorapidities of the $`b`$ quarks, and the angular distance $`\mathrm{\Delta }\varphi `$ between their directions in the transverse plane. The first two parameters have similar distributions in both samples. The $`\mathrm{\Delta }\varphi `$ distribution is shown in fig. 26 for the three different mechanisms.
For what concerns gluon splitting, the distribution is slightly peaked at small $`\mathrm{\Delta }\varphi `$. The angle between the two $`b`$-quarks produced by the gluon-fusion mechanism has a peak at $`\mathrm{\Delta }\varphi \pi `$, as expected, since in the process $`ggb\overline{b}`$ the b-quarks are produced back-to-back in the transverse plane. The last distribution corresponds to the flavour excitation production mechanism, for which the back-to-back topology is preferred. We can conclude that the total $`\mathrm{\Delta }\varphi `$ distributions of sample A and sample B are slightly different. Some care should be taken about this, as it could affect the estimated efficiency of selection cuts.
### 4.4 Pythia 6.125
We have studied the same problem using the new 6.125 version Pythia. We have generated two new samples A and B with Pythia 6.125, and we have found the following results:
* Sample A: the $`b\overline{b}`$ production cross section is $`\sigma _{b\overline{b}}^{tot}`$=220 $`\mu `$b. Gluon fusion contributes $``$ 47 $`\mu `$b and total gluon splitting gives 173 $`\mu `$b.
* Sample B: $`b\overline{b}`$ production cross section is $`\sigma _{b\overline{b}}^{tot}`$=465 $`\mu `$b. In this case, gluon fusion is $``$ 51 $`\mu `$b and total gluon splitting is $``$ 193 $`\mu `$b. The contribution from the flavour excitation is about 221 $`\mu `$b.
Also the way Pythia 6.125 generates $`b\overline{b}`$ pairs depends on the $`\widehat{P_t}`$ cut-off. The difference with Pythia 5.75 values, are due to different total cross section in the two versions.
### 4.5 Interpretation
Many of the features of Pythia illustrated in this section are easily explained<sup>7</sup><sup>7</sup>7T. Sjöstrand and E. Norrbin, private communication.. It turns out that Pythia treats differently processses with a low and high $`\widehat{p}_T^{\mathrm{min}}`$. The limit is related to the scale of multiple interactions, which is fixed to 2 GeV in the older versions, and was made energy dependent in Pythia 6, being 3.2 GeV at the LHC energy. When $`\widehat{p}_T^{\mathrm{min}}`$ is above this scale, the hard process is selected according to conventional matrix elements. Below this scale, the hardest interaction is instead taken from the naive jet cross section multiplied by a “Sudakov style” form factor, that represents the probability that higher $`p_T`$ interactions did not take place in the rest of the event. Since this procedure implies the computation of all parton-parton scattering processes, the choice was made to exclude from it the incoming $`b`$ and $`c`$ components, to save time in the computation. This feature is no longer considered useful in modern times, the computers being much faster. Thus, in Pythia 6.138, also the $`b`$ and $`c`$ processes will be implemented in the low $`p_T`$ mode.
The difference in the total cross section in Pythia 5.7 and 6.1 have a physical origin, since 6.1 uses newer parton distributions that, according to HERA data, are more singular in the small $`x`$ region.
The authors of Pythia recommend the following procedure for the generation of $`b`$ events. Parton fusion and flavour excitation can be generated separately; the relevant massive matrix elements are used for parton fusion, and one can go to the limit $`p_T0`$ with this process. Gluon splitting cannot be generated separately: all hard processes must be generated, excluding parton fusion and flavour excitation, and one should look for the heavy flavour. Multiple interactions are there switched off, in order to avoid a double-counting of the jet cross section. This is adequate for the study of the $`b`$ production properties, but clearly does not fully represent the structure of the underlying event. In future Pythia versions, when flavour excitation is included in the minimum bias machinery with multiple interaction, this latter should offer an almost equivalent alternative, but still without the correct mass treatment of the parton fusion process near threshold. Other limitations still remain from complex problems related to the treatment of beam remnants; therefore, flavour excitation is only enabled for the hardest interaction in the multiple-interaction scenario.
A sample of commented code is included below. By using different flags (MEKIND=0,1,2) three samples will be generated: parton fusion, flavour excitation and gluon splitting.
```
INTEGER KFINTMP(-40:40)
C... Multiple interactions switched off
MSTP(81)=0
PARP(81)=0.D0
PARP(82)=0.D0
C... Maximum virtuality in ISR is PARP(67)*Q**2
C... This is the default in Pythia 6.137
PARP(67)=1.D0
C... Choose heavy quark (bottom=5, charm=4)
MASSIVE=5
C... Helper variable
HQMASS=PMAS(MASSIVE,1)
C... Choose the kind of heavy quark production:
C... MEKIND is a local variable set to 0, 1 or 2
IF (MEKIND==0) THEN ! Massive matrix elements
MSEL=MASSIVE
ELSE IF (MEKIND==1) THEN ! Flavour excitation
MSEL=1
CKIN(3)=HQMASS
CKIN(5)=CKIN(3)
ELSE IF (MEKIND==2) THEN ! Gluon splitting (ISR, FSR)
MSEL=1
CKIN(3)=HQMASS
CKIN(5)=CKIN(3)
END IF
C... More restrictive cuts can be put here.
C... Example, 100 events in total.
NEVENTS=100
C *** EVENT LOOP ***
IF (MEKIND==1) NEVENTS=NEVENTS/2
C.... Loop over incoming partons
DO ISIDE=1,2
IF (MEKIND/=1.AND.ISIDE==1) THEN
GOTO 100
ELSE IF (MEKIND==1) THEN
C... Only for flavour excitation:
C... Make backup copy of KFIN array
DO IKF=-40,40
KFINTMP(IKF)=KFIN(ISIDE,IKF)
C... Remove all incoming partons:
KFIN(ISIDE,IKF)=0
END DO
C... Select only b/bbar as incoming partons:
KFIN(ISIDE, MASSIVE)=1
KFIN(ISIDE,-MASSIVE)=1
END IF
DO IEV=1,NEVENTS
C... Generate an event
CALL PYEVNT
C... For gluon splitting, remove events with HQ in the hard interaction
C... to avoid double counting:
IF (MEKIND==2) THEN
DO I=5,8
IF (ABS(K(I,2))==MASSIVE) GOTO 50
END DO
END IF
C... Analysis...
50 END DO
C... Print statistics
CALL PYSTAT(1)
C... Restore KFIN matrix:
IF (MEKIND==1.AND.ISIDE==1) THEN
DO IKF=-40,40
KFIN(ISIDE,IKF)=KFINTMP(IKF)
END DO
END IF
100 END DO
```
## 5 ASYMMETRIES<sup>8</sup><sup>8</sup>8Section coordinators: E. Norrbin and R. Vogt
### 5.1 Introduction
Sizeable leading particle asymmetries between e.g. $`\mathrm{D}^{}`$ and $`\mathrm{D}^+`$ have been observed in several fixed target experiments . It is of interest to investigate to what extent these phenomena translate to bottom production and higher energies. No previous experiment has observed asymmetries for bottom hadrons due to limited statistics or other experimental obstacles. Bottom asymmetries are in general expected to be smaller than for charm because of the larger bottom mass, but there is no reason why they should be absent. In the fixed target experiment HERA-B, bottom asymmetries could very well be large even at central rapidities, but the conclusion of the present study is that asymmetries at the LHC are likely to be small. In the following we study possible asymmetries between $`\mathrm{B}`$ and $`\overline{\mathrm{B}}`$ hadrons at the LHC within the Lund string fragmentation model and the intrinsic heavy quark model .
In the string fragmentation model , the perturbatively produced heavy quarks are colour connected to the beam remnants. This gives rise to beam-drag effects where the heavy hadron can be produced at larger rapidities than the heavy quark. The extreme case in this direction is the collapse of a small string, containing a heavy quark and a light beam remnant valence quark of the proton, into a single hadron. This gives rise to flavour correlations which are observed as asymmetries. Thus, in the string model, there can be coalescence between a perturbatively produced bottom quark and a light quark in the beam remnant producing a leading bottom hadron.
There is also the possibility to have coalescence between the light valence quarks and bottom quarks already present in the proton, because the wave function of the proton can fluctuate into Fock configurations containing a $`\mathrm{b}\overline{\mathrm{b}}`$ pair, such as $`|\mathrm{uudb}\overline{\mathrm{b}}`$. In these states, two or more gluons are attached to the bottom quarks, reducing the amplitude by $`𝒪(\alpha _\mathrm{s}^2)`$ relative to parton fusion . The longest-lived fluctuations in states with invariant mass $`M`$ have a lifetime of $`𝒪(2P_{\mathrm{lab}}/M^2)`$ in the target rest frame, where $`P_{\mathrm{lab}}`$ is the projectile momenta. Since the comoving bottom and valence quarks have the same rapidity in these states, the heavy quarks carry a large fraction of the projectile momentum and can thus readily combine to produce bottom hadrons with large longitudinal momenta. Such a mechanism can then dominate the hadroproduction rate at large $`x_\mathrm{F}`$. This is the underlying assumption of the intrinsic heavy quark model , in which the wave function fluctuations are initially far off shell. However, they materialize as heavy hadrons when light spectator quarks in the projectile Fock state interact with the target .
In both models the coalescence probability is largest at small relative rapidity and rather low transverse momentum where the invariant mass of the $`\overline{\mathrm{Q}}\mathrm{q}`$ system is small, enhancing the binding amplitude. One exception is at very large $`p_{}`$, where the collapse of a scattered valence quark with a $`\overline{\mathrm{b}}`$ quark from the parton shower is also possible, giving a further (small) source of leading particle asymmetries in the string model.
### 5.2 Lund String Fragmentation
Before describing the Lund string fragmentation model, some words on the perturbative heavy quark production mechanisms included in the Monte Carlo event generator Pythia used in this study is in order. We study $`\mathrm{pp}`$ events with one hard interaction because events with no hard interaction are not expected to produce heavy flavours and events with more than one hard interaction — multiple interactions — are beyond the scope of this initial study and presumably would not influence the asymmetries. After the hard interaction is generated, parton showers are added, both to the initial (ISR) and final (FSR) state. The branchings in the shower are taken to be of lower virtualities than the hard interaction introducing a virtuality (or time) ordering in the event. This approach gives rise to several heavy quark production mechanisms, which we will call pair creation, flavour excitation and gluon splitting. The names may be somewhat misleading since all three classes create pairs at $`\mathrm{g}\mathrm{Q}\overline{\mathrm{Q}}`$ vertices, but it is in line with the colloquial nomenclature. The three classes are characterized as follows.
The hard subprocess is one of the two LO parton fusion processes $`\mathrm{gg}\mathrm{Q}\overline{\mathrm{Q}}`$ or $`\mathrm{q}\overline{\mathrm{q}}\mathrm{Q}\overline{\mathrm{Q}}`$. Parton showers do not modify the production cross sections, but only shift kinematics. For instance, in the LO description, the $`\mathrm{Q}`$ and $`\overline{\mathrm{Q}}`$ have to emerge back-to-back in azimuth in order to conserve momentum, while the parton shower allows a net recoil to be taken by one or several further partons.
A heavy flavour from the parton distribution of one beam particle is put on mass shell by scattering against a parton of the other beam, i.e. $`\mathrm{Qq}\mathrm{Qq}`$ or $`\mathrm{Qg}\mathrm{Qg}`$. When the $`\mathrm{Q}`$ is not a valence flavour, it must come from a branching $`\mathrm{g}\mathrm{Q}\overline{\mathrm{Q}}`$ of the parton-distribution evolution. In most current sets of parton-distribution functions, heavy-flavour distributions are assumed to vanish for virtuality scales $`Q^2<m_\mathrm{Q}^2`$. The hard scattering must therefore have a virtuality above $`m_\mathrm{Q}^2`$. When the initial-state shower is reconstructed backwards , the $`\mathrm{g}\mathrm{Q}\overline{\mathrm{Q}}`$ branching will be encountered, provided that $`Q_0`$, the lower cutoff of the shower, obeys $`Q_0^2<m_\mathrm{Q}^2`$. Effectively the processes therefore become at least $`\mathrm{gq}\mathrm{Q}\overline{\mathrm{Q}}\mathrm{q}`$ or $`\mathrm{gg}\mathrm{Q}\overline{\mathrm{Q}}\mathrm{g}`$, with the possibility of further emissions. In principle, such final states could also be obtained in the above pair-creation case, but the requirement that the hard scattering must be more virtual than the showers avoids double counting.
A $`\mathrm{g}\mathrm{Q}\overline{\mathrm{Q}}`$ branching occurs in the initial- or final-state shower but no heavy flavours are produced in the hard scattering. Here the dominant $`\mathrm{Q}\overline{\mathrm{Q}}`$ source is gluons in the final-state showers since time-like gluons emitted in the initial state are restricted to a smaller maximum virtuality. Except at high energies, most initial state gluon splittings instead result in flavour excitation, already covered above. An ambiguity of terminology exists with initial-state evolution chains where a gluon first branches to $`\mathrm{Q}\overline{\mathrm{Q}}`$ and the $`\mathrm{Q}`$ later emits another gluon that enters the hard scattering. From an ideological point of view, this is flavour excitation, since it is related to the evolution of the heavy-flavour parton distribution. From a practical point of view, however, we choose to classify it as gluon splitting, since the hard scattering does not contain any heavy flavours.
In summary, the three classes above are then characterized by having 2, 1 or 0, respectively, heavy flavours in the final state of the LO hard subprocess. Another way to proceed is to add next-to-leading order (NLO) perturbative processes, i.e the $`𝒪(\alpha _\mathrm{s}^3)`$ corrections to the parton fusion . However, with our currently available set of calculational tools, the NLO approach is not so well suited for exclusive Monte Carlo studies where hadronization is added to the partonic picture.
Flavour excitation and gluon splitting give significant contributions to the total b cross section at LHC energies and thus must be considered when this is of interest, see the following. However, NLO calculations probably do a better job on the total b cross section itself (while, for the lighter $`c`$ quark, production in parton showers is so large that the NLO cross sections are more questionable). The shapes of single heavy quark spectra are not altered as much as the correlations between $`\mathrm{Q}`$ and $`\overline{\mathrm{Q}}`$ when extra production channels are added. Similar observations have been made when comparing NLO to LO calculations . Likewise, asymmetries between single heavy quarks are also not changed much by adding further production channels, so for simplicity we consider only the pair creation process here.
After an event has been generated at the parton level we add fragmentation to obtain a hadronic final state. We use the Lund string fragmentation model. Its effects on charm production were described in . Here we only summarize the main points.
In the string model, confinement is implemented by spanning strings between the outgoing partons. These strings correspond to a Lorentz-invariant description of a linear confinement potential with string tension $`\kappa 1`$ GeV/fm. Each string piece has a colour charge at one end and its anticolour at the other. The double colour charge of the gluon corresponds to it being attached to two string pieces, while a quark is only attached to one. A diquark is considered as being in a colour antitriplet representation, and thus behaves (in this respect) like an antiquark. Then each string contains a colour triplet endpoint, a number (possibly zero) of intermediate gluons and a colour antitriplet end. An event will normally contain several separate strings, especially at high energies where $`\mathrm{g}\mathrm{q}\overline{\mathrm{q}}`$ splittings occur frequently in the parton shower.
The string topology can be derived from the colour flow of the hard process with some ambiguity arising from colour-suppressed terms. Consider e.g. the LO process $`\mathrm{gg}\mathrm{b}\overline{\mathrm{b}}`$ where two distinct colour topologies are possible. Representing the proton remnant by a $`\mathrm{u}`$ quark and a $`\mathrm{ud}`$ diquark (alternatively $`\mathrm{d}`$ plus $`\mathrm{uu}`$), one possibility is to have the three strings $`\mathrm{b}`$$`\mathrm{ud}`$, $`\overline{\mathrm{b}}`$$`\mathrm{u}`$ and $`\mathrm{u}`$$`\mathrm{ud}`$, fig. 27, and the other is identical except the $`\mathrm{b}`$ is instead connected to the $`\mathrm{ud}`$ diquark of the other proton because the initial state is symmetric.
Once the string topology has been determined, the Lund string fragmentation model can be applied to describe the nonperturbative hadronization. To first approximation, we assume that the hadronization of each colour singlet subsystem, i.e. string, can be considered separately from that of all the other subsystems. Presupposing that the fragmentation mechanism is universal, i.e. process-independent, the good description of $`\mathrm{e}^+\mathrm{e}^{}`$ annihilation data should carry over. The main difference between $`\mathrm{e}^+\mathrm{e}^{}`$ and hadron–hadron events is that the latter contain beam remnants which are colour-connected with the hard-scattering partons.
Depending on the invariant mass of a string, practical considerations lead us to distinguish the following three hadronization prescriptions:
In the ideal situation, each string has a large invariant mass. Then the standard iterative fragmentation scheme, for which the assumption of a continuum of phase-space states is essential, works well. The average multiplicity of hadrons produced from a string increases linearly with the string ‘length’, which means logarithmically with the string mass. In practice, this approach can be used for all strings above some cutoff mass of a few GeV.
If a string is produced with a small invariant mass, perhaps only a single two-body final state is kinematically accessible. In this case the standard iterative Lund scheme is not applicable. We call such a low-mass string a cluster and consider its decay separately. When kinematically possible, a $`\mathrm{Q}`$$`\overline{\mathrm{q}}`$ cluster will decay into one heavy and one light hadron by the production of a light $`\mathrm{q}\overline{\mathrm{q}}`$ pair in the colour force field between the two cluster endpoints with the new quark flavour selected according to the same rules as in normal string fragmentation. The $`\overline{\mathrm{q}}`$ cluster end or the new $`\mathrm{q}\overline{\mathrm{q}}`$ pair may also denote a diquark. In the latest version of Pythia, anisotropic decay of a cluster has been introduced, where the mass dependence of the anisotropy has been matched to string fragmentation.
This is the extreme case of cluster decay, where the string mass is so small that the cluster cannot decay into two hadrons. It is then assumed to collapse directly into a single hadron which inherits the flavour contents of the string endpoints. The original continuum of string/cluster masses is replaced by a discrete set of hadron masses, mainly $`\mathrm{B}`$ and $`\mathrm{B}^{}`$ (or the corresponding baryon states). This mechanism plays a special rôle since it allows flavour asymmetries favouring hadron species that can inherit some of the beam-remnant flavour contents. Energy and momentum is not conserved in the collapse so that some energy-momentum has to be taken from, or transferred to, the rest of the event. In the new version, a scheme has been introduced where energy and momentum are shuffled locally in an event.
We assume that the nonperturbative hadronization process does not change the perturbatively calculated total rate of bottom production. By local duality arguments , we further presume that the rate of cluster collapse can be obtained from the calculated rate of low-mass strings. In the process $`\mathrm{e}^+\mathrm{e}^{}\mathrm{c}\overline{\mathrm{c}}`$ local duality suggests that the sum of the $`\mathrm{J}/\psi `$ and $`\psi ^{}`$ cross sections approximately equal the perturbative $`\mathrm{c}\overline{\mathrm{c}}`$ production cross section in the mass interval below the $`\mathrm{D}\overline{\mathrm{D}}`$-threshold. Similar arguments have also been proposed for $`\tau `$ decay to hadrons and shown to be accurate. In the current case, the presence of other strings in the event also allows soft-gluon exchanges to modify parton momenta as required to obtain the correct hadron masses. Traditional factorization of short- and long-distance physics would then also preserve the total bottom cross section. Local duality and factorization, however, do not specify how to conserve the overall energy and momentum of an event when a continuum of $`\overline{\mathrm{b}}\mathrm{d}`$ masses is to be replaced by a discrete $`\mathrm{B}^0`$. In practice, however, the different possible hadronization mechanisms do not affect asymmetries much. The fraction of the string-mass distribution below the two particle threshold effectively determines the total rate of cluster collapse and therefore the asymmetry.
The cluster collapse rate depends on several model parameters. The most important ones are listed here with the Pythia parameter values that we have used. The Pythia parameters are included in the new default parameter set in Pythia 6.135 and later versions.
* Quark masses The quark masses affect the threshold of the string-mass distribution. Changing the quark mass shifts the string-mass threshold relative to the fixed mass of the lightest two-body hadronic final state of the cluster. Smaller quark masses imply larger below-threshold production and an increased asymmetry. The new default masses are PMAS(1)$`=m_\mathrm{u}=`$ PMAS(2)$`=m_\mathrm{d}=`$ 0.33D0, PMAS(3)$`=m_\mathrm{s}=`$ 0.5D0, PMAS(4)$`=m_\mathrm{c}=`$ 1.5D0 and PMAS(5)$`=m_\mathrm{b}=`$ 4.8D0.
* Width of the primordial $`k_{}`$ distribution. If the incoming partons are given small $`p_{}`$ kicks in the initial state, asymmetries can appear at larger $`p_{}`$ since the beam remnants are given compensating $`p_{}`$ kicks, thus allowing collapses at larger $`p_{}`$. The new parameters are PARP(91)=1.D0 and PARP(93)=5.D0.
* Beam remnant distribution functions (BRDF). When a gluon is picked out of the proton, the rest of the proton forms a beam remnant consisting, to first approximation, of a quark and a diquark. How the remaining energy and momentum should be split between these two is not known from first principles. We therefore use different parameterizations of the splitting function and check the resulting variations. We find significant differences only at large rapidities where an uneven energy-momentum splitting tend to shift bottom quarks connected to a beam remnant diquark more in the direction of the beam remnant, hence giving rise to asymmetries at very large rapidities. We use an intermediate scenario in this study, given by MSTP(92)=3.
* Threshold behaviour between cluster decay and collapse. Consider a $`\mathrm{b}\overline{\mathrm{d}}`$ cluster with an invariant mass at, or slightly above, the two particle threshold. Should this cluster decay to two hadrons or collapse into one? In one extreme point of view, a $`\mathrm{B}\pi `$ pair should always be formed when above this threshold, and never a single $`\mathrm{B}`$. In another extreme, the two-body fraction would gradually increase at a succession of thresholds: $`\mathrm{B}\pi `$, $`\mathrm{B}^{}\pi `$, $`\mathrm{B}\rho `$, $`\mathrm{B}^{}\rho `$, etc., where the relative probability for each channel is given by the standard flavour and spin mixture in string fragmentation. In our current default model, we have chosen to steer a middle course by allowing two attempts (MSTJ(17)=2) to find a possible pair of hadrons. Thus a fraction of events may collapse to a single resonance also above the $`\mathrm{B}\pi `$ threshold, but $`\mathrm{B}\pi `$ is effectively weighted up. If a large number of attempts had been allowed (this can be varied using the free parameter MSTJ(17)), collapse would only become possible for cluster masses below the $`\mathrm{B}\pi `$ threshold.
The colour connection between the produced heavy quarks and the beam remnants in the string model gives rise to an effect called beam remnant drag. In an independent fragmentation scenario the light cone energy momentum of the quark is simply scaled by some factor picked from a fragmentation function. Thus, on average the rapidity is conserved in the fragmentation process. This is not necessarily so in string fragmentation, where both string ends contribute to the four-momentum of the produced heavy hadron. If the other end of the string is a beam remnant, the hadron will be shifted in rapidity in the direction of the beam remnant resulting in an increase in $`|y|`$. This beam-drag is shown qualitatively in fig. 28, where the rapidity shift is shown as a function of rapidity and transverse momentum. This shift is not directly accessible experimentally, only indirectly as a discrepancy between the shape of perturbatively calculated quark distributions and the data.
### 5.3 Intrinsic Heavy Quarks
The wavefunction of a hadron in QCD can be represented as a superposition of Fock state fluctuations, e.g. $`|n_\mathrm{V}`$, $`|n_\mathrm{V}\mathrm{g}`$, $`|n_\mathrm{V}\mathrm{Q}\overline{\mathrm{Q}}`$, …components where $`n_\mathrm{V}\mathrm{uud}`$ for a proton. When the projectile scatters in the target, the coherence of the Fock components is broken and the fluctuations can hadronize either by uncorrelated fragmentation as for leading twist production or coalescence with spectator quarks in the wavefunction . The intrinsic heavy quark Fock components are generated by virtual interactions such as $`\mathrm{gg}\mathrm{Q}\overline{\mathrm{Q}}`$ where the gluons couple to two or more projectile valence quarks. Intrinsic $`\mathrm{Q}\overline{\mathrm{Q}}`$ Fock states are dominated by configurations with equal rapidity constituents so that, unlike sea quarks generated from a single parton, the intrinsic heavy quarks carry a large fraction of the parent momentum .
The frame-independent probability distribution of an $`n`$–particle $`\mathrm{b}\overline{\mathrm{b}}`$ Fock state is
$`{\displaystyle \frac{dP_{\mathrm{ib}}^n}{dx_i\mathrm{}dx_n}}=N_n{\displaystyle \frac{\delta (1_{i=1}^nx_i)}{(m_h^2_{i=1}^n(\widehat{m}_i^2/x_i))^2}},`$ (26)
where $`\widehat{m}_i^2=k_{,i}^2+m_i^2`$ is the effective transverse mass of the $`i^{\mathrm{th}}`$ particle and $`x_i`$ is the light-cone momentum fraction. The probability, $`P_{\mathrm{ib}}^n`$, is normalized by $`N_n`$ and $`n=5`$ for baryon production from the $`|n_\mathrm{V}\mathrm{b}\overline{\mathrm{b}}`$ configuration. The delta function conserves longitudinal momentum. The dominant Fock configurations are closest to the light-cone energy shell and therefore the invariant mass, $`M^2=_i\widehat{m}_i^2/x_i`$, is minimized. Assuming $`\stackrel{}{k}_{,i}^2`$ is proportional to the square of the constituent quark mass, we choose $`\widehat{m}_\mathrm{q}=0.45`$ GeV, $`\widehat{m}_\mathrm{s}=0.71`$ GeV, and $`\widehat{m}_\mathrm{b}=5`$ GeV .
The $`x_\mathrm{F}`$ distribution for a single bottom hadron produced from an $`n`$-particle intrinsic bottom state can be related to $`P_{\mathrm{ib}}^n`$ and the inelastic $`\mathrm{pp}`$ cross section by
$`{\displaystyle \frac{\sigma _{\mathrm{ib}}^H(\mathrm{pp})}{dx_\mathrm{F}}}={\displaystyle \frac{dP_H}{dx_\mathrm{F}}}\sigma _{\mathrm{pp}}^{\mathrm{in}}{\displaystyle \frac{\mu ^2}{4\widehat{m}_\mathrm{b}^2}}\alpha _s^4(M_{\mathrm{b}\overline{\mathrm{b}}}).`$ (27)
The probability distribution is the sum of all contributions from the $`|n_\mathrm{V}\mathrm{b}\overline{\mathrm{b}}`$ and the $`|n_\mathrm{V}\mathrm{b}\overline{\mathrm{b}}\mathrm{q}\overline{\mathrm{q}}`$ configurations with $`\mathrm{q}=\mathrm{u}`$, $`\mathrm{d}`$, and $`\mathrm{s}`$ and includes uncorrelated fragmentation and coalescence, as described below, when appropriate . The factor of $`\mu ^2/4\widehat{m}_\mathrm{b}^2`$ arises from the soft interaction which breaks the coherence of the Fock state. We take $`\mu ^20.1`$ GeV<sup>2</sup> . The intrinsic charm probability, $`P_{\mathrm{ic}}^5=0.31`$%, was determined from analyses of the EMC charm structure function data . The intrinsic bottom probability is scaled from the intrinsic charm probability by the square of the transverse masses, $`P_{\mathrm{ib}}=P_{\mathrm{ic}}(\widehat{m}_\mathrm{c}/\widehat{m}_\mathrm{b})^2`$. The intrinsic bottom cross section is reduced relative to the intrinsic charm cross section by a factor of $`\alpha _\mathrm{s}^4(M_{\mathrm{b}\overline{\mathrm{b}}})/\alpha _\mathrm{s}^4(M_{\mathrm{c}\overline{\mathrm{c}}})`$ . Taking these factors into account, we obtain $`\sigma _{\mathrm{ib}}^5(pN)7`$ nb at 14 TeV.
There are two ways of producing bottom hadrons from intrinsic $`\mathrm{b}\overline{\mathrm{b}}`$ states. The first is by uncorrelated fragmentation. If we assume that the $`\mathrm{b}`$ quark fragments into a $`\mathrm{B}`$ meson, the $`\mathrm{B}`$ distribution is
$`{\displaystyle \frac{dP_{\mathrm{ib}}^{nF}}{dx_\mathrm{B}}}={\displaystyle 𝑑z\underset{i=1}{\overset{n}{}}dx_i\frac{dP_{\mathrm{ib}}^n}{dx_1\mathrm{}dx_n}\frac{D_{B/b}(z)}{z}\delta (x_\mathrm{B}zx_\mathrm{b})},`$ (28)
These distributions are assumed for all intrinsic bottom production by uncorrelated fragmentation with $`D_{H/b}(z)=\delta (z1)`$. At low $`p_{}`$, this approximation should not be too bad, as seen in fixed target production .
If the projectile has the corresponding valence quarks, the bottom quark can also hadronize by coalescence with the valence spectators. The coalescence distributions are specific for the individual bottom hadrons. It is reasonable to assume that the intrinsic bottom Fock states are fragile and can easily materialize into bottom hadrons in high-energy, low momentum transfer reactions through coalescence. The coalescence contribution to bottom hadron production is
$`{\displaystyle \frac{dP_{\mathrm{ib}}^{nC}}{dx_H}}={\displaystyle \underset{i=1}{\overset{n}{}}dx_i\frac{dP_{\mathrm{ib}}^n}{dx_1\mathrm{}dx_n}\delta (x_Hx_{H_1}\mathrm{}x_{H_{n_\mathrm{V}}})}.`$ (29)
where the coalescence function is simply a delta function combining the momentum fractions of the quarks in the Fock state configuration that make up the valence quarks of the final-state hadron.
Not all bottom hadrons can be produced from the minimal intrinsic bottom Fock state configuration, $`|n_\mathrm{V}\mathrm{b}\overline{\mathrm{b}}`$. However, coalescence can also occur within higher fluctuations of the intrinsic bottom Fock state. For example, in the proton, the $`\mathrm{B}^{}`$ and $`\mathrm{\Xi }_\mathrm{b}^0`$ can be produced by coalescence from $`|n_\mathrm{V}\mathrm{b}\overline{\mathrm{b}}\mathrm{u}\overline{\mathrm{u}}`$ and $`|n_\mathrm{V}\mathrm{b}\overline{\mathrm{b}}\mathrm{s}\overline{\mathrm{s}}`$ configurations. These higher Fock state probabilities can be obtained using earlier results on $`\psi \psi `$ pair production . If all the measured $`\psi \psi `$ pairs arise from $`|n_\mathrm{V}\mathrm{c}\overline{\mathrm{c}}\mathrm{c}\overline{\mathrm{c}}`$ configurations, $`P_{\mathrm{icc}}4.4\%P_{\mathrm{ic}}`$ . It was found that the probability of a $`|n_\mathrm{V}\mathrm{c}\overline{\mathrm{c}}\mathrm{q}\overline{\mathrm{q}}`$ state was then $`P_{\mathrm{icq}}=(\widehat{m}_\mathrm{c}/\widehat{m}_\mathrm{q})^2P_{\mathrm{icc}}`$ . If we then assume $`P_{\mathrm{ibq}}=(\widehat{m}_\mathrm{c}/\widehat{m}_\mathrm{b})^2P_{\mathrm{icq}}`$, we find that
$`P_{\mathrm{ibq}}\left({\displaystyle \frac{\widehat{m}_\mathrm{c}}{\widehat{m}_\mathrm{b}}}\right)^2\left({\displaystyle \frac{\widehat{m}_\mathrm{c}}{\widehat{m}_\mathrm{q}}}\right)^2P_{\mathrm{icc}},`$ (30)
leading to $`P_{\mathrm{ibu}}=P_{\mathrm{ibd}}70.4\%P_{\mathrm{ib}}`$ and $`P_{\mathrm{ibs}}28.5\%P_{\mathrm{ib}}`$. To go to still higher configurations, one can make similar assumptions. However, as more partons are included in the Fock state, the coalescence distributions soften and approach the fragmentation distributions, eventually producing bottom hadrons with less momentum than uncorrelated fragmentation from the minimal $`\mathrm{b}\overline{\mathrm{b}}`$ state if a sufficient number of $`\mathrm{q}\overline{\mathrm{q}}`$ pairs are included. There is then no longer any advantage to introducing more light quark pairs into the configuration—the relative probability will decrease while the potential gain in momentum is not significant. Therefore, we consider production by fragmentation and coalescence from the minimal state and the next higher states with $`\mathrm{u}\overline{\mathrm{u}}`$, $`\mathrm{d}\overline{\mathrm{d}}`$ and $`\mathrm{s}\overline{\mathrm{s}}`$ pairs.
The probability distributions entering Eq. (27) for $`\mathrm{B}^0`$ and $`\overline{\mathrm{B}}^0`$ are
$`{\displaystyle \frac{dP_{\mathrm{B}^0}}{dx_\mathrm{F}}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{1}{10}}{\displaystyle \frac{dP_{\mathrm{ib}}^{5F}}{dx_\mathrm{F}}}+{\displaystyle \frac{1}{4}}{\displaystyle \frac{dP_{\mathrm{ib}}^{5C}}{dx_\mathrm{F}}}\right)+{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{1}{10}}{\displaystyle \frac{dP_{\mathrm{ibu}}^{7F}}{dx_\mathrm{F}}}+{\displaystyle \frac{1}{5}}{\displaystyle \frac{dP_{\mathrm{ibu}}^{7C}}{dx_\mathrm{F}}}\right)`$ (31)
$`+{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{1}{10}}{\displaystyle \frac{dP_{\mathrm{ibd}}^{7F}}{dx_\mathrm{F}}}+{\displaystyle \frac{2}{5}}{\displaystyle \frac{dP_{\mathrm{ibd}}^{7C}}{dx_\mathrm{F}}}\right)+{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{1}{10}}{\displaystyle \frac{dP_{\mathrm{ibs}}^{7F}}{dx_\mathrm{F}}}+{\displaystyle \frac{1}{5}}{\displaystyle \frac{dP_{\mathrm{ibs}}^{7C}}{dx_\mathrm{F}}}\right)`$
$`{\displaystyle \frac{dP_{\overline{\mathrm{B}}^0}}{dx_\mathrm{F}}}`$ $`=`$ $`{\displaystyle \frac{1}{10}}{\displaystyle \frac{dP_{\mathrm{ib}}^{5F}}{dx_\mathrm{F}}}+{\displaystyle \frac{1}{10}}{\displaystyle \frac{dP_{\mathrm{ibu}}^{7F}}{dx_\mathrm{F}}}+{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{1}{10}}{\displaystyle \frac{dP_{\mathrm{ibd}}^{7F}}{dx_\mathrm{F}}}+{\displaystyle \frac{1}{8}}{\displaystyle \frac{dP_{\mathrm{ibd}}^{7C}}{dx_\mathrm{F}}}\right)+{\displaystyle \frac{1}{10}}{\displaystyle \frac{dP_{\mathrm{ibs}}^{7F}}{dx_\mathrm{F}}}.`$ (32)
See Ref. for more details and the probability distributions of other bottom hadrons.
### 5.4 Model predictions
In this section we present some results from both models. Figure 29 shows the asymmetry between $`\mathrm{B}^0`$ and $`\overline{\mathrm{B}}^0`$ as a function of $`y`$ for several $`p_{}`$ cuts in the string model. The asymmetry is essentially zero for central rapidities and increases slowly with rapidity. When the kinematical limit is approached, the asymmetry changes sign for small $`p_{}`$ because of the drag-effect since $`\mathrm{b}`$-quarks are often connected to diquarks from the proton beam remnant, fig. 27, thus producing $`\overline{\mathrm{B}}^0`$ hadrons which are shifted more in rapidity than $`\mathrm{B}^0`$. Cluster collapse, on the other hand, tend to enhance the production of leading particles (in this case $`\mathrm{B}^0`$) so the two mechanisms give rise to asymmetries with different signs. Collapse is the main effect at small rapidities while eventually at very large $`y`$, the drag effect dominates.
In Table 12 we study the parameter dependence of the asymmetry by looking at the integrated asymmetry for different kinematical regions using three different parameter sets:
* Set 1 is the new default as presented in section 5.2.
* Set 2 The same as Set 1 except it uses simple counting rules in the beam remnant splitting, i.e. each quark get on average one third of the beam remnant energy-momentum.
* Set 3 The old parameter set, before fitting to fixed-target data, is included as a reference. This set is characterized by current algebra masses, lower intrinsic $`k_{}`$, and an uneven sharing of beam remnant energy-momentum.
We see that in the central region the asymmetry is generally very small whereas for forward (but not extremely forward) rapidities and moderate $`p_{}`$ the asymmetry is around 1–2%. In the very forward region at small $`p_{}`$, drag asymmetry dominates which can be seen from the change in sign of the asymmetry. The asymmetry is fairly stable under moderate variations in the parameters even though the difference between the old and new parameter sets (Set 1 and 3) are large in the central region. Set 1 typically gives rise to smaller asymmetries.
The cross sections for all intrinsic bottom hadrons are given as a function of $`x_\mathrm{F}`$ in fig. 30. The bottom baryon distributions are shown in fig. 30(a). The $`\mathrm{\Lambda }_\mathrm{b}^0`$ ($`\mathrm{\Sigma }_\mathrm{b}^0`$) distributions are the largest and most forward peaked of all the distributions. The $`\mathrm{\Sigma }_\mathrm{b}^{}`$ is the smallest and the softest, similar to that of the bottom-strange mesons and baryons shown in fig. 30(b). The different coalescence probabilities assumed for hadrons from the $`|\mathrm{uudb}\overline{\mathrm{b}}\mathrm{s}\overline{\mathrm{s}}`$ configuration have little real effect on the shape of the cross section, dominated by independent fragmentation. Of the $`\mathrm{B}`$ mesons shown in fig. 30(c), the $`\mathrm{B}^+`$ and $`\mathrm{B}^0`$ cross sections are the largest since both can be produced from the 5 particle configuration. The $`\mathrm{B}^{}`$ and $`\overline{\mathrm{B}}^0`$ distributions are virtually identical. We note that the $`x_\mathrm{F}`$ distributions of other bottom hadrons not included in the figure would be similar to the bottom-strange hadrons since they would be produced by fragmentation only.
The $`x_\mathrm{F}`$ distribution for final-state hadron $`H`$ is the sum of the leading-twist fusion and intrinsic bottom components,
$`{\displaystyle \frac{d\sigma _{hN}^H}{dx_\mathrm{F}}}={\displaystyle \frac{d\sigma _{\mathrm{lt}}^H}{dx_\mathrm{F}}}+{\displaystyle \frac{d\sigma _{\mathrm{ib}}^H}{dx_\mathrm{F}}}.`$ (33)
The intrinsic bottom cross sections from Section 5.3 are combined with a leading twist calculation using independent fragmentation where drag effects are not included. The leading twist results have been smoothed and extrapolated to large $`x_\mathrm{F}`$ <sup>9</sup><sup>9</sup>footnotetext: Thanks to J. Klay at UC Davis for extending the curves to large $`x_\mathrm{F}`$. to facilitate a comparison with the intrinsic bottom calculation. The resulting total $`\mathrm{B}^0`$ and $`\overline{\mathrm{B}}^0`$ distributions are shown in fig. 31, along with the corresponding asymmetry. Note that since the intrinsic heavy quark $`p_{}`$ distributions are more steeply falling than the leading twist, we only consider $`p_{}<5`$ GeV. The distributions are drawn to emphasize the high $`x_\mathrm{F}`$ region where the distributions differ. The asymmetry is $`0.1`$ at $`x_\mathrm{F}0.25`$, corresponding to $`y6.5`$. Therefore, intrinsic bottom should not be a significant source of asymmetries.
### 5.5 Summary
To summarize, we have studied possible production asymmetries between $`\mathrm{b}`$ and $`\overline{\mathrm{b}}`$ hadrons, especially $`\mathrm{B}^0`$ and $`\overline{\mathrm{B}}^0`$, as predicted by the Lund string fragmentation model and the intrinsic heavy quark model. We find negligible asymmetries for central rapidities and large $`p_{}`$ (in general, less than 1%). For some especially favoured kinematical ranges such as $`y>3`$ and $`5<p_{}<10`$ GeV the collapse asymmetry could be as high as 1–2%. Intrinsic bottom becomes important only for $`x_\mathrm{F}>0.25`$ and $`p_{}<5`$ GeV, corresponding to $`y>6.5`$.
## 6 QUARKONIUM PRODUCTION<sup>10</sup><sup>10</sup>10Section coordinators: M. Krämer, F. Maltoni, M.A. Sanchis-Lozano
The production of charmonium and bottomonium states at high-energy colliders has been the subject of considerable interest during the past few years. New experimental results from $`p\overline{p}`$, $`ep`$ and $`e^+e^{}`$ colliders have become available, some of which revealed dramatic shortcomings of earlier quarkonium production models. In theory, progress has been made on the factorization between the short distance physics of heavy-quark creation and the long-distance physics of bound state formation. The colour-singlet model has been superseded by a consistent and rigorous framework, based on the use of non-relativistic QCD (NRQCD) , an effective field theory that includes the so-called colour-octet mechanisms. On the other hand, the colour evaporation model of the early days of quarkonium physics has been revived . However, despite the recent theoretical and experimental developments the range of applicability of the different approaches is still subject to debate, as is the quantitative verification of factorization. Because the quarkonium mass is still not very large with respect to the QCD scale, in particular for the charmonium system, non-factorizable corrections may not be suppressed enough, if the quarkonium is not part of an isolated jet, and the expansions in NRQCD may not converge very well. In this situation a global analysis of various processes is mandatory in order to assess the importance of different quarkonium production mechanisms, as well as the limitations of a particular theoretical framework (for reviews on different quarkonium production processes see e.g. .) By the time the LHC starts operating, new experimental data from the Tevatron and HERA as well as theoretical progress, e.g. in the calculation of higher-order corrections, will have significantly improved the present picture and will allow more precise predictions than what is possible at present. In the following, we will therefore focus on the general phenomenological implications of the NRQCD approach for quarkonium production at the LHC, rather than aiming at a detailed and comprehensive numerical analysis. Based on the information provided by the present Tevatron data we will derive predictions for observables crucial for future LHC analyses, such as differential cross sections and quarkonium polarization.
In the NRQCD approach, the cross section for producing a quarkonium state $`H`$ at a hadron collider can be expressed as a sum of terms, each of which factors into a short-distance coefficient and a long-distance matrix element:
$$\text{d}\sigma (pp/p\overline{p}H+X)=\underset{n}{}\text{d}\widehat{\sigma }(pp/p\overline{p}Q\overline{Q}[n]+x)𝒪^H[n],$$
(34)
where $`n`$ denotes the colour, spin and angular momentum state of an intermediate $`Q\overline{Q}`$ pair. The short-distance cross section $`\text{d}\widehat{\sigma }`$ can be calculated perturbatively in the strong coupling $`\alpha _s`$. The NRQCD matrix elements $`𝒪^H[n]`$ (see for their definition) are related to the non-perturbative transition probabilities from the $`Q\overline{Q}`$ state $`n`$ into the quarkonium $`H`$. They scale with a definite power of the intrinsic heavy-quark velocity $`v`$ . ($`v^20.3`$ for charmonium and $`v^20.1`$ for bottomonium.) The general expression (34) is thus a double expansion in powers of $`\alpha _s`$ and $`v`$.
The NRQCD formalism implies that so-called colour-octet processes associated with higher Fock state components of the quarkonium wave function must contribute to the cross section. Heavy quark pairs that are produced at short distances in a colour-octet state can evolve into a physical quarkonium through radiation of soft gluons at late times in the production process, when the quark pair has already expanded to the quarkonium size. Such a possibility is ignored in the colour-singlet model, where only those heavy quark pairs that are produced in the dominant Fock state (i.e. in a colour-singlet state and with the spin and angular momentum quantum numbers of the meson) are assumed to form a physical quarkonium. The most profound theoretical evidence that the colour-singlet model is incomplete comes from the presence of infrared divergences in the production cross sections and decay rates of $`P`$-wave states. Within the NRQCD approach, this problem finds its natural solution since the infrared singularities are factored into a colour-octet operator matrix element . While colour-octet contributions are needed for a consistent description of $`P`$-wave quarkonia, they are phenomenologically even more important for $`S`$-wave states like $`J/\psi `$ or $`\mathrm{{\rm Y}}`$. According to the velocity scaling rules, colour-octet matrix elements for the production of $`S`$-wave quarkonia are suppressed by a factor $`v^4`$ compared to the leading colour-singlet contributions. However, as discussed in some detail below, colour-octet processes can become significant if the short-distance cross section for producing $`Q\overline{Q}`$ in a colour-octet state is enhanced.
The production of $`S`$-wave charmonium in $`p\overline{p}`$ collisions at the Tevatron has attracted considerable attention and has stimulated much of the recent theoretical development in quarkonium physics. The CDF collaboration has measured cross sections for the production of $`J/\psi `$ and $`\psi (2S)`$ states not coming from $`B`$ or radiative $`\chi `$ decays, for a wide range of transverse momenta $`5\text{GeV}\mathrm{\Gamma }<p_t(\psi )\mathrm{\Gamma }<\mathrm{\hspace{0.33em}20}\text{GeV}`$ . Surprisingly, the experimental cross sections were found to be orders of magnitudes larger than the theoretical expectation based on the leading-order colour-singlet model . This result is particularly striking because the data extends out to large transverse momenta where the theoretical analysis is rather clean. The shortcoming of the colour-singlet model can be understood by examining a typical Feynman diagram contributing to the leading-order parton cross section, fig. 32(a).
At large transverse momentum, the two internal quark propagators are off-shell by $`p_t^2`$ so that the parton differential cross section scales like $`\text{d}\sigma /\text{d}p_t^21/p_t^8`$, as indicated in the figure. On the other hand, when $`p_t2m_c`$ the quarkonium mass can be considered small and the inclusive charmonium cross section is expected to scale like any other single-particle inclusive cross section $`1/p_t^4`$. The dominant production mechanism for charmonium at sufficiently large $`p_t`$ must thus be via fragmentation , the production of a parton with large $`p_t`$ which subsequently decays into charmonium and other partons. A typical fragmentation contribution to colour-singlet $`J/\psi `$ production is shown in fig. 32(b). While the fragmentation contributions are of higher order in $`\alpha _s`$ compared to the fusion process fig. 32(a), they are enhanced by a power $`p_t^4/(2m_c)^4`$ at large $`p_t`$ and can thus overtake the fusion contribution at $`p_t2m_c`$. When colour-singlet fragmentation is included, the $`p_t`$ dependence of the theoretical prediction is in agreement with the Tevatron data but the normalization is still underestimated by about an order of magnitude , indicating that an additional fragmentation contribution is still missing. It is now generally believed that gluon fragmentation into colour-octet $`{}_{}{}^{3}S_{1}^{}`$ charm quark pairs , as shown in fig. 32(c), is the dominant source of $`J/\psi `$ and $`\psi (2S)`$ at large $`p_t`$ at the Tevatron. The probability of forming a $`J/\psi `$ particle from a pointlike $`c\overline{c}`$ pair in a colour-octet $`{}_{}{}^{3}S_{1}^{}`$ state is given by the NRQCD matrix element $`𝒪^{J/\psi }[{}_{}{}^{3}S_{1}^{(8)}]`$ which is suppressed by $`v^4`$ relative to the non-perturbative factor of the leading colour-singlet term. However, this suppression is overcompensated for by the gain in two powers of $`\alpha _s/\pi `$ in the short-distance cross section for producing colour-octet $`{}_{}{}^{3}S_{1}^{}`$ charm quark pairs as compared to colour-singlet fragmentation. At $`𝒪(v^4)`$ in the velocity expansion, two additional colour-octet channels have to be included, fig. 32(d), which do not have a fragmentation interpretation at order $`\alpha _s^3`$ but which become significant at moderate $`p_t2m_c`$ . The importance of the $`{}_{}{}^{1}S_{0}^{(8)}`$ and $`{}_{}{}^{3}P_{J}^{(8)}`$ contributions cannot be estimated from naive power counting in $`\alpha _s`$ and $`v`$ alone, but rather follows from the dominance of $`t`$-channel gluon exchange, forbidden in the leading-order colour-singlet cross section.
The different contributions to the $`J/\psi `$ transverse momentum distribution are compared to the CDF data in fig. 33.
As mentioned above, the colour-singlet model at lowest order in $`\alpha _s`$ fails dramatically when confronted with the experimental results. When colour-singlet fragmentation is included, the prediction increases by more than an order of magnitude at large $`p_t`$, but it still falls below the data by a factor of $`30`$. The CDF results on charmonium production can be explained by including the leading colour-octet contributions and adjusting the unknown non-perturbative parameters to fit the data. Numerically one finds the non-perturbative matrix elements to be of $`𝒪(10^2\text{GeV}^3)`$, see Table 13, perfectly consistent with the $`v^4`$ suppression expected from the velocity scaling rules. Similar conclusions can be drawn for $`\psi (2S)`$ production at the Tevatron.
The analysis of the CDF data alone, although very encouraging, does not strictly prove the phenomenological relevance of colour-octet contributions because free parameters have to be introduced to fit the data. However, if factorization holds the non-perturbative matrix elements, Table 13, are universal and can be used to make predictions for various processes and observables. Besides a global analysis of different reactions, the measurement of quarkonium cross sections at the LHC will be crucial to assess the importance of the individual production mechanisms and to test factorization. In fig. 34 we have collected the cross section predictions for direct $`J/\psi `$ and $`\psi (2S)`$ production as well as the production of $`J/\psi `$ from radiative $`\chi `$ decays at the LHC.
The theoretical curves include the statistical errors in the extraction of the NRQCD matrix elements \[Table 13\]. There are, however, additional theoretical uncertainties which might affect the prediction, but which have not yet been fully quantified. In particular the determination of the $`𝒪_8^\psi [{}_{}{}^{1}S_{0}^{}]`$ and $`𝒪_8^\psi [{}_{}{}^{3}P_{0}^{}]`$ matrix elements ($`\psi `$ denoting $`J/\psi `$ or $`\psi (2S)`$) from the Tevatron data is very sensitive to effects that modify the shape of the charmonium $`p_t`$ distribution at relatively small $`p_t\mathrm{\Gamma }<\mathrm{\hspace{0.33em}8}`$ GeV. Those effects include the small-$`x`$ behaviour of the gluon distribution , the evolution of the strong coupling , as well as systematic effects inherent in NRQCD, such as the inaccurate treatment of the energy conservation in the hadronization of the colour-octet $`c\overline{c}`$ pairs . Moreover, higher-order QCD corrections are expected to play an important role, as discussed in more detail below. The cross sections collected in fig. 34 should thus not be viewed as firm NRQCD predictions but will be refined as more experimental and theoretical information on charmonium production becomes available over the next few years.
The inclusion of higher-order QCD corrections is required to reduce the theoretical uncertainty and to allow a more precise prediction of the LHC cross sections. Next-to-leading order (NLO) calculations for quarkonium production at hadron colliders are presently available only for total cross sections . Significant higher-order corrections to differential distributions are expected from the strong renormalization and factorization scale dependence of the leading-order results . Moreover, the NLO colour-singlet cross section includes processes like $`g+gQ\overline{Q}[{}_{}{}^{3}S_{1}^{(1)}]+g+g`$ which are dominated by $`t`$-channel gluon exchange and scale as $`\alpha _s^4(2m_Q)^2/p_t^6`$. At $`p_t2m_Q`$ their contribution is enhanced with respect to the the leading-order cross section, fig. 32(a), which scales as $`\alpha _s^3(2m_Q)^4/p_t^8`$. This is born out by preliminary studies which include part of the NLO hadroproduction cross section and by the complete calculation of NLO corrections to the related process of quarkonium photoproduction . The NLO colour-singlet cross section may be comparable in size to the colour-octet $`{}_{}{}^{1}S_{0}^{}`$ and $`{}_{}{}^{3}P_{J}^{}`$ processes, which scale as $`\alpha _s^3v^4(2m_Q)^2/p_t^6`$ (see fig. 32(d)), and affect the determination of the corresponding NRQCD matrix elements from the Tevatron data. A full NLO analysis is however needed before quantitative conclusions can be drawn.
Another source of potentially large higher-order corrections is the multiple emission of soft or almost collinear gluons from the initial state partons. These corrections, as well as effects related to intrinsic transverse momentum, are expected to modify the shape of the transverse momentum distribution predominantly at relatively low values of $`p_t\mathrm{\Gamma }<\mathrm{\hspace{0.33em}2}m_Q`$. Initial state radiation can be partially summed in perturbation theory , but so far only total cross sections have been considered in the literature . An estimate of the effect on the transverse momentum distribution should be provided by phenomenological models where a Gaussian $`k_t`$-smearing is added to the initial state partons. The result of these calculations not only depends on the average $`k_t`$, which enters as a free parameter, but also on the details of how the smearing is implemented. Moreover, a lower cut-off has to be provided which regulates the divergences at $`p_t=0`$. Using the NLO calculation for the total cross section , one can obtain the rough estimate that perturbative Sudakov effects should be confined below $`p_t12`$ GeV for both charmonium and bottomonium production at Tevatron energies. Qualitatively, the inclusion of $`k_t`$-smearing leads to an enhancement of the short distance cross section at small $`p_t`$, which results in smaller values for the fits of the $`𝒪_8^\psi [{}_{}{}^{1}S_{0}^{}]`$ and $`𝒪_8^\psi [{}_{}{}^{3}P_{0}^{}]`$ NRQCD matrix elements . The actual size of the effect, however, turns out to be very different for the two models studied in the literature.
An alternative approach to treat the effect of initial state radiation is by means of Monte Carlo event generators which include multiple gluon emission in the parton shower approximation. Comprehensive phenomenological analyses have been carried out for charmonium production at the Tevatron and at the LHC using the event generator Pythia supplemented by the leading colour-octet processes . The inclusion of initial state radiation as implemented in Pythia leads to an enhancement of the short-distance cross section. The size of the effect is significantly larger than for the Gaussian $`k_t`$-smearing mentioned above, and it extents out to large $`p_t`$. Consequently, the $`𝒪_8^\psi [{}_{}{}^{1}S_{0}^{}]`$ and $`𝒪_8^\psi [{}_{}{}^{3}P_{0}^{}]`$ NRQCD matrix elements estimated from the Monte Carlo analysis of the Tevatron cross sections are significantly lower than the ones listed in Table 13 (see for details).<sup>12</sup><sup>12</sup>12Support is added to the Monte Carlo extraction of the NRQCD matrix elements by analyses of $`J/\psi `$ production in inelastic $`\gamma p`$-scattering and $`B`$ decays , which seem to prefer small values of $`𝒪_8^\psi [{}_{}{}^{1}S_{0}^{}]`$ and $`𝒪_8^\psi [{}_{}{}^{3}P_{0}^{}]`$. Figure 35 shows the individual contributions to the direct $`J/\psi `$ cross section at the LHC as estimated with the Pythia Monte Carlo .
Note that for consistency the curves are based on the NRQCD matrix elements extracted from the Monte Carlo analysis of the Tevatron data rather than on the values of the leading-order fit listed in Table 13. One observes that the final prediction is consistent with the result presented in fig. 34 within errors. The extrapolation of the Tevatron fits to LHC energies seems rather insensitive to the details of the underlying theoretical description, and different approaches yield similar predictions for the LHC cross sections as long as the appropriate NRQCD matrix elements are used. The Monte Carlo implementation should therefore represent a convenient and reliable tool for the experimental simulation of quarkonium production processes at the LHC.
A crucial test of the NRQCD approach to charmonium production at hadron colliders is the analysis of $`J/\psi `$ and $`\psi (2S)`$ polarization at large transverse momentum. Recall that at large $`p_t`$, $`\psi `$ production should be dominated by gluon fragmentation into a colour-octet $`{}_{}{}^{3}S_{1}^{}`$ charm quark pair, fig. 32(c). When $`p_t2m_c`$ the fragmenting gluon is effectively on-shell and transverse. The intermediate $`c\overline{c}`$ pair in the colour-octet $`{}_{}{}^{3}S_{1}^{}`$ state inherits the gluon’s transverse polarization and so does the quarkonium, because the emission of soft gluons during hadronization does not flip the heavy quark spin at leading order in the velocity expansion. Consequently, at large transverse momentum one should observe transversely polarized $`J/\psi `$ and $`\psi (2S)`$ . The polarization can be measured through the angular distribution in the decay $`\psi l^+l^{}`$, given by $`\text{d}\mathrm{\Gamma }/\text{d}\mathrm{cos}\theta 1+\alpha \mathrm{cos}^2\theta `$, where $`\theta `$ denotes the angle between the lepton three-momentum in the $`\psi `$ rest frame and the $`\psi `$ three-momentum in the lab frame. Pure transverse polarization implies $`\alpha =1`$. Corrections to this asymptotic limit due to spin-symmetry breaking and higher order fragmentation contributions have been estimated to be small . The dominant source of depolarization comes from the colour-octet fusion diagrams, fig. 32(d), which are important at moderate $`p_t`$. Still, at $`𝒪(v^4)`$ in the velocity expansion, the polar angle asymmetry $`\alpha `$ can be unambiguously calculated within NRQCD in terms of the three non-perturbative matrix elements \[Table 13\] that have been determined from the unpolarized cross section. In fig. 36 we display the theoretical prediction for $`\alpha `$ in $`\psi (2S)`$ production at the Tevatron as function of the $`\psi (2S)`$ transverse momentum.
No transverse polarization is expected at $`p_t5`$ GeV, but the angular distribution is predicted to change drastically as $`p_t`$ increases. A preliminary measurement from CDF does not support this prediction, but the experimental errors are too large to draw definite conclusions. A similar picture emerges from the analysis of $`J/\psi `$ polarization , where, however, the theoretical analysis is complicated by the fact that the data sample still includes $`J/\psi `$ that have not been produced directly but come from decays of higher excited states .
Polarization measurements are crucial to discriminate the NRQCD approach from the colour evaporation model, where the cross section for a specific charmonium state is given as a universal fraction of the inclusive $`c\overline{c}`$ production cross section integrated up to the open charm threshold. In general, the assumption of a single universal long-distance factor is too restrictive. It implies a universal $`\sigma (\chi _c)/\sigma (J/\psi )`$ ratio, which is not supported by the comparison of charmonium production in hadron-hadron and photon-hadron collisions. Still, since the colour evaporation model allows colour-octet charm quark pairs from gluon fragmentation to hadronize into charmonium, it can describe the $`p_t`$ distribution of the Tevatron data . In contrast to the NRQCD approach, however, the colour evaporation model predicts charmonium to be produced unpolarized. The model assumes unsuppressed gluon emission from the $`c\overline{c}`$ pair during hadronization which randomizes spin and colour. This assumption is clearly wrong in the heavy quark limit where spin symmetry is at work and soft gluon emission does not flip the heavy quark spin. Nonetheless, since the charm quark mass is not very large with respect to the QCD scale, the applicability of heavy quark spin symmetry to charmonium physics has to be tested by confronting the NRQCD polarization signature with experimental data.
To definitely resolve the issue of quarkonium polarization, a high-statistics measurement extending out to large transverse momentum will be necessary. Such a measurement can be carried out at the LHC, where one expects a polarization pattern similar to that predicted for the Tevatron, see fig. 37.
The absence of a substantial fraction of transverse polarization in $`\psi `$ production at large $`p_t`$ would represent a serious problem for the application of the NRQCD factorization approach to the charmonium system and might indicate that the charm quark mass is not large enough for a nonrelativistic approach to work in all circumstances.
The application of NRQCD should be on safer grounds for the bottomonium system. As $`v^20.1`$ for bottomonium, higher-order terms in the velocity expansion (in particular colour-octet contributions) are expected to be less relevant than in the case of charmonium. Cross sections for the production of $`\mathrm{{\rm Y}}`$ states have been measured at the Tevatron in the region $`p_t\mathrm{\Gamma }<\mathrm{\hspace{0.33em}20}\text{GeV}`$ . The leading-order colour-singlet model predictions underestimate the data, the discrepancy being, however, much less significant than in the case of charmonium. Given the large theoretical uncertainties in the cross section calculation, in particular at small $`p_t\mathrm{\Gamma }<M_\mathrm{{\rm Y}}`$, the need for colour-octet contribution is not yet as firmly established as for charmonium production. The inclusion of both next-to-leading order corrections and the summation of soft gluon radiation is required to obtain a realistic description of the $`\mathrm{{\rm Y}}`$ cross section in the $`p_t`$-range probed by present data. Such calculations have not yet been performed, and we have therefore not attempted a systematic fit of the bottomonium NRQCD matrix elements. Our predictions for the $`\mathrm{{\rm Y}}`$ cross section at the LHC, figs. 38,39, are based on a simple choice of the non-perturbative input parameters \[Table 14\] which is consistent with the present experimental information from the Tevatron.
The cross sections should thus not be regarded as firm predictions of NRQCD but rather as order-of-magnitude estimates. The expected theoretical progress and more experimental information will allow a more precise prediction in the near future.
The impact of initial state gluon radiation on the $`\mathrm{{\rm Y}}`$ cross sections at the Tevatron has been estimated by adding a Gaussian $`k_t`$-smearing as discussed previously in the context of charmonium production. An average $`k_t3`$ GeV and a $`K`$-factor $`3`$ are found to bring the leading-order colour-singlet cross section in line with the experimental $`\mathrm{{\rm Y}}(1S,2S)`$ data at $`p_t\mathrm{\Gamma }<M_\mathrm{{\rm Y}}`$ . Similar results have been obtained within a Monte Carlo analysis , leading to significantly lower fit values for the colour-octet NRQCD matrix elements than those determined from a leading-order calculation \[Table 14\]. Moreover, the Monte Carlo results imply that no feeddown from $`\chi `$ states produced through colour-octet $`{}_{}{}^{3}S_{1}^{}`$ $`b\overline{b}`$ states is needed to describe the inclusive $`\mathrm{{\rm Y}}`$ cross section, in contrast to what is found at leading-order. The calculation of next-to-leading order corrections and a systematic treatment of soft gluon radiation within perturbation theory are required to resolve these issues. Figure 40 shows the inclusive $`\mathrm{{\rm Y}}(1S)`$ cross section at the LHC as obtained from the Monte Carlo calculation .
The curves are based on the NRQCD matrix elements extracted from the Monte Carlo analysis of the Tevatron data . As in the case of charmonium production, one observes that the final LHC prediction is consistent with the leading-order result presented in fig. 38 within errors.
Let us finally present the polarization pattern predicted for direct $`\mathrm{{\rm Y}}(1S)`$ production at the LHC, fig. 41, based on the NRQCD matrix elements of Table 14.
Higher-order corrections to the gluon fragmentation function will lead to a small reduction of the transverse polarization at large $`p_t`$ and should be included once precise data become available. If the charmonium mass is indeed not large enough for a nonrelativistic expansion to be reliable, the onset of transverse $`\mathrm{{\rm Y}}`$ polarization at $`p_tM_\mathrm{{\rm Y}}`$ may become the single most crucial test of the NRQCD factorization approach.
In summary, we have discussed some of the phenomenological implications of the NRQCD approach for quarkonium production at the LHC and presented ’state-of-the-art’ predictions for $`\psi `$ and $`\mathrm{{\rm Y}}`$ differential cross sections and polarization.<sup>13</sup><sup>13</sup>13Other processes that have been studied in the literature include quarkonium production in association with photons or electroweak bosons , as well as $`\eta `$ and $`P`$-wave quarkonium production . Among the theoretical issues that need to be addressed in the future are the calculation of higher-order QCD corrections, the summation of higher-order terms in the velocity expansion and quantitative insights in the effect of higher-twist contributions. Besides a global analysis of different production processes and observables at various colliders, quarkonium physics at the LHC will play a crucial role to assess the importance of colour-octet processes and to conclusively test the applicability of non-relativistic QCD and heavy-quark spin symmetry to the charmonium and bottomonium systems.
## 7 PROSPECTS FOR $`b`$ PRODUCTION MEASUREMENTS AT THE LHC<sup>14</sup><sup>14</sup>14Section coordinator: M. Smizanska and P. Vikas
Of the existing and currently proposed accelerator facilities, the LHC will yield the largest rate of $`b`$ quarks. A well defined program for $`b`$ production investigations, and the development of dedicated detection strategies optimised for ATLAS, CMS and LHCb, are required for the succesful exploitation of the rich LHC potential. After an introduction summarising the main physics motivations, we review the detector and trigger features relevant for $`b`$ production in the LHC experiments. The kinematic ranges accessible to the three experiments are then described. Theoretical motivations and possible measurement methods are presented for single $`b`$ quark properties, correlations in $`b`$ production, multiple heavy flavour production, polarization, and charge asymmetry effects in $`B`$-hadron production in $`pp`$ interactions. Based on earlier performance studies, the potential for these measurements is estimated and some preliminary results are presented. We conclude with a summary of the present status of the preparations for $`b`$-production studies.
### 7.1 Introduction
While many LHC studies have been devoted to $`B`$-decays, $`b`$ production has not yet been directly addressed. Even though $`b`$ decay investigations will provide some information on the production, at the discussions of this workshop it became clear that they are not sufficient to cover all aspects of production.
Heavy quark production in high energy hadronic collisions is important for the study of Quantum Chromodynamics (QCD). Nowadays, QCD is recognized as a well established and solid theory. If disagreements between the theoretical predictions and the experimental data are found, they will suggest the lack of understanding only of a particular production mechanism. In many cases these disagreements may be attributed to a too slow convergence of the perturbation series. In other cases, there may be important contributions from nonperturbative effects. Strictly speaking, the production measurements are not going to test the principles of QCD, but rather to outline the boundaries, where the predictions of perturbation theory provide an adequate description and exhaust all the visible effects. In this context, it will be certainly useful to test as many different processes as possible.
We present below some examples of such processes and observables, which can potentially be studied in the LHC experiments. Besides testing QCD, there exist other motivations to understand production properties; for instance, as a control of the systematics in CP violation. Double $`b`$ pair production is also a background in some channels of Higgs detection for LHC . Measurements of the $`b`$ production by ATLAS and CMS in the initial years of low luminosity running will also be used to optimise the trigger selections at high luminosity for rare $`B`$ decays.
### 7.2 Detector and trigger characteristics relevant for $`b`$-production
The ATLAS, CMS and LHCb detectors and triggers are described in detail elsewhere . Even though the signal-to-noise ratio for $`b`$ events is higher at LHC than at lower energy hadron machines, only about $`1\%`$ of the non diffractive inelastic collisions will produce $`b`$-quark pairs. Events with $`B`$ hadrons can be distinguished from other inelastic $`pp`$ interactions by the presence of leptons, of secondary vertices and particles with high $`p_T`$. Each of the three experiments will have several levels of triggers to efficiently select the interesting events containing $`B`$ hadrons while maintaining manageable trigger rates. The information from the muon detectors and the electromagnetic and hadronic calorimeters will be used by the lowest level trigger in all the three experiments. In LHCb the lowest level trigger performs a pile-up veto followed by soft cuts on first level trigger objects like muons ($`p_T>1`$ GeV), electrons ($`E_T>2.1`$ GeV) or hadron clusters ($`E_T>2.4`$ GeV) reducing the trigger rate to 1 MHz. The more time-consuming operations, like vertex reconstruction and using information from the RICH for particle identification, will be performed by the higher level triggers. The final event rate expected from LHCb is $``$ 200 Hz. ATLAS and CMS are central detectors for high $`p_T`$ physics designed to operate at high luminosities. The low-level trigger objects have higher $`p_T`$ limits than in LHCb: single muons $`p_T>6(7)`$ GeV in ATLAS(CMS) or dimuon triggers with a minimal $`p_T`$ of each muon in the interval $`(36)`$ GeV in ATLAS and $`(24)`$ GeV in CMS . However, thanks to the higher luminosity, despite the higher $`p_T`$ thresholds, they will have statistics comparable to LHCb in many exclusive channels. Simulations done on both experiments have demonstrated that at a luminosity of $`10^{33}cm^2s^1`$, in spite of 2-3 pileup events on the average accompanying the $`b`$ event in the same bunch crossing, $`B`$-decays can be triggered on and further cleanly separated from background in off-line reconstruction .
### 7.3 Kinematic ranges
The central detectors ATLAS and CMS will cover the pseudorapidity region $`|\eta |<2.5`$; the more forward LHCb is optimised for $`1.8<\eta <4.9`$. The overlap between the experiments is less then a unit of pseudorapidity, in the region $`1.8<\eta <2.5`$. The low transverse momentum cutoffs in each experiment are limited mainly by the admissible low-level trigger rates. In the statistically dominant channels, ATLAS and CMS will be efficient for $`B`$-hadrons with $`p_T10`$ GeV and LHCb for $`p_T>2`$ GeV. The domains of the Bjorken $`x`$ variable for different values of the $`b`$ quark transverse momentum $`p_T`$ are given in fig. 42 for two situations: when both the $`b`$ and $`\overline{b}`$ are in a fiducial volume of a detector; and when only one of them is there. It is clear that in all three LHC experiments the sampled range of $`x`$ is contained within the region already covered by HERA . For comparison, the analogous distribution is calculated for CDF conditions (fig. 43).
### 7.4 Single $`b`$ quark production
#### 7.4.1 Theoretical motivations
The inclusive differential cross section $`d\sigma /dp_Td\eta `$, where $`p_T`$ and $`\eta `$ are the transverse momentum and the pseudorapidity of the $`b`$ or $`\overline{b}`$ quarks, provide the basic information on $`b`$ production. As discussed in the previous section, next-to-leading order (NLO) calculations give a cross section lower than CDF and D0 data by a factor of $`2.4`$ . However, the shape of the $`p_T`$ distribution is well reproduced by LO+NLO predictions, by a semihard model of the BFKL type and also by Pythia . In the region of high $`p_T`$, the effects of higher order contributions are taken into account by means of the resummation technique . In ref. LO+NLO contributions are included together with the resummation of all terms of order $`\alpha _s^k\mathrm{ln}^k(p_T/m_b)`$ and $`\alpha _s^{k+)}\mathrm{ln}^k(p_T/m_b)`$. These contributions change the shape of the $`p_T`$ spectrum. Thus measurements of high $`p_T`$ single $`b`$-spectra may be considered as a dedicated test for the QCD resummation technique.
#### 7.4.2 Measurement possibilities
Experiments can measure the doubly differential cross section $`d\sigma /dp_Td\eta `$, where $`p_T`$ and $`\eta `$ are the transverse momentum and the pseudorapidity of a $`B`$ hadron, or of a jet associated with a $`B`$ hadron, or only one of the decay products of a $`B`$ hadron (for example $`J/\psi `$ or $`\mu `$). From these experimentally measured quantities, the $`d\sigma /dp_Td\eta `$ of the parent $`b`$-quark can be extracted, using appropriate models of hadronization and decay.
The determination of the absolute value of the cross section is also important. Three independent measurements (ATLAS, CMS and LHCb) can be done at the same energy. The determination of the absolute cross sections is always difficult, since it requires a precise understanding of the luminosity, of the trigger and reconstruction efficiency and of the background contributions. Several techniques of luminosity measurement are under study. It appears that precisions of $`3\%`$ could be achieved . The overlap in the detection phase space of ATLAS, CMS and LHCb, in the region $`1.8<\eta <2.5`$ and $`p_T>10`$ GeV, can be used for cross-checks.
#### 7.4.3 Exclusive channels
From trigger and offline studies and the present experience with CDF it is known that the three LHC experiments can provide high statistics samples of some exclusive $`B`$-decay channels cleanly separated from the background. The statistically dominant channels are those containing $`J/\psi \mu ^+\mu ^{}`$ ($`B_dJ/\psi K^0`$, $`B_dJ/\psi K^{}`$, $`B^\pm J/\psi K^\pm `$ and $`B_s^0J/\psi \varphi `$), which are also needed for CP violation studies. Moreover, LHCb will cleanly separate large statistics of purely hadronic exclusive decays, where the dominant ones are $`B_dD^{}\pi ^+`$ and $`B_dD^{}a_1^+`$. With these processes one can cover the differential $`p_T`$ cross section measurements starting approximately from $`p_T>10`$ GeV for ATLAS and CMS and $`p_T>2`$ GeV for LHCb respectively. The numbers of these events after three years of run at luminosities of $`10^{33}cm^2s^1`$ for ATLAS and CMS and five years at $`210^{32}cm^2s^1`$ for LHCb, are shown in fig. 44 as a function of a minimal transverse momentum of the $`B`$-hadron $`p_T`$.
#### 7.4.4 Inclusive $`bJ/\psi X`$channels
The inclusive channels $`bJ/\psi X`$ can be used to extend the available statistics for production measurements to high transverse momenta (fig. 44).
A preliminary study from CMS shows that for $`p_T^{J/\psi }300`$ GeV, which corresponds to $`p_T^b550`$ GeV, a b-tagging efficiency of $`50\%`$ can be achieved with a $`J/\psi `$ mass and decay length reconstruction. This will give a signal-to-noise ratio of $`2.5`$ taking into account the prediction for prompt $`J/\psi `$ production of ref. .
In ATLAS a study has been done for events $`bJ/\psi (\mu \mu )X`$in which the $`p_T`$ of the $`b`$ quark was chosen larger than 50 GeV. In particular, it was shown that the mass resolution of the $`J/\psi `$ will not be degraded due to events in which a signal reconstructed in the muon system is wrongly associated to a non-muon track in the inner detector.
#### 7.4.5 Inclusive $`b`$-jet production
Another method for $`b`$ production studies discussed at the workshop was based on inclusive $`b`$-jet reconstruction. In both ATLAS and CMS this technique was developed for the Higgs search . The $`b`$-jet cross section is expected to be a small fraction (close to or larger than 2% for jets with $`E_T`$ larger than about 20 GeV) of the single-jet cross section . If this method is to be used for single $`b`$ quark production, it will require prescaling of the trigger for the lower $`p_T`$ region or a cut on very high transverse momenta ($`p_T>150`$ GeV), to reduce the huge rate of non-$`b`$ QCD background .
Figure 45 shows the preliminary results of the CMS $`b`$-tagging efficiency and mistagging probability for high $`E_T`$ jets using the technique described in . The study demonstrates that for tagging efficiencies of $`35\%55\%`$ the mistagging probability is better than $`2\%`$ up to $`E_T200`$ GeV. Beyond that, the $`b`$-tagging efficiency and mistagging probability deteriorate significantly. The algorithm will be further optimised, possibly including lepton identification.
The method of $`b`$ cross section determination based on inclusive $`b`$-jet identification will be heavily dependent on the precise understanding of the non $`b`$-jet rejection factors. Further feasibility studies on this method are necessary.
### 7.5 Correlations in $`b`$ production
#### 7.5.1 Theoretical motivations
As discussed in Section 2, the overall normalisation of the production cross section, as well as the normalisation of the inclusive $`b`$ spectra, remain uncertain within a factor $`2`$ because of inherent theoretical uncertainties. Therefore the measurement of these values does not provide a sringent test of NLO contributions. The expected correlations between the $`b`$ and the $`\overline{b}`$ quarks can be computed in leading and next-to-leading order . The shapes of two-particle distributions are sensitive to the NLO contribution, and thus can be used for these tests. In particular, distributions in the following quantities involving both the $`b`$ and $`\overline{b}`$ quark can be considered: the relative azimuthal distance $`\mathrm{\Delta }\varphi (b\overline{b})<1`$, the pair invariant mass, the pair transverse momentum and the pair rapidity.
#### 7.5.2 Measurement possibilities
The choice of the decay channels is driven by the requirement that the acceptance should not vanish when the $`b`$ and the $`\overline{b}`$ are close in phase space. The goal is to avoid isolation cuts in both trigger and offline algorithms requiring a large separation between the decay products of the $`B`$ and of the $`\overline{B}`$. The processes under consideration are based on the reconstruction of a $`J/\psi `$ originating from the displaced vertex of a $`B`$-hadron, and of an additional lepton coming from the semileptonic decay of the associated $`\overline{B}`$ hadron. For example, in the ATLAS experiment, for an integrated luminosity of $`30\mathrm{f}b^1`$, approximately $`510^5`$ such events are expected, with the exclusively reconstructed $`B`$-decays containing the $`J/\psi `$ (Table 15). CDF and D0 measurements showed that $`b\overline{b}`$ pairs are mostly produced back-to-back . However the region most sensitive to differences between the models is $`\mathrm{\Delta }\varphi (b\overline{b})<1`$ rad, where only $`14\%`$ of the events are expected . The statistics may possibly be increased using the semi-inclusive decays $`b\overline{b}J/\psi X`$accompanied by a lepton (Table 15). As an example we quote recent studies in ATLAS , performed using simulated events with $`B_dJ/\psi K^0`$, $`J/\psi \mu ^+\mu ^{}`$. They indicate that the signal events can indeed be reconstructed in cases when the difference of azimuthal angles between the $`J/\psi `$ and the other muon is small (Fig. 46). It is important to note that no selection cuts requiring model dependent corrections were necessary.
The study can be extended to events with $`J/\psi \mu ^+\mu ^{}`$accompanied by an electron and for $`J/\psi e^+e^{}`$ combined with a muon or an electron. Using all these combinations of leptons will allow the measurement of the same variables by different detectors, leading to an improved control of systematic errors.
### 7.6 Multiple heavy quark pair production
#### 7.6.1 Theoretical motivations
At present, only the leading order calculation, $`O(\alpha _s^4)`$, is available for the $`bb\overline{b}\overline{b}`$ production cross section. The effects of higher order corrections can only be estimated using the event generator Pythia 5.7. Since the predictions of Pythia appear to be about a factor of 10 above the leading order analytical calculations , further theoretical studies are needed.
#### 7.6.2 Measurement possibilities
Events with four $`b`$ quarks can be identified in several ways, the most appropriate one depending upon the context. As a background to Higgs search, the requirement is four $`b`$ jets in the fiducial volume. For the purpose of testing QCD predictions on double $`b`$ production, it may be sufficient to reconstruct events with three $`b`$ quarks in the fiducial volume. For three $`b`$ quarks with $`p_T>10`$ GeV and $`|\eta |<2.5`$, Pythia gives a cross section of $`140`$ nb, which corresponds to 140 events produced per second. Despite this large number, it will be necessary to define features allowing on-line selection of these events in the presence of huge non-$`b`$ and single $`b`$ backgrounds.
As a source of an incorrect tag in CP violation measurements, the relevant $`bb\overline{b}\overline{b}`$ events are those with two $`b`$ hadrons, produced with the same flavour charge, identified in the fiducial volume, while two other $`B`$-hadrons, produced with the opposite flavour, are not detected. A direct measurement of this case could use reconstructed charged mesons or baryons, which are self-tagging. However the expected statistics of these events is insufficient. In fact, double $`b`$ production is expected to be only a minor source of wrong tags . Techniques exist to determine the wrong-tag rate from all processes regardless of its origin, which does not need to be identified .
Similar to the case of double $`b`$ production, the production of doubly heavy hadrons, such as the $`B_c+b+c`$, $`\mathrm{\Xi }_{bc}+b+c`$, etc., refers to an $`O(\alpha _s^4)`$ lowest order QCD process, and also provides a test of perturbative QCD calculations. The question of higher order QCD contributions and, probably, nonperturbative contributions is still open . The total production rate in this case will not be indicative enough to establish the role of different production mechanisms, and the measurement should instead concentrate on the specific event topologies. In particular, various correlations between particles carrying charm and bottom may be of importance.
Measurement possibilities are under investigation for the channel $`B_c^{()}J/\psi \pi `$ . A list of possible semileptonic and nonleptonic $`B_c`$ decays can be found for instance in. The decays $`B_c^{()}J/\psi \mu \nu `$, $`B_c^{()}J/\psi \rho ^+`$, $`B_c^{()}J/\psi K^+`$, $`B_c^{()}J/\psi D_s^+`$ are other potentially interesting modes.
### 7.7 Other measurements
$`B`$ hadrons with non-zero spin can be polarized perpendicular to their production plane. Polarization measurements of $`b`$ hadrons produced in nucleon fragmentation could clarify the problems of different polarization models that failed to reproduce the existing data on strange hyperon production . In particular, information about the quark mass dependence of polarization effects could be obtained. For symmetry reasons, in $`pp`$ collisions this polarization vanishes at zero rapidity, so that the expected observed polarization in ATLAS and CMS will be smaller than in the more forward LHCb. Using the method of helicity analyses of cascade decay $`\mathrm{\Lambda }_b^0\mathrm{\Lambda }^0J/\psi `$ the $`\mathrm{\Lambda }_b^0`$ polarization can be measured in ATLAS with a precision better than 0.016 . Another approach to $`\mathrm{\Lambda }_b^0`$ polarisation measurement, using the same decay channel can be found in .
In proton-proton collisions a charge production asymmetry of $`b`$ hadrons is expected. The asymmetry is defined as the difference of production probabilities of a $`B`$ hadron and its antiparticle. From the theoretical point of view, the asymmetries can provide information on the effects of soft dynamics during the fragmentation and hadronization (i.e., on the soft interactions between the produced $`b`$ quark and the remnants of the disrupted proton). The relevant physical effects are expected to be unimportant in the central rapidity region covered by ATLAS and CMS. In the more forward region of LHCb the asymmetry may rise to a few percent. A detailed theoretical discussion of this issues are given in Section 5.
Any production asymmetry is always measured in the presence of a CP violation asymmetry originating from $`B`$-hadron decays. In some cases these two effects are expected to be of the same order. This is the case, for instance, in the channels $`B_dJ/\psi K^{}(K^+\pi ^{})`$, $`B^+J/\psi K^+`$ and $`\mathrm{\Lambda }_b^0\mathrm{\Lambda }^0J/\psi `$, which are expected to have a small CP violation ($`<1\%`$). A way of estimating the relative size of these two effects may be based on the fact that the production asymmetry varies with the transverse momentum and the rapidity of produced $`b`$-quark, while the decay asymmetry should remain the same. Measurements of such small effects will require good understanding of the possible instrumental detection asymmetries.
### 7.8 Conclusions
The properties of $`b`$ production at the LHC can be measured by the three experiments, which are complementary in phase space. The small overlap region will allow a cross check on the cross section normalization. The kinematic conditions are such that Bjorken $`x`$ values sampled in b-production are above 10<sup>-5</sup>, a region lower than at the Tevatron, but already covered by HERA. Differential cross section measurements using exclusive $`B`$ hadron decays will be most important at small $`p_T`$ values. At high $`p_T`$ values and for correlations and multiple heavy flavour production measurements, the statistics can be increased by semi-inclusive $`B`$ decays containing $`J/\psi `$. Possible methods using $`b`$-jets require further study, to control the non-$`b`$ QCD background. The enormous LHC statistics will also allow to study the production polarization and charge production asymmetries.
## 8 TUNING OF MULTIPLE INTERACTIONS GENERATED BY PYTHIA<sup>16</sup><sup>16</sup>16Section coordinators: P. Bartalini, O. Schneider.
### 8.1 Introduction
The track multiplicity distribution as well as the transverse momentum distribution of charged particles in proton-proton interactions (the so-called minimum bias events) affect the performances of the low level triggers and the detector occupancy of the LHCb experiment . They should therefore be modelled reliably in Monte Carlo programs. In particular, at LHC energies, multiple interactions play an important role, and should not be neglected.
In Section 8.2 we examine the multiple interaction models available in Pythia to describe the event structure in hadron-hadron collisions. In Section 8.3 we select a compilation of homogeneous data at different energies suitable to tune the multiple interactions parameters of Pythia; the tuning procedure is presented in Section 8.4. We use the phenomenological extrapolations at LHC energy in order to get the predictions for the track multiplicity and the transverse momentum distributions in minimum bias and $`b\overline{b}`$ events; these are reported in Section 8.5.
### 8.2 Multiple interaction models
The multiple interactions scenario is needed to describe the multiplicity observables at hadron colliders and is also supported by direct observation . The basic assumption is that several parton-parton interactions can occur within a single hadron-hadron collision. Four different models are available in Pythia. The main parameter of these models, $`P_{T_{min}}`$, is the minimum transverse momentum of the parton-parton collisions; it effectively controls the average number of parton-parton interactions and hence the average track multiplicity. The differences between the four models, which mainly affect the shape of the multiplicity distribution, are the following:
* Model 1 (default in Pythia)
All the hadron collisions are equivalent (as opposed to model 3 and 4 below) and all the parton-parton interactions are independent; the $`P_{T_{min}}`$ parameter represents an abrupt cut-off.
* Model 2
Same as Model 1 but with a continuous turn-off of the cross section at $`P_{T_{min}}`$.
* Model 3
Same as Model 2, but hadronic matter in the colliding hadrons is distributed according to a Gaussian shape, and a varying impact parameter between the two hadrons is assumed.
* Model 4
Same as Model 3 but the matter distribution is described by two concentric Gaussian distributions.
The varying impact parameter models (Models 3 and 4) were developed to fit the UA5 data . A recent study performed by the CDF collaboration concludes that a varying impact parameter model (Model 3) is also preferred to describe the underlying tracks in $`b`$ events produced at the Tevatron.
In the absence of published results on multiplicity distributions in minimum bias events at the Tevatron, we compare again the predictions of Model 1 and Model 3 with the UA5 data, using the final charged multiplicity distribution from the full $`\overline{p}p`$ data sample collected by UA5 at $`\sqrt{s}=546`$ GeV , a recent version of Pythia<sup>18</sup><sup>18</sup>18Version 6.125 was used for all the studies reported here., and a modern set of parton distribution functions, CTEQ4L , tuned on both HERA and Tevatron data. For this comparison, the main multiple interactions parameter ($`P_{T_{min}}`$) is tuned in each model to reproduce the mean value of the measured charged multiplicity distribution in not single diffractive events<sup>19</sup><sup>19</sup>19In this paper we define as “non single-diffractive event” any inelastic hadron-hadron interaction that cannot be regarded as a single diffractive event; in the framework of the Pythia hadronic interactions, the “non single-diffractive” sample includes the $`22`$ partonic processes and the double diffractive hadron-hadron interactions. ($`<N_{ch}>=29.4\pm 0.3\pm 0.9`$). We obtain $`P_{T_{min}}=1.63\pm 0.02`$ for Model 1 and $`P_{T_{min}}=1.97\pm 0.03`$ for Model 3. The shapes of the multiplicity distributions are compared in fig. 47. It is clear that Model 3 is preferred over Model 1 to describe the UA5 data. In particular the shape of the tail at high multiplicities is reproduced much better by Model 3. The UA5 results are corrected for the lower efficiency expected on double diffractive events. Therefore the simulated samples include the generation of all kind of non single-diffractive events. The uncertainty in the diffractive cross sections relative to the partonic ones can affect the observed discrepancies between the data and the Pythia predictions in the low multiplicity region<sup>20</sup><sup>20</sup>20The $`\overline{p}p`$ cross sections predicted by Pythia at $`\sqrt{s}=546`$ GeV are $`30.7`$ mb for the partonic processes and $`5.3`$ mb for the double diffractive processes..
### 8.3 Mean charged multiplicity at $`𝜼\mathbf{=}\mathrm{𝟎}`$
In order to produce realistic Pythia predictions for the multiplicity observables in the LHC environment, it is necessary to take into account the energy dependence of the $`P_{T_{min}}`$ parameter. Unfortunately there are not many published data concerning the charged multiplicity distribution in minimum bias events at hadron colliders. On the other hand there are some data available relative to the average charged multiplicity in non single-diffractive events, in particular for the central pseudo-rapidity region. Therefore, to study the energy dependence of the $`P_{T_{min}}`$ parameter at generator level, we consider an homogeneous sample of corrected average charged multiplicity measurements at six different center-of-mass energies ($`\sqrt{s}=`$ $`50`$, $`200`$, $`546`$, $`630`$, $`900`$ and $`1800`$ GeV) in the pseudo-rapidity region $`|\eta |<0.25`$ . The energy dependence of $`dN_{ch}/d\eta `$ at $`\eta =0`$ is shown in fig. 48a together with the fit of a quadratic function of $`\mathrm{ln}(s)`$ proposed in Reference ; using this fit to extrapolate at LHC energy would predict $`dN_{ch}/d\eta =6.11\pm 0.29`$ at $`\eta =0`$.
### 8.4 Tuning of the multiple interaction parameter $`𝑷_{𝑻_{𝒎𝒊𝒏}}`$
The average charged multiplicity measurements performed on non single-diffractive data in $`\overline{p}p`$ collisions and described in Section 8.3 are used to tune the main multiple interaction parameter in Pythia, $`P_{T_{min}}`$. We generate non single-diffractive events. At each value of $`\sqrt{s}`$, the $`P_{T_{min}}`$ parameter is adjusted to reproduce the average multiplicity measured in the data. The uncertainty on the tuned value of $`P_{T_{min}}`$ reflects the uncertainty on the data. However, the tuned parameters depend on other aspects of the Pythia simulation: in particular the effects of various choices for the multiple interaction model and the parton distribution functions are investigated. For simplicity, the results of these studies are shown only for some representative settings:
* as an example of pre-HERA parton distribution functions we consider the CTEQ2L set used by default in Pythia versions 5.7, but recently retracted by their authors;
* as an example of post-HERA parton distribution functions we consider the CTEQ4L and GRV94L sets, the latter being the new default in Pythia versions 6.1.
This study is restricted to Models 1 and 3 for multiple interactions (see Section 8.2).
The results of the tuning procedure are shown in fig. 48b: in each case, $`P_{T_{min}}`$ appears to be monotonically increasing as a function of $`\sqrt{s}`$. This dependence is much more pronounced for the post-HERA parton distribution functions regardless of the choice for the multiple interactions model.
It was shown in Reference that the post-HERA parton distribution functions imply an energy-dependent $`P_T`$ cut-off. This is heuristically motivated by the Lipatov-like dependence of the gluonic parton distribution function in the small-$`x`$ limit:
$$xg(x,Q^2)\mathrm{constant}\times x^ϵ\mathrm{for}x0$$
(35)
with $`ϵ0.08`$, while the pre-HERA parton distribution functions give a reduced charge screening effect and consequently a less sensitive running of the $`P_T`$ cut-off. This is heuristically motivated by the Regge-like dependence of the gluonic parton distribution function in the small-$`x`$ limit:
$$xg(x,Q^2)\mathrm{constant}\mathrm{for}x0.$$
(36)
In order to extrapolate $`P_{T_{min}}`$ at LHC energy, one needs to find a reasonable function to fit the tuned $`P_{T_{min}}`$ values as a function of $`\sqrt{s}`$ for the different parton distribution functions and multiple interactions models; a four degree of freedom fit is performed using the following exponential form, inspired by the recent implementations added in Pythia since version 6.120 :
$$P_{T_{min}}(\sqrt{s})=P_{T_{min}}^{\mathrm{LHC}}\left(\frac{\sqrt{s}}{14\mathrm{TeV}}\right)^{2ϵ}.$$
(37)
The fitted functions are superimposed on fig. 48b and the results obtained for the fitted parameters $`ϵ`$ and $`P_{T_{min}}^{\mathrm{LHC}}`$ are given in Table 16. This quantitative analysis demonstrates that the power law expressed in Equation 37 holds for values of $`ϵ`$ between $`0.08`$ and $`0.10`$ if post-HERA parton distribution functions are used, and for somewhat smaller values of $`ϵ`$ ($`0.05`$) for the pre-HERA parton distribution functions.
### 8.5 Pythia predictions at LHC energy
Figures 49a and b show multiplicity and pseudorapidity distributions for charged particles predicted by Pythia at LHC with CTEQ4L and Model 3 for multiple interactions. The value of $`P_{T_{min}}`$ at LHC is obtained by an extrapolation
$$P_{T_{min}}^{\mathrm{LHC}}=P_{T_{min}}^{\mathrm{Tevatron}}\left(\frac{14\mathrm{TeV}}{1.8\mathrm{TeV}}\right)^{2ϵ}$$
(38)
where $`P_{T_{min}}^{\mathrm{Tevatron}}`$ is the $`P_{T_{min}}`$ value tuned at the Tevatron energy of $`1.8\mathrm{TeV}`$. For the parameter $`ϵ`$, we adopt the results given in Table 16. It is important to note that the predictions $`dN_{ch}/d\eta =6.30\pm 0.42`$ (for $`ϵ=0.0870.005`$) at $`\eta =0`$ are consistent with the phenomenological fit displayed in fig. 48a ($`dN_{ch}/d\eta =6.11\pm 0.29`$).
In order to demonstrate the importance of the correct $`P_{T_{min}}`$ extrapolation, figs. 49a and b also show results obtained by assuming $`P_{T_{min}}^{\mathrm{LHC}}=P_{T_{min}}^{\mathrm{Tevatron}}`$, i.e. $`ϵ=0`$ not supported by the data as demonstrated in Section 8.4. The multiplicity distribution has a tail at high multiplicities and $`dN_{ch}/d\eta `$ at $`\eta =0`$ is not consistent with that obtained from the phenomenological fit.
Figure 50 shows the same distributions as fig. 49, but for the CTEQ2L parton distribution functions. It is interesting to note that, once the extrapolation of $`P_{T_{min}}`$ is properly done, there is no large difference between the multiplicity and pseudorapidity distributions obtained with different structure functions.
Figure 51a-d compare the Pythia predictions for the multiplicity and transverse momentum distributions in the LHCb angular acceptance ($`1.8<\eta <4.9`$) between minimum bias and $`b\overline{b}`$ events<sup>21</sup><sup>21</sup>21The $`b\overline{b}`$ events are selected among the minimum bias events.. These predictions are obtained with CTEQ4L, multiple interactions Model 3 and the proper $`P_{T_{min}}`$ extrapolation. They show clear differences between minimum bias and $`b\overline{b}`$ events, in particular higher average multiplicity and transverse momentum for $`b\overline{b}`$ events.
Pythia predictions with multiple interactions Model 1 for the multiplicity and transverse momentum distributions are shown in fig. 52 for minimum bias and $`b\overline{b}`$ events. Compared to the results obtained with Model 3, less significant differences between minimum bias and $`b\overline{b}`$ events is observed. In Section 8.2 we have stressed that a varying impact parameter model for multiple interactions (i.e. Model 3) is needed to describe the charged track multiplicity in hadron-hadron interactions. There are arguments in favour of adopting a multiple interactions model with varying impact parameter to describe the heavy flavour production at hadron colliders , though there are no experimental data at low transverse momentum and high pseudorapidity (i.e. in the LHCb acceptance region). A more detailed discussion on the effect of multiparton interactions in $`b\overline{b}`$ events at LHCb can be found in Reference .
### 8.6 Conclusions
Comparisons between Pythia and experimental data demonstrate that, in order to reproduce the charged track multiplicity spectrum, a varying impact parameter model has to be adopted.
The varying impact parameter models predict sensitive differences in multiplicity and $`P_T`$ distribution between light and heavy flavour events.
The running of the $`P_T`$ cut-off parameter in Pythia multiple interactions is mandatory. Predictions made at LHC energy with a fixed $`P_T`$ cut-off tuned at lower energies overestimate the multiplicity observables. Taking into account the running of the $`P_T`$ parameter is even more important if post-HERA parton distribution functions are used.
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# Limitations of the number selfconsistent Random Phase Approximation
## 1 Introduction
The random phase approximation (RPA) and its quasiparticle generalization (QRPA) have been widely used in the last decades to study electromagnetic transitions and beta decays in medium and heavy nuclei . The proton-neutron quasiparticle random phase approximation (pn-QRPA) has been extensively employed in the description of single and double beta decays in vibrational nuclei. However the RPA develops a collapse, i.e. it presents imaginary eigenvalues for strengths beyond a critical value of the force .
A whole family of extensions of the QRPA, called renormalized QRPA (RQRPA) are known that do not develop any collapse by implementing the Pauli principle in a consistent way, beyond the simplest quasiboson approximation . However, in its simplest versions there is a violation of the non energy weighted Ikeda sum rule . Calculations to determine the amount of the violation and some improvements to the RQRPA, in order to restore the sum rule, have been presented . It has been shown that treating simultaneously BCS and QRPA equations one can fulfil the Ikeda sum rule for the Fermi case when a schematic model is used .
In recent articles we studied the expectation values of the particle and quasiparticle numbers, the particle number fluctuations and the number of particle pairs with J=0, T=1 and J=1, T=0 in the ground state of <sup>76</sup>Ge as a function of the residual proton-neutron interaction, using realistic Hilbert spaces . We found an important amount of particle number fluctuations in the RQRPA ground state beyond the QRPA collapse, pointing out a source of uncertainty in the RQRPA results. The analysis of the number of pairs showed that the isoscalar-isovector phase transition found in exact calculations is causing the QRPA collapse and is missed in the RQRPA formalism.
In the present work we go a step further, studying <sup>76</sup>Ge with a renormalized RPA where at the same time the mean field is changed by minimizing the energy and fixing the number of particles in the correlated ground state. While particle number fluctuations are smaller than in the previous cases, they still exhibit a clear increase after the point of collapse. More remarkably, the pairing gap is strongly reduced in comparison with its experimental value, and the Ikeda sum rule is violated. Both results cast serious doubts about the double beta decay half-lives estimated using the QRPA extensions, in particular for those nuclei where standard QRPA calculations collapse . Together with the conclusions obtained in , we get a clear picture of the limitations associated with the QRPA extensions, which at the end are the same common sense ask for: you cannot allow the residual proton-neutron interaction to dictate the composition of the ground state wave function without missing contact with actual nuclei, even if the formalism allows you to overcome the collapse.
In Section 2 the renormalized gap and number equations are introduced, whose relationship with the RQRPA equations is shown in Section 3. Some relevant expectation values and the Ikeda sum rule are discussed in Section 4, results for <sup>76</sup>Ge are presented in Section 5 and the conclusions in Section 6.
## 2 Renormalized Gap and Number Equations
In this section the gap and number equations are obtained minimizing the Hamiltonian expectation value $`0|H|0`$ in the correlated vacuum $`|0`$. The Hamiltonian is
$$H=H_p+H_n+H_{p,n}.$$
(1)
The first two terms refer to the proton and neutron hamiltonians
$$H_t=\underset{𝐭}{}(e_t\lambda )a_𝐭^{}a_𝐭+\frac{1}{2}\underset{𝐭^{}s}{}<𝐭_1𝐭_2|V|𝐭_3𝐭_4>a_{𝐭_1}^{}a_{𝐭_2}^{}a_{𝐭_4}a_{𝐭_3},$$
(2)
where the single particle energies are denoted by $`e_t`$, the chemical potential by $`\lambda `$, and the last term corresponds to the proton-neutron interaction
$$H_{p,n}=\underset{𝐩,𝐩^{},𝐧,𝐧^{}}{}<𝐩,𝐧|V|𝐩^{},𝐧^{}>a_𝐩^{}a_𝐧^{}a_𝐧^{}a_𝐩^{}.$$
(3)
The subscripts $`𝐭(t)`$ stand for $`𝐩(p)`$(protons) or $`𝐧(n)`$ (neutrons), being $`𝐭t,m_t`$, with $`t\{n_t,l_t,j_t\}`$ and $`m_tm_j`$.
Through the Bogoliubov transformation
$$\alpha _𝐭^{}=u_ta_𝐭^{}v_ta_{\overline{𝐭}},$$
(4)
with $`a_{\overline{𝐭}}=(1)^{t+m_t}a_{t,mt}`$, we get
$$H=U+\underset{p}{}2\mathrm{\Omega }_pϵ_p\widehat{𝒩}_p+\underset{n}{}2\mathrm{\Omega }_nϵ_n\widehat{𝒩}_n+H_{22}+H_{40}+H_{04}+H_{13}+H_{31},$$
(5)
being
$`\widehat{𝒩}_{tt^{}}{\displaystyle \frac{\left[\alpha _t^{}\alpha _{\overline{t^{}}}\right]^0}{\sqrt{2\mathrm{\Omega }_t}}},\widehat{𝒩}_t\widehat{𝒩}_{t=t^{}},\mathrm{\Omega }_t{\displaystyle \frac{(2j_t+1)}{2}},`$
and where $`U`$ and the quasiparticle energies $`ϵ_t`$ are defined as
$`U={\displaystyle \underset{t}{}}\left[2\mathrm{\Omega }_tv_t^2(e_t\lambda )+{\displaystyle \underset{t^{}}{}}\sqrt{\mathrm{\Omega }_t^{}\mathrm{\Omega }_t}v_t^2v_t^{}^2F(tt,t^{}t^{},0)\mathrm{\Omega }_tu_tv_t\mathrm{\Delta }_t\right],`$ (6)
$`ϵ_t=\left[e_t\lambda +{\displaystyle \underset{t^{}}{}}\sqrt{{\displaystyle \frac{\mathrm{\Omega }_t^{}}{\mathrm{\Omega }_t}}}v_t^{}^2F(tt,t^{}t^{},0)\right](v_t^2u_t^2)+2\mathrm{\Delta }_tu_tv_t`$ (7)
being
$`\mathrm{\Delta }_t=1/2{\displaystyle \underset{t^{}}{}}\sqrt{{\displaystyle \frac{\mathrm{\Omega }_t^{}}{\mathrm{\Omega }_t}}}u_t^{}v_t^{}G(tt,t^{}t^{},0),`$ (8)
the ’gap’ and F, G the usual particle-hole(PH) and particle-particle(PP) coupled two-particle matrix elements. The terms $`H_{nm}`$ in (5) destroy $`m`$ quasiparticles and creates $`n`$ quasiparticles, respectively.
To obtain the quasiparticle mean field occupations $`v_i`$ the ground state energy $`0|H|0`$ is minimized, i.e. it is asked that
$`{\displaystyle \frac{}{v}}0|H|0=0,`$ (9)
subject to the constrictions
$$0|\widehat{N}|0=N,0|\widehat{Z}|0=Z,$$
(10)
and the normalization $`u_i^2+v_i^2=1`$, being $`\widehat{N}`$ and $`\widehat{Z}`$ the neutron and proton particle number operators respectively.
The standard BCS procedure treats protons and neutrons separately and in Hamiltonian (5) only pairing interactions, refering in general to like particles interacting through the J=0 channel, are included . In this case the ground state energy is just U as defined in Eq. (6) and $`|0`$ is the BCS ground state. The residual interactions, either between like particles or protons and neutrons, represented by the terms $`H_{nm}`$ in (5), are usually included in a second step, most commonly using the QRPA .
One of the drawbacks of this procedure is that Eq. (10) is only enforced for the BCS vacuum. When residual interactions are present the expectation values of the particle number operators do not coincide with the actual number. The SCRPA is designed to overcome this difficulty by solving Eq (9) and (10) using the RPA vacuum.
Full selfconsistency requires to consider the proton-neutron interaction contribution in the minimization, but the system of equations which describes this problem is nonlinear and rather complicated, and has only be implemented in schematic models . In order to perform calculations in realistic Hilbert spaces we will include in this first step only the like-particles part of the Hamiltonian, while the ground state $`|0`$ will be sensitive to proton-neutron interaction through the RPA equations, following a philosophy close to Ref . From here on we will refer to this approximation as SRQRPA. We are absolutely aware that our treatment is not fully self-consistent. We are just meeting the requirements of Eq. (10) for the RPA vacuum, and including the modifications in the mean field due to pn-correlations at the lowest level. However, as is shown below, even this mild modifications have very important effects in the observables of the system.
Following , but including only the modifications of the gap equations due to use of the RPA vacuum instead of the BCS vacuum, i.e. not taking into account the proton-neutron residual interaction explicitly, we arrive to the modified gap equation
$`2(\overline{e_t}\lambda )v_tu_t(u_t^2v_t^2)\stackrel{~}{\mathrm{\Delta }}_t=0,`$ (11)
where
$`\overline{e}_t=e_t+{\displaystyle \underset{t^{}}{}}\sqrt{{\displaystyle \frac{\mathrm{\Omega }_t^{}}{\mathrm{\Omega }_t}}}v_t^{}^2F(tt,t^{}t^{},0))`$ (12)
is the single particle energy corrected by the self-energy and
$`\stackrel{~}{\mathrm{\Delta }}_t=1/2{\displaystyle \underset{t^{}}{}}\sqrt{{\displaystyle \frac{\mathrm{\Omega }_t^{}}{\mathrm{\Omega }_t}}}u_t^{}v_t^{}G(tt,t^{}t^{},0)(120|\widehat{𝒩}_t^{}|0),`$ (13)
the ’renormalized’ gap. Notice that the pairing interaction appears in the gap equation renormalized by the presence of proton-neutron residual interactions, which introduce a finite number of quasiparticles in the RPA vacuum . For the sake of simplicity we have dropped some higher order terms in Eq. (11) which come from the interaction between like particles connecting the many quasiparticle components of the correlated ground state. The definition of the renormalized gap (13) is unique at this level of approximation . This is an important definition to be kept in mind, because the factor $`(120|\widehat{𝒩}_t^{}|0)`$ will play a definite role in suppressing the gap when proton-neutron residual interactions are large.
What remains is to couple the renormalized gap problem with the RQRPA equations.
## 3 RQRPA
The nuclear excited states are constructed as
$$|\lambda JM\mathrm{\Omega }^{}(\lambda JM)|0,$$
(14)
$$\mathrm{\Omega }^{}(\lambda JM)=\underset{pn}{}\left[X_{pn}(\lambda J)A_{pn}^{}(JM)Y_{pn}(\lambda J)A_{pn}(\overline{JM})\right],$$
(15)
where
$$\widehat{A}_{pn}^{}(JM)[\alpha _p^{}\alpha _n^{}]^{JM}D_{pn}^{1/2},D_{pn}\left(10|\widehat{𝒩}_p+\widehat{𝒩}_n|0\right),$$
(16)
are the renormalized two-quasiparticle proton-neutron creation operators, which satisfy
$$0|[A_{pn}(JM),A_{p^{}n^{}}^{}(J^{}M^{})]|0=\delta _{pp^{}}\delta _{nn^{}}\delta _{JJ^{}}\delta _{MM^{}}.$$
(17)
Here $`|0`$ is the RPA correlated ground state, defined by the condition
$$\mathrm{\Omega }(\lambda JM)|0=0.$$
(18)
Each amplitudes $`X_{pn}(\lambda J)`$ and $`Y_{pn}(\lambda J)`$ are associated with the excitation energy $`\omega _{\lambda J}`$ of the $`\lambda `$-th state with angular momentum $`J`$. They are the eigenvectors and eigenvalues, respectively, of the RPA equations
$`\left(\begin{array}{cc}A(J)\hfill & B(J)\hfill \\ B^{}(J)\hfill & A^{}(J)\hfill \end{array}\right)\left(\begin{array}{c}X(\lambda J)\hfill \\ Y(\lambda J)\hfill \end{array}\right)=\omega _{\lambda J}\left(\begin{array}{c}X(\lambda J)\hfill \\ Y(\lambda J)\hfill \end{array}\right),`$ (25)
where
$`A_{pn,p^{}n^{}}(J)`$ $`=`$ $`(ϵ_p+ϵ_n)\delta _{pp^{}}\delta _{nn^{}}+D_{pn}^{1/2}U_{pn,p^{}n^{}}^{}(J)D_{p^{}n^{}}^{1/2},`$
$`B_{pn,p^{}n^{}}(J)`$ $`=`$ $`D_{pn}^{1/2}U_{pn,p^{}n^{}}^{}(J)D_{p^{}n^{}}^{1/2},`$ (26)
with
$`U_{pn,p^{}n^{}}^{}(J)`$ $`=`$ $`G(pn,p^{}n^{},JM)(u_pu_nu_p^{}u_n^{}+v_pv_nv_p^{}v_n^{})`$
$`+`$ $`F(pn,p^{}n^{},JM)(u_pv_nu_p^{}v_n^{}+v_pu_nv_p^{}u_n^{}),`$
$`U_{pn,p^{}n^{}}^{}(J)`$ $`=`$ $`G(pn,p^{}n^{},JM)(u_pu_nv_p^{}v_n^{}+v_pv_nu_p^{}u_n^{})`$ (27)
$`+`$ $`F(pn,p^{}n^{},JM)(v_pu_nu_p^{}v_n^{}+u_pv_nv_p^{}u_n^{}).`$
The RPA equations (25) depend explicitly on the mean field occupations $`v_i`$ through the A and B matrices defined in (26) and (27). At the same time, the quasiparticle occupations $`v_i`$ depend on the RPA amplitudes $`X,Y`$, on which depend also $`|0`$, through the renormalized gap equation (11) and number equations (10).
## 4 Expectation values in the correlated RPA vacuum
Using the quasi-boson approximation
$$[A_{pn}(JM),A_{p^{}n^{}}^{}(J^{}M^{})]0|[A_{pn}(JM),A_{p^{}n^{}}^{}(J^{}M^{})]|0=\delta _{pp^{}}\delta _{nn^{}}\delta _{JJ^{}}\delta _{MM^{}},$$
(28)
the RPA ground state defined by Eq.(18) can be written as
$$|0=N_0e^{\widehat{S}}|BCS,$$
(29)
with the BCS vacuum defined by the property
$$\alpha _t|BCS=0,$$
(30)
and
$`\widehat{S}={\displaystyle \frac{1}{2}}{\displaystyle \underset{pnp^{}n^{}J}{}}\sqrt{(2J+1)}C(J)_{pnp^{}n^{}}\left[A_{pn}^{}(J)A_{p^{}n^{}}^{}(J)\right]^0,`$ (31)
$`C(J)_{pnp^{}n^{}}={\displaystyle \underset{\lambda }{}}Y(J)_{pn,\lambda }^{}X(J)_{\lambda ,p^{}n^{}}^^1.`$
The quasiparticle occupations $`0|\widehat{𝒩}_t|0`$ evaluated using the QRPA vacuum (29) have the explicit form
$`0|\widehat{𝒩}_p|0`$ $`=`$ $`{\displaystyle \underset{\lambda Jn^{}}{}}{\displaystyle \frac{(2J+1)}{2\mathrm{\Omega }_p}}D_{pn^{}}|Y(J)_{pn^{},\lambda }|^2,`$
$`0|\widehat{𝒩}_n|0`$ $`=`$ $`{\displaystyle \underset{\lambda Jp^{}}{}}{\displaystyle \frac{(2J+1)}{2\mathrm{\Omega }_n}}D_{p^{}n}|Y(J)_{p^{}n,\lambda }|^2,`$ (32)
The mean particle numbers of protons and neutrons are
$`Z0|\widehat{𝒵}|0=2{\displaystyle \underset{p}{}}\mathrm{\Omega }_pv_p^2+2{\displaystyle \underset{p}{}}\mathrm{\Omega }_p(u_p^2v_p^2)0|\widehat{𝒩}_p|0,`$ (33)
$`N0|\widehat{𝒩}|0=2{\displaystyle \underset{n}{}}\mathrm{\Omega }_nv_n^2+2{\displaystyle \underset{n}{}}\mathrm{\Omega }_n(u_n^2v_n^2)0|\widehat{𝒩}_n|0,`$
When the BCS vacuum is used, the second term in each of these expressions vanish and one is left with the usual number equation. When the residual interaction is present these terms have a relevant contribution to the particle number . They are included in the present work in a selfconsistent way by modifying the mean field occupations $`v_p,v_n`$ to obtain the correct number of particles in the correlated RPA vacuum.
The Fermi transition operators are written in terms of the pair creation and annihilation operators as
$$\tau ^{}=\underset{pn}{}[v_pu_n[\alpha _p^{}\alpha _n^{}]^{00}+u_pv_n[\alpha _p\alpha _n]^{00}],\tau ^+=\{\tau ^{}\}^{}$$
(34)
which are not the exact expressions because they are missing the scattering terms which create a proton (neutron) quasiparticle and annihilates a neutron (proton) quasiparticle.
The total strengths $`S_\pm `$ associated with these transition operators are
$$S_\pm =\underset{\lambda }{}|\lambda J=0|\tau ^\pm |0|^2.$$
(35)
The Ikeda sum rule states that, when the exact operators $`\tau ^\pm `$ are used, and when the set of states $`|\lambda `$ is a complete one, including all the states in the odd-odd nuclei which can be connected with the ground state $`|0`$ of the even-even nuclei through the transition operators, then
$$S_{}S_+=NZ$$
(36)
However, in the present case the Fermi transition operators are truncated, and the strengths difference has a more complicated form
$$S_{}S_+=0|\widehat{N}\widehat{Z}|0\underset{pn}{}(u_n^2v_p^2)0|\widehat{𝒩}_n\widehat{𝒩}_p|0.$$
(37)
The first term is equivalent to $`NZ`$ due to the constrictions (10). The second term gives rise to the violation of the Ikeda sum rule (36). When the BCS vacuum is employed, these term has no contribution. For this reason the standard pn-QRPA fulfils the Ikeda sum rule. There are some special cases in which the proton and neutron quasiparticle occupations in the ground state are equal, and this term also vanishes, as it was found in the SO(5) model and in the single mode model of the RQRPA . Only when the sum in (32) is restricted to $`J=0`$, forcing $`j_p=j_n`$ and $`\mathrm{\Omega }_p=\mathrm{\Omega }_n`$, it is true that $`0|\widehat{𝒩}_n\widehat{𝒩}_p|0=0`$ for each level. In the present case, when a realistic Hilbert space is employed and the sum over different J’s is included in (32), those occupations are not equal and the Ikeda sum rule is violated. It is worth to mention that, while in general the quasiparticle occupations for protons and neutrons are different, the expectation value of the total number of quasiparticles $`N_{qp}_𝐭0|\widehat{𝒩}_t|0`$ is the same, as can be easily seen making the sum in Eq. (32).
## 5 Results
In the present work we study Fermi beta excitations in <sup>76</sup>Ge. We adopt a $`\delta `$-type residual interaction already used previously , and our Hilbert space has six single particle energy levels, including all the single-particle orbitals from oscillator shells $`3\mathrm{}\omega `$ plus $`1g_{9/2}`$ and $`1g_{7/2}`$ from the $`4\mathrm{}\omega `$ oscillator shell. They were obtained using a Coulomb-corrected Wood-Saxon potential. Their numerical values for <sup>76</sup>Ge are tabulated in Table 1 of ref. . We include $`J^\pi =0^\pm ,..,3^\pm `$ in the sums in Eq.(32).
In order to describe the dependence of the various observables on the proton-neutron residual interaction we use the parameter
$$\mathrm{s}=\frac{v_{\mathrm{𝐩𝐧}}^{pp}}{v^{pair}},$$
which is the ratio between the coupling constant in the proton-neutron $`J=0`$ particle-particle channel and the pairing force constant $`v^{pair}=(v_{\mathrm{𝐩𝐩}}^{pair}+v_{\mathrm{𝐧𝐧}}^{pair})/2`$. We compare results obtained within the usual QRPA with those coming from the RQRPA (no mean field modification), and the SRQRPA.
In Fig. 1 we show the lowest $`J^\pi =0^+`$ excitation energies $`\omega `$ (in MeV) in <sup>76</sup>As, calculated from the <sup>76</sup>Ge ground state, for the QRPA (dotted lines), RQRPA (dashed lines) , and SRQRPA (full lines) approximations, as a function of the residual interaction parameter $`\mathrm{s}`$. It can be clearly seen that the excitation energy goes to zero, and collapses around $`\mathrm{s}=2.5`$ for the QRPA, while in the RQRPA and SRQRPA formalisms the collapse is avoided, although the excitation energies predicted after the collapse in these two approaches differ by a factor two.
In Figure 2 the sum rule $`S_{}S_+`$, normalized to the Ikeda value $`NZ`$, is presented as a function of $`\mathrm{s}`$ with the same convention of Fig. 1. While in the QRPA the Ikeda sum rule is always fulfiled, in the other cases it is violated. Notice that, having the correct number of protons and neutrons in the correlated ground state, the SRQRPA departs less from $`NZ`$ that the RQRPA, but the departure is anyway noticeably. It is important to point out that, in the present calculations, even when $`\mathrm{s}=0`$ the other proton-neutron residual interaction channels, for $`J0`$, are present and have a finite value. For this reason there is a violation of the Ikeda sum rule at $`\mathrm{s}=0`$.
To describe in more detail the above mentioned point we present in Figure 3 the total number of quasiparticles as a function of $`\mathrm{s}`$. When only $`J=0`$ states are included in the sum in Eq. (32) the number of quasiparticles goes to zero when $`\mathrm{s}`$ is zero. In the other three cases the residual interaction in the other channels is present and increase this number.
The fluctuations in the particle number $`\mathrm{\Delta }N\sqrt{0|(\widehat{N}N)^2|0}`$ represent a warning about the region of applicability of the different QRPA models . They are shown in Fig. 4a (for protons) and Fig. 4b (for neutrons) as a function of $`\mathrm{s}`$. It is interesting to note that the pure BCS fluctuations $`\mathrm{\Delta }N=2\sqrt{_t\mathrm{\Omega }_t^2u_t^2v_t^2}`$ diminish when the pn residual interaction increases, due to the sharpening in the particle distribution when the mean field is varied to keep the number of particles in their right value. In the SRQRPA these fluctuations diminish slightly at first, dominated by this mean field effect, to increase in a noticeably way after the QRPA collapse. The QRPA results for the total quasiparticle number and the particle number fluctuations explode for $`s2.5`$ when only the $`J^\pi =0^+`$ contribution is included in Eq. (32). When all the angular momentum contributions are taken into account the $`J^\pi =1^+`$ component collapses around $`s=2.3`$, and these expectation values diverge.
The single particle occupation numbers $`v_p^2,v_n^2`$ are presented in Figure 5 as a function of the single particle energies. While only six levels are in use, the curves represent their smoothed values. It can be seen that both protons and neutrons have their distributions sharpened, i.e. the larger $`\mathrm{s}`$ the more they resemble a step distribution, although this effect is more evident for protons.
The behavior of the pairing energy gap as a function of $`\mathrm{s}`$ is shown in Fig. 6a (for protons) and Fig. 6b (for neutrons). As advanced in the previous section, the pairing gaps, while different for the different single particle levels, all behave the same way. They show a strong reduction when $`\mathrm{s}`$ increases, down to one half of the original value. This is not a curiosity: the pairing energy gap is observable, and any attempt to properly describe the beta decay using the QRPA extensions must keep the descriptive power of the simplest formalism, the QRPA, where the pairing constants are adjusted to reproduce the observed gap.
## 6 Conclusions
A study of the beta decay of <sup>76</sup>Ge using the SRQRPA was presented. The mean field gap and number equations were introduced and solved together with the RQRPA equations, but the selfconsistency was only included to guarantee the good number of particles in the correlated ground state.
The Ikeda sum rule was studied in some detail. It is known to be fulfiled in standard pn-QRPA calculations , violated in RQRPA ones and recovered when self-consistent RQRPA calculations are performed in simple models. In the present work it was shown that the Ikeda sum rule is violated when a realistic Hilbert space is used in spite of using the number selfconsistent QRPA approach.
The most remarkable phenomenon found in the calculation was the strong reduction of the pairing gap. In the presence of the proton-neutron residual interaction the mean field changes needed to have the correct number of particles generate sharper distributions for the occupations of single particle states, visible when plotting $`v^2`$ as a function of the single particle energies for different values of $`\mathrm{s}`$. At the same time, it strongly reduces the energy gap, which can be as small as a half of the observed value.
Having in mind the various merits of the renormalized and self-consistent extensions of the quasiparticle random phase approximation, we conclude that it must be used with extreme caution in the region where the standard QRPA collapse. Not only the particle number fluctuations have a clear increase, but observable quantities depart from their measured values.
## 7 Acknowledgements
This work was supported in part by Conacyt, México.
Figures caption
Figure 1. Lowest $`J^\pi =0^+`$ excitation energies in <sup>76</sup>As, calculated from the <sup>76</sup>Ge ground state, for the QRPA, RQRPA, and SRQRPA approximations, as function of the residual interaction parameter $`\mathrm{s}`$.
Figure 2. Fermi sum rule $`S_{}S_+`$, normalized to the Ikeda value $`NZ`$, as function of $`\mathrm{s}`$ with the same line conventions used in fig. 1.
Figure 3. Total quasiparticle number (equal for protons and neutrons) evolution as function of $`\mathrm{s}`$. The result for the QRPA case when the sums in Eq. (26) are restricted to $`J^\pi =0^+`$ is also shown.
Figure 4. Proton (a) and neutron (b) number fluctuations obtained using the BCS, QRPA with only $`J^\pi =0^+`$ and all the $`J`$’s in Eq. (32), RQRPA and SRQRPA descriptions.
Figure 5. Single particle occupation numbers $`v_p^2`$ (thick lines), and $`v_n^2`$ (thin lines) as a smoothed function of the single particle energies $`e`$, for the values $`s=0.,2.5`$, and $`3.1`$ of the residual interaction parameter.
Figure 6. Evolution of the pairing energy gap as a function of $`\mathrm{s}`$ for protons (a) and for neutrons (b).
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# CP violating phase and quark mixing angles from flavour permutational symmetry breaking
## I Introduction
In this paper we review some previous work \- and give some new results on the functional relations between quark masses and quark mixing angles and CP-violating phases occurring in the phenomenological parametrizations of the quark mixing matrix. These functional relations result from the equivalence under a rephasing of the quark fields of the phenomenological parametrizations and the parametrization derived from the breaking of the flavour permutational symmetry.
In the Standard Electroweak theory of particle interactions, quark flavour mixing is described by means of a unitary mixing matrix $`𝐕`$. The measurables of this matrix, which are invariant under a rephasing of the quark fields, are the moduli of its elements, i. e., the quantities $`|𝐕_{ij}|`$. In the case of three families, unitarity of $`𝐕`$ constrains the number of independent moduli to four. In consequence, phenomenological parametrizations of the quark mixing matrix expressed in terms of four parameters were introduced without taking the possible functional relations between the quark masses and the flavour mixing parameters into account. Kobayashi and Maskawa originally chose as independent parameters three rotation angles and one CP-violating phase. A number of parametrizations of this kind, different from the original Kobayashi-Maskawa form have been proposed , one of the most commonly used is the “standard” parametrization advocated by the Particle Data Group . From a mathematical point of view, two different parametrizations of the $`3\times 3`$ unitary, quark mixing matrix containing four suitably defined, independent parameters are equivalent if the moduli of corresponding entries are equal. In such scheme, both, the masses of the quarks as well as the mixing angles and the CP-violating phase occur in the theory as free, independent parameters.
In contrast, the elements of the quark mixing matrix $`𝐕^{th}`$, derived in two previous papers from a simple ansatz on the breaking of the flavour permutational symmetry -, are explicit functions of the four quark mass ratios $`m_u/m_t`$, $`m_c/m_t`$, $`m_d/m_b`$, $`m_s/m_b`$, and only two, free, linearly independent parameters, namely, the symmetry breaking parameter $`Z^{1/2}`$ and the CP-violating phase $`\mathrm{\Phi }`$. The numerical values of $`Z^{1/2}`$ and $`\mathrm{\Phi }`$ which characterize the preferred symmetry breaking pattern are extracted from a $`\chi ^2`$ fit of the theoretical expressions $`|𝐕_{ij}^{th}|`$ to the experimentally determined values of the moduli of the elements of the mixing matrix $`|𝐕_{ij}^{exp}|`$. It is remarkable that the quality of the best fit of $`𝐕^{th}`$ to the experimental data is as good as the quality of the fit of the phenomenological parametrizations $`𝐕^{PDG}`$ or $`𝐕^{KM}`$ to the same data. More precisely, when the best set of parameters of each parametrization is used, the moduli of corresponding entries in $`𝐕^{th}`$ and $`𝐕^{PDG}`$ or $`𝐕^{KM}`$ are numerically equal and give an equally good representation of the experimentally determined values of the moduli of the mixing matrix $`|𝐕_{ij}^{exp}|`$. Hence, we are justified in writing
$`|𝐕_{ij}^{th}|=|𝐕_{ij}^{PDG}|or|𝐕_{ij}^{th}|=|𝐕_{ij}^{KM}|,`$ (1)
even though $`𝐕^{th}`$ has only two free, real linearly independent adjustable parameters, while the number of adjustable parameters in $`𝐕^{PDG}`$ or $`𝐕^{KM}`$ is four.
In what follows, it will be shown that by means of suitable rephasing of the quark fields $`𝐕^{th}`$ may be changed into new, phase transformed forms $`\stackrel{~}{𝐕}^{th}`$ or $`\widehat{𝐕}^{th}`$ such that all their matrix elements are numerically equal to the corresponding entries in $`𝐕^{PDG}`$ or $`𝐕^{KM}`$, both in modulus and phase. Once this equality is established, we derive exact explicit analitycal expressions for the mixing angles and the CP-violating phase of the two phenomenological parametrizations $`𝐕^{PDG}`$ and $`𝐕^{KM}`$, in terms of the quark mass ratios $`m_u/m_t`$, $`m_c/m_t`$, $`m_d/m_b`$, $`m_s/m_b`$, the flavour symmetry breaking parameter $`Z^{1/2}`$ and the CP-violating phase $`\mathrm{\Phi }`$.
The plan of this paper is as follows: In Sec. II, we introduce some basic concepts and fix the notation by way of a very brief sketch of the group theoretical derivation of mass matrices with a modified Fritzsch texture. Sect. III is devoted to the derivation of exact, explicit expressions for the elements of the mixing matrix $`V_{ij}^{th}`$ in terms of the quark mass ratios and the parameters $`Z^{1/2}`$ and $`\mathrm{\Phi }`$ characterizing the symmetry breaking pattern. In Sec. IV, the phase equivalence of $`𝐕^{th}`$ and the phenomenological parametrrization $`𝐕^{PDG}`$ and $`𝐕^{KM}`$ is established. The equations that define the rephasing transformation connecting $`𝐕^{th}`$ and $`𝐕^{PDG}`$ are solved in Sec. V. Sections VI and VII are devoted to the derivation of explicit expressions for the mixing parameters $`\mathrm{sin}\theta _{12}`$, $`\mathrm{sin}\theta _{23}`$, $`\mathrm{sin}\theta _{13}`$ and the CP violating phase $`\delta _{13}`$ of $`𝐕^{PDG}`$ as funtions of the quark mass ratios and the parameters $`Z^{1/2}`$ and $`\mathrm{\Phi }`$. In Sec. VIII, the equations that define the rephasing transformation relating $`𝐕^{th}`$ and $`𝐕^{KM}`$ are solved. Explicit expressions for the mixing parameters $`\mathrm{sin}\theta _1`$, $`\mathrm{sin}\theta _2`$, $`\mathrm{sin}\theta _3`$ and the CP violating phase $`\delta _{KM}`$ of $`𝐕^{KM}`$ as funtions of the quark mass ratios and the symmetry breaking parameters are obtained in Sec. IX and X. Our paper ends in Sec. XI with a summary of results and some conclusions.
## II Mass matrices from the breaking of $`S_L(3)S_R(3)`$
In the Standard Model, analogous fermions in different generations, say $`u,c`$ and $`t`$ or $`d,s`$ and $`b`$, have completely identical couplings to all gauge bosons of the strong, weak and electromagnetic interactions. Prior to the introduction of the Higgs boson and mass terms, the Lagrangian is chiral and invariant with respect to any permutation of the left and right quark fields. The introduction of a Higgs boson and the Yukawa couplings give mass to the quarks and leptons when the gauge symmetry is spontaneously broken. The quark mass term in the Lagrangian, obtained by taking the vacuum expectation value of the Higgs field in the quark Higgs coupling, gives rise to quark mass matrices $`𝐌_𝐝`$ and $`𝐌_𝐮`$,
$$_Y=\overline{𝐪}_{d,L}𝐌_d𝐪_{d,R}+\overline{𝐪}_{u,L}𝐌_u𝐪_{u,R}+h.c.$$
(2)
In this expression, $`𝐪_{d,L,R}(x)`$ and $`𝐪_{u,L,R}(x)`$ denote the left and right quark $`d`$\- and $`u`$-fields in the current or weak basis, $`𝐪_q(x)`$ is a column matrix, its components $`𝐪_{q,k}(x)`$ are the quark Dirac fields, $`k`$ is the flavour index. In this basis, the charged hadronic currents,
$$J_\mu \overline{q}_{u,L}\gamma _\mu q_{d,L},$$
(3)
are not changed if both, the d-type and the u-type fields are transformed with the same unitary matrix.
### A Modified Fritzsch texture
A number of authors -, \[-\] have pointed out that realistic quark mass matrices result from the flavour permutational symmetry $`S_L(3)S_R(3)`$ and its spontaneous or explicit breaking. The group $`S(3)`$ treats three objects symmetrically, while the hierarchical nature of the mass matrices is a consequence of the representation structure $`\mathrm{𝟏}\mathrm{𝟐}`$ of $`S(3)`$, which treats the generations differently. Under exact $`S_L(3)S_R(3)`$ symmetry, the mass spectrum for either up or down quark sectors consists of one massive particle in a singlet irreducible representation and a pair of massless particles in a doublet irreducible representation, the corresponding quark mass matrix with the exact $`S_L(3)S_R(3)`$ symmetry will be denoted by $`𝐌_{3q}`$. In order to generate masses for the first and second families, we add the terms $`𝐌_{2q}`$ and $`𝐌_{1q}`$ to $`𝐌_{3q}`$. The term $`𝐌_{2q}`$ breaks the permutational symmetry $`S_L(3)S_R(3)`$ down to $`S_L(2)S_R(2)`$ and mixes the singlet and doublet representation of $`S(3)`$. $`𝐌_{1q}`$ transforms as the mixed symmetry term in the doublet complex tensorial representation of $`S_{diag}(3)S_L(3)S_R(3)`$. Putting the first family in a complex representation will allow us to have a CP violating phase in the mixing matrix. Then, in a symmetry adapted basis , $`𝐌_q`$ takes the form
$$\begin{array}{ccc}\hfill M_q& =& m_{3q}\left[\left(\begin{array}{ccc}0& A_qe^{i\varphi _q}& 0\\ A_qe^{i\varphi _q}& 0& 0\\ 0& 0& 0\end{array}\right)+\left(\begin{array}{ccc}0& 0& 0\\ 0& \mathrm{}_q+\delta _q& B_q\\ 0& B_q& \mathrm{}_q\delta _q\end{array}\right)\right]\hfill \\ & +& m_{3q}\left(\begin{array}{ccc}0& 0& 0\\ 0& 0& 0\\ 0& 0& 1\mathrm{}_q\end{array}\right)=m_{3q}\left(\begin{array}{ccc}0& A_qe^{i\varphi _q}& 0\\ A_qe^{i\varphi _q}& \mathrm{}_q+\delta _q& B_q\\ 0& B_q& 1\delta _q\end{array}\right).\hfill \end{array}$$
(4)
From the strong hierarchy in the masses of the quark families, $`m_{3q}>>m_{2q}>m_{1q}`$, we expect $`1\delta _q`$ to be very close to unity. The entries in the mass matrix may be readily expressed in terms of the mass eigenvalues $`(m_{1q},m_{2q},m_{3q})`$ and the small parameter $`\delta _q`$. Computing the invariants of $`M_q`$, $`trM_q`$, $`trM_{q}^{}{}_{}{}^{2}`$ and $`detM_q`$, we get
$`A_q^2=\stackrel{~}{m}_{1q}\stackrel{~}{m}_{2q}(1\delta _q)^1,\mathrm{}_q=\stackrel{~}{m}_{2q}\stackrel{~}{m}_{1q},`$ (5)
(6)
$`B_q^2=\delta _q((1\stackrel{~}{m}_{1q}+\stackrel{~}{m}_{2q}\delta _q)\stackrel{~}{m}_{1q}\stackrel{~}{m}_{2q}(1\delta _q)^1),`$ (7)
where $`\stackrel{~}{m}_{1q}=m_{1q}/m_{3q}`$ and $`\stackrel{~}{m}_{2q}=m_{2q}/m_{3q}`$.
If each possible symmetry breaking pattern is now characterized by the ratio
$$Z_{q}^{}{}_{}{}^{1/2}=B_q/(\mathrm{}_q+\delta _q),$$
(8)
the small parameter $`\delta _q`$ is obtained as the solution of the cubic equation
$$\delta _q\left[(1+\stackrel{~}{m}_{2q}\stackrel{~}{m}_{1q}\delta _q)(1\delta _q)\stackrel{~}{m}_{1q}\stackrel{~}{m}_{2q}\right]Z_q(\stackrel{~}{m}_{2q}+\stackrel{~}{m}_{1q}+\delta _q)^2=0,$$
(9)
which vanishes when $`Z_q`$ vanishes. An exact explicit expression for $`\delta _q`$ as function of the quark mass ratios and $`Z_q`$ is given in . An approximate solution to Eq. (9) for $`\delta _q(Z_q)`$, valid for small values of $`Z_q`$ ($`Z_q10`$), is
$`\delta _q\left(Z_q\right){\displaystyle \frac{Z_q\left(\stackrel{~}{m}_{2q}\stackrel{~}{m}_{1q}\right)^2}{\left(1\stackrel{~}{m}_{1q}\right)\left(1+\stackrel{~}{m}_{2q}\right)+2Z_q\left(\stackrel{~}{m}_{2q}\stackrel{~}{m}_{1q}\right)(1+\frac{1}{2}(\stackrel{~}{m}_{2q}\stackrel{~}{m}_{1q}))}}.`$ (10)
### B Symmetry breaking pattern
In the symmetry adapted basis, the matrix $`𝐌_{\mathrm{𝟐}𝐪}`$, written in term of $`Z_q^{1/2}`$, takes the form
$$𝐌_{2q}=m_{3q}\left(\stackrel{~}{m}_{2q}+\stackrel{~}{m}_{1q}+\delta _q\right)\left(\begin{array}{ccc}0& 0& 0\\ 0& 1& Z_q^{1/2}\\ 0& Z_q^{1/2}& 1\end{array}\right),$$
(11)
when $`Z_q^{1/2}`$ vanishes, $`𝐌_{2q}`$ is diagonal and there is no mixing of singlet and doublet representations of $`S(3)`$. Therefore, in the symmetry adapted basis, the parameter $`Z_q^{1/2}`$ is a measure of the amount of mixing of singlet and doublet irreducible representations of $`S_{diag}(3)S_L(3)S_R(3)`$.
We may easily give a meaning to $`Z_q^{1/2}`$ in terms of permutations. From Eqs. (2) and (11), we notice that the symmetry breaking term in the Yukawa Lagrangian, $`\overline{𝐪}_L𝐌_{2q}𝐪_R`$ is a functional of only two fields: $`\frac{1}{\sqrt{3}}\left(q_2(X)+\sqrt{2}q_3(X)\right)`$ and $`\frac{1}{\sqrt{3}}\left(\sqrt{2}q_2(X)+q_3(X)\right)`$. Under the permutation of these fields, $`\overline{𝐪}_L𝐌_{2q}𝐪_R`$ splits into the sum of an antisymmetric term $`\overline{𝐪}_L𝐌_{2q}^A𝐪_R`$ which changes sign, and a symmetric term $`\overline{𝐪}_L𝐌_{2q}^S𝐪_R`$, which remains invariant,
$$𝐌_{2q}=\frac{2}{9}m_{3q}\left\{a\left(\begin{array}{ccc}0& 0& 0\\ 0& 1& \sqrt{8}\\ 0& \sqrt{8}& 1\end{array}\right)+2b\left(\begin{array}{ccc}0& 0& 0\\ 0& 1& \frac{1}{\sqrt{8}}\\ 0& \frac{1}{\sqrt{8}}& 1\end{array}\right)\right\},$$
(12)
where $`a=(\delta _q\mathrm{}_q)(\sqrt{2}Z_q^{1/2}\frac{1}{2})`$ and $`b=(\delta _q\mathrm{}_q)(\frac{\sqrt{2}}{2}Z_q^{1/2}+2)`$. It is evident that there is a corresponding decomposition of the mixing parameter $`Z_q^{1/2}`$,
$$Z_{q}^{}{}_{}{}^{1/2}=N_{Aq}Z_A^{1/2}+N_{Sq}Z_S^{1/2}$$
(13)
with
$$1=N_{Aq}+N_{Sq},$$
(14)
where $`Z_A^{1/2}=\sqrt{8}`$ is the mixing parameter of the matrix $`𝐌_{2q}^A`$, and $`Z_S^{1/2}=\frac{1}{\sqrt{8}}`$ is the mixing parameter of $`𝐌_{2q}^S`$. In this way, a unique linear combination of $`Z_A^{1/2}`$ and $`Z_S^{1/2}`$ is associated to the simmetry breaking pattern characterized by $`Z_q^{1/2}`$. Thus, the different symmetry breaking patterns defined by $`𝐌_{2q}`$ for different values of the mixing parameter $`Z_q^{1/2}`$ are labeled in terms of the irreducible representations of the group $`\stackrel{~}{S}(2)`$ of permutations of the two fields in $`\overline{𝐪}_L𝐌_{2q}𝐪_R`$. The pair of numbers $`(N_A,N_S)`$ enters as a convenient mathematical label of the symmetry breaking pattern without introducing any assumption about the actual pattern of $`S_L(3)S_R(3)`$ symmetry breaking realized in nature.
### C The Jarlskog invariant
The Jarlskog invariant, $`J`$, may be computed directly from the commutator of the mass matrices
$$J=\frac{det\{i[𝐌_u,𝐌_d]\}}{2F}$$
(15)
where
$$F=(1+\stackrel{~}{m}_c)(1\stackrel{~}{m}_u)(\stackrel{~}{m}_c+\stackrel{~}{m}_u)(1+\stackrel{~}{m}_s)(1\stackrel{~}{m}_d)(\stackrel{~}{m}_s+\stackrel{~}{m}_d).$$
(16)
Substitution of the expression (4) for $`𝐌_u`$ and $`𝐌_d`$, in Eq. (15), with $`Z_u^{1/2}=Z_d^{1/2}=Z^{1/2}`$ gives
$`J`$ $`={\displaystyle \frac{Z\sqrt{\frac{\stackrel{~}{m}_u/\stackrel{~}{m}_c}{1\delta _u}}\sqrt{\frac{\stackrel{~}{m}_d/\stackrel{~}{m}_s}{1\delta _d}}sin\mathrm{\Phi }}{(1+\stackrel{~}{m}_c)(1\stackrel{~}{m}_u)(1+\stackrel{~}{m}_u/\stackrel{~}{m}_c)(1+\stackrel{~}{m}_s)(1\stackrel{~}{m}_d)(1+\stackrel{~}{m}_d/\stackrel{~}{m}_s)}}`$ (17)
$`\times `$ $`\{[(\mathrm{}_u+\delta _u)(1\delta _d)(\mathrm{}_d+\delta _d)(1\delta _u)]^2\left({\displaystyle \frac{\stackrel{~}{m}_u\stackrel{~}{m}_c}{1\delta _u}}\right)(\mathrm{}_d+\delta _d)^2`$ (18)
$``$ $`\left({\displaystyle \frac{\stackrel{~}{m}_d\stackrel{~}{m}_s}{1\delta _d}}\right)(\mathrm{}_u+\delta _u)^2+2\sqrt{{\displaystyle \frac{\stackrel{~}{m}_u\stackrel{~}{m}_c}{1\delta _u}}}\sqrt{{\displaystyle \frac{\stackrel{~}{m}_d\stackrel{~}{m}_s}{1\delta _d}}}(\mathrm{}_u+\delta _u)(\mathrm{}_d+\delta _d)cos\mathrm{\Phi }\}.`$ (19)
where $`\mathrm{}_q`$ and $`\delta _q`$ are defined in Eqs. (5) and (9). In this, way, an exact closed expression for $`J`$ in terms of the quark mass ratios, the CP violating phase $`\mathrm{\Phi }`$, and the parameter $`Z`$ that characterizes the symmetry breaking pattern is derived.
## III The Mixing Matrix
The Hermitian mass matrix $`𝐌_q`$ may be written in terms of a real symmetric matrix $`\overline{𝐌}_q`$ and a diagonal matrix of phases $`𝐏_q`$ as follows
$$𝐌_q=𝐏_q\overline{𝐌}_q𝐏_{q}^{}{}_{}{}^{},$$
(20)
The real symmetric matrix $`\overline{𝐌}_q`$ may be brought to a diagonal form by means of an orthogonal transformation
$$\overline{𝐌}_q=𝐎_q𝐌_{q,diag}𝐎_q^T,$$
(21)
where
$$𝐌_{q,diag}=m_{3q}diag[\stackrel{~}{m}_{1q},\stackrel{~}{m}_{2q},1],$$
(22)
with subscripts $`1,2,3`$ refering to $`u,c,t`$ in the u-type sector and $`d,s,b`$ in the d-type sector.
After diagonalization of the mass matrices $`𝐌_q`$, one obtains the mixing matrix $`𝐕^{th}`$ as
$$𝐕^{th}=𝐎_{u}^{}{}_{}{}^{T}𝐏^{ud}𝐎_d,$$
(23)
where $`𝐏^{ud}`$ is the diagonal matrix of relative phases
$$𝐏^{ud}=diag[1,e^{i\mathrm{\Phi }},e^{i\mathrm{\Phi }}],$$
(24)
and
$$\mathrm{\Phi }=(\varphi _u\varphi _d).$$
(25)
The orthogonal matrix $`𝐎_q`$is given by
$$\begin{array}{ccc}\hfill 𝐎_q=\left(\begin{array}{ccc}(\stackrel{~}{m}_{2q}f_1/D_1)^{1/2}& (\stackrel{~}{m}_{1q}f_2/D_2)^{1/2}& (\stackrel{~}{m}_{1q}\stackrel{~}{m}_{2q}f_3/D_3)^{1/2}\\ ((1\delta _q)\stackrel{~}{m}_{1q}f_1/D_1)^{1/2}& ((1\delta _q)\stackrel{~}{m}_{2q}f_2/D_2)^{1/2}& ((1\delta _q)f_3/D_3)^{1/2}\\ (\stackrel{~}{m}_{1q}f_2f_3/D_1)^{1/2}& (\stackrel{~}{m}_{2q}f_1f_3/D_2)^{1/2}& (f_1f_2/D_3)^{1/2}\end{array}\right),& & \end{array}$$
(26)
where
$`\mathrm{f}_1=1\stackrel{~}{m}_{1q}\delta _q,\mathrm{f}_2=1+\stackrel{~}{m}_{2q}\delta _q,\mathrm{f}_3=\delta _q,`$ (27)
$`D_1=(1\delta _q)\left(1\stackrel{~}{m}_{1q}\right)\left(\stackrel{~}{m}_{2q}+\stackrel{~}{m}_{1q}\right),`$ (28)
$`D_2=(1\delta _q)\left(1+\stackrel{~}{m}_{2q}\right)\left(\stackrel{~}{m}_{2q}+\stackrel{~}{m}_{1q}\right),`$ (29)
$`D_3=(1\delta _q)\left(1+\stackrel{~}{m}_{2q}\right)\left(1\stackrel{~}{m}_{1q}\right).`$ (30)
In these expressions, $`\delta _u`$ and $`\delta _d`$ are, in principle, functions of the quark mass ratios and the parameters $`Z_u^{1/2}`$ and $`Z_d^{1/2}`$ respectively. However, in we found that keeping $`Z_u^{1/2}`$ and $`Z_d^{1/2}`$ as free, independent parameters gives rise to a continuous ambiguity in the fitting of $`|V_{ij}^{th}|`$ to the experimental data. To avoid this ambiguity we further assumed that the up and down mass matrices are generated following the same symmetry breaking pattern, that is,
$$Z_u^{1/2}=Z_d^{1/2}=Z^{1/2}.$$
(31)
Then, from Eqs. (23) - (31) all matrix elements in $`𝐕^{th}`$ may be written in terms of four quark mass ratios and only two free, real parameters: the parameter $`Z^{1/2}`$ which characterizes the symmetry breaking pattern in the u- and d-sectors and the CP violating phase $`\mathrm{\Phi }`$. The computation of $`V_{ij}^{th}`$ is quite straightforward. Here, we will give, in explicit form, only those elements of $`𝐕^{th}`$ which will be of use later. From Eqs. (23)-(31) we obtain,
$`V_{us}^{th}`$ $`=`$ $`\left({\displaystyle \frac{\stackrel{~}{m}_c\left(1\stackrel{~}{m}_u\delta _u\right)\stackrel{~}{m}_d\left(1+\stackrel{~}{m}_s\delta _d\right)}{\left(1\delta _u\right)\left(1\stackrel{~}{m}_u\right)\left(\stackrel{~}{m}_c+\stackrel{~}{m}_u\right)\left(1\delta _d\right)\left(1+\stackrel{~}{m}_s\right)\left(\stackrel{~}{m}_s+\stackrel{~}{m}_d\right)}}\right)^{1/2}`$ (32)
$`+`$ $`\left({\displaystyle \frac{\stackrel{~}{m}_u\stackrel{~}{m}_s}{\left(1\stackrel{~}{m}_u\right)\left(\stackrel{~}{m}_c+\stackrel{~}{m}_u\right)\left(\stackrel{~}{m}_d+\stackrel{~}{m}_s\right)}}\right)^{1/2}\{\left({\displaystyle \frac{\left(1\stackrel{~}{m}_u\delta _u\right)\left(1+\stackrel{~}{m}_s\delta _d\right)}{\left(1+\stackrel{~}{m}_s\right)}}\right)^{1/2}`$ (33)
$`+`$ $`\left({\displaystyle \frac{\left(1+\stackrel{~}{m}_c\delta _u\right)\delta _u\left(1\stackrel{~}{m}_d\delta _d\right)\delta _d}{(1\delta _u)(1\delta _d)(1+\stackrel{~}{m}_s)}}\right)^{1/2}\}e^{i\mathrm{\Phi }}`$ (34)
$`V_{ub}^{th}`$ $`=`$ $`\left({\displaystyle \frac{\stackrel{~}{m}_c(1\stackrel{~}{m}_u\delta _u)}{(1\delta _u)(1\stackrel{~}{m}_u)(\stackrel{~}{m}_c+\stackrel{~}{m}_u)}}{\displaystyle \frac{\stackrel{~}{m}_d\stackrel{~}{m}_s\delta _d}{(1\delta _d)(1+\stackrel{~}{m}_s)(1\stackrel{~}{m}_d)}}\right)^{1/2}`$ (35)
$`+`$ $`\{\left({\displaystyle \frac{\stackrel{~}{m}_u(1+\stackrel{~}{m}_c\delta _u)\delta _u(1\stackrel{~}{m}_d\delta _d)(1+\stackrel{~}{m}_s\delta _d)}{(1\delta _u)(1\stackrel{~}{m}_u)(\stackrel{~}{m}_c+\stackrel{~}{m}_u)(1\delta _d)(1+\stackrel{~}{m}_s)(1\stackrel{~}{m}_d)}}\right)^{1/2}`$ (36)
$`+`$ $`\left({\displaystyle \frac{\stackrel{~}{m}_u(1\stackrel{~}{m}_u\delta _u)\delta _d}{(1\stackrel{~}{m}_u)(\stackrel{~}{m}_c+\stackrel{~}{m}_u)(1+\stackrel{~}{m}_s)(1\stackrel{~}{m}_d)}}\right)^{1/2}\}e^{i\mathrm{\Phi }}`$ (37)
$`V_{cs}^{th}`$ $`=`$ $`\left({\displaystyle \frac{\stackrel{~}{m}_u\left(1+\stackrel{~}{m}_c\delta _u\right)\stackrel{~}{m}_d\left(1+\stackrel{~}{m}_s\delta _d\right)}{\left(1\delta _u\right)\left(1+\stackrel{~}{m}_c\right)\left(\stackrel{~}{m}_c+\stackrel{~}{m}_u\right)\left(1\delta _d\right)\left(1+\stackrel{~}{m}_s\right)\left(\stackrel{~}{m}_s+\stackrel{~}{m}_d\right)}}\right)^{1/2}`$ (38)
$`+`$ $`\{\left({\displaystyle \frac{\stackrel{~}{m}_c\delta _u\left(1\stackrel{~}{m}_u\delta _u\right)\stackrel{~}{m}_s\delta _d\left(1\stackrel{~}{m}_d\delta _d\right)}{\left(1\delta _u\right)\left(1+\stackrel{~}{m}_c\right)\left(\stackrel{~}{m}_c+\stackrel{~}{m}_u\right)\left(1\delta _d\right)\left(1+\stackrel{~}{m}_s\right)\left(\stackrel{~}{m}_s+\stackrel{~}{m}_d\right)}}\right)^{1/2}`$ (39)
$`+`$ $`\left({\displaystyle \frac{\stackrel{~}{m}_c\left(1+\stackrel{~}{m}_c\delta _u\right)\stackrel{~}{m}_s\left(1+\stackrel{~}{m}_s\delta _d\right)}{\left(1+\stackrel{~}{m}_c\right)\left(1\stackrel{~}{m}_u\right)\left(1+\stackrel{~}{m}_s\right)\left(1\stackrel{~}{m}_d\right)}}\right)^{1/2}\}e^{i\mathrm{\Phi }}`$ (40)
and
$`V_{cb}^{th}`$ $`=`$ $`\left({\displaystyle \frac{\stackrel{~}{m}_u(1+\stackrel{~}{m}_c\delta _u)}{(1\delta _u)(1+\stackrel{~}{m}_c)(\stackrel{~}{m}_c+\stackrel{~}{m}_u)}}{\displaystyle \frac{\stackrel{~}{m}_d\stackrel{~}{m}_s\delta _d}{(1\delta _d)(1+\stackrel{~}{m}_s)(1\stackrel{~}{m}_d)}}\right)^{1/2}`$ (41)
$`+`$ $`\{\left({\displaystyle \frac{\stackrel{~}{m}_c(1\stackrel{~}{m}_u\delta _u)\delta _u(1\stackrel{~}{m}_d\delta _d)(1+\stackrel{~}{m}_s\delta _d)}{(1\delta _u)(1+\stackrel{~}{m}_c)(\stackrel{~}{m}_c+\stackrel{~}{m}_u)(1\delta _d)(1+\stackrel{~}{m}_s)(1\stackrel{~}{m}_d)}}\right)^{1/2}`$ (42)
$`+`$ $`\left({\displaystyle \frac{\stackrel{~}{m}_c(1+\stackrel{~}{m}_c\delta _u)}{(\stackrel{~}{m}_c+\stackrel{~}{m}_u)(1+\stackrel{~}{m}_c)}}{\displaystyle \frac{\delta _d}{(1+\stackrel{~}{m}_s)(1\stackrel{~}{m}_d)}}\right)^{1/2}\}e^{i\mathrm{\Phi }}.`$ (43)
### A The “best” symmetry breaking pattern
In order to find the actual pattern of $`S_L(3)S_R(3)`$ symmetry breaking realized in nature,we made a $`\chi ^2`$ fit of the exact expressions for the moduli of the entries in the mixing matrix, $`|V_{ij}^{th}|`$, the Jarlskog invariant, $`J^{th}`$, and the three inner angles of the unitarity triangle, $`\alpha ^{th}`$, $`\beta ^{th}`$ and $`\gamma ^{th}`$, to the experimentally determined values of $`|V_{ij}^{exp}|`$, $`J^{exp}`$, $`\alpha ^{exp}`$, $`\beta ^{exp}`$ and $`\gamma ^{exp}`$. A detailed account of the fitting procedure is given in . Here, we will give only a brief relation of the main points in the fitting procedure.
For the purpose of calculating quark mass ratios and computing the mixing matrix, it is convenient to give all quark masses as running masses at some common energy scale , . In the present calculation, following Peccei , Fritzsch and the Ba-Bar book , we used the values of the running quark masses evaluated at $`\mu =m_t`$.
$`m_u`$ $`=`$ $`3.25\pm 0.9MeVm_c=760\pm 29.5MeVm_t=171.0\pm 12GeV`$ (44)
$`m_d`$ $`=`$ $`4.4\pm 0.64MeVm_s=100\pm 6MeVm_b=2.92\pm 0.11GeV`$ (45)
These values, with the exception of $`m_s`$, $`m_c`$ and $`m_b`$, were taken from the work of Fusaoka and Koide see also Fritzsch and Leutwyler . The values of $`m_c(m_t)`$ and $`m_b(m_t)`$ were obtained by rescaling to $`\mu =m_t`$ the recent calculations of $`m_c(m_c)`$ and $`m_b(m_b)`$ by Pineda and Ynduráin and Ynduráin . The value of $`m_s`$ agrees with the latest determination made by the ALEPH collaboration from a study of $`\tau `$ decays involving kaons .
We kept the mass ratios $`\stackrel{~}{m}_c=m_c/m_t`$ , $`\stackrel{~}{m}_s=m_s/m_b`$ and $`\stackrel{~}{m}_d=m_d/m_b`$ fixed at their central values
$`\stackrel{~}{m}_c=0.0044,\stackrel{~}{m}_s=0.034and\stackrel{~}{m}_d=0.0015,`$ (46)
but we took the value
$`\stackrel{~}{m}_u=0.000036,`$ (47)
which is close to its upper bound. We found the following best values for $`\mathrm{\Phi }`$ and $`Z^{1/2}`$,
$`\mathrm{\Phi }^{}=90^{},Z^{1/2}={\displaystyle \frac{1}{2}}\left[Z_S^{1/2}Z_A^{1/2}\right]=\sqrt{{\displaystyle \frac{81}{32}}}.`$ (48)
corresponding to a value of $`\chi ^20.32`$. The values of the parameters $`\delta _u(Z)`$ and $`\delta _d(Z)`$ obtained from (46), (47) and (48) are
$`\delta _u(Z^{1/2})=0.000048,\delta _d(Z^{1/2})=0.00228.`$ (49)
Before proceeding to give the numerical results for the mixing matrix $`𝐕^{th}`$, it will be convenient to stress the following points:
1. The masses of the lighter quarks are the less well determined, while the moduli of the entries in $`|V_{ij}^{exp}|`$ with the largest error bars, namely $`|V_{ub}|`$ and $`|V_{td}|`$, are the most sensitive to changes in the ratios $`m_u/m_c`$ and $`m_d/m_s`$ respectively. Hence, the quality of the fit of $`|V_{ij}^{th}|`$ to $`|V_{ij}^{exp}|`$ is good $`(\chi ^20.5)`$ even if relatively large changes in the masses of the lighter quarks are made. The sensitivity of $`|V_{ub}|`$ and $`|V_{td}|`$ to changes in $`m_u/m_c`$ and $`m_d/m_s`$ respectively, is reflected in the shape of the unitarity triangle which changes appreciably when the masses of the ligther quarks change within their uncertainty ranges. The best simultaneous $`\chi ^2`$ fit of $`|V_{ij}^{th}|`$, $`J^{th}`$, and $`\alpha ^{th}`$, $`\beta ^{th}`$ and $`\gamma ^{th}`$, to the experimentally determined quantities was obtained when the ratio $`\stackrel{~}{m}_u=m_u/m_t`$ is taken close to its upper bound, as given in (47). Furthermore, the chosen high value of $`\stackrel{~}{m}_u`$ gives for the ratio $`|V_{ub}/V_{cb}|`$ the value
$`{\displaystyle \frac{|V_{ub}|}{|V_{cb}|}}\sqrt{{\displaystyle \frac{m_u}{m_c}}}=0.09\pm 0.025`$ (50)
in very good agreement with its latest world average \[, , \].
2. As the energy scale changes, say from $`\mu =m_t`$ to $`\mu =1GeV`$, the running quark masses change appreciably, but since the masses of light and heavy quarks increase almost in the same proportion, the resulting dependence of the quark mass ratios on the energy scale is very weak. When the energy scale changes from $`\mu =m_t`$ to $`\mu =1GeV`$, $`\stackrel{~}{m}_u`$ and $`\stackrel{~}{m}_d`$ decrease by about $`25\%`$ and $`\stackrel{~}{m}_c`$ and $`\stackrel{~}{m}_s`$ also decrease but by less than $`16\%`$.
3. In view of the previous considerations, a reasonable range of values for the running quark mass ratios, evaluated at $`\mu =m_t=171GeV`$, would be as follows
$`0.000022`$ $`\stackrel{~}{m}_u`$ $`0.000037`$ (51)
$`0.0043`$ $`\stackrel{~}{m}_c`$ $`0.0046`$ (52)
$`0.0013`$ $`\stackrel{~}{m}_d`$ $`0.0017`$ (53)
$`0.032`$ $`\stackrel{~}{m}_s`$ $`0.036`$ (54)
The results of the $`\chi ^2`$ fit of the theoretical expressions for $`|V_{ij}^{th}|`$, $`J^{th}`$, $`\alpha ^{th}`$, $`\beta ^{th}`$ and $`\gamma ^{th}`$ to the experimentally determined quantities is as follows:
The quark mixing matrix computed from the theoretical expresion $`𝐕^{th}`$ with the numerical values of quark mass ratios given in (46)and (47) and the corresponding best values of the symmetry breaking parameter, $`Z^{1/2}=\sqrt{81/32}`$, and the CP-violating phase, $`\mathrm{\Phi }^{}=90^{}`$, is
$$V^{th}=\left(\begin{array}{ccc}0.9749e^{i1^{}}& 0.2225e^{i157^{}}& 0.0036e^{i85^{}}\\ 0.2223e^{i113^{}}& 0.9742e^{i89^{}}& 0.040e^{i90^{}}\\ 0.0086e^{i270^{}}& 0.0392e^{i270^{}}& 0.9992e^{i90^{}}\end{array}\right)$$
(55)
In order to have an estimation of the sensivity of our numerical results to the uncertainty in the values of the quark mass ratios, we computed the range of values of the matrix of moduli $`|V_{ij}^{th}|`$, corresponding to the range of values of the mass ratios given in (51), but keeping $`\mathrm{\Phi }`$ and $`Z^{1/2}`$ fixed at the values $`\mathrm{\Phi }^{}=90^{}`$ and $`Z^{1/2}=\sqrt{81/32}`$. The result is
$$|𝐕^{th}|=\left(\begin{array}{ccc}0.97350.9771& 0.21510.2263& 0.00280.0040\\ 0.21510.2263& 0.97260.9764& 0.0370.043\\ 0.00780.0093& 0.0360.042& 0.99910.9993\end{array}\right),$$
(56)
which is to be compared with the experimentally determined values of the matrix of moduli ,
$$|𝐕^{exp}|=\left(\begin{array}{ccc}0.97420.9757& 0.2190.226& 0.0020.005\\ 0.2190.225& 0.97340.9749& 0.0370.043\\ 0.0040.014& 0.0350.043& 0.99900.9993\end{array}\right).$$
(57)
As is apparent from (55), (56) and (57), the agreement between computed and experimental values of all entries in the mixing matrix is very good. The estimated range of variation in the computed values of the moduli of the four entries in the upper left corner of the matrix $`|𝐕^{th}|`$ is larger than the error band in the corresponding entries of the matrix of the experimentally determined values of the moduli $`|𝐕^{exp}|`$. The estimated range of variation in the computed values of the entries in the third column and the third row of $`|V_{ij}^{th}|`$ is comparable with the error band of the corresponding entries in the matrix of experimentally determined values of the moduli, with the exception of the elements $`|V_{ub}^{th}|`$ and $`|V_{td}^{th}|`$ in which case the estimated range of variation due to the uncertainty in the values of the quark mass ratios is significantly smaller than the error band in the experimentally determined value of $`|V_{ub}^{exp}|`$ and $`|V_{td}^{exp}|`$.
The value obtained for the Jarlskog invariant is
$$J^{th}=3.0\times 10^5$$
(58)
in good agreement with the value $`|J^{exp}|=(3.0\pm 1.3)\times 10^5sin\delta `$ obtained from current data on CP violation in the $`K^{}\overline{K}^{}`$ mixing system and the $`B^{}\overline{B}^{}`$ mixing system .
For the inner angles of the unitarity triangle, we found the following central values:
$`\alpha =83.6^{}\beta =23.2^{}\gamma =73.2^{}.`$ (59)
An estimation of the range of values of the three inner angles of the unitarity triangle, compatible with the experimental information on the absolute values of the matrix elements $`𝐕^{exp}`$, is given by Mele and Ali . According to this authors, $`75^{}\alpha 121^{}`$, $`16^{}\beta 34^{}`$, and $`38^{}\gamma 81^{}`$. We see that the central value of $`\beta `$ obtained in this work is close to the lower limit according to Mele , while $`\gamma `$ is close to the upper limit given by Mele and $`\alpha `$ is in the allowed range given by these authors.
## IV Phase equivalence of $`𝐕^{th}`$ and the phenomenological parametrizations $`𝐕^{PDG}`$ and $`𝐕^{KM}`$
In the mass basis, the quark charged currents take the form
$$J_c^\mu =\frac{g}{\sqrt{2}}\overline{q}_{Li}^u\gamma ^\mu V_{ij}q_{Lj}^d.$$
(60)
A redefinition of the phases of the quark fields which leaves $`J_c^\mu `$ invariant, will change the arguments of $`V_{ij}`$ but leave the moduli $`|V_{ij}|`$ invariant,
$`V_{ij}\stackrel{~}{V}_{ij}=e^{i\chi _i^u}V_{ij}e^{i\chi _j^d}.`$ (61)
Hence, the invariant measurables of the quark mixing matrix are the moduli of its elements, i.e., the quantities $`|V_{ij}|`$, and the Jarlskog invariant J. But even J up to a sign is a function of the moduli . In the case of three families, unitarity of $`𝐕`$ constrains the number of independent moduli to four .
In consequence, in the current literature, $`𝐕`$, is usually parametrized in terms of four independent parameters. Kobayashi and Maskawa originally chose as independent parameters three rotation angles and one CP-violating phase. A number of parametrizations of this kind, different from the original Kobayashi-Maskawa form have been proposed . One of the most commonly used is the “standard” parametrization advocated by the Particle Data Book . From a mathematical point of view, all parametrizations of the flavour mixing matrix containing four, suitably defined, independent parameters are equivalent.
In contrast, the parametrization $`𝐕^{\mathrm{𝐭𝐡}}`$, derived in the previous sections from the breaking of the flavour permutational symmetry, has only two free, linearly independent parameters, namely, the symmetry breaking parameter $`Z^{1/2}`$ and the CP-violating phase $`\mathrm{\Phi }`$. When the best values of the parameters $`Z^{1/2}`$ and $`\mathrm{\Phi }`$ are used, the mixing matrix $`𝐕^{\mathrm{𝐭𝐡}}`$ reproduces the central values of all experimentally determined quantities, that is, the moduli $`|V_{ij}^{exp}|`$, the Jarlskog invariant $`J^{exp}`$ and the three inner angles, $`\alpha `$, $`\beta `$ and $`\gamma `$ of the unitarity triangle. The quality of the fit of $`𝐕^{\mathrm{𝐭𝐡}}`$ to the experimental data is as good as the quality of the fit of the phenomenological parametrizations $`𝐕^{\mathrm{𝐏𝐃𝐆}}`$ and $`𝐕^{\mathrm{𝐊𝐌}}`$ to the same data. More precisely, when the best set of adjustable parameters of each parametrization $`𝐕^{\mathrm{𝐭𝐡}}`$, $`𝐕^{\mathrm{𝐏𝐃𝐆}}`$ and $`𝐕^{\mathrm{𝐊𝐌}}`$, obtained from a $`\chi ^2`$ fit to the same experimental data, the moduli of corresponding entries in the matrices $`𝐕^{\mathrm{𝐭𝐡}}`$, $`𝐕^{\mathrm{𝐏𝐃𝐆}}`$ and $`𝐕^{\mathrm{𝐊𝐌}}`$, are numerically equal and give an equally good representation of the experimentally determined values $`|𝐕_{}^{\mathrm{𝐞𝐱𝐩}}{}_{\mathrm{𝐢𝐣}}{}^{}|`$. From this observation, it follows that the symmetry derived $`𝐕^{\mathrm{𝐭𝐡}}`$ and the phenomenological parametrizations $`𝐕^{\mathrm{𝐏𝐃𝐆}}`$ and $`𝐕^{\mathrm{𝐊𝐌}}`$ should be equivalent up to a rephasing of the quark fields in the mass representation. In the following, we will show that it is possible to derive new theoretical mixing matrices $`\stackrel{~}{𝐕}^{\mathrm{𝐭𝐡}}`$ and $`\widehat{𝐕}^{th}`$, related to $`𝐕^{\mathrm{𝐭𝐡}}`$ by biunitary rephasing transformations, and such that all corresponding entries in $`\stackrel{~}{𝐕}^{\mathrm{𝐭𝐡}}`$ and $`𝐕^{\mathrm{𝐏𝐃𝐆}}`$ or $`\widehat{𝐕}^{\mathrm{𝐭𝐡}}`$ and $`𝐕^{\mathrm{𝐊𝐌}}`$ are equal in modulus and phase. From here, we will obtain exact, explicit expressions for the three mixing angles and the CP-violating phase ocurring in $`𝐕^{\mathrm{𝐏𝐃𝐆}}`$ and $`𝐕^{\mathrm{𝐊𝐌}}`$ as functions of the quark mass ratios and the parameters $`Z^{1/2}`$ and $`\mathrm{\Phi }^{}`$ characterizing the preferred symmetry breaking pattern.
## V Phase equivalence of $`𝐕^{th}`$ and $`𝐕^{PDG}`$
The standard parametrization of the mixing matrix recomended by the Particle Data Group is written in terms of three mixing angles $`\theta _{12},\theta _{23},\theta _{13}`$ and one CP violating phase $`\delta _{13}`$,
$$𝐕^{PDG}=\left(\begin{array}{ccc}c_{12}c_{13}& s_{12}c_{13}& s_{13}e^{i\delta _{13}}\\ s_{12}c_{23}c_{12}s_{23}s_{13}e^{i\delta _{13}}& c_{12}c_{23}s_{12}s_{23}s_{13}e^{i\delta _{13}}& s_{23}c_{13}\\ s_{12}s_{23}c_{12}c_{23}s_{13}e^{i\delta _{13}}& c_{12}s_{23}s_{12}c_{23}s_{13}e^{i\delta _{13}}& c_{23}c_{13}\end{array}\right)$$
(62)
where $`c_{ij}=\mathrm{cos}\theta _{ij}`$ and $`s_{ij}=\mathrm{sin}\theta _{ij}`$.
The range of values of the experimentally determined moduli in $`|V_{ij}^{exp}|`$, as given by Caso et al , corresponds to $`90\%`$ confidence limits on the range of values of the mixing angles of
$$0.219s_{12}0.226,$$
(63)
$$0.037s_{23}0.043,$$
(64)
$$0.002s_{13}0.005.$$
(65)
Seven of the nine absolute values of the CKM entries have been measured directly, by tree level processes. A range of values for the four parameters, $`s_{12},s_{23},s_{13}`$ and $`\delta _{13}`$ which is consistent with the seven direct measurements and the experimentally determined values of the moduli of $`|𝐕|^{exp}`$ , is given by Nir
$$0.2173s_{12}0.2219,$$
(66)
$$0.0378s_{23}0.0412,$$
(67)
$$0.00237s_{13}0.00395,$$
(68)
$`c_{13}`$ is known to deviate from unity only in the sixth decimal place \[, \].
The CP violating phase $`\delta _{13}`$, at present, is not constrained by direct measurements. However, the measurements of CP violation in $`K`$ decays force $`\delta _{13}`$ to lie in the range
$$0\delta _{13}\pi .$$
(69)
The standard parametrization $`𝐕^{PDG}`$ was introduced without taking the possible functional relations between the quark masses and the flavour mixing parameters into account. In contrast, these functional relations are explicitly exhibited in the theoretical expressions, $`V_{ij}^{th}`$, derived in the previous sections. Furthermore, we have seen that, when the best values of the parameters $`Z^{1/2}`$ and $`\mathrm{\Phi }`$ are used, the mixing matrix $`𝐕^{th}`$ reproduces the central values of all experimentally determined quantities, that is, the moduli $`|V_{ij}^{exp.}|`$, the Jarlskog invariant $`J^{exp.}`$ and the three inner angles, $`\alpha `$, $`\beta `$ and $`\gamma `$, of the unitarity triangle . Since the two parametrizations reproduce the same set of experimental data equally well, we are justified in writing
$$|V_{ij}^{th}|=|V_{ij}^{PDG}|=|V_{ij}^{exp}|.$$
(70)
We cannot simply equate $`𝐕^{th}`$ and $`𝐕^{PDG}`$ because the arguments of corresponding matrix elements in the two parametrizations are not equal
$$arg(V_{ij}^{th})arg(V_{ij}^{PDG}).$$
(71)
This difference is of no physical consequence, it reflects the freedom in choosing the unobservable phases of the quark fields in the mass representation. A redefinition of the phases of the quark fields which leaves $`J_c^\mu `$ invariant, will change the argument of $`V_{ij}^{th}`$ but leave the moduli $`|V_{ij}^{th}|`$ invariant,
$`V_{ij}^{th}\stackrel{~}{V}_{ij}^{th}=e^{i\chi _i^u}V_{ij}^{th}e^{i\chi _j^d}.`$ (72)
The phases $`\chi _i^u`$ and $`\chi _j^d`$ ocurring in Eq. (61) will be determined from the requirement that corresponding entries in $`\stackrel{~}{𝐕}^{th.}`$ and $`𝐕^{PDG}`$ be equal,
$$|V_{ij}^{th}|e^{i(w_{ij}^{th}(\chi _i^u\chi _j^d))}=|V_{ij}^{PDG}|e^{iw_{ij}^{PDG}},$$
(73)
in this expression $`w_{ij}^{th}`$ and $`w_{ij}^{PDG}`$ are the arguments of $`V_{ij}^{th}`$ and $`V_{ij}^{PDG}`$ respectively. Since the moduli $`|V_{ij}^{th}|`$ and $`|V_{ij}^{PDG}|`$ are equal, the arguments of the entries in the two parametrizations are related by the set of nine equations
$$\chi _i^u\chi _j^d=w_{ij}^{th}w_{ij}^{PDG}.$$
(74)
The set of Eqs. (74) relate the differences of the unobservable quark field phases to the differences of the arguments of corresponding entries in $`𝐕^{th}`$ and $`𝐕^{PDG}`$.
We notice that the set of Eqs. (74) is overdetermined. In the left hand side of Eqs. (74) there are nine differences of unobsevable quark field phases formed from only six different unknown quark field phases. Since only differences of phases may be determined, the phases themselves are defined only up to an additive constant which may be fixed by giving the value of one of the unknown quark phases. Therefore, in Eqs. (74) there are nine equations to determine five unknowns. This is possible only if a set of 4 consistency conditions is satisfied. The consistency conditions are non trivial relations expressing the five non-vanishing arguments $`w_{ij}^{PDG}`$ of $`V_{ij}^{PDG}`$ in terms only of the known arguments $`w_{ij}^{th}`$ of $`V_{ij}^{th}`$.
A first consistency condition may be derived from the equality of the determinants of $`\stackrel{~}{𝐕}^{\mathrm{𝐭𝐡}}`$ and $`V^{PDG}`$. From the definition of the rephasing transformation, Eqs. (61) and (73), it follows that
$`det𝐕^{th}=det\left[𝐗_u^{}𝐕^{PDG}𝐗_d\right],`$ (75)
in this expression $`𝐗_u`$ and $`𝐗_d`$ are the diagonal unitary matrices of phases ocurring in Eq. (61). The determinant of $`𝐕^{PDG}`$ is one, hence,
$$det\left[𝐗_u^{}𝐕^{PDG}𝐗_d\right]=e^{i_{i=1}^3\left(\chi _i^{(u)}\chi _i^{(d)}\right)}.$$
(76)
Similarly, from the definition of $`𝐕^{th}`$, Eq. (23), we get
$$det𝐕^{th}=det\left[𝐎_u^T𝐏^{ud}𝐎_d\right]=det\left(𝐎_u^T𝐎_d\right)det𝐏^{ud},$$
(77)
the determinant of the orthogonal matrices is one, and the determinant of the diagonal matrix of phases $`𝐏^{ud}`$ is $`e^{i2\mathrm{\Phi }}`$. Taking for $`\mathrm{\Phi }`$ the best value $`\mathrm{\Phi }^{}=\pi /2`$, we obtain
$$det𝐕^{th}=e^{i2\mathrm{\Phi }^{}}=e^{i\pi }.$$
(78)
Substitution of Eq. (76) and Eq. (78) in Eq. (75) gives
$$\underset{i=1}{\overset{3}{}}\left(\chi _i^{(u)}\chi _i^{(d)}\right)=2\mathrm{\Phi }^{}=\pi .$$
(79)
This phase relation guarantees the equality of the determinants of $`\stackrel{~}{𝐕}^{th}`$ and $`𝐕^{PDG}`$.
The sum of the unobservable quark field phases ocurring in the left hand side of Eq. (79) may be computed from Eqs. (74),
$$\underset{i=1}{\overset{3}{}}\left(\chi _i^{(u)}\chi _i^{(d)}\right)=\underset{i=1}{\overset{3}{}}w_{ii}^{th}w_{22}^{PDG}.$$
(80)
Now, we eliminate the unobservable quark field phases between Eq. (79) and Eq. (80), to get,
$$w_{22}^{PDG}=\underset{i=1}{\overset{3}{}}w_{ii}^{th}2\mathrm{\Phi }^{}.$$
(81)
This, relation shows that $`\mathrm{arg}(V_{22}^{PDG})`$ is uniquely determined $`(mod2\pi )`$ in terms of the arguments of the entries in $`𝐕^{th}`$.
A set of consistency conditions for the solution of Eqs. (74) may be derived in a similar way by elimination of the quark field phases. From Eqs. (74), differences of phases of the same quark field type, say $`\left(\chi _j^{(d)}\chi _j^{}^{(d)}\right)`$, may be computed from Eqs. (74) in at least three different ways. This redundancy implies the existence of non-trivial relations among the arguments of the entries of the two parametrizations. For example, from Eqs. (74), the difference $`\left(\chi _2^{(u)}\chi _3^{(d)}\right)\left(\chi _2^{(u)}\chi _2^{(d)}\right)`$ gives
$$\chi _2^{(d)}\chi _3^{(d)}=w_{23}^{th}w_{22}^{th}+w_{22}^{PDG},$$
(82)
and the difference $`\left(\chi _1^{(u)}\chi _3^{(d)}\right)\left(\chi _1^{(u)}\chi _2^{(d)}\right)`$ gives
$$\chi _2^{(d)}\chi _3^{(d)}=w_{13}^{th}w_{12}^{th}+\delta _{13}.$$
(83)
If the phase difference $`(\chi _2^{(d)}\chi _3^{(d)})`$ is eliminated between Eqs. (82) and (83) we get
$$\delta _{13}w_{22}^{PDG}=w_{12}^{th}w_{13}^{th}w_{22}^{th}+w_{23}^{th}.$$
(84)
Using the same elimination procedure for all possible combinations $`\left(\chi _i^{(u)}\chi _j^{(d)}\right)\left(\chi _i^{(u)}\chi _j^{}^{(d)}\right)`$ we derive a set of nine equations, only four of which are linearly independent. One of these is Eq. (84), for the other three we may take
$$w_{21}^{PDG}+w_{22}^{PDG}=w_{11}^{th}w_{12}^{th}w_{21}^{th}+w_{22}^{th},$$
(85)
$$w_{31}^{PDG}w_{32}^{PDG}=w_{11}^{th}+w_{12}^{th}+w_{31}^{th}w_{32}^{th},$$
(86)
and
$$w_{22}^{PDG}+w_{32}^{PDG}=w_{22}^{th}+w_{23}^{th}+w_{32}^{th}w_{33}^{th}.$$
(87)
Since, in $`𝐕^{PDG}`$ there are five entries with non-vanishing arguments, namely, $`w_{13}^{PDG}=\delta _{13},w_{21}^{PDG},w_{22}^{PDG},w_{31}^{PDG}`$ and $`w_{32}^{PDG}`$, we require still one more equation relating the arguments of the entries of the two parametrizations. This is obtained from the phase relations between the determinants of the two matrices, $`𝐕^{th}`$ and $`𝐕^{PDG}`$.
With the help of Eq. (81) we solve Eqs. (84)-(87) for all the other non-vanishing arguments of $`𝐕^{PDG}`$
$$\delta _{13}=w_{11}^{th}+w_{12}^{th}w_{13}^{th}+w_{23}^{th}+w_{33}^{th}2\mathrm{\Phi }^{}$$
(88)
$$w_{21}^{PDG}=w_{21}^{th}+w_{12}^{th}+w_{33}^{th}2\mathrm{\Phi }^{}$$
(89)
$$w_{31}^{PDG}=w_{31}^{th}+w_{12}^{th}+w_{23}^{th}2\mathrm{\Phi }^{}$$
(90)
$$w_{32}^{PDG}=w_{32}^{th}+w_{23}^{th}+w_{11}^{th}2\mathrm{\Phi }^{}.$$
(91)
In this way, we have shown that the arguments $`w_{ij}^{PDG}`$ of $`V_{ij}^{PDG}`$ are uniquely determined (mod $`2\pi `$) by the arguments $`w_{ij}^{th}`$ of $`V_{ij}^{th}`$.
We now return to the question of the quark field phases and the phase transformation from $`V_{ij}^{th}`$ to $`V_{ij}^{PDG}`$. Substitution of Eqs. (81)-(91) into Eqs. (74), gives the differences of the quark field phases explicitly in terms of the known arguments $`w_{ij}^{th}`$ of $`V_{ij}^{th}`$. The quark field phases themselves are determined only up to a common additive constant. Since the quark field phases are unobservable, without loss of generality, we may fix one of them, and solve for the others. In this way, if we set $`\chi _1^d=0`$, we get
$`\chi _1^d`$ $`=`$ $`0^{},`$ (92)
$`\chi _2^d`$ $`=`$ $`w_{11}^{th}w_{12}^{th},`$ (93)
$`\chi _3^d`$ $`=`$ $`w_{23}^{th}w_{33}^{th}w_{12}^{th}+2\mathrm{\Phi }^{},`$ (94)
$`\chi _1^u`$ $`=`$ $`w_{11}^{th},`$ (95)
$`\chi _2^u`$ $`=`$ $`w_{12}^{th}w_{33}^{th}+2\mathrm{\Phi }^{},`$ (96)
$`\chi _3^u`$ $`=`$ $`w_{23}^{th}w_{12}^{th}+2\mathrm{\Phi }^{}.`$ (97)
Then, the diagonal matrices of phases required to compute the phase transformed $`\stackrel{~}{𝐕}^{th}`$ are
$$𝐗_u=diag[e^{iw_{11}^{th}},e^{i(w_{12}^{th}w_{33}^{th}+2\mathrm{\Phi }^{})},e^{i(w_{23}^{th}w_{12}^{th}+2\mathrm{\Phi }^{})}]$$
(98)
and
$$𝐗_d=diag[1,e^{i(w_{11}^{th}w_{12}^{th})},e^{i(w_{12}^{th}w_{23}^{th}w_{33}^{th}+2\mathrm{\Phi }^{})}].$$
(99)
Hence, with the help of Eqs. (88)-(91), we verify that
$$𝐗_𝐮^{}𝐕^{th}𝐗_𝐝=𝐕^{PDG}$$
(100)
is satisfied as an identity, provided that $`|V_{ij}^{th}|=|V_{ij}^{PDG}|`$.
## VI The mixing angles of the Standard Parametrization $`𝐕^{PDG}`$
The standard parametrization $`V_{ij}^{PDG}`$ advocated by the PDG and the symmetry derived parametrization $`V_{ij}^{th}`$, give an equally good representation of the experimentally determined values of the moduli of the entries in the quark mixing matrix $`|V_{ij}^{exp}|`$ . Hence, we may write
$$|V_{ij}^{th}|=|V_{ij}^{PDG}|,$$
(101)
even though $`V_{ij}^{th}`$ has only two adjustable parameters $`(Z^{1/2},\mathrm{\Phi })`$ while the number of adjustable parameters in $`V_{ij}^{PDG}`$ is four, namely, $`(\theta _{12},\theta _{23},\theta _{13},\delta _{13})`$. All entries in $`|V_{ij}^{th}|`$ are explicit functions of the four quark mass ratios $`(m_u/m_t,m_c/m_t,m_d/m_b,m_s/m_b)`$ and the two parameters $`Z^{1/2}`$ and $`\mathrm{\Phi }`$. The equality of the moduli of corresponding entries of the two parametrizations will allow us to derive explicit expressions for the mixing angles in terms of the four quark mass ratios $`(m_u/m_t,m_c/m_t,m_d/m_b,m_s/m_b)`$ and the parameters $`Z^{1/2}`$ and $`\mathrm{\Phi }`$.
From the equality of $`|V_{13}^{th}|`$ and $`|V_{13}^{PDG}|`$, it follows that
$$\mathrm{sin}\theta _{13}=|V_{ub}^{th}|,$$
(102)
if we take $`|V_{ub}^{th}|`$ from (35), and we set $`\mathrm{\Phi }`$ and $`Z^{1/2}`$ equal to their best values $`\mathrm{\Phi }^{}=\pi /2`$ and $`Z^{1/2}=\sqrt{\frac{81}{32}}`$, we get
$`\mathrm{sin}\theta _{13}`$ $`=`$ $`\{{\displaystyle \frac{\stackrel{~}{m}_c(1\stackrel{~}{m}_u\delta _u^{})\stackrel{~}{m}_d\stackrel{~}{m}_s\delta _d^{}}{(1\delta _u^{})(1\stackrel{~}{m}_u)(\stackrel{~}{m}_c+\stackrel{~}{m}_u)(1\delta _d^{})(1+\stackrel{~}{m}_s)(1\stackrel{~}{m}_d)}}`$ (103)
$`+`$ $`[\left({\displaystyle \frac{\stackrel{~}{m}_u(1\stackrel{~}{m}_u\delta _u^{})\delta _d^{}}{(1\stackrel{~}{m}_u)(\stackrel{~}{m}_c+\stackrel{~}{m}_u)(1+\stackrel{~}{m}_s)(1\stackrel{~}{m}_d)}}\right)^{1/2}`$ (104)
$``$ $`\left({\displaystyle \frac{\stackrel{~}{m}_u(1+\stackrel{~}{m}_c\delta _u^{})\delta _u^{}(1\stackrel{~}{m}_d\delta _d^{})(1+\stackrel{~}{m}_s\delta _d^{})}{(1\delta _u^{})(1\stackrel{~}{m}_u)(\stackrel{~}{m}_c+\stackrel{~}{m}_u)(1\delta _d^{})(1+\stackrel{~}{m}_s)(1\stackrel{~}{m}_d)}}\right)^{1/2}]^2\}^{\frac{1}{2}}`$ (105)
The computation of $`\mathrm{sin}\theta _{23}`$ is slightly more involved. From Eq. (62) and the equality of $`|V_{ij}^{th}|`$ and $`|V_{ij}^{PDG}|`$, we obtain
$`\mathrm{sin}\theta _{23}={\displaystyle \frac{|V_{cb}^{PDG}|}{\sqrt{1|V_{ub}^{PDG}|^2}}}={\displaystyle \frac{|V_{cb}^{th}|}{\sqrt{1|V_{ub}^{th}|^2}}}.`$ (106)
Substitution of the expressions (41) and (35) with $`\mathrm{\Phi }^{}=\pi /2`$ and $`Z^{1/2}=\sqrt{\frac{81}{32}}`$ for $`|V_{cb}^{th}|`$ and $`|V_{ub}^{th}|`$ in Eq. (106) gives
$`\mathrm{sin}\theta _{23}`$ $`=`$ $`\sqrt{{\displaystyle \frac{1\stackrel{~}{m}_u}{1+\stackrel{~}{m}_c}}}\{\stackrel{~}{m}_u(1+\stackrel{~}{m}_c\delta _u^{})\stackrel{~}{m}_d\stackrel{~}{m}_s\delta _d^{}+[\sqrt{(1\delta _u^{})\stackrel{~}{m}_c(1+\stackrel{~}{m}_c\delta _u^{})(1\delta _d^{})\delta _d^{}}`$ (107)
$``$ $`\sqrt{\stackrel{~}{m}_c(1\stackrel{~}{m}_u\delta _u^{})\delta _u^{}(1\stackrel{~}{m}_d\delta _d^{})(1+\stackrel{~}{m}_s\delta _d^{})}]^2\}^{1/2}`$ (108)
$`\times `$ $`\{[\sqrt{(1\delta _u^{})\stackrel{~}{m}_u(1\stackrel{~}{m}_u\delta _u^{})(1\delta _d^{})\delta _d^{}}`$ (109)
$``$ $`\sqrt{\stackrel{~}{m}_u(1+\stackrel{~}{m}_c\delta _u^{})\delta _u^{}(1\stackrel{~}{m}_d\delta _d^{})(1+\stackrel{~}{m}_s\delta _d^{})}]^2`$ (110)
$`+`$ $`(1\delta _u^{})(1\stackrel{~}{m}_u)(\stackrel{~}{m}_c+\stackrel{~}{m}_u)(1\delta _d^{})(1+\stackrel{~}{m}_s)(1\stackrel{~}{m}_d)\stackrel{~}{m}_c(1\stackrel{~}{m}_u\delta _u^{})\stackrel{~}{m}_d\stackrel{~}{m}_s\delta _d^{}\}^{1/2}.`$ (111)
Similarly from Eq. (62) and the equality of $`|V_{12}^{th}|`$ and $`|V_{12}^{PDG}|`$, we obtain
$`\mathrm{sin}\theta _{12}={\displaystyle \frac{|V_{us}^{PDG}|}{\sqrt{1|V_{ub}^{PDG}|^2}}}={\displaystyle \frac{|V_{us}^{th}|}{\sqrt{1|V_{ub}^{th}|^2}}}.`$ (112)
Then, substitution of the expressions (32) and (35) for $`|V_{us}^{th}|`$ and $`|V_{ub}^{th}|`$ in (112) gives
$`\mathrm{sin}\theta _{12}`$ $`=`$ $`\sqrt{{\displaystyle \frac{1\stackrel{~}{m}_d}{\stackrel{~}{m}_s+\stackrel{~}{m}_d}}}\{\stackrel{~}{m}_c(1\stackrel{~}{m}_u\delta _u^{})\stackrel{~}{m}_d(1+\stackrel{~}{m}_s\delta _d^{})`$ (113)
$`+`$ $`[\sqrt{(1\delta _u^{})\stackrel{~}{m}_u(1\stackrel{~}{m}_u\delta _u^{})(1\delta _d^{})\stackrel{~}{m}_s(1+\stackrel{~}{m}_s\delta _d^{})}`$ (114)
$`+`$ $`\sqrt{\stackrel{~}{m}_u(1+\stackrel{~}{m}_c\delta _u^{})\delta _u^{}\stackrel{~}{m}_s(1\stackrel{~}{m}_d\delta _d^{})\delta _d^{}}]^2\}^{1/2}`$ (115)
$`\times `$ $`\{[\sqrt{(1\delta _u^{})\stackrel{~}{m}_u(1\stackrel{~}{m}_u\delta _u^{})(1\delta _d^{})\delta _d^{}}`$ (116)
$``$ $`\sqrt{\stackrel{~}{m}_u(1+\stackrel{~}{m}_c\delta _u^{})\delta _u^{}(1\stackrel{~}{m}_d\delta _d^{})(1+\stackrel{~}{m}_s\delta _d^{})}]^2`$ (117)
$`+`$ $`(1\delta _u^{})(1\stackrel{~}{m}_u)(\stackrel{~}{m}_c+\stackrel{~}{m}_u)(1\delta _d^{})(1+\stackrel{~}{m}_s)(1\stackrel{~}{m}_d)`$ (118)
$``$ $`\stackrel{~}{m}_c(1\stackrel{~}{m}_u\delta _u^{})\stackrel{~}{m}_d\stackrel{~}{m}_s\delta _d^{}\}^{1/2}.`$ (119)
The computed values for $`\mathrm{sin}\theta _{12},\mathrm{sin}\theta _{23}`$ and $`\mathrm{sin}\theta _{13}`$ corresponding to the best $`\chi ^2`$ fit of $`|V_{ij}^{th}|`$, $`J^{th}`$ and $`\alpha ^{th}`$, $`\beta ^{th}`$ and $`\gamma ^{th}`$ to the experimentally determined quantities $`|V_{ij}^{exp}|`$, $`J^{exp}`$ and the three inner angles of the unitarity triangle $`\alpha ^{exp}`$, $`\beta ^{exp}`$ and $`\gamma ^{exp}`$ are obtained when the numerical values of $`|V_{us}^{th}|`$, $`|V_{ub}^{th}|`$ and $`|V_{cb}^{th}|`$ computed from Eqs.(32), (35), (41) and given in Eq. (55) are substituted in to Eqs. (102), (106) and (112). In this way, we get
$`\mathrm{sin}\theta _{12}^{}=0.222,`$ (120)
$`\mathrm{sin}\theta _{23}^{}=0.040,`$ (121)
$`\mathrm{sin}\theta _{13}^{}=0.0036.`$ (122)
The numerical value of $`cos\theta _{13}^{}`$ deviates from unity in the sixth decimal place.
We notice that the numerical values of the mixing angles computed from quark masses and the best values of the symmetry breaking parameters coincide almost exactly with the central values of the experimentally determined quantities, as could be expected from Eq. (101). This observation is interesting because, in the case of three families, the most general form of the mixing matrix has at most four free, independent parameters which could be four independent moduli or three mixing angles and one phase as occurs in $`𝐕^{PDG}`$. The symmetry derived $`𝐕^{th}`$ has only two free, real independent parameters. In spite of that, the quality of the fit of $`𝐕^{th}`$ to the experimental data is as good as the quality of the fit of $`𝐕^{PDG}`$ to the same data. The predictive power of $`𝐕^{th}`$ implied by this fact originates in the flavour permutational symmetry of the Standard Model and the assumed symmetry breaking pattern from which the texture in the quark mass matrices and $`𝐕^{th}`$ were derived.
## VII The Cp violating phase $`\delta _{13}`$
The CP violating phase $`\delta _{13}`$ of the standard parametrization $`𝐕^{PDG}`$ of the quark mixing matrix is given in Eq. (88) in terms of the arguments $`w_{ij}^{th}`$ of five entries in the theoretical expression for $`V_{ij}^{th}`$ and the corresponding CP violating phase $`\mathrm{\Phi }`$. Taking from Eq. (55) the numerical values of the arguments $`w_{ij}^{th}`$ and setting $`\mathrm{\Phi }`$ equal to the best value $`\mathrm{\Phi }^{}=\pi /2`$, we obtain the numerical value of $`\delta _{13}`$ corresponding to the best fit of $`|V_{ij}^{th}|`$ to the experimental data
$$\delta _{13}^{}=73.2^{}.$$
(123)
This predicted value of $`\delta _{13}`$ is very close to the numerical value of the third inner angle $`\gamma `$, of the unitarity triangle. The difference may readily be computed in terms of the arguments $`w_{ij}^{th}`$. From the expression for $`\gamma `$
$$\gamma =arg\left[\frac{V_{cb}^{}V_{cd}}{V_{ub}^{}V_{ud}}\right]$$
(124)
we get
$$\gamma =w_{11}^{th}w_{13}^{th}w_{21}^{th}+w_{23}^{th}+\pi $$
(125)
which, when compared with the expression (88) for $`\delta _{13}`$ gives
$$\gamma =\delta _{13}(w_{12}^{th}+w_{21}^{th}+w_{33}^{th}2\mathrm{\Phi }^{}\pi ).$$
(126)
Taking from Eq. (55) the numerical values of the arguments corresponding to the best values $`\mathrm{\Phi }^{}=90^{}`$ and $`Z^{1/2}=\sqrt{\frac{81}{32}}`$, we obtain
$$(w_{12}^{th}+w_{21}^{th}+w_{33}^{th}2\mathrm{\Phi }^{}\pi )=0.04^{}.$$
(127)
This is, indeed, a very small number, and justifies the approximation
$$\gamma \delta _{13}^{}.$$
(128)
According to this, the value of $`|\gamma |`$ computed from quark mass ratios and the best values of the parameters $`Z^{1/2}`$ and $`\mathrm{\Phi }^{}`$ is $`|\gamma |=73.2^{}`$, in agreement with the bounds extracted from the precise measurements of the $`B_d^0`$ oscillation frequency and the measurements of the rates of the exclusive hadronic decays $`B^\pm \pi K`$ and the CP averaged $`B^\pm \pi ^\pm \pi ^0`$ . Exact explicit expressions for the CP violating phase $`\delta _{13}`$ in terms of the four quark mass ratios and the parameters $`Z^{1/2}`$ and $`\mathrm{\Phi }^{}`$ may readily be found; such an expression could be derived from Eq. (88) in terms of the arguments of five matrix elements of $`𝐕^{th}`$. However, a simpler expression, involving only four matrix elements of $`𝐕^{th}`$ may be obtained from the Jarlskog invariant $`J`$.
The Jarlskog invariant may be written in terms of four matrix elements of $`𝐕`$ as
$$J=\mathrm{𝐼𝑚}[V_{12}V_{23}V_{13}^{}V_{22}^{}].$$
(129)
Since $`J`$ is an invariant, its value is independent of the particular parametrization of $`𝐕`$. If we write the right hand side of Eq. (129) in terms of the standard parametrization $`𝐕^{PDG}`$, we obtain
$$\mathrm{sin}\delta _{13}=\frac{J^{th}}{s_{12}s_{13}s_{23}c_{12}c_{13}^2c_{23}}.$$
(130)
The terms in the denominator in the right hand side of this expression were written in Eqs. (102), (106) and (112) in terms of the moduli $`|V_{12}^{th}|`$, $`|V_{13}^{th}|`$ and $`|V_{23}^{th}|`$. Hence,
$$s_{12}s_{13}s_{23}c_{12}c_{13}^2c_{23}=\frac{|V_{12}^{th}||V_{13}^{th}||V_{23}^{th}|[(1|V_{13}^{th}|^2|V_{12}^{th}|^2)(1|V_{13}^{th}|^2|V_{23}^{th}|^2)]^{1/2}}{1|V_{13}^{th}|^2}.$$
(131)
Substitution of Eq. (131) in Eq. (130) gives
$$\mathrm{sin}\delta _{13}=\frac{(1|V_{13}^{th}|^2)\mathrm{𝐼𝑚}[V_{12}^{th}V_{23}^{th}V_{13}^{th}V_{22}^{th}]}{|V_{12}^{th}||V_{13}^{th}||V_{23}^{th}|\sqrt{(1|V_{13}^{th}|^2|V_{12}^{th}|^2)(1|V_{13}^{th}|^2|V_{23}^{th}|^2)}},$$
(132)
the right hand side of this equation may be written in terms of the quark mass ratios and the symmetry breaking parameters $`Z^{1/2}`$ and $`\mathrm{\Phi }^{}`$ with the help of Eqs. (103), (107) and (113).
A simpler expression which leads to a very accurate approximation for $`\delta _{13}`$ is obtained from Eq. (132) if the matrix elements in the square brackets are written as modulus and argument, and use is made of the unitarity of $`𝐕^{th}`$ to simplify the denominator,
$$\mathrm{sin}\delta _{13}=\frac{(1|V_{13}^{th}|^2)|V_{22}^{th}|\mathrm{sin}(w_{12}^{th}+w_{23}^{th}w_{13}^{th}w_{22}^{th})}{|V_{11}^{th}||V_{33}^{th}|}.$$
(133)
Explicit expressions for the arguments $`w_{12}^{th},w_{23}^{th},w_{13}^{th}`$ and $`w_{22}^{th}`$ in terms of the quark mass ratios may be derived from Eqs. (32)-(41) setting $`Z^{1/2}`$ and $`\mathrm{\Phi }`$ equal to their best values $`Z^{1/2}=\sqrt{\frac{81}{32}}`$ and $`\mathrm{\Phi }^{}=\pi /2`$, we get
$$w_{us}^{th}=\pi \mathrm{tan}^1\left(\sqrt{\frac{\stackrel{~}{m}_u\stackrel{~}{m}_s}{\stackrel{~}{m}_c\stackrel{~}{m}_d}}\left[\sqrt{(1\delta _u^{})(1\delta _d^{})}+\sqrt{\delta _u^{}\delta _d^{}\frac{(1+\stackrel{~}{m}_c\delta _u^{})(1\stackrel{~}{m}_d\delta _d^{})}{(1\stackrel{~}{m}_u\delta _u^{})(1+\stackrel{~}{m}_s\delta _d^{})}}\right]\right),$$
(134)
$$w_{cb}^{th}=\pi \mathrm{tan}^1\left(\sqrt{\frac{\stackrel{~}{m}_c}{\stackrel{~}{m}_u\stackrel{~}{m}_d\stackrel{~}{m}_s}}\left[\sqrt{(1\delta _u^{})(1\delta _d^{})}\sqrt{\frac{\delta _u^{}(1\stackrel{~}{m}_u\delta _u^{})(1\stackrel{~}{m}_d\delta _d^{})(1+\stackrel{~}{m}_s\delta _d^{})}{\delta _d^{}(1+\stackrel{~}{m}_c\delta _u^{})}}\right]\right),$$
(135)
$$w_{ub}^{th}=\mathrm{tan}^1\left(\sqrt{\frac{\stackrel{~}{m}_u}{\stackrel{~}{m}_c\stackrel{~}{m}_d\stackrel{~}{m}_s}}\left[\sqrt{(1\delta _u^{})(1\delta _d^{})}\sqrt{\frac{\delta _u^{}(1+\stackrel{~}{m}_c\delta _u^{})(1\stackrel{~}{m}_d\delta _d^{})(1+\stackrel{~}{m}_s\delta _d^{})}{\delta _d^{}(1\stackrel{~}{m}_u\delta _u^{})}}\right]\right),$$
(136)
$$w_{cs}^{th}=\mathrm{tan}^1\left(\sqrt{\frac{\stackrel{~}{m}_c\stackrel{~}{m}_s}{\stackrel{~}{m}_u\stackrel{~}{m}_d}}\left[\sqrt{(1\delta _u^{})(1\delta _d^{})}+\sqrt{\delta _u^{}\delta _d^{}\frac{(1\stackrel{~}{m}_u\delta _u^{})(1\stackrel{~}{m}_d\delta _d^{})}{(1+\stackrel{~}{m}_c\delta _u^{})(1+\stackrel{~}{m}_s\delta _d^{})}}\right]\right).$$
(137)
Computing the second factor in square brackets in the leading order of magnitude, we get
$$w_{us}^{th}\pi \mathrm{tan}^1\left(\sqrt{\frac{\stackrel{~}{m}_u\stackrel{~}{m}_s}{\stackrel{~}{m}_c\stackrel{~}{m}_d}}\right),$$
(138)
$$w_{cb}^{th}\pi \mathrm{tan}^1\left(\sqrt{\frac{\stackrel{~}{m}_c}{\stackrel{~}{m}_u\stackrel{~}{m}_d\stackrel{~}{m}_s}}\left[(1\sqrt{\frac{\delta _u^{}}{\delta _d^{}}})\right]\right),$$
(139)
$$w_{ub}^{th.}\mathrm{tan}^1\left(\sqrt{\frac{\stackrel{~}{m}_u}{\stackrel{~}{m}_c\stackrel{~}{m}_d\stackrel{~}{m}_s}}\left[(1\sqrt{\frac{\delta _u^{}}{\delta _d^{}}})\right]\right),$$
(140)
and
$$w_{cs}^{th}\mathrm{tan}^1\left(\sqrt{\frac{\stackrel{~}{m}_c\stackrel{~}{m}_s}{\stackrel{~}{m}_u\stackrel{~}{m}_d}}\right).$$
(141)
The modulus $`|V_{ub}^{th}|`$ has already been expressed in terms of quark mass ratios and the parameters characterizing the symmetry breaking pattern $`Z^{1/2}`$ and $`\mathrm{\Phi }^{}`$, in Eqs. (102) and (103). Similar expressions for the other moduli ocurring in Eq. (133) may also be given
$`|V_{ud}|`$ $`=`$ $`\left({\displaystyle \frac{\stackrel{~}{m}_c\left(1\stackrel{~}{m}_u\delta _u^{}\right)\stackrel{~}{m}_s\left(1\stackrel{~}{m}_d\delta _d^{}\right)}{\left(1\delta _u^{}\right)\left(1\stackrel{~}{m}_u\right)\left(\stackrel{~}{m}_c+\stackrel{~}{m}_u\right)\left(1\delta _d^{}\right)\left(1\stackrel{~}{m}_d\right)\left(\stackrel{~}{m}_s+\stackrel{~}{m}_d\right)}}\right)^{1/2}`$ (142)
$`\times `$ $`\left\{1+{\displaystyle \frac{\stackrel{~}{m}_u\stackrel{~}{m}_d}{\stackrel{~}{m}_c\stackrel{~}{m}_s}}\left[\left((1\delta _u^{})(1\delta _d^{})\right)^{1/2}+\left(\delta _u^{}\delta _d^{}{\displaystyle \frac{(1+\stackrel{~}{m}_c\delta _u^{})(1+\stackrel{~}{m}_s\delta _d^{})}{(1\stackrel{~}{m}_u\delta _u^{})(1\stackrel{~}{m}_d\delta _d^{})}}\right)^{1/2}\right]^2\right\},`$ (143)
$`|V_{cs}|`$ $`=`$ $`\left({\displaystyle \frac{\stackrel{~}{m}_c\left(1+\stackrel{~}{m}_c\delta _u^{}\right)\stackrel{~}{m}_s\left(1+\stackrel{~}{m}_s\delta _d^{}\right)}{\left(1+\stackrel{~}{m}_c\right)\left(\stackrel{~}{m}_c+\stackrel{~}{m}_u\right)\left(1+\stackrel{~}{m}_s\right)\left(\stackrel{~}{m}_s+\stackrel{~}{m}_d\right)}}\right)^{1/2}`$ (144)
$`\times `$ $`\{[1+\left({\displaystyle \frac{\delta _u^{}\delta _d^{}\left(1\stackrel{~}{m}_u\delta _u^{}\right)\left(1\stackrel{~}{m}_d\delta _d^{}\right)}{\left(1\delta _u^{}\right)\left(1\delta _d^{}\right)\left(1+\stackrel{~}{m}_c\stackrel{~}{m}_u\right)\left(1+\stackrel{~}{m}_s\delta _d^{}\right)}}\right)^{1/2}]^2`$ (145)
$`+`$ $`{\displaystyle \frac{\stackrel{~}{m}_u\stackrel{~}{m}_d}{\stackrel{~}{m}_c\stackrel{~}{m}_s}}{\displaystyle \frac{1}{(1\delta _u^{})(1\delta _d^{})}}\}^{1/2},`$ (146)
$`|V_{tb}|`$ $`=`$ $`\left[{\displaystyle \frac{(1\stackrel{~}{m}_u\delta _u^{})(1+\stackrel{~}{m}_c\delta _u^{})(1\stackrel{~}{m}_d\delta _d^{})(1\stackrel{~}{m}_s\delta _d^{})}{(1\delta _u^{})(1+\stackrel{~}{m}_c)(1\stackrel{~}{m}_u)(1\delta _d^{})(1+\stackrel{~}{m}_s)(1\stackrel{~}{m}_d)}}\right]^{1/2}`$ (147)
$`\times `$ $`\{[1+\left({\displaystyle \frac{\delta _u^{}\delta _d^{}(1\delta _u^{})(1\delta _d^{})}{(1+\stackrel{~}{m}_c\delta _u^{})(1\stackrel{~}{m}_u\delta _u^{})(1\stackrel{~}{m}_d\delta _d^{})(1+\stackrel{~}{m}_s\delta _d^{})}}\right)^{1/2}]^2`$ (148)
$`+`$ $`{\displaystyle \frac{\stackrel{~}{m}_u\stackrel{~}{m}_c\delta _u^{}\stackrel{~}{m}_d\stackrel{~}{m}_s\delta _d^{}}{(1\stackrel{~}{m}_u\delta _u^{})(1+\stackrel{~}{m}_c\delta _u^{})(1+\stackrel{~}{m}_s\delta _d^{})(1\stackrel{~}{m}_d\delta ^{})}}\}^{1/2}.`$ (149)
Computing in the leading order of magnitude, the first factor in the right hand side of Eq. (133) gives
$$\frac{(1|V_{13}^{th}|^2)|V_{22}^{th}|}{|V_{11}^{th}||V_{33}^{th}|}\frac{(1\delta _u^{})(1\stackrel{~}{m}_u)(1\delta _d^{})(1\stackrel{~}{m}_d)}{(1\stackrel{~}{m}_u\delta _u^{})(1\stackrel{~}{m}_d\delta _d^{})}\left(1\frac{\stackrel{~}{m}_u}{\stackrel{~}{m}_d}(\sqrt{\delta _d^{}}\sqrt{\delta _u^{}})\right).$$
(150)
Inserting in to Eq. (150) the numerical values of the mass ratios and $`\sqrt{\delta _d^{}}\sqrt{\delta _u^{}}=0.04`$, we find that the right hand side of Eq.(150) differs from one in the third decimal place,
$$\frac{(1|V_{13}^{th}|^2)|V_{22}^{th}|}{|V_{11}^{th}||V_{33}^{th}|}1.$$
(151)
Therefore,
$$\mathrm{sin}\delta _{13}^{}\mathrm{sin}(w_{us}^{th}+w_{cb}^{th}w_{ub}^{th}w_{cs}^{th}),$$
(152)
taking the numerical values of the argument in the right hand side of Eq. (152) from (55), we obtain
$$\delta _{13}^{}73^{},$$
(153)
in agreement with Eq. (123). The approximate expression Eq. (152) for $`\mathrm{sin}\delta _{13}^{}`$ could also be derived from Eq. (84) if $`w_{22}^{PDG}`$ is neglected. Computing $`w_{22}^{PDG}`$ from Eq. (81) and (134)-(137), we obtain $`w_{22}^{PDG}=0.0018^{}`$ which shows that Eq. (152) is a very good approximation. Since Eq. (81) was derived from the the phase relations expressing the arguments of $`V_{ij}^{PDG}`$ in terms of those of $`V_{ij}^{th}`$, while Eq. (133) was derived from the expression Eq. (132) for the Jarlskog invariant, the agreement found between Eqs. (81) and (133)- (152) is a consistency check of our formalism.
## VIII Phase equivalence of $`𝐕^{th}`$ and the Kobayashi-Maskawa parametrization $`𝐕^{KM}`$
The quark mixing matrix was parametrized by Kobayashi and Maskawa in terms of the three mixing angles, $`\theta _1`$, $`\theta _2`$, and $`\theta _3`$, and one CP violating phase $`\delta _{13}`$
$$𝐕^{KM}=\left(\begin{array}{ccc}c_1& s_1c_3& s_1s_3\\ s_1c_2& c_1c_2c_3s_2s_3e^{i\delta _{KM}}& c_1c_2s_3+s_2s_3e^{i\delta _{KM}}\\ s_1s_2& c_1s_2c_3c_2s_3e^{i\delta _{KM}}& c_1s_2s_3c_2s_3e^{i\delta _{KM}}\end{array}\right)$$
(154)
where $`c_i=\mathrm{cos}\theta _i`$ and $`s_i=\mathrm{sin}\theta _i`$.
It is readily verified that
$$det𝐕^{KM}=e^{i\delta _{KM}}$$
(155)
As discussed in section IV, the parametrization $`𝐕^{th}`$ derived from the breaking of the flavour symmetry and the Kobayashi-Maskawa parametrization $`𝐕^{KM}`$, give an equally good representation of the values of the moduli $`|𝐕_{ij}^{exp}|`$ of the mixing matrix even if $`𝐕^{th}`$ has only two free, linearly independent parameters. Hence, the two parametrizations are equivalent up to a rephasing of the quark fields. Therefore, we may define a phase transformed $`\widehat{𝐕}^{th}`$, such that all entries in $`\widehat{𝐕}_{ij}^{th}`$ are numerically equal to the corresponding entries in $`𝐕_{ij}^{KM}`$ both in modulus and phase,
$$\widehat{𝐕}^{th}=𝐗_{KM}^u𝐕^{th}𝐗_{KM}^d=𝐕^{KM}.$$
(156)
The diagonal matrices
$$𝐗_{KM}^u=diag[e^{i\varphi _1^u},e^{i\varphi _2^u}e^{i\varphi _3^u}]$$
(157)
and
$$𝐗_{KM}^d=diag[e^{i\varphi _1^d},e^{i\varphi _2^d}e^{i\varphi _3^d}]$$
(158)
are determined from the equality of $`\widehat{𝐕}^{th}`$ and $`𝐕_{ij}^{KM}`$,
$$\varphi _i^u\varphi _j^d=w_{ij}^{th}w_{ij}^{KM}.$$
(159)
where
$$w_{ij}^{KM}=arg(𝐕_{ij}^{KM}).$$
(160)
In the left hand side of Eqs. (159) there are nine differences of phases formed from only six unobservables quark field phases. Differences of phases of the same quark field type, say $`(\varphi _j^d\varphi _j^{}^d)`$, may be computed from Eqs. (159) in at least three different ways. Elimination of the unobservables differences of quark field phases between these expressions gives a set of consistency conditions relating the known $`w_{ij}^{th}`$ and $`w_{ij}^{KM}`$.
$$w_{12}^{KM}=\pi ,$$
(161)
$$w_{13}^{KM}=\pi ,$$
(162)
$$w_{22}^{KM}=w_{11}^{th}w_{12}^{th}w_{21}^{th}+w_{22}^{th}+\pi ,$$
(163)
$$w_{23}^{KM}=w_{11}^{th}w_{13}^{th}w_{21}^{th}+w_{23}^{th}+\pi ,$$
(164)
$$w_{32}^{KM}=w_{11}^{th}w_{12}^{th}w_{31}^{th}+w_{32}^{th}+\pi ,$$
(165)
$$w_{33}^{KM}=w_{11}^{th}w_{13}^{th}w_{31}^{th}+w_{33}^{th}+\pi ,$$
(166)
when the expressions (161-166) are substituted in Eqs. (159) we obtain the differences of the quark field phases as functions of the arguments $`w_{ij}^{th}`$ of $`𝐕_{ij}^{th}`$. The quark field phases themselves are determined only up to a common additive constant which is fixed by giving the value of one of them. If we set $`\chi _i^u=0`$, we get
$$𝐗_{\mathrm{𝐊𝐌}}^𝐮=diag[1,e^{i(w_{11}^{th}w_{21}^{th})},e^{i(w_{11}^{th}w_{31}^{th})}]$$
(167)
and
$$𝐗_{\mathrm{𝐊𝐌}}^𝐝=diag[e^{iw_{11}^{th}},e^{i(w_{12}^{th}\pi )},e^{i(w_{13}^{th}\pi )}].$$
(168)
With the help of (167) and (168) it is readily verified that (156) is satisfied as an identity.
## IX The CP violating phase $`\delta _{KM}`$ of the Kobayashi-Maskawa parametrization.
The CP violating phase $`\delta _{KM}`$ is implicitly determined by the equations (161 \- 166). An explicit expression for $`\delta _{KM}`$ in terms of the arguments $`w_{ij}^{th}`$ of $`𝐕_{ij}^{th}`$ may be obtained from the identity
$$det[𝐕^{KM}]=det[\widehat{𝐕}^{th}],$$
(169)
and
$$det[\widehat{𝐕}^{th}]=e^{i_{i=1}^3(\varphi _i^u\varphi _i^d)}e^{i2\mathrm{\Phi }^{}}.$$
(170)
Substitution of (155) and (170) in (169) gives
$$\delta _{KM}=\underset{i=1}{\overset{3}{}}(\varphi _i^d\varphi _i^u)+2\mathrm{\Phi }^{}\pi $$
(171)
Now, from Eqs. (159) and (161 \- 166),
$$\delta _{KM}=w_{ud}^{th}w_{us}^{th}w_{ub}^{th}w_{cd}^{th}w_{td}^{th}.$$
(172)
This expression gives the CP-violating phase $`\delta _{KM}`$ of the Kobayashi-Maskawa parametrization, $`𝐕^{KM}`$, in terms of the known arguments of the symmetry derived parametrization $`𝐕^{th}`$.
Taking the numerical values of the arguments $`w_{ij}^{th}`$ ocurring in (172) from Eq. (55), we obtain
$$\delta _{KM}=96.4^{}$$
(173)
or
$$\delta _{KM}=\pi 83.6^{}.$$
(174)
The value of the inner angle $`\alpha `$ of the unitarity triangle found in our $`\chi ^2`$ fit of $`𝐕^{th}`$ to the experimentally determined values of the moduli $`|𝐕_{ij}^{exp}|`$ is
$$\alpha =83.6^{}.$$
(175)
Therefore, within the acuracy of our best fit to the experimental data
$$\delta _{KM}\pi \alpha .$$
(176)
This is only an approximation, to derive an exact relation we compute $`\alpha `$ from
$$\alpha =\mathrm{arg}\left(\frac{V_{ub}^{}V_{ud}}{V_{tb}^{}V_{td}}\right),$$
(177)
substitution of $`𝐕_{ij}^{th}`$ for $`𝐕_{ij}`$ in Eq. (177) gives
$$\alpha =\pi w_{ub}^{th}+w_{ud}^{th}+w_{tb}^{th}w_{td}^{th}.$$
(178)
Comparing (178) with (172) gives
$$\delta _{KM}=\alpha (w_{us}^{th}+w_{cd}^{th}w_{tb}^{th}+\pi ).$$
(179)
Taking the numerical values of the $`w_{ij}^{th}`$ from (55) we get
$$w_{us}^{th}+w_{cd}^{th}+w_{tb}^{th}+\pi =3\pi +0.04^{}.$$
(180)
Hence,
$`\delta _{KM}=\pi \alpha 0.04^{}mod(2\pi )`$ (181)
which shows that (176) is an excellent approximation to the numerical value of $`\delta _{KM}`$.
In passing, we notice that
$$\alpha =w_{tb}^{KM}$$
(182)
is an exact result.
## X The mixing angles in the parametrization of Kobayashi-Maskawa.
Once it is established that when the best set of parameters of each parametrization, $`𝐕^{KM}(\theta _1,\theta _2,\theta _3,\delta _{KM})`$ and $`𝐕^{th}(\stackrel{~}{m}_i,Z^{1/2},\mathrm{\Phi })`$ is used, the moduli of corresponding entries in the two parametrizations are numerically equal, we may identify corresponding entries and solve for the mixing angles $`\theta _1,\theta _2`$ and $`\theta _3`$ in terms of the quark mass ratios and the parameters $`\mathrm{\Phi }`$ and $`Z`$.
From the identification
$$\mathrm{cos}\theta _1=|𝐕_{ud}^{th}|,$$
(183)
we obtain
$`\mathrm{cos}\theta _1`$ $`=`$ $`\left({\displaystyle \frac{(1\stackrel{~}{m}_u\delta _u^{})(1\stackrel{~}{m}_d\delta _d^{})}{(1\delta _u^{})(1\stackrel{~}{m}_u)(1\stackrel{~}{m}_d)(1\delta _d^{})(1+\stackrel{~}{m}_u/\stackrel{~}{m}_c)(1+\stackrel{~}{m}_d/\stackrel{~}{m}_s)}}\right)^{1/2}`$ (184)
$`\times `$ $`\left[1+{\displaystyle \frac{\stackrel{~}{m}_u\stackrel{~}{m}_d}{\stackrel{~}{m}_c\stackrel{~}{m}_s}}\left(\left((1\delta _u^{})(1\delta _d^{})\right)^{1/2}+\left(\delta _u^{}\delta _d^{}{\displaystyle \frac{(1+\stackrel{~}{m}_c\delta _u^{})(1+\stackrel{~}{m}_s\delta _d^{})}{(1\stackrel{~}{m}_u\delta _u^{})(1+\stackrel{~}{m}_d\delta _d^{})}}\right)^{1/2}\right)^2\right]^{1/2}`$ (185)
The angle $`\theta _2`$ is obtained from the identification
$$\mathrm{sin}\theta _1\mathrm{sin}\theta _2=|𝐕_{td}^{th}|$$
(186)
which gives
$$\mathrm{sin}\theta _2=\frac{|𝐕_{td}^{th}|}{\mathrm{sin}\theta _1}$$
(187)
where
$`\mathrm{sin}\theta _1=\sqrt{1|𝐕_{ud}^{th}|^2},`$ (188)
$`|𝐕_{ud}^{th}|`$ is given in Eqs. (183) and (184), and
$`𝐕_{td}^{th}`$ $`=`$ $`\left({\displaystyle \frac{\stackrel{~}{m}_u\stackrel{~}{m}_c\delta _u\stackrel{~}{m}_s(1\stackrel{~}{m}_d\delta _d)}{(1\delta _u)(1+\stackrel{~}{m}_c)(1\stackrel{~}{m}_u)(1\delta _d)(1\stackrel{~}{m}_d)(\stackrel{~}{m}_s+\stackrel{~}{m}_d)}}\right)^{1/2}`$ (189)
$`+`$ $`\{\left({\displaystyle \frac{(1\stackrel{~}{m}_u\delta _u)(1+\stackrel{~}{m}_c\delta _u)\stackrel{~}{m}_d(1+\stackrel{~}{m}_s\delta _u)\delta _d}{(1\delta _u)(1+\stackrel{~}{m}_c)(1\stackrel{~}{m}_u)(1\delta _d)(1\stackrel{~}{m}_d)(\stackrel{~}{m}_s+\stackrel{~}{m}_d)}}\right)^{1/2}`$ (190)
$`+`$ $`\left({\displaystyle \frac{\delta _u\stackrel{~}{m}_d(1\stackrel{~}{m}_d\delta _d)}{(1+\stackrel{~}{m}_c)(1\stackrel{~}{m}_u)(1\stackrel{~}{m}_d)(\stackrel{~}{m}_s+\stackrel{~}{m}_d)}}\right)^{1/2}\}e^{i\mathrm{\Phi }}`$ (191)
Computing in the leading order, we obtain
$`\mathrm{sin}\theta _2`$ $``$ $`\sqrt{{\displaystyle \frac{1+m_u/m_c}{1+\frac{m_um_s}{m_cm_d}}}}Z^{1/2}\{{\displaystyle \frac{\stackrel{~}{m}_s\stackrel{~}{m}_d}{\sqrt{(1\stackrel{~}{m}_d)(1+\stackrel{~}{m}_s)+2Z(\stackrel{~}{m}_s\stackrel{~}{m}_d)(1+\frac{(\stackrel{~}{m}_s\stackrel{~}{m}_d)}{2})}}}`$ (192)
$``$ $`{\displaystyle \frac{\stackrel{~}{m}_c\stackrel{~}{m}_u}{\sqrt{(1\stackrel{~}{m}_u)(1+\stackrel{~}{m}_c)+2Z(\stackrel{~}{m}_c\stackrel{~}{m}_u)(1+\frac{(\stackrel{~}{m}_c\stackrel{~}{m}_u)}{2})}}}\}`$ (193)
similarly the mixing angle $`\theta _3`$ may be obtained from the identification
$$\mathrm{sin}\theta _1\mathrm{sin}\theta _3=|𝐕_{ub}^{th}|$$
(194)
Then,
$$\mathrm{sin}\theta _3=\frac{|𝐕_{ub}^{th}|}{\sqrt{1|𝐕_{ud}^{th}|^2}}$$
(195)
where $`𝐕_{ub}^{th}`$ is given in Eq. (35) and $`\mathrm{sin}\theta _1`$ is given in Eq. (188). Computing in the leading order of magnitude, we find,
$`\mathrm{sin}\theta _3`$ $``$ $`\sqrt{{\displaystyle \frac{1+m_d/m_s}{1+\frac{m_dm_c}{m_sm_u}}}}Z^{1/2}\{{\displaystyle \frac{\stackrel{~}{m}_s\stackrel{~}{m}_d}{\sqrt{(1\stackrel{~}{m}_d)(1+\stackrel{~}{m}_s)+2Z(\stackrel{~}{m}_s\stackrel{~}{m}_d)(1+\frac{(\stackrel{~}{m}_s\stackrel{~}{m}_d)}{2})}}}`$ (196)
$``$ $`{\displaystyle \frac{\stackrel{~}{m}_c\stackrel{~}{m}_u}{\sqrt{(1\stackrel{~}{m}_u)(1+\stackrel{~}{m}_c)+2Z(\stackrel{~}{m}_c\stackrel{~}{m}_u)(1+\frac{(\stackrel{~}{m}_c\stackrel{~}{m}_u)}{2})}}}\}`$ (197)
The exact values computed for $`\mathrm{sin}\theta _1`$, $`\mathrm{sin}\theta _2`$ and $`\mathrm{sin}\theta _3`$ corresponding to the best fit of $`𝐕_{ij}^{th}`$, $`J^{th}`$ and $`\alpha ^{th}`$, $`\beta ^{th}`$ and $`\gamma ^{th}`$ to the experimentally determined quantities are obtained from Eqs. (183, 187) and (195) when the numerical values of $`𝐕_{ud}^{th}`$, $`𝐕_{td}^{th}`$ and $`𝐕_{ub}^{th}`$ computed from Eqs. (142), (184) and (189) are substituted. We obtain
$$\mathrm{sin}\theta _1=0.2225$$
(198)
$$\mathrm{sin}\theta _2=0.0384$$
(199)
$$\mathrm{sin}\theta _3=0.0162$$
(200)
and, from Eq. (173),
$$\mathrm{sin}\delta _{KM}=0.9938.$$
(201)
If the mixing angles $`\theta _1`$$`,\theta _2`$ and $`\theta _3`$ are restricted to lie in the first quadrant, the corresponding numerical values of the angles are
$`\theta _1=12.86^{}\theta _2=2.2^{}and\theta _3=0.93^{}`$ (202)
which together with the numerical value of the CP-violating phase $`\delta _{KM}`$ found in the previous section, $`\delta _{KM}=96.4^{}`$, gives the best set of values of mixing parameters in the Kobayashi-Maskawa parametrization of the mixing matrix corresponding to the best set of parameters $`\mathrm{\Phi }^{}=90^{}`$ and $`Z^{1/2}=\sqrt{81/32}`$ of the flavour symmetry breaking derived parametrization $`𝐕^{th}`$ of the mixing matrix.
As expected, from the way they were obtained, the numerical values of the moduli $`|𝐕_{ij}^{KM}|`$ computed with the help of these numerical values of the mixing angles $`\theta _1`$, $`\theta _2`$, $`\theta _1`$, and $`\delta _{KM}`$ reproduce the central values of the experimentally determined $`|𝐕_{ij}^{exp}|`$, given in Caso $`\mathrm{𝑒𝑡}\mathrm{𝑎𝑙}`$.
In the Kobayashi-Maskawa parametrization, the Jarlskog invariant is
$$J=s_1^2s_2s_3c_1c_2c_3s_\delta $$
(203)
The corresponding numerical value is
$$J=3.0\times 10^5$$
(204)
in excellent agrement with $`J^{exp}`$.
## XI Summary and Conclusions
In the Standard Electroweak model of particle interactions, both, the masses of the quarks as well as the mixing parameters and the CP-violating phase appear as free independent parameters. In consequence, phenomenological parametrizations of the quark mixing matrix were introduced without taking the possible functional relations between the quark masses and the flavour mixing parameters into account. These functional relations are explicitly exhibited in the theoretical quark mixing matrix $`𝐕^{th}`$ derived from the breaking of the flavour permutational symmetry in previous works , and reviewed in sections II and III of this paper.
In this work, we explicitly exhibit the phase equivalence of the theoretical mixing matrix, $`𝐕^{th}`$, and two phenomenological parametrizations: the original Kobayashi-Maskawa , $`𝐕^{KM}`$, and the standard parametrization $`𝐕^{PDG}`$ advocated by the Particle Data Group , , which are written in terms of three mixing angles and one CP-violating phase.
The equality of the moduli of corresponding entries in $`𝐕^{th}`$ and $`𝐕^{KM}`$ or $`𝐕^{PDG}`$ allowed us to derive, exact, explicit expressions for the mixing angles and the CP-violating phase of the two phenomenological parametrizations as functions of four quark mass ratios $`m_u/m_t`$, $`m_c/m_t`$, $`m_d/m_b`$, $`m_s/m_b`$, and the flavour symmetry breaking parameters: $`Z^{1/2}`$ and $`\mathrm{\Phi }`$.
The numerical values of the mixing parameters of the PDG advocated standard parametrization $`\mathrm{sin}\theta _{12}^{},\mathrm{sin}\theta _{23}^{},\mathrm{sin}\theta _{13}^{}`$ computed from the quark mass ratios and the best values of the parameters $`Z^{1/2}`$ and $`\mathrm{\Phi }^{}`$, coincide almost exactly with the central values of the same mixing parameters determined from a fit to the experimental data , as could be expected from the phase equivalence of $`𝐕^{th}`$ and $`𝐕^{PDG}`$, and the good agreement of $`|𝐕_{ij}^{th}|`$ with the corresponding experimentally determined values $`|𝐕_{ij}^{exp}|`$.
Similarly, from the equality of the moduli of corresponding entries in $`𝐕^{th}`$ and $`𝐕^{KM}`$, we derived exact, explicit expressions for the Kobayashi-Maskawa mixing parameters, $`\mathrm{sin}\theta _1,\mathrm{sin}\theta _2`$ and $`\mathrm{sin}\theta _3`$ as functions of the four quark mass ratios and the flavour symmetry breaking parameters $`Z^{1/2}`$ and $`\mathrm{\Phi }`$. As in the previous case, the numerical values of the mixing parameters, $`\mathrm{sin}\theta _1^{},\mathrm{sin}\theta _2^{}`$ and $`\mathrm{sin}\theta _3^{}`$ and the CP-violating phase computed from our expressions and the quark mass ratios and the best values, $`Z^{1/2}`$ and $`\mathrm{\Phi }^{}`$, of the symmetry breaking parameters, are such that the numerical values of $`|𝐕_{ij}^{KM}|`$ reproduce almost exactly the central values of the experimentally determined values of the moduli, $`|𝐕_{ij}^{exp}|`$, and the Jarlskog invariant, $`J^{exp}`$, as given in Caso $`\mathrm{𝑒𝑡}.\mathrm{𝑎𝑙}.`$ .
The numerical values of the CP-violating phase $`\delta _{13}=arg\left(𝐕_{ub}^{PDG}\right)`$ computed from quark mass ratios and the best values of the parameters $`Z^{1/2}`$ and $`\mathrm{\Phi }^{}`$ is $`\delta _{13}=73.2^{}`$ in good agreement with the bounds extracted from the precise measurements of the $`B_d^o`$ oscillations frequency and the measurements of the rates of the exclusive hadronic decays $`B^\pm \pi K`$ . The numerical values of the CP-violating phase $`\delta _{KM}`$ of the Kobayashi-Maskawa parametrization computed from quark mass ratios and the best values of $`Z^{1/2}`$ and $`\mathrm{\Phi }`$ is $`\delta _{KM}^{}=96.4^{}`$. It may be worth to remark that $`\delta _{KM}`$ is not a small number as it is sometimes assumed -.
In conclusion, the three mixing angles and the CP-violating phase which appear in the phenomenological parametrization of the quark mixing matrix as free, linearly independent parameters are expressed in this work as functions of four quark mass ratios and two flavours symmetry breaking parameters $`Z^{1/2}`$ and $`\mathrm{\Phi }`$. The numerical values of the mixing angles and CP-violating phase computed from quark mass ratios and the best values of the symmetry breaking parameters $`Z^{1/2}=\sqrt{81/32}`$ and $`\mathrm{\Phi }^{}=90^{}`$ coincide almost exactly with the central values of the same mixing parameters determined from a fit to the experimental data. The predictive power of the theoretical quark mixing matrix $`𝐕^{th}`$ implied by this fact has its origin in the flavour permutational symmetry of the Standard Model and the assumed symmetry breaking pattern from which the texture in the quark mass matrices and the quark mixing matrix $`𝐕^{th}`$ were derived.
## Acknowledgments
We are indebted to Prof. M. D. Scadron for some useful discussions and critical remarks. This work was partially supported by DGAPA-UNAM under contract No. PAPIIT-IN125298 and by CONACYT (México) under contract 32238E.
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# Untitled Document
SPIN-2000/07
hep-th/0003005
Determinism and Dissipation in
Quantum Gravity
Gerard ’t Hooft
Institute for Theoretical Physics
University of Utrecht, Princetonplein 5
3584 CC Utrecht, the Netherlands
and
Spinoza Institute
Postbox 80.195
3508 TD Utrecht, the Netherlands
e-mail: g.thooft@phys.uu.nl
internet: http://www.phys.uu.nl/~thooft/
Abstract
Without invalidating quantum mechanics as a principle underlying the dynamics of a fundamental theory, it is possible to ask for even more basic dynamical laws that may yield quantum mechanics as the machinery needed for its statistical analysis. In conventional systems such as the Standard Model for quarks and leptons, this would lead to hidden variable theories, which are known to be plagued by problems such as non-locality. But Planck scale physics is so different from field theories in some flat background space-time that here the converse may be the case: we speculate that causality and locality can only be restored by postulating a deterministic underlying theory. A price to be paid may be that the underlying theory exhibits dissipation of information.
1. Introduction In the opening lecture of this School$`^\text{1}`$ it was explained that the physical degrees of freedom of a black hole are distributed on its horizon in such a way that there appears to be one Boolean degree of freedom per unit of surface area of size
$$A_0=4\mathrm{ln}\mathrm{\hspace{0.17em}2}L_{\mathrm{Planck}}^2.$$
$`(1.1)`$
These are the degrees of freedom that appear in the equations of motion, c.q. the Schrödinger equation for a black hole. In turn, according to the principles of General Relativity, these same equations of motion should apply to what happens in the (nearly) flat space-time experienced by an in-going observer. This led us to the formulation of ”the Holographic Principle”, according to which the physical degrees of freedom that describe all physical events in some definite region of space and time, can be mapped onto a two-dimensional surface in such a way that there is exactly one Boolean degree of freedom per unit of surface of size $`A_0`$.
One striking aspect of this observation is that the same region of space-time could be used to describe a different black hole with its horizon some place else. Apparently, one may choose any member of an infinite set of possible surfaces to project the physical degrees of freedom onto it. Also, one should be able to map directly from one surface onto another.
It now appears to be a satisfactory feature of string theories$`^\text{2}`$ and certain semi-perturbative extensions of them, that they manage to reproduce this holographic principle$`^\text{3}`$. This principle implies that any complete theory combining quantum mechanics with gravity should exhibit an upper limit to the total number of independent quantum states that increases exponentially with the surface area of a system, rather than its volume.
All this raises a number of important questions. First, what does locality mean for such theories? And how can notions such as causality, unitarity, and local Lorentz invariance make sense if there is no trace of ‘locality’ left? In this lecture, a theory is developed that will not postulate the quantum states as being its central starting point, but rather classical, deterministic degrees of freedom$`^\text{4}`$. Quantum states, being mere mathematical devices enabling physicists to make statistical predictions, may turn out to be derived concepts, with a not strictly locally formulated definition. Once it is realized that quantum states may be non-local, derived concepts, it is natural to consider making one more step.
In the past, many versions of hidden variable theories were dismissed by a majority of researchers for two reasons: one reason was that these theories did not seem to work properly in the sense that counter examples could be constructed using eigenstates of certain symmetries: rotation symmetry, isospin symmetry, and so on. The second reason was that there appeared to be no need for such theories.
We observe that most of the familiar symmetries are absent at the Planck scale. There are clearly no conservation laws such as isospin, and, in the absence of true locality one cannot rotate any system with respect to a reference system, since they do not decouple. Constructing counter examples to hidden variable theories then becomes a lot harder. Also we claim that some relaxed version of quantum mechanics could be extremely helpful in bypassing the holographic principle; locality could be restored, and once again causality could be reconciled with (a weaker version of) local Lorentz invariance, or general coordinate invariance. We suspect that such steps may be needed in constructing logically coherent theories for Planck scale physics.
Thus, one may have one or several motivations for a reconsideration of “hidden variable” theories:
$`i.`$Einstein’s wish for “reality”. At the time this is written, the quantum mechanical doctrine, according to which all physical states form a Hilbert space and are controlled by non-commuting operators, is fully taken for granted in theoretical physics. No return to a more deterministic description of “reality” is considered necessary; to the contrary, string theorists often give air to their suspicion that the real world is even crazier than quantum mechanics. One might however complain that the description of what really constitutes concepts such as space, time, matter, causality, and the like, is becoming increasingly and uncomfortably obscure. By some physicists this may be regarded as an inescapable course of events, with which we shall have to learn to live, but others such as this author strongly prefer a more complete description of the notion of reality. We admit however that at scales relevant to atomic physics, or even the Standard Model, there is no direct logical inconsistency to be found in Quantum Mechanics.
$`ii.`$“Quantum Cosmology”. How would one describe the ‘wave function of the universe’? An extremely important example of a quantum cosmological model, is a model of gravitating particles in 1 time, 2 space dimensions$`^\text{5}`$. Here, a complete formalism for the quantum version at first sight seems to be straightforward$`^\text{6}`$, but when it comes to specifying exact details, one discovers that we cannot rigorously define what quantum mechanical amplitudes are, what it means when it is claimed that “the universe will collapse with such-and-such probability”, what and where the observers are, what they are made of, and so on. Eventually, one would have to admit that at cosmological scales, any experiment is done only once, yielding answers either ‘yes’ or ‘no’, but never probability distributions. The cosmological wave function is a dubious notion. Quite conceivably, quantum mechanics as we know it only refers to repeatable experiments at small time and distance scales. The true laws of physics are about certainties, not probabilities.
Note that, since the entire hamiltonian of the universe is exactly conserved, the “wave function of the universe” would be in an exact eigenstate of the hamiltonian, and therefore, the usual Schrödinger equation is less appropriate than the description of the evolution in the so-called Heisenberg representation. Quantum states are space-time independent, but operators may depend on space-time points — although only if the location of these space-time points can be defined in a coordinate-free manner! Note that, besides energy, also total momentum and angular momentum of the universe must be conserved (and they too must be zero).$`p_(`$
iii. Even at a local scale (i.e. not cosmological), there are problems that we could attribute to a clash with Quantum Mechanics. Apart from the question of the cosmological principle, these are:
\- the non-renormalizability of gravity;
\- the fact that the gravitational action (the Einstein-Hilbert action) is not properly bounded in Euclidean space, while the Maxwell and Yang-Mills actions are. This is related to the fundamental instability of the gravitational force.
\- topologically non-trivial quantum fluctuations. They could destroy the causal coherence of any theory. Perhaps most such fluctuations may have to be outlawed, as they would also require the boundary conditions to fluctuate into topologically non-trivial ones.
\- black holes cause the most compelling conflicts with local quantum mechanics.
\- there still is the mystery of the cosmological constant. It appears to require a reconsideration not only of physical principles at the Planck scale, but also at cosmological scales, since we are dealing here with an infrared divergence that appears to be cancelled out in a way that requires new physics.
2. States and probabilities Imagine a universe in which different “states” are possible. A state may be characterized any possible way:
$$|\psi =|\stackrel{}{x}_1,\stackrel{}{\sigma }_1,\text{},\begin{array}{c}\text{your name and}\\ \text{telephone}\\ \text{number}\end{array},$$
$`(2.1)`$
There are two ways in which we can introduce the equations of motion, and what we mean by ‘state’ depends on that. At first sight it seems most natural to postulate that a state evolves in time. So, at a given time $`t`$ we have a state, and there is some equation that tells us what this state looks like at later (or maybe also at earlier) times $`t^{}`$. We refer to this as the Schrödinger picture. But when we realize that the notion of time may depend on the clocks used, or more generally, that it requires the introduction of Cauchy surfaces, we might opt for a different notion of “state”. A state is then defined as time-independent. The universe is in a particular state, and in this state all observables may depend on time; the time variable is then linked to the observable, $`𝒪(t)`$. This is the Heisenberg picture. It is familiar from quantum mechanics, where now observables are said to be time-dependent, obeying an evolution equation of the form $`\mathrm{d}𝒪(t)/\mathrm{d}t=i[𝒪(t),H]`$. We will frequently switch between Schrödinger and Heisenberg pictures.
Let us begin with the Schrödinger picture. For simplicity then we take time to be discrete (in the Heisenberg picture, the question whether time is discrete or continuous would not be so important). A simple example is a universe that can only be in three different states, and we have a prescription for the evolution as follows:
This of course is a completely deterministic universe. Nevertheless, it may be useful to introduce the Hilbert space spanned by these three states, so as to enable us to handle the evolution statistically. In this space, the one-time-step evolution operator would be
$$U=\left(\begin{array}{ccc}0& 0& 1\\ 1& 0& 0\\ 0& 1& 0\end{array}\right),$$
$`(2.2)`$
a unitary matrix. Suppose we also consider states of the form
$$|\psi =\alpha |1+\beta |2+c|3,$$
$`(2.3)`$
then after one time step, we would have
$$|\psi _{t+1}=\left(\begin{array}{ccc}0& 0& 1\\ 1& 0& 0\\ 0& 1& 0\end{array}\right)|\psi _t=U(t,t+1)|\psi _t.$$
$`(2.4)`$
We may define the probability $`P(i)`$ for being in a state $`|i`$ as
$$P(1)=|\alpha |^2;P(2)=|\beta |^2;P(3)=|\gamma |^2,$$
$`(2.5)`$
and observe that probability conservation corresponds to unitarity of the evolution matrix $`U`$.
Let us now turn to a basis in which $`U`$ is diagonal:
$$U\left(\begin{array}{ccc}1& 0& 0\\ 0& e^{2\pi /3}& 0\\ 0& 0& e^{2\pi /3}\end{array}\right).$$
$`(2.6)`$
One then may write
$$U=\mathrm{exp}\left(iH\delta t\right),$$
$`(2.7)`$
where the unit time interval $`\delta t`$ will often be taken to be one. So, we can take as our ‘hamiltonian’:
$$H=\left(\begin{array}{ccc}0& & \\ & 2\pi /3& \\ & & +2\pi /3\end{array}\right).$$
$`(2.8)`$
Note that this is the hamiltonian of an atom with a magnetic moment in a homogeneous magnetic field, although it should of course also be noted that all eigenvalues of $`H`$ are only well-defined modulo integral multiples of $`2\pi `$.
All that was needed for the above manipulations is that the deterministic law used as a starting point was time-reversible, otherwise the matrix $`U`$ would not have been unitary. Thus, we make our first important observation:$`^\text{7}`$ any deterministic, time-reversible system can be described using a quantum mechanical Hilbert space, where states obey a Schrödinger equation, and where the absolute squares of the coefficients of the wave functions represent probabilities. Prototype examples are clocks that count periodically over $`N`$ different states — they can be mapped onto atoms with spin $`j`$ in a magnetic field, $`N=2j+1`$.
The converse is not true: most simple quantum mechanical systems do not allow a deterministic interpretation. If there would only be a finite and small number of states, it is rather easy to read off when a deterministic interpretation may be allowed, but in an infinite volume limit this is not at all very straightforward, in particular when non-local transformations are allowed, as seems to be the case in gravitational physics such as strings and black holes! Can a mapping condition be formulated?
If, in a quantum theory, in its Heisenberg formulation, a complete set of operators $`𝒪(t)`$ can be found that mutually commute at all times,
$$[𝒪(t),𝒪(t^{})]=0,(t,t^{}),$$
$`(2.9)`$
then the theory may be said to be deterministic. A set of operators is complete if any given set of eigenvalues $`O(t)`$ unambiguously specifies a basis element of Hilbert space. We may take such operators to define expectation values of ‘truly existing’ observables. In honor of J. Bell, we call these operators ‘beables’. The basis generated this way will be called a primordial basis. It could be that different complete sets of beables can be found, one set not commuting with another, so that we have different choices for the primordial basis. In that case we will have several competing ‘theories’ for the ontological facts described by our equations.
Operators $`P`$ that do not commute with the beables of the theory, such as the evolution operator $`U(t,t^{})`$ of Eq. (2.2), will be called ‘changeables’.
Quantum physicists in search of a deterministic theory have the assignment: find a complete set of beables for the Standard Model, or a modified version of it. Equivalently, find a primordial basis. In the model of an atom spinning in a magnetic field, we succeeded in doing just that. The beables are the matrices diagonal in the primordial basis $`\{|1,|2,|3\}`$ of Eq. (2.3).
3. ‘Neutrinos’
There is a more interesting system, actually realized to some extent in the real world, for which a primordial basis can be constructed.$`^\text{4}`$ Consider massless, non-interacting chiral fermions in four space-time dimensions. We can think of neutrinos, although of course real neutrinos deviate slightly from the ideal model described here.
First, take the first-quantized theory. The hamiltonian for a Dirac particle is
$$H=\stackrel{}{\alpha }\stackrel{}{p}+\beta m,\{\alpha _i,\alpha _j\}=2\delta _{ij},\{\alpha _i,\beta \}=0,\beta ^2=1.$$
$`(3.1)`$
Taking $`m=0`$, we can limit ourselves to the subspace projected out by the operator $`\frac{1}{2}(1+\gamma _5)`$, at which point the Dirac matrices become two-dimensional. The Dirac equation then reads
$$H=\stackrel{}{\sigma }\stackrel{}{p},$$
$`(3.2)`$
where $`\sigma _{1,\mathrm{\hspace{0.17em}2},\mathrm{\hspace{0.17em}3}}`$ are the Pauli matrices. We now consider the following candidates for ‘beables’:
$$\{\widehat{p},\widehat{p}\stackrel{}{\sigma },\widehat{p}\stackrel{}{x}(t)\},$$
$`(3.3)`$
where $`\widehat{p}`$ stands for $`\pm \stackrel{}{p}/|p|`$, with the sign such that $`\widehat{p}_x>0`$. We do not directly specify the sign of $`\stackrel{}{p}`$.
Writing $`p_j=i\frac{}{x_j}`$, one readily checks that these three operators commute, and that they continue to do so at all times. Indeed, the first two are constants of the motion, whereas the last one evolves into
$$\widehat{p}\stackrel{}{x}(t)=\widehat{p}\stackrel{}{x}(0)+\widehat{p}\stackrel{}{\sigma }t.$$
$`(3.4)`$
The fact that these operators form a complete set is also easy to verify: in momentum space, $`\widehat{p}`$ determines the orientation; let us take this to be the $`z`$ direction. Then, in momentum space, the absolute value of $`p`$, as well as its sign, are identified with its $`z`$-component, and it is governed by the operator $`i/p_z=x_z=\widehat{p}\stackrel{}{x}`$. The spin is defined in the $`z`$-direction by $`\widehat{p}\stackrel{}{\sigma }`$.
Mathematically, these equations appear to describe a plane, or a flat membrane, moving in orthogonal direction with the speed of light. Given the orientation (without its sign) $`\widehat{p}`$, the coordinate $`\widehat{p}\stackrel{}{x}`$ describes its distance from the origin, and the variable $`\widehat{p}\stackrel{}{\sigma }`$ specifies in which of the two possible orthogonal directions the membrane is moving. Note that, indeed, this operator flips sign under $`180^{}`$ rotations, as it is required for a spin $`\frac{1}{2}`$ representation. This, one could argue, is what a neutrino really is: a flat membrane moving in the orthogonal direction with the speed of light. But we’ll return to that later: the theory can be further improved (see the end of Sect. 4).
We do note, of course, that in the description of a single neutrino, the hamiltonian is not bounded from below, as one would require. In this very special model, there is a remedy to this, and it is precisely Dirac’s second quantization procedure. We consider a space with an infinite number of these membranes, running in all of the infinitely many possible directions $`\widehat{p}\stackrel{}{\sigma }`$. In order to get the situation under control, we introduce a temporary cut-off: in each of the infinitely many possible directions $`\widehat{p}`$, we assume that the membranes sit in a discrete lattice of possible positions. The lattice length $`a`$ may be as small as we please. Furthermore, consider a box with length $`L`$, being as large as we please. The first-quantized neutrino then has a finite number of energy levels, between $`\pi /a`$ and $`+\pi /a`$. The state we call ‘vacuum state’, has all negative energy levels filled and all positive energy levels empty. All excited states now have positive energy. Since the Dirac particles do not interact, their numbers are exactly conserved, and the collection of all observables (3.3) for all Dirac particles still correspond to mutually commuting operators.
In this very special model we thus succeed in producing a complete primordial basis, generated by operators that commute with one another at all times (beables), whereas the hamiltonian is bounded from below. We consider this to be an existence proof, but it would be more satisfying if we could have produced a less trivial model. Unfortunately, our representation of neutrinos as infinite, strictly flat membranes, appears to be impossible to generalize so as to introduce mass terms and/or interactions.
Further attempts at obtaining more realistic models which are fundamentally quantum mechanical yet allow for a deterministic interpretation failed because it did not appear to be possible to create a hamiltonian that is bounded below, so that a very special state can be selected, being the lowest eigenstate, which can be identified unambiguously as the vacuum state. In the Schrödinger picture, two classes of models may be considered: the ones with continuous time, and the ones with discrete time. If time is continuous, and a set of beables $`q_i(t)`$ is found, a natural choice for the hamiltonian would be
$$H=\underset{i}{}p_if_i(𝐪);\dot{q}_i=f_i(𝐪),$$
$`(3.5)`$
where $`p_i`$ are the ordinary momentum operators associated to $`q_i`$. We see immediately that, in spite of the quantum mechanical notation, the $`q_i`$ evolve in a deterministic manner. But we also see that, regardless the choice for the functions $`f_i`$, this hamiltonian can never be bounded below.
Time does not have to be continuous. It would suffice if a set of beables could be defined to form a sufficiently dense lattice in space-time, and this brings us to cellular automaton models. In these models, only a finite number of possible states is needed in a given spacelike volume. One might hope that then also the Hamilton density would be a finite-dimensional matrix, so that the existence of a ground state might follow, in certain cases. Again, however, there is a problem. Let now
$$q_i(t+1)=f_i\left(𝐪(t)\right),$$
$`(3.6)`$
then this defines uniquely the evolution operator over integer time steps:
$$U(t+1,t)=e^{iH},$$
$`(3.7)`$
but then this defines the eigenvalues of $`H`$ only modulo $`2\pi `$. Again, the notion of a lowest eigenstate is questionable.
In the Heisenberg picture, the dimensionality of a limit cycle does not change if we replace the time variable by one with smaller time steps, or even a continuous time. Working with a continuous time variable then has the advantage that the associated operator, the hamiltonian, is unambiguous in that case.
The neutrino example of this section does show that in some cases a bounded hamiltonian may yet exist (here, it is the continuum limit that singles out a special choice for $`H`$ with unambiguous eigenvalues), so one may hope that more is possible, but something drastically new may be needed.
4. Information loss
The new ingredient needed might be information loss$`^\text{8}`$. At first sight this is surprising. One would have thought that, with information loss, the evolution operator will no longer be unitary, and hence no quantum mechanical interpretation is allowed. Consider however a model universe with 4 elements. They will be indicated not by Dirac brackets, but as (1), (2), (3) and (4). The evolution law is as depicted in Fig. 1.
Fig. 1. Transition rule with information loss.
If we would associate basis elements of a Hilbert space to each of these states, the evolution operator would come out as
$$U(t+1,t)=\left(\begin{array}{cccc}0& 0& 1& 0\\ 1& 0& 0& 1\\ 0& 1& 0& 0\\ 0& 0& 0& 0\end{array}\right),$$
$`(4.1)`$
and this would not be a unitary operator. Of course, the reason why the operator is not unitary is that the evolution rule (2.1) is not time reversible. After a short lapse of time, only states (1), (2) and (3) can be reached. In this simple example, it is clear that one should simply erase state (4), and treat the upper left $`3\times 3`$ part of Eq. (4.1) as the unitary evolution matrix. Thus, the quantum system corresponding to the new evolution law is three-dimensional, not four-dimensional, and it seems to be trivial to restore time-reversibility and unitarity.
In more complicated non-time-reversible evolving systems, however, the ‘genuine’ quantum states and the false ones (the ones that cannot be reached from the far past) are actually quite difficult to distinguish. For this reason, we introduce the notion of equivalence classes. Two states are called equivalent if, after some finite time interval, they evolve into the same state. The system described above has three equivalence classes,
$$E_1=\{(1),(4)\},E_2=\{(2)\},E_3=\{(3)\}.$$
$`(4.2)`$
Quantum states will now be identified with these equivalence classes. In our example, in the Schrödinger picture, we have three basis elements, $`|1=E_1;|2=E_2;|3=E_3`$. In terms of these objects $`E_1,E_2,E_3`$, one has an evolution operator
$$U=\left(\begin{array}{ccc}0& 0& 1\\ 1& 0& 0\\ 0& 1& 0\end{array}\right),$$
$`(4.3)`$
and a hamiltonian operator $`H`$ can be defined such that $`U=e^{iH}`$. Our model universe would be in an eigenstate of this hamiltonian.
An extreme example of a situation where equivalence classes must occur is the black hole. We imagine a theory with classical (that is, not quantum mechanical) general relativistic field equations. Classical black holes may result as solutions. Since this classical system surely has information loss, the equivalence classes $`E_i`$ will each comprise all possible initial states that result in the same black hole (with the same mass, charge and angular momentum) after collapse. It is important to realize that the equivalence classes may be much smaller than the classes of primordial ‘states’. Their definition is not local, in the sense that two states that differ at different locations may belong to the same equivalence class. One might suspect that this could explain the apparent need for non-locality in conventional attempts at constructing hidden variable theories.
The physical distinction between theories with information loss and theories without information loss is not very clear in models with a small number of distinct states such as the example given above. After all, one may simply ignore the ‘unreachable’ states and notice that the universe ends up in a limit cycle that is indistinguishable from what a universe without information loss would do. It is the fact that the true universe is far too large ever to end up in a limit cycle, and the fact that we wish to understand small regions of this universe, nearly but not quite decoupled from the rest of the world, which make the introduction of information loss and equivalence classes non-trivial. The most acute problem to be addressed is how to create a hamiltonian that can be viewed as the integral over a hamiltonian density that is bounded below and has (more or less) local commutation rules.
What we address next is the question how to construct the laws of physics for a universe that is essentially open, i.e. it consists of smaller parts glued together. These smaller parts form an infinite 3 (space-)dimensional world. We begin with gluing two small pieces together.
If left alone, a tiny segment of the universe may be assumed to enter into a limit cycle. In fact, it may have several different cycles to choose from, depending on the initial state. In the Schrödinger picture, each of these cycles forms a sequence of a large number of states, and there are numerous different energy eigenstates that the system can choose from. But, as stated before, if time is not an extrinsically defined coordinate, it is meaningless to consider time dependence. We then prefer the Heisenberg picture, where we have only one state for each cycle. In this state it is the observables that take a periodic sequence of values, but only if time could be defined at all. In any case, the dimensionality of Hilbert space is then determined by the number of different possible cycles.
If we have two adjacent segments of the universe, the relative time coordinate is well-defined and important. It may or may not be defined as an arbitrary real number. A way to introduce coupling is as follows. We introduce a cyclic time coordinate on each segment (which could be seen as a reintroduction of the Schrödinger picture). Now assume that the relative speed of the time evolution of the two adjacent segments is determined by local dynamics. Using general relativistic notation, one would have
$$\frac{\mathrm{d}}{\mathrm{d}t}\left(|\psi _1|\psi _2\right)=\left(\sqrt{g^{00}(1)}\frac{}{t}|\psi _1\right)|\psi _2+\sqrt{g^{00}(2)}|\psi _1\frac{}{t}|\psi _2.$$
$`(4.4)`$
Here, $`\sqrt{g^{00}}`$ may be interpreted as the ‘gravitational potential field’, and the ratio $`\sqrt{g^{00}(2)}/\sqrt{g^{00}(1)}`$ is the ‘gravitational field strength’.
Take the case that each segment has just one limit cycle. Each now has a periodic ‘time’ variable $`q[\mathrm{\hspace{0.17em}0},\mathrm{\hspace{0.17em}1})`$, but we also take an external time variable $`t`$, so we have $`q_1(t)`$ and $`q_2(t)`$. Let the undisturbed time derivatives be
$$\mathrm{d}q_i/\mathrm{d}t\dot{q}_i(t)=v_i,$$
$`(4.5)`$
so that the (undisturbed) periods are $`T_i=1/v_i`$. The really relevant quantity is the ratio
$$\mathrm{d}q_1(t)/\mathrm{d}t:\mathrm{d}q_2(t)/\mathrm{d}t=\mathrm{d}q_1/\mathrm{d}q_2.$$
In Fig. 2, this is the slope of the trajectories for the solutions.
Fig. 2. Flow chart of a continuum model with two periodic variables, $`q_1`$ and $`q_2`$. In this example, there are two stable limit cycles, $`A`$ and $`B`$, representing the two ‘quantum states’ of this ‘universe’. In between, there are two orbits that would be stable in the time-reversed model.
The Schrödinger Hilbert space is spanned by the states $`|q_1,q_2`$, and our formal hamiltonian is
$$H=v_1p_1+v_2p_2;p_j=i/q_j.$$
$`(4.6)`$
In this case, even the zero-energy states span an infinite Hilbert space, so, in the Heisenberg picture, the product universe has an infinity of possible states.
Information loss is now introduced by adding a tiny perturbation that turns the flow equations into a non-Jacobian one:
$$v_1v_1^0+\epsilon f(q_1,q_2);v_2v_2^0+\epsilon g(q_1,q_2).$$
$`(4.7)`$
The effect of these extra terms can vary a lot, but in the generic case, one expects the following (assuming $`\epsilon `$ to be a sufficiently tiny number):
Let the ratio $`v_1^0/v_2^0`$ be sufficiently close to a rational number $`N_1/N_2`$. Then, at specially chosen initial conditions there may be periodic orbits, with period
$$P=v_1^0/N_1=v_2^0/N_2,$$
$`(4.8)`$
where now $`v_1^0`$ and $`v_2^0`$ have been tuned to exactly match the rational ratio — possible deviations are absorbed into the perturbation terms. Nearby these stable orbits, there are non-periodic orbits, which in general will converge into any one of the stable ones, see Fig. 2. After a sufficiently large lapse of time, we will always be in one of the stable orbits, and all information concerning the extent to which the initial state did depart from the stable orbit, is washed out. Of course, this only happens if the Jacobian of the evolution, the quantity $`_i(/q_i)\dot{q}_i`$, departs from unity. Information loss of this sort normally does not occur in ordinary particle physics, although of course it is commonplace in macroscopic physics, such as the flow of liquids with viscosity (see Sect. 6).
The stable orbits now represent our new equivalence classes (note that, under time reversal, there are other stable orbits in between the previous ones). Most importantly, we find that the equivalence classes will form a discrete set, in a model of this sort, most often just a finite set, so that, back in the Heisenberg picture, our ‘universe’ will be just in a finite number of distinct quantum states.
Generalizing this model to the case of more than two periodic degrees of freedom is straightforward. We see that, if the flow equations are allowed to be sufficiently generic (no constraints anywhere on the values of the Jacobians), then distinct stable limit orbits will arise. There is only one parameter that remains continuous, which is the global time coordinate. If we insert $`H|\psi =0`$ for the entire universe, then the global time coordinate is no longer physically meaningful, as it obtains the status of an unobservable gauge degree of freedom.
In the above models, what we call ‘quantum states’, coincides with Poincaré limit cycles of the universe. We repeat, just because our model universes are so small, we were able to identify these. When we glue tiny universes together to obtain larger and hence more interesting models, we get much longer Poincaré cycles, but also much more of them. Eventually, in practice, sooner or later, one has to abandon the hope of describing complete Poincaré cycles, and replace them by the more practical definitions of equivalence classes. At that point, when one combines mututally weakly interacting universes, the effective quantum states are just multiplied into product Hilbert spaces.
Our introduction of information loss and equivalence classes sheds new light on the neutrino model introduced in Sect. 3. In that section, we concluded that neutrinos seem to ‘be’ infinite, flat, sheets, which may eventually become untenable when space-time curvature and other interactions or mass terms are taken into account. Now, we have another option. The sheets are not the primordial beables but they are the equivalence classes. We could have that neutrinos are more-or-less conventional, but classical, point particles, with auxiliary velocity vectors $`𝐩`$ of unit length. The rules of motion are now such that the velocity in the direction of $`𝐩`$ is rigorously fixed to be $`c`$, but the velocities in the transverse directions are chaotic and not re-traceable due to information loss. In such a picture it is no longer impossible to imagine tiny deviations from the rule to incorporate interactions and tiny mass terms.
5. Harmonic and anharmonic oscillators.
What we have so far is a strategy. We still have the question how a model, either with or without information loss, can emerge in such a way that the total hamiltonian is the integral of a Hamilton density bounded from below, so that we can understand the chaotic nature of our vacuum. The neutrino model was one very special case. Now let us concentrate on the most elementary building block for bosonic fields in Nature: the harmonic oscillator, with possible disturbances.
A classical harmonic oscillator may be described by the equations
$$\dot{x}=y;\dot{y}=x.$$
$`(5.1)`$
The (Schrödinger) states are then all sets $`|x,y`$, and the classical equation (5.1) is generated by the ‘hamiltonian’
$$H=yp_xxp_y.$$
$`(5.2)`$
This, of course, is not bounded below, the price to be paid for writing a classical system quantum mechanically.
We may rewrite however,
$$\begin{array}{cc}\hfill H& =1/4(y+p_x)^21/4(yp_x)^21/4(x+p_y)^2+1/4(xp_y)^2\hfill \\ & =H_1H_2\hfill \\ \hfill H_1& =\frac{1}{2}P_1^2+\frac{1}{2}Q_1^2;H_2=\frac{1}{2}P_2^2+\frac{1}{2}Q_2^2;\hfill \end{array}$$
$`(5.3)`$
with
$$\begin{array}{cc}\hfill P_1& =\frac{1}{\sqrt{2}}(p_x+y);P_2=\frac{1}{\sqrt{2}}(x+p_y);\hfill \\ \hfill Q_1& =\frac{1}{\sqrt{2}}(xp_y);Q_2=\frac{1}{\sqrt{2}}(yp_x).\hfill \end{array}$$
$`(5.4)`$
The new variables $`P_i`$ and $`Q_i`$ obey the usual commutation rules:
$$[P_i,Q_j]=i\delta _{ij};[P_i,P_j]=[Q_i,Q_j]=0,$$
$`(5.5)`$
so that the two parts of the hamiltonian (5.3) commute: $`[H_1,H_2]=0`$. Thus, we found that although the classical harmonic oscillator has an unbounded quantum hamiltonian, we can decompose it into two genuine quantum hamiltonians, both of the familiar harmonic oscillator type and both bounded from below, but with a minus sign in between. All we now have to do is to ‘postulate’ that, say, $`H_2|\psi =\frac{1}{2}\mathrm{}\omega `$ is imposed as a constraint equation on all states. Then we have a quantum harmonic oscillator.
This, however, is not the true solution to our problem. The splitting (5.3) is far too arbitrary. Where does the constraint $`H_2=\frac{1}{2}`$ come from, and how can it survive interactions? How do we couple two or more such oscillators without spoiling the constraint? Let us analyze the reason why the splitting (5.3) is possible.
The classical harmonic oscillator has two conserved operators, besides the hamiltonian itself. These are: — the radius $`r`$ of the orbit: $`\varrho ^2=x^2+y^2`$, and — the dilatation operator $`D=xp_x+yp_yi`$; classically, after all, the periodic solutions of the oscillator behave identically after a scale transformation $`x,y\lambda x,\lambda y;p_x,p_y\lambda ^1p_x,\lambda ^1p_y`$. We easily check that
$$[H,x^2+y^2]=0;[H,D]=0.$$
$`(5.6)`$
It is due to these operators that we can construct $`H_1`$ and $`H_2`$:
$$\begin{array}{cc}\hfill H_1& =\frac{1}{4\varrho ^2}(\varrho ^2+H)^2+\frac{1}{4\varrho ^2}(D+i)^2;\hfill \\ \hfill H_2& =\frac{1}{4\varrho ^2}(\varrho ^2H)^2+\frac{1}{4\varrho ^2}(D+i)^2.\hfill \end{array}$$
$`(5.7)`$
We see immediately that $`H_1`$ and $`H_2`$ commute, since $`D`$ and $`\varrho `$ commute with $`H`$ (they do not commute with each other, but that does no harm), and we see that $`H=H_1H_2`$.
But we also see that, in this notation, the construction is fairly arbitrary. The contribution of the dilatation operator is unnecessary. We may perform the more interesting splitting:
$$\begin{array}{cc}\hfill \text{if}[\varrho ,H]=0& \text{then}H=H_1H_2\hfill \\ & H_1=\frac{1}{4\varrho ^2}(\varrho ^2+H)^2;\hfill \\ & H_2=\frac{1}{4\varrho ^2}(\varrho ^2H)^2.\hfill \end{array}$$
$`(5.8)`$
Now this generalizes to any function $`\varrho ^2`$ of the coordinates $`x`$ and $`y`$ that happens to be conserved, $`[\varrho ^2,H]=0`$. So, in general, let
$$H=𝐩𝐟(𝐪);[H,\varrho ^2(𝐪)]=0,$$
$`(5.9)`$
and
$$\begin{array}{cc}\hfill H=H_1H_2& ;\hfill \\ \hfill H_{1,2}=\frac{1}{4\varrho ^2}(\varrho ^2\pm H)^2& ;[H_1,H_2]=0.\hfill \end{array}$$
$`(5.10)`$
Then introduce as a constraint:
$$H_2|\psi 0,$$
$`(5.11)`$
this implies that
$$HH_1\varrho ^20.$$
$`(5.12)`$
Fig. 3. Stable orbits, a) for the harmonic oscillator, b) an anhormonic oscillator. After switching on a dissipative term, the regions in between these trajectories will have only non-periodic solutions, tending towards the stable attractors.
The subclass of states obeying the constraint (5.11) obeys the Schrödinger equation $`\frac{\mathrm{d}}{\mathrm{d}t}|\psi =iH_1|\psi `$, with a hamiltonian $`H_1`$ bounded from below. But now we can motivate the constraint by introducing information loss! This goes as follows. Let us assume that, in a deterministic system, there are stable orbits, separated by regions where solutions are non-periodic, attracted to the stable attractors. First, we mimic the system in terms of a non-dissipative model, where all orbits would be stable. It is then described by Eq. (5.9). Let the periods of these orbits be functions $`T(\varrho )`$ of $`\varrho `$. The constraint (5.11), (5.12) implies $`H=\varrho ^2`$, or
$$e^{iHT}|\psi =|\psi .$$
$`(5.13)`$
These are the states for which the periods $`T(\varrho )`$ obey
$$\varrho ^2T(\varrho )=2\pi n,nZ.$$
$`(5.14)`$
Thus, the constraint appears to correspond to limiting oneself to the stable orbits only. Note that, with Eq. (5.14), the hamiltonian can obtain any kind of eigenvalue spectrum, as opposed to the equidistant lines of the harmonic oscillator.
6. Conclusions In this lecture we investigated classical, deterministic, dissipative models, and we found that, in general, they develop distinct stable orbits. The mathematics for analyzing these models requires that we first introduce non-dissipative equations, which allow a formalism using quantum mechanical notation, but, without dissipation, it cannot be understood why the hamiltonian would be bounded from below. Then we find that dissipation imposes constraints on the solutions, which appear to provide bounded hamiltonians. It is remarkable that dissipation also leads to an apparent quantization of the orbits, and this quantization indeed resembles the quantum structure seen in the real world.
The next step, yet to be taken, is to couple infinite numbers of dissipating oscillators to form models of quantum field theories. This may appear to be a very difficult task, but we do notice that in classical general relativity black hole formation is inevitable, and black holes indeed absorb information. This would imply that the distance scale at which dissipation plays a role must be the Planck scale.
At scales between the Standard Model and the Planck scale, the introduction of dissipation would be a new approach. Perhaps we can find models resembling Navier-Stokes liquids with viscosity. In non-relativistic models, the dimensionality of a viscosity $`\eta `$ is given by
$$[\eta /\varrho ]=[\mathrm{cm}^2/\mathrm{sec}],$$
$`(6.1)`$
where $`\varrho `$ is the mass density of the fluid. In a relativistic theory, where there is a fixed unit of velocity $`c=1`$, the cm and the sec have the same dimensionality, so now $`\eta /\varrho `$ has the dimensionality of a length. It is tempting to take this to be the Planck length. We may take strictly continuous fields, which however at distance scales $`\mathrm{}`$ small compared to the Planck scale are totally controlled by viscosity. There, at small Reynolds number$`^\text{9}`$,
$$R=\varrho u\mathrm{}/\eta ,$$
$`(6.2)`$
where $`u`$ are the typical velocities, the field distributions show no further structure, but at distance scales large compared to the Planck scale, one may expect ‘turbulence’, i.e., chaotic behaviour, for which we propose the introduction of apparently quantum mechanical techniques in order to describe the statistics. One then may invoke the renormalization group in order to reach the length scales of the Standard Model, and make contact with the real world.
It is far too early to ask for tangible results and firm testable predictions of the approach that we have in mind. A very indirect prediction may perhaps be made. We conjecture that the apparently quantum mechanical nature of our world is due to the statistics of fluctuations that occur at the Planck scale, in terms of a regime of completely deterministic dynamics. This would entail that all quantum mechanical effects should be reproducible in some deterministic model, including all machinations with what is usually called a ‘quantum computer’. As is well-known, quantum computers, if they can be constructed, will be able to do computations no ordinary computer can accomplish. This would be a contradiction with our claim that it can be mimicked using ordinary computers. However, we are unable to mimic any quantum system we like. Interactions are essential and unavoidable. It is also these interactions that cause unwanted decoherence in a quantum computer. Experimenters are trying to create devices in which the ideal situation is approached as well as possible. I now claim that it will be impossible to shield these devices from the unwanted interactions, so that the ideal quantum computer can never be built. More precisely:
No quantum computer can ever be built that can outperform a classical computer if the latter would have its components and processing speed scaled to Planck units.
Because the Planck units are extremely tiny, this still leaves lots of room for quantum computers to do miracles, but eventually, there will be a limit. An ideal quantum computer that would consist of millions of parts, could, in principle, factor integers with millions of digits into prime numbers. It appears that present programs can factor a number of $`N`$ digits using memories of the order of $`10^{\sqrt{N\mathrm{log}N}}`$ cells, in $`10^{\sqrt{N\mathrm{log}N}}`$ steps. Perhaps a reasonable computer takes $`10^{120}`$ Planck volumes; that would limit the factorizable numbers to $`10^{4000}`$ or so. Thus we predict that even a quantum computer will not be able to exceed such limits in practice.
References
1. G. ’t Hooft, The Holographic Principle, Opening Lecture, Erice, August 1999. See also: G. ’t Hooft, Dimensional reduction in quantum gravity. In Salamfestschrift: a collection of talks, World Scientific Series in 20th Century Physics, vol. 4, Eds. A. Ali, J. Ellis and S. Randjbar-Daemi (World Scientific, 1993), THU-93/26, gr-qc/9310026; Black holes and the dimensionality of space-time, in Proceedings of the Symposium “The Oskar Klein Centenary”, 19-21 Sept. 1994, Stockholm, Sweden. Ed. U. Lindström, World Scientific 1995, p. 122; L. Susskind, J. Math. Phys. 36 (1995) 6377, hep-th/9409089.
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# 1 Introduction
## 1 Introduction
In a previous Letter it has been argued that a non-trivial $`|\theta ^{ind}`$ -vacua may be formed in heavy ion collisions<sup>1</sup><sup>1</sup>1 In what follows we often omit the label “ind” for the induced $`\theta `$ vacuum state. We hope this state $`|\theta ^{ind}`$ will not be confused with $`|\theta ^{fundamental}`$ which is zero in our world and which can not be changed in QCD. The simplest way to visualize $`|\theta ^{ind}`$ is to assume that right after the QCD phase transition the flavor singlet phase of the chiral condensate is non-zero in a macroscopically large domain. This phase is identified with $`\theta ^{ind}`$. This identification is a direct consequence of the transformation properties of the fundamental QCD lagrangian under $`U(1)_A`$ rotations by which the chiral singlet phase can be rotated away at the cost of the appearence of the induced $`\theta ^{ind}`$ ( see Sections 2.5 and 3.2 for detail explanations of the adopted terminology. See also ).. As is known, in the infinite volume limit and in thermal equilibrium the $`|\theta `$ vacuum state is the absolutely stable ground state of a new world with new physics quite different from ours. In particular, P and CP symmetries are strongly violated in this world. In spite of the fact that the $`|\theta `$ vacuum state has a higher energy (see below) this state is stable due to the superselection rule: There is no any gauge invariant observable $`𝒜`$ in QCD, which would communicate between two different worlds, $`\theta ^{}|𝒜|\theta \delta (\theta \theta ^{})`$. Therefore, there are no transitions between these worlds. The problem of why $`\theta =0`$ in our world is known as the strong CP problem. The best solution of this problem, which has passed the test of time, is the introduction of the Peccei-Quinn symmetry and the corresponding pseudo-goldstone particle known as the axion, -,. The axion solution of the strong CP problem suggests that the $`\theta `$ parameter of QCD is promoted to the dynamical axion field $`\theta \frac{a(x)}{f_a}`$ with very small coupling constant, $`1/f_a1`$. Once $`\theta `$ becomes the dynamical field, it automatically relaxes to $`\theta =0`$ as the lowest energy state.
Of course, we do not expect that such a stable $`|\theta `$ state can be produced in heavy ion collisions: However, we do expect that locally, for a short period of time due to the non-equilibrium dynamics after the QCD phase transition, a large domain with wrong $`|\theta `$ direction may be formed. This provides us with a unique opportunity to study a new state of matter. Numerical calculations of the evolution of the chiral fields after the QCD phase transition support this expectation. An obvious question is how this $`|\theta `$ vacuum state can be observed? To be more concrete, the question we are addressing in this letter is formulated in the following way: “What kind of observable should be measured at RHIC to confirm that a new ground state (and, therefore, a new set of excitations accompanying this ground state) has been created? ”. The most natural answer to this question is that: If the axions do exist, they can be produced during the relaxation of $`|\theta 0^{ind}`$ to the lowest possible state with $`\theta =0`$ as suggested in ,. However, the axion production rate under the conditions which can be achieved at RHIC will be too low in comparison with the limit already achieved from the astrophysical and cosmological considerations. The second obvious choice for the observable which can be a signal of creation of the $`\theta `$ vacuum state is some CP odd correlation as recently suggested in <sup>2</sup><sup>2</sup>2In the original papers, the authors discussed not the induced $`|\theta ^{ind}`$ vacua which is the subject of the present paper, but rather some metastable (even in the infinite volume limit) vacua which may (or may not) exist within the given QCD parameters $`m_i,N_c,N_f,0|\overline{\mathrm{\Psi }}_i\mathrm{\Psi }_i|0`$ as we know them. It is not the goal of this paper to comment on whether such metastable states are likely /unlikely to exist in Nature; or likely /unlikely to be produced at RHIC. Rather we want to emphasize that the CP odd observables suggested in are sensitive to any $`CP`$ odd physics, not necessary related to the metastable states introduced in.. However, as Roberto Peccei noticed recently<sup>3</sup><sup>3</sup>3 I am thankful to Roberto for the discussions on this subject, which actually motivated the present search for an alternative signal on the decay products in the $`|\theta ^{ind}`$ vacua background. the signal could be considerably washed out by re-scattering of the pions and their interactions in the final states. Indeed, as is well-known (see e.g. reviews), the final state interactions will produce complex phases which mimic true CP-odd effects. In practice it is quite difficult to overcome the difficulty of separating a true CP violation from its simulation due to the final state interactions. For example, it is well known that the final state interactions is the most difficult obstacle in the study of CP violation in $`K^0(\overline{K^0})\pi ^\pm \mu ^{}\nu `$ decays by measuring the muon transverse polarization $`\xi _i`$ (defined by the correlation $`(\stackrel{}{p_\pi }\times \stackrel{}{p_\mu }\stackrel{}{\xi })`$). Fortunately, the final state interactions in the charged kaon decays $`K^+(K^{})\pi ^0\mu ^\pm \nu `$ are strongly suppressed, and therefore a measurement of this correlation in the ongoing experiment at KEK on the level $`>10^6`$ would imply a new source of CP violation. In general, one should expect that similar mimicry of true CP violation due to the final state interactions also occurs in the heavy ion collisions. The main goal of this letter is to offer some new observable which:
a) can be measured at RHIC;
b) will be the signal of the new physics of the $`\theta `$ vacua;
c) does not suffer from the deficiencies mentioned above.
Before I explain the main idea of the construction of such an observable, I would like to explain the low-energy physics in the $`\theta `$ world.
## 2 Low Energy Physics of the $`\theta `$ World
In this section (except subsection 2.5) we study the phenomenology of the world when $`\theta ^{fund}0`$. It will give us some useful hints about CP odd physics in the unusual environment with $`\theta 0`$.
The starting point of our analysis is the low energy effective lagrangian which reproduces the anomalous conformal and chiral Ward Identities. The corresponding construction in the large $`N_c`$ limit has been known for a long time,. The generalization of the construction of ref., for finite $`N_c`$ was given in ,, and we shall use formulae from the papers ,. However, we should remark at the very beginning that all local properties of the effective lagrangians for finite and infinite $`N_c`$ are very much the same. Small quantitative differences in local physics between the description of , on the one hand and the description of, on the other hand, do not alter the qualitative results which follow. First I discuss the properties of the unusual $`\theta `$ vacua itself. Discussions of the properties of the Goldstone bosons in the $`\theta `$ vacuum background will follow after that.
### 2.1 Effective chiral lagrangian and properties of $`\theta `$ vacua
In the effective Lagrangian approach, the light matter fields are described by the unitary matrix $`U_{ij}`$ corresponding to the phases of the chiral condensate $`0|\overline{\mathrm{\Psi }}_L^i\mathrm{\Psi }_R^j|0=|0|\overline{\mathrm{\Psi }}_L\mathrm{\Psi }_R|0|U_{ij}`$. In terms of $`U`$ the effective potential for arbitrary $`\theta `$ takes the form ,:
$$V(U,\theta )=E\mathrm{cos}\left[\frac{1}{N_c}(\theta i\mathrm{log}DetU)\right]\frac{1}{2}Tr(MU+M^+U^+),$$
(1)
where $`M`$ is the diagonal mass matrix for quark masses defined with their condensates ($`M_i=diag(m_i\overline{\mathrm{\Psi }}_i\mathrm{\Psi }_i)`$ ) and $`E=b\alpha _s/(32\pi )G^2`$, with $`b=11N_c/32N_f/3`$, is the vacuum energy (“cosmological” term ) required by the conformal anomaly. Expanding the cosine (this corresponds to the expansion in $`1/N_c`$) we recover the result of at lowest order in $`1/N_c`$ together with the constant term $`E`$:
$$V(U,\theta ,N_c\mathrm{})==E\frac{\nu ^2_{YM}}{2}(\theta i\mathrm{log}DetU)^2\frac{1}{2}Tr(MU+M^+U^+)+\mathrm{},$$
(2)
where we used the fact that in the large $`N_c`$ limit $`E(\frac{1}{N_c})^2=\nu ^2_{YM}`$ where $`\nu ^2_{YM}<0`$ is the topological susceptibility in pure YM theory. Corrections in $`1/N_c`$ stemming from Eq.(1) constitute a new result of ref.. First of all, let me review the picture of the vacuum structure for $`\theta 0`$ stemming from the effective potential (1). In the next subsection I review some properties of the Goldstone modes in the $`\theta 0`$ background. To study the vacuum properties it is convenient to parametrize the fields $`U`$ as $`U=diag(\mathrm{exp}i\varphi _q)`$ such that the potential (1) takes the form:
$$V=E\mathrm{cos}\left(\frac{1}{N_c}\theta \frac{1}{N_c}\varphi _i\right)M_i\mathrm{cos}\varphi _i,$$
(3)
The minimum of this potential is determined by the following equation:
$$\frac{1}{N_c}\mathrm{sin}\left(\frac{1}{N_c}\theta \frac{1}{N_c}\varphi _i\right)=\frac{M_i}{E}\mathrm{sin}\varphi _i,i=1,\mathrm{},N_f.$$
(4)
At lowest order in $`1/N_c`$ this equation coincides with that of . For general values of $`M_i/E`$, it is not possible to solve Eq.(4) analytically. However, in the realistic case $`\epsilon _u,\epsilon _d1,\epsilon _s1`$ where $`\epsilon _i=\frac{N_cM_i}{E}`$, the approximate solution can be found:
$`\mathrm{sin}\varphi _u`$ $`=`$ $`{\displaystyle \frac{m_d\mathrm{sin}\theta }{[m_u^2+m_d^2+2m_um_d\mathrm{cos}\theta ]^{1/2}}}+O(\epsilon _u,\epsilon _d),`$
$`\mathrm{sin}\varphi _d`$ $`=`$ $`{\displaystyle \frac{m_u\mathrm{sin}\theta }{[m_u^2+m_d^2+2m_um_d\mathrm{cos}\theta ]^{1/2}}}+O(\epsilon _u,\epsilon _d),`$ (5)
$`\mathrm{sin}\varphi _s`$ $`=`$ $`O(\epsilon _u,\epsilon _d).`$
This solution coincides with the one of Ref. to leading order in $`\epsilon _u,\epsilon _d`$. In what follows for the numerical estimates and for simplicity we shall use the $`SU(2)`$ limit $`m_u=m_dm_s`$ where the solution (2.1) can be approximated as:
$`\varphi _u{\displaystyle \frac{\theta }{2}},\varphi _d{\displaystyle \frac{\theta }{2}},\varphi _s0,0\theta <\pi `$
$`\varphi _u{\displaystyle \frac{\theta +2\pi }{2}},\varphi _d{\displaystyle \frac{\theta 2\pi }{2}},\varphi _s0,\pi \theta <2\pi .`$ (6)
Once solution (2.1) is known, one can calculate the vacuum energy and topological charge density $`Q=0|\frac{\alpha _s}{8\pi }G\stackrel{~}{G}|0`$ as a function of $`\theta `$. In the limit $`m_u=m_dm,\overline{d}d=\overline{u}u\overline{q}q`$ one has:
$`V_{vac}(\theta )V_{vac}(\theta =0)2m|\overline{q}q|(|\mathrm{cos}{\displaystyle \frac{\theta }{2}}|1)`$
$`\theta |{\displaystyle \frac{\alpha _s}{8\pi }}G\stackrel{~}{G}|\theta ={\displaystyle \frac{V_{vac}(\theta )}{\theta }}=m|0|\overline{q}q|0|\mathrm{sin}{\displaystyle \frac{\theta }{2}},`$ (7)
As expected, the $`\theta `$ dependence appears only in combination with $`m`$ and goes away in the chiral limit. One can also calculate the chiral condensate $`\overline{\mathrm{\Psi }}_L^i\mathrm{\Psi }_R^i`$ in the $`\theta `$ vacua using solution (2.1) for vacuum phases:
$$\theta |\overline{q}q|\theta =\mathrm{cos}\frac{\theta }{2}0|\overline{q}q|0_{\theta =0},\theta |\overline{q}i\gamma _5q|\theta =\mathrm{sin}\frac{\theta }{2}0|\overline{q}q|0_{\theta =0}$$
(8)
Both expressions (2.1) and (8) explicitly demonstrate that P, CP symmetries are explicitly broken in the $`\theta `$ vacuum background. This has very important phenomenological consequences which will be discussed in the next section. The crucial point for the analysis presented above is the fact that the $`\theta `$ parameter appears in the effective lagrangian (1) only in the combination, $`(\theta i\mathrm{log}DetU)`$. This is a direct consequence of the transformation properties of the chiral fields under $`U(1)_A`$ rotations. To convince the reader that (1) does indeed represent the anomalous effective low energy Lagrangian, three of its most salient features are listed below :
i) Eq. (1) correctly reproduces the Witten- Di Vecchia -Veneziano effective chiral Lagrangian in the large $`N_c`$ limit;
ii) it reproduces the anomalous conformal and chiral Ward identities of $`QCD`$;
iii) it reproduces the known dependence in $`\theta `$ up to small corrections. Accordingly, it leads to the correct $`2\pi `$ periodicity of observables.
### 2.2 Goldstone bosons in $`\theta `$ world
In this subsection I would like to describe some unusual properties of the Goldstone modes in the unusual $`\theta `$ environment. I refer the interested readers to the original papers , for a detailed description of the approach. Here I present some old/new results which can be easily derived from the general expressions given in ref.. To study the properties of the pseudo-goldstone bosons, we parametrize the chiral matrix (1) in the form
$$U=U_0\mathrm{exp}\left[i\sqrt{2}\frac{\pi ^a\lambda ^a}{f_\pi }+i\frac{2}{\sqrt{N_f}}\frac{\eta ^{}}{f_\pi }\right],$$
(9)
where $`\lambda ^a`$ are Gell-Mann matrices of $`SU(N_f=3)`$, $`\pi ^a`$ is the octet of pseudo-goldstone bosons, and $`f_\eta ^{}f_\pi 133MeV`$. In this formula $`U_0`$ solves the minimization equations for the effective potential (1), and the fields $`\pi ^a,\eta ^{}`$ all have vanishing vacuum expectation values. In what follows we shall use the approximate solution (2.1) for the vacuum phases $`\varphi _i`$. A simple calculation of the mass matrix in the $`SU(2)`$ limit when $`m_u=m_dm`$ yields the following result for an arbitrary value of $`\theta `$ (for the general case $`m_um_d`$ see):
$`m_\pi ^2={\displaystyle \frac{4}{f_\pi ^2}}(M_q|\mathrm{cos}{\displaystyle \frac{\theta }{2}}|),m_\eta ^2={\displaystyle \frac{4}{3f_\pi ^2}}(M_q|\mathrm{cos}{\displaystyle \frac{\theta }{2}}|+2M_s)`$
$`m_\eta ^{}^2=4{\displaystyle \frac{EN_f}{f_\pi ^2N_c^2}}+{\displaystyle \frac{4}{N_f}}{\displaystyle \frac{1}{f_\pi ^2}}\left(2M_q|\mathrm{cos}{\displaystyle \frac{\theta }{2}}|+M_s\right),`$ (10)
$`m_{\eta ,\eta ^{}}^2=m_{\eta ^{},\eta }^2=2\sqrt{{\displaystyle \frac{2}{3N_f}}}{\displaystyle \frac{1}{f_\pi ^2}}(2M_q|\mathrm{cos}{\displaystyle \frac{\theta }{2}}|2M_s))`$
In this isospin limit $`m_u=m_dm`$, the $`\pi ^0`$ does not mix with $`\eta `$ or $`\eta ^{}`$, and $`m_{\pi ^0}`$ coincides with the physical mass of the $`\pi ^0`$. At the same time $`\eta `$ and $`\eta ^{}`$ do mix with each other. The mixing is determined by the following matrix:
$`m_{\eta \eta ^{}}^2=\left(\begin{array}{cc}\frac{4}{3f_\pi ^2}(2m_s\overline{s}s+m\overline{q}q|\mathrm{cos}\frac{\theta }{2}|)& \frac{4\sqrt{2}}{3f_\pi ^2}(m_s\overline{s}sm\overline{q}q|\mathrm{cos}\frac{\theta }{2}|)\hfill \\ \frac{4\sqrt{2}}{3f_\pi ^2}(m_s\overline{s}sm\overline{q}q|\mathrm{cos}\frac{\theta }{2}|)& \frac{4}{3f_\pi ^2}(m_s\overline{s}s+2m\overline{q}q|\mathrm{cos}\frac{\theta }{2}|)+4\frac{EN_f}{f_\pi ^2N_c^2}\hfill \end{array}\right)`$ (13)
which coincides with an accuracy $`O(m^2)`$ with the matrix given by Veneziano at $`\theta =0`$
$`m_{\eta \eta ^{}}^2=\left(\begin{array}{cc}\frac{1}{3}(4m_K^2m_\pi ^2)& \frac{2\sqrt{2}}{3}(m_K^2m_\pi ^2)\hfill \\ \frac{2\sqrt{2}}{3}(m_K^2m_\pi ^2)& \frac{2}{3}m_K^2+\frac{1}{3}m_\pi ^2+\frac{\chi }{N_c}\hfill \end{array}\right).`$ (16)
The only difference between eq.(13) and Eq.(16) is that the topological susceptibility $`\chi `$ in pure YM theory in the latter is substituted by the term proportional to the gluon condensate in real QCD in the former.
A few remarks are in order. First, from Eq.(2.2) one could naively think that the $`\pi `$ meson becomes massless at $`\theta =\pi `$. This is a consequence of our approximation $`m_u=m_d`$. For $`m_um_d`$ the $`\pi `$ meson never becomes massless, and in the vicinity $`\theta \pi `$ one should use the exact formula $`m_\pi ^2=\frac{2}{f_\pi ^2}(M_u\mathrm{cos}(\varphi _u)+M_d\mathrm{cos}(\varphi _d))`$ with exact solutions (2.1) for phases $`\varphi _u,\varphi _d`$ ( see). However, the observation that the masses of the pseudo-goldstone bosons get corrections of the order of $`1`$ in the $`\theta 0`$ background is absolutely correct. In particular, the $`\pi `$ meson could easily have a mass of, let us say, $`95MeV`$ for a certain value of $`\theta `$ rather than $`135MeV`$ from the Particle Data Booklet. Masses of all other hadrons are also influenced by the $`\theta `$ vacua. However, this influence is not as profound as for pseudo-goldstone bosons because all effects due to the $`\theta 0`$ background are proportional to $`m`$ and must vanish in the chiral limit. For pseuodo-Goldstone bosons and the $`\eta ^{}`$ meson this effect can be calculated exactly in the limit $`m0`$:
$`{\displaystyle \frac{m_\eta ^{}^2(\theta =0)m_\eta ^{}^2(\theta 0)}{m_\eta ^{}^2(\theta =0)}}={\displaystyle \frac{8}{3f_\pi ^2m_\eta ^{}^2}}m|\overline{q}q|\left(1|\mathrm{cos}{\displaystyle \frac{\theta }{2}}|\right)>0`$
$`{\displaystyle \frac{m_\eta ^2(\theta =0)m_\eta ^2(\theta 0)}{m_\eta ^2(\theta =0)}}={\displaystyle \frac{4}{3f_\pi ^2m_\eta ^2}}m|\overline{q}q|\left(1|\mathrm{cos}{\displaystyle \frac{\theta }{2}}|\right)>0`$ (17)
$`{\displaystyle \frac{m_K^2(\theta =0)m_K^2(\theta 0)}{m_K^2(\theta =0)}}={\displaystyle \frac{2}{f_\pi ^2m_K^2}}m|\overline{q}q|\left(1|\mathrm{cos}{\displaystyle \frac{\theta }{2}}|\right)>0`$
For all other heavy hadrons one can not present analogous calculations; however, one should expect a similar, quite substantial effect
$$\frac{m\overline{q}q}{f_\pi ^2}(100MeV)^2$$
for all hadrons.
My first remark was about masses (spectrum), see discussions above; my second remark is about quantum numbers. Namely, I would like to argue that all hadrons in the $`\theta `$-vacuum environment cease to carry definite P, CP parities but instead become mixtures of the states with different quantum numbers. This remark becomes quite obvious if one remembers that the $`\theta `$ vacuum ground state with the chiral and gluon condensates given by formulae (8) and (2.1) correspondingly, is not invariant under these symmetries. Therefore, in general one should expect that all excitations in this background are not eigenstates of P and CP parities. Therefore: in the presence of a non-zero angle $`\theta `$, the pseudo-goldstone bosons cease to be the pure pseudoscalars and acquire scalar components. This mixing of states with different parities has been known since where it was explicitly demonstrated that the charmonium levels with quantum numbers $`0^+`$ and $`0^{}`$ do mix. In this case the calculations are under theoretical control (due to the large mass $`m_c`$ of the charmed quark), and the result for the mixing is expressed in terms of the CP-odd gluon condensate $`\theta |\frac{\alpha _s}{4\pi }G\stackrel{~}{G}|\theta =2m|0|\overline{q}q|0|\mathrm{sin}\frac{\theta }{2}`$, see (2.1). It was also demonstrated in ref. that the $`\eta \pi \pi `$ decay is allowed in the $`\theta `$ vacuum background. This decay also can be interpreted as a result of mixing pseudo-Goldstone bosons $`0^{}`$ with a scalar particle $`0^+`$ which can easily decay to two $`\pi `$ mesons. Due to the phenomenological importance of this result as a definite signature of the produced $`|\theta ^{ind}`$ state in heavy ion collisions, we derive the decay rate for $`\mathrm{\Gamma }(\eta \pi \pi )`$ in a separate subsection below. We express the rate $`\mathrm{\Gamma }(\eta \pi \pi )`$ exclusively in terms of the chiral vacuum condensate (8). At small $`\theta `$ our formula coincides with the one presented long ago in ref..
### 2.3 $`\eta \pi \pi `$ decay
As is known, this decay is strongly suppressed in our world with $`\theta =0`$. Indeed, CP parity for the initial state is $`()`$, while for the final state it is $`(+)`$. This is the main reason for a very small full width $`\mathrm{\Gamma }1.2KeV`$ for $`\eta `$ meson which can decay strongly only to $`3\pi `$ mesons with a strong isospin suppression (three $`\pi `$ mesons can not be in the $`I=0`$ state) as well as a phase volume suppression. The starting point of our calculations is the following matrix element:
$$A(\eta \pi \pi )=2\pi |m_u\overline{u}um_d\overline{d}d|\eta .$$
(18)
Our goal is to estimate this matrix element in the limit of vanishing $`4`$-momenta of the $`\pi `$ and $`\eta `$ mesons. Therefore, one can use the standard PCAC-technique. The result is:
$$A(\eta \pi \pi )=2\pi |m_u\overline{u}um_d\overline{d}d|\eta \frac{2}{\sqrt{6}}\frac{(m_u+m_d)\theta |\overline{u}i\gamma _5u+\overline{d}i\gamma _5d|\theta }{f_\pi ^3}$$
$$\frac{8}{\sqrt{6}}\frac{m\mathrm{sin}\frac{\theta }{2}0|\overline{q}q|0_{\theta =0}}{f_\pi ^3},$$
(19)
where at the last step we used formula (8) for the CP odd chiral condensate which is not zero in the $`\theta `$ world. In this derivation we neglect all corrections proportional to $`m_um_d`$ having in mind that at the very end we shall use the approximate solution (2.1) for the chiral phases and masses (2.2) rather than the exact expression (2.1). Now we are ready to estimate the decay rate:
$$\mathrm{\Gamma }(\eta \pi \pi )=\frac{|A(\eta \pi \pi )|^2}{16\pi m_\eta (\theta )}\left(1\frac{4m_\pi ^2(\theta )}{m_\eta ^2(\theta )}\right)^{\frac{1}{2}}$$
$$\frac{2}{3\pi m_\eta (\theta )}\left(\frac{m\mathrm{sin}\frac{\theta }{2}0|\overline{q}q|0}{f_\pi ^3}\right)^2\left(1\frac{4m_\pi ^2(\theta )}{m_\eta ^2(\theta )}\right)^{\frac{1}{2}},$$
(20)
with all masses to be calculated in the $`\theta `$ vacua (2.2). Formula (20) reduces to the corresponding expression of ref. in the limit $`\theta 0`$ when the standard PCAC relation $`m|0|\overline{q}q|0|\frac{1}{4}m_\pi ^2f_\pi ^2`$ is substituted into eq.(20). For numerical estimates we neglect the $`\theta `$ dependence of masses and arrive at the following decay rate:
$$\mathrm{\Gamma }(\eta \pi \pi )0.5MeV(\mathrm{sin}\frac{\theta }{2})^2,$$
(21)
which essentially determines the full width of the $`\eta `$ meson in the $`\theta 0`$ world.
### 2.4 $`\eta ^{}\pi \pi `$ decay
One can carry out similar calculations for the $`\eta ^{}\pi \pi `$ decay as well. This decay is also strongly suppressed for the same reasons discussed above for the $`\eta \pi \pi `$ decay. The full width of $`\eta ^{}`$ meson (from particle data booklet), $`\mathrm{\Gamma }(\eta ^{})0.2MeV`$, is very small in comparison with the “normal” width $`(100MeV)`$ which one could expect for a strongly interacting particle with mass $`1GeV`$. This suppression is due to the small phase volume of the only allowed strong decay, $`\eta ^{}\eta \pi \pi `$.
The starting point of our calculations of $`\mathrm{\Gamma }(\eta ^{}\pi \pi )`$ with $`\theta 0`$ is the following matrix element
$$A(\eta ^{}\pi \pi )=2\pi |m_u\overline{u}um_d\overline{d}d|\eta ^{}.$$
(22)
Our goal is to estimate this matrix element in the limit of vanishing $`4`$-momenta of the $`\pi `$ mesons. We use the standard PCAC-technique as we did above for $`\eta \pi \pi `$ decay. The result of this calculation can be presented in terms of the following $`\eta ^{}`$ matrix element:
$$A(\eta ^{}\pi \pi )=2\pi |m_u\overline{u}um_d\overline{d}d|\eta ^{}\frac{(m_u+m_d)0|\overline{u}u+\overline{d}d|\eta ^{}_\theta }{f_\pi ^2}.$$
(23)
In our world with $`\theta =0`$ this matrix element is obviously zero. However, as discussed above, due to the mixing of states with different parities in the $`\theta 0`$ background this matrix element is expected to be nonzero in the $`\theta 0`$ vacuum state<sup>4</sup><sup>4</sup>4 it must vanish in $`\theta 0`$ limit though. .
The simplest way to evaluate the $`\eta ^{}`$ matrix element (23) is to make use of our knowledge of the vacuum condensates (2.1, 8)in the $`\theta 0`$ world. One can differentiate $`\theta |\frac{\alpha _s}{8\pi }G\stackrel{~}{G}|\theta `$ with respect to $`m`$ in order to derive the following low energy theorem which will be used for the estimation of the matrix element (23) in what follows.
$`({\displaystyle \frac{}{m_u}}+{\displaystyle \frac{}{m_d}})\theta |{\displaystyle \frac{\alpha _s}{8\pi }}G\stackrel{~}{G}|\theta =i\underset{q0}{lim}{\displaystyle 𝑑xe^{iqx}\theta |(\overline{u}u(x)+\overline{d}d(x)),\frac{\alpha _s}{8\pi }G\stackrel{~}{G}(0)|\theta }=`$
$`={\displaystyle \frac{1}{2}}0|\overline{u}u+\overline{d}d|0\mathrm{sin}{\displaystyle \frac{\theta }{2}},`$ (24)
where at the last step we used our knowledge of the $`\theta `$ dependence for the vacuum condensate $`\theta |\frac{\alpha _s}{8\pi }G\stackrel{~}{G}|\theta =m0|\overline{q}q|0\mathrm{sin}\frac{\theta }{2}`$, see eq. (2.1). The same relation can be obtained in a different way, by differentiating $`\theta |\overline{q}q|\theta `$ with respect to $`\theta `$. Indeed,
$`{\displaystyle \frac{}{\theta }}\theta |\overline{u}u+\overline{d}d|\theta =i\underset{q0}{lim}{\displaystyle 𝑑xe^{iqx}\theta |(\overline{u}u(0)+\overline{d}d(0)),\frac{\alpha _s}{8\pi }G\stackrel{~}{G}(x)|\theta }=`$
$`={\displaystyle \frac{1}{2}}0|\overline{u}u+\overline{d}d|0\mathrm{sin}{\displaystyle \frac{\theta }{2}},`$ (25)
where at the last step we used our knowledge of the $`\theta `$ dependence for the vacuum condensate $`\theta |\overline{q}q|\theta =\mathrm{cos}\frac{\theta }{2}0|\overline{q}q|0`$, see eq. (8). From this low-energy theorem one can easily estimate the matrix element (23) we are interested in. Indeed, using the standard dispersion relations and keeping in the imaginary part the contribution of the $`\eta ^{}`$ meson only ( justified in the large $`N_c`$ limit), we arrive at the following estimate for the $`\eta ^{}`$ matrix element:
$$0|\frac{\alpha _s}{8\pi }G\stackrel{~}{G}|\eta ^{}_\theta \eta ^{}|\overline{u}u+\overline{d}d|0_\theta =m_\eta ^{}^2(\theta )\frac{1}{2}0|\overline{u}u+\overline{d}d|0\mathrm{sin}\frac{\theta }{2}.$$
(26)
In the chiral limit one can neglect the $`\theta `$\- dependence in the mass term $`m_\eta ^{}^2(\theta )m_\eta ^{}^2(\theta =0)+0(m)`$ and in the matrix element $`0|\frac{\alpha _s}{8\pi }G\stackrel{~}{G}|\eta ^{}_\theta 0|\frac{\alpha _s}{8\pi }G\stackrel{~}{G}|\eta ^{}_{\theta =0}+0(m)`$. Therefore, for estimation purposes, one can use the corresponding values calculated for the $`\theta =0`$ vacuum.
$$0|\overline{u}u+\overline{d}d|\eta ^{}_\theta =\frac{\sqrt{3}}{f_\pi (0.50.8)}0|\overline{u}u+\overline{d}d|0\mathrm{sin}\frac{\theta }{2}$$
(27)
where we used the estimate of the matrix element $`0|\frac{\alpha _s}{4\pi }G\stackrel{~}{G}|\eta ^{}\frac{(0.50.8)}{\sqrt{3}}f_\pi m_\eta ^{}^2`$ from ref.. Substituting this expression into (23), we arrive at the following estimate for the amplitude of $`\eta ^{}\pi \pi `$:
$$A(\eta ^{}\pi \pi )\frac{4\sqrt{3}}{(0.50.8)}\frac{m\mathrm{sin}\frac{\theta }{2}0|\overline{q}q|0_{\theta =0}}{f_\pi ^3},$$
(28)
which is the same order of magnitude as $`A(\eta \pi \pi )`$ (19). This value for the amplitude translates to the following estimate for the decay rate
$$\mathrm{\Gamma }(\eta ^{}\pi \pi )\frac{|A(\eta ^{}\pi \pi )|^2}{16\pi m_\eta ^{}}\frac{6}{\pi m_\eta ^{}}\left(\frac{m\mathrm{sin}\frac{\theta }{2}0|\overline{q}q|0}{f_\pi ^3}\right)^22MeV(\mathrm{sin}\frac{\theta }{2})^2,$$
(29)
which is almost an order of magnitude larger than the full width of the $`\eta ^{}`$ meson in the $`\theta =0`$ world. What is more important, this width is exclusively due to the CP odd decay $`\mathrm{\Gamma }(\eta ^{}\pi \pi )`$.
### 2.5 Forming of the induced $`\theta `$ vacua at RHIC
Now we want to argue that the induced $`|\theta ^{ind}`$ vacuum state with an effective $`\theta ^{ind}0`$ ( which was treated so far as a pure theoretical construction) can, nevertheless, be experimentally produced and studied in the real world in the relativistic heavy ion collider. Of course, we do not expect that a stable $`|\theta ^{ind}`$ state can be produced in heavy ion collisions; however, we do expect that a relatively stable and sufficiently large domain with a wrong induced $`|\theta ^{ind}`$ orientation may be formed due to the non-equilibrium dynamics after the QCD phase transition. If this is indeed the case, such a study gives us a unique opportunity to mimic the physics of the world when the fundamental $`\theta ^{fund}0`$. This situation is believed to have occurred during the QCD phase transition in the development of the early Universe.
The idea is very similar to the old idea of creation of the Disoriented Chiral Condensate (DCC) in heavy ion collisions ,, (see also nice review for a discussion of DCC as an example of an out of equilibrium phase transition). DCC refers to regions of space (interior) in which the chiral condensate points in a different direction from that of the ground state (exterior), and separated from the latter by a hot shell of debris. Our starting point is the conjecture that a classically large domain with a nonzero $`U(1)_A`$ chiral phase may be formed in a heavy ion collisions, i.e. we assume that the expectation value for the chiral condensate in a sufficiently large domain has, in general, a nonzero $`U(1)`$ phase: $`\overline{\mathrm{\Psi }}_L\mathrm{\Psi }_Re^{i\varphi ^{singlet}}|\overline{\mathrm{\Psi }}\mathrm{\Psi }|`$. This phase is identified with $`\theta ^{ind}`$. Such an identification is a direct consequence of the transformation properties of the fundamental QCD lagrangian under $`U(1)_A`$ rotations when the chiral singlet phase can be rotated away at the cost of the introduction of the induced $`\theta `$, $`\theta ^{ind}=N_f\varphi ^{singlet}`$. The production of non-trivial $`|\theta ^{ind}`$-vacua would occur in much the same way as discussed above for DCC. The new element is that in addition to chiral fields differing from their true vacuum values the induced $`\theta `$-parameter, which is zero in the real world, becomes effectively nonvanishing in the macroscopically large domain. Of course, once such a conjecture is made ( the classically large domain with a nonzero $`U(1)_A`$ chiral phase is formed) one has to check that this assumption is self-consistent. Namely, one should check that (initially) randomly distributed phases $`\varphi _i`$ in the $`\theta ^{ind}`$-induced background will, indeed, relax to the same non-zero value $`\varphi _i\frac{\theta _{ind}}{N_f}`$ in a sufficiently short period of time. The numerical study of ref. suggests this is indeed the case.
We want to pause here to explain the terminology we adopt regarding the induced $`|\theta ^{ind}`$ vacua. Of course the fundamental parameter $`\theta `$ which enters the fundamental QCD lagrangian is zero in our world. One should remember, however, that the fundamental $`\theta =0`$ is a combination of two pieces: the first part is related to the original term $`\theta Q=\theta \frac{\alpha _s}{8\pi }G^a\stackrel{~}{G}^a`$; the second part is related to the mass matrix $``$ which, in general, may have an arbitrary phase. The $`U(1)_A`$ rotation of the $`\mathrm{\Psi }`$ fields brings this matrix to a canonical form, producing however an additional contribution to the $`\theta `$ term<sup>5</sup><sup>5</sup>5By the way, the famous strong CP problem is formulated in these terms as follows: Why do these two contributions cancel each other with precision better than $`10^9`$?. Our induced $`\theta ^{ind}`$ can be considered as the second term related to the rotation of the $`\mathrm{\Psi }`$ fields. After this point we can apply the same philosophy as for DCC. The chiral fields $`\varphi _i`$ are allowed to take random values as discussed above, and they begin to roll toward the true solution $`\overline{\varphi _u}\overline{\varphi _d}\theta /2,\overline{\varphi _s}0`$ (2.1) and of course overshoot it. The situation is very similar to what was described for the DCC with the only difference that in general we expect an arbitrary $`|\theta ^{ind}`$-disoriented state to be created in heavy ion collisions, not necessarily the $`|\theta =0`$ state. It is assumed that the rapid expansion of the high energy shell leaves behind an effectively zero temperature region in the interior which is isolated from the true vacuum. The high temperature non-equilibrium evolution is very suddenly stopped, or “quenched”, leaving the interior region in a non-equilibrium initial state that then begins to evolve according to (almost) zero temperature Lagrangian dynamics. Therefore starting from an initial non-equilibrium state we can study the behavior of the chiral fields using the zero temperature equations of motion. These equations are non-linear and cannot be solved analytically but the numerical analysis can be done and has been presented in ref., where the equations of motion for the phases of the chiral condensate in the theory with two quark flavors have been studied:
$$\ddot{\varphi }_i^2\varphi _i+\gamma \dot{\varphi _i}+\frac{d}{d\varphi _i}V(\varphi _j,\theta )=0i=u,d.$$
(30)
Here $`^2`$ is a three dimensional spatial derivative and the potential is given in (3). Emission of pions and expansion of the domain will contribute to the damping, $`\gamma `$, as might other processes. The contribution of these unknown effects was simulated in ref. by including a damping term with a reasonable value for the damping constant, $`\gamma \mathrm{\Lambda }_{QCD}200MeV`$. The results of this numerical study can be summarized in the following way:
a)The $`|\theta ^{ind}`$ domain does form. The production of a non-trivial $`\theta `$ is indicated by the fact that the chiral fields relax to constant non-zero values on a time scale over which spatial oscillations of the fields vanish. In other words, the zero mode settles down to a non-zero value $`\frac{\theta }{2}`$. At the same time, all higher momentum modes vanish extremely rapidly and are negligible long before the zero mode settles down to its equilibrium value.
b)The formation of a non-perturbative condensate is also supported by observation of the phenomenon of coarsening, i.e. the phenomenon of amplification of the zero mode as time increases. The effect of coarsening as well as the formation of a nonperturbative condensate is very similar to earlier discussions in ref..
c)The formation of a non-perturbative condensate is also supported by the test of volume -independence of our results. Namely, our quasi zero mode approaches the same value irrespective of the total spatial volume indicating that this really is a nonperturbative condensate that we are evolving to. If it were not we would expect the value of the coefficient to decrease when volume of the system increases.
## 3 Signatures and observables of the produced $`|\theta ^{ind}`$ vacua
Having discussed the properties of the $`\theta `$ vacuum state itself, pseudo-Goldstone bosons as its lightest excitations, and having assumed that this unusual $`|\theta ^{ind}`$ ground state of matter is produced and will live a sufficiently long time $`\tau (10100F)`$, we are now in position to answer the question formulated in the Introduction: “What should we measure to see all these interesting phenomena?”. In what follows we suggest (based on our previous discussions regarding the mass-shifts and other properties of the Goldstone modes) a few simple observables which (hopefully) can be measured and which uniquely signal that the induced $`|\theta ^{ind}`$ vacuum state indeed has been produced.
Before we begin, a few general remarks are in order. First of all, the produced $`|\theta ^{ind}`$ is a relatively short-lived state with respect to electromagnetic and weak interactions and a relatively long-lived state with respect to strong interactions with $`\tau (20MeV)^1(2MeV)^1`$. After the time $`\tau `$, the shell separating $`|\theta ^{ind}`$ world from our world with $`\theta =0`$, breaks down<sup>6</sup><sup>6</sup>6For numerical estimates presented below we use $`\tau (5MeV)^1`$.. All hadrons which have been formed in the $`|\theta ^{ind}`$ background will suddenly find themselves in the new vacuum state (with $`\theta =0`$). They have no choice but to transform to the asymptotic states of this new (for them) $`\theta =0`$ world. The dynamics of this strong transformation is quite complicated and beyond the scope of the present work. However, it is absolutely clear that the only signal which might be relevant for the study of $`|\theta ^{ind}`$-physics is a signature which originates during the time $`\tau `$ while $`|\theta ^{ind}`$ exists. If something happens after that time, it can not be relevant for a study of the $`|\theta ^{ind}`$ vacua. Our second remark is that a signal must safely penetrate through the debris of the hot shell to deliver the information about the induced $`|\theta ^{ind}`$ state and not about something else. Essentially this constraint implies that the decay products originated during the time $`\tau `$ are better to be photons, electrons or muons, rather than the strongly interacting hadrons. In this case the effect related to the electromagnetic rather than strong interactions would be suppressed by a factor $`\tau /\tau _{em}\alpha 1`$, where $`\tau _{em}`$ is a typical time scale for an electromagnetic process. Based on these general remarks and previous discussions (Section 2) regarding the mass-shifts (2.2) and parity-mixing of the hadrons in the $`\theta `$ background, we list below and comment on some signatures which (hopefully) can be measured and would signal the $`|\theta ^{ind}`$ vacua creation in Heavy Ion Collisions.
### 3.1 Signatures from the hadron decays in the $`|\theta ^{ind}`$ background
a. We start from the analysis of $`\eta \pi \pi `$ as the most important signature of the produced $`|\theta ^{ind}`$ vacua. As is known, this decay is strongly suppressed in our world due to the opposite CP parities of the initial ($`\eta `$) and final states ($`2\pi `$). In the previous section we have estimated the width (21) for the decay $`\mathrm{\Gamma }(\eta \pi \pi )0.5MeV(\mathrm{sin}\frac{\theta }{2})^2`$ it turns out to be much larger than the full width $`\mathrm{\Gamma }(\eta )1.18\pm 0.11KeV`$ of the $`\eta `$ meson in our world with $`\theta =0`$. One should expect that a noticeable fraction, $`0.5MeV(\mathrm{sin}\frac{\theta }{2})^2\times \tau 0.1`$, of all produced $`\eta `$ mesons in the $`\theta ^{ind}`$ background would decay to $`2\pi `$ before the shell separating the $`\theta ^{ind}`$ background and our world with $`\theta =0`$ breaks down. If the transition of these $`2\pi `$ mesons (produced in the $`\theta ^{ind}`$ background ) to our $`|\theta =0`$ background is an adiabatic process, then one should expect these $`2\pi `$ mesons become the asymptotic states of our world with the corresponding masses and quantum numbers $`0^{}`$. However, the effect that these $`\pi `$ mesons were produced from $`\eta `$ decay in the $`|\theta ^{ind}`$ background should not be washed out and the corresponding resonance behavior in $`\pi \pi `$ invariant mass should be seen in the spectrum. We should remark at this point that a signal is expected only for low momentum pions, because $`|\theta ^{ind}`$ formation is caused by amplification of low momentum modes. As well, as we mentioned in the previous section, $`|\theta ^{ind}`$ will be different for each given event, such that the corresponding average over the large number of events is zero. However, a measurement on an event by event basis should produce a non-zero result. We should also remark here that mass shifts for Goldstone bosons in $`|\theta `$ background and $`|\theta `$ background are the same. In both cases masses are lower than in our vacuum state $`|\theta =0`$ (see (2.2)). Our last remark: the fact that hadron properties are different in an unusual environment (when temperature $`T`$ and/or chemical potential $`\mu `$ are not zero) is not a new idea (see e.g.). However, $`|\theta ^{ind}`$ vacua is a quite special environment: in this background the hadrons do not posses the definite quantum numbers, like $`P`$ parity. Therefore, their properties could be drastically different from the ones which we know from the Particle Data Booklet.
b. One could repeat the same comments regarding $`\eta ^{}\pi \pi `$ decay as a very profound signature of the produced $`|\theta ^{ind}`$ vacua. As is known, this decay is also strongly suppressed in our world due to the opposite CP parities of the initial ($`\eta ^{}`$) and final states ($`2\pi `$). Besides that, the full width of $`\eta ^{}`$ is quite small $`\mathrm{\Gamma }(\eta ^{})0.2MeV`$ in our world with $`\theta =0`$. In this case the effect is even more profound than in the previously discussed $`\eta \pi \pi `$ decay because a larger fraction of all produced $`\eta ^{}`$ mesons in the $`\theta ^{ind}`$ background would decay to $`2\pi `$ before the shell separating the $`\theta ^{ind}`$ background from our world with $`\theta =0`$ breaks down. Indeed, from eq.(29)we estimate this fraction on the level $`2MeV(\mathrm{sin}\frac{\theta }{2})^2\tau 1/3`$.
c. As we discussed above, the electromagnetic, rather than the strong processes could provide a much better signature for studying the $`|\theta ^{ind}`$ induced vacua. Of course, the relevant effects would be suppressed by a factor $`\tau /\tau _{em}1`$, where $`\tau _{em}`$ is a typical time scale for an electromagnetic process, but, the signal is expected to be sufficiently clean due to the free penetrating of leptons and photons through the hot debris of the hadrons. In particular, the mass of $`\eta ^{}`$ meson in $`|\theta ^{ind}`$ induced background (2.2) can be independently measured from $`\eta ^{}\rho (\omega )\gamma `$ decay along with $`\eta ^{}2\pi `$ decay as described above. The corresponding rate is expected to be on the same level as in our world with $`\theta =0`$, i.e. $`\mathrm{\Gamma }(\eta ^{}\rho (\omega )\gamma )0.06MeV`$. Therefore, a quite noticeable portion of the produced $`\eta ^{}s`$ will decay to $`\rho (\omega )\gamma `$. As usual, we only take into account decays which were originated while the $`|\theta ^{ind}`$ vacua existed, therefore the effect is of the order $`\tau 0.06MeV10^2`$. One could discuss similar decays such as $`\omega \pi ^0(\eta )\gamma `$, $`\varphi \pi ^0(\eta )\gamma `$, $`\eta ^{}(\eta ,\pi ^0)2\gamma `$, etc. with the following general conclusion: the hadron mass shifts and the lack of definite $`P,CP`$ quantum numbers would lead, in general, to some deviation from the standard pattern. In particular, the positions of the resonances as well as the polarization properties of photons emitted from these hadrons, are expected to be quite different from what one could anticipate normally in our world with $`\theta =0`$.
### 3.2 Signatures from the $`|\theta ^{ind}`$ decay
We have been discussing so far the signatures originating from the hadrons ( which are excitations in the $`|\theta ^{ind}`$ background). Now we want to discuss a signature which may occur when the ground state $`|\theta ^{ind}`$ decays by itself. We should pause here to recall some relevant properties of the $`|\theta ^{ind}`$ state.
As explained in the previous sections we expect that, in general, an induced $`|\theta ^{ind}`$\- vacuum state would be created in heavy ion collisions, similar to the creation of the disoriented chiral condensate with an arbitrary isospin direction. It should be a large domain ($`L10100F`$) with a wrong $`\theta ^{ind}0`$ orientation which will mimic the physics of the world with fundamental $`\theta ^{fund}0`$. It is expected that this macroscopically large domain would be quite stable (with respect to the strong interactions) with $`\tau (20MeV)^1(2MeV)^1`$. This state can be understood as the appearance of a macroscopically large ($`10100F`$) domain where the flavor- singlet phase of the chiral condensate is strongly correlated. This singlet phase is identified with $`\theta ^{ind}`$ phase. Indeed, this identification is a direct consequence of the transformation properties of the fundamental QCD lagrangian under $`U(1)_A`$ rotations when the chiral singlet phase can be rotated away at the cost of the appearance of the induced $`\theta ^{ind}`$ term. A different name for the same object would be a “ zero (spatially - independent) mode of the $`\eta ^{}`$ field” which exactly corresponds to the flavor-singlet spatial-independent part of the chiral condensate phase. However, we prefer to use the term $`|\theta ^{ind}`$ because the symbol $`\eta ^{}`$ is usually associated with the $`\eta ^{}`$ meson (excitation) and not with a classical constant field (condensate) in a large domain<sup>7</sup><sup>7</sup>7We would like to make a comment regarding this terminology. In Ref. T.D.Lee considered the possibility of T violation in strong interactions due to the $`\eta ^{}`$ condensate. It is equivalent, in a modern context, to the induced $`\theta ^{ind}0`$ angle in the large domain as explained above. For example, his $`\eta ^{}`$ condensate induces a nonzero electric dipole moment of the neutron, see also comments by D. Kharzeev, R.D.Pisarski, and M.Tytgat in Physics Today..
Having made these short remarks regarding the $`|\theta ^{ind}`$ state, we are in position to answer on the question: “ What happens to the induced $`|\theta ^{ind}`$ state when it blows apart?”. The obvious answer is that it will mainly decay to the neutral pseudo-Goldstone bosons with low momentum modes $`|\stackrel{}{k}|L^1`$. A detailed analysis of the spectrum and other properties of this decay will be presented elsewhere, but here I want to make some estimates based on the simple energetic considerations. As we discussed earlier (2.1), the vacuum energy density in the $`\theta `$ state is greater than in the $`\theta =0`$ vacuum state by the amount:
$$\mathrm{\Delta }E=E_{vac}(\theta )E_{vac}(\theta =0)=2m|\overline{q}q|(1|\mathrm{cos}\frac{\theta }{2}|).$$
When $`|\theta ^{ind}`$ state blows apart, the energy associated with this background will be released as the Goldstone bosons carrying the total energy:
$$\mathrm{\Delta }ϵ_\theta 2m_q|\overline{q}q|(1|\mathrm{cos}\frac{\theta }{2}|)V20\left(\frac{V}{(10F)^3}\right)GeV,$$
(31)
where $`V=L^3`$ is $`3d`$ volume of the $`|\theta ^{ind}`$ background measured in $`Fermi`$. Therefore, given the $`cm`$ energy $`\sqrt{s}=40TeV`$, only a small fraction:
$$\rho \frac{\mathrm{\Delta }ϵ_\theta }{\sqrt{s}}\frac{20(\frac{V}{(10F)^3})GeV}{40TeV}10^3(\frac{V}{(10F)^3})$$
of the total collision energy will be released through the decay of the induced $`|\theta ^{ind}`$ state. Two of the most profound features of the Goldstone bosons produced in this decay are the following: 1. Due to the fact that the $`|\theta ^{ind}`$ formation is caused by amplification of low momentum modes, we expect the spectrum of the Goldstone bosons from this decay to be strongly enhanced at low $`|\stackrel{}{k}|L^1`$; 2. Due to the nonzero value for the topological density $`\theta |\frac{\alpha _s}{8\pi }G\stackrel{~}{G}|\theta =\frac{V_{vac}(\theta )}{\theta }=m_q|0|\overline{q}q|0|\mathrm{sin}\frac{\theta }{2}`$, see (2.1), one could expect that some $`P,CP`$ correlations would appear in this background, as discussed in . However, as we argued earlier this signal may be suppressed by the final state interactions.
Therefore we suggest a somewhat different signature to observe the formation of the induced $`|\theta ^{ind}`$ domain. As mentioned above, an appropriate signal have to be of the electromagnetic origin to be able to deliver the information on the local conditions at its emission site. We argue below that such a signal would be the excess of photons and/or leptons which will be produced through the decay of the induced $`|\theta ^{ind}`$ state with a spectrum which is strongly peacked at very low $`|\stackrel{}{k}|L^1`$ and which demonstrates some unusual polarization properties. The electromagnetic energy which we are about to discuss originates from two different sources: 1. the direct photons/leptons which are produced from the electromagnetic vacuum energy stored in the domain$`V`$; 2. the indirect photons/leptons which are produced from the neutral pseudo -Goldstone bosons $`\pi ^02\gamma ,\eta 2\gamma ,\eta ^{}2\gamma ,\eta 3\pi ^0\gamma ^{}s`$ etc. These neutral pseudo-Goldstone bosons are essentially the decay products of the $`|\theta ^{ind}`$ state constructed from diagonal $`\varphi _i`$ fields (3).
First of all, let us discuss the direct photon/lepton production. We want to argue that there is an electromagnetic vacuum energy stored in the region where the induced nonperturbative vacuum condensates have nonzero values (2.1,8). The first hint (that the expectation value for the electromagnetic field might not be zero ) comes from analysis of the effective lagrangian (3) including the term $`\mathrm{\Delta }V_{em}`$ related to the electromagnetic anomaly,
$$\mathrm{\Delta }V_{em}=\frac{\alpha }{4\pi }N_c\left(\underset{i}{}\varphi _iQ_i^2\right)F_{\mu \nu }\stackrel{~}{F_{\mu \nu }},$$
(32)
where $`\alpha =1/137`$ is the fine structure constant, $`F_{\mu \nu }`$ is electromagnetic field, $`Q_i`$ are the quark charges. Phases, $`\varphi _i`$, in this formula are defined in the same way as before (3). One usually uses the expression (32) for the description of the $`\pi 2\gamma ,\eta 2\gamma ,\eta ^{}2\gamma `$ decays. In this case the $`\varphi _i`$ phases are nothing but the standard Goldstone bosons (with appropariate normalization $`\varphi \frac{\pi }{f_\pi }`$). However, the effective lagrangian (3) together with the additional term (32) has a perfectly physical meaning even when the $`\varphi _i`$ fields describe the vacuum state itself, not necessarilly Goldstone excitations ($`\pi `$ mesons). This was exactly the point of the original papers , and our discussion in section (2.1). In this case, a nonzero $`\varphi _i`$ in a macroscopically large (but finite) domain might imply a nonzero value for the correlation $`\frac{\alpha }{4\pi }N_cF_{\mu \nu }\stackrel{~}{F_{\mu \nu }}`$ over the same region. In other words, due to the electromagnetic interaction $`F_{\mu \nu }`$ with the phase of the quark condensate $`\varphi `$, the electromagnetic field also becomes correlated over the same region.
Another, more direct hint which also suggests that the electromagnetic field has a large scale correlation in the presence of the quark condensate, is based on the observation that the operators $`\frac{\alpha N_c}{4\pi }F_{\mu \nu }F_{\mu \nu }`$ and $`m\overline{\mathrm{\Psi }}\mathrm{\Psi }`$ do mix. This mixing is absolutely irrelevant for QCD itself, where much more important gluon operator such as $`G_{\mu \nu }^2`$ exists. Nevertheless, such a $`F^2m\overline{\mathrm{\Psi }}\mathrm{\Psi }`$ mixing implies that in QCD where the quark fields condense, the electromagnetic field (due to the interaction with these condensed quarks $`m\overline{\mathrm{\Psi }}\mathrm{\Psi }`$) also becomes a correlated field in the nonperturbative QCD vacuum background. This electromagnetic vacuum energy together with the standard (infinite) perturbative vacuum energy (due to the zero vacuum fluctuations ) merely causes an unobservable shift of the vacuum energy. This effect can be safely ignored in all cases except the one where a macroscopically large domain is produced and has ground state properties different from our vacuum. In this case, the electromagnetic vacuum energy can be released through the production of photons when the induced $`|\theta ^{ind}`$ states falls apart.
The simplest way estimate the scale of this effect is to estimate the difference between $`F_{\mu \nu }F_{\mu \nu }`$ in our world with $`\theta =0`$ and with $`\theta 0`$. As was explained above, due to the mixing of the operators $`F_{\mu \nu }F_{\mu \nu }`$ and $`m\overline{\mathrm{\Psi }}\mathrm{\Psi }`$, the problem is reduced to the analysis of the chiral condensate as a function of $`\theta `$, which was discussed earlier, see eq.(8). Therefore, for approximate solution (2.1) we arrive at the following estimate:
$$|\theta |\frac{1}{4}F_{\mu \nu }F_{\mu \nu }|\theta ||0|\frac{1}{4}F_{\mu \nu }F_{\mu \nu }|0||\frac{3\alpha }{2\pi }(Q_u^2+Q_d^2)\mathrm{ln}(\mu L)\left(1\mathrm{cos}(\frac{\theta }{2})\right)m_q\overline{q}q|,$$
(33)
where $`\mu `$ is the normalization point, and $`L10F`$ is a typical dimensional parameter of the system. The sign in the formula is easy to understand: the electromagnetic vacuum energy for the $`|\theta 0`$ ground state is higher than for the $`|\theta =0`$ state because it is proportional to the vacuum contribution of the chiral condensate $`m_q\overline{q}q`$, which has the lowest value for $`\theta =0`$ state, see (2.1).
As we discussed earlier, when the $`|\theta ^{ind}`$ state blows apart, the energy associated with this background will be released mainly as the neutral pseudo-Goldstone bosons carrying the total energy given by eq.(31) such that only a small fraction, $`\rho 10^3`$, of the total collision energy will be released through the decay of the induced $`|\theta ^{ind}`$ state. Our estimate (33) essentially says that only a small fraction of this energy will be released as electromagnetic energy through production of direct photons and leptons.
$$\rho _{em}|\frac{3\alpha }{2\pi }(Q_u^2+Q_d^2)\mathrm{ln}(\mu L)|\rho 10^6.$$
(34)
One should say the formula (34) is probably an overestimation of the effect because we assumed in the derivation that along with the formation of the induced $`|\theta ^{ind}`$ state, the electromagnetic fluctuations also settle down as prescribed by eq.(33) during the same, relatively short period of time $`\tau 10F`$. This may or may not be the case. A further analysis is needed to study a time-scale required to form the electromagnetic correlation (33) after the chiral quark condensate is formed. In any event, this effect is relatively small and will not play an important role in the discussion which follows.
Now we want to discuss the second (much more important) source of the electromagnetic energy which will be released when $`|\theta ^{ind}`$ state blows apart. This source is related to the total energy of the background (31), rather than a small electromagnetic fraction (33) of it. Indeed, as we mentioned earlier, the $`|\theta ^{ind}`$ domain itself will eventually decay to the neutral pseudo-Goldstone bosons $`\pi ^0,\eta ,\eta ^{}`$, the fields which essentially form the $`|\theta ^{ind}`$ domain, see (3). It is important to note that these neutral $`\pi ^0,\eta ,\eta ^{}`$ bosons will mainly decay to photons/dileptons such that large fraction of the background energy (31) eventually will be released as the electromagnetic energy. Indeed, almost all produced $`\pi ^0`$ will decay to photons/dileptons; more than 70 percent of $`\eta `$-mesons will do the same through the decays $`\eta 2\gamma ,\eta 3\pi ^0`$. A noticeable fraction of $`\eta ^{}`$ mesons will also transform to the photons/dileptons through the decays $`\eta ^{}2,\gamma ,\eta ^{}\pi ^0\pi ^0\eta `$. Therefore, we expect that a considerable part of the vacuum energy (31) stored in the $`|\theta ^{ind}`$ domain will be eventually released as it indirect photons/dileptons.
A detailed discussion of the spectrum of the produced particles will be presented elsewhere, but now let me mention that the most profound feature of the photons/dileptons produced in the decay of $`|\theta ^{ind}`$ domain, is quite similar to what we discussed above for the pseudo-Goldstone boson production themselves. Namely, due to the fact that the $`|\theta ^{ind}`$ formation is caused by amplification of low momentum modes, we expect the spectrum of the produced particles to be strongly enhanced at low $`|\stackrel{}{k}|L^1`$. We also expect that the produced particles may have some unusual polarization (P, CP-odd) properties due to the nonzero value of the condensates $`G_{\mu \nu }\stackrel{~}{G_{\mu \nu }}`$ and $`\overline{\mathrm{\Psi }}i\gamma _5\mathrm{\Psi }`$ inside the domain $`V`$. In particular, one could expect that the number of low-energy ($`kL^1`$) right- and left- handed photons will be different.
To conclude this section regarding the electromagnetic energy stored in the nontrivial background, we would like to make the following remark. The electromagnetic $`F_{\mu \nu }^2`$ is always correlated with the quark condensate. Therefore, the corresponding energy may be released when the quark condensate in a large domain (due to unusual conditions) is different from its vacuum value. The $`|\theta ^{ind}`$ state is only one example where this may happen. Another example would be an unusual environment with nonzero temperature and chemical potential $`\mu `$ where, as is well known, the chiral condensate may have a different magnitude. In this case the stored electromagnetic energy will also be released through the direct photon/lepton production as discussed above. The only difference in the estimate (33) would be a replacement
$$\left(1\mathrm{cos}(\frac{\theta }{2})\right)|m_q\overline{q}q|\left[|m_q\overline{q}q_{(\mu =T=0)}||m_q\overline{q}q_{(\mu T0)}|\right],$$
(35)
at least for sufficiently small $`T`$ and $`\mu `$.
Let me conclude this work with the following speculative conjecture. The longstanding problem regarding the strong enhancement for invariant di-lepton masses below the $`\rho `$ mass reported by various collaborations at CERN may be intimately related to the low-energy indirect photon/dilepton production described above. Indeed, as we argued above, a noticeable fraction of the total energy (31) will be eventually released as electromagnetic energy. The main feature of these photons/dileptons is a strong enhancement of the low-energy modes because the corresponding spectrum of the neutral pseudo-Goldstone bosons is peaked at low $`kL^1`$. Detailed calculations are needed before this speculative conjecture can get further support. At least, the standard theoretical analysis (see e.g. a nice summary in ref. and references therein for details) fails to reproduce the enhancement in the data. Let me conclude on this optimistic note/speculation.
## 4 Acknowledgments
I wish to thank Dima Kharzeev, Larry McLerran, Robert Pisarski and Edward Shuryak for valuable comments. I am also grateful to Adrian Melissinos for useful discussions regarding the feasibility of the experiment. I also would like to thank the members of BNL’s theory group for their interest in this work. The research was supported in part by Canadian NSERC.
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# Spectral statistics of chaotic systems with a point-like scatterer
## Abstract
The statistical properties of a Hamiltonian $`H_0`$ perturbed by a localized scatterer are considered. We prove that when $`H_0`$ describes a bounded chaotic motion, the universal part of the spectral statistics are not changed by the perturbation. This is done first within the random matrix model. Then it is shown by semiclassical techniques that the result is due to a cancellation between diagonal diffractive and off-diagonal periodic-diffractive contributions. The compensation is a very general phenomenon encoding the semiclassical content of the optical theorem.
In quantum systems, the chaotic or disordered nature of the classical motion is reflected in the statistical properties of the high lying eigenvalues and eigenvectors. For example, the spectral statistics of ballistic cavities are universal for energy ranges that are small compared to the inverse time of flight through the system. These universal properties are well described by random matrix theory (RMT) .
Consider now a perturbation imposed to a chaotic system. On the classical side, the dynamics is structurally stable, in the sense that a generic smooth perturbation leaves the dynamics chaotic. We are interested in the quantum mechanical effects of a particular class of perturbations that are non-classical. If the unperturbed motion is described by the Hamiltonian $`H_0`$ acting in an $`N`$-dimensional Hilbert space, we consider Hamiltonians of the form
$$H=H_0+\lambda N|vv|,$$
(1)
where $`|v`$ is a fixed vector. $`N`$ is included in the perturbation for future convenience. The eigenvalues $`\{\omega _i\}`$ of $`H`$ satisfy the equation
$$\underset{k}{}\frac{|v_k|^2}{\omega ϵ_k}=\frac{1}{\lambda N},$$
(2)
with $`\{ϵ_k\}`$ the eigenvalues of $`H_0`$ and $`v_k=\phi _k|v`$ the amplitudes of $`|v`$ in the eigenbasis of $`H_0`$.
Rank-one perturbations like in Eq.(1-2) appear in several contexts. The most common one is when a local short-range impurity or point scatterer is added to the system . The physical consequences of such a perturbation were studied for Fermi gases , in the context of RMT and for ballistic motion of particles in regular and chaotic cavities. Another context is in the physics of many body problems. In the basis of $`H_0`$, the perturbation simply reads $`v_kv_l^{}`$ and is said to be separable. In a mean field approach, it is the simplest model leading to collective modes via the residual interaction . Although our results are general, we adopt for simplicity the language of the localized point scatterer.
A local perturbation is purely wave-mechanical. For a system with $`f`$ degrees of freedom, it represents a modification of the dynamics in a volume $`(2\pi \mathrm{})^f`$ in phase space, that tends to zero in the semiclassical limit. For example, the addition of a point scatterer in a ballistic cavity leaves invariant the classical motion while quantum mechanically it induces wave effects like diffraction. The modifications of the eigenvalues by the perturbation are described by Eq.(2). A statistical analysis of the perturbed spectrum was done for the energy levels of a regular integrable cavity with a point scatterer in . It was demonstrated that a short range repulsion between the eigenvalues, different from RMT, is induced by the perturbation, thus considerably modifying the initial Poisson distribution. More recently, M. Sieber has studied, using semiclassical techniques, the modifications by a point scatterer of the spectral statistics of chaotic systems. He showed that diffractive orbits produce finite contributions which may induce deviations with respect to the random matrix model. Whether this deviation really exists for chaotic systems, or on the contrary if there are other (non-diagonal) semiclassical contributions that cancel the purely diffractive terms is the question we answer here.
We prove by two different approaches, namely a purely statistical model and a semiclassical calculation, that a local perturbation produces no deviations with respect to RMT. At the first place, assuming that the unperturbed eigenvalues and eigenvectors components in Eq.(2) are distributed according to RMT, i.e. their joint probability density are given by the standard formulae
$$P(\{ϵ_k\})\underset{i>j}{}|ϵ_iϵ_j|^\beta ,$$
(3)
and
$$P(\{v_k\})=\underset{i}{}\left(\frac{\beta N}{2\pi }\right)^{1\beta /2}\mathrm{exp}(\beta N|v_i|^2/2),$$
(4)
we show that the joint probability density for the perturbed eigenvalues is exactly the same as the distribution of the unperturbed ones
$$P(\{\omega _k\})\underset{i>j}{}|\omega _i\omega _j|^\beta .$$
(5)
Here, $`\beta =1`$ (resp. $`2`$) for systems with (resp. without) time-reversal symmetry. In the second place, and to complete the analysis, a semiclassical calculation of the spectral form factor is considered. It is expressed as a double sum over all the periodic and diffractive orbits of the system. The latter are orbits that are scattered by the perturbation. In the diagonal contribution of the diffractive orbits was obtained (cf Eq.(13) below). We compute the off-diagonal contribution coming from the interference of periodic and diffractive orbits, and find that this contribution exactly cancels the diagonal diffractive term. We thus recover the statistics of RMT. The basic physical ingredient at the basis of this cancellation is the unitarity of quantum scattering processes, i.e. conservation of the flux scattered by the impurity. Although our semiclassical result is less general than Eq.(5) – it is only valid for the short time behavior of a two-point function –, it applies to systems with arbitrary diffraction coefficient not expressible as a Hamiltonian of the form (1).
In chaotic and disordered system the local universal fluctuations of the spectrum are essentially related to the properties of the Jacobian (3). We ignore here problems related to the confinement of the eigenvalues, that are of minor importance for our purpose. We start the proof of Eq.(5) from the joint distribution function of both the old and new eigenvalues, obtained in Ref.
$$P(\{ϵ_i\},\{\omega _j\})\frac{_{i>j}(ϵ_iϵ_j)(\omega _i\omega _j)}{_{i,j}|ϵ_i\omega _j|^{1\beta /2}}\mathrm{e}^{\rho _i(\omega _iϵ_i)},$$
with $`\rho =\beta /2\lambda `$. We restrict for simplicity to $`\lambda >0`$ ($`\lambda <0`$ is treated in the same manner). Eq.(2) imposes the restrictions $`ϵ_i\omega _iϵ_{i+1}`$ (trapping). The distribution for the perturbed eigenvalues, $`\omega _i`$, is then defined as
$`P(\{\omega _i\})`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\omega _1}}dϵ_1{\displaystyle _{\omega _1}^{\omega _2}}dϵ_2\mathrm{}{\displaystyle _{\omega _{N1}}^{\omega _N}}dϵ_NP(\{ϵ_i\},\{\omega _j\})`$ (6)
$``$ $`\mathrm{e}^{\rho _i\omega _i}{\displaystyle \underset{i>j}{}}(\omega _i\omega _j)W(\beta ,\rho ),`$ (7)
with
$$W(\beta ,\rho )=_{\mathrm{}}^{\omega _1}\frac{\mathrm{e}^{\rho ϵ_1}\mathrm{d}ϵ_1}{F(ϵ_1)}\mathrm{}_{\omega _{N1}}^{\omega _N}\frac{\mathrm{e}^{\rho ϵ_N}\mathrm{d}ϵ_N}{F(ϵ_N)}\underset{i>j}{}(ϵ_iϵ_j),$$
and $`F(ϵ)=_j|ϵ\omega _j|^{1\beta /2}`$. Using the standard expression for the Vandermonde determinant in $`W`$ and integrating term by term one obtains
$$W(\beta ,\rho )=det\left[I_j^{(i1)}\right]_{i,j=1,\mathrm{},N},$$
(8)
where $`I_j^{(i)}=_\rho ^iI_j`$ is the i-$`th`$ derivative with respect to $`\rho `$ of
$$I_j=I_j^{(0)}=_{\omega _{j1}}^{\omega _j}\frac{\mathrm{e}^{\rho ϵ}\mathrm{d}ϵ}{F(ϵ)}.$$
(9)
For $`j=1`$, $`\omega _{j1}=\mathrm{}`$.
It is straightforward to check that the $`I_j`$’s satisfy, for any $`j`$, the following differential equation
$$\left[\underset{i}{}(_\rho \omega _i)+\frac{\beta }{2\rho }\underset{i}{}\underset{j(i)}{}(_\rho \omega _j)\right]I_j=0.$$
(10)
This differential equation allows to write
$$I_j^{(N)}=\underset{i=0}{\overset{N1}{}}a_iI_j^{(i)},$$
with some coefficients $`a_i`$. $`W(\beta ,\rho )`$ as defined in Eq.(8) is the Wronskian of this equation. It then follows that
$$_\rho W=a_{N1}W,$$
(11)
where in our case $`a_{N1}=_i\omega _i\beta N/2\rho `$. Integration of Eq.(11) leads to
$$W(\beta ,\rho )=\frac{W_0}{\rho ^{\beta N/2}}\mathrm{exp}(\rho \underset{i}{}\omega _i).$$
When this result is replaced in Eq.(6) one gets
$$P(\{\omega _i\})W_0\underset{i>j}{}(\omega _i\omega _j).$$
(12)
$`W_0`$ is an integration factor, independent of $`\rho `$. We compute it from the asymptotic behavior of $`I_j^{(i)}`$ when $`\rho \mathrm{}`$. In this limit the integral in Eq.(9) may be evaluated explicitly
$$\underset{\rho \mathrm{}}{lim}I_j\frac{\mathrm{e}^{\rho \omega _j}}{\rho ^{\beta /2}_{ij}|\omega _i\omega _j|^{1\beta /2}}.$$
To leading order, $`I_j^{(i)}=\omega _j^iI_j`$ and from Eq.(8) one gets
$$W_0\underset{i>j}{}\frac{|\omega _i\omega _j|^\beta }{(\omega _i\omega _j)}.$$
From this equation and Eq.(12) we recover the random matrix distribution function Eq.(5). A chaotic system coupled to the environment through a one-channel antenna has been considered in . The model is equivalent to Eq.(2) but with imaginary $`\lambda `$. For $`\lambda \mathrm{}`$ the imaginary part of the perturbed energies is small and Eq.(5) is obtained. Our method, which takes explicit care of the trapping problem, allows to prove this result for arbitrary $`\lambda `$.
In real physical systems agreement with random matrix theory is observed in a limited range. This universal behavior concerns correlations over energy ranges that are small compared to $`h/T_{\mathrm{min}}`$, with $`T_{\mathrm{min}}`$ the typical period of the shortest periodic orbit. The above random matrix calculation establishes that the universal part of the spectrum is not changed by the presence of the scatterer. On the other hand, the non-universal behavior of the correlation functions occurring at scales of the order of, or larger than, $`h/T_{\mathrm{min}}`$ are modified by the scattering center, since new diffractive orbits are introduced .
Let us now turn to a semiclassical treatment of the spectral correlations. These are based on trace formulae expansions of the density of states $`d(\omega )=_k\delta (\omega \omega _k)`$, written as a sum of smoothed plus oscillatory terms $`d=\overline{d}+d^{(\mathrm{osc})}`$. We characterize the correlations by the spectral form factor defined as
$$K(\tau )=_{\mathrm{}}^{\mathrm{}}\frac{\mathrm{d}\eta }{\overline{d}}d^{\text{(osc)}}\left(E+\frac{\eta }{2}\right)d^{\text{(osc)}}\left(E\frac{\eta }{2}\right)\mathrm{exp}\left(2\pi i\eta \tau \overline{d}\right).$$
(13)
The average indicated by brackets is taken over an energy window containing many quantum levels but whose size is small compared to $`E`$. We again consider a fully chaotic system with a point-like scatterer. In the geometrical theory of diffraction $`d^{(\mathrm{osc})}=d_p^{(\mathrm{osc})}+d_d^{(\mathrm{osc})}`$, where $`d_p^{(\mathrm{osc})}`$ and $`d_d^{(\mathrm{osc})}`$ are expressed as interferent sums over periodic and diffractive orbits, respectively .
$$d_{p,d}^{(\mathrm{osc})}(E)=\underset{p,d}{}A_{p,d}\mathrm{exp}(i\frac{S_{p,d}(E)}{\mathrm{}}i\frac{\pi }{2}\mu _{p,d}),$$
(14)
with
$`A_p={\displaystyle \frac{T_p}{2\pi \mathrm{}|det(M_p1)|^{1/2}}}`$ (15)
$`A_d={\displaystyle \frac{T_d𝒟(\stackrel{}{n},\stackrel{}{n}^{})\mathrm{e}^{i\pi (f+1)/4}|detN|^{1/2}}{4\pi \mathrm{}k(2\pi \mathrm{})^{(f1)/2}}}.`$ (16)
$`S_{p,d}(E)`$ is the action of the periodic (resp. diffractive) orbits, $`T_{p,d}`$ denotes their period, $`M_p`$ is the monodromy matrix of the periodic orbit, $`N`$ is the matrix $`N_{ij}=^2S_d/y_iy_j`$ (where $`\stackrel{}{y}`$ are coordinates orthogonal to the diffractive trajectory). $`𝒟(\stackrel{}{n},\stackrel{}{n}^{})`$ is the scattering amplitude of the scattering center located at $`\stackrel{}{x_0}`$ with incoming $`\stackrel{}{n}`$ and outgoing $`\stackrel{}{n}^{}`$ directions, defined in terms of the perturbed ($`G`$) and unperturbed ($`G_0`$) Green’s functions by the relation
$$G(\stackrel{}{x},\stackrel{}{x}^{})=G_0(\stackrel{}{x},\stackrel{}{x}^{})+\frac{\mathrm{}^2}{2m}G_0(\stackrel{}{x},\stackrel{}{x_0})𝒟(\stackrel{}{n},\stackrel{}{n}^{})G_0(\stackrel{}{x_0},\stackrel{}{x}^{}).$$
Using the properties of the periodic orbits of chaotic systems, the diagonal contribution of $`d_p^{(\mathrm{osc})}`$ in Eq.(13) gives the short-time random matrix result $`K_p(\tau )=(2/\beta )\tau `$ . The one-scattering contribution of the diffractive orbits in the same approximation is
$$K_d(\tau )=\frac{\tau ^2}{8\beta \pi ^2}\left(\frac{k}{2\pi }\right)^{2f4}\sigma ,$$
(17)
with $`k`$ the modulus of the wave vector at the impurity and $`\sigma `$ its total cross section
$$\sigma =|𝒟(\stackrel{}{n},\stackrel{}{n}^{})|^2d\mathrm{\Omega }d\mathrm{\Omega }^{}.$$
(18)
($`\mathrm{d}\mathrm{\Omega }`$ is the solid angle element). For simplicity, we restrict the calculations to one scattering event (multiple scattering may be considered likewise).
Our purpose is to compute the off-diagonal cross-term coming from the product of $`d_p^{(\mathrm{osc})}`$ and $`d_d^{(\mathrm{osc})}`$ in Eq.(13). The semiclassical expression for this contribution is
$$K_{pd}(\tau )=\frac{2\pi \mathrm{}}{\overline{d}}\underset{p,d}{}A_pA_d^{}\mathrm{exp}(i(S_pS_d)/\mathrm{})\delta (T\frac{T_p+T_d}{2})+c.c..$$
(19)
After energy smoothing, $`K_{pd}`$ has significant contributions only from orbits with close actions $`S_pS_d`$ (having therefore approximately the same period). Pairs of orbits satisfying this condition may be constructed by considering the neighborhood of the forward scattering orbits. To each periodic orbit passing nearby the scatterer $`𝒪`$ we associate an “almost periodic” diffractive orbit that is similar to the periodic orbit but comes back to $`𝒪`$ with a slightly different momentum. In Eq.(19) the double sum now involves all the possible pairs of trajectories constructed this way. Consider a surface of section that includes $`𝒪`$ and is transversal to the momentum of the periodic orbit when it comes nearby $`𝒪`$. Let coordinates measured from $`𝒪`$ and momenta in the plane be denoted by $`(\stackrel{}{q},\stackrel{}{p})`$. Consider all the periodic orbits of period $`T`$ that cut the section through a differential element $`\mathrm{d}^{f1}q\mathrm{d}^{f1}p`$ located at a distance $`\stackrel{}{q}`$ from $`𝒪`$. The difference of action between these periodic orbits and the diffractive orbits associated to them as mentioned above is
$$S_pS_d=(1/2)Q_{ij}q_iq_j,$$
(20)
with
$$Q_{ij}=^2S/q_iq_j+^2S/q_i^{}q_j+^2S/q_i^{}q_j^{},$$
and $`\stackrel{}{q}`$ ($`\stackrel{}{q}^{}`$) are initial (resp. final) coordinates on the surface of section. Moreover, one can show that
$$|detQ|=|det(M_p1)detN|\mathrm{cos}^2\theta ,$$
(21)
where $`\theta `$ is the angle between the normal to the surface of section and the momentum of the diffractive orbit.
By generalizing arguments used in the derivation of the Hannay - Ozorio de Almeida sum rule one can prove the following sum rule
$$\underset{p}{}\frac{\delta (TT_p)\chi (\stackrel{}{q}_p,\stackrel{}{p}_p)}{|det(M_p1)|}=\frac{\mathrm{d}^{f1}q\mathrm{d}^{f1}p\chi (\stackrel{}{q},\stackrel{}{p})}{\mathrm{\Sigma }},$$
(22)
where $`\chi (\stackrel{}{q},\stackrel{}{p})`$ is a test function defined on the surface of section and $`(\stackrel{}{q}_p,\stackrel{}{p}_p)`$ are the coordinates of the points at which the periodic orbit $`p`$ crosses the surface of section. $`\mathrm{\Sigma }=\mathrm{d}^f𝐱\mathrm{d}^f𝐩\delta (EH(𝐱,𝐩))`$ is the total phase-space volume at energy $`E`$. From Eq. (19), using Eqs.(20) and (22), we have
$`K_{pd}=`$ $`{\displaystyle \frac{\overline{d}\tau ^2\mathrm{e}^{i\pi (f+1)/4}}{\beta k(2\pi \mathrm{})^{(f3)/2}\mathrm{\Sigma }}}{\displaystyle \sqrt{|det(M_p1)||detN|}}`$ (24)
$`\times 𝒟^{}(\stackrel{}{n},\stackrel{}{n})\mathrm{e}^{(i/2\mathrm{})Q_{ij}q_iq_j}\mathrm{d}^{f1}q\mathrm{d}^{f1}p+c.c.`$
Integrating the quadratic form in the exponent, taking into account Eq.(21), using the semiclassical density of states $`\overline{d}=\mathrm{\Sigma }/(2\pi \mathrm{})^f`$ and the fact that the differential element for the momenta may be written $`\mathrm{d}^{f1}p=(\mathrm{}k)^{f1}\mathrm{cos}\theta \mathrm{d}\mathrm{\Omega }`$, one obtains the final expression
$$K_{pd}(\tau )=\frac{\tau ^2}{2\pi \beta }\left(\frac{k}{2\pi }\right)^{f2}i\left[𝒟^{}(\stackrel{}{n},\stackrel{}{n})𝒟(\stackrel{}{n},\stackrel{}{n})\right]d\mathrm{\Omega }.$$
(25)
This is the result for the cross-term contribution. Note that it depends only on $`𝒟(\stackrel{}{n},\stackrel{}{n})`$, which translates the fact that in general interference terms between periodic and diffractive orbits can be large only for the diffraction in the forward direction.
To make contact with Eq.(17) we use a general relation valid for the elastic scattering on a finite range potential. The conservation of the flux scattered by the scattering center imposes a relation between the imaginary part of the scattering amplitude and the scattering cross section. This is the well known optical theorem , that in $`f`$ dimensions takes the form
$$i\left[𝒟^{}(\stackrel{}{n},\stackrel{}{n})𝒟(\stackrel{}{n},\stackrel{}{n})\right]=\frac{1}{4\pi }\left(\frac{k}{2\pi }\right)^{f2}|𝒟(\stackrel{}{n},\stackrel{}{n}^{})|^2d\mathrm{\Omega }^{}.$$
Combining this relation with Eq. (25) one gets our final result
$$K_{pd}(\tau )=K_d(\tau ).$$
(26)
The interference between periodic and diffractive orbits exactly cancels the diagonal contribution of the diffractive orbits, Eq.(17). We recover from semiclassical methods, at least for a two-point function and short times, the RMT result.
The two basic elements producing the cancellation are the sum rule (22) and the optical theorem. Only the former is characteristic of chaotic systems, the latter being very general. The present semiclassical results may be extended by similar methods to multiple scattering events. In a wider context, it should be mentioned that this is one of the rare cases in which a calculation of off-diagonal contributions (whose role is essential in producing the correct result) is done explicitly for chaotic systems.
We have concentrated on the fluctuation properties of eigenvalues of chaotic systems, and have demonstrated that they are unchanged by a local perturbation. This applies to high lying states, were the statistical hypotheses hold. On the opposite extreme, a local perturbation may lead to important modifications of the properties of the ground state of the system. Take for example a negative $`\lambda `$. According to Eq.(2), each perturbed eigenvalue remains trapped by two unperturbed ones, except the ground state. The energy of the ground state may diminish arbitrarily with increasing $`|\lambda |`$ and, as can easily be shown, the associated wavefunction becomes more and more localized at the impurity. In our considerations we have ignored the presence of this “collective” mode.
The authors are grateful for many useful discussions with O. Bohigas, M. Saraceno, M. Sieber and U. Smilansky. After completion of this manuscript we became aware of related semiclassical results obtained by M. Sieber.
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# Conditional indifference and conditional preservation
## Introduction
Within the last years, the propositional limitations of classical belief revision have been overcome piece by piece. For instance, Boutilier (?) investigated belief revision within a modal framework, and Williams (?) proposed *transmutation schemas for knowledge systems*. In general, *epistemic states* have moved into the center of interest as representations of belief states of some individual or intelligent agent at a given time. Besides propositional or first-order facts, reflecting certain knowledge, they may contain assumptions, preferences among beliefs, and, as a crucial ingredient, conditional knowledge. They may be represented in different ways, e.g. by gradings of plausibility or disbelief, by making use of epistemic entrenchment, or as a probability distribution.
Epistemic states provide an excellent framework to study *iterative revisions* which are important to build fully dynamic systems<sup>1</sup><sup>1</sup>1An interesting approach to iterated revisions of *belief sets* was proposed quite recently by Lehmann, Magidor and Schlechta in (?).. While propositional AGM theory only observes the results of revisions, considering epistemic states under change allows one to focus on the mechanisms underlying that change, taking conditional beliefs as *revision policies* explicitly into account (cf. (???)). The connection between epistemic states, $`\mathrm{\Psi }`$, (iterative) revision operators, $``$, and conditionals, $`(B|A)`$, is established by the *Ramsey test*
$$\mathrm{\Psi }(B|A)\text{iff}\mathrm{\Psi }AB.$$
(1)
Hence revising epistemic states does not only mean to deal with propositional beliefs – it also requires studying how conditional beliefs are changed.
Darwiche and Pearl (?) rephrased the AGM postulates for epistemic states. Applying the *minimal change paradigm* of propositional belief revision in that extended framework, as Goldszmidt and Boutilier did in (?), however, may produce unintuitive results (?). So, Darwiche and Pearl (?) advanced four postulates as a cautious approach to describe *principles of conditional preservation* when revising epistemic states by *propositional* beliefs. In (??), we extended their approach in considering revisions of epistemic states by *conditionals*. We proposed a set of axioms outlining conditional revisions which are in accordance with some fundamental postulates of revisions (like, for instance, *success*) and with propositional AGM theory, and which preserve conditional beliefs in observing conditional interactions. These interactions were specified by two newly introduced relations between conditionals, that of *subconditionality* and that of *perpendicularity* (see (?)).
Earlier, in (?), we defined a *principle of conditional preservation* in quite a different, namely probabilistic, framework. This principle there was based on the algebraic notion of conditional structures, made use of group theoretical means and postulated the numerical values of the given probability distribution to follow the *conditional structures* of worlds.
In this paper, we will bring together both approaches to conditional preservation, the qualitative one and the quantitative one, in the semi-quantitative framework of ordinal conditional functions. We will rephrase the probabilistic principle of conditional preservation of (?) for ordinal conditional functions, and we will show, that this quantitative principle of conditional preservation implies the corresponding qualitative postulates of (?). Actually, the quantitative version is much stronger than the qualitative one, dealing with *sets of conditionals* instead of only one revising conditional, and handling interactions of conditionals of arbitrary complexity. Although numerical in nature, the principle is based on a symbolic representation of conditional influences on worlds, called *conditional structures*. Numbers (rankings or probabilities) are only considered as manifestations of complex conditional interactions. The representation of conditional structures by the aid of group theory provides a rich methodological framework to study conditionals in belief revision and nonmonotonic reasoning.
In the following section, we fix some notations and describe the relationship between conditionals and epistemic states. Then we introduce the crucial notion of conditional structures which the property of conditional indifference is based upon. By applying the concept of indifference to revision functions, we obtain a precise formalization of the principle of conditional preservation, and we characterize ordinal conditional functions observing this principle. Further on, we compare our approach to Goldszmidt, Morris & Pearl’s system-Z and system-$`Z^{}`$ in several examples. Finally, we bring together qualitative and quantitative approaches to the principle of conditional preservation, proving the formalization given here to be a most fundamental one. A summary and an outlook conclude this paper.
## Conditionals and epistemic states
We consider a propositional language $``$ over a finite alphabet $`a,b,c\mathrm{}`$. Let $`\mathrm{\Omega }`$ denote the set of possible worlds for $``$, i.e. $`\mathrm{\Omega }`$ is a complete set of interpretations of $``$. Throughout this paper, we will write $`\overline{A}`$ instead of $`\neg A`$, and $`AB`$ instead of $`AB`$, for formulas $`A,B`$.
Conditionals $`(B|A)`$ represent statements of the form “*If A then B*”, expressing a relationship between two (propositional) formulas $`A`$, the *antecedent* or *premise*, and $`B`$, the *consequent*. $`()`$ denotes the set of all conditionals $`(B|A)`$ with $`A,B`$. A conditional $`(B|)`$ with a tautological antecedent is taken to correspond to its (propositional) consequent, $`B`$. $`(D|C)`$ is called a *subconditional* of $`(B|A)`$, written as
$$(D|C)(B|A),$$
(2)
iff $`CDAB`$ and $`C\overline{D}A\overline{B}`$. Typically, subconditionals arise by strengthening the antecedent of a conditional, e.g. $`(b|ac)`$ is a subconditional of $`(b|a)`$, $`(b|ac)(b|a)`$.
Each possible world $`\omega \mathrm{\Omega }`$ either confirms $`(B|A)`$, in case that $`\omega AB`$, or refutes it, if $`\omega A\overline{B}`$, or does not even satisfy its premise, $`\omega \vDash ̸A`$, and so is of no relevance for it. Although conditionals are evaluated with respect to worlds, they cannot really be accepted (as entities) in single, isolated worlds. To validate conditionals, we need richer epistemic structures than plain propositional interpretations, at least to compare different worlds with regard to their relevance for a conditional (see, for example, (???)). Epistemic states as representations of cognitive states of intelligent agents provide an adequate framework for conditionals.
An epistemic notion that turned out to be of great importance both for conditionals and epistemic states, in particular in the context of belief revision, is that of *plausibility*: conditionals are supposed to represent plausible conclusions, and plausibility relations on formulas or worlds, respectively, guide AGM-revisions of belief sets and of epistemic states (???).
As Spohn (?) emphasized, however, it is not enough to consider the qualitative ordering of propositions according to their plausibility – also relative distances between degrees of plausibility should be taken into account. So he introduced ordinal conditional functions $`\kappa `$ (*OCF’s, ranking functions*) (?) from worlds to ordinals such that some worlds are mapped to the minimal element $`0`$. Here, we will simply assume that OCF’s are functions $`\kappa :\mathrm{\Omega }\{0,\mathrm{}\}`$ from the set of worlds to the natural numbers, extended by $`0`$ and $`\mathrm{}`$. They specify non-negative integers as degrees of plausibility – or, more precisely, as degrees of *disbelief* – for worlds. The smaller $`\kappa (\omega )`$ is, the more plausible the world $`\omega `$ appears, and what is believed (for certain) in the epistemic state represented by $`\kappa `$ is described precisely by the set $`\text{Mod}(\kappa ):=\{\omega \mathrm{\Omega }\kappa (\omega )=0\}`$. For propositional formulas $`A,B`$, we set $`\kappa (A)=\mathrm{min}\{\kappa (\omega )\omega A\}`$, so that $`\kappa (AB)=\mathrm{min}\{\kappa (A),\kappa (B)\}`$. In particular, $`0=\mathrm{min}\{\kappa (A),\kappa (\overline{A})\}`$, so that at least one of $`A`$ or $`\overline{A}`$ is considered mostly plausible. A proposition $`A`$ is believed iff $`\kappa (\overline{A})>0`$, which is denoted by $`\kappa A`$. A conditional $`(B|A)()`$ may be assigned a degree of plausibility via $`\kappa (B|A)=\kappa (AB)\kappa (A)`$. Each OCF $`\kappa `$ induces a (propositional) AGM-revision operator $``$ by setting $`\text{Mod}(\kappa A)=\mathrm{min}_\kappa (\text{Mod}(A))`$ (see (?)). The Ramsey test (1) then reads $`\kappa (B|A)`$ iff $`\kappa AB`$. This is in accordance with the plausibility relation imposed by $`\kappa `$, as the following lemma shows:
###### Lemma 1
Let $`(B|A)`$ be a conditional in $`()`$, let $`\kappa `$ be an ordinal conditional function. Then $`\kappa (B|A)`$ (by applying the Ramsey test) iff $`\kappa (AB)<\kappa (A\overline{B})`$.
So $`\kappa `$ accepts a conditional (via the Ramsey test) iff $`AB`$ is more plausible than $`A\overline{B}`$. The proof of this lemma is immediate.
## Conditional structures
By observing the behavior of worlds with respect to it, each conditional $`(B|A)`$ can be considered as a generalized (namely three-valued) indicator function on worlds:
$$(B|A)(\omega )=\{\begin{array}{cc}\hfill 1:& \omega AB\hfill \\ \hfill 0:& \omega A\overline{B}\hfill \\ \hfill u:& \omega \overline{A}\hfill \end{array}$$
(3)
where $`u`$ stands for undefined (??)). Intuitively, incorporating a conditional as a plausible conclusion in an epistemic state means to make – at least some – worlds confirming the conditional more plausible than the worlds refuting it. In this sense, conditionals to be learned have effects on possible worlds (more exactly, on their degrees of plausibility), shifting them appropriately to establish the intended plausible relationship. (3) then provides a classification of worlds for achieving this: On confirming worlds $`\omega AB`$, i.e. $`(B|A)(\omega )=1`$, $`(B|A)`$ possibly has a positive effect, while on refuting worlds $`\omega A\overline{B}`$, $`(B|A)`$ possibly has a negative effect; the effects on worlds $`\omega `$ with $`(B|A)(\omega )=u`$ is unclear. Which worlds will actually be shifted depends on the chosen revision procedure – for the conditional, all worlds in either of the partitioning sets are indistinguishable.
When we consider (finite) sets of conditionals $`=\{(B_1|A_1),\mathrm{},(B_n|A_n)\}()`$, we have to modify the representation (3) appropriately to identify the effect of each conditional in $``$ on worlds in $`\mathrm{\Omega }`$. This leads to introducing the functions $`\sigma _i=\sigma _{(B_i|A_i)}`$ below (see (4)) which generalize (3) by replacing the numbers $`0`$ and $`1`$ by abstract symbols. Moreover, we will make use of a group structure to represent the joint impact of conditionals on worlds.
To each conditional $`(B_i|A_i)`$ in $``$ we associate two symbols $`𝐚_i^+,𝐚_i^{}`$. Let
$$_{}=𝐚_1^+,𝐚_1^{},\mathrm{},𝐚_n^+,𝐚_n^{}$$
be the free abelian group with generators $`𝐚_1^+,𝐚_1^{},\mathrm{},𝐚_n^+,𝐚_n^{}`$, i.e. $`_{}`$ consists of all elements of the form $`(𝐚_1^+)^{r_1}(𝐚_1^{})^{s_1}\mathrm{}(𝐚_n^+)^{r_n}(𝐚_n^{})^{s_n}`$ with integers $`r_i,s_i`$ (the ring of integers). Each element of $`_{}`$ can be identified by its exponents, so that $`_{}`$ is isomorphic to $`^{2n}`$ (?). The commutativity of $`_{}`$ corresponds to the fact that the conditionals in $``$ shall be effective simultaneously, without assuming any order of application. So our way of dealing with conditionals is a symmetric, homogeneous one – we do not need (user-defined) priorities among conditionals. Note that, although we will speak of *multiplication* and *products* in $`_{}`$, the generators of $`_{}`$ are merely juxtaposed, like words.
For each $`i,1in`$, we define a function $`\sigma _i:\mathrm{\Omega }_{}`$ by setting
$$\sigma _i(\omega )=\{\begin{array}{cc}\hfill 𝐚_i^+\text{if}& (B_i|A_i)(\omega )=1\hfill \\ \hfill 𝐚_i^{}\text{if}& (B_i|A_i)(\omega )=0\hfill \\ \hfill 1\text{if}& (B_i|A_i)(\omega )=u\hfill \end{array}$$
(4)
$`\sigma _i(\omega )`$ represents the manner in which the conditional $`(B_i|A_i)`$ applies to the possible world $`\omega `$. The neutral element $`1`$ of $`_{}`$ corresponds to the non-applicability of $`(B_i|A_i)`$ in case that the antecedent $`A_i`$ is not satisfied. The function $`\sigma _{}:\mathrm{\Omega }_{}`$,
$$\sigma _{}(\omega )=\underset{1in}{}\sigma _i(\omega )=\underset{\genfrac{}{}{0pt}{}{1in}{\omega A_iB_i}}{}𝐚_i^+\underset{\genfrac{}{}{0pt}{}{1in}{\omega A_i\overline{B}_i}}{}𝐚_i^{}$$
describes the all-over effect of $``$ on $`\omega `$. $`\sigma _{}(\omega )`$ is called *(a representation of) the conditional structure of $`\omega `$ with respect to $``$*. For each world $`\omega `$, $`\sigma _{}(\omega )`$ contains at most one of each $`𝐚_i^+`$ or $`𝐚_i^{}`$, but never both of them because each conditional applies to $`\omega `$ in a well-defined way. The next lemma (which is easy to prove) shows that this property characterizes conditional structure functions:
###### Lemma 2
Let $`\sigma :\mathrm{\Omega }`$ be a map from the set of worlds $`\mathrm{\Omega }`$ to the free abelian group $`=𝐚_1^+,𝐚_1^{},\mathrm{},𝐚_n^+,𝐚_n^{}`$ generated by $`𝐚_1^+,𝐚_1^{},\mathrm{},𝐚_n^+,𝐚_n^{}`$, such that $`\sigma (\omega )`$ contains at most one of each $`𝐚_i^+`$ or $`𝐚_i^{}`$, for each world $`\omega \mathrm{\Omega }`$. Then there is a set of conditionals $``$ with $`\text{card}()n`$ such that $`\sigma =\sigma _{}`$.
Example 3 Let $`=\{(c|a),(c|b)\}`$, where $`a,b,c`$ are atoms, and let $`_{}=𝐚_1^+,𝐚_1^{},𝐚_2^{},𝐚_2^{}`$. We associate $`𝐚_1^\pm `$ with the first conditional, $`(c|a)`$, and $`𝐚_2^\pm `$ with the second one, $`(c|b)`$. The following table shows the values of the function $`\sigma _{}`$ on worlds $`\omega \mathrm{\Omega }`$:
$$\begin{array}{cccc}& & & \\ \omega \hfill & \sigma _{}(\omega )\hfill & \omega \hfill & \sigma _{}(\omega )\hfill \\ & & & \\ abc\text{}\hfill & 𝐚_1^+𝐚_2^+\hfill & \overline{a}bc\hfill & 𝐚_2^+\hfill \\ ab\overline{c}\hfill & 𝐚_1^{}𝐚_2^{}\hfill & \overline{a}b\overline{c}\hfill & 𝐚_2^{}\hfill \\ a\overline{b}c\hfill & 𝐚_1^+\hfill & \overline{a}\overline{b}c\hfill & 1\hfill \\ a\overline{b}\overline{c}\hfill & 𝐚_1^{}\hfill & \overline{a}\overline{b}\overline{c}\hfill & 1\hfill \end{array}$$
$`abc`$ confirms both conditionals, so its conditional structure is represented by $`𝐚_1^+𝐚_2^+`$. This corresponds to the product (in $`_{}`$) of the conditional structures of the worlds $`\overline{a}bc`$ and $`a\overline{b}c`$. Two worlds, namely $`\overline{a}\overline{b}c`$ and $`\overline{a}\overline{b}\overline{c}`$, are not affected at all by the conditionals in $``$.
The logical structure of antecedents and consequents of the conditionals in $``$ does not really matter, nor do logical relationships between the conditionals. All that we need is a conditional’s partitioning property on the set of worlds (cf. (3) and (4)). $`\sigma _{}`$ labels each world appropriately and allows us to compare different worlds with respect to the impact the conditionals in $``$ exert on them. The following example illustrates that also multiple copies of worlds may be necessary to relate conditional structures:
Example 4 Consider the set $`=\{(d|a),(d|b),(d|c)\}`$ of conditionals using the atoms $`a,b,c,d`$. Let $`𝐚_1^\pm ,𝐚_2^\pm ,𝐚_3^\pm `$ be the group generators associated with $`(d|a),(d|b)`$, $`(d|c)`$, respectively. Then we have
$`\sigma _{}(ab\overline{c}d)\sigma _{}(a\overline{b}cd)\sigma _{}(\overline{a}bcd)`$ $`=`$ $`(𝐚_1^+𝐚_2^+)(𝐚_1^+𝐚_3^+)(𝐚_2^+𝐚_3^+)`$
$`=(𝐚_1^+)^2(𝐚_2^+)^2(𝐚_3^+)^2`$ $`=`$ $`(𝐚_1^+𝐚_2^+𝐚_3^+)^2`$
$`=`$ $`\sigma _{}(abcd)^2.`$
Here two copies of $`abcd`$, or of its structure, respectively, are necessary to match the product of the conditional structures of $`ab\overline{c}d,a\overline{b}cd`$ and $`\overline{a}bcd`$.
To compare worlds adequately with respect to their conditional structures, we take the worlds $`\omega \mathrm{\Omega }`$ as formal generators of the free abelian group
$$\widehat{\mathrm{\Omega }}:=\omega \omega \mathrm{\Omega }$$
$`\widehat{\mathrm{\Omega }}`$ consists of all products $`\widehat{\omega }=\omega _{1}^{}{}_{}{}^{r_1}\mathrm{}\omega _{m}^{}{}_{}{}^{r_m}`$, with $`\omega _1,\mathrm{},\omega _m\mathrm{\Omega }`$, and $`r_1,\mathrm{}r_m`$ integers. Introducing such a “multiplication between worlds” is nothing but a technical means to comply with the multiplicative structure the effects of conditionals impose on worlds. As in $`_{}`$, multiplication in $`\widehat{\mathrm{\Omega }}`$ actually means juxtaposition. In (?), where we first developed these ideas, we considered multi-sets of worlds (corresponding to elements in $`\widehat{\mathrm{\Omega }}`$ with only positive exponents) and calculated the conditional structure of such a multi-set as the *conditional weight* it is carrying. Making use of arbitrary elements of $`\widehat{\mathrm{\Omega }}`$ as group elements, however, provides a much more convenient and elegant framework to deal with conditional structures. We will usually write $`{\displaystyle \frac{\omega _1}{\omega _2}}`$ instead of $`\omega _1\omega _2^1`$.
Now $`\sigma _{}`$ may be extended to $`\widehat{\mathrm{\Omega }}`$ in a straightforward manner by setting
$$\sigma _{}(\widehat{\omega })=\sigma _{}(\omega _1)^{r_1}\mathrm{}\sigma _{}(\omega _m)^{r_m},$$
yielding a *homomorphism of groups* $`\sigma _{}:\widehat{\mathrm{\Omega }}_{}`$. For $`\widehat{\omega }=\omega _{1}^{}{}_{}{}^{r_1}\mathrm{}\omega _{m}^{}{}_{}{}^{r_m}\widehat{\mathrm{\Omega }}`$, we obtain
$`\sigma _{}(\omega _{1}^{}{}_{}{}^{r_1}\mathrm{}\omega _{m}^{}{}_{}{}^{r_m})=`$
$`{\displaystyle \underset{1in}{}}(𝐚_i^+)^{_{k:\sigma _i(\omega _k)=𝐚_i^+}r_k}{\displaystyle \underset{1in}{}}(𝐚_i^{})^{_{k:\sigma _i(\omega _k)=𝐚_i^{}}r_k},`$
as a representation of its conditional structure. The exponent of $`𝐚_i^+`$ in $`\sigma _{}(\widehat{\omega })`$ indicates the number of worlds in $`\widehat{\omega }`$ which confirm the conditional $`(B_i|A_i)`$, each world being counted with its multiplicity, and in the same way, the exponent of $`𝐚_i^{}`$ indicates the number of worlds that are in conflict with $`(B_i|A_i)`$.
By investigating suitable elements of $`\widehat{\mathrm{\Omega }}`$, it is possible to isolate the (positive or negative) net impacts of conditionals in $``$ , as the following example illustrates:
Example 5 (continued) In Example 2 above, we have
$$\sigma _{}(\frac{abc}{\overline{a}bc})=\frac{𝐚_1^+𝐚_2^+}{𝐚_2^+}=𝐚_1^+$$
So $`{\displaystyle \frac{abc}{\overline{a}bc}}`$ reveals the positive net impact of the conditional $`(c|a)`$ within $``$, symbolized by $`𝐚_1^+`$.
Similarly, in Example 2, the element $`{\displaystyle \frac{ab\overline{c}\overline{d}\overline{a}\overline{b}c\overline{d}}{a\overline{b}c\overline{d}}}`$ isolates the negative net impact of the second conditional, $`(d|b)`$:
$`\sigma _{}\left({\displaystyle \frac{ab\overline{c}\overline{d}\overline{a}\overline{b}c\overline{d}}{a\overline{b}c\overline{d}}}\right)={\displaystyle \frac{𝐚_1^{}𝐚_2^{}𝐚_3^{}}{𝐚_1^{}𝐚_3^{}}}=𝐚_2^{}`$.
The following example is taken from (?, p. 68f):
Example 6 Consider the set $``$ consisting of the following conditionals:
$$\begin{array}{ccc}r_1:\hfill & (f|b)\hfill & \text{Birds fly.}\hfill \\ r_2:\hfill & (b|p)\hfill & \text{Penguins are birds.}\hfill \\ r_3:\hfill & (\overline{f}|p)\hfill & \text{Penguins do not fly.}\hfill \\ r_4:\hfill & (w|b)\hfill & \text{Birds have wings.}\hfill \\ r_5:\hfill & (a|f)\hfill & \text{Animals that fly are airborne.}\hfill \end{array}$$
In Table 1, we list the conditional structures of all possible worlds; this table will be helpful in the sequel.
Having the same conditional structure defines an equivalence relation $`_{}`$ on $`\widehat{\mathrm{\Omega }}`$:
$$\widehat{\omega }_1_{}\widehat{\omega }_2\text{iff}\sigma _{}(\widehat{\omega }_1)=\sigma _{}(\widehat{\omega }_2).$$
(5)
Those elements of $`\widehat{\mathrm{\Omega }}`$ that are balanced with respect to the effects of conditionals in $``$ are contained in the *kernel of $`\sigma _{}`$*, $`\text{ker}\sigma _{}=\{\widehat{\omega }\widehat{\mathrm{\Omega }}\sigma _{}(\widehat{\omega })=1\}`$. $`\text{ker}\sigma _{}`$ does not depend on the chosen representation of conditional structures by symbols in $`_{}`$ and thus, it is an invariant of $``$ (?).
Often, besides the conditionals explicitly given in $``$, implicit normalizing constraints have to be taken into account, like, e.g. $`\kappa ()=0`$ for ordinal conditional functions. This can be achieved by focusing on equivalence with repect to $`\sigma _{}`$. Since $`\sigma _{}`$ simply counts the generators occurring in $`\widehat{\omega }`$, two elements $`\widehat{\omega }_1=\omega _1^{r_1}\mathrm{}\omega _m^{r_m},\widehat{\omega }_2=\nu _1^{s_1}\mathrm{}\nu _p^{s_p}\widehat{\mathrm{\Omega }}`$ are $`\sigma _{}`$-equivalent, $`\widehat{\omega }_1_{}\widehat{\omega }_2`$, iff $`_{1jm}r_j=_{1kp}s_k`$. This means, $`\widehat{\omega }_1_{}\widehat{\omega }_2`$ iff they both are a (cancelled) product of the same number of generators, each generator being counted with its corresponding exponent.
## Conditional indifference
To study conditional interactions, we now focus on the behavior of OCF’s $`\kappa :\mathrm{\Omega }\{0,\mathrm{}\}`$ with respect to the multiplication in $`\widehat{\mathrm{\Omega }}`$. Each such function may be extended to a homomorphism, $`\kappa :\widehat{\mathrm{\Omega }}_+(,+)`$, by setting
$$\kappa (\omega _{1}^{}{}_{}{}^{r_1}\mathrm{}\omega _{m}^{}{}_{}{}^{r_m})=r_1\kappa (\omega _1)+\mathrm{}+r_m\kappa (\omega _m),$$
where $`\widehat{\mathrm{\Omega }}_+`$ is the subgroup of $`\widehat{\mathrm{\Omega }}`$ generated by the set $`\mathrm{\Omega }_+:=\{\omega \mathrm{\Omega }\kappa (\omega )\mathrm{}\}`$. This allows us to analyze numerical relationships holding between different $`\kappa (\omega )`$. Thereby, it will be possible to elaborate the conditionals whose structures $`\kappa `$ follows, that means, to determine sets of conditionals $`()`$ with respect to which $`\kappa `$ is indifferent:
Definition 7 Suppose $`\kappa :\mathrm{\Omega }\{0,\mathrm{}\}`$ is an OCF, and $`()`$ is a set of conditionals such that $`\kappa (A)\mathrm{}`$ for all $`(B|A)`$. $`\kappa `$ is *indifferent with respect to* $``$ iff the following two conditions hold:
1. If $`\kappa (\omega )=\mathrm{}`$ then there is $`(B|A)`$ such that $`\sigma _{(B|A)}(\omega )1`$ and $`\kappa (\omega ^{})=\mathrm{}`$ for all $`\omega ^{}`$ with $`\sigma _{(B|A)}(\omega ^{})=\sigma _{(B|A)}(\omega )`$.
2. $`\kappa (\widehat{\omega }_1)=\kappa (\widehat{\omega }_2)\text{whenever }\sigma _{}(\widehat{\omega }_1)=\sigma _{}(\widehat{\omega }_2)`$ for $`\widehat{\omega }_1_{}\widehat{\omega }_2\widehat{\mathrm{\Omega }}_+`$.
If $`\kappa `$ is indifferent with respect to $`()`$, then it does not distinguish between different elements $`\widehat{\omega }_1,\widehat{\omega }_2`$ with the same conditional structure with respect to $``$. Normalizing constraints are taken into account by observing $`_{}`$-equivalence. Conversely, any deviation $`\kappa (\widehat{\omega })0`$ can be explained by the conditionals in $``$ acting on $`\widehat{\omega }`$ in a non-balanced way. Condition (i) in Definition Conditional indifference is necessary to deal with worlds $`\omega \mathrm{\Omega }_+`$. Conditional indifference, as defined in Definition Conditional indifference, captures interactions of conditionals of arbitrary depth by making use of the homomorphism induced by $`\kappa `$. It also respects, however, indifference on the superficial level of the function $`\kappa `$ itself:
###### Lemma 8
If the ordinal conditional function $`\kappa `$ is indifferent with respect to $``$, then $`\sigma _{}(\omega _1)=\sigma _{}(\omega _2)`$ implies $`\kappa (\omega _1)=\kappa (\omega _2)`$ for all worlds $`\omega _1,\omega _2\mathrm{\Omega }`$.
The next theorem gives a simple criteria to check conditional indifference with ordinal conditional functions. Moreover, it provides an intelligible schema to construct conditional indifferent functions.
###### Theorem 9
An OCF $`\kappa `$ is indifferent with respect to a set $`=\{(B_1|A_1),\mathrm{},(B_n|A_n)\}`$ of conditionals iff $`\kappa (A_i)\mathrm{}`$ for all $`i,1in`$, and there are rational numbers $`\kappa _0,\kappa _i^+,\kappa _i^{}`$, $`1in`$, such that for all $`\omega \mathrm{\Omega }`$,
$$\kappa (\omega )=\kappa _0+\underset{\genfrac{}{}{0pt}{}{1in}{\omega A_iB_i}}{}\kappa _i^++\underset{\genfrac{}{}{0pt}{}{1in}{\omega A_i\overline{B_i}}}{}\kappa _i^{}$$
(6)
*Sketch of proof.* According to Lemma 8, the equivalence relation (5) provides a rough classification of the worlds in $`\mathrm{\Omega }`$ with respect to the conditionals in $``$. Obtaining a representation of the form (6) then amounts to checking the solvability of a linear equational system. The proof of this theorem is very similar to the proof of the analogous theorem for probabilistic representation of knowledge given in (?). $`\mathrm{}`$
## The principle of conditional preservation
Minimality of change is a crucial paradigm for belief revision, and a “principle of conditional preservation” is to realize this idea of minimality when conditionals are involved in change. Minimizing absolutely the changes in conditional beliefs, as in (?), is an important proposal to this aim, but it does not always lead to intuitive results (?). The idea we will develop here rather aims at *preserving the conditional structure of knowledge* within an epistemic state which we assume to be represented by an OCF $`\kappa `$.
We just explained what it means for an OCF $`\kappa `$ to follow the structure imposed by $``$ on the set of worlds by introducing the notion of conditional indifference (cf. Definition Conditional indifference). Pursuing this approach further in the framework of belief revision, a revision of $`\kappa `$ by simultaneously incorporating the conditionals in $``$, $`\kappa ^{}=\kappa `$, can be said to preserve the conditional structure of $`\kappa `$ with respect to $``$ if the *relative change function* $`\kappa ^{}\kappa `$ is indifferent with respect to $``$<sup>2</sup><sup>2</sup>2Why just $`\kappa ^{}\kappa `$? First, it is more accurate than, e.g., $`\mathrm{max}\{0,\kappa ^{}\kappa \}`$, in the sense of taking differences in degrees of plausibility seriously. Second, it makes use of the conditional “$``$”, considering revision as a generalized conditional operation; for more details, see (?).. Taking into regard the worlds $`\omega `$ with $`\kappa (\omega )=\mathrm{}`$ appropriately, this gives rise to the following definitions:
Definition 10 Let $`\kappa `$ be an OCF, and let $``$ be a finite set of conditionals. Let $`\kappa ^{}=\kappa `$ denote the result of revising $`\kappa `$ by $``$. Presuppose further<sup>3</sup><sup>3</sup>3Note that success $`\kappa ^{}`$ is not compulsory for conditional indifference and conditional preservation. We only presuppose $`\kappa ^{}(A)\mathrm{}`$ for all $`(B|A)`$ in order to exclude pathological cases. that $`\kappa ^{}(A)\mathrm{}`$ for all $`(B|A)`$.
1. $`\kappa ^{}`$ is called *$`\kappa `$-consistent* iff $`\kappa (\omega )=\mathrm{}`$ implies $`\kappa ^{}(\omega )=\mathrm{}`$.
2. $`\kappa ^{}`$ is *indifferent with respect to $``$ and $`\kappa `$* iff $`\kappa ^{}`$ is $`\kappa `$-consistent and the following two conditions hold:
1. If $`\kappa ^{}(\omega )=\mathrm{}`$ then $`\kappa (\omega )=\mathrm{}`$, or there is $`(B|A)`$ such that $`\sigma _{(B|A)}(\omega )1`$ and $`\kappa ^{}(\omega ^{})=\mathrm{}`$ for all $`\omega ^{}`$ with $`\sigma _{(B|A)}(\omega ^{})=\sigma _{(B|A)}(\omega )`$.
2. $`(\kappa ^{}\kappa )(\widehat{\omega }_1)=(\kappa ^{}\kappa )(\widehat{\omega }_2)\text{whenever }\sigma _{}(\widehat{\omega }_1)=\sigma _{}(\widehat{\omega }_2)`$ and $`\widehat{\omega }_1_{}\widehat{\omega }_2`$ for $`\widehat{\omega }_1,\widehat{\omega }_2\widehat{\mathrm{\Omega }}_+^{}`$, where $`\widehat{\mathrm{\Omega }}_+^{}=\omega \mathrm{\Omega }\kappa ^{}(\omega )\mathrm{}`$.
The principle of conditional preservation is now realized as an indifference property:
Definition 11 Let $`\kappa `$ be an OCF, and let $``$ be a finite set of conditionals. A revision $`\kappa ^{}=\kappa `$ satisfies the *principle of conditional preservation* iff $`\kappa ^{}`$ is indifferent with respect to $``$ and $`\kappa `$.
So $`\kappa `$ satisfies the principle of conditional preservation if any change in plausibility is clearly and unambigously induced by $``$. The next theorem characterizes revisions of ordinal conditional functions that satisfy the principle of conditional preservation. The theorem is obvious by observing Theorem 9.
###### Theorem 12
Let $`\kappa ,\kappa ^{}`$ be OCF’s, and let $`=\{(B_1|A_1),`$ $`\mathrm{},(B_n|A_n)\}`$ be a (finite) set of conditionals in $`()`$. A revision $`\kappa ^{}=\kappa `$ satisfies the principle of conditional preservation iff $`\kappa ^{}(A_i)\mathrm{}`$ for all $`i,1in`$, and there are numbers $`\kappa _0,\kappa _i^+,\kappa _i^{},1in`$, such that for all $`\omega \mathrm{\Omega }`$,
$$\kappa ^{}(\omega )=\kappa (\omega )+\kappa _0+\underset{\genfrac{}{}{0pt}{}{1in}{\omega A_iB_i}}{}\kappa _i^++\underset{\genfrac{}{}{0pt}{}{1in}{\omega A_i\overline{B_i}}}{}\kappa _i^{}$$
(7)
Comparing Theorems 9 and 12 with one another, we see that an OCF $`\kappa `$ is indifferent with respect to a finite set of conditionals $``$ iff it can be taken as a revision $`\kappa _0`$ satisfying the principle of conditional preservation, where $`\kappa _0(\omega )=0`$ for all $`\omega \mathrm{\Omega }`$ is the *uniform ordinal conditional function*.
Up to now, we have not yet taken the *success condition* $`\kappa ^{}`$ into regard, postulating that the revised OCF in fact represents the conditionals in $``$.
Definition 13 Let $`\kappa ,\kappa ^{}`$ be OCF’s, and let $``$ be a set of conditionals. $`\kappa ^{}=\kappa `$ is called a *c-revision* iff $`\kappa ^{}`$ and $`\kappa ^{}`$ satisfies the principle of conditional preservation. $`\kappa `$ is called a *c-representation* of $``$, iff $`\kappa `$ and $`\kappa `$ is indifferent with respect to $``$.
Theorems 9 and 12 provide simple schemes to construct c-revisions and c-representations. The numbers $`\kappa _0,\kappa _i^+,\kappa _i^{},1in`$, then have to be chosen appropriately to ensure that $`\kappa ^{()}`$ is an ordinal conditional function, and such that $`\kappa ^{()}(AB)<\kappa ^{()}(A\overline{B})`$ for all conditionals $`(B|A)`$ (cf. Lemma 1). In the special case that $`\kappa `$ is a representation of $``$, we obtain the following corollary by some easy calculations:
###### Corollary 14
Let $`=\{(B_1|A_1),\mathrm{},(B_n|A_n)\}`$ be a (finite) set of conditionals in $`()`$, and let $`\kappa `$ be an OCF.
$`\kappa `$ is a c-representation of $``$ iff $`\kappa (A_i)\mathrm{}`$ for all $`i,1in`$, and there are numbers $`\kappa _0,\kappa _i^+,\kappa _i^{},1in`$, such that for all $`\omega \mathrm{\Omega }`$,
$$\kappa (\omega )=\kappa _0+\underset{\genfrac{}{}{0pt}{}{1in}{\omega A_iB_i}}{}\kappa _i^++\underset{\genfrac{}{}{0pt}{}{1in}{\omega A_i\overline{B_i}}}{}\kappa _i^{},$$
(8)
and
$`\kappa _i^{}\kappa _i^+`$ $`>`$ $`\underset{\omega A_iB_i}{\mathrm{min}}\left({\displaystyle \underset{\genfrac{}{}{0pt}{}{ji}{\omega A_jB_j}}{}}\kappa _j^++{\displaystyle \underset{\genfrac{}{}{0pt}{}{ji}{\omega A_j\overline{B}_j}}{}}\kappa _j^{}\right)`$
$``$ $`\underset{\omega A_i\overline{B}_i}{\mathrm{min}}\left({\displaystyle \underset{\genfrac{}{}{0pt}{}{ji}{\omega A_jB_j}}{}}\kappa _j^++{\displaystyle \underset{\genfrac{}{}{0pt}{}{ji}{\omega A_j\overline{B}_j}}{}}\kappa _j^{}\right)`$
*Proof.* $`\kappa `$ is a c-representation of $`R`$ iff $`\kappa `$ and $`\kappa `$ is indifferent with respect to $``$. From Theorem 9, we obtain representation (8). Due to Lemma 1, $`\kappa `$ iff $`\kappa (A_iB_i)<\kappa (A_i\overline{B}_i)`$, i.e. iff
$`\underset{\omega A_iB_i}{\mathrm{min}}\kappa _0+{\displaystyle \underset{\genfrac{}{}{0pt}{}{1jn}{\omega A_jB_j}}{}}\kappa _j^++{\displaystyle \underset{\genfrac{}{}{0pt}{}{1jn}{\omega A_j\overline{B}_j}}{}}\kappa _j^{}`$
$`<\underset{\omega A_i\overline{B}_i}{\mathrm{min}}\kappa _0+{\displaystyle \underset{\genfrac{}{}{0pt}{}{1jn}{\omega A_jB_j}}{}}\kappa _j^++{\displaystyle \underset{\genfrac{}{}{0pt}{}{1jn}{\omega A_j\overline{B}_j}}{}}\kappa _j^{},`$
which is equivalent to
$`\underset{\omega A_iB_i}{\mathrm{min}}\kappa _i^++{\displaystyle \underset{\genfrac{}{}{0pt}{}{ji}{\omega A_jB_j}}{}}\kappa _j^++{\displaystyle \underset{\genfrac{}{}{0pt}{}{ji}{\omega A_j\overline{B}_j}}{}}\kappa _j^{}`$
$`<\underset{\omega A_i\overline{B}_i}{\mathrm{min}}\kappa _i^{}+{\displaystyle \underset{\genfrac{}{}{0pt}{}{ji}{\omega A_jB_j}}{}}\kappa _j^++{\displaystyle \underset{\genfrac{}{}{0pt}{}{ji}{\omega A_j\overline{B}_j}}{}}\kappa _j^{}.`$
This shows (14). $`\mathrm{}`$
The difference $`\kappa _i^{}\kappa _i^+`$, or the right hand side of (14), respectively, measures the effort needed to establish the i-th conditional. To calculate suitable constants $`\kappa _i^+,\kappa _i^{}`$, we apply the following heuristics: To establish conditional beliefs, one can make confirming worlds more plausible (if required, which amounts to choose $`\kappa _i^+0`$), or refuting worlds less plausible (if required, which means $`\kappa _i^{}0`$). The normalizing constant $`\kappa _0`$ then has to be chosen appropriately to ensure that actually an OCF is obtained.
We prefer the second alternative, presupposing
$$\kappa _i^{}0,\text{ and }\kappa _i^+=0\text{for}\mathrm{\hspace{0.33em}1}in$$
(10)
Then (14) reduces to
$$\kappa _i^{}>\underset{\omega A_iB_i}{\mathrm{min}}\underset{\genfrac{}{}{0pt}{}{ji}{\omega A_j\overline{B}_j}}{}\kappa _j^{}\underset{\omega A_i\overline{B}_i}{\mathrm{min}}\underset{\genfrac{}{}{0pt}{}{ji}{\omega A_j\overline{B}_j}}{}\kappa _j^{}$$
(11)
for $`1in`$. If there are worlds $`\omega `$ with neutral conditional structure, $`\sigma _{}(\omega )=1`$, we may set $`\kappa _0=0`$. So, we obtain a c-representation of $``$ via
$$\kappa (\omega )=\underset{\genfrac{}{}{0pt}{}{1in}{\omega A_i\overline{B}_i}}{}\kappa _i^{},$$
(12)
where the $`\kappa _i^{}`$ have to satisfy (11).
Example 15 We will use Corollary 14 and the heuristics (10) to obtain a c-representation (12) of the conditionals $`=\{r_1,\mathrm{},r_5\}`$ of Example 2. To calculate constants $`\kappa _i^{}`$ according to (11), Table 1 proves to be helpful. We only have to focus on the $`𝐚^{}`$-labels of worlds, and we obtain
$`\kappa _5^{},\kappa _4^{},\kappa _1^{}>0,`$
$`\kappa _3^{}>\mathrm{min}\{\kappa _1^{},\kappa _2^{}\},\kappa _2^{}>\mathrm{min}\{\kappa _1^{},\kappa _3^{}\}`$
So we set
$$\kappa _5^{}=\kappa _4^{}=\kappa _1^{}=1,\kappa _2^{}=\kappa _3^{}=2.$$
(13)
Since there are also worlds $`\omega `$ with $`\sigma _{}(\omega )=1`$ (cf. Table 1), we set $`\kappa _0=0`$. So we obtain a c-representation of $`r_1,\mathrm{},r_5`$ by
$$\kappa (\omega )=\underset{\genfrac{}{}{0pt}{}{1i5}{\omega A_i\overline{B}_i}}{}\kappa _i^{},$$
(14)
where the $`A_i,B_i`$’s are the antecedents and consequents of the rules $`r_i`$ and the $`\kappa _i^{}`$ are defined as in (13) for $`1i5`$ (see also Table 2 in Example A comparison with system-Z and system-$`Z^{}`$ below).
## A comparison with system-Z and system-$`Z^{}`$
A well-known method to represent a (finite) set $`=\{r_i=(B_i|A_i)1in\}`$ of conditionals by an OCF is to apply the *system-Z* of Goldszmidt and Pearl (??). The corresponding ranking function $`\kappa ^z`$ is given by
$$\kappa ^z(\omega )=\{\begin{array}{c}0,\text{if }\omega \text{ does not falsify any }r_i,\hfill \\ 1+\underset{\genfrac{}{}{0pt}{}{1in}{\omega A_i\overline{B}_i}}{\mathrm{max}}Z(r_i),\text{otherwise}\hfill \end{array}$$
(15)
where $`Z`$ is an ordering on $``$ observing the (logical) interactions of the conditionals (for a detailed description of $`Z`$, see, for instance, (?)). $`\kappa ^z`$ assigns to each world $`\omega `$ the lowest possible rank admissible with respect to the constraints in $``$. Comparing (15) with (6), we see that in general, $`\kappa ^z`$ is *not* a c-representation of $``$, since in its definition (15), *maximum* is used instead of *summation* (see Example A comparison with system-Z and system-$`Z^{}`$ below). The numbers $`Z(r_i)`$, however, may well serve to define appropriate constants $`\kappa _i^{}`$ in (6). Setting $`\kappa _0=\kappa _i^+=0`$, and $`\kappa _i^{}=Z(r_i)+1`$ for $`1in`$, we obtain from $`Z`$ a c-representation $`\kappa _c^z`$ of $``$ via
$$\kappa _c^z(\omega )=\{\begin{array}{c}0,\text{if }\omega \text{ does not falsify any }r_i,\hfill \\ \underset{\genfrac{}{}{0pt}{}{1in}{\omega A_i\overline{B}_i}}{}(Z(r_i)+1),\text{otherwise}.\hfill \end{array}$$
(16)
An even more sophisticated representation is obtained by combining the system-Z approach with the principle of maximum entropy (*ME-principle*), yielding system-$`Z^{}`$ (?). The corresponding $`Z^{}`$-rankings of the conditionals in $``$ have to satisfy the following equation (see equation (16) in (?, p. 225))
$`Z^{}(r_i)+\underset{\omega A_i\overline{B}_i}{\mathrm{min}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{ji}{\omega A_j\overline{B}_j}}{}}Z^{}(r_j)`$ (17)
$`=1+\underset{\omega A_iB_i}{\mathrm{min}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{ji}{\omega A_j\overline{B}_j}}{}}Z^{}(r_j),`$
and $`\kappa ^{}`$ is then calculated by
$$\kappa ^{}(\omega )=\underset{\omega A_i\overline{B}_i}{}Z^{}(r_i)$$
(18)
(see equation (18) in (?, p. 225)). For so-called *minimal-core sets* – these are sets $``$ allowing each conditional to be separable from the other rules by restricting conditional interactions –, a procedure is given to calculate $`Z^{}`$-rankings in (?).
Like our method, system-$`Z^{}`$ makes use of summation instead of maximization, as in system-Z. And equations (17) determining the $`Z^{}`$-rankings look similar to our inequality constraints (14). More exactly, if we follow the heuristics (10) and set $`Z^{}(r_i)=\kappa _i^{}`$, then system-$`Z^{}`$ turns out to be a special instance of our more general scheme in Corollary 14. In particular, system-$`Z^{}`$ yields a c-representation.
This similarity is not accidental – the ME-principle not only provides a powerful base for system-$`Z^{}`$, but also influenced the idea of conditional indifference presented in this paper. In (?), we characterized the ME-principle by four axioms, one of which was the *postulate of conditional preservation*. Conditional preservation for probability functions there was realized in full analogy to that for OCF’s defined here. So both c-representations and ME-distributions comply with a fundamental principle for representing conditionals, and it is this principle of conditional preservation (or principle of conditional indifference, respectively) that is responsible for a peculiar thoroughness and accuracy when incorporating conditionals. Since we realized this principle completely in a semi-quantitative setting, we did not have to refer to probabilities and to ME-distributions, and we were able to formalize the acceptance conditions, (14), in a purely qualitative manner.
We will illustrate our method by various examples which are taken from (?) and (?) to allow a direct comparison with system-Z and system-$`Z^{}`$.
Example 16 Consider once again the conditionals $`r_1,\mathrm{},r_5`$ from Example 2. Here we have $`Z(r_1)=Z(r_4)=Z(r_5)=0`$ and $`Z(r_2)=Z(r_3)=1`$ (for the details, see (?, p. 69)). By setting $`\kappa _i^{}=Z(r_i)+1`$ for $`1i5`$, we obtain the same constants, (13), as in Example 14. Furthermore, by applying the procedure *Z-rank* in (?), we calculate $`Z^{}(r_1)=Z^{}(r_4)=Z^{}(r_5)=1`$ and $`Z^{}(r_2)=Z^{}(r_3)=2`$, therefore also $`Z^{}(r_i)=\kappa _i^{}`$, $`1i5`$. So, in this example, $`\kappa _c^z`$ from (16) actually is the OCF from (14) and coincides with $`\kappa ^{}`$. For instance,
$`\kappa _c^z(\overline{p}\overline{b}\overline{f}wa)`$ $`=`$ $`0,`$
$`\kappa _c^z(pb\overline{f}wa)`$ $`=`$ $`\kappa _1^{}=1,`$
$`\kappa _c^z(p\overline{b}fw\overline{a})`$ $`=`$ $`\kappa _2^{}+\kappa _3^{}+\kappa _5^{}=5.`$
In Table 2, we list the ranks of all possible worlds, first computed by system-Z, according to (15), and then computed as a c-representation, $`\kappa _c^z`$, of $``$, according to (16).
This table reveals clearly that $`\kappa ^z`$ is not a c-representation of $``$: Associating symbols $`𝐚_i^+,𝐚_i^{}`$ with the conditionals $`r_i`$ in $``$, $`1i5`$, respectively, we obtain
$$\sigma _{}\left(\frac{pbfwa\overline{p}bfw\overline{a}}{pbfw\overline{a}\overline{p}bfwa}\right)=\frac{𝐚_1^+𝐚_2^+𝐚_3^{}𝐚_4^+𝐚_5^+𝐚_1^+𝐚_4^+𝐚_5^{}}{𝐚_1^+𝐚_2^+𝐚_3^{}𝐚_4^+𝐚_5^{}𝐚_1^+𝐚_4^+𝐚_5^+}=1,$$
but $`\kappa ^z\left({\displaystyle \frac{pbfwa\overline{p}bfw\overline{a}}{pbfw\overline{a}\overline{p}bfwa}}\right)`$ $`=\kappa ^z(pbfwa)+\kappa ^z(\overline{p}bfw\overline{a})\kappa ^z(pbfw\overline{a})\kappa ^z(\overline{p}bfwa)=2+12=10`$.
What is the actual benefit of this formal principle of conditional preservation? Comparing $`\kappa ^z(\omega )`$ to $`\kappa _c^z(\omega )`$, we see that $`\kappa _c^z`$ is more fine-grained. Therefore, it represents more conditionals. Consider, e.g., the conditional $`(w|pb\overline{f}a)`$ – does a non-flying, but airborne penguin possess wings or not? $`\kappa ^z`$ does not know, we have $`\kappa ^z(pb\overline{f}aw)=\kappa ^z(pb\overline{f}a\overline{w})=1`$. On the other hand, $`\kappa _c^z`$ accepts this conditional: $`\kappa _c^z(pb\overline{f}aw)=1<2=\kappa _c^z(pb\overline{f}a\overline{w})`$. To show that this is more than pure speculation, consider the conditional structures of $`pb\overline{f}aw`$ and $`pb\overline{f}a\overline{w}`$: $`\sigma _{}(pb\overline{f}aw)=𝐚_1^{}𝐚_2^+𝐚_3^+𝐚_4^+`$, $`\sigma _{}(pb\overline{f}a\overline{w})=𝐚_1^{}𝐚_2^+𝐚_3^+𝐚_4^{}`$. Thus, except for $`r_4`$, both worlds behave exactly the same with respect to the conditionals in $``$, but, due to $`r_4`$, our penguin is supposed to have wings since it is a bird.
Similar arguments apply when considering the conditionals $`(a|pbf),(a|pbfw),(w|pbf),(w|pbfa)`$: $`\kappa ^z`$ is totally indifferent when confronted with flying super-penguins - it assigns the same degree of plausibility, 2, to any of the involved worlds. So, it accepts neither of these conditionals, whereas $`\kappa _c^z`$ accepts all of them.
Finally, let us consider the conditional $`(a|bfw)`$. It is accepted by $`\kappa ^z`$, as well as by $`\kappa _c^z`$, as may easily be checked. But what happens when the variable $`p`$ is taken into account? $`\kappa ^z`$ establishes $`(a|\overline{p}bfw)`$, but is undecided with respect to $`(a|pbfw)`$. By contrast, $`\kappa _c^z`$ not only accepts both of these conditionals, but also establishes them with equal strength:
$$\kappa _c^z\left(\frac{pbfwa}{pbfw\overline{a}}\right)=\kappa _c^z\left(\frac{\overline{p}bfwa}{\overline{p}bfw\overline{a}}\right)=1.$$
(19)
This is a simple consequence of the principle of conditional preservation in this case, since $`\sigma _{}\left({\displaystyle \frac{pbfwa\overline{p}bfw\overline{a}}{pbfw\overline{a}\overline{p}bfwa}}\right)=1`$ (see above). (19) is justified because each of the involved quotients has the same conditional structure, i.e. shows the same behavior, with respect to $``$. By considering arbitrary group elements in $`\text{ker}\sigma _{}`$, even very complicated interrelationships between degrees of strength associated with conditionals can be observed.
Therefore, using System-Z and the $`\mathrm{max}`$-operator means to establish conditionals only on a superficial level, whereas obeying the principle of conditional preservation ensures that conditional knowledge is propagated thoroughly and deeply in plausibility structures of epistemic states. This well-behavedness with respect to subconditionals is also observed in (?). It is illustrated by the next example, too:
Example 17 Consider the conditionals
$$\begin{array}{ccc}r_1:\hfill & (f|s)\hfill & \text{Swedes are fair-haired.}\hfill \\ r_2:\hfill & (t|s)\hfill & \text{Swedes are tall.}\hfill \end{array}$$
We apply Corollary 14 and heuristics (10) and calculate $`\kappa _1^{},\kappa _2^{}0`$. So we set $`\kappa _1^{}=\kappa _2^{}=1`$, and we obtain a c-representation, $`\kappa `$, of the form (12). Then not only the subconditional $`(f|st)`$ of $`(f|s)`$ is accepted, but also the subconditional $`(f|s\overline{t})`$, because $`\kappa (sf\overline{t})=1<2=\kappa (s\overline{f}\overline{t})`$.
Goldszmidt, Morris and Pearl (?) compared this situation to the one where instead of $`r_1,r_2`$, merely the conditional $`(ft|s)`$ is learned. Here, only the first subconditional, $`(f|st)`$, is accepted, but not the second one, $`(f|s\overline{t})`$. This becomes intelligible by considering conditional structures: $`sf\overline{t}`$ and $`s\overline{f}\overline{t}`$ both show the same behavior with respect to $`(ft|s)`$ (namely, they refute it), while $`sft`$ and $`s\overline{f}t`$ show different behaviors.
Observing conditional structures emphasizes once again that in general, the joint integration of conditionals cannot be achieved by learning only one conditional – conditionals resist to propositional treatment. In our framework, each conditional constitutes an independent piece of knowledge.
Finally, let us consider an example that cannot be dealt with by system-$`Z^{}`$ in a straightforward manner because the involved set of conditionals is not a minimal-core set:
Example 18 Let $``$ consist of the rules
$$r_1:(b|a),r_2:(c|b),r_3:(c|a).$$
We list the conditional structures in Table 3 to make argumentation easier.
$``$ is not a minimal-core set in the sense of (?) because $`r_3`$ is only refuted by worlds that also refute either $`r_1`$ or $`r_2`$. So equation (17) is not solvable to yield $`Z^{}`$-rankings.
For our approach, however, dealing with $``$ is no problem: Using (10) and (11), we obtain
$$\kappa _1^{},\kappa _2^{}>0,\kappa _3^{}>0\mathrm{min}\{\kappa _1^{},\kappa _2^{}\}$$
We set $`\kappa _1^{}=\kappa _2^{}=1`$ and $`\kappa _3^{}=0`$, and we obtain by (12) an appropriate ranking function, $`\kappa `$ (see Table 4).
Actually, $`r_3`$ seems to be redundant since $`\kappa _3^+=\kappa _3^{}=0`$, and it is in fact already established by $`r_1`$ and $`r_2`$. But what about the subconditionals $`(c|ab)`$ and $`(c|a\overline{b})`$ of $`r_3`$? While the first one is accepted, due to the impact of $`r_2`$, the second subconditional is not, since $`r_3`$ is taken to be redundant (see the conditional structures in Table 3).
To ensure the thorough propagation of conditional knowledge to subconditionals, we have to postulate $`\kappa _i^{}>0`$ in (10), in order to protect the influence of each conditional against numerical cancellations. With $`\kappa _1^{}=\kappa _2^{}=\kappa _3^{}=1`$, we obtain the OCF $`\kappa _1`$ in Table 4, which accepts both subconditionals of $`r_3`$, as desired.
## Unifying qualitative and quantitative approaches
We defined the principle of conditional preservation as an indifference property of the revised ranking function (cf. Definition The principle of conditional preservation) – conditional preservation means to maintain numerical relationships. Therefore, it appears here as an essentially quantitative notion. Nevertheless, it is applied to ranking functions which are used in a qualitative, or at least, semi-quantitative setting. Moreover, note that conditional indifference is based on the notion of conditional structures which are represented in a completely symbolic way. So, what is the connection between the principle of conditional preservation introduced here, and the approaches to conditional preservation in a purely qualitative framework, as were proposed in (?) and in (?)?
In (?), we advanced a set of postulates apt to guide the revisions of epistemic states, $`\mathrm{\Psi }`$, by conditional beliefs, $`(B|A)`$. There, the idea of conditional preservation was put in formal terms by making use of two relations, *subconditionality*, $``$ (see (2)), and *perpendicularity*, $`\underset{¯}{}`$, on the set of conditionals:
$$(D|C)\underset{¯}{}(B|A)\text{iff either }CAB,\text{ or }A\overline{B},\text{ or }\overline{A}.$$
If $`(D|C)\underset{¯}{}(B|A)`$, then for all worlds which $`(D|C)`$ may be applied to, $`(B|A)`$ has the same effect and thus, it yields no further partitioning – $`(B|A)`$ is *irrelevant* for $`(D|C)`$.
Conditional preservation then was described in (?) by the following three postulates:
If $`(D|C)\underset{¯}{}(B|A)`$ then $`\mathrm{\Psi }(D|C)`$ iff
$`\mathrm{\Psi }(B|A)(D|C)`$.
If $`(D|C)(B|A)`$ and $`\mathrm{\Psi }(D|C)`$ then
$`\mathrm{\Psi }(B|A)(D|C)`$.
If $`(D|C)(\overline{B}|A)`$ and $`\mathrm{\Psi }(B|A)(D|C)`$ then $`\mathrm{\Psi }(D|C)`$.
These axioms cover the postulates of Darwiche and Pearl in (?) (see (??)) and support their intuitive ideas with more formal arguments. Applying Definition The principle of conditional preservation to the case $`=\{(B|A)\}`$, the principle of conditional preservation reduces to postulate
$$(\kappa ^{}\kappa )(\omega _1)=(\kappa ^{}\kappa )(\omega _2)\text{if}(B|A)(\omega _1)=(B|A)(\omega _2)$$
for the revised function $`\kappa ^{}=\kappa `$ (cf. also (?)). It can then be shown that such a revision also satisfies the postulates (CR5)-(CR7) stated above (?).
This means that the principle of conditional preservation, phrased in its full (numerical) complexity in this paper, also covers the approaches to conditional preservation proposed in a qualitative framework.
To guarantee the thorough propagation of conditional knowledge to subconditionals, we may supplement (CR5)-(CR7) with another postulate:
If $`(D|C)(B|A)`$ and $`\mathrm{\Psi }\vDash ̸(\overline{D}|C)`$, then
$`\mathrm{\Psi }(B|A)(D|C)`$.
(CR8) clearly exceeds the paradigm of conditional preservation, in favor of imposing conditional structure as long as there are no conflicts.
Axioms (CR5)-(CR8) only deal, however, with revising an epistemic state by one single conditional. A topic of our ongoing research is to generalize them so as to apply for incorporating sets of conditionals, too.
## Summary and outlook
In this paper, we presented an approach to realize the idea of conditional preservation in knowledge representation and belief revision in a very comprehensive way. We introduced a formal notion of conditional structure of worlds, and then we phrased a *principle of conditional preservation* for (revisions of) ordinal conditional functions as indifference with respect to these structures. We presented simple schemes to construct ordinal conditional functions observing this principle, and we compared our approach to system-Z and system-$`Z^{}`$. Finally, we showed that the principle formalized here also covers other, qualitative approaches to conditional preservation.
In (?), we developed these ideas in an even more general setting, with so-called *conditional valuation functions* representing epistemic states. Ordinal conditional functions are special instances of these functions, as well as probability distributions and possibility distributions. So the principle of conditional preservation, as was presented here, can be applied to nearly all major (semi-)quantitative representation forms for epistemic states. Actually, this paper continues and extends work begun in (?).
The algebraic notion of *conditional structures* of worlds played a crucial part for conditional preservation here. Its representation via group theory does not only provide an elegant methodological framework for handling conditionals. In (?), we present an approach to solve the “inverse representation problem” – which conditionals are most appropriate to represent an epistemic state? – by making extensive use of group theoretical means to elaborate conditional structures.
Acknowledgements
I am grateful to three anonymous referees whose comments helped me to improve the presentation of my results.
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# Fate of Kaluza-Klein Bubble *footnote **footnote *Preprint numbers: CGPG-99/12-8, RESCEU-6/00, DAMTP-2000-30, hep-th/0003066. This article will appear in Phys. Rev. D.
## I Introduction
It is likely that superstring or M-theory governs the physics of gravity or space-time in higher energy stages. Such theories are naturally formulated in the higher dimensions than four. We expect a plausible scenario that such a higher dimensional space-time somehow evolves to the stable four dimensional space-time according to the history of the Universe . The so-called brane world scenario is the most actively being investigated along to this line. This scenario is motivated by Horava and Witten’s theory which shows that an eleven dimensional supergravity theory on the orbihold $`𝐑^{10}\times S^1/Z_2`$ is related to the ten-dimensional $`E_8\times E_8`$ heterotic string theory. Therein the matters are confined to the ten-dimensional space-time (three-brane) and gravitons are propagating in the full eleven dimensions. The brane world space-time should be stable.
Although the brane world scenario may be plausible at the reduction from eleven to ten dimensions, the space-time will be still compactified to four dimensions in the normal Calabi-Yau’s way. Regarding to these full scenario of the compactification, the stability of the space-time becomes the important issue to be investigated. The positive energy theorem guarantees the stability of the four dimensional asymptotically flat space-time in the framework of general relativity . Surprisingly, the existence of the extra dimensions can drastically change the situation. Witten showed that the five dimensional Minkowski space-time decays into the so-called Kaluza-Klein (KK) bubble space-time unless we assume the existence of the elementary fermion related to supersymmetry , of which existence we can not expect generally. This also may indicate that the ‘bubble’ appears somewhere at the bulk or on the brane in the brane world scenario and disturbed the three-brane where we are living.
The metric of the KK bubble space-time given by Witten is written as
$`ds_5^2`$ $`=`$ $`r^2dt^2+\left(1{\displaystyle \frac{r_0^2}{r^2}}\right)d\chi ^2+\left(1{\displaystyle \frac{r_0^2}{r^2}}\right)^1dr^2+r^2\mathrm{cosh}^2td\mathrm{\Omega }^2,`$ (1)
where the $`\chi `$-direction will be compactified and $`rr_0`$, and $`d\mathrm{\Omega }^2=d\theta ^2+\mathrm{sin}^2\theta d\phi ^2`$. In general case, the metric has a conical singularity at $`r_0`$. However, if we carefully take a periodicity along the $`\chi `$-direction, the metric can be regularized. More precisely to see this, we write the metric near $`r=r_0`$ as
$`ds_5^2`$ $``$ $`r_0^2dt^2+{\displaystyle \frac{2(rr_0)}{r_0}}d\chi ^2+{\displaystyle \frac{r_0}{2(rr_0)}}dr^2+r_0^2\mathrm{cosh}^2td\mathrm{\Omega }^2`$ (2)
$`=`$ $`r_0^2dt^2+2r_0\left[R^2d\left({\displaystyle \frac{\chi }{r_0}}\right)^2+dR^2\right]+r_0^2\mathrm{cosh}^2td\mathrm{\Omega }^2,`$ (3)
where $`R=\sqrt{rr_0}`$. Then we realize that the period should be set to be $`\chi _p=2\pi r_0`$. As one can see, the ‘boundary of bubble’ located at $`r=r_0`$ expands rapidly like $`\mathrm{cosh}t`$ and the space-time does not have naked singularities. Here, we remind you that the total energy is zero. We are imaging the boundary of the space-time, $`r=r_0`$, as the surface of the ‘bubble’.
Interestingly Brill and Pfister gave an initial data which has the negative total energy related to the size of the compactified dimension. The space-time with the negative energy may be favor in the aspect of energetics. One year later, Brill and Horowitz gave an initial data in a simple way. (We will briefly review their construction in Sec. II.) Contrasted to the ‘Witten bubble’, their solution has arbitrary negative energy regardless of the size of the compactified space. This is too far from our intuition that the negative energy is proportional to the Casimir energy due to the boundary effect of the compactified space. Therefore it is difficult to imagine the classical evolution after the vacuum decay.
Corley and Jacobson discussed the subsequent evolution of Brill-Horowitz’s initial data. They found that the positive acceleration of the bubble’s surface area for the negative mass bubbles, and they conjectured that KK bubble with negative energy cannot collapse. However, their study is not sufficient to conclude the final fate of the bubble, as they already mentioned, because they considered only the initial behavior of the time-symmetric data and they did not do any dynamical studies.
In this paper, we report our numerical analysis on this final fate problem of KK bubble, especially of the negative mass bubble. We start our numerical simulation from the Brill-Horowitz’s initial data, and evolve the space-time using the standard Arnowitt-Deser-Misner formulation (but 4+1 dimensional decomposition). We will show that the space-time initially behaves as Corley-Jacobson’s analysis, and expands forever, although the acceleration will be negative. Despite of the expanding, we will observe the irregular behavior of the curvature invariant.
This paper is organized as follows. In Sec. II, we give a brief review of Brill-Horowitz’s construction of their initial data. In Sec. III, we describe numerical method and equations. The results of our simulations are shown in in Sec. IV. Finally, we summarize our results in Sec. V.
## II Brill-Horowitz’s Initial Data
In this section we briefly review Brill and Horowitz’s argument . Let us consider an initial slice with $`K_{ij}=0`$ in five dimensional vacuum space-times, where $`K_{ij}`$ is the extrinsic curvature of a four dimensional spacelike hypersurface. In this slice the Hamiltonian constraint equation becomes $`{}_{}{}^{(4)}R=0`$, where $`{}_{}{}^{(4)}R`$ is the four dimensional Ricci scalar. Here one can easily see that the Euclidean Reissner-Nordstrom metric with imaginary ‘charge’ $`iq`$ satisfies the Hamiltonian constraint equation, because the ‘energy-momentum’ tensor of the four dimensional Maxwell field is traceless. The metric of the hypersurface is given by
$`{}_{}{}^{(4)}g=U(r)d\chi ^2+{\displaystyle \frac{dr^2}{U(r)}}+r^2d\mathrm{\Omega }^2,`$ (4)
where $`U(r)=12m/rq^2/r^2`$ and $`rr_+:=m+\sqrt{m^2+q^2}`$. In the same way as the previous Witten’s example, the metric is approximately written as
$`{}_{}{}^{(4)}g{\displaystyle \frac{4}{U^{}(r_+)}}\left[R^2d\left({\displaystyle \frac{U^{}(r_+)\chi }{2}}\right)^2+dR^2\right]+r_+^2d\mathrm{\Omega }^2,`$ (5)
near $`r=r_+`$, where $`R=\sqrt{rr_+}`$. To avoid a conical singularity at $`r=r_+`$, we assume the period $`\chi _p=4\pi /U^{}(r_+)=2\pi r_+^2/(r_+m)`$ along the $`\chi `$-direction.
The total energy is evaluated as $`E=m/2`$. The $`m`$ is arbitrary parameter and $`q`$ determines the size of the compactified space. So the total energy can be arbitrary negative.
## III Field equations and our numerical method
To know the final fate of the KK bubbles with the negative energy, we study the subsequent time evolution for a long time numerically. We apply 4+1 decomposition of space-time along to the Arnowitt-Deser-Misner formulation for the actual time integrations. We describe equations and basic numerical techniques in this section.
The metric of the full space-time is assumed to be
$`ds_5^2`$ $`=`$ $`N(r,t)^2dt^2+e^{2a(r,t)}U(r)d\chi ^2+e^{2b(r,t)}U(r)^1dr^2+r^2e^{2c(r,t)}d\mathrm{\Omega }^2,`$ (6)
where $`U(r)=12m/rq^2/r^2`$, $`N`$ is the lapse function, and the metric components $`a,b`$ and $`c`$ are now time dependent. The evolution equations of the four-metric $`\gamma _{ij}`$ and the extrinsic curvature $`K_{ij}`$ become Here, for simplicity, we tacitly supposed the boundary condition so that the location of the bubble is ‘fixed’ under the variation of the action. As a result we obtain the 5-dimensional vacuum Einstein equation and can show the consistent result given in Appendix A. Since the Cauchy development of the initial data cannot cover all region outside the bubble, one may be able to consider another boundary conditions, which might be artificial.
$`\dot{K}_j^i`$ $`=`$ $`N({}_{}{}^{(4)}R_{j}^{i}+KK_j^i){}_{}{}^{(4)}D_{}^{i}{}_{}{}^{(4)}D_{j}^{}N,`$ (7)
$`\dot{\gamma }_{ij}`$ $`=`$ $`2NK_{ij},`$ (8)
a dot denotes the time derivative, and $`{}_{}{}^{(4)}R_{j}^{i}`$ and $`{}_{}{}^{(4)}D_{}^{i}`$ denote four dimensional Ricci curvature and the covariant derivative, respectively. For the reader’s convenience, we write down several terms in (7) for the metric (6) as:
$`{}_{}{}^{(4)}R_{\chi }^{\chi }`$ $`=`$ $`e^{2b}\left[\left(a^2a^{\prime \prime }2a^{}c^{}{\displaystyle \frac{2a^{}}{r}}+a^{}b^{}\right)U+\left({\displaystyle \frac{3}{2}}a^{}+{\displaystyle \frac{1}{2}}b^{}c^{}{\displaystyle \frac{1}{r}}\right)U^{}{\displaystyle \frac{1}{2}}U^{\prime \prime }\right],`$ (9)
$`{}_{}{}^{(4)}R_{r}^{r}`$ $`=`$ $`e^{2b}\left[\left(a^2a^{\prime \prime }2c^22c^{\prime \prime }{\displaystyle \frac{4c^{}}{r}}+a^{}b^{}+2b^{}c^{}+{\displaystyle \frac{2b^{}}{r}}\right)U+\left({\displaystyle \frac{3}{2}}a^{}+{\displaystyle \frac{1}{2}}b^{}c^{}{\displaystyle \frac{1}{r}}\right)U^{}{\displaystyle \frac{1}{2}}U^{\prime \prime }\right],`$ (10)
$`{}_{}{}^{(4)}R_{\theta }^{\theta }`$ $`=`$ $`{}_{}{}^{(4)}R_{\phi }^{\phi }=e^{2b}\left[\left(2c^2c^{\prime \prime }{\displaystyle \frac{4c^{}}{r}}{\displaystyle \frac{1}{r^2}}+c^{}b^{}+{\displaystyle \frac{b^{}}{r}}a^{}c^{}{\displaystyle \frac{a^{}}{r}}\right)U+\left(c^{}{\displaystyle \frac{1}{r}}\right)U^{}+{\displaystyle \frac{e^{2b2c}}{r^2}}\right],`$ (11)
and
$`{}_{}{}^{(4)}D_{\chi }^{}{}_{}{}^{(4)}D_{}^{\chi }N`$ $`=`$ $`e^{2b}\left(a^{}+{\displaystyle \frac{1}{2}}{\displaystyle \frac{U^{}}{U}}\right)UN^{},`$ (12)
$`{}_{}{}^{(4)}D_{r}^{}{}_{}{}^{(4)}D_{}^{r}N`$ $`=`$ $`e^{2b}\left[\left(N^{\prime \prime }b^{}N^{}\right)U+{\displaystyle \frac{1}{2}}N^{}U^{}\right],`$ (13)
$`{}_{}{}^{(4)}D_{\theta }^{}{}_{}{}^{(4)}D_{}^{\theta }N`$ $`=`$ $`e^{2b}\left(c^{}+{\displaystyle \frac{1}{r}}\right)N^{}U,`$ (14)
where a dash denotes the derivative on $`r`$.
We start our simulation from the initial data of Brill-Horowitz’s momentarily static solution, such as
$`a(r,0)=b(r,0)=c(r,0)`$ $`=`$ $`0,`$ (15)
$`K_\chi ^\chi (r,0)=K_r^r(r,0)=K_\theta ^\theta (r,0)`$ $`=`$ $`0.`$ (16)
The numerical region is taken as $`r_+rr_e`$, where $`r_+:=m+\sqrt{m^2+q^2}`$ is the location of the bubble at the initial data and $`r_e`$ is the numerical outer boundary. We stress from the construction that the Kaluza-Klein bubble space-time is restricted in $`r_+r\mathrm{}`$. We apply the Robin boundary condition at $`r=r_e`$ such as all the components fall off as they form an asymptotically flat spacetime. At the inner boundary $`r=r_+`$, we use the fact that both $`a`$ and $`b`$ evolve synchronously as we describe in the Appendix, and use both the evolution equation for tr$`K`$,
$$\dot{K}=NK_{ij}K^{ij}{}_{}{}^{(4)}D_{}^{i}{}_{}{}^{(4)}D_{i}^{}N,$$
(17)
where we used the Hamiltonian constraint equation, and the momentum constraint equation,
$${}_{}{}^{(4)}D_{j}^{}K_r^j{}_{}{}^{(4)}D_{r}^{}K=0,$$
(18)
so as the system evolves properly.
In order to specify the lapse function, $`N`$, we apply both the geodesic slicing condition, $`N=1`$, and the maximal slicing condition, $`K=0`$, which equation becomes \[directly from Eq. (17)\]
$${}_{}{}^{(4)}\mathrm{\Delta }N=NK_{ij}K^{ij}.$$
(19)
This elliptic equation is solved using the incomplete Cholsky conjugate gradient method. The outer boundary for $`N`$ is set again as asymptotically flat, and the inner boundary at $`r_+`$ for solving (19) we lineally extrapolate 4-metric components.
We apply Brailovskaya integration scheme (a second order predictor-corrector method) for the time evolution. The accuracy of the calculation is checked by monitoring the violation of the Hamiltonian constraint equation. The numerical code passed convergence tests, and the results shown in this paper are all obtained with acceptable accuracy.
## IV Results
### A Acceleration of the bubble surface
We first check whether our code reveals the initial behavior discussed by Corley and Jacobson . We calculate the area of the bubble,
$$A(t)=4\pi g_{\theta \theta }(r_+,t),$$
(20)
together with its time derivative $`\dot{A}`$, and its acceleration,
$$\ddot{A}=4\pi \ddot{g}_{\theta \theta }=4\pi r^2[2N\dot{K}_\theta ^\theta 2\dot{N}K_\theta ^\theta +4N^2(K_\theta ^\theta )^2]e^{2c}.$$
(21)
We first show this acceleration in Fig.1(a)(b), since this was the quantity discussed by Corley and Jacobson. The Figs. 1(a) and (b) are of the geodesic slicing condition and of the maximal slicing condition, respectively. We fix the charge $`q`$ and varied $`m`$ from negative to positive values. Except for the transition at $`m=0`$, as one can see later, our result is not qualitatively sensitive under changes of $`m/q`$. Under both slicing conditions, we see that the negative mass bubble start expanding (positive $`\ddot{A}`$) initially, yet will soon be in de-accelerating phase (negative $`\ddot{A}`$), while the positive mass bubble keep accelerating all the way in Figs. 1(a), and in the region in Figs. 1(b). More precisely, for the positive mass cases in Figs. 1(b), we observe from the numerical results that the acceleration will reach and stay at a positive value in the final stage, even if it goes negative for a short time, which is happen to quite small positive mass cases.
Such an initial behavior (for both positive and negative mass bubbles) does agree with Corley and Jacobson’s analysis (we remark that their analysis was under the geodesic slicing condition). However, the turning behavior into de-accelarating phase could not find in their analysis. The de-accelarating does not mean collapsing feature directly. Actually, upto we stop our time evolution, the numerical data of the area, (20), monotonically increases \[Figs. 1(c)\], while its velocity goes down for negative mass bubbles \[Figs. 1(d)\]. However, from this facts, we can not say that negative mass bubbles will expand forever, because we can see the blow-ups of the Riemann invariant and collapsing lapse behavior as we show next. (We had to stop time evolution for negative mass bubble case when we face the blow-ups of the Riemann invariant.)
### B Collapse of lapse
Since we found that the time integration using the maximal slicing condition survives long term time evolution than that of the geodesic slicing condition, we will show only the results of the maximal slicing condition hereafter.
The maximal slicing condition is known as a robust gauge condition for singularity avoidance (or, exactly speaking, avoiding the vanishing of the volume elements of the associated Eulerian observers) . This is because the lapse will go quite small value in the strong gravitational field. Contrary, we may guess whether the space-time will collapse or not by monitoring the lapse function.
We plot the lapse function, $`N`$, in Fig.2. Fig.2(a) is the lapse function at the bubble surface, $`r=r_+`$, versus time. We see the lapse evolves small value for the case of negative mass bubble space-time. The lines end at the time when the violation of the constraint equation begin growing. From above standard behavior of the maximal sliced lapse functions, we may say that the negative mass bubble space-time is ‘collapsing’ in some senses. Fig. 2(b) is snapshots of $`N`$ at several time for the case of negative mass bubble space-time.
### C Riemann invariant
In order to confirm our guess of the ‘collapsing’ behavior of the negative mass bubble, we calculated the Kretchman invariant (Riemann invariant) $`R_{ijkl}R^{ijkl}`$ of both 4 and 5-dimensional Riemann tensor. We see both blow up in the cases of negative mass bubbles. We plotted a typical behavior of the invariant as a function of $`r`$ and $`t`$, in Fig.3. These lines suggest that the possibility of the formation of singularity in the final phase of evolution.
### D Apparent horizons
In order to confirm whether a black hole is formed or not in such a case, we check the appearance of the apparent horizon. The location of the apparent horizon is given by the position that the expansion rate, of the outgoing null geodesic congruence turns into the negative. The appearance of the apparent horizon indicates the existence of the event horizon.
The definition of the apparent horizon might not be unique in our five-dimensional space-time because it depends on the dimension of the space-time which the null geodesic congruence runsIn usual Kaluza-Klein picture, it is natural that the null geodesic congruence runs in four dimensional part.. If the null congruence propagating in full five dimensional space-time, the expansion rate, $`{}_{}{}^{(4)}\theta _{+}^{}`$, is given by
$`{}_{}{}^{(4)}\theta _{+}^{}`$ $`=`$ $`{}_{}{}^{(4)}_{a}^{}s^a{}_{}{}^{(4)}K+s^as_{}^{b}{}_{}{}^{(4)}K_{ab}`$ (22)
$`=`$ $`(a^{}+2c^{}+{\displaystyle \frac{2}{r}}+{\displaystyle \frac{U^{}}{2U}})\sqrt{U}e^b(K_\chi ^\chi +K_\theta ^\theta +K_\phi ^\phi )`$ (23)
where $`s^a=(0,1/\sqrt{g_{rr}},0,0)`$ is a outer pointing vector in our spatial four metric. On the other hand, if the null is confined to non-compactified four dimensions, the expansion rate, $`{}_{}{}^{(3)}\theta _{+}^{}`$, is given similarly by
$`{}_{}{}^{(3)}\theta _{+}^{}`$ $`=`$ $`{}_{}{}^{(3)}_{a}^{}s^a{}_{}{}^{(3)}K+s^as_{}^{b}{}_{}{}^{(3)}K_{ab}`$ (24)
$`=`$ $`2(c^{}+{\displaystyle \frac{1}{r}})\sqrt{U}e^b(K_\theta ^\theta +K_\phi ^\phi ).`$ (25)
We analyzed both $`{}_{}{}^{(3)}\theta _{+}^{}`$ and $`{}_{}{}^{(4)}\theta _{+}^{}`$ in our evolving space-time. Surprisingly, in all cases (positive and negative mass bubbles), both expansions remain positive definite everywhere as we show an example in Fig.4. These suggest us no-appearance of apparent horizons.
## V Discussion
We numerically studied the dynamical evolution of the Brill and Horowitz’s initial data which can have the negative energy. As the zero energy Witten’s bubble space-time, we show that the ‘bubbles’ with negative energy will expand by mean of area upto the time we stop the simulations. At first glance this result supports Corley and Jacobson’s conjecture. However, from the facts that the curvature invariant blows up, and no appearance of the apparent horizon, we suspect that a formation of a naked singularity as the final fate of Kaluza-Klein negative energy bubble <sup>§</sup><sup>§</sup>§ An anonymous referee of this article pointed out the similarity of the positive and negative mass bubble results. However, from the results we obtained, we believe that there are qualitative differences between positive and negative mass cases in their dynamical behaviors. (We remark that our simulations are only up to a finite time in order to keep the resolution against the expansion of the spacetime.) . Hence, we may have to consider seriously the decay problem from the Kaluza-Klein vacuum to the Witten-type ‘bubble’ space-time. Possible resolution to this may be given by assuming the supersymmetry which may forbid the decay , or by constructing quantum gravity theory which may smooth out singularities as normally been expected.
Although the negative mass bubbles are expanding, we obtained the result that the bubble spacetime terminates at the singularity. At first glance, they are incompatible, because the naive picture, which the expanding keeps regularity. However, the picture may be based on the Raychaudhri-type equation and the equation does not hold in the present case. Moreover, the area cannot properly describes whether the system will collapse or not. Properly speaking, we need the proper radius from the center which is absent in the present case.
Finally, we would like to comment on the so-called brane world scenario . The brane world is motivated by the reduction from the M-theory to the $`E_8\times E_8`$ heterotic superstring theory. This reduction drastically changes the picture of the reduction because ‘matters’ are confined to the ten dimensions and gravitons are propagating in the eleven-dimensions. Here we call the timelike hypersurface, where matters are confined, by ‘brane’. A plausible history of the compactification are still quite actively discussed recently and we do not reach the consensus at this moment. Apart from this compactification scenario, the reduction from ten to four-dimensions follows the well-known Kaluza-Klein type or Calabi-Yau compactification. Thus, our present analysis is basically applicable to the space-time on the brane, because we supposed the well-known Kaluza-Klein compactification. More precisely, we may say that the Witten-type Kaluza-Klein ‘bubble’ space-time on the 4-brane will be reduced from at least 6 dimensional space-time. As was recently reported , the effective Einstein equations on the brane are different from the normal Einstein equations. Therefore, it might be worth re-asking what is the final fate of Kaluza-Klein bubble if we describe the space-time by such a modified Einstein equation when we take a brane world scenario.
## Acknowledgements
TS is grateful to Gary Gibbons and DAMTP relativity group for their hospitality. HS appreciates Abhay Ashtekar and Pablo Laguna for helpful comments, and CGPG group for their hospitality. HS also thank Ethan Honda for describing his numerical method. Numerical computations were performed using machines at CGPG. Both authors are supported by the Japan Society for the Promotion of Science (JSPS) as research fellows abroad.
## A The dynamical equations at the surface of the bubble
In this appendix, we show that $`\dot{a}=\dot{b}`$ at the location of the bubble, $`r=r_+`$, during the time evolution.
The explicit expression of Eqs (7) and (8) are given by
$`\left({\displaystyle \frac{\dot{a}}{N}}\right)^{}`$ $`=`$ $`Ne^{2b}\left[\left(a^2a^{\prime \prime }2a^{}c^{}{\displaystyle \frac{2a^{}}{r}}+a^{}b^{}\right)U+\left({\displaystyle \frac{3}{2}}a^{}+{\displaystyle \frac{1}{2}}b^{}c^{}{\displaystyle \frac{1}{r}}\right)U^{}{\displaystyle \frac{1}{2}}U^{\prime \prime }\right]`$ (A2)
$`{\displaystyle \frac{1}{N}}(\dot{a}+\dot{b}+2\dot{c})\dot{a}+e^{2b}\left(a^{}U+{\displaystyle \frac{1}{2}}U^{}\right)N^{},`$
$`\left({\displaystyle \frac{\dot{b}}{N}}\right)^{}`$ $`=`$ $`Ne^{2b}\left[\left(a^2a^{\prime \prime }2c^22c^{\prime \prime }{\displaystyle \frac{4c^{}}{r}}+a^{}b^{}+2b^{}c^{}+{\displaystyle \frac{2b^{}}{r}}\right)U+\left({\displaystyle \frac{3}{2}}a^{}+{\displaystyle \frac{1}{2}}b^{}c^{}{\displaystyle \frac{1}{r}}\right)U^{}{\displaystyle \frac{1}{2}}U^{\prime \prime }\right]`$ (A4)
$`{\displaystyle \frac{1}{N}}(\dot{a}+\dot{b}+2\dot{c})\dot{b}+e^{2b}\left[(N^{\prime \prime }b^{}N^{})U+{\displaystyle \frac{1}{2}}N^{}U^{}\right],`$
$`\left({\displaystyle \frac{\dot{c}}{N}}\right)^{}`$ $`=`$ $`Ne^{2b}\left[\left(2c^2c^{\prime \prime }{\displaystyle \frac{4c^{}}{r}}{\displaystyle \frac{1}{r^2}}+c^{}b^{}+{\displaystyle \frac{b^{}}{r}}a^{}c^{}{\displaystyle \frac{a^{}}{r}}\right)U+\left(c^{}{\displaystyle \frac{1}{r}}\right)U^{}+{\displaystyle \frac{e^{2b2c}}{r^2}}\right]`$ (A6)
$`{\displaystyle \frac{1}{N}}(\dot{a}+\dot{b}+2\dot{c})\dot{c}+e^{2b}\left(c^{}+{\displaystyle \frac{1}{r}}\right)N^{}U.`$
At $`r=r_+`$, we can truncate $`U`$, since $`U=0`$. By adding a suffix $`+`$ for the variables which is evaluated at $`r=r_+`$, the above equations become
$`\left({\displaystyle \frac{\dot{a}_+}{N_+}}\right)^{}`$ $`=`$ $`N_+e^{2b_+}\left[\left({\displaystyle \frac{3}{2}}a_+^{}+{\displaystyle \frac{1}{2}}b_+^{}c_+^{}{\displaystyle \frac{1}{r_+}}\right)U_+^{}{\displaystyle \frac{1}{2}}U_+^{\prime \prime }\right]`$ (A8)
$`{\displaystyle \frac{1}{N_+}}(\dot{a}_++\dot{b}_++2\dot{c}_+)\dot{a}_++{\displaystyle \frac{1}{2}}e^{2b_+}U_+^{}N_+^{},`$
$`\left({\displaystyle \frac{\dot{b}_+}{N_+}}\right)^{}`$ $`=`$ $`N_+e^{2b_+}\left[\left({\displaystyle \frac{3}{2}}a_+^{}+{\displaystyle \frac{1}{2}}b_+^{}c_+^{}{\displaystyle \frac{1}{r_+}}\right)U_+^{}{\displaystyle \frac{1}{2}}U_+^{\prime \prime }\right]`$ (A10)
$`{\displaystyle \frac{1}{N_+}}(\dot{a}_++\dot{b}_++2\dot{c}_+)\dot{b}_++{\displaystyle \frac{1}{2}}e^{2b_+}U_+^{}N_+^{}.`$
Subtracting Eq. (A10) from Eq. (A8), we obtain
$`\left({\displaystyle \frac{\dot{a}_+\dot{b}_+}{N_+}}\right)^{}=(\dot{a}_++\dot{b}_++2\dot{c}_+){\displaystyle \frac{\dot{a}_+\dot{b}_+}{N_+}}.`$ (A11)
Therefore we get
$`{\displaystyle \frac{\dot{a}_+\dot{b}_+}{N_+}}=Ae^{(a_++b_++2c_+)},`$ (A12)
where $`A`$ is a constant. Since the initial conditions, (15) and (16), imply $`a_+=b_+`$, which implies $`A=0`$. Therefore at the boundary, $`r=r_+`$, we can set $`a_+=b_+`$ even after the long time integration.
Figure Captions
Fig.1
Acceleration of the bubble surface, (21), versus time. The figures (a) and (b) are of the geodesic slicing condition and of the maximal slicing condition, respectively. In both figures, we see that the negative mass bubble will soon be in collapse phase, although they start expanding initially. We set $`q=1`$ (hereafter for all figures).
For the case of evolutions with the maximum slicing condition, we plot (c) area of the bubble surface $`A`$, (20), versus time, and (d) the velocity of the bubble surface, $`dA/dt`$.
Fig.2
The lapse function, $`N`$, are plotted (the solutions of maximal slicing condition). The figure (a) is $`N`$ at $`r=r_+`$ versus time. We see the lapse evolves small value for the case of negative mass bubble space-time. The figure (b) is snapshots of $`N`$ at several times for a case of negative mass bubble ($`m=0.4`$) space-time.
Fig.3
(a) Typical snapshots of the Riemann invariant $`R_{ijkl}R^{ijkl}`$ of 4-dimensional Riemann curvature for the case of negative mass bubble ($`m=0.4`$) are plotted. Only the region near the bubble surface is drawn. (b) The Riemann invariant $`R_{ijkl}R^{ijkl}`$ of 5-dimensional Riemann curvature are plotted as a function of time. We see blow-ups in the cases of negative mass bubbles (we cut the display range at $`10^5`$). The values are evaluated at a point right from the bubble surface (that is, at $`r_++\mathrm{\Delta }r`$).
Fig.4
A typical sample of the outgoing null expansion rate $`{}_{}{}^{(3)}\theta _{+}^{}`$ and $`{}_{}{}^{(4)}\theta _{+}^{}`$ are plotted for the case of negative mass ($`m=0.4`$) bubbles.
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# References
Fermi Gas in Harmonic oscillator potentials
X.X. Yi<sup>1,3</sup> J.C.Su<sup>2</sup>
<sup>1</sup>Institute of Theoretical Physics, Northeast Normal University,
Changchun 130024, China
<sup>2</sup>Department of Physics, Jilin University,
Changchun 130023, China
<sup>3</sup>Institute of Theoretical Physics, Academia Sinica, Peking 100080, China
Abstract
Assuming the validity of grand canonical statistics, we study the properties of a spin-polarized Fermi gas in harmonic traps. Universal forms of Fermi temperature $`T_F`$, internal energy $`U`$ and the specific heat per particle of the trapped Fermi gas are calculated as a function of particle number, and the results compared with those of infinite number particles.
PACS numbers:03.75.Fi,05.30.Fk
The ideal Fermi gas is an old and well-understood problem since the non-interaction Fermi gas is a good zeroth-order approximation for many familiar systems. Like the trapped degenerate atomic gases \[1-3\] that provide exciting opportunities for the manipulation and quantitative study of quantum statistical effects, the behavior of trapped Fermi gas also merits attention, both as a degenerate quantum system in its own right and as a possible precursor to a paired Fermi condensate at low temperature though it perhaps is not as dramatic as the phase transition associated with bosons.
The trapped atomic gases reported in Ref.\[1-3\] are dilute. The effects of predominantly short-range atom-atom interactions are therefore week. For dilute spin-polarized Fermi gases, the s-wave scattering amplitude vanishes due to the antisymmetry of the many-Fermi wave function, and the p-wave scattering is small at low energy. Therefore the dilute Fermi gas is usually treated as an ideal Fermi gas.
Unlike the usual system in most textbooks of statistical mechanics, e.g., however, the number of trapped atoms in Ref.\[1-3\] is finite (about $`10^9`$). For finite number of trapped fermions, the applicability of the usual thermodynamical calculations and the neglect of the ground energy are no longer available. For these relatively low numbers (compared to $`10^{23}`$), the effects caused by the above approximation are nonvanishing.
Harmonic traps provide a particularly simple realization of the confined Fermi system. In this trap, Butts and Rokhsar calculated the spatial and momentum distributions for trapped fermions using Thomas-Fermi approximation. The spatial distribution of the trapped cloud provides an explicit visualization of a real-space ”Fermi sea” and the momentum distribution, unlike the spatial one, is isotropic.
In this paper, we calculate the Fermi temperature, internal energy and the specific heat per particle of the finite number trapped Fermi gas. We indeed find marked differences from the usual treatments: a correction to the Fermi temperature, internal energy and the specific heat per particle of the finite number trapped fermions, a constraint to the number of trapped fermions for given trap frequencies.
Consider $`N`$ spin-polarized fermions of mass $`m`$ moving in anisotropic oscillator potentials. The single-particle levels are familiar
$$E_N=E(n_x,n_y,n_z)=\underset{i=x,y,z}{}\mathrm{}\omega _i(n_i+\frac{1}{2}),$$
(1)
where $`\omega _i(i=x,y,z)`$is the trap frequencies in the $`i`$ direction, $`n_i`$ is non-negative integer. Note that each oscillator state is assumed to be filled with a single fermion, since only one spin orientation is confined by the magnetic trap. Under real experimental conditions of trapped atomic gases, the temperature is high on the scale of the trap level spacing, namely, $`k_BT>>\mathrm{}\omega _i(i=x,y,z).`$ Therefore, within the canonical ensemble, the partition function
$$Q(\beta )=\underset{N}{}e^{\beta E_N}=\underset{n_x,n_y,n_z=0}{\overset{\mathrm{}}{}}e^{\beta (n_x\omega _x+n_y\omega _y+n_z\omega _z)}=\underset{i=x,y,z}{}\frac{1}{1e^{\beta \omega _i}}$$
(2)
of such a system without interactions can be expanded as follows:
$$Q(\beta )a_2\beta ^3+a_1\beta ^2+a_0\beta ^1+a_1+O(\beta )$$
(3)
with
$`a_2`$ $`=`$ $`{\displaystyle \frac{1}{\omega _x\omega _y\omega _z}}`$
$`a_1`$ $`=`$ $`{\displaystyle \frac{1}{2}}a_2(\omega _x+\omega _y+\omega _z)`$ (4)
$`a_0`$ $`=`$ $`{\displaystyle \frac{1}{12}}a_2(\omega _x^2+\omega _y^2+\omega _z^2+3\omega _x\omega _y+3\omega _y\omega _z+3\omega _z\omega _z)`$
$`a_1`$ $`=`$ $`{\displaystyle \frac{1}{8}}+{\displaystyle \frac{1}{24}}(\omega _x/\omega _y+\omega _x/\omega _z+\omega _y/\omega _z+\omega _y/\omega _x+\omega _z/\omega _x+\omega _z/\omega _y),`$ (5)
where, $`\beta =k_BT`$, $`k_B`$ is the Boltzmann constant, and the contribution of the ground state was singled away for special treatment. On the other hand, by making use of
$$Q(\beta )=\underset{N}{}e^{\beta E_N}=_{\epsilon _0}^{\mathrm{}}e^{\beta E}\rho (E)𝑑E$$
(6)
the partition function can be calculated equivalently, where $`\rho (E)`$ is the density of states, $`\epsilon _0`$ stands for the energy of the ground state. Comparing with eq.(3), it is proved that the density of states $`\rho (E)`$ takes a form:
$$\rho (E)=b_0+b_1E+b_2E^2+\mathrm{}.$$
(7)
This expansion can’t contain the terms $`E^r`$ with $`r`$ being negative or non-integer, since they become zero after comparison with the direct calculation given by eq.(3). Substitution eq.(7) into eq.(6), one can easily find:
$$Q(\beta )=(b_0+b_1\epsilon _0+b_2\epsilon _0^2)\beta ^1+(b_12b_2\epsilon _0)\beta ^2+2b_2\beta ^3+\mathrm{}$$
(8)
comparison of eq.(8) with the direct calculation(3) shows that
$$b_0=a_0\epsilon _0a_1\epsilon _0^2a_2,b_1=a_1+\epsilon _0a_2,b_2=0.5a_2$$
(9)
There are two additional terms $`b_0`$ and $`b_1E`$ in the density of states in comparison with the result of Butts, as we see in following, these additional terms will result in the shift of the thermodynamic quantities.
Now, let us consider a system of $`N`$ noninteracting fermions that the population $`N(E_i)`$ of a state with energy $`E_i`$ is given by the Fermi-Dirac distribution
$$N(E_i)=\frac{1}{e^{\beta (E_i\mu )}+1}$$
(10)
Here, we set the statistical weights corresponding the state $`E_i`$, $`g_i=1`$. $`\mu `$ stands for the chemical potential, which is determined by the constraint that the total number of the particles in the system is $`N`$:
$$N=\underset{i}{\overset{\mathrm{}}{}}N(E_i).$$
(11)
At zero temperature the Fermi-Dirac distribution factor is unity for energies less than the Fermi energy $`E_F=\mu (T=0,N)`$, and zero otherwise. A straightforward integration of Eq.(11) gives:
$$E_F^3+\frac{3b_1}{2b_2}E_F^2+3\frac{b_0}{b_2}E_F\frac{3}{b_2}N^{^{}}=0$$
(12)
with $`N^{^{}}=N+\frac{1}{3}b_2\epsilon _0^2+\frac{1}{2}b_1\epsilon _0^2+b_0\epsilon _0`$. The terms with $`b_0`$ and $`b_1`$ are corrections from finite number effects to Fermi energy. In the large number particle limit, $`E_F=(\frac{3}{b_2}N^{^{}})^{1/3}`$. This just is the result of Ref.. Some words of caution are now in order. Fermi energy $`E_F`$ is determined by eq.(12), i.e. $`E_F`$ is a solution of eq.(12). From the mathematical point of view, the eq.(12) has a only real solution
$$E_F(\frac{3}{b_2}N^{^{}})^{1/3}=\frac{1}{4}\frac{b_1}{b_2}\frac{1}{2}\frac{b_0}{b_2}$$
(13)
in the situation of the number of trapped fermior satisfying
$$N^{^{}}<N_{max},N_{max}=[(\frac{p}{3})^{3/2}+q_0]b_0/3,$$
(14)
where
$$p=\frac{3b_1^2}{4b_2^2}+3\frac{b_0}{b_2},q_0=\frac{b_1^3}{4b_2^3}\frac{3b_0b_1}{2b_2^2}.$$
(15)
Otherwise, the eq.(12) have three real solution, but only the positive smallest one meets the requirements for Fermi energy from the physical point of view. We would like to point out that the results presented here are valid for $`k_BT>>\mathrm{}\omega _i(i=x,y,z).`$ If this condition is broken, we should make use of a numerical method to study the effects of finite particles.
For finite temperature, neglecting the effect of zero-point energy, the eq.(11) gives:
$$N=b_2k_B^3T^3f_3(z)+b_1k_B^2T^2f_2(z)+b_0k_BTf_1(z)$$
(16)
where $`f_n(z)`$ stands for the Fermi integral, which is given that
$$f_n(z)=\frac{1}{\mathrm{\Gamma }(n)}_0^{\mathrm{}}\frac{x^{n1}}{z^1x^n+1}𝑑x,\mathrm{\Gamma }(n)=_0^{\mathrm{}}e^xx^{n1}𝑑x.$$
$`z`$ is a fugacity. The last two terms in eq.(16) are due to the effects of finite particles. If one introduces a temperature
$$T_F^0=\frac{1}{k_B}(\frac{N}{b_2f_3(z_F)})^{1/3},$$
(17)
which denotes the Fermi temperature of infinite number fermions trapped in anisotropic oscillator potentials, i.e. $`k_BT_F^0=\mu `$. Then the Fermi temperature $`T_F`$ defined as the same as in most textbooks for finite number fermions in the case of $`\epsilon _0=0`$ takes a form
$$T_F=T_F^0[1\frac{1}{3}\frac{b_1}{b_2}\frac{1}{k_BT_F^0}\frac{f_2(z_F)}{f_3(z_F)}\frac{2}{3}\frac{b_0}{b_2}\frac{1}{(k_B^2T_F^0)^2}\frac{f_1(z_F)}{f_3(z_F)}],$$
(18)
where $`z_F=z|_{T=T_F^0}`$. The last two terms in eq.(18) are correction to Fermi energy from the finite number effect, it shows that the Fermi temperature for finite trapped fermions is lower than that in the case of infinite trapped fermions.
Generally speaking, the specific heat is more interest from the experimental point of view, since the low-temperature behavior of specific heat $`c`$ is generally treated as the hallmark of onset of phase transition . It is well known that the specific heat can be derived from the internal energy, which is expressed by
$$\beta U=\beta _0^{\mathrm{}}\frac{E\rho (E)}{exp[\beta (E\epsilon _0)]+1}𝑑E$$
(19)
Here we dropped the ground state energy, which is less than $`2\epsilon _0`$ and remains constant. Substituting eq.(7) into eq.(19), one easily finds
$$U=\frac{b_2}{\beta ^4}f_4(z)+\frac{b_1}{\beta ^3}f_3(z)+\frac{b_0}{\beta ^2}f_2(z),$$
(20)
the first term remains unchanged in the limit of very large number particles, while the last two terms close to zero. The specific heat per particle of the trapped fermions is given that
$$c=[\frac{1}{N}\frac{U}{T}]_N,$$
(21)
straightforward calculation gives
$`c`$ $`=`$ $`{\displaystyle \frac{1}{N}}\{4b_2k_B^4T^3f_4(z)+3b_1k_B^3T^2f_3(z)+2b_0k_B^2Tf_2(z)`$ (22)
$`{\displaystyle \frac{3}{2}}b_2k_B^4T^3{\displaystyle \frac{f_{3/2}(z)f_3(z)}{f_{1/2}(z)}}{\displaystyle \frac{3}{2}}b_1k_B^3T^2{\displaystyle \frac{f_{3/2}(z)f_2(z)}{f_{1/2}(z)}}{\displaystyle \frac{3}{2}}b_0k_B^2T{\displaystyle \frac{f_{3/2}(z)f_1(z)}{f_{1/2}(z)}}\}`$
It is evident that the specific heat per particle has increased due to the finite particle effects. The results are illustrated in Fig.1. As state in most textbooks of statistical mechanics, for ideal Fermi gas in a 3D box, the specific heat increases proportionally with temperature $`T`$ for $`T<<T_F`$, whereas it closes to a constant $`1.5k_B`$ for $`T\mathrm{}`$. For trapped fermions, however, the specific heat increases predominantly with $`T`$ for $`T<<T_F`$, and close to $`3k_B`$ with $`T\mathrm{}`$. Physically, the difference between the case of 3D box and that of trapped fermions is the confinment.This leads to the difference in their energy levels. Consequently, different energy levels would result in different dependence of specific heat on temperature.
In conclusion, we have discussed the behaviors of fermions trapped in an anisotropic oscillator potentials, it was shown that corrections due to the effect of finite fermions and the ground state energy are small, whether they are observable depends on the experimental parameters. We believe the results will be useful to them who is working experimentally in this area.
ACKNOWLEDGMENTS:
One of us(X.X.Yi) would like to thank prof. C.P.Sun for his helpful discussions.
Fig.1 The specific heat per particle vs. temperature. The dots show the situation for a very large number of particles($`N=10^{23}`$), and the dashed line that for a smaller number($`10^8`$). The trap frequencies are $`\omega _x=500Hz,\omega _y=600Hz,\omega _z=800Hz.`$
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# Interacting six-dimensional topological field theories
## 1 Introduction
Recently, L. Baulieu and P. West introduced a six-dimensional topological model of Witten-type involving $`2`$-form potentials . In the sequel, the gauge-fixing procedure and twist in this model have been studied in more detail in reference ; these authors also determined vector supersymmetry (VSUSY-) transformations which represent an additional symmetry of the model.
The goal of the present paper is to discuss two different generalizations of the free Abelian model to interacting models. The first consists of coupling the Abelian $`2`$-form potentials to a non-Abelian Yang-Mills field by virtue of a Chern-Simons term, as suggested in reference . The corresponding action represents a six-dimensional topological version of the Chapline-Manton term appearing in the action for ten-dimensional supergravity coupled to super Yang-Mills theory . The second generalization consists of considering non-Abelian (charged) $`2`$-form potentials which are coupled to a Yang-Mills connection, following the lines of reference .
Before discussing these generalizations, we summarize the free Abelian model while using differential forms to simplify the notation (section 2). Our paper concludes with some comments concerning possible extensions of these six-dimensional topological models. We note that all of our considerations concern the classical theory (tree-level).
## 2 Free Abelian model
The arena is a compact pseudo-Riemannian $`6`$-manifold $`_6`$ and the basic fields are Abelian $`2`$-form potentials $`B_2`$ and $`B_2^c`$ which are independent of each other. From the associated curvature $`3`$-forms $`G_3=dB_2`$ and $`G_3^c=dB_2^c`$, one can construct the classical action
$$\mathrm{\Sigma }_{cl}=G_3G_3^c.$$
(1)
Here and in the following, the integrals are understood as integrals of $`6`$-forms over $`_6`$ and the wedge product symbol is always omitted.
### 2.1 Symmetries
The action (1) is invariant under the ordinary gauge transformations
$$\delta _{\lambda _1}B_2=d\lambda _1,\delta _{\lambda _1}B_2^c=0,$$
(2)
which represent a reducible symmetry in the present case, and it is invariant under the shift- (or topological Q-) symmetry
$$\delta _{\lambda _2}B_2=\lambda _2,\delta _{\lambda _2}B_2^c=0,$$
(3)
which also represents a reducible symmetry. In equations (2),(3) and in the following, it is understood that the “$`c`$-conjugated” equations also hold, e.g. $`c`$-conjugation of equation (2) gives $`\delta _{\lambda _1^c}B_2^c=d\lambda _1^c,\delta _{\lambda _1^c}B_2=0`$.
In the sequel, we will describe the infinitesimal symmetries in a BRST-framework and we will derive the BRST-transformations from a horizontality condition (Russian formula) . Thus, we introduce a series of ghost fields associated with the reducible gauge transformations (2) and we collect them in a generalized $`2`$-form,
$$\stackrel{~}{B}_2=B_2+V_1^1+m^2.$$
(4)
Here, the upper and lower indices denote the ghost-number and form degree, respectively. The total degree of a field is the sum of its ghost-number and form degree, and all commutators $`[,]`$ are graded with respect to this total degree. The BRST-differential $`s`$, which describes both the ordinary gauge transformations and the shift-transformations, is combined with the exterior derivative $`d`$ in a single operator,
$`\stackrel{~}{d}=d+s,`$ (5)
which is nilpotent by virtue of the relations $`d^2=s^2=[s,d]=0`$. Thus, the generalized field strength
$$\stackrel{~}{G}_3\stackrel{~}{d}\stackrel{~}{B}_2$$
(6)
satisfies the generalized Bianchi identity
$$\stackrel{~}{d}\stackrel{~}{G}_3=0.$$
(7)
The BRST-transformations of the classical and ghost fields<sup>1</sup><sup>1</sup>1Here, “classical” fields are not opposed to quantum fields, but simply refer to the fields appearing in the classical action. are now obtained from the horizontality condition
$$\stackrel{~}{G}_3=G_3+\psi _2^1+\phi _1^2+\varphi ^3,$$
(8)
which involves a series of ghosts associated with the shift-symmetry. By inserting the field expansions (4) and (8) into relations (6) and (7), we obtain the $`s`$-variations
$$\begin{array}{ccccccc}\hfill sB_2& =& \psi _2^1dV_1^1\hfill & ,& \hfill s\psi _2^1& =& d\phi _1^2\hfill \\ \hfill sV_1^1& =& \phi _1^2dm^2\hfill & ,& \hfill s\phi _1^2& =& d\varphi ^3\hfill \\ \hfill sm^2& =& \varphi ^3\hfill & ,& \hfill s\varphi ^3& =& 0\hfill \end{array}$$
and $`sG_3=d\psi _2^1`$. Since the field expansions (4) and (8) have not been truncated, the obtained BRST-transformations are nilpotent by construction .
### 2.2 Gauge-fixing
Let us briefly review the gauge-fixing procedure while using differential forms. To start with, we consider the shift degrees of freedom for the fields $`B_2`$ and $`B_2^c`$. These are fixed by imposing a self-duality condition relating the corresponding field strengths:
$$G_3=G_3^c.$$
(9)
This relation is equivalent to imposing a self-duality condition on $`G_{}d(B_2B_2^c)`$ and an anti-self-duality condition on $`G_+d(B_2+B_2^c)`$. Henceforth, relation (9) is analogous to the self-duality condition for the curvature $`2`$-form $`F`$ in four-dimensional topological Yang-Mills theory .
With the help of a BRST-doublet $`(\chi _3^1,H_3)`$, i.e.
$$s\chi _3^1=H_3,sH_3=0,$$
(10)
the constraint (9) can be implemented in the gauge-fixing action:
$$\mathrm{\Sigma }_{sd}=s\{\chi _3^1(G_3+G_3^c)\}.$$
(11)
Since the shift-symmetry represents a reducible symmetry, it is necessary to re-iterate the gauge-fixing procedure for the action $`\mathrm{\Sigma }_{cl}+\mathrm{\Sigma }_{sd}`$: this leads to the introduction of the anti-ghosts and Lagrange multipliers of tables 2 and 2, all of which have been arranged in Batalin-Vilkovisky pyramids. These fields again represent BRST-doublets:
$$\begin{array}{ccccccc}\hfill s\varphi ^3& =& \eta ^2\hfill & ,& \hfill s\eta ^2& =& 0\hfill \\ \hfill s\phi _1^2& =& \eta _1^1\hfill & ,& \hfill s\eta _1^1& =& 0\hfill \\ \hfill s\chi ^1& =& \eta ^2\hfill & ,& \hfill s\eta ^2& =& 0.\hfill \end{array}$$
(12)
In summary, the shift-invariance of the classical action is fixed by virtue of the gauge-fixing action
$$\mathrm{\Sigma }_Q=\mathrm{\Sigma }_{sd}+s\left\{\phi _1^2d\psi _2^1+\varphi ^3d\phi _1^2+\chi ^1d\phi _1^2CC\right\},$$
(13)
where $`CC`$ stands for the $`c`$-conjugated expressions.
The reducible gauge symmetry (2) is fixed in a similar way: one considers the usual gauge condition $`dB_2=0`$ and re-iterates the gauge-fixing procedure. This leads to the introduction of the series of anti-ghosts and multipliers presented in tables 4 and 4, the $`s`$-variations being given by
$$\begin{array}{ccccccc}\hfill sm^2& =& \beta ^1\hfill & ,& \hfill s\beta ^1& =& 0\hfill \\ \hfill sV_1^1& =& b_1\hfill & ,& \hfill sb_1& =& 0\hfill \\ \hfill sn& =& \beta ^1\hfill & ,& \hfill s\beta ^1& =& 0.\hfill \end{array}$$
(14)
Thus, the ordinary gauge degrees of freedom of $`B_2`$ are fixed by the functional
$$\mathrm{\Sigma }_{og}=s\{V_1^1dB_2+m^2dV_1^1+ndV_1^1CC\}$$
(15)
and the complete gauge-fixed action of the model reads
$$\mathrm{\Sigma }=\mathrm{\Sigma }_{cl}+\mathrm{\Sigma }_Q+\mathrm{\Sigma }_{og}.$$
(16)
### 2.3 Vector supersymmetry
Due to the fact that we are considering a topological model of Witten-type, one expects the complete gauge-fixed action to admit a VSUSY. At the infinitesimal level, the VSUSY-transformations are described by the operator $`\delta _\tau `$ where $`\tau \tau ^\mu _\mu `$ is a constant, $`s`$-invariant vector field of ghost-number zero<sup>2</sup><sup>2</sup>2In order to avoid technical complications related to the global geometry, we limit the considerations of this section to flat space-time.. The variation $`\delta _\tau `$ acts as an antiderivation which lowers the ghost-number by one unit and which anticommutes with $`d`$. The operators $`s`$ and $`\delta _\tau `$ satisfy a graded algebra of Wess-Zumino type,
$$[s,\delta _\tau ]=_\tau ,$$
(17)
where $`_\tau [i_\tau ,d]`$ denotes the Lie derivative along the vector field $`\tau `$ and $`i_\tau `$ the interior product by $`\tau `$. We will simply refer to the relation (17) as the SUSY-algebra.
The $`\delta _\tau `$-variations of all fields can be determined by applying the general procedure introduced in reference . To start with, we derive the VSUSY-transformations of the classical and ghost fields by expanding the so-called $`0`$-type symmetry conditions
$$\delta _\tau \stackrel{~}{B}_2=0,\delta _\tau \stackrel{~}{G}_3=_\tau \stackrel{~}{B}_2.$$
(18)
with respect to the ghost-number. We thus obtain
$$\begin{array}{ccccccc}\hfill \delta _\tau B_2& =& 0\hfill & ,& \hfill \delta _\tau \psi _2^1& =& _\tau B_2\hfill \\ \hfill \delta _\tau V_1^1& =& 0\hfill & ,& \hfill \delta _\tau \phi _1^2& =& _\tau V_1^1\hfill \\ \hfill \delta _\tau m^2& =& 0\hfill & ,& \hfill \delta _\tau \varphi ^3& =& _\tau m^2.\hfill \end{array}$$
(19)
The $`\delta _\tau `$-variations of the anti-ghosts are found by requiring the $`\delta _\tau `$-invariance of the total action (16) and by applying the commutation relations (17). Finally, the VSUSY-transformations of all multipliers follow from the ones of the corresponding anti-ghosts by imposing the algebra (17) for all of them:
$$\begin{array}{ccccccc}\hfill \delta _\tau \chi _3^1& =& 0\hfill & ,& \hfill \delta _\tau H_3& =& _\tau \chi _3^1\hfill \\ \hfill \delta _\tau \varphi ^3& =& 0\hfill & ,& \hfill \delta _\tau \eta ^2& =& _\tau \varphi ^3\hfill \\ \hfill \delta _\tau \phi _1^2& =& 0\hfill & ,& \hfill \delta _\tau \eta _1^1& =& _\tau \phi _1^2\hfill \\ \hfill \delta _\tau \chi ^1& =& _\tau n\hfill & ,& \hfill \delta _\tau \eta ^2& =& _\tau \chi ^1+_\tau \beta ^1\hfill \\ \hfill \delta _\tau m^2& =& _\tau \varphi ^3\hfill & ,& \hfill \delta _\tau \beta ^1& =& _\tau m^2_\tau \eta ^2\hfill \\ \hfill \delta _\tau V_1^1& =& _\tau \phi _1^2\hfill & ,& \hfill \delta _\tau b_1& =& _\tau V_1^1_\tau \eta _1^1\hfill \\ \hfill \delta _\tau n& =& 0\hfill & ,& \hfill \delta _\tau \beta ^1& =& _\tau n.\hfill \end{array}$$
(20)
Thus, it is by construction that the total action is $`\delta _\tau `$-invariant and that the SUSY-algebra is fulfilled off-shell for all fields of the model. Our results coincide with those found in reference by other methods. We refer to the latter work for the relation of VSUSY to a twist of a supersymmetric field theory.
## 3 Abelian model with Chern-Simons term
The authors of reference considered the interaction of the Abelian $`2`$-forms $`B_2`$ and $`B_2^c`$ with a non-Abelian Yang-Mills (YM) connection $`A`$ by virtue of a Chern-Simons term with coupling constant $`\lambda `$,
$$\mathrm{\Omega }_3(A)=\lambda \mathrm{tr}(AdA+\frac{2}{3}AAA).$$
(21)
The proposed action reads<sup>3</sup><sup>3</sup>3We do not include the ordinary YM-action $`\mathrm{tr}(FF)`$ as in reference since it depends on the metric and therefore destroys the topological nature of the model.
$$\widehat{\mathrm{\Sigma }}_{cl}=\left(G_3\mathrm{\Omega }_3\right)\left(G_3^c\mathrm{\Omega }_3\right).$$
(22)
This functional represents a six-dimensional topological version of the expression $`\mathrm{tr}(G_3\mathrm{\Omega }_3)(G_3\mathrm{\Omega }_3)`$ which appears in the action for ten-dimensional supergravity coupled to super YM .
The equations of motion for $`A`$ and $`B_2`$ (or $`B_2^c`$) have the form
$$F(G_3G_3^c)=0\mathrm{and}\mathrm{tr}(FF)=0,$$
where $`F=dA+\frac{1}{2}[A,A]`$ denotes the curvature $`2`$-form associated to $`A`$. The latter equations imply $`F=0`$ (i.e. the same equation of motion as in the three-dimensional Chern-Simons theory).
### 3.1 Symmetries
The action (22) is not anymore invariant under the shift $`\delta B_2=\psi _2^1`$. However, it is invariant under the YM-gauge transformations
$`\delta A`$ $`=`$ $`Dc(dc+[A,c])`$
$`\delta B_2`$ $`=`$ $`\lambda \mathrm{tr}(cdA)=\delta B_2^c,`$ (23)
which leave $`G_3\mathrm{\Omega }_3`$ and $`G_3^c\mathrm{\Omega }_3`$ invariant.
The BRST-transformations of $`A`$ and of the YM-ghost read
$$sA=Dc,sc=\frac{1}{2}[c,c].$$
(24)
They follow from the horizontality condition $`\stackrel{~}{F}=F`$, where
$$\stackrel{~}{F}=\stackrel{~}{d}\stackrel{~}{A}+\frac{1}{2}[\stackrel{~}{A},\stackrel{~}{A}],\stackrel{~}{A}=A+c.$$
(25)
The generalized Chern-Simons form
$$\stackrel{~}{\mathrm{\Omega }}_3=\lambda \mathrm{tr}(\stackrel{~}{A}\stackrel{~}{d}\stackrel{~}{A}+\frac{2}{3}\stackrel{~}{A}\stackrel{~}{A}\stackrel{~}{A})$$
can be expanded with respect to the ghost-number,
$$\stackrel{~}{\mathrm{\Omega }}_3=\mathrm{\Omega }_3+\mathrm{\Omega }_2^1+\mathrm{\Omega }_1^2+\mathrm{\Omega }^3,$$
(26)
which provides the well-known solution of the descent equations (e.g. see )
$$\begin{array}{ccccccc}\hfill \mathrm{\Omega }_3& =& \lambda \mathrm{tr}(AdA+\frac{2}{3}AAA)\hfill & ,& \hfill s\mathrm{\Omega }_3+d\mathrm{\Omega }_2^1& =& 0\hfill \\ \hfill \mathrm{\Omega }_2^1& =& \lambda \mathrm{tr}(cdA)\hfill & ,& \hfill s\mathrm{\Omega }_2^1+d\mathrm{\Omega }_1^2& =& 0\hfill \\ \hfill \mathrm{\Omega }_1^2& =& \lambda \mathrm{tr}(ccA)\hfill & ,& \hfill s\mathrm{\Omega }_1^2+d\mathrm{\Omega }^3& =& 0\hfill \\ \hfill \mathrm{\Omega }^3& =& \lambda \mathrm{tr}(\frac{1}{3}ccc)\hfill & ,& \hfill s\mathrm{\Omega }^3& =& 0.\hfill \end{array}$$
(27)
We now use this result to discuss the $`B`$-field sector. The generalized field strength of $`B_2`$ is defined as before (i.e. $`\stackrel{~}{G}_3=\stackrel{~}{d}\stackrel{~}{B}_2`$ with $`\stackrel{~}{B}_2=B_2+V_1^1+m^2`$), but the horizontality condition (8) of the free model is now replaced by the horizontality condition
$$\stackrel{~}{G}_3=G_3+\mathrm{\Omega }_2^1+\mathrm{\Omega }_1^2+\mathrm{\Omega }^3.$$
(28)
Expansion with respect to the ghost-number yields the $`s`$-variations
$`sB_2`$ $`=`$ $`dV_1^1+\mathrm{\Omega }_2^1`$
$`sV_1^1`$ $`=`$ $`dm^2+\mathrm{\Omega }_1^2`$
$`sm^2`$ $`=`$ $`\mathrm{\Omega }^3,`$ (29)
where the explicit expressions for the $`\mathrm{\Omega }_p^q(A,c)`$ were given in equations (27). Furthermore, substitution of (28) in $`\stackrel{~}{d}\stackrel{~}{G}_3=0`$ leads to $`sG_3=d\mathrm{\Omega }_2^1`$ (and reproduces some of the descent equations (27)).
The BRST-transformations (24) and (29) leave the classical action (22) invariant.
### 3.2 Gauge-fixing
In the YM-sector, the gauge symmetry is fixed in the standard way,
$$\mathrm{\Sigma }_{gf}^A=s\mathrm{tr}\left\{\overline{c}dA\right\}=\mathrm{tr}\left\{bdA\overline{c}dDc\right\},$$
(30)
where we made use of a BRST-doublet $`(s\overline{c}=b,sb=0)`$.
In the $`B`$-sector, the local symmetry $`\delta B_2=dV_1^1`$ is fixed as for the free model, i.e. by introducing the gauge-fixing functional $`\mathrm{\Sigma }_{gf}^B\mathrm{\Sigma }_{og}`$ given by equation (15).
In summary, the complete action of the interacting model reads $`\mathrm{\Sigma }_{int}=\widehat{\mathrm{\Sigma }}_{cl}+\mathrm{\Sigma }_{gf}^B+\mathrm{\Sigma }_{gf}^A`$.
### 3.3 VSUSY
Due to the absence of shift-symmetries in the present model, the only possible choice for VSUSY-transformations is given by the so-called $`\mathrm{}`$-type symmetry conditions . However, the derivation of $`\delta _\tau `$-variations for all fields is substantially more complicated in the present case since the SUSY-algebra only closes on-shell. Therefore, we will not further elaborate on this point here .
## 4 Non-Abelian model
Consider a YM-connection $`A`$ and a $`2`$-form potential $`B_2`$, both with values in a given Lie algebra. The field strength of $`B_2`$ is now defined by
$$G_3=DB_2dB_2+[A,B_2]$$
(31)
and it satisfies the second Bianchi identity $`DG_3=[F,B_2]`$, where $`F=dA+\frac{1}{2}[A,A]`$ denotes the YM-curvature.
A natural generalization of the action (1) for the Abelian potentials is given by
$`\mathrm{\Sigma }_{cl}`$ $`=`$ $`{\displaystyle \mathrm{tr}\left\{G_3G_3^cF[B_2,B_2^c]\right\}}.`$ (32)
Neither $`B_2`$ nor $`A`$ propagate in this model (very much like $`A`$ in four-dimensional topological YM-theory).
### 4.1 Symmetries
Following reference , we now spell out all local symmetries of the functional (32). As in the previously discussed models, one considers the generalized gauge fields
$$\stackrel{~}{A}=A+c,\stackrel{~}{B}_2=B_2+V_1^1+m^2$$
(33)
and the associated generalized field strengths
$$\stackrel{~}{F}=\stackrel{~}{d}\stackrel{~}{A}+\frac{1}{2}[\stackrel{~}{A},\stackrel{~}{A}],\stackrel{~}{G}_3=\stackrel{~}{D}\stackrel{~}{B}_2\stackrel{~}{d}\stackrel{~}{B}_2+[\stackrel{~}{A},\stackrel{~}{B}_2].$$
(34)
The BRST-transformations in the YM-sector can be summarized by the following horizontality condition which involves ghost fields for the shifts of $`A`$:
$$\stackrel{~}{F}=F+\psi _1^1+\phi ^2.$$
(35)
From this relation and the generalized Bianchi identity $`\stackrel{~}{D}\stackrel{~}{F}=0`$, we obtain
$$\begin{array}{ccccccc}\hfill sA& =& \psi _1^1Dc\hfill & ,& \hfill s\psi _1^1& =& D\phi ^2[c,\psi _1^1]\hfill \\ \hfill sc& =& \phi ^2\frac{1}{2}[c,c]\hfill & ,& \hfill s\phi ^2& =& [c,\phi ^2]\hfill \end{array}$$
and $`sF=D\psi _1^1[c,F]`$.
In the $`B`$-sector, the $`s`$-variations follow from the horizontality condition
$$\stackrel{~}{G}_3=G_3+\psi _2^1+\phi _1^2+\varphi ^3$$
(36)
and the generalized Bianchi identity $`\stackrel{~}{D}\stackrel{~}{G}_3=[\stackrel{~}{F},\stackrel{~}{B}_2]`$: they read
$$\begin{array}{ccccccc}\hfill sB_2& =& \psi _2^1DV_1^1[c,B_2]\hfill & ,& \hfill s\psi _2^1& =& D\phi _1^2[c,\psi _2^1]+[F,m^2]+[\psi _1^1,V_1^1]+[\phi ^2,B_2]\hfill \\ \hfill sV_1^1& =& \phi _1^2Dm^2[c,V_1^1]\hfill & ,& \hfill s\phi _1^2& =& D\varphi ^3[c,\phi _1^2]+[\psi _1^1,m^2]+[\phi ^2,V_1^1]\hfill \\ \hfill sm^2& =& \varphi ^3[c,m^2]\hfill & ,& \hfill s\varphi ^3& =& [c,\varphi ^3]+[\phi ^2,m^2]\hfill \end{array}$$
and $`sG_3=D\psi _2^1[c,G_3]+[F,V_1^1]+[\psi _1^1,B_2]`$. The action (32) is inert under the BRST-transformations (4.1),(4.1) which are nilpotent by construction.
### 4.2 Gauge-fixing
The shift- and ordinary gauge symmetry in the $`B`$-sector are fixed as in the free Abelian model, except that all fields are now Lie algebra-valued. Thus, the total gauge-fixing action in the $`B`$-sector is given by
$$\mathrm{\Sigma }^B=\mathrm{\Sigma }_Q^B+\mathrm{\Sigma }_{og}^B,$$
(37)
with
$`\mathrm{\Sigma }_Q^B`$ $`=`$ $`s{\displaystyle }\mathrm{tr}\{\chi _3^1(G_3+G_3^c)+[\phi _1^2d\psi _2^1+\varphi ^3d\phi _1^2+\chi ^1d\phi _1^2CC]\}`$ (38)
$`\mathrm{\Sigma }_{og}^B`$ $`=`$ $`s{\displaystyle \mathrm{tr}\left\{V_1^1dB_2+m^2dV_1^1+ndV_1^1CC\right\}},`$
The BRST-transformations of the anti-ghost and multiplier fields are given by equations (10),(12) and (14).
In the YM-sector, the gauge-fixing can be done along the lines of four-dimensional topological YM-theory. However, the familiar four-dimensional self-duality condition for the curvature form $`F`$ does not make sense in six dimensions and it has to be generalized by introducing a $`4`$-form $`T_4`$ which is invariant under some maximal subgroup of $`SO(6)`$ : the self-duality condition can then be written as
$$F=\mathrm{\Omega }_2F\mathrm{with}\mathrm{\Omega }_2T_4.$$
(39)
This constraint is implemented by the gauge-fixing action
$$\mathrm{\Sigma }_{sd}^A=s\mathrm{tr}\left\{\chi _4^1(F\mathrm{\Omega }_2F)\right\},$$
(40)
where $`\chi _4^1`$ belongs to a BRST-doublet $`(s\chi _4^1=H_4,sH_4=0)`$. The residual gauge symmetries can then be fixed as in topological YM-theory by using a linear gauge-fixing term ,
$$\mathrm{\Sigma }^A=\mathrm{\Sigma }_{sd}^A+s\mathrm{tr}\left\{\overline{\varphi }^2d\psi _1^1+\overline{c}dA\right\},$$
(41)
which involves the BRST-doublets $`(\overline{\varphi }^2,\eta ^1)`$ and $`(\overline{c},b)`$. In summary, the total gauge-fixed action is given by $`\mathrm{\Sigma }=\mathrm{\Sigma }_{cl}+\mathrm{\Sigma }^B+\mathrm{\Sigma }^A`$.
### 4.3 VSUSY
In order to derive the VSUSY-transformations, one considers $`0`$-type symmetry conditions in the $`A`$\- and $`B`$-sectors:
$$\begin{array}{ccccccc}\hfill \delta _\tau \stackrel{~}{A}& =& 0\hfill & ,& \hfill \delta _\tau \stackrel{~}{F}& =& _\tau \stackrel{~}{A}\hfill \\ \hfill \delta _\tau \stackrel{~}{B}_2& =& 0\hfill & ,& \hfill \delta _\tau \stackrel{~}{G}_3& =& _\tau \stackrel{~}{B}_2.\hfill \end{array}$$
(42)
By expanding with respect to the ghost-number and substituting the horizontality conditions (35),(36), one finds
$$\delta _\tau (A,c)=0,\delta _\tau (B_2,V_1^1,m^2)=0.$$
(43)
and
$$\begin{array}{ccccccc}\hfill \delta _\tau \psi _1^1& =& _\tau A\hfill & ,& \hfill \delta _\tau \psi _2^1& =& _\tau B_2\hfill \\ \hfill \delta _\tau \phi ^2& =& _\tau c\hfill & ,& \hfill \delta _\tau \phi _1^2& =& _\tau V_1^1\hfill \\ & & & & \hfill \delta _\tau \varphi ^3& =& _\tau m^2.\hfill \end{array}$$
(44)
The $`\delta _\tau `$-variations for the BRST-doublets occurring in the gauge-fixing action of the $`B`$-sector are given by equations (20) and those in the YM-sector read
$$\begin{array}{ccccccc}\hfill \delta _\tau \overline{c}& =& _\tau \overline{\varphi }^2\hfill & ,& \hfill \delta _\tau b& =& _\tau \overline{c}+_\tau \eta ^1\hfill \\ \hfill \delta _\tau \chi _4^1& =& 0\hfill & ,& \hfill \delta _\tau H_4& =& _\tau \chi _4^1\hfill \\ \hfill \delta _\tau \overline{\varphi }^2& =& 0\hfill & ,& \hfill \delta _\tau \eta ^1& =& _\tau \overline{\varphi }^2.\hfill \end{array}$$
The total action is invariant under the given VSUSY-transformations which satisfy the VSUSY-algebra off-shell.
## 5 Concluding remarks
A possible generalization of the non-Abelian model consists of the addition of a $`BF`$-term $`\mathrm{tr}(B_4F)`$ (see and references therein for such models in arbitrary dimensions): such a term breaks the invariance under shifts of $`A`$. Other possible extensions are the inclusion of the topological invariant $`\mathrm{\Omega }_2\mathrm{tr}(FF)`$ (where $`\mathrm{\Omega }_2`$ is a closed $`2`$-form), which leads to “nearly topological” field theories , or the addition of a term $`\mathrm{tr}(FDZ_3)`$ involving a $`3`$-form $`Z_3`$. The resulting models can be discussed along the lines of the present paper .
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# PM/00-12 TTP00-04 hep-ph/0003151 March 2000 Light Hybrid Mesons in QCD
## 1 INTRODUCTION
Since the discovery of QCD, it has been emphasized that exotic mesons beyond the standard octet, exist as a consequence of the non-perturbative aspects of quantum chromodynamics (QCD). Since the understanding of the nature of the $`\eta ^{}`$ , a large amount of theoretical efforts have been furnished in the past and pursued at present for predicting the spectra of the exotics using different QCD-like models such as the flux tube , the bags , the quark and constituent gluon models. In this paper, we present new developments of the analysis of the hybrid mesons using QCD spectral sum rules (QSSR) à la SVZ (for a review, see e.g. ) by including the radiative perturbative corrections and the short distance tachyonic gluon mass effects which mimic the ones of UV renormalons . In this sense our results are an update of earlier results. Our predictions for the masses will be compared with the lattice results and the recent experimental candidates .
## 2 QCD SPECTRAL SUM RULES (QSSR)
### Description of the method
Since its discovery in 1979, QSSR has proved to be a powerful method for understanding the hadronic properties in terms of the fundamental QCD parameters such as the QCD coupling $`\alpha _s`$, the (running) quark masses and the quark and/or gluon QCD vacuum condensates. The description of the method has been often discussed in the literature, where a pedagogical introduction can be, for instance, found in the book . In practice (like also the lattice), one starts the analysis from the two-point correlator (standard notations):
$$\mathrm{\Pi }_{V/A}^{\mu \nu }(q^2)id^4xe^{iqx}0|𝒯𝒪_{V/A}^\mu (x)\left(𝒪_{V/A}^\nu (0)\right)^{}|0=\left(g^{\mu \nu }q^2q^\mu q^\nu \right)\mathrm{\Pi }_{V/A}^{(1)}(q^2)+q^\mu q^\nu \mathrm{\Pi }_{V/A}^{(0)}(q^2),$$
(1)
built from the hadronic local currents $`𝒪_\mu ^{V/A}(x)`$:
$$𝒪_V^\mu (x):g\overline{\psi }_i\lambda _a\gamma _\nu \psi _jG_a^{\mu \nu }:,𝒪_A^\mu (x):g\overline{\psi }_i\lambda _a\gamma _\nu \gamma _5\psi _jG_a^{\mu \nu }:$$
(2)
which select the specific quantum numbers of the hybrid mesons; A and V refer respectively to the vector and axial-vector currents. The invariant $`\mathrm{\Pi }^{(1)}`$ and $`\mathrm{\Pi }^{(0)}`$ refer to the spin one and zero mesons. One exploits, in the sum rule approaches, the analyticity property of the correlator which obeys the well-known Källen–Lehmann dispersion relation:
$$\mathrm{\Pi }_{V/A}^{(1,0)}(q^2)=_0^{\mathrm{}}\frac{dt}{tq^2iϵ}\frac{1}{\pi }\text{Im}\mathrm{\Pi }_{V/A}^{(1,0)}+\mathrm{}$$
(3)
where … represent subtraction terms which are polynomials in the $`q^2`$-variable. In this way, the $`sumrule`$ expresses in a clear way the duality between the integral involving the spectral function Im$`\mathrm{\Pi }_{V/A}^{(1,0)}(t)`$ (which can be measured experimentally), and the full correlator $`\mathrm{\Pi }_{V/A}^{(1,0)}(q^2)`$. The latter can be calculated directly in the QCD in the Euclidean space-time using perturbation theory (provided that $`q^2+m^2`$ ($`m`$ being the quark mass) is much greater than $`\mathrm{\Lambda }^2`$), and the Wilson expansion in terms of the increasing dimensions of the quark and/or gluon condensates which simulate the non-perturbative effects of QCD.
### Beyond the usual SVZ expansion
Using the Operator Product Expansion (OPE) , the two-point correlator reads for $`m=0`$:
$$\mathrm{\Pi }_{V/A}^{(1,0)}(q^2)\underset{D=0,2,\mathrm{}}{}\frac{1}{\left(q^2\right)^{D/2}}\underset{dimO=D}{}C(q^2,\nu )𝒪(\nu ),$$
where $`\nu `$ is an arbitrary scale that separates the long- and short-distance dynamics; $`C`$ are the Wilson coefficients calculable in perturbative QCD by means of Feynman diagrams techniques; $`𝒪(\nu )`$ are the quark and/or gluon condensates of dimension $`D`$. In the massless quark limit, one may expect the absence of the terms of dimension 2 due to gauge invariance. However, it has been emphasized recently that the resummation of the large order terms of the perturbative series, and the effects of the higher dimension condensates due e.g. to instantons, can be mimiced by the effect of a tachyonic gluon mass $`\lambda `$ which generates an extra $`D=2`$ term not present in the original OPE. Its presence might be understood from the analogy with the short distance linear part of the QCD potential <sup>2</sup><sup>2</sup>2Some evidence of this term is found from the lattice analysis of the static quark potential , though the extraction of the continuum result needs to be clarified.. The strength of this short distance mass has been estimated from the $`e^+e^{}`$ data to be :
$$\frac{\alpha _s}{\pi }\lambda ^2(0.060.07)\mathrm{GeV}^2,$$
(4)
which leads to the value of the square of the (short distance) string tension: $`\sigma \frac{2}{3}\alpha _s\lambda ^2[(400\pm 20)\mathrm{MeV}]^2`$ in an (unexpected) good agreement with the lattice result of about $`[(440\pm 38)\mathrm{MeV}]^2`$. In addition to Eq. (4), the strengths of the vacuum condensates having dimensions $`D6`$ are also under good control, namely:
* $`\overline{s}s/\overline{d}d0.7\pm 0.2`$ from the meson and baryon systems ;
* $`\alpha _sG^2(0.07\pm 0.01)\mathrm{GeV}^4`$ from sum rules of $`e^+e^{}I=1\mathrm{hadrons}`$ and heavy quarkonia , and from the lattice ;
* $`g\overline{\psi }\lambda _a/2\sigma ^{\mu \nu }G_{\mu \nu }^a\psi (0.8\pm 0.1)\mathrm{GeV}^2\overline{\psi }\psi ,`$ from the baryons , light mesons and the heavy-light mesons ;
* $`\alpha _s\overline{u}u^25.8\times 10^4\mathrm{GeV}^6`$ from $`e^+e^{}I=1\mathrm{hadrons}`$ ;
* $`g^3G^31.2\mathrm{GeV}^2\alpha _\mathrm{s}\mathrm{G}^2`$ from dilute gaz instantons .
### Spectral function
In the absence of the complete data, the spectral function is often parametrized using the “naïve” duality ansatz:
$`{\displaystyle \frac{1}{\pi }}\text{Im}\mathrm{\Pi }_{V/A}^{(1,0)}(t)2M_H^4f_H^2\delta (tM_H^2)+\mathrm{`}\mathrm{`}\mathrm{QCD}\mathrm{continuum}\mathrm{"}\times \theta (\mathrm{t}\mathrm{t}_\mathrm{c}),`$ (5)
which has been tested using $`e^+e^{}`$ and $`\tau `$-decay data, to give a good description of the spectral integral in the sum rule analysis; $`f_H`$ (an analogue to $`f_\pi `$) is the hadron’s coupling to the current; $`2n`$ is the dimension of the correlator; while $`t_c`$ is the QCD continuum’s threshold.
### Form of the sum rules and optimization procedure
Among the different sum rules discussed in the literature , we shall be concerned with the following Laplace sum rule (LSR) and its ratio <sup>3</sup><sup>3</sup>3FESR or $`\tau `$-like sum rules are complement to the Laplace sum rules and will be used if necessary, though the final results are independent on the form of the sum rules used.:
$$_n^{(1,0)}(\tau )=_0^{\mathrm{}}𝑑tt^n\text{exp}(t\tau )\frac{1}{\pi }\text{Im}\mathrm{\Pi }_{V/A}^{(1,0)}(t),_n\frac{d}{d\tau }\mathrm{log}_n,(n0).$$
(6)
The advantage of the Laplace sum rules with respect to the previous dispersion relation is the presence of the exponential weight factor which enhances the contribution of the lowest resonance and low-energy region accessible experimentally. For the QCD side, this procedure has eliminated the ambiguity carried by subtraction constants, arbitrary polynomial in $`q^2`$, and has improved the convergence of the OPE by the presence of the factorial dumping factor for each condensates of given dimensions. The ratio of the sum rules is a useful quantity to work with, in the determination of the resonance mass, as it is equal to the meson mass squared, in the usual duality ansatz parametrization. As one can notice, there are “a priori” two free external parameters $`(\tau ,t_c)`$ in the analysis. The optimized result will be (in principle) insensitive to their variations. In some cases, the $`t_c`$-stability is not reached due to the too naïve parametrization of the spectral function. In order to restore the $`t_c`$-stability of the results one can either fix the $`t_c`$-values by the help of FESR (local duality) or improve the parametrization of the spectral function by introducing threshold effects with the help of chiral perturbation theory The results discussed below satisfy these stability criteria.
## 3 QCD EXPRESSION OF THE TWO-POINT FUNCTION
A QCD analysis of the two-point function have been done in the past by different groups , where (unfortunately) the non-trivial QCD expressions were wrong leading to some controversial predictions . The final correct QCD expression is given in . In this paper, we extend the analysis by taking into account the non-trivial $`\alpha _s`$ correction and the effect of the new $`1/q^2`$ term not taken into account into the SVZ expansion. The corrected QCD expressions of the correlator are given in to lowest order of perturbative QCD but including the contributions of the condensates of dimensions lower or equal than six. The new terms appearing in the OPE are presented in the following <sup>4</sup><sup>4</sup>4 The results described below \[Eqs. (7) to (9)\] have been obtained with the help of program packages GEFICOM (see, e.g. ) and MINCER written in FORM . More details on the derivation of these results will be published elsewhere. Note that the results for the NLO radiative corrections are derived in neglecting some possible mixings of our operators with those containing more $`\gamma `$-matrixes like $`g\overline{\psi }_i\lambda _a\gamma _\mu \sigma _{\nu \lambda }\psi _jG_a^{\nu \lambda }`$ which could in principle mix with $`𝒪_V^\mu `$. We expect that effects due to the mixings will be small. :
* The perturbative QCD expression including the NLO radiative corrections reads:
$`{\displaystyle \frac{1}{\pi }}\text{Im}\mathrm{\Pi }_{V/A}^{(1)}(t)_{pert}`$ $`=`$ $`{\displaystyle \frac{\alpha _s}{60\pi ^3}}t^2\left\{1+{\displaystyle \frac{\alpha _s}{\pi }}\left[{\displaystyle \frac{121}{16}}{\displaystyle \frac{257}{360}}n_f+\left({\displaystyle \frac{35}{36}}{\displaystyle \frac{n_f}{6}}\right)\mathrm{log}{\displaystyle \frac{\nu ^2}{t}}\right]\right\}`$ (7)
$`{\displaystyle \frac{1}{\pi }}\text{Im}\mathrm{\Pi }_{V/A}^{(0)}(t)_{pert}`$ $`=`$ $`{\displaystyle \frac{\alpha _s}{120\pi ^3}}t^2\left\{1+{\displaystyle \frac{\alpha _s}{\pi }}\left[{\displaystyle \frac{1997}{432}}{\displaystyle \frac{167}{360}}n_f+\left({\displaystyle \frac{35}{36}}{\displaystyle \frac{n_f}{6}}\right)\mathrm{log}{\displaystyle \frac{\nu ^2}{t}}\right]\right\}`$
* The anomalous dimension of the current can be easily deduced to be:
$$\nu \frac{\mathrm{d}}{\mathrm{d}\nu }𝒪_V^\mu =\frac{16}{9}\frac{\alpha _s}{\pi }𝒪_V^\mu .$$
(8)
* The lowest order correction due to the (short distance) tachyonic gluon mass reads:
$`{\displaystyle \frac{1}{\pi }}\text{Im}\mathrm{\Pi }_{V/A}^{(1)}(t)_\lambda `$ $`=`$ $`{\displaystyle \frac{\alpha _s}{60\pi ^3}}{\displaystyle \frac{35}{4}}\lambda ^2t`$
$`{\displaystyle \frac{1}{\pi }}\text{Im}\mathrm{\Pi }_{V/A}^{(0)}(t)_\lambda `$ $`=`$ $`{\displaystyle \frac{\alpha _s}{120\pi ^3}}{\displaystyle \frac{15}{2}}\lambda ^2t`$ (9)
* The (corrected) contributions of the dimension-four and -six terms reads in the limit $`m^2=0`$ :
$`\mathrm{\Pi }_V^{(1)}(q^2)_{NP}`$ $`=`$ $`{\displaystyle \frac{1}{9\pi }}\left[\alpha _sG^2+8\alpha _sm\overline{\psi }\psi \right]\mathrm{log}{\displaystyle \frac{q^2}{\nu ^2}}`$
$`+{\displaystyle \frac{1}{q^2}}\left[{\displaystyle \frac{16\pi }{9}}\alpha _s\overline{\psi }\psi ^2+{\displaystyle \frac{1}{48\pi ^2}}g^3G^3{\displaystyle \frac{83}{432}}{\displaystyle \frac{\alpha _s}{\pi }}mg\overline{\psi }G\psi \right]`$
$`\mathrm{\Pi }_A^{(0)}(q^2)_{NP}`$ $`=`$ $`\left[{\displaystyle \frac{1}{6\pi }}\left[\alpha _sG^2+8\alpha _sm\overline{\psi }\psi \right]+{\displaystyle \frac{11}{18}}{\displaystyle \frac{\alpha _s}{\pi }}{\displaystyle \frac{1}{q^2}}mg\overline{\psi }G\psi +𝒪\left({\displaystyle \frac{1}{q^2}}\right)\right]\mathrm{log}{\displaystyle \frac{q^2}{\nu ^2}},`$ (10)
where one can notice the miraculous cancellation of the $`\mathrm{log}`$-coefficient of the dimension-six condensates for $`\mathrm{\Pi }_V^{(1)}`$.
## 4 PROPERTIES OF LIGHT HYBRIDS
### The $`\stackrel{~}{\rho }(1^+)`$
The experimental (resp. theoretical) situation has been reviewed in (resp. ). The sum rule analysis of the spectrum is based on the 2-point correlator $`\mathrm{\Pi }(q^2)_{V/A}`$ associated to the hybrid currents.
* One expects, from different QCD-like approaches, that the lightest exotic state is the one with the quantum numbers $`1^+`$ <sup>5</sup><sup>5</sup>5From QCD spectral sum rules, we also expect to a good approximation that the $`1^{}`$ is almost degenerate with the $`1^+`$.. From the analysis of the moments $`_{0,1}`$, we notice that the effect of the perturbative corrections (slightly decrease) and of the new dimension-two contribution (slightly increase) are almost negligible. This means that perturbation theory expansion in $`\alpha _s`$ converges well. The main uncertainties come from the value of $`t_c`$ because the result does not show $`t_c`$ stability. The $`\tau `$-stability of $`_0`$ also disappears if one considers the value of the subtraction constant proposed in :
$$q^2\mathrm{\Pi }_V^{(1)}|_{q^2=0}\frac{16\pi }{9}\alpha _s\overline{\psi }\psi ^2,$$
(11)
as it would cancel the effect of $`\overline{\psi }\psi ^2`$ appearing in the OPE. However, we shall check “a posteriori” that by approximating it with the sum rule estimated quantity 2$`M_{\stackrel{~}{\rho }}^4f_{\stackrel{~}{\rho }}^2`$ from the following Eq. (13), this result is inaccurate and, therefore, we shall not include this term in our analysis. An independent measurement of this quantity, e.g. on the lattice is required.
* Using the different QCD input parameters given previously and the value $`\mathrm{\Lambda }=(0.35\pm 0.05)`$ GeV, the positivity ($`t_c\mathrm{}`$) of the $`_0`$ moment leads to the rigorous upper bound:
$$M_{\stackrel{~}{\rho }}1.9\text{GeV},$$
(12)
which excludes some of range spanned by the quenched lattice estimates of $`(1.9\pm 0.2)`$ GeV .
* For reasonnable finite values of $`t_c3.5`$ GeV<sup>2</sup> (beginning of $`\tau `$-stability) to 4.5 GeV<sup>2</sup> as also fixed by the Finite Energy Sum Rule constraints , we obtain at the stability point $`\tau (0.50.6)`$ GeV<sup>-2</sup> of $`_0`$, the common solution of $`_0`$ and $`_1`$:
$$M_{\stackrel{~}{\rho }}(1.61.7)\text{GeV},f_{\stackrel{~}{\rho }}(2550)\text{MeV},(M_{\stackrel{~}{\rho }^{}}\sqrt{t_c})M_{\stackrel{~}{\rho }}200\mathrm{MeV},$$
(13)
where the $`\stackrel{~}{\rho }^{}`$ is the radial excitation. One can consider this result as an improvement of the available sum rule results ranging from 1.4 to 2.1 GeV, . Though, we cannot absolutely exclude the presence of the $`1.41.6`$ GeV experimental candidates , we expect from your analysis that this observed state is a hybrid which can have a small $`\overline{q}q\overline{q}q`$ component through mixing.
* $`\stackrel{~}{\rho }^{}`$-$`\stackrel{~}{\rho }`$ mass-splitting is much smaller than $`M_\rho ^{}M_\rho 700`$ MeV, and can signal a rich population of $`1^+`$ states above 1.6 GeV.
* The hadronic widths have been computed in . Given our new values of the mass and decay constant, the updated values are:
$`\mathrm{\Gamma }(\stackrel{~}{\rho }\rho \pi )`$ $``$ $`274\text{MeV},\mathrm{\Gamma }(\stackrel{~}{\rho }\gamma \pi )3\text{MeV},`$
$`\mathrm{\Gamma }(\stackrel{~}{\rho }\eta ^{}\pi )`$ $``$ $`3\text{MeV},\mathrm{\Gamma }(\stackrel{~}{\rho }\pi \pi ,\overline{K}K,\eta _8\eta _8)𝒪(m_q^2).`$ (14)
* One can measure the $`SU(3)`$ breakings and the mass of the $`\stackrel{~}{\varphi }(\overline{s}s)`$ from the difference of the ratio of moments, which gives :
$$M_{\stackrel{~}{\varphi }}^2M_{\stackrel{~}{\rho }}^2\frac{20}{3}\overline{m}_s^2\frac{160\pi ^2}{9}m_s\overline{s}s\tau 0.3\text{GeV}^2M_{\stackrel{~}{\varphi }}(1.71.8)\text{GeV}.$$
(15)
The quenched lattice results are in the range of ($`2.0\pm 0.2`$) GeV , which is slightly higher than our result.
### The $`\stackrel{~}{\eta }(0^{})`$
Similar analysis can be done for the pseudoscalar channel. In this case, the most convenient sum rule to work with is $`_1`$, which presents both $`\tau `$ and $`t_c`$ stabilities. Using the previous QCD input parameters, stabilities are reached for $`\tau 0.12`$ GeV<sup>-2</sup> and $`t_c7.8`$ GeV<sup>2</sup> showing again that the mass-splitting between the radial excitation and the ground state is tiny. At these values one obtains <sup>6</sup><sup>6</sup>6 In the case $`\lambda =0`$, i.e. of the ordinary SVZ expansion, the stability is obtained for $`\tau 0.15`$ GeV<sup>-2</sup> at which, one can deduce $`M_{\stackrel{~}{\eta }}3.2`$ GeV$`\sqrt{t_c}`$.:
$$M_{\stackrel{~}{\eta }}2.8\mathrm{GeV}\sqrt{t_c},$$
(16)
which we consider as an update of the previous results in . Again the effects of the correction terms are small. The relatively higher value of the mass of the $`0^{}`$ meson than the one of the $`1^+`$, is mainly due to the relative strength of the perturbative and non-perturbative terms.
## 5 CONCLUSIONS
There are some progresses in the long run study and experimental search for the exotics. Before some definite conclusions, one still needs improvments of the present data, and some improved lattice unquenched estimates which should complement the QCD spectral sum rule (QSSR) results. In this paper we have updated previous sum rule analysis of the light hybrids by including the perturbative radiative corrections and the new effect due to the (short distance) tachyonic gluon mass not included in the original SVZ expansion. However, these effects are negligible which are reassuring for the validity of the approximation used. Our result which is $`M_{\stackrel{~}{\rho }}(1.61.7)\text{GeV}`$ can be renconciled with the existence of the $`1^+`$ states at $`(1.41.6)`$ GeV seen recently in hadronic machines (BNL and Crystal Barrel) , but in the same time predicts the existence of a $`1^{}`$ hybrid almost degenerate with the $`1^+`$, and which could manifest in $`e^+e^{}`$ hadrons by mixing with the radial excitations of the $`\rho `$ and $`\omega `$ mesons. In addition, the relatively low value of the continuum threshold indicates that we expect a rich population of (axial-) vector hybrids in the region above 1.8 GeV.
In our analysis, the $`0^{}`$ mass is about 2.8 GeV, which is in the range of the different charmonium states, such that it could mix with these charmonium states as well. Moreover, the small splitting between the continuum threshold and the lowest ground state indicates that rich population of pseudoscalar hybrids is expected in the 3 GeV mass range.
Light hybrid mesons ($`1^{}`$ and $`0^+`$) might be (partly) responsible of the anomalous behaviour of the $`e^+e^{}`$ hadrons cross section observed in the region below 4 GeV.
## Acknowledgements
We are grateful to Alexei Pivovarov for his collaboration on early stage of the present work as well as for numerous discussions and a good advice. K. Ch. would like to thank Matthias Steinhauser for his help with prgramming Feynman rules for hybrid currents; his work was supported by DFG under Contract Ku 502/8-1. S.N. thanks the KEK Theory Group and K. Hagiwara for their hospitality.
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# Static and dynamic spin correlations in the spin-glass phase of slightly-doped La2-xSrxCuO4
## I Introduction
The phase diagram of La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub> shows that the magnetic state changes dramatically with Sr doping. The parent material La<sub>2</sub>CuO<sub>4</sub> exhibits three-dimensional (3D) long-range antiferromagnetic (AF) order below $``$325 K. When a small fraction of La is replaced by Sr, which corresponds to hole-doping, the 3D AF order disappears and the low temperature magnetic phase is replaced by a disordered magnetic phase in which commensurate two-dimensional (2D) short-range AF fluctuations are observed. In crucible-grown samples the commensurate fluctuations develop a static component at low temperatures, signalling the onset of spin-glass order.
In superconducting samples, an essential feature is that the magnetic correlations become incommensurate (IC). Detailed studies on the hole concentration dependence of the low energy magnetic excitations have been performed by Yamada $`et`$ $`al.`$ They find that the incommensurability ($`\delta `$) is almost linear with hole concentration ($`x`$) with $`\delta x`$ below $`x`$0.12. Recently, static magnetic ordering has been observed in superconducting La<sub>1.88</sub>Sr<sub>0.12</sub>CuO<sub>4</sub> with the magnetic onset temperature near $`T_c`$. The elastic magnetic peaks are observed at the same IC positions as those of the magnetic inelastic peaks. A model that describes this behavior is that of stripe ordering of spin and charge (hole) density waves as observed in La<sub>2-y-x</sub>Nd<sub>y</sub>Sr<sub>x</sub>CuO<sub>4</sub>. In this case the charge and, concomitantly, spin stripes run approximately along the $`a_{\mathrm{tetra}}`$ or $`b_{\mathrm{tetra}}`$ axis; we label this the collinear stripe phase.
Thus, the magnetism and the transport properties in the doped La<sub>2</sub>CuO<sub>4</sub> system are intimately related. In the insulating phase at low hole concentrations, spin-glass behavior is observed and there are strong quasi-elastic commensurate spin fluctuations; dynamic IC spin fluctuations persist in the superconducting phase. It has been known for some time that in La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub> the instantaneous magnetic correlations change from being commensurate to IC at the insulator-to-superconductor boundary. Recently, Wakimoto $`et`$ $`al.`$ have found that in a sample grown with the crucible-free traveling solvent floating zone technique which results in purer crystals the static magnetic correlations at low temperature are also IC in the insulating spin-glass La<sub>1.95</sub>Sr<sub>0.05</sub>CuO<sub>4</sub>. They have examined the intensity profiles and have shown that there are only 2 satellite peaks along $`b_{\mathrm{ortho}}`$ while in superconducting compounds the IC peaks are located parallel to both the $`a_{\mathrm{tetra}}`$ and $`b_{\mathrm{tetra}}`$ axes. These magnetic correlations in La<sub>1.95</sub>Sr<sub>0.05</sub>CuO<sub>4</sub> are consistent with diagonal charge stripes, in which the stripes run along the $`a_{\mathrm{ortho}}`$ axis. Actually, such diagonal stripes have been predicted theoretically. Diagonal stripes are also reported experimentally in insulating La<sub>2-x</sub>Sr<sub>x</sub>NiO<sub>4</sub>. We emphasize, however, that only a one-dimensional spin modulation has been observed in La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub> to-date; any associated charge ordering has not yet been detected. These results lead to the important conclusion that the static magnetic spin modulation changes from diagonal to collinear at $`x=0.055\pm 0.005`$, coincident with the insulator-to-superconductor transition.
A fundamental question is whether or not the diagonal one-dimensional IC magnetic correlations persist throughout the spin-glass phase down to the critical concentration of $`x`$=0.02 for 3D Néel ordering. The present neutron scattering study clarifies this point and yields important new information on the concentration dependence of the incommensurability. Especially, we find that at the lowest concentration within the context of the stripe model the inferred charge density is $``$1 hole/Cu as in La<sub>2-x</sub>Sr<sub>x</sub>NiO<sub>4</sub>.
Another important point is to clarify the nature of the magnetic excitations in the diagonal IC state. Intensive studies of the inelastic magnetic spectra in insulating La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub> were performed by Keimer $`et`$ $`al.`$ They studied the energy and temperature dependences of the $`Q`$-integrated susceptibility. However, the $`Q`$-dependence of the excitation spectra was not discussed since the detailed peak profile was not known. Matsuda $`et`$ $`al.`$ also studied the inelastic magnetic spectra in La<sub>1.98</sub>Sr<sub>0.02</sub>CuO<sub>4</sub>, in which the energy and temperature dependences of the excitation spectra were measured. They only discussed the results qualitatively since the detailed peak profile could not be clarified. Now that the static magnetic correlations have been elucidated, the excitation spectra can be analyzed qualitatively. Specifically, we found that the magnetic correlations change from being incommensurate to commensurate at $`\omega `$7 meV and $`T`$70 K, indicating a characteristic energy for the IC structure of 6-7 meV.
## II Experimental details
The single crystal of La<sub>1.976</sub>Sr<sub>0.024</sub>CuO<sub>4</sub> was grown by the traveling solvent floating zone (TSFZ) method. The crystal was annealed in an Ar atmosphere at 900 C for 24 h. The dimensions of the rod shaped crystal were $``$5$`\mathrm{\Phi }\times `$25 mm<sup>3</sup>. The lattice constants were $`a_{\mathrm{ortho}}`$=5.349 Å, $`b_{\mathrm{ortho}}`$=5.430 Å ($`b/a`$1.015), and $`c`$=13.151 Å at 10 K. From the universal relation for the spin-glass transition temperature, the tetragonal-to-orthorhombic structural transition temperature, and the orthorombicity $`b/a`$, the effective hole concentration was estimated to be 0.024$`\pm `$0.003.
The neutron scattering experiments were carried out on the cold neutron three-axis spectrometer HER and the thermal neutron three-axis spectrometer TOPAN installed at JRR-3M at the Japan Atomic Energy Research Institute (JAERI). The horizontal collimator sequences were guide-open-S-80-80 with the fixed incident neutron energy $`E_\mathrm{i}`$=5 meV at HER and 30-30-S-60-60 with the fixed final neutron energy $`E_\mathrm{f}`$=14.7 meV at TOPAN. Contamination from higher-order beams was effectively eliminated using Be filters at HER and PG filters at TOPAN. The single crystal, which was oriented in the $`(HK0)_{\mathrm{ortho}}`$ or $`(H0L)_{\mathrm{ortho}}`$ scattering plane, was mounted in a closed cycle refrigerator. In this paper, we use the low temperature orthorhombic phase ($`Bmab`$) notation $`(h,k,l)_{\mathrm{ortho}}`$ to express Miller indices.
The crystal has a twin structure and there exist two domains. The two domains are estimated to be equally distributed from the ratio of the nuclear Bragg peak intensities from both domains. Figure 1A shows the scattering geometry in the $`(HK0)`$ scattering plane. The filled triangles correspond to the (1,0,0) and (0,1,0) Bragg points from domain A while the open triangles denote the (1,0,0) and (0,1,0) Bragg points from domain B.
## III Results and Discussion
### A Static properties
Below $``$40 K elastic magnetic peaks develop and at low temperatures the peaks are clearly resolved at the IC positions (1, $`\pm ϵ`$, 0) and (0, 1$`\pm ϵ`$, 0) with $`ϵ`$0.023. This corresponds to the same diagonal one-dimensional spin modulation observed in La<sub>1.95</sub>Sr<sub>0.05</sub>CuO<sub>4</sub> which has $`ϵ`$0.064. The open and filled circles in Fig. 1A correspond to the IC magnetic peaks from the two domains in the $`(HK0)`$ zone, respectively.
Figures 1B-1D show transverse and longitudinal elastic scans around (1,0,0) and (0,1,0). Two peaks are observed in the transverse scan A while one intense peak together with a weak shoulder on the low-$`h`$ side is observed in the longitudinal scans B and C. The instrumental resolution at (1,0,0) can be estimated from higher-order reflections, which in turn are measured by removing the Be filters. As illustrated in Figs. 1B and 1C, the magnetic peaks are much broader than the resolution along both $`h`$ and $`k`$. It should be noted that, by contrast, the magnetic peaks in the superconducting state of La<sub>1.88</sub>Sr<sub>0.12</sub>CuO<sub>4</sub> are all resolution limited, indicating that the magnetic correlation length in that system is quite large in the CuO<sub>2</sub> plane.
Figure 2 shows the $`L`$-dependence of the magnetic elastic peaks at (1,0,$`L`$) and (0,0.975,$`L`$), respectively, at 10 K. The background estimated from the high temperature data (100 K) has been subtracted so that the remaining signal is purely magnetic. These scans probe the magnetic correlations along the $`c`$ axis between neighboring 2D antiferromagnetically correlated planes. Broad peaks are observed at $`(1,0,even)`$ and (0,$``$1,$`odd`$) which coincide with the magnetic Bragg peak positions in pure La<sub>2</sub>CuO<sub>4</sub>. However, the magnetic intensity at $`(1,0,even)`$ initially increases with increasing $`L`$, implying a cluster spin-glass model as in La<sub>1.98</sub>Sr<sub>0.02</sub>CuO<sub>4</sub>: the spin system forms antiferromagnetically correlated clusters which have randomly different spin directions in the CuO<sub>2</sub> plane although the propagation vector of the AF order is along $`a_{\mathrm{ortho}}`$ in each cluster.
The solid lines in Figs. 1B-1D are the results of fits to a convolution of the resolution function with 3D squared Lorentzians. The two intense peaks in Fig. 1B originate primarily from the magnetic signals at (1,$`\pm ϵ`$,0) in domain A while the weak shoulder in Fig. 1C originates from magnetic signals at (0,1$`ϵ`$,$`\pm `$1) in domain B. The relatively intense peaks at (0,1$`\pm ϵ`$,0) occur because of the short correlation length along the $`c`$ axis, which in turn makes the (0,1$`\pm ϵ`$,$`L`$), with $`L`$ odd, magnetic peaks broad along the $`c`$ axis as shown in Fig. 2. The instrumental resolution function is also elongated along the $`c`$ axis so that the magnetic signals are effectively integrated. The observed data are fitted with $`\xi _a^{}`$=94.9$`\pm `$4.0 Å, $`\xi _b^{}`$=39.9$`\pm `$1.3 Å, $`\xi _c^{}`$=3.15$`\pm `$0.08 Å, and $`ϵ`$=0.0232$`\pm `$0.0004 r.l.u., where $`\xi _a^{}`$, $`\xi _b^{}`$, and $`\xi _c^{}`$ represent the inverse of the half width at half maxima of the elastic peak widths in $`Q`$ along the $`a`$, $`b`$, and $`c`$ axes, respectively. The calculation reproduces the observed profiles quite well. The error bars represent one standard deviation statistical error limits for the assumed lineshape. The true error limits, indicating possible systematic errors are much larger. The static correlation length perpendicular to the CuO<sub>2</sub> plane is 3.15 Å, which is much less than the distance between nearest-neighbor CuO<sub>2</sub> planes ($`c/2`$6.5 Å), indicating that the static magnetic correlations are almost two-dimensional. The between-plane correlation length is shorter than those in La<sub>1.48</sub>Nd<sub>0.4</sub>Sr<sub>0.12</sub>CuO<sub>4</sub> ($`\xi _c0.55c`$ at 1.38 K) and La<sub>1.775</sub>Sr<sub>0.225</sub>NiO<sub>4</sub> ($`\xi _c1.06c`$ at 10 K).
The solid lines in Fig. 2 show the calculated profiles using 3D squared Lorentzian profiles convoluted with the instrumental resolution function. The parameters determined above are held fixed and only the overall scale factor has been adjusted. In order to reproduce the $`L`$-dependence of the $`(1,0,even)`$ intensity, the cluster spin-glass model, as described above, has been used in the calculation. The calculation describes the observed profiles reasonably well. The slight deviation at large $`L`$ in both the $`(1,0,L)`$ and $`(0,0.975,L)`$ scans probably reflects a decrease at large $`L`$ of the magnetic form factor, which has been assumed to be constant in the calculation.
The diagonal magnetic stripe model thus provides a good description of the data; this is one of the most significant results of this study. The incommensurability $`ϵ`$ corresponds to the inverse modulation period of the spin density wave. Here, $`ϵ`$ is defined in orthorhombic notation so that $`ϵ=\sqrt{2}\times \delta `$ where $`\delta `$ is defined in tetragonal units. As shown in Fig. 3, $`\delta `$ follows the linear relation $`\delta =x`$ reasonably well over the range 0.03$`x`$0.12 which spans the insulator-superconductor transition. In a charge stripe model this corresponds to a constant charge per unit length in both the diagonal and collinear stripe phases, or equivalently, 0.7 and 0.5 holes per Cu respectively because of the $`\sqrt{2}`$ difference in Cu spacings in the diagonal and collinear geometries. Our value for $`x`$=0.024 definitely deviates from the $`\delta =x`$ line and instead appears to be close to $``$1 hole/Cu as in La<sub>2-x</sub>Sr<sub>x</sub>NiO<sub>4</sub> where there is $``$1 hole/Ni along the diagonal stripes. This suggests that as the hole concentration is decreased, in the context of the stripe model, the hole concentration evolves progressively from $``$0.5 hole/Cu at $`x`$=0.12 to 1 hole/Cu at $`x`$=0.024. This behavior is very different from that in La<sub>2-x</sub>Sr<sub>x</sub>NiO<sub>4</sub>, where the hole density is $``$1 hole/Ni over a wide range of hole concentrations in the insulating phase albeit at rather larger hole densities. We note that Machida and Ichioka predict 1 hole/Cu throughout the diagonal stripe phase.
The static spin correlation lengths, which are derived from the inverse peak widths in $`Q`$, in La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub> ($`x`$=0.02, 0.024, and 0.05) are summarized in Table 1. With increasing hole concentration, the peak widths both parallel and perpendicular to the CuO<sub>2</sub> plane rapidly broaden. In La<sub>1.95</sub>Sr<sub>0.05</sub>CuO<sub>4</sub>, an $`L`$-scan shows that the peak width along $`L`$ is much broader than that in La<sub>1.976</sub>Sr<sub>0.024</sub>CuO<sub>4</sub>, indicating that the magnetic correlations are more two-dimensional.
There are at least two possible origins of the finite correlation lengths in the CuO<sub>2</sub> plane for the static order in the spin-glass La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub>. The first is that the lengths simply measure the spin decoherence distance of the AF spin clusters. The second is that the disorder originates primarily from a random distribution of stripe spacings and orientations as discussed by Tranquada $`et`$ $`al`$ It is noted that a rotation of the stripe orientation away from the $`a_{\mathrm{tetra}}`$ and $`b_{\mathrm{tetra}}`$ axes is observed in La<sub>2</sub>CuO<sub>4+y</sub> (Ref. ) and La<sub>1.88</sub>Sr<sub>0.12</sub>CuO<sub>4</sub> (Ref. ). Further experiments and theoretical calculations will be required to choose between these possibilities.
### B Magnetic excitations
As mentioned in the previous section, static magnetic correlations in the spin-glass phase has been revealed. We now consider how the inelastic magnetic correlations behave in the spin-glass phase of La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub>. Neutron inelastic scattering measurements were performed in La<sub>1.976</sub>Sr<sub>0.024</sub>CuO<sub>4</sub>. The measurements were performed in the ($`H0L`$) scattering plane. We first carried out the measurements in the ($`HK0`$) scattering plane, in which the elastic IC magnetic peaks can be resolved reasonably well since the instrumental resolution is narrower than the intrinsic peak widths as shown in Fig. 4(a). However, with increasing transfer energy, the resolution becomes worse and the magnetic intensity decreases due to the structure factor. Because of these two effects, it is very difficult to follow the energy dependence of the excitation spectra. In the ($`H0L`$) scattering plane, since the resolution is elongated perpendicular to $`H`$, as shown in Fig. 4(b), the IC peak cannot be well resolved. However, as we will show later, the intrinsic peak configuration can be estimated by fitting to model functions. The advantage of the measurements in the ($`H0L`$) scattering plane is that the measurements can be performed with increased scattering intensity.
First, we performed constant-$`Q`$ scans around ($`\pi ,\pi `$) in order to study the magnetic anisotropy in La<sub>1.976</sub>Sr<sub>0.024</sub>CuO<sub>4</sub>. Figure 5(a) shows the result of constant-$`Q`$ scans at (0.992,0,0.7), which corresponds to just the midpoint between (1,0,0.7) and (0,1,0.7). The background intensities are measured at (1.15,0,0.7). The magnetic intensity decreases gradually with increasing energy, indicating that the magnetic excitation spectrum is gapless in La<sub>1.976</sub>Sr<sub>0.024</sub>CuO<sub>4</sub>. For comparison, constant-$`Q`$ scans in pure La<sub>2</sub>CuO<sub>4</sub> with the same spectrometer condition are shown in Fig. 5(b). An excitation gap due to the out-of-plane anisotropy is found at $``$5 meV, which is consistent with that observed previously.
Figure 6 shows constant-$`\omega `$ scans in the ($`H0L`$) scattering plane at various energies and temperatures. A sharp excitation peak is centered at $`H`$=1 at $`\omega `$=3 meV and $`T`$=10 K. On the other hand, the peak position shifts progressively to lower $`H`$ at higher energies and temperatures.
We speculate that this behavior may be explained as follows. At 3 meV, the magnetic peaks exist at IC positions as observed in the elastic scans in the ($`HK0`$) scattering zone as shown in Fig. 4(a). Since the instrumental resolution elongated vertically integrates the magnetic signal around $`H`$=1 very effectively in the $`(H0L)`$ scattering plane, as shown in Fig. 4(b), a sharp and intense peak centered at $`H`$=1 is observed while a weak tail is found at lower $`H`$. The solid line in the Fig. 6(a) is the result of a calculation assuming that the magnetic peaks are located at exactly the same positions as those determined from the elastic measurements in the ($`HK0`$) scattering plane. As excitation energy is increased, the peak separation appears to to become smaller. The excitation spectrum at 6 meV is fitted to 3D Lorentzians. The fitting parameters are the peak separation $`ϵ`$, the isotropic inverse inelastic peak width in the CuO<sub>2</sub> plane $`\xi _{ab}^{\prime \prime }`$, and the amplitude. The inverse inelastic peak width along the c-axis $`\xi _c^{\prime \prime }`$ is fixed at 3.15 Å, which is the same as $`\xi _c^{}`$. The solid line in the Fig. 6(b) is the result of a fit to the 3D Lorentzians with $`ϵ`$= 0.0096$`\pm `$0.0068 r.l.u. and $`\xi _{ab}^{\prime \prime }`$=48$`\pm `$12 Å. Finally, the magnetic correlations appears to become commensurate and isotropic at 9 meV as shown in Fig. 4(c). In this case, there exist two equi-intense peaks at (0,1,-0.6) ($`H`$=0.985) and (1,0,-0.6) ($`H`$=1) resulting in one broad peak located at $`H`$0.99. The solid line in Fig. 6(c) is the result of a fit to the 3D Lorentzians with $`\xi _c^{\prime \prime }`$=3.15 Å (fixed) and $`\xi _{ab}^{\prime \prime }`$=42$`\pm `$13 Å. $`ϵ`$ is fixed at 0 r.l.u. since it becomes very close to 0 r.l.u. even when fitted. From these results, we conclude that magnetic correlations change from being incommensurate to commensurate between $``$6 and $``$7.5 meV.
The commensurate magnetic correlations at higher energies are similar to those observed in pure La<sub>2</sub>CuO<sub>4</sub> above the gap energy $``$5 meV as shown above. The excitation spectrum in pure La<sub>2</sub>CuO<sub>4</sub> at 6 meV is shown in Fig. 6(f). The spectrum is consistent with the peak configuration as shown in Fig. 4(c) since the magnetic correlations are commensurate and isotropic. The peak width is resolution-limited, indicating that the correlation length is very long in the CuO<sub>2</sub> plane. This is in striking contrast to the situation in La<sub>1.976</sub>Sr<sub>0.024</sub>CuO<sub>4</sub>. The solid line in Fig. 6(f) is the result of a calculation assuming the 3D Lorentzians with $`\xi _{ab}^{\prime \prime }`$=700 Å (fixed), $`\xi _c^{\prime \prime }`$=1 Å (fixed), and $`ϵ`$= 0 r.l.u. (fixed).
The temperature dependence of the excitation spectra in La<sub>1.976</sub>Sr<sub>0.024</sub>CuO<sub>4</sub> is quite similar to the energy dependence. A sharp excitation peak centered at $`H`$=1 at low temperatures shifts to lower $`H`$ with increasing temperature. The solid line in the Fig. 6(d) is the result of a fit to an assumed 3D Lorentzian line-shape with $`ϵ`$= 0.0113$`\pm `$0.0068 r.l.u. and $`\xi _{ab}^{\prime \prime }`$=60$`\pm `$11 Å. $`\xi _c^{\prime \prime }`$ is fixed at 1 Å since the system becomes magnetically 2D above $``$40 K. $`ϵ`$ and $`\xi _{ab}^{\prime \prime }`$ do not change sensitively with changes in $`\xi _c^{\prime \prime }`$ in the fitting. The solid line in Fig. 6(e) is the result of a fit to a 3D Lorentzian with $`\xi _{ab}^{\prime \prime }`$=62$`\pm `$10 Å and $`\xi _c^{\prime \prime }`$=1 Å (fixed). $`ϵ`$ is fixed at 0 r.l.u. since it becomes very close to 0 r.l.u. even when fitted. From these results, we conclude that the magnetic correlations change from being incommensurate to commensurate between 55 and 70 K.
Figure 7 represents a summary of the neutron inelastic measurements in La<sub>1.976</sub>Sr<sub>0.024</sub>CuO<sub>4</sub>. The open and filled circles signify that the magnetic correlations are diagonal IC and commensurate, respectively. The diagonal IC phase exists below $`\omega `$7 meV and $`T`$70 K ($``$6 meV). This result indicates that the characteristic energy for the diagonal IC structure is 6-7 meV.
## IV Summary
In brief, we find that a short range static one-dimensional diagonal spin modulation exists at low temperatures across the entire insulating spin-glass region in La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub>. Further, within the context of a spin and charge stripe model the charge density per unit length is almost constant for all values of $`x`$, but shows a significant deviation near the spin-glass 3D Néel boundary suggesting stability of diagonal stripes with 1 hole/Cu at low $`x`$.
The magnetic excitation spectra suggest that magnetic correlations change from diagonal incommensurate to commensurate at $`\omega `$7 meV and $`T`$70 K. Above these energy and temperature the magnetic correlations are similar to those in pure La<sub>2</sub>CuO<sub>4</sub> although the range of order in the CuO<sub>2</sub> plane is much shorter in La<sub>1.976</sub>Sr<sub>0.024</sub>CuO<sub>4</sub>.
## Acknowledgments
We would like to thank A. Aharony, K. Katsumata, and K. Machida for stimulating discussions. This study was supported in part by the U.S.-Japan Cooperative Program on Neutron Scattering, by a Grant-in-Aid for Scientific Research from the Japanese Ministry of Education, Science, Sports and Culture, by a Grant for the Promotion of Science from the Science and Technology Agency, and by CREST. Work at Brookhaven National Laboratory was carried out under Contract No. DE-AC02-98CH10886, Division of Material Science, U.S. Department of Energy. The research at MIT was supported by the National Science Foundation under Grant No. DMR97-04532 and by the MRSEC Program of the National Science Foundation under Award No. DMR98-08941.
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# References
Classical phase space and statistical
mechanics of identical particles
USITP-00-02
OSLO-TP 2-00
March-2000
T.H. Hansson<sup>1</sup>, S.B. Isakov<sup>2</sup>, J.M. Leinaas<sup>3</sup>, and U. Lindström<sup>4</sup>
<sup>1,4</sup> Department of Physics, University of Stockholm, Box 6730,
S-11385 Stockholm, Sweden
<sup>2,3</sup>Department of Physics, University of Oslo, P.O. Box 1048 Blindern, N-0316 Oslo, Norway
ABSTRACT
Starting from the quantum theory of identical particles, we show how to define a classical mechanics that retains information about the quantum statistics. We consider two examples of relevance for the quantum Hall effect: identical particles in the lowest Landau level, and vortices in the Chern-Simons Ginzburg-Landau model. In both cases the resulting classical statistical mechanics is shown to be a nontrivial classical limit of Haldane’s exclusion statistics. e-mail: <sup>1</sup> hansson@physto.se, <sup>2</sup> serguei.isakov@fys.uio.no, <sup>3</sup> j.m.leinaas@fys.uio.no, <sup>4</sup> ul@physto.se
<sup>2</sup> Present address: NumeriX Corporation, 546 Fifth Avenue, New York, NY 10036, USA
<sup>1,4</sup> Supported by the Swedish Natural Science Research Council
<sup>3</sup> Supported by the Norwegian Research Council
1. Introduction
Particle statistics is usually considered to be a quantum effect. It is expressed through the symmetry of the wave function of a system of identical particles and does not appear, normally, as an interaction in the Hamiltonian. In two space dimensions it can be represented as a special kind of interaction, but being of the Aharonov-Bohm type, it does not give rise to any force on the particles.
Thus, at the level of classical trajectories of individual particles, there is no difference between identical and non-identical particles. There is however one place in the classical description of particles where their indistinguishability is important, namely in the statistical mechanics. There the trajectories of individual particles no longer matters, but the volume of the available phase space is important for thermodynamical quantities. Indistinguishability is introduced by dividing the phase space volume of $`N`$ non-identical particles with the factor $`N!`$. This reduction is essential to give the correct expression for the entropy and thus to resolve Gibbs’ paradox.
The reduction in phase space is readily understood. If the particles are indistinguishable all configurations that can be related by a permutation of the particles correspond to one and the same physical configuration. This single configuration for identical particles is then represented as $`N!`$ different configurations in the case of distinguishable particles. The configuration space of indistinguishable point particles is therefore derived from the space of distinguishable particles by an identification of equivalent points.
The identification of points implies that the configuration space of a system of identical particles is not everywhere a smooth manifold, there are singularities corresponding to the situation where two or more particles occupy the same point in space. Such a point is a geometrical singularity, a point of infinite curvature. For the phase space the situation is similar, the identification of points introduce singularities, although in general these singularities do not have the same simple geometrical interpretation as in configuration space.
The quantum description of identical particles can be introduced in terms of wave functions, or alternatively in terms of path integrals, defined on the configuration space with identifications . The presence of singularities then are important, since it divides the continuous paths into different classes, depending on how they evolve around the singularities. Such classes can be associated with different phase factors. There is only one characteristic phase factor for each system of identical particles, corresponding to an exchange of two particles, and this factor identifies the statistics. Viewed in this way, the statistics parameter associated with the particles labels inequivalent quantizations of the classical system. Thus, the statistics parameter appears in the quantization of the system and is not present in the classical description of the particles.
In this paper we will to discuss an alternative approach to the classical description of identical particles. This does not mean that we consider the standard description of point particles referred to above as being in any sense incorrect. However, we would like to stress that starting from the quantum theory there are different possibilities for describing the corresponding classical system, and we would like to examine one where the statistics parameter is present also at the classical level. As discussed in the paper we may view this as a non-standard way of taking the classical limit.
The way we introduce the classical description is to consider, in a general form, a coherent state representation of the quantum system. We assume the coherent states to be determined by a set of particle coordinates, and we further assume the time evolution (in the low energy regime) to a good approximation to be described simply by the motion of these coordinates. There is a manifold defined by the set of possible coordinates and a natural phase space structure inherited from the full quantum description. This phase space is a smooth manifold, even when the particle coordinates coincide, and the reduction corresponding to the factor $`1/N!`$ does not have to be introduced by hand, but appears naturally when calculating the phase space volume.
To clarify this idea, we consider the case of a harmonic oscillator coherent state representation of a system of identical point particles in some detail. We show how the classical description introduced in this way distinguishes between bosons, fermions, and in general anyons , and we calculate the available phase space volume for the case of $`N`$ identical particles in a finite volume. The classical statistics parameter is then identified as the phase space volume occupied by each of the particle present in the system. Viewed in this way the description has the character of a classical analogue of the quantum exclusion statistics introduced some time ago by Haldane . We examine this correspondence in some detail by considering the statistical mechanics of our classical system.
The description we use is not restricted to systems of point particles. We illustrate this by considering vortex solutions of the Chern-Simons Ginzburg-Landau theory. The manifold defined by the $`N`$-vortex configurations has a natural phase space structure, and although this can not be fully determined, the phase space volume can be calculated and the statistics parameter identified. This particle description of vortices is closely related to a description of vortices in the (relativistic) abelian Higgs model previously discussed by Samols, Manton and others although in their case the vortex manifold is identified as a configuration space rather than as a phase space.
2. Classical phase space from the quantum description
In this section we consider a general quantum system and a subset of states $`|\psi _x`$, which is indexed by a set of coordinates $`x=\{x_1,x_2,\mathrm{},x_N\}`$. These may be the coordinates of a system of (identical) particles or the coordinates of an $`N`$ soliton configuration, but we do not have to be more specific at this point. We only assume that the wave function evolves smoothly with a change of these coordinates, and that it is symmetric under an interchange of any pair of the $`N`$ coordinates. We furthermore assume, that in the regime of interest (typically at low energies), the time evolution of the system, to a good approximation, can be described (up to a phase factor) as a time evolution of the coordinates only. This means that it makes sense to consider the restricted (constrained) system where the evolution of the system is projected to the manifold defined by the normalized states $`|\psi _x`$. Since the physical states correspond to rays in the Hilbert space, i.e. to state vectors defined up to a complex factor, we consider the classical $`N`$-particle space $``$, derived from the quantum description, to be defined by the normalized states $`|\psi _x`$ only up to such a phase factor. It is the phase space structure of the space $``$ which will be of importance for our discussion.
The Schrödinger equation of the quantum system can be derived from the Lagrangian,
$`L=i\mathrm{}\psi |\dot{\psi }\psi |H|\psi ,`$ (2.1)
where $`H`$ is the Hamiltonian of the system, and the Lagrangian of the constrained system is obtained from this by restricting $`|\psi `$ to the subset of states $`|\psi _x`$. Expressed in terms of the coordinates $`x`$, it has the generic form
$`L=\dot{x}_iA_i(x)V(x),`$ (2.2)
where $`A_i`$ is the Berry connection
$`A_i=i\mathrm{}\psi _x|_i\psi _x,`$ (2.3)
$`_i`$ denotes the partial derivative with respect to $`x_i`$,<sup>1</sup><sup>1</sup>1 We use a shorthand notation by treating $`x_i`$ as a single parameter. In reality the phase space space for each particle will be multi-dimensional. and the potential $`V`$ is the expectation value of the Hamiltonian in the state $`|\psi _x`$. The equation of motion derived from the Lagrangian is,
$`f_{ij}\dot{x}_j=_iV,`$ (2.4)
with
$`f_{ij}=_iA_j_jA_i.`$ (2.5)
Under the general condition that $`f_{ij}`$ is an everywhere invertible matrix (which in particular means that the space $``$ is even-dimensional), a Poisson bracket can be defined and a symplectic structure introduced on $``$. The bracket has the form
$`\{A,B\}=(f^1)_{ij}_iA_jB,`$ (2.6)
and the equation of motion can then be written as
$`\dot{x}_i=\{x_i,V\}.`$ (2.7)
The corresponding symplectic form is,
$`\omega ={\displaystyle \frac{1}{2}}f_{ij}dx_idx_j,`$ (2.8)
and in particular this determines the phase space volume. Thus, under the general conditions mentioned, a classical phase space can be derived from the quantum description. Note, however, that it is a generalized phase space in the sense that no configuration space has been identified.
The symplectic structure of the manifold $``$ has a simple geometric interpretation. It is defined as the imaginary part of the scalar product in the tangent space of $``$, which is obtained by projection from the Hilbert space of the full quantum system and can be written as
$`f_{ij}=2\mathrm{}\mathrm{}\{D_i\psi _x|D_j\psi _x\},`$ (2.9)
with $`D_i`$ the projected derivative,
$`|D_i\psi _x=|_i\psi _x|\psi _x\psi _x|_i\psi _x.`$ (2.10)
Written in this form, it is manifest that the symplectic form, defining the classical kinetic energy, only depends on the properties of the projected subspace. The real part of the scalar product gives another, related structure on $``$, which can be interpreted as a metric ,
$`g_{ij}=2\mathrm{}\mathrm{}\{D_i\psi _x|D_j\psi _x\}.`$ (2.11)
This construction provides a natural way to introduce a metric on the phase space, and makes it possible to discuss its geometry.
At this point we will introduce an assumption about the geometrical structure of $``$, which leads to a simplification in the discussion to follow. We assume $``$ to be a Kähler manifold. This has several technical implications, but we will only use that there is a complex structure on $``$, such that the symplectic and metric structures referred to above are the antisymmetric and symmetric part of the same complex Kähler metric. In terms of complex coordinates on $``$, we then get the following expressions for the symplectic form and the metric,
$`\omega =f_{\overline{z}_iz_j}d\overline{z}_idz_j,`$
$`ds^2=2if_{\overline{z}_iz_j}d\overline{z}_idz_j,`$ (2.12)
with
$`f_{\overline{z}_iz_j}`$ $`=`$ $`_{\overline{z}_i}A_j_{z_j}A_{\overline{i}}.`$ (2.13)
This tensor can further be expressed as
$`f_{\overline{z}_iz_j}=i_{\overline{z}_i}_{z_j}K(z,\overline{z}),`$ (2.14)
where $`K(z,\overline{z})`$ is the Kähler potential.
The condition that $``$ is a Kähler manifold is satisfied when the state vectors which define this manifold are, up to normalization, analytic functions of $`z_i`$,
$`|\psi _z=𝒩(\overline{z},z)|\varphi _z,`$ (2.15)
where $`|\varphi _z`$ denotes the analytic part of the state vector and $`𝒩(\overline{z},z)`$ is the normalization factor. The vector potentials are then given by
$`A_i`$ $`=`$ $`i\mathrm{}_{z_i}\mathrm{ln}\overline{𝒩}(\overline{z},z),`$
$`A_{\overline{i}}`$ $`=`$ $`i\mathrm{}_{\overline{z}_i}\mathrm{ln}𝒩(\overline{z},z),`$ (2.16)
and the Kähler potential is related in a simple way to the normalization factor,
$`K(\overline{z},z)=\mathrm{}\mathrm{ln}|𝒩(\overline{z},z)|^2.`$ (2.17)
3. Coherent states of identical particles
We now illustrate the general discussion by considering coherent states of the one dimensional harmonic oscillator, or equivalently, charged particles moving in two dimensions in the presence of a strong magnetic field that restricts the available states to the lowest Landau level. In this example we can explicitly derive the metric and symplectic structure, and show that they can be obtained from a Kähler potential.
We first define the coherent states for bosons and fermions and derive the corresponding classical mechanics. There is no unambiguous way to define coherent states for anyons, but we will use a construction which is very natural in this context, and again study the corresponding classical mechanics.
Since we are interested in the statistical mechanics, we also want to calculate the N-particle phase space volumes in the different cases. For this, it is necessary to start with a finite volume, and then take the thermodynamic limit. There are two obvious ways to confine the system, either by a potential, or by restricting the motion to a compact surface. In this section we shall consider the latter and study particles moving on a sphere. The case of a harmonic confining potential is treated in the Appendix.
A. Bosons and fermions in the plane
We shall use the notation of ref. and define a coherent state by translations of a minimum uncertainty reference state $`|0`$. The translation operators, $`D(z)`$, form a unitary and irreducible representation of the Heisenberg-Weyl group, and in the following we shall use the following explicit representation in terms of creation and annihilation operators,
$`D(z)=e^{za^{}\overline{z}a}=e^{\frac{1}{2}\overline{z}z}e^{za^{}}e^{\overline{z}a},`$ (3.1)
where $`[a,a^{}]=1`$, and $`z`$ is a dimensionless complex coordinate. In addition to the obvious relations $`D(z)^{}=D(z)`$ and $`D(0)=1`$, we shall need the following multiplication rule,
$`D(z_1)D(z_2)=e^{\frac{1}{2}(\overline{z}_1z_2\overline{z}_2z_1)}D(z_1+z_2).`$ (3.2)
The coherent states are now defined by,
$`|z=D(z)|0=e^{\frac{1}{2}\overline{z}z}e^{za^{}}|0.`$ (3.3)
with a reference state $`|0`$ which is annihilated by $`a`$. For convenience we shall use a notation where the normalized coherent states are labeled by $`z`$ only, although the normalization factor also depend on $`\overline{z}`$. This is to distinguish the coherent states from the position eigenstates $`|z,\overline{z}`$, and should lead to no confusion. From (3.2) we immediately get the overlap between two coherent states,
$`z_1|z_2=0|D^{}(z_1)D(z_2)|0=0|D(z_1)D(z_2)|0=e^{\frac{1}{2}(\overline{z}_1z_1+\overline{z}_2z_2)+\overline{z}_1z_2}.`$ (3.4)
An unsymmetrized basis of N-particle coherent states is defined by
$`|𝐳=|z_1,z_2,\mathrm{}z_N`$ $`=`$ $`D_1(z_1)D_2(z_2)\mathrm{}D_N(z_N)|0,0,\mathrm{}0`$
$`=`$ $`e^{\left(\frac{1}{2}_{i=1}^N\overline{z}_iz_i\right)}e^{\left(_{i=1}^Nz_ia_i^{}\right)}|\mathrm{𝟎}.`$
Normalized N-particle coherent Bose and Fermi states are symmetric and antisymmetric linear combinations respectively,
$`|𝐳,\pm =|z_1,z_2,\mathrm{}z_N_\pm `$ $`=`$ $`𝒩(z_i,\overline{z}_i){\displaystyle \frac{1}{\sqrt{N}!}}{\displaystyle \underset{P}{}}\eta _P^\pm e^{z_{i_P}a_i^{}}|\mathrm{𝟎},`$ (3.6)
where sum is over permutations, $`P`$, and the sum in the exponent over the index i is suppressed. The permutation factor, $`\eta _P`$, equals 1 for bosons and $`\pm 1`$ for fermions depending on whether the permutation is even or odd. Note that all dependence on $`\overline{z}_i`$ is in the normalization factor $`𝒩`$, which is given by,
$`|𝒩(𝐳,\overline{𝐳})|^2={\displaystyle \underset{P,P^{}}{}}\eta _P\eta _P^{}{\displaystyle \frac{1}{N!}}\mathrm{𝟎}|e^{\overline{z}_{j_P^{}}a_j}e^{z_{i_P}a_i^{}}|\mathrm{𝟎}={\displaystyle \underset{P}{}}\eta _Pe^{\overline{z}_{i_P}z_i}.`$ (3.7)
Following the general discussion in the previous section, we write the classical phase space Lagrangian (2.1) for the N-body system as,
$`L(𝐳,\overline{𝐳})=𝐳,\pm |i\mathrm{}_t\widehat{H}|𝐳,\pm ,`$ (3.8)
where $`\widehat{H}`$ is the quantum Hamiltonian. In the following we shall use the harmonic oscillator as a simple illustration, i.e. we take,
$`\widehat{H}=\mathrm{}\omega {\displaystyle \underset{i=1}{\overset{N}{}}}(a_i^{}a_i+{\displaystyle \frac{1}{2}}).`$ (3.9)
Using (3.8), (3.9) and (3.6) we get,
$`L(𝐳,\overline{𝐳})=i\mathrm{}(\dot{\overline{z}_i}_{\overline{z}_i}\mathrm{ln}𝒩\dot{z_i}_{z_i}\mathrm{ln}\overline{𝒩})\mathrm{}\omega z_i_{z_i}\mathrm{ln}|𝒩|^2{\displaystyle \frac{1}{2}}N\mathrm{}\omega .`$ (3.10)
By varying with respect to $`\overline{z}_i`$ one easily verifies that the equation of motion is that of a harmonic oscillator, i.e. $`\dot{z}_i=i\omega z_i`$. We can rewrite (3.10) on the standard form (2.2),
$`L(z_i,\overline{z}_i)={\displaystyle \frac{1}{2}}\left(A_{\overline{z}_i}\dot{\overline{z}_i}+A_{z_i}\dot{z_i}\right)V(z_i,\overline{z}_i),`$ (3.11)
and using the phase convention $`𝒩=\overline{𝒩}`$ the potentials $`A_z`$ and $`A_{\overline{z}}`$ and $`V`$ can all be obtained from the Kähler potential (2.17),
$`V(𝐳,\overline{𝐳})`$ $`=`$ $`\omega z_i_{z_i}K(𝐳,\overline{𝐳}),`$
$`A_i(𝐳,\overline{𝐳})`$ $`=`$ $`{\displaystyle \frac{i}{2}}_{z_i}K(𝐳,\overline{𝐳}),`$ (3.12)
$`A_{\overline{i}}(𝐳,\overline{𝐳})`$ $`=`$ $`{\displaystyle \frac{i}{2}}_{\overline{z}_i}K(𝐳,\overline{𝐳}),`$
with
$`K(𝐳,\overline{𝐳})=\mathrm{}\mathrm{ln}|𝒩(𝐳,\overline{𝐳})|^2=\mathrm{}{\displaystyle \underset{P}{}}\eta _Pe^{\overline{z}_{i_P}z_i}.`$ (3.13)
To get some insight into the meaning of these expressions we first analyze the two body case. Expressed in center of mass and relative coordinates, $`Z=\frac{1}{2}(z_1+z_2)`$ and $`z=z_1z_2`$, the Kähler potential reads,
$`K(Z,z,\overline{Z},\overline{z})_\pm =2\overline{Z}Z+\mathrm{ln}\left[e^{\frac{1}{2}\overline{z}z}\pm e^{\frac{1}{2}\overline{z}z}\right],`$ (3.14)
where $`\pm `$ refers to bosons and fermions respectively. The corresponding Lagrangians for the relative coordinate are now obtained using (3.11) - (3.14)
$`L_B(\overline{z},z)`$ $`=`$ $`{\displaystyle \frac{i\mathrm{}}{4}}(\overline{z}\dot{z}\dot{\overline{z}}z)\mathrm{tanh}{\displaystyle \frac{\overline{z}z}{2}}{\displaystyle \frac{\mathrm{}\omega }{2}}\overline{z}z\mathrm{tanh}{\displaystyle \frac{\overline{z}z}{2}},`$ (3.15)
$`L_F(\overline{z},z)`$ $`=`$ $`{\displaystyle \frac{i\mathrm{}}{4}}(\overline{z}\dot{z}\dot{\overline{z}}z)\mathrm{coth}{\displaystyle \frac{\overline{z}z}{2}}{\displaystyle \frac{\mathrm{}\omega }{2}}\overline{z}z\mathrm{coth}{\displaystyle \frac{\overline{z}z}{2}}.`$
Note that fermionic Lagrangian is singular in the limit of small $`r=\sqrt{\overline{z}z}`$, i.e. when the particles come close together. This is however of no physical significance, since the singular piece is a total time derivative,
$`\underset{r0}{lim}L_F={\displaystyle \frac{i\mathrm{}}{2}}(\overline{z}\dot{z}\dot{\overline{z}}z){\displaystyle \frac{1}{\overline{z}z}}={\displaystyle \frac{i\mathrm{}}{2}}_t\mathrm{ln}(z/\overline{z}),`$ (3.16)
which can be absorbed in $`𝒩`$ as a pure phase (relaxing the reality condition), or equivalently, as a pure gauge term in $`A_z`$ and $`A_{\overline{z}}`$. From (2.14) we get for the symplectic two form,
$`f_{\overline{z}z}^B`$ $`=`$ $`{\displaystyle \frac{i\mathrm{}}{2}}\left(\mathrm{tanh}{\displaystyle \frac{\overline{z}z}{2}}+{\displaystyle \frac{\overline{z}z}{2\mathrm{cosh}^2\frac{\overline{z}z}{2}}}\right)\stackrel{r\mathrm{}}{}{\displaystyle \frac{i\mathrm{}}{2}},`$ (3.17)
$`f_{\overline{z}z}^F`$ $`=`$ $`{\displaystyle \frac{i\mathrm{}}{2}}\left(\mathrm{coth}{\displaystyle \frac{\overline{z}z}{2}}{\displaystyle \frac{\overline{z}z}{2\mathrm{sinh}^2\frac{\overline{z}z}{2}}}\right)\stackrel{r\mathrm{}}{}{\displaystyle \frac{i\mathrm{}}{2}},`$
where $`z=re^{i\varphi }`$. Note that although the Kähler potential is singular in the fermion case, the metric is well defined.
In both cases, in the limit of large $`r`$, we retain the naive flat metric, $`ds^2=\mathrm{}dzd\overline{z}`$. Since we refer to relative coordinates this expression is reduced by a factor 2, as compared to the case of a single particle<sup>2</sup><sup>2</sup>2 Our normalization is such, that for a harmonic oscillator hamiltonian with $`m=\omega =1`$, $`z`$ is related to the usual coordinates and momenta by, $`z=(x+ip)/\sqrt{2\mathrm{}}`$, so $`\mathrm{}dzd\overline{z}=[(dx)^2+(dp)^2]/2`$.. (Note that the appearance of $`\mathrm{}`$ in these classical expressions is due to the use of the dimensionless coordinate $`z`$. As a more natural variable in the classical description we may use the dimensional coordinate $`\stackrel{~}{z}=\sqrt{\mathrm{}}z`$ which then would remove $`\mathrm{}`$ from the expressions.)
The more interesting limit is that of small r,
$`ds^2\stackrel{r0}{}\mathrm{}[\rho ^2d\theta ^2+d\rho ^2]`$ $`(bosons),`$ (3.18)
$`ds^2\stackrel{r0}{}{\displaystyle \frac{\mathrm{}}{3}}[\rho ^2d\theta ^2+d\rho ^2]`$ $`(fermions),`$
where we have changed variables to $`\rho =r^2/2`$ and $`\theta =2\varphi `$. We note that in these variables the metric has the standard flat-metric form in polar coordinates. Thus, the new angular variable is restricted to an interval of $`2\pi `$ (for $`\rho =0`$ to be a regular point), and therefore $`\varphi `$ is restricted to an interval of $`\pi `$. This has implications also for large separation of the particles, where the space of relative coordinates has the geometry of a cone rather than that of a plane. This is similar to the situation for the configuration space of two identical particles when this is constructed by identification of physically equivalent points . However, in the present case the space is geometrically smooth for small separation and the reduction in volume (essentially by a factor 2) compared to that of non-identical particles appears naturally from the metric of the space and not through the identification (by hand) of equivalent points.
The phase space of bosons and fermions has the same flat metric for large separation. However, for small separation it is smooth but different in the two cases. This leads to a correction to the phase space volume coming from the short-distance behaviour which is different for bosons and fermions. If we fix the maximal separation, $`R`$ <sup>3</sup><sup>3</sup>3 R is again measured in dimensionsless units, just as z, of the particles, the volume of the interior of the selected region is determined from the form of the potential $`A_z`$ and $`A_{\overline{z}}`$ on the boundary alone. With the description still restricted to relative coordinates, the phase space volume becomes<sup>4</sup><sup>4</sup>4 The phase space volume given by (3.19) is identical to the Berry phase associated with the interchange of the two particles. The close relation between phase space volume and Berry phase is the analogue of the two well-known aspects of quantum statistics: the symmetry of the wave function on one hand and the Pauli exclusion principle on the other.
$`V={\displaystyle f_{\overline{z}z}𝑑\overline{z}}dz=\left[{\displaystyle A_z𝑑z}+{\displaystyle A_{\overline{z}}𝑑\overline{z}}\right].`$ (3.19)
This gives in the two cases
$`V_B={\displaystyle \frac{1}{2}}\mathrm{}\pi R^2,V_F={\displaystyle \frac{1}{2}}\mathrm{}(\pi R^22\pi ).`$ (3.20)
The difference in phase space volume is a manifestation of the the difference in statistics in the classical description. This is readily understood from a semiclassical description where the number of states in a (single particle) phase space is identical to the volume divided by $`h`$. (The factor $`1/2`$ in (3.20) follows from the fact that we refer to relative coordinates, with the angular integration in (3.19) running from $`0`$ to $`\pi `$. The volume of one particle-space with the position of the second particle fixed would not include this factor.) In the following we will simply take the reduction in phase space volume due to the presence of the other particle as defining the statistics parameter in the classical description. We will then examine how the phase space volume of an $`N`$-particle system depends on this parameter, not only for bosons and fermions but also for intermediate values of the statistics parameter.
B. Anyons in the plane
It is well known that the harmonic oscillator coherent states states are identical to the maximally localized states of charged particles in a magnetic field projected to the lowest Landau level (LLL). The translation from harmonic oscillator to particle in the LLL is as follows
$`a`$ $`=`$ $`{\displaystyle \frac{i}{\sqrt{2}}}(\stackrel{~}{\mathrm{\Pi }}_xi\stackrel{~}{\mathrm{\Pi }}_y)\mathrm{},`$ (3.21)
$`z`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(R_xiR_y){\displaystyle \frac{1}{\mathrm{}}},`$
where (in the symmetric gauge) $`\stackrel{~}{\mathrm{\Pi }}_i=p_i\frac{1}{2}eBϵ_{ij}\widehat{x}_j`$ are the generators of magnetic translations, $`R_i`$ the guiding center coordinates, and $`\mathrm{}=(\mathrm{}/eB)^{\frac{1}{2}}`$ the magnetic length. This relation just expresses that the configuration space of charged particles in the LLL is mathematically equivalent to the phase space of a particle in one dimension. The re-interpretation of the coherent states as describing particles in the lowest Landau level is helpful in two respects. We can in a simple way generalize the coherent state representation of bosons and fermions to fractional statistics, i.e. to anyons in the lowest Landau level. We can also more easily introduce a finite volume and take the correct thermodynamic limit. Although there is no natural way to restrict particle motion in one dimension to a finite phase space, charged particles moving in a finite area penetrated by a constant magnetic field, makes perfect sense. In particular we can (in theory) study anyons moving on a compactified space, like a sphere. This problem will be studied below, but first we generalize our coherent state formalism to the case of fractional statistics.
In complex coordinates, an N-body anyon wave function has the form,
$`\mathrm{\Psi }^\nu (𝐳,\overline{𝐳})={\displaystyle \underset{i<j}{}}\left({\displaystyle \frac{\overline{z}_i\overline{z}_j}{z_iz_j}}\right)^{\nu /2}\mathrm{\Psi }_B(𝐳,\overline{𝐳}),`$ (3.22)
where $`\mathrm{\Psi }_B`$ is a totally symmetric function. In general very little is known about anyonic energy eigenstates for $`N>2`$. Exceptions are the LLL anyon states in a magnetic field which are of the form ,
$`\mathrm{\Psi }_𝐦^\nu (𝐳,\overline{𝐳})={\displaystyle \underset{i<j}{}}(\overline{z}_i\overline{z}_j)^\nu e^{\frac{1}{2}𝐳\overline{𝐳}}S_𝐦(\overline{𝐳}),`$ (3.23)
where $`𝐦=(m_1,\mathrm{}m_N)`$, $`m_i`$ integer, and
$`S_𝐦(𝐳)=𝒩_𝐦𝒮{\displaystyle \underset{𝐦}{}}z_i^{m_i}.`$ (3.24)
$`𝒮`$ is the symmetrization operator and $`𝒩_𝐦`$ a normalization constant. We now recall that the fermion and boson coherent states in (3.6), up to a normalization factor, is nothing but the projection on the lowest Landau level of the appropriately (anti)symmetrized position eigenstates,<sup>5</sup><sup>5</sup>5 This can easily be verified by explicit calculation, and it is also natural since the coherent states are the minimum uncertainty states centered around $`𝐳`$.
$`|𝐳,\pm =C_\pm (𝐳,\overline{𝐳})P_{LLL}|𝐳,\overline{𝐳},\pm ,`$ (3.25)
which implies that any bosonic of fermionic N-body wave function in the lowest Landau level can be expressed as
$`\mathrm{\Psi }^\pm (𝐳,\overline{𝐳})=𝐳,\overline{𝐳},\pm |P_{LLL}|\mathrm{\Psi }=C_\pm (𝐳,\overline{𝐳})^1𝐳,\pm |\mathrm{\Psi }.`$ (3.26)
In particular, if we are given a complete set $`\mathrm{\Psi }_𝐦^\pm (𝐳,\overline{𝐳})`$ of such LLL wave functions, we can reconstruct the N-particle coherent states by,
$`|𝐳,\pm `$ $`=`$ $`C_\pm (𝐳,\overline{𝐳}){\displaystyle \underset{𝐦}{}}|𝐦𝐦|𝐳,\overline{𝐳},\pm `$
$`=`$ $`C_\pm (𝐳,\overline{𝐳}){\displaystyle \underset{𝐦}{}}\overline{\mathrm{\Psi }}_𝐦^\pm (𝐳,\overline{𝐳})|𝐦,`$
and using that the states $`|𝐦`$ are normalized, we can express the normalization $`C_\pm (𝐳,\overline{𝐳})`$ as
$`|C_\pm (𝐳,\overline{𝐳})|^2={\displaystyle \underset{𝐦}{}}|\mathrm{\Psi }_𝐦^\pm (𝐳,\overline{𝐳})|^2=𝐳,\overline{𝐳}\pm |P_{LLL}|𝐳,\overline{𝐳},\pm .`$ (3.28)
The same procedure can now be applied to anyons in the LLL provided that we define the coherent states by the projection of the position eigenstates on the LLL. This definition is not unique, but can be shown to have several good properties . By substituting (3.23) in the relations corresponding to (S0.Ex11) and (3.28) and redefining the normalization constant by an exponential and a Jastrow factor which is common for all the wave functions, we get
$`|𝐳,\nu =𝒩_\nu (𝐳,\overline{𝐳}){\displaystyle \underset{𝐦}{}}|𝐦S_𝐦(𝐳),`$ (3.29)
with
$`|𝒩_\nu (𝐳,\overline{𝐳})|^2={\displaystyle \underset{𝐦}{}}|S_𝐦|^2=e^{\overline{𝐳}𝐳}{\displaystyle \underset{i<j}{}}|z_iz_j|^{2\nu }𝐳,\overline{𝐳},\nu |P_{LLL}|𝐳,\overline{𝐳},\nu .`$ (3.30)
These expressions are the anyonic counterparts to (3.6) and (3.7) in the case of bosons and fermions, and the derivation of the classical mechanics follows mutatis mutandis. The expressions are of course much more complicated than in the boson or fermion case, and there is no known analytic expression for the Kähler potential except in the case of two particles, where the polynomial part of the wavefunction in the relative coordinate $`z`$ is given by
$`S_m(z)={\displaystyle \frac{z^{2m+\nu }}{\sqrt{\pi 2^{2m}\mathrm{\Gamma }(2m+1+\nu )}}}.`$ (3.31)
The Kähler potential can then be calculated from (3.30) and expressed in terms of a generalized hypergeometric function,
$`K(\overline{z},z)=\mathrm{}\mathrm{ln}\left[{\displaystyle \frac{1}{\pi \mathrm{\Gamma }(1+\nu )}}F_{12}(1;{\displaystyle \frac{1}{2}}+{\displaystyle \frac{\nu }{2}},1+{\displaystyle \frac{\nu }{2}};{\displaystyle \frac{(\overline{z}z)^2}{16}})\right].`$ (3.32)
The large $`r`$ limit can be obtained from the properties of the hypergeometric function, and coincides with the result for bosons and fermions. The small $`r`$ limit can be read off directly from the leading term in (3.30),
$`\underset{r0}{lim}K(\overline{z},z)=\mathrm{}\nu \mathrm{ln}\overline{z}z+{\displaystyle \frac{\mathrm{}}{2(1+\nu )(2+\nu )}}\left({\displaystyle \frac{\overline{z}z}{2}}\right)^2+const,`$ (3.33)
giving the asymptotic metric
$`ds^2\stackrel{r0}{}={\displaystyle \frac{2\mathrm{}}{(1+\nu )(2+\nu )}}[\rho ^2d\theta ^2+d\rho ^2],`$ (3.34)
which interpolates smoothly between the expressions (3.18) for bosons and fermions. (Note that the phase space metric found in this way does not have the same physical dimension as the metric of the space in which the anyons move. Thus, there is a scale factor $`\mathrm{}^2/\mathrm{}`$ between the two metrics, where $`\mathrm{}`$ is the magnetic length.)
In spite of the complicated form of the Kähler potential $`K(𝐳,\overline{𝐳})`$, there is a simple property that immediately follows if we write it in the following form using the second identity in (3.30),
$`K(𝐳,\overline{𝐳})=\mathrm{}[\overline{𝐳}𝐳2\nu {\displaystyle \underset{i<j}{}}\mathrm{ln}|z_iz_j|+\mathrm{ln}𝐳,\overline{𝐳},\nu |P_{LLL}𝐳,\overline{𝐳},\nu ].`$ (3.35)
For a translationally invariant state (corresponding to a constant magnetic field and no external potential) only the first term can depend on the CM coordinate $`Z`$. Equivalently, we consider the limit $`z_i=Z`$ where the positions of all the particles coincide, and get
$`K(Z,\overline{Z})=N\mathrm{}\overline{Z}Z,`$ (3.36)
corresponding to a single particle of charge $`N`$, as expected. Below we shall see that this relation is altered when the particles are moving on a sphere, and the corresponding expression will allow us to calculate the pertinent phase space volume for particles with different statistics.
B. Identical particles on a sphere
In this section we will calculate the $`N`$-particle phase space volume for particles confined to a finite region. We are interested in the dependence of the volume on the particle statistics. This extends the previous discussion of the two-particle case and makes it possible to derive the (classical) statistical mechanics of the particles. As a convenient regularization of the system size, we consider a phase space with spherical geometry<sup>6</sup><sup>6</sup>6 The most natural choice for a closed two-surface seems to be the torus, since the sphere has a finite curvature. There is, however, a technical difficulty in generalizing the results in the plane to a torus in that the magnetic translation operators are not well defined for general translations due to the periodicity conditions . .
A charged particle moving on a unit sphere, penetrated by $`2j`$ units of magnetic flux has a total angular momentum $`J=j+L`$ where $`L`$ is the orbital angular momentum. The lowest Landau level corresponds to $`L=0`$ and has a $`2j+1`$ degeneracy . Again we can construct coherent states by acting on an arbitrary minimal uncertainty state, that we shall take to be $`|j,j`$, with the appropriate group elements of SU(2). We describe the sphere by stereographic projection and use a dimensonless complex coordinate, $`z`$, related to the polar angles by, $`z=\mathrm{tan}(\theta /2)e^{i\varphi }`$. It is easy to show that this $`z`$ is translated into the dimensionless $`z`$ introduced earlier in (3.21) by the substitution $`z\sqrt{\frac{1}{2j}}z`$. For ease of notation, we shall make this substitution only in the final expressions.
The SU(2) generators in the $`j`$-representation, and the corresponding coherent states are given by,
$`D(z)=e^{zJ_+}e^{\eta J_0}e^{\overline{z}J_{}},`$ (3.37)
where $`J_i`$ satisfy the standard angular momentum commutation relations and $`\eta =\mathrm{ln}(1+\overline{z}z)`$. The $`D(z)`$:s satisfy a multiplication rule similar to (3.2), but here it is sufficient to know the overlap,
$`z|w=[(1+\overline{z}z)(1+\overline{w}w)]^j(1+\overline{z}w)^{2j}.`$ (3.38)
Note that in the limit $`j=R^2/l^2\mathrm{}`$ corresponding to a large radius, or a strong magnetic field, we recover,
$`z|we^{\frac{1}{2}\overline{z}z\frac{1}{2}\overline{w}w+\overline{z}w},`$ (3.39)
where we have rescaled $`z\sqrt{\frac{1}{2j}}z`$. We can now immediately take over the results (3.6) and (3.7) for the bosonic and fermionic states in the plane ,
$`|𝐳,\pm =𝒩(𝐳,\overline{𝐳}){\displaystyle \frac{1}{\sqrt{N}!}}{\displaystyle \underset{P}{}}\eta _P^\pm e^{z_{i_P}J_+^i}|\mathrm{𝟎},`$ (3.40)
and their normalization
$`|𝒩(𝐳,\overline{𝐳})|^2={\displaystyle \underset{P,P^{}}{}}\eta _P\eta _P^{}{\displaystyle \frac{1}{N!}}\mathrm{𝟎}|e^{\overline{z}_{j_P^{}}J_{}^j}e^{z_{i_P}J_+^i}|\mathrm{𝟎}.`$ (3.41)
Using (3.37) and (3.38) we can easily calculate the relevant overlap,
$`\mathrm{𝟎}|e^{\overline{z}J_{}^j}e^{wJ_+^i}|\mathrm{𝟎}=\delta _{ij}(1+\overline{z}w)^{2j},`$ (3.42)
so
$`|𝒩(𝐳,\overline{𝐳})|^2={\displaystyle \underset{P}{}}\eta _P{\displaystyle \underset{i}{}}(1+\overline{z}_{i_P}z_i)^{2j}.`$ (3.43)
For the case of N coinciding bosons, $`z_i=z`$, we immediately get the following Kähler potential
$`K(z,\overline{z})=N\mathrm{}2j\mathrm{ln}\left(1+{\displaystyle \frac{\overline{z}z}{2j}}\right),`$ (3.44)
and the corresponding metric
$`ds^2={\displaystyle \frac{2N\mathrm{}}{(1+\frac{\overline{z}z}{2j})^2}}dzd\overline{z},`$ (3.45)
is just $`N`$ times that of a sphere. Following Manton , we can use this result to obtain the the volume of the N-boson phase space<sup>7</sup><sup>7</sup>7In a formula for the volume of the moduli space for vortices moving on an arbitrary genus Riemannian surface is derived making extensive use of theorems from complex geometry. The essential observation is that the submanifold spanned by the configurations of N coinciding bosons is a complex curve of degree N in the manifold $`CP_N`$. The metric (3.45) immediately gives the volume corresponding to this complex curve, which can be shown to be N times the area, $`A`$, obtained by integrating the fundamental two form, $`\omega `$, in (S0.Ex1) over a complex line. It then follows from a general theorem for Kähler manifolds that the total volume is given by,
$`V={\displaystyle \frac{1}{N!}}(A)^N.`$ (3.46)
In our case, using (3.44) and (S0.Ex1), we get
$`A=\mathrm{}{\displaystyle _{sph}}\omega =h2j={\displaystyle \frac{\mathrm{}4\pi R^2}{l^2}}=e\mathrm{\Phi },`$ (3.47)
which shows that the area is $`h`$ times the number of flux quanta $`\varphi _0=h/e`$ which penetrate the sphere, or equivalently, $`\mathrm{}`$ times the area in units of $`l^2`$.
In the fermion case, we first note from (3.43) that the normalization can be expressed as a determinant,
$`|𝒩(𝐳,\overline{𝐳})|^2=N!det(1+\overline{z}_iz_j)^{2j}.`$ (3.48)
In this case we cannot directly put the particles on top of each other (that would give a diverging $`𝒩`$), so we instead put $`z_i=z+\delta _i`$ and consider the limit $`\delta _i0`$. Expanding in $`\delta _i`$, we get,
$`det(1+\overline{z}_iz_j)^{2j}(1+\overline{z}z)^{2jN}detM_{ij},`$ (3.49)
where
$`M_{ij}=1+{\displaystyle \frac{2j}{(1+\overline{z}z)}}(\overline{z}\delta _i+z\overline{\delta }_j).`$ (3.50)
The matrix $`M`$ is hermitian, so the determinant is real and furthermore it has zeros for $`\delta _i=\delta _j`$. These properties together with power counting is sufficient to establish
$`detM_{ij}=C\left({\displaystyle \frac{\overline{z}z}{(1+\overline{z}z)^2}}\right)^{\frac{N(N1)}{2}}{\displaystyle \underset{i<j}{}}|\delta _i\delta _j|^2,`$ (3.51)
where $`C`$ is a $`z`$ and $`\delta `$ independent constant. Combining (3.49) and (3.51) we get, up to a constant, the Kähler potential
$`K(z,\overline{z})`$ $`=`$ $`\mathrm{ln}(1+\overline{z}z)^{2jNN(N1)}+\mathrm{ln}(\overline{z}z)^{\frac{N(N1)}{2}}`$
$``$ $`2jN\left(1{\displaystyle \frac{N1}{2j}}\right)\mathrm{ln}\left(1+{\displaystyle \frac{\overline{z}z}{2j}}\right){\displaystyle \frac{1}{2}}N(N1)\mathrm{ln}(\overline{z}z).`$
The last term can be removed by a so called Kähler gauge transformation, and corresponds to redefining the normalization constant by an analytic factor that does not contribute to the metric. The first term differs from the boson case (3.44) only by the “reduction” factor $`1(N1)/2j`$ which equals one for $`N=1`$, corresponding to a single fermion, and zero for $`N=2j+1`$ corresponding to a filled Landau level. Using the same argument as in the boson case, we get the phase space volume,
$`V_F={\displaystyle \frac{1}{N!}}\left(A(N1)h\right)^N.`$ (3.53)
Again the interpretation is clear - the available phase space for one particular fermion is reduced with an amount $`h`$ by each of the other particles present in the system. This is consistent with the semi-classical interpretation, where each quantum state is associated with a phase space volume $`h`$. Note that there is a maximum number of particles allowed, $`N=2j+1`$, in which case the phase space volume (3.53) vanishes. This corresponds to the situation where all the lowest angular momentum states are filled, i.e. to a filled lowest Landau level. Thermodynamically this state is interpreted as being incompressible, as we will discuss further in Sect. 5.
Finally we consider the case of anyons. A complete set of LLL anyon wave functions on the sphere corresponding to (3.23) in the plane is given by,
$`\mathrm{\Psi }_𝐦^\nu (𝐳,\overline{𝐳})={\displaystyle \underset{i<j}{}}(\overline{z}_i\overline{z}_j)^\nu {\displaystyle \underset{i}{}}(1+z_i\overline{z}_i)^{j+\frac{\nu }{2}(N1)}S_𝐦(\overline{𝐳}).`$ (3.54)
Following the steps leading from (3.23) to (3.30), we get the anyonic version for (3.35),
$`K(𝐳,\overline{𝐳})`$ $`=`$ $`\mathrm{}\left(j+{\displaystyle \frac{\nu }{2}}(N1)\right){\displaystyle \underset{i}{}}\mathrm{ln}\left(1+{\displaystyle \frac{z_i\overline{z}_i}{2j}}\right)`$
$``$ $`2\mathrm{}\nu {\displaystyle \underset{i<j}{}}\mathrm{ln}|z_iz_j|+\mathrm{}\mathrm{ln}𝐳,\overline{𝐳},\nu |P_{lll}𝐳,\overline{𝐳},\nu .`$
Note that in this case we do not have any explicit expression for the Kähler potential corresponding to (3.41) in the case of bosons and fermions. We can, however, again deduce the metric corresponding to configurations with coinciding anyons, i.e. $`z_i=z`$, from the above expression, again noting that the complicated overlap in the last term must be independent of $`z`$, this time because of rotational invariance. The Kähler potential becomes,
$`K(z,\overline{z})=\mathrm{}N2j\left(1{\displaystyle \frac{\nu (N1)}{2j}}\right)\mathrm{ln}\left(1+{\displaystyle \frac{2\overline{z}z}{j}}\right)+\mathrm{}`$ (3.56)
which again smoothly interpolates between the bosonic ($`\nu =0`$) and fermionic ($`\nu =1`$) results,
$`V_\nu ={\displaystyle \frac{1}{N!}}\left(A\nu (N1)h\right)^N.`$ (3.57)
The expressions we have found above for the $`N`$-particle phase space volume demonstrates a ”classical exclusion principle”. Thus, each new particle introduced in the system will find the available volume reduced by $`\alpha =\nu h`$ relative to the previous one. The quantity $`\alpha `$, i.e. the reduction in phase space volume, can be taken as defining the classical statistics parameter of the particles. In the present case it is simply the (dimensionless) quantum statistics parameter $`\nu `$ multiplied with Planck’s constant $`h`$. In other cases such a classical statistics parameter may be possible to define in terms of reduced phase space volume, even if there is no underlying point particle description. In the next section we will study such an example.
4. Vortex statistics
In the previous sections we have discussed how a phase space description, with a classical statistics parameter, can be derived from a (constrained) quantum description. In this section we will consider a somewhat different system; a classical field theory with soliton solutions. The system we have in mind is the Chern-Simons Ginzburg-Landau (CSGL) theory, originally introduced as a field theory for the quantum Hall effect , with vortices (quasi-particles) as soliton solutions. In a certain approximation the dynamics can be described in terms of vortex coordinates alone, and a phase space description can be derived from the full theory. In this description the vortices will be associated with a non-trivial classical statistics parameter, and we will show that the value of this parameter agrees with the value of the (quantum) fractional statistics parameter usually associated with quasi-particles of the quantum Hall effect. In this derivation we will make use of the close connection which exists between the CSGL Lagrangian and the Lagrangian of the relativistic abelian Higgs model discussed by Samols, Manton and others . The metric of the vortex space is the same in these two cases and it is Kähler . We make use of the results of Manton for the volume of the $`N`$-vortex space, to determine the classical statistics parameter of the vortices.
The field theory Lagrangian, which describes fields in a 2+1-dimensional space time, is
$`L={\displaystyle d^2x[i\mathrm{}\varphi ^{}D_0\varphi \frac{\mathrm{}^2}{2m}|\stackrel{}{D}\varphi |^2\frac{\lambda }{4}(|\varphi |^2\rho _0)^2+\mu \mathrm{}ϵ^{\mu \nu \rho }a_\mu _\nu a_\rho ]},`$ (4.1)
where $`\varphi `$ is a complex matter field, $`a_\mu `$ a Chern-Simons field,<sup>8</sup><sup>8</sup>8 Although the theory is non-relativistic, we choose to use a relativistic notation for the CS field with $`a_0=a^0`$, $`a_i=a^i`$ (i=1,2) and $`b=ϵ_{ij}_ia^j`$. $`m`$ a mass parameter, $`\lambda `$ the interaction strength, $`\rho _0`$ the preferred density of the system and $`\mu `$ a statistics parameter. In this classical field theory $`\mathrm{}`$ only plays the role of a dimensional parameter. For the original Laughlin quantum Hall states described by the model the statistics parameter takes the values $`\mu =1/[4\pi (2k+1)]`$ with integer $`k`$. $`D_0`$ and $`\stackrel{}{D}`$ denote covariant derivatives
$`D_0`$ $`=`$ $`{\displaystyle \frac{}{t}}+ia_0,`$
$`\stackrel{}{D}`$ $`=`$ $`i\stackrel{}{A},`$ (4.2)
with
$`\stackrel{}{A}=\stackrel{}{a}+{\displaystyle \frac{e}{\mathrm{}}}\stackrel{}{A}_{ext}.`$ (4.3)
$`\stackrel{}{A}_{ext}`$ describes a constant external magnetic field, $`B_{ext}`$, which we assume is adjusted to fit the parameter $`\rho _0`$ so that the ground state is described by a constant field $`\varphi `$ of density $`\rho _0`$ with vanishing effective magnetic field, $`B=b+\frac{e}{\mathrm{}}B_{ext}=0`$. The physical interpretation is that the system is at (or close to) the center point of a quantum Hall plateau.
It is convenient to change to dimensionless form,
$`L={\displaystyle d^2x[i\varphi ^{}D_0\varphi \frac{1}{2}|\stackrel{}{D}\varphi |^2\frac{\stackrel{~}{\lambda }}{4}(|\varphi |^21)^2+ϵ^{\mu \nu \rho }a_\mu _\nu a_\rho ]},`$ (4.4)
where the new Lagrangian is obtained from the original one by the substitutions,
$`\varphi \sqrt{\rho _0}\varphi ,`$
$`A_0{\displaystyle \frac{\mathrm{}\rho _0}{\mu m}}A_0,\stackrel{}{A}\sqrt{{\displaystyle \frac{\rho _0}{\mu }}}\stackrel{}{A},`$
$`\stackrel{}{r}\sqrt{{\displaystyle \frac{\mu }{\rho _0}}}\stackrel{}{r},t{\displaystyle \frac{\mu m}{\mathrm{}\rho _0}}t,`$
$`L\mathrm{}^2{\displaystyle \frac{\rho _0}{m}}L.`$ (4.5)
The single dimensionless parameter in the rescaled Lagrangian is $`\stackrel{~}{\lambda }=\frac{\mu m}{\mathrm{}^2}\lambda `$.
There is no independent dynamics associated with the Chern-Simons field, since variation with respect to $`a_0`$ gives a relation between magnetic field and (charge) density, which can be written as,
$`B+({\displaystyle \frac{1}{2}}\rho B_{ext})=0,`$ (4.6)
with $`\rho =|\varphi |^2`$. We assume $`B_{ext}=1/2`$ to give $`\rho =1`$ and $`B=0`$ in the ground state (all fields in dimensionless form). For finite energy configurations these values are reached asymptotically. With this assumption the constraint equation (4.6) for $`B`$ is rewritten as,
$`B+{\displaystyle \frac{1}{2}}(\rho 1)=0.`$ (4.7)
For stationary states the energy can be expressed as,
$`E={\displaystyle d^2x[\frac{1}{2}|D_+\varphi |^2+\frac{1}{2}B+(\stackrel{~}{\lambda }1)B^2]}.`$ (4.8)
with $`D_+\varphi =(D_1+iD_2)\varphi `$. The energy has the lower bound
$`E{\displaystyle d^2x\frac{1}{2}B}=N\pi ,`$ (4.9)
where $`N`$ is a non-negative integer. For the special value $`\stackrel{~}{\lambda }=1`$ of the coupling the lower bound can be saturated. The field $`\varphi `$ then satisfies the linear differential equation
$`D_+\varphi =0.`$ (4.10)
The two equations (4.7) and (4.10) define (for $`\stackrel{~}{\lambda }=1`$) stationary vortex configurations with $`N`$ vortices of equal circulation <sup>9</sup><sup>9</sup>9 We have chosen $`B_{ext}`$ to be positive. With opposite sign the lower bound is saturated with vortices of opposite circulation.. Since the configurations (for fixed $`N`$) are degenerate in energy, the vortices can be regarded as non-interacting. Gauge-equivalent configurations may naturally be considered as physically equivalent, and the vortex space can then be identified with the (moduli) space obtained from the space of field configurations after the identification of gauge-equivalent configurations . A point in vortex space is identified by the set of $`N`$ (unordered) vortex coordinates which corresponds to zeros of the field $`\varphi `$.
For values of $`\stackrel{~}{\lambda }`$ close to $`1`$ the low energy configurations correspond to slowly moving vortices. A meaningful approximation is then to impose (4.7) and (4.10) as constraints on the field configurations. The constrained fields describe a system of weakly interacting vortices. In the $`a_0=0`$ gauge the Lagrangian takes the form,
$`L={\displaystyle d^2x[i\varphi ^{}\dot{\varphi }+\dot{\stackrel{}{A}}\times (\stackrel{}{A}\stackrel{}{A}_{ext})(\stackrel{~}{\lambda }1)B^2]}N\pi .`$ (4.11)
Due to (4.7) this Lagrangian is invariant (up to a total time derivative) under time dependent gauge transformations, $`\stackrel{}{A}\stackrel{}{A}+\xi ,\varphi e^{i\xi }\varphi `$, and can therefore be interpreted as the vortex Lagrangian defined on the space of gauge-equivalent field configurations.
It is useful to consider the two fields $`\varphi `$ and $`\stackrel{}{A}`$ as components of a complex two-component field ,
$`u=\left(\begin{array}{c}\varphi \\ A\end{array}\right),`$ (4.12)
with $`A`$ as the complex field $`A=A_1+iA_2`$. A hermitian scalar product is introduced as
$`u|v={\displaystyle d^2xu^{}v}.`$ (4.13)
With this notation, up to total time derivatives the kinetic term in the Lagrangian (i.e. the part with time derivatives) can be written as
$`T=iuu_{ext}|\dot{u},`$ (4.14)
with $`u_{ext}=(0,A_{ext})`$.
The vortex configurations are described by a set of vortex coordinates $`x=\{x_1,x_2,\mathrm{},x_n\}`$, with $`x_i`$ corresponding to two real or one complex coordinate of vortex $`i`$. The precise form of the multi-vortex configurations for given coordinates is not known, but their existence is . The kinetic term for these configurations can be written
$`T=𝒜_i(x)\dot{x}_i,`$ (4.15)
where the new vector potential is
$`𝒜_i`$ $`=`$ $`iuu_{ext}|_iu`$ (4.16)
$`=`$ $`iu|_iu_iu_{ext}|u,`$
and where the differentiation now is with respect to the vortex coordinates. The corresponding field tensor which defines the symplectic form and the phase space structure of the vortex space is
$`_{ij}=_i𝒜_j_j𝒜_i=2\mathrm{}_iu|_ju.`$ (4.17)
The vortex space is the space of gauge equivalent field configurations which satisfy (4.7) and (4.10). Let us re-consider this gauge equivalence in terms of the complex fields $`u`$. The infinitesimal gauge transformations have the form
$`\delta _vu=(i\chi \varphi ,2_{\overline{z}}\chi ),`$ (4.18)
with $`\chi `$ a real function, so the vectors $`\delta _vu`$ define a real vector space. This can be extended to a complex vector space if $`\chi `$ is allowed to be complex. We will refer to the corresponding transformations as complex gauge transformations, with $`\mathrm{\Gamma }(u)`$ as the projection onto this complex subspace, and we will refer to the directions defined by the gauge transformations (4.18) as vertical directions. The orthogonal directions (horizontal directions) are defined by variations in the field $`u`$,
$`\delta _hu=(\delta \varphi ,\delta A),`$ (4.19)
which satisfy $`\delta _vu|\delta _hu=0`$, implying
$`2_z\delta A+i\varphi ^{}\delta \varphi =0.`$ (4.20)
The real part is
$`\delta \stackrel{}{A}+{\displaystyle \frac{i}{2}}(\varphi ^{}\delta \varphi \varphi \delta \varphi ^{})=0,`$ (4.21)
and the imaginary part is
$`\delta B+{\displaystyle \frac{1}{2}}\delta \rho =0.`$ (4.22)
Changes in the fields which follows from variations in the vortex coordinates will automatically satisfy the second equation, (4.22), due to the constraint (4.7). The first equation, (4.21), can be satisfied provided we make use of the freedom to include real gauge transformations in the variations of the fields. Thus, we may assume both these equations, or equivalently (4.20) to be satisfied when $`\delta A`$ and $`\delta \varphi `$ are replaced by the corresponding derivatives with respect to vortex positions.
We introduce $`\mathrm{\Pi }(u)=I\mathrm{\Gamma }(u)`$ as the projection onto the horizontal directions and $`D_i=\mathrm{\Pi }_i`$ as the projected derivative. With the assumption that the gauge condition (4.21) is satisfied the vector potential $`𝒜_i`$ can be written as
$`𝒜_i=iuu_{ext}|D_iu.`$ (4.23)
It transforms under a (complex) gauge transformation $`\chi `$ as
$`𝒜_i𝒜_{}^{}{}_{i}{}^{}=𝒜_i_i\mathrm{\Theta }.`$ (4.24)
with
$`\mathrm{\Theta }={\displaystyle d^2x[\varphi ^{}\varphi +2i_{\overline{z}}(A^{}A_{ext}^{})]\chi }.`$ (4.25)
This means that $`_{ij}`$ is invariant under complex gauge transformations and can be written as
$`_{ij}=2\mathrm{}D_iu|D_ju,`$ (4.26)
since the difference between $`_iu`$ and $`D_iu`$ can (locally) be eliminated by a (vortex-position dependent) gauge transformation.
We may consider the quantity
$`\eta _{ij}=D_iu|D_ju`$ (4.27)
as defining a hermitian metric tensor for the vortex space. It is obtained from the scalar product (4.13) by projection on the horizontal directions. The tensor $`_{ij}`$, which is derived from the kinetic part of the Lagrangian, and which defines the phase space structure of the vortex space, can now be identified with the imaginary part of this metric tensor.
The metric (4.27) is relevant also for the relativistic abelian Higgs model, as we shall now demonstrate. The Lagrangian of this model has the form
$`L={\displaystyle d^2x[\frac{1}{2}(D_\mu \varphi )^{}D^\mu \varphi \frac{1}{4}F_{\mu \nu }F^{\mu \nu }\frac{2\stackrel{~}{\lambda }1}{8}(|\varphi |^21)^2]}.`$ (4.28)
It is quadratic in time derivatives and has the Chern-Simons field replaced by a Maxwell field. The energy has the same lower bound (4.9) as the GLCS model, and for $`\stackrel{~}{\lambda }=1`$ this lower limit is saturated if both the equations (4.7) and (4.10), known as the Bogomolny equations, are satisfied. Thus, if these equations are used to define the $`N`$-vortex space, the vortex space is the same for the two models. However, the kinematics, as defined by the kinetic part of the Lagrangian is not the same in the two cases. The non-relativistic model is linear in time derivatives, which means that the vortex space has the character of a phase space, while the relativistic model is quadratic in time derivatives, which gives the vortex space the character of a configuration space.
Expressed in the $`A_0=0`$ gauge and constrained by the Bogomolny equations, the Lagrangian of the Higgs model has the form
$`L^{}=T^{}V,`$ (4.29)
with
$`T^{}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle d^2x[\dot{\varphi }^{}\dot{\varphi }+\dot{\stackrel{}{A}}\dot{\stackrel{}{A}}]},`$
$`V`$ $`=`$ $`{\displaystyle d^2x(\stackrel{~}{\lambda }1)B^2}+N\pi ,`$ (4.30)
and the fields constrained by Gauss’ law
$`\dot{\stackrel{}{A}}+{\displaystyle \frac{i}{2}}(\varphi ^{}\dot{\varphi }\dot{\varphi }^{}\varphi )=0.`$ (4.31)
The potential $`V`$ is the same, but the kinetic term $`T^{}`$ is different from that of the CSGL model. We note that the constraint (4.31) corresponds to the real part (4.21) of the condition for motion in horizontal direction. Also the imaginary part (4.22) is satisfied due to the Bogomolny equations. The constraint on the motion, given by Gauss’ law, can be expressed in terms of the projection $`\mathrm{\Pi }`$ and the field $`u`$ as
$`T^{}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\dot{u}|\mathrm{\Pi }|\dot{u}`$ (4.32)
$`=`$ $`{\displaystyle \frac{1}{2}}\dot{x}_i\dot{x}_jD_iu|D_ju.`$
The kinetic term $`T^{}`$ is different, but related to the kinetic term $`T`$ of the CSGL model. $`T^{}`$ is determined by the real part, whereas $`T`$ is determined (up to a gauge transformation) by the imaginary part of the same hermitian metric (4.27). This metric has been examined in some detail by Samols for the case of the abelian Higgs model. It is a Kähler metric, which is conveniently expressed in terms of complex vortex coordinates as
$`ds^2=2i_{\overline{z}_iz_j}d\overline{z}_idz_j,`$ (4.33)
with
$`_{\overline{z}_iz_j}=i\left[D_{\overline{z}_i}u|D_{z_j}uD_{z_j}u|D_{\overline{z}_i}u\right].`$ (4.34)
The corresponding Kähler 2-form is
$`\omega =_{\overline{z}_iz_j}d\overline{z}_idz_j.`$ (4.35)
The Kähler form determines the symplectic structure and the volume of the vortex space. This (2$`N`$-dimensional) volume is the same whether the vortex space is considered as a configuration space (i.e. with the volume determined by the real part of the metric) or as a phase space (with the volume determined by the imaginary part), and has been calculated by Manton for $`N`$ vortices on a sphere. The result is <sup>10</sup><sup>10</sup>10 Manton’s result has been changed with a factor $`(2\pi )^N`$ to fit the definition of the metric in this paper. We also express the $`N`$-vortex volume in terms of the single-vortex volume $`A`$ rather than the area of the sphere. This gives a factor $`N1`$ instead of $`N`$ in the second term of Eq.(4.36).
$`V_N={\displaystyle \frac{1}{N!}}(A8\pi ^2(N1))^N,`$ (4.36)
with $`A`$ as the volume (area) of the one-vortex space. The volume (4.36) has the same form as discussed in Sect. 3 for $`N`$ identical particles with non-trivial classical statistics. With the vortex space interpreted as a phase space, the statistics of the vortices can be extracted from the reduction term $`8\pi ^2N`$. In order to do so correctly we have to re-write the volume (4.36) in dimensional form. The phase space volume is determined by the Lagrangian (4.4), and as follows from the transformations (LABEL:rescale) the phase space dimensions are correctly re-introduced by the substitutions $`A\mu \mathrm{}A`$ and $`V_N(\mu \mathrm{})^NV_N`$. In dimensional form the expression for the $`N`$-vortex volume is
$`V_N={\displaystyle \frac{1}{N!}}(A4\pi \mu h(N1))^N,`$ (4.37)
and the classical statistics parameter as determined by the reduction in available phase space due to the presence of other vortices therefore is
$`\alpha =4\pi \mu hgh.`$ (4.38)
We can interpret $`g`$, the classical parameter divided by $`h`$, as the dimensionless quantum statistics parameter. The value $`g=4\pi \mu `$ agrees with the value of the statistics parameter as determined from Berry phase calculations with Laughlin wave functions , or from the properties of vortices in the CSGL model .
5. Statistical mechanics
A. Entropy and pressure of identical particles
We have already emphasized that even if the classical equations of motion, and thus the classical dynamics, does not depend on the classical statistics parameter $`\alpha `$, the statistical mechanics (and thus the thermodynamics), does. In this section we demonstrate this by first calculating the entropy and pressure in the two model systems considered in sections 3 and 4, charged particles in the lowest Landau level, and vortices in the CSGL model respectively. The results of this calculation fit nicely into the general framework of “fractional exclusion statistics” for particles with degenerate energy levels, and we briefly review the basics of this topic before further discussing our results.
We assume the interaction strength to be $`\stackrel{~}{\lambda }=1`$ for the CSGL theory. The vortex system is then degenerate in energy. That is also the case for a system of anyons in the lowest Landau level. Thus, both these systems have the special property that the energy does not depend on the state, but only on the number of particles. This means that the statistical mechanics is determined by the phase space volume $`V_N`$, which has been determined in the previous sections, and by the energy $`E_N`$. The classical partition function is simply the total number of states, $`V_N/h^N`$ multiplied with the Boltzmann factor, i.e.
$`Z_N={\displaystyle \frac{V_N}{h^N}}e^{\beta E_N},`$ (5.1)
The following simple expressions for the free energy $`F`$ and the entropy $`S`$, immediately follow,
$`F`$ $`=`$ $`E_NT\mathrm{ln}(V_N/h^N),`$ (5.2)
$`S`$ $`=`$ $`\mathrm{ln}(V_N/h^N).`$
where is Boltzmann’s constant is set to unity. The pressure is usually defined by $`P=(F/𝒱)_T`$, where $`𝒱`$ is the volume of real space, but in the systems we have considered the real two-dimensional space where the particles or vortices move is proportional to the phase space, so we simply define the pressure as,
$`P=\left({\displaystyle \frac{F}{A}}\right)_T=T{\displaystyle \frac{\mathrm{ln}(V_N/h^N)}{A}},`$ (5.3)
where $`A=V_1`$ is the phase space volume for a single particle. Substituting the results (3.57) or (4.37), we get
$`S`$ $`=`$ $`N\mathrm{ln}(1\alpha \rho )+N\mathrm{ln}{\displaystyle \frac{A}{h}}N\mathrm{ln}N+N,`$ (5.4)
$`\beta P`$ $`=`$ $`{\displaystyle \frac{\rho }{1\alpha \rho }},`$ (5.5)
where $`\alpha =\nu h`$ or $`gh`$, and where we have introduced the classical phase space density $`\rho =N/A`$ and neglected the difference between $`N`$ and $`N1`$, which is irrelevant in the thermodynamic limit.
The expression (5.5) shows that there is a maximum density $`\rho =1/\alpha `$ allowed by the system, which corresponds to an infinite pressure and therefore to an incompressible state. For the phase space volume this means $`V_N=0`$, i.e. there is no available phase space volume for any new particle added to the system. For the anyon system this situation corresponds to a completely filled Landau level. What is unusual about this is that the blocking, which can be interpreted as representing a generalized Pauli principle, shows up not only in the quantum but also in the classical description of the system.
B. Exclusion statistics and the classical limit
The generalization of the Pauli exclusion principle introduced by Haldane , usually called exclusion statistics, states that in the presence of particles in a set of given quantum states, the number of available one-particle states for any new particle added to the system is reduced. More precisely, the addition of $`\mathrm{\Delta }N`$ particles changes the number of available states, $`d_N`$ according to
$`\mathrm{\Delta }d_N=g\mathrm{\Delta }N,`$ (5.6)
where $`g`$ is the exclusion statistics parameter. The statistical weight, or number of states available for the full $`N`$-particle system, is given by the formula
$`W_N={\displaystyle \frac{(G+(1g)(N1))!}{N!(GgN(1g))!}},`$ (5.7)
where $`G`$ is the number of single-particle states. Clearly (5.6) and (5.7) reduces to standard expressions when $`g=0`$ (for bosons, with no exclusion) and $`g=1`$ (for fermions, with total exclusion). There exist some (theoretical) realizations of exclusion statistics for particles in one dimension (i.e. with two-dimensional phase space) for $`g`$ different from these two values. One particular case is the system of anyons confined to the lowest Landau level, which we have already considered . In that case the exclusion statistics parameter $`g`$ is identical to the anyon statistics parameter $`\nu `$.
The statistical mechanics of particles with exclusion statistics can be derived from the statistical weight (5.7) when the total energy can be written as a sum of single-particle energies and (5.7) is applied separately to (single-particle) energy levels . The result for the entropy is
$`S`$ $`=`$ $`{\displaystyle \underset{k}{}}D_k\{[1+(1g)n_k]\mathrm{ln}[1+(1g)n_k]`$
$`+`$ $`(1gn_k)\mathrm{ln}(1gn_k)n_k\mathrm{ln}n_k\},`$
where the sum runs over single-particle energy states. $`D_k`$ is the degeneracy of the $`k`$-th level and the quantum distribution function, $`n_k`$, is the average occupation number of the state $`k`$.
Since each quantum state occupies the phase space volume $`h^𝒟`$, with $`2𝒟`$ the dimension of the single-particle phase space, we can relate $`n`$ and $`\rho `$ in the semiclassical limit by $`n=\rho h^𝒟`$. In the Boltzmann limit, $`h0`$ and $`n0`$, all dependence on $`g`$ in (S0.Ex22) goes away. If we, however, define the classical physics by the double limit $`h0`$, $`g\mathrm{}`$ and $`gh^𝒟\alpha `$, where $`\alpha `$ is interpreted as a classical statistics parameter<sup>11</sup><sup>11</sup>11 Such a way of taking the classical limit is well-known from other contexts. Thus, a charged particle can in the quantum mechanical description be characterized by a dimensionless charge $`g=q/\sqrt{\mathrm{}c}`$, where $`q`$ is the physical charge. (For $`q=e`$ we have $`g^2=4\pi \alpha `$, with $`\alpha `$ the fine structure constant.) With $`g`$ fixed the charge $`q`$ depends on $`\mathrm{}`$ and vanishes in the limit $`\mathrm{}0`$. However, if the classical limit is taken as $`\mathrm{}0,g\mathrm{}`$ with $`g\sqrt{\mathrm{}c}q`$ the dimensional charge $`q`$ survives the classical limit. (S0.Ex22) gets a nontivial limit of
$`S={\displaystyle \underset{k}{}}D_kh^𝒟\left[\rho _k\mathrm{ln}(1\alpha \rho _k)\rho _k\mathrm{ln}(\rho _kh)+\rho _k\right].`$ (5.9)
If we further specialize to the case of fully degenerate states in a two-dimensional phase space, where the sum is simply replaced by the total number of available single-particle states, $`G=A/h`$, and where $`\rho _k`$ is replaced by $`N/A`$, we exactly regain (5.4). This demonstrates that the classical statistical mechanics discussed in the previous section can be regarded as a special limit of exclusion statistics, different from the Boltzmann limit.
An alternative way to see the correspondence is to start from the the equation of state for exclusion statistics particles with the same energy,
$`\beta P={\displaystyle \frac{G}{𝒱}}\mathrm{ln}\left(1+{\displaystyle \frac{n}{1gn}}\right),`$ (5.10)
where $`n=N/G`$. Introducing the density $`\stackrel{~}{\rho }=N/𝒱`$ and taking the double limit defined above we get,
$`\beta P={\displaystyle \frac{\stackrel{~}{\rho }}{1\alpha \stackrel{~}{\rho }\frac{𝒱}{V_1}}}.`$ (5.11)
If we identify the the physical volume $`𝒱`$ with the one particle phase space volume $`A`$, so that $`\stackrel{~}{\rho }=\rho `$, we reproduce (5.5).
In the Appendix it is shown that even in a non-degenerate case, with particles in a harmonic oscillator potential, the classical statistical mechanics, defined as in Sect. 5A, coincides with that of exclusion statistics when the classical limit is taken in the way discussed above. Clearly what is important for the connection with exclusion statistics is the two defining relations (5.6) and (5.7) which determine the number of states in the system. In the classical description they are represented by the expressions for the phase space volume, and by taking the limit $`h0,gh^𝒟\alpha `$ it is straightforward to demonstrate that (5.6) and (5.7) reproduce the expressions for the phase space volume derived in Sects. 3 and 4.
6. Discussion
In this paper we have described a way to encode the particle statistics in the classical Lagrangian of a many particle system. The important point is that the Lagrangian includes more information about the system than just the classical equation of motion. It also gives information about the volume of the phase space, which in the quantum description corresponds to the number of states. If the $`N`$-particle volume can be determined as a function of the single particle volume, a classical statistics parameter can be defined as the reduction in available phase space volume for one particle by the presence of the others. Viewed in this way this classical statistics can be regarded as an analogue of exclusion statistics. In the specific examples we have considered, this relation can be made more specific and the classical statistical mechanics derived from this can be seen as a special way to take classical limit of exclusion statistics.
To make this idea more precise we have considered cases where the classical mechanics can be derived from the quantum description by constraining the motion in Hilbert space to (generalized) coherent states. For bosons, fermions and even anyons with a two-dimensional phase space the Lagrangian can be derived and the phase space volume can be calculated. The dimensional, classical statistics parameter, defined as the volume occupied by each particle present, in these cases are simply the dimensionless quantum statistics parameter multiplied with Planck’s constant $`h`$. In another example, vortices in CSGL theory, there is no such underlying point particle description, but a similar classical Lagrangian can be found and the classical statistics parameter can be related to the coupling of the Chern-Simons term.
There are several interesting questions raised by this description:
$``$ Is the ”classical fermion” description useful in some cases? This description would correspond to retaining the fermions’ ability to occupy phase space, but otherwise treat them as classical particles. (A classical electron would then be characterized both by a charge and a (classical) statistics parameter.) Can the description give a useful approximation for other objects, like vortices in superfluids or superconductors?
$``$ In the examples we have studied the phase space is two-dimensional, but the formalism (like for exclusion statistics) does not seem to depend in any crucial way on dimension. Are there non-trivial higher-dimensional examples? (Fermions in two and three dimensions can certainly be represented like this.)
$``$ What about quantizing such a classical theory? In the cases we have studied, with a Kähler metric defined on phase space, a quantum description can presumably be derived in a unique way by use of analyticity properties. When regarded as a ”re-quantization” of the system, how does it relate to the original quantum description. What would in particular the quantum description of the CSGL vortices be?
All these questions seem to merit further investigation.
Appendix: Harmonic oscillator potential
In this Appendix we consider the statistical mechanics of particles in a harmonic oscillator potential. The particles are ”classical anyons” in the sense discussed in Sect. 3, i.e. the one derived from quantum anyons in the lowest Landau level. The system can also be interpreted as a coherent state representation of particles in a one-dimension harmonic oscillator potential, in a form interpolating between bosons and fermions. We calculate the partition function of the $`N`$-particle system and show that this is related to the partition function of a (quantum) system of particles with exclusion statistics in a harmonic oscillator potential by the same correspondence as obtained in Sect. 4.
The wave functions of the lowest Landau level have the form
$`\psi (𝐳,\overline{𝐳})={\displaystyle \underset{i<j}{}}(\overline{z}_i\overline{z}_j)^\nu f(\overline{𝐳})e^{\frac{1}{2}\overline{𝐳}𝐳},`$ (A.1)
with $`f(\overline{𝐳})f(z_1,\mathrm{},z_N)`$ as a general anti-analytic function of the complex particle coordinates. It is assumed to be symmetric in the variables. We introduce analytic basis vectors by
$`𝐳|\psi =f(\overline{𝐳}).`$ (A.2)
The basis vectors $`|z_1,\mathrm{},z_N`$ are not normalized, but we assume $`\nu `$ to be chosen such that they are regular and non-vanishing at points of coincidence of particle positions. Normalized vectors are introduced by
$`|\psi _{𝐳,\overline{𝐳}}`$ $`=`$ $`𝒩_{𝐳,\overline{𝐳}}|𝐳,`$
$`|𝒩_{𝐳,\overline{𝐳}}|^2`$ $`=`$ $`𝐳|𝐳.`$ (A.3)
Defined in this way $`|𝒩_{𝐳,\overline{𝐳}}|^2`$ is a regular function with no zeros anywhere in $`N`$-particle space, and the Kähler potential $`K=\mathrm{ln}|𝒩|^2`$ is a regular function everywhere.
The Hamiltonian depends on two frequencies, the cyclotron frequency $`\omega _c`$ determined by the external magnetic field and the frequency $`\omega _0`$ of the additional harmonic oscillator potential. When acting on the anti-analytic part $`f(\overline{𝐳})`$ of the wave functions of LLL, the Hamiltonian has the form
$`H`$ $`=`$ $`\mathrm{}(\omega _t\omega _c){\displaystyle \underset{i}{}}\overline{z}_i_{\overline{z}_i}+\mathrm{}\omega _t\left[{\displaystyle \frac{\nu }{2}}N(N1)+{\displaystyle \frac{N}{2}}\right]`$ (A.4)
$`=`$ $`\mathrm{}\omega {\displaystyle \underset{i}{}}\overline{z}_i_{\overline{z}_i}+V_N^0,`$ (A.5)
with $`\omega _t=\sqrt{\omega _c^2+\omega _0^2}`$, $`\omega =\omega _t\omega _c`$ and $`V_N^0`$ the quantum mechanical ground state energy
$`V_N^0=\mathrm{}\omega _t\left[{\displaystyle \frac{\nu }{2}}N(N1)+{\displaystyle \frac{N}{2}}\right].`$ (A.6)
For a system of particles in a one-dimensional harmonic oscillator potential the Hamiltonian is essentially the same, except that it depends on a single frequency $`\omega _0`$,
$`H=\mathrm{}\omega _0{\displaystyle \underset{i}{}}\overline{z}_i_{\overline{z}_i}+\mathrm{}\omega _0\left[{\displaystyle \frac{\nu }{2}}N(N1)+{\displaystyle \frac{N}{2}}\right].`$ (A.7)
Thus the difference between these two cases is only an overall $`N`$-dependent shift of the energy spectrum.
The energy of the classical description is determined by the matrix elements of the Hamiltonian (A.4),
$`V(𝐳)`$ $`=`$ $`𝐳|H|𝐳|𝒩_{𝐳,\overline{𝐳}}|^2`$
$`=`$ $`\left[\left\{\mathrm{}\omega {\displaystyle \underset{i}{}}\overline{z}_i_{\overline{z}_i}+\mathrm{}\omega _t\left[{\displaystyle \frac{\nu }{2}}N(N1)+{\displaystyle \frac{N}{2}}\right]\right\}|𝒩_{𝐳,\overline{𝐳}}|^2\right]|𝒩_{𝐳,\overline{𝐳}}|^2`$
$`=`$ $`\mathrm{}\omega {\displaystyle \underset{i}{}}\overline{z}_i_{\overline{z}_i}\mathrm{ln}|𝒩_{𝐳,\overline{𝐳}}|^2+\mathrm{}\omega _t\left[{\displaystyle \frac{\nu }{2}}N(N1)+{\displaystyle \frac{N}{2}}\right],`$
and the $`N`$-particle partition function is
$`Z_N={\displaystyle \frac{1}{h^N}}{\displaystyle \frac{\omega ^N}{N!}e^{\beta V}},`$ (A.9)
where $`\omega `$ is the symplectic form
$`\omega =f_{\overline{z}_iz_j}d\overline{z}_idz_j,`$ (A.10)
with
$`f_{\overline{z}_iz_j}=i\mathrm{}_{\overline{z}_i}_{z_j}\mathrm{ln}|𝒩_{𝐳,\overline{𝐳}}|^2.`$ (A.11)
The form of the energy makes it possible to evaluate the integrals in the expression for the partition function. We write it as
$`Z_N={\displaystyle \frac{e^{\beta V_N^0}}{\pi ^NN!}}ϵ_{ij\mathrm{}k}{\displaystyle d^2z_1\mathrm{}d^2z_N[_{\overline{z}_i}_{z_1}\mathrm{ln}|𝒩_{𝐳,\overline{𝐳}}|^2\mathrm{}_{\overline{z}_k}_{z_N}\mathrm{ln}|𝒩_{𝐳,\overline{𝐳}}|^2]}`$
$`\times \mathrm{exp}\{\beta [\mathrm{}\omega {\displaystyle \underset{i}{}}z_i_{z_i}\mathrm{ln}|𝒩_{𝐳,\overline{𝐳}}|^2]\}.`$
The partition function can be rewritten as
$`Z_N={\displaystyle \frac{1}{(\beta \mathrm{}\omega )}}{\displaystyle \frac{e^{\beta V_N^0}}{\pi ^NN!}}ϵ_{ij\mathrm{}k}{\displaystyle }d^2z_1\mathrm{}d^2z_N{\displaystyle \frac{1}{z_1}}_{\overline{z}_i}\left[\mathrm{ln}\right|𝒩_{𝐳,\overline{𝐳}}|^2\mathrm{}`$
$`\times _{\overline{z}_k}_{z_N}\mathrm{ln}|𝒩_{𝐳,\overline{𝐳}}|^2\mathrm{exp}\{\beta (\mathrm{}\omega {\displaystyle \underset{i}{}}z_i_{z_i}\mathrm{ln}|𝒩_{𝐳,\overline{𝐳}}|^2)\}],`$
and by use of the identity
$`{\displaystyle \frac{1}{z_1}}_{\overline{z}_i}=_{\overline{z}_i}{\displaystyle \frac{1}{z_1}}\pi \delta (z_1)\delta _{i1}`$ (A.14)
the integration over $`z_1`$ can be performed
$`Z_N={\displaystyle \frac{\pi }{\beta \mathrm{}\omega }}{\displaystyle \frac{e^{\beta V_N^0}}{\pi ^NN!}}ϵ_{j\mathrm{}k}{\displaystyle }d^2z_2\mathrm{}d^2z_N[_{\overline{z}_j}_{z_2}\mathrm{ln}|𝒩_{𝐳,\overline{𝐳}}|^2\mathrm{}`$
$`\times _{\overline{z}_k}_{z_N}\mathrm{ln}|𝒩_{𝐳,\overline{𝐳}}|^2]\mathrm{exp}\{\beta [\omega {\displaystyle \underset{i}{}}z_i_{z_i}\mathrm{ln}|𝒩_{𝐳,\overline{𝐳}}|^2]\}.`$
The $`(N1)`$-particle integral in this expression is of the same form as the original $`N`$-particle integral, and by repeating the procedure $`N`$ times we get the following simple expression for the partition function
$`Z_N`$ $`=`$ $`{\displaystyle \frac{1}{(\beta \mathrm{}\omega )^NN!}}e^{\beta V_N^0}`$ (A.16)
$`=`$ $`{\displaystyle \frac{1}{(\beta \mathrm{}\omega )^NN!}}\mathrm{exp}\{\beta \mathrm{}\omega _t\left[{\displaystyle \frac{\nu }{2}}N(N1)+{\displaystyle \frac{N}{2}}\right]\}.`$
The classical expression for the partition function can be compared with the partition function of the quantum system
$`Z_N=Tre^{\beta H},`$ (A.17)
with $`H`$ given by (A.4). This expression is easily evaluated, since it can be written as
$`Z_N`$ $`=`$ $`e^{\beta V_N^0}{\displaystyle \underset{l_1=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{l_2=l_1}{\overset{\mathrm{}}{}}}\mathrm{}{\displaystyle \underset{l_N=l_{N1}}{\overset{\mathrm{}}{}}}e^{\beta \mathrm{}\omega \underset{i}{}l_i}`$ (A.18)
$`=`$ $`e^{\beta V_N^0}[{\displaystyle \underset{n=1}{\overset{N}{}}}(1e^{n\mathrm{}\beta \omega }]^1.`$
This expression shows that in the limit $`\mathrm{}0`$, with $`\mathrm{}\nu `$ fixed, the partition function (A.18) of the quantum system coincides with the classical partition function (S0.Ex31). (Note however that the classical function depends on $`\mathrm{}`$ explicitly, not only through the statistics factor $`\alpha =h\nu `$, due to the contribution from the ground state energy.)
It is well known that the system of particles in the lowest Landau level can be regarded as a special realization of exclusion statistics , and the correspondence between the two partition functions discussed here is therefore essentially the same as the correspondence between the classical statistical mechanics and the statistical mechanics of particles with exclusion statistics discussed in Sect. 4. If we use the harmonic oscillator as a volume regulator the relation between the discussion in this Appendix and in Sect. 4 becomes even more direct. The thermodynamic limit is here taken by interpreting the limit $`\omega _00`$ in a specific way . For the quantum case the harmonic oscillator regulator has been used in , and the expressions for the entropy and equation of state of anyons in the LLL were found in this way. Due to the correspondence between the quantum and classical descriptions, the thermodynamic limit of the classical functions with the harmonic oscillator regularization, will be identical to the corresponding functions of Sect. 4. That is what should be expected, since for the thermodynamic limit it should be of no significance whether volume regularization is done by confinement to a sphere or by confinement in a harmonic oscillator potential.
Acknowledgment
Serguei Isakov acknowledges the support received through a NATO Science Fellowship granted by the Norwegian Research Council. S.I. also appreciates warm hospitality of NORDITA during his stay there in the summer of 1999, where part of this work was done.
Three of us (T.H.H., S.I. and J.M.L.) would also like to thank the Department of physics, Norwegian University of Science and Technology (Trondheim) for the inspiring atmosphere of the Workshop on low-dimensional physics in June 1999.
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# Correlation energy of the pairing Hamiltonian
## I Introduction
An important goal of nuclear theory is to predict nuclear binding energies. The mean-field approximation certainly provides a good starting point, but the correlation energy associated with nearly degenerate configurations can be of the order of several MeV. Since one would like an order of magnitude better accuracy for global modeling, the correlations must be treated with some care. The correlations associated with broken mean-field symmetries, namely center-of-mass localization, deformations, and pairing, are especially important and need to be singled out for special attention. The study of such correlation effects has attracted much interest, see Refs. for recent citations. Popular computational methods <sup>*</sup><sup>*</sup>*One may also ask whether simplified exact solutions may be of use for realistic applications. The Richardson solution for a pairing Hamiltonian is useful in this respect, but it requires a state-independent pairing interaction, which may not be realistic enough. That solution has recently been generalized to separable pairing interactions, but it is not yet clear whether the computational simplicity is preserved. include the generator coordinate method (GCM), variation-with-projection methods (e.g., VAP) , and the RPA .
In previous work, we have advocated the RPA approach for deformations and for center-of-mass localization. The correlation energy is calculated from the RPA formula
$$E_{corr}=\frac{1}{2}\left(\underset{i}{}\mathrm{}\omega _iTr(A)\right),$$
(1)
where $`\omega _i`$ is the frequency of the RPA phonon and $`A`$ is the $`A`$ matrix in the RPA equations. This method requires solving the RPA equations for each nucleus, which is computationally easy if the interaction is assumed to be separable. On other other hand, methods such as VAP and GCM are rather complicated to apply. One should also remember that a global theory requires a method that is not only systematic but also preserves the relative computational simplicity of the mean field approximation. In this respect, the Lipkin-Nogami method is a popular candidate for a computationally easy method to go beyond the mean field pairing theory .
In this paper, we continue our study of the accuracy of the RPA correlation formula (1), going on to pairing correlations. There is a long history of the RPA treatment of pairing correlations. An early study by Bang and Krumlinde showed that the RPA formula reproduces the exact correlation energy rather well in a schematic model. Kyotoku et al. compared the leading $`1/N`$ behavior of different methods including the Lipkin-Nogami method. They found that only RPA gave the exact coefficient in the condensed phase. The RPA method has in fact been used in realistic models of deformed nuclei. The RPA correlation in the normal phase was studied in Ref. using the self-consistent version of RPA. But perhaps surprisingly in view of the large literature, we have not found any studies that specifically compare the RPA with the computationally attractive alternative methods, testing the behavior across shell closures and for odd $`N`$ systems. It is important that the method should not introduce spurious discontinuities when the mean field solution changes character; otherwise separation energies could not be reliably calculated .
To investigate the applicability of Eq. (1) and study the effects of symmetry breaking restoration, in this paper we employ a well-known schematic two-level pairing model, and compare its exact solutions to the approximations that are computationally attractive. The details of the RPA correlation energy are given in the next section. They include both the quasi-particle RPA extension of the BCS theory and the RPA for the pairing Hamiltonian when the strength is too weak for the mean field approximation to support the BCS solution (the “pairing-vibration” regime). The specific application of the various methods to the two-level model is given in the following section III.
## II Pairing Hamiltonian and the RPA
The pairing Hamiltonian and the BCS solution are well-known and we just summarize the equations to confirm the standard notation. In this paper we consider a Hamiltonian with an arbitrary single-particle term but a pairing interaction whose strength is state-independent,
$$H=\underset{j}{}ϵ_j\widehat{N}_jG\underset{j,j^{}}{}A_j^{}A_j^{},$$
(2)
where $`ϵ_j`$ is a single-particle energy. The number operator $`\widehat{N}_j`$ for the shell $`j`$ is given by
$$\widehat{N}_j=\underset{m>0}{}\left(a_{jm}^{}a_{jm}+a_{j\overline{m}}^{}a_{j\overline{m}}\right).$$
(3)
Here, $`a_{j\overline{m}}`$ is the annihilation operator for the time-reversal state and is given by $`a_{j\overline{m}}=()^{jm}a_{jm}`$. The pairing operator $`A_j^{}`$ in the Hamiltonian (2) is given by
$$A_j^{}=\underset{m>0}{}a_{jm}^{}a_{j\overline{m}}^{},$$
(4)
and $`A_j`$ is the Hermitian conjugate of $`A_j^{}`$. The exact solutions of the Hamiltonian were obtained a long time ago by Richardson and Sherman. In Appendix A, we solve the Richardson equation for a simple, but non-trivial case where there are two pairs in a single $`j`$-shell.
The BCS theory is the mean field solution to Eq. (2). We remind the reader that it is derived variationally from a ground state wave function of the form
$$|BCS=\underset{j,m>0}{}(u_j+v_ja_{jm}^{}a_{j\overline{m}}^{})|,$$
(5)
where $`v_j`$ is the usual pair occupation amplitude and $`u_j`$ satisfies $`u_j^2+v_j^2=1`$. This wave function is appropriate for even-$`N`$ systems, $`N`$ being the number of particles in a system. The extension to odd-$`N`$ systems is given at the end of this subsection. Important quantities in the variational solution are the chemical potential $`\lambda `$ and the pairing gap $`\mathrm{\Delta }=G_j\mathrm{\Omega }_ju_jv_j`$, $`\mathrm{\Omega }_j=(2j+1)/2`$ being the pair degeneracy of the $`j`$-shell. When the pairing strength $`G`$ is small, the variational equations only have the trivial solution, $`v_j,u_j=0`$ or 1 and $`\mathrm{\Delta }=0`$. The ground state wave function is thus given by $`|HF=_{j,m}a_{jm}^{}|`$ and the ground state energy is obtained as
$$E_{HF}=HF|H|HF=\underset{j:occupied}{}\mathrm{\Omega }_j(2ϵ_jG).$$
(6)
When the strength of the pairing interaction $`G`$ is larger than some critical value $`G_{crit}`$, the variational equations have a non-trivial solution $`\mathrm{\Delta }0`$ and the ground state energy is given by
$$E_{BCS}=BCS|H|BCS=2\underset{j}{}\mathrm{\Omega }_jϵ_jv_j^2\frac{\mathrm{\Delta }^2}{G}G\underset{j}{}\mathrm{\Omega }_jv_j^4.$$
(7)
For odd $`N`$ systems, one of the particles in a system does not form a pair and blocks a level. When the $`k`$-th level is blocked, the HF energy, the BCS energy, and the pairing gap are modified to
$`E_{HF}`$ $`=`$ $`{\displaystyle \underset{j:occupied}{}}\stackrel{~}{\mathrm{\Omega }}_j(2ϵ_jG)+ϵ_k,`$ (8)
$`E_{BCS}`$ $`=`$ $`2{\displaystyle \underset{j}{}}\stackrel{~}{\mathrm{\Omega }}_jϵ_jv_j^2{\displaystyle \frac{\mathrm{\Delta }^2}{G}}G{\displaystyle \underset{j}{}}\stackrel{~}{\mathrm{\Omega }}_jv_j^4+ϵ_k,`$ (9)
$`\mathrm{\Delta }`$ $`=`$ $`G{\displaystyle \underset{j}{}}\stackrel{~}{\mathrm{\Omega }}_ju_jv_j,`$ (10)
respectively. Here, $`\stackrel{~}{\mathrm{\Omega }}_j=\mathrm{\Omega }_j`$ for $`jk`$, and $`\stackrel{~}{\mathrm{\Omega }}_k=\mathrm{\Omega }_k\delta _{N,odd}`$. The chemical potential $`\lambda `$ is determined so that
$$\underset{j}{}\stackrel{~}{\mathrm{\Omega }}_jv_j^2+1=N,$$
(11)
is satisfied. This formalism is referred to as the blocked-BCS theory.
### A Random phase approximation
We next introduce the random phase approximation to compute the correlation energy. Let us first consider the RPA in the superfluid phase (QRPA). We refer to Refs. for the formulation. The result is the well-known RPA matrix equation
$$\left(\begin{array}{cc}A& B\\ B& A\end{array}\right)\left(\begin{array}{c}X\\ Y\end{array}\right)=\mathrm{}\omega \left(\begin{array}{c}X\\ Y\end{array}\right),$$
(12)
where the matrices $`A`$ and $`B`$ are given by
$`A_{ij}`$ $`=`$ $`2E_i\delta _{i,j}G\sqrt{\stackrel{~}{\mathrm{\Omega }}_i}\sqrt{\stackrel{~}{\mathrm{\Omega }}_j}(u_i^2u_j^2+v_i^2v_j^2),`$ (13)
$`B_{ij}`$ $`=`$ $`G\sqrt{\stackrel{~}{\mathrm{\Omega }}_i}\sqrt{\stackrel{~}{\mathrm{\Omega }}_j}(u_i^2v_j^2+v_i^2u_j^2),`$ (14)
respectively. Here, $`\stackrel{~}{\mathrm{\Omega }}_j`$ is defined in the previous subsection, and $`E_j=\sqrt{(ϵ_j\lambda Gv_j^2)^2+\mathrm{\Delta }^2}`$ is the quasi-particle energy. The RPA excitation operator $`Q^{}`$ is given in terms of quasi-particle operators
$`\alpha _{jm}^{}`$ $`=`$ $`u_ja_{jm}^{}v_ja_{j\overline{m}},`$ (15)
$`\alpha _{j\overline{m}}^{}`$ $`=`$ $`u_ja_{j\overline{m}}^{}+v_ja_{jm},`$ (16)
as
$$Q^{}=\underset{j}{}\left(X_j\underset{m}{}\alpha _{jm}^{}\alpha _{j\overline{m}}^{}Y_j\underset{m}{}\alpha _{j\overline{m}}\alpha _{jm}\right)/\sqrt{\stackrel{~}{\mathrm{\Omega }}_j}.$$
(17)
The QRPA correlation energy is given by Eq. (1) with the $`A`$ matrix given by Eq. (13). Note that the second term on the right hand side of Eq. (1) is not $`_i2E_i`$ but $`Tr(A)`$. This point was overlooked in the integral approaches to the correlation energy in Refs. .
We next consider the RPA in a normal-fluid phase (pp-RPA), which describes pairing vibrations. The equation for the pp-RPA can be obtained from the QRPA equation by setting $`\mathrm{\Delta }=0`$ and $`v_h=u_p=1`$, where $`p`$ and $`h`$ denote particle and hole states, respectively. As for the chemical potential $`\lambda `$, there is no definite value for it in the normal fluid phase since it can be anywhere between the highest occupied and the lowest unoccupied levels. Notice, however, that the pairing vibration describes the ground state in the $`N\pm 2`$ systems and the role of the chemical potential is therefore just to shift the energies by an amount $`\pm 2\lambda `$. After removing these trivial energy shifts, one obtains
$`A_{pp^{}}`$ $`=`$ $`2ϵ_p\delta _{p,p^{}}G\sqrt{\stackrel{~}{\mathrm{\Omega }}_p}\sqrt{\stackrel{~}{\mathrm{\Omega }}_p^{}},`$ (18)
$`A_{hh^{}}`$ $`=`$ $`2(ϵ_hG)\delta _{h,h^{}}G\sqrt{\stackrel{~}{\mathrm{\Omega }}_h}\sqrt{\stackrel{~}{\mathrm{\Omega }}_h^{}},`$ (19)
$`B_{ph}`$ $`=`$ $`G\sqrt{\stackrel{~}{\mathrm{\Omega }}_p}\sqrt{\stackrel{~}{\mathrm{\Omega }}_h},`$ (20)
$`A_{ph}`$ $`=`$ $`B_{pp^{}}=B_{hh^{}}=0.`$ (21)
This $`2(N_p+N_h)\times 2(N_p+N_h)`$ dimensional matrix equation, where $`N_p`$ is the number of unoccupied shells and $`N_h`$ is that of occupied shells, can be decoupled into two $`(N_p+N_h)\times (N_p+N_h)`$ matrix equations as
$$\left(\begin{array}{cc}A& B\\ B& C\end{array}\right)\left(\begin{array}{c}X^a\\ Y^a\end{array}\right)=\mathrm{}\omega _a\left(\begin{array}{c}X^a\\ Y^a\end{array}\right),$$
(22)
and
$$\left(\begin{array}{cc}C& B\\ B& A\end{array}\right)\left(\begin{array}{c}X^r\\ Y^r\end{array}\right)=\mathrm{}\omega _r\left(\begin{array}{c}X^r\\ Y^r\end{array}\right).$$
(23)
Here, the matrix $`C`$ is defined as $`C_{hh^{}}=A_{hh^{}},C_{ph}=0`$, and RPA amplitudes are given by $`X_p^a=X_p,Y_h^a=Y_h,X_h^r=X_h`$, and $`Y_p^r=Y_p`$. The first equation (22) describes the ground state of the $`N+2`$ system and is referred to as the addition mode, while the second equation (23) describes the $`N2`$ system and is called the removal mode. Noticing that Eq. (1) is a half of a sum of the difference between the RPA and the Tamm-Dancoff approximation (TDA) frequencies for each mode of excitation , we find the correlation energy for the addition mode to be
$$E_{corr}^a=\frac{1}{2}\left(\underset{i}{}\mathrm{}\omega _{ai}Tr(A)\right),$$
(24)
while that for the removal mode is
$$E_{corr}^r=\frac{1}{2}\left(\underset{i}{}\mathrm{}\omega _{ri}Tr(C)\right).$$
(25)
The total correlation energy is the sum of these, $`E_{corr}=E_{corr}^a+E_{corr}^r`$.
Typically, one finds that an RPA mode goes to zero frequency at the mean-field phase transition. This is not the case for Eqs. (22) and (23), thus giving a discontinuity in the frequencies at the transition. However, as we have mentioned above, this is an artifact of an awkward choice of the chemical potential $`\lambda `$, and it will be the case that the sum of an addition and removal mode goes to zero.
### B Lipkin-Nogami method
An alternative way to restore the broken gauge symmetry of the BCS approximation is to carry out the number projection of the BCS wave function. The Lipkin-Nogami method provides an approximate way for number projection. It was first invented by Lipkin , and was developed by Nogami and his collaborators . Because of its relative simplicity, it has been widely applied .
In the Lipkin-Nogami method, the expectation value $`H\lambda \widehat{N}\lambda _2\widehat{N}^2_{BCS}`$ is varied. The resultant equations to be solved are given by
$`{\displaystyle \frac{2}{G}}={\displaystyle \underset{j}{}}{\displaystyle \frac{\stackrel{~}{\mathrm{\Omega }}_j}{\sqrt{\stackrel{~}{ϵ}_j+\mathrm{\Delta }^2}}},`$ (26)
$`2{\displaystyle \underset{j}{}}\stackrel{~}{\mathrm{\Omega }}_jv_j^2+\delta _{N,odd}=N,`$ (27)
$`\lambda _2={\displaystyle \frac{G}{4}}\left\{{\displaystyle \frac{\left(_i\stackrel{~}{\mathrm{\Omega }}_iu_i^3v_i\right)\left(_i\stackrel{~}{\mathrm{\Omega }}_iu_iv_i^3\right)_i\stackrel{~}{\mathrm{\Omega }}_iu_i^4v_i^4}{\left(_i\stackrel{~}{\mathrm{\Omega }}_iu_i^2v_i^2\right)^2_i\stackrel{~}{\mathrm{\Omega }}_iu_i^4v_i^4}}\right\},`$ (28)
where $`\stackrel{~}{ϵ}_i,v_i`$, and $`u_i`$ are defined as
$`\stackrel{~}{ϵ}_i`$ $`=`$ $`ϵ_i+(4\lambda _2G)v_i^2\lambda ,`$ (29)
$`v_i^2`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(1{\displaystyle \frac{\stackrel{~}{ϵ}_i}{\sqrt{\stackrel{~}{ϵ}_i^2+\mathrm{\Delta }^2}}}\right),`$ (30)
$`u_i^2`$ $`=`$ $`1v_i^2,`$ (31)
respectively. We use the blocked-Lipkin-Nogami prescription for odd $`N`$ systems. One of the characteristic features of the Lipkin-Nogami method is that the pairing gap $`\mathrm{\Delta }`$ has a finite value even in the weak $`G`$ limit where $`\mathrm{\Delta }`$ is zero in the BCS approximation. The ground state energy in the Lipkin-Nogami method is given by
$$E_{LN}=2\underset{j}{}\stackrel{~}{\mathrm{\Omega }}_jϵ_jv_j^2\frac{\mathrm{\Delta }^2}{G}G\underset{j}{}\stackrel{~}{\mathrm{\Omega }}_jv_j^44\lambda _2\underset{j}{}\stackrel{~}{\mathrm{\Omega }}_ju_j^2v_j^2+ϵ_k\delta _{N,odd}.$$
(32)
## III Two-level pairing model
We now apply the above equations to a schematic two-level model. This model was first introduced in Ref., and has been used in the literature to test several approximations . We label the lower and the higher levels 1 and 2, taking $`ϵ_1=ϵ/2,ϵ_2=ϵ/2`$. The Hamiltonian can be numerically diagonalised using the quasi-spin formalism . The basis states are denoted by $`|S_1S_{10};S_2S_{20}`$, where $`S_i`$ and $`S_{i0}`$ are defined by $`S_i=(\mathrm{\Omega }_i\nu _i)/2`$ and $`S_{i0}=(N_i\mathrm{\Omega }_i)/2`$, respectively. The latter takes a value of $`S_i,S_i+1,\mathrm{},S_i`$. $`\nu _i`$ is the seniority quantum number. For the ground state of even-$`N`$ systems, the total seniority $`\nu _1+\nu _2`$ is 0, while it is 1 for odd-$`N`$ systems. The matrix elements of the Hamiltonian read
$`S_1^{}S_{10}^{};S_2^{}S_{20}^{}|H|S_1S_{10};S_2S_{20}=\delta _{S_1,S_1^{}}\delta _{S_2,S_2^{}}\{ϵ(S_{20}S_{10}{\displaystyle \frac{\mathrm{\Omega }_1\mathrm{\Omega }_2}{2}})\delta _{S_{10},S_{10}^{}}\delta _{S_{20},S_{20}^{}}`$ (33)
$`G\left(S_1(S_1+1)S_{10}(S_{10}+1)+S_2(S_2+1)S_{20}(S_{20}+1)\right)\delta _{S_{10},S_{10}^{}}\delta _{S_{20},S_{20}^{}}`$ (34)
$`G\sqrt{(S_1(S_1+1)S_{10}(S_{10}1)}\sqrt{S_2(S_2+1)S_{20}(S_{20}+1)}\delta _{S_{10},S_{10}^{}+1}\delta _{S_{20},S_{20}^{}1}`$ (35)
$`G\sqrt{(S_1(S_1+1)S_{10}(S_{10}+1)}\sqrt{S_2(S_2+1)S_{20}(S_{20}1)}\delta _{S_{10},S_{10}^{}1}\delta _{S_{20},S_{20}^{}+1}\}.`$ (36)
We first consider a symmetric two-level problem i.e., $`\mathrm{\Omega }_1=\mathrm{\Omega }_2=\mathrm{\Omega }`$, and assume the number of fermionic particle is $`N=2\mathrm{\Omega }`$. Equations for several approximations can be solved analytically for such a system. The gap equation in the BCS theory leads to a pairing gap of
$$\mathrm{\Delta }=\sqrt{G^2\mathrm{\Omega }^2\frac{\stackrel{~}{ϵ}^2}{4}},$$
(37)
together with
$`v_1^2`$ $`=`$ $`u_2^2={\displaystyle \frac{1}{2}}\left(1+{\displaystyle \frac{\stackrel{~}{ϵ}}{2G\mathrm{\Omega }}}\right)v^2,`$ (38)
$`u_1^2`$ $`=`$ $`v_2^2={\displaystyle \frac{1}{2}}\left(1{\displaystyle \frac{\stackrel{~}{ϵ}}{2G\mathrm{\Omega }}}\right)u^2.`$ (39)
$`\lambda `$ $`=`$ $`{\displaystyle \frac{\stackrel{~}{ϵ}ϵ}{2}}Gv^2,`$ (40)
where $`\stackrel{~}{ϵ}`$ is defined as $`\stackrel{~}{ϵ}=2\mathrm{\Omega }ϵ/(2\mathrm{\Omega }1)`$.
From Eq. (37), one can find the critical strength of the phase transition to be $`G_{crit}=ϵ/(2\mathrm{\Omega }1)`$. For $`G`$ larger than $`G_{crit}`$, the system is in the superfluid phase, and the matrices $`A`$ and $`B`$ for the QRPA equation are given by (see Eqs. (13) and (14))
$`A_{11}`$ $`=`$ $`A_{22}=G\mathrm{\Omega }{\displaystyle \frac{\mathrm{\Delta }^2}{2G\mathrm{\Omega }^2}}+{\displaystyle \frac{\mathrm{\Delta }^2}{2G\mathrm{\Omega }}},`$ (41)
$`A_{12}`$ $`=`$ $`A_{21}={\displaystyle \frac{\mathrm{\Delta }^2}{2G\mathrm{\Omega }}},`$ (42)
$`B_{11}`$ $`=`$ $`B_{22}={\displaystyle \frac{\mathrm{\Delta }^2}{2G\mathrm{\Omega }^2}}+{\displaystyle \frac{\mathrm{\Delta }^2}{2G\mathrm{\Omega }}},`$ (43)
$`B_{12}`$ $`=`$ $`B_{21}=G\mathrm{\Omega }{\displaystyle \frac{\mathrm{\Delta }^2}{2G\mathrm{\Omega }}},`$ (44)
The solutions of the QRPA equation are $`\mathrm{}\omega =0`$ and $`\sqrt{4\mathrm{\Delta }^22\mathrm{\Delta }^2/\mathrm{\Omega }}`$, from which the correlation energy is computed as
$$E_{corr}=\sqrt{\mathrm{\Delta }^2\mathrm{\Delta }^2/2\mathrm{\Omega }}G\mathrm{\Omega }+\frac{\mathrm{\Delta }^2}{2G\mathrm{\Omega }^2}\frac{\mathrm{\Delta }^2}{2G\mathrm{\Omega }}.$$
(45)
When the strength of the pairing interaction $`G`$ is smaller than $`G_{crit}`$, the system is in the normal fluid phase. The mean field energy is then evaluated as $`E_{HF}=ϵ\mathrm{\Omega }G\mathrm{\Omega }`$. For the present two-level model, the $`A,B`$, and $`C`$ matrices for the pp-RPA are just numbers and are given by $`A=ϵG\mathrm{\Omega }`$, $`B=G\mathrm{\Omega }`$, and $`C=ϵG\mathrm{\Omega }+2G`$, respectively. The frequencies for the addition and the removal modes are then found to be
$`\mathrm{}\omega _a`$ $`=`$ $`G+\sqrt{ϵ+G}\sqrt{ϵ+G2G\mathrm{\Omega }},`$ (46)
$`\mathrm{}\omega _r`$ $`=`$ $`G+\sqrt{ϵ+G}\sqrt{ϵ+G2G\mathrm{\Omega }},`$ (47)
respectively. The total correlation energy is thus obtained as
$$E_{corr}=\sqrt{ϵ+G}\sqrt{ϵ+G2G\mathrm{\Omega }}(ϵG\mathrm{\Omega }+G).$$
(48)
Figure 1 shows the RPA frequencies for each modes of excitation as a function of $`G`$, for $`\mathrm{\Omega }`$ = 8. As we noted in the previous section, we see that because of the chemical potential the pp-RPA frequencies do not match with the QRPA frequencies at the critical point of phase transition from a normal-fluid to a superfluid phases. Figure 2 compares the ground state energy obtained by several methods. The solid line is the exact solution obtained by numerically diagonalizing the Hamiltonian. The dashed line is the ground state energy in the mean field BCS approximation. It considerably deviates from the exact solution through the entire range of $`G`$ shown in the figure. The dot-dashed line takes into account the RPA correlation energy in addition to the mean field energy. It reproduces very well the exact solutions, except in the vicinity of the critical point of the phase transition. Around the critical point, one would need to compute the correlation energy with some care, using e.g., the self-consistent RPA discussed in Ref. which removes the cusp behaviour of the RPA frequency around the critical strength.
It is interesting to compare the present mean field plus RPA approach with the Lipkin-Nogami method. The equations for the Lipkin-Nogami method for the two-level model were solved in Ref. . The results are given by
$`\mathrm{\Delta }^2`$ $`=`$ $`(G\mathrm{\Omega })^2(1\stackrel{~}{\kappa }^2),`$ (49)
$`v_1^2`$ $`=`$ $`u_2^2=(1+\stackrel{~}{\kappa })/2,`$ (50)
$`u_1^2`$ $`=`$ $`v_2^2=(1\stackrel{~}{\kappa })/2,`$ (51)
$`4\lambda _2G`$ $`=`$ $`{\displaystyle \frac{2G\mathrm{\Omega }\stackrel{~}{\kappa }^2}{(2\mathrm{\Omega }1)(1\stackrel{~}{\kappa }^2)}},`$ (52)
where $`\stackrel{~}{\kappa }`$ is the physical solution, which satisfies $`0\stackrel{~}{\kappa }1`$, of an equation
$$2(1\mathrm{\Omega })\stackrel{~}{\kappa }^3+(2\mathrm{\Omega }1)\kappa \stackrel{~}{\kappa }^2+(2\mathrm{\Omega }1)\stackrel{~}{\kappa }(2\mathrm{\Omega }1)\kappa =0,$$
(53)
and $`\kappa `$ is defined as $`ϵ/2G\mathrm{\Omega }`$. The ground state energy in the Lipkin-Nogami method given by Eq. (32) is denoted by the thin solid line in Fig. 2. Although it reproduces the exact results for large values of $`G`$, it deviates significantly from them for small values. This behaviour is consistent with the numerical observation in Ref. as well as the result of Ref. where it was found that the Lipkin-Nogami method is only correct in the limit of a strong pairing force. In marked contrast, the ground state energy in the BCS plus QRPA coincides with the exact solution at the leading order of an expansion in 1/$`\mathrm{\Omega }`$ for any strength of the pairing interaction as long as it is larger than the critical value .
From Fig. 2, it is unclear whether the Lipkin-Nogami method or RPA are more accurate for purposes of a global theory of binding. So we now consider a more realistic situation, varying the particle number $`N`$ rather than the interaction strength $`G`$. We consider the paring energy in oxygen isotopes, taking the neutron 1p and 2s-1d shells as the lower and higher levels of the two-level model. The pair degeneracy $`\mathrm{\Omega }`$ thus reads $`\mathrm{\Omega }_1=3`$ and $`\mathrm{\Omega }_2=6`$, and the number of particle in a system is given by $`N=A82`$ for the <sup>A</sup>O nucleus. We assume that the energy difference between the two levels $`ϵ`$ is given by $`ϵ=41A^{1/3}`$ and the pairing strength $`G=23/A`$. The upper panel of Fig. 3 shows the ground state energy as a function of $`A`$. In order to match with the experimental data for the <sup>16</sup>O nucleus, we have added a constant $`72.8`$ MeV to the Hamiltonian for all the isotopes. The exact solutions are denoted by the filled circles. The deviation from the BCS approximation (the dashed line) is around 2 MeV for even $`A`$ systems and it is around 1.2 MeV for odd $`A`$ systems. This value varies within about 0.5 MeV along the isotopes and shows relatively strong $`A`$ dependence. One can notice that the RPA approach (the dot-dashed line) reproduces quite well the exact solutions. On the contrary, the Lipkin-Nogami approach (the thin solid line) is much less satisfactory and shows a different $`A`$ dependence from the exact results. The pairing gaps $`\mathrm{\Delta }`$ in the BCS approximation and in the Lipkin-Nogami method are shown separately in the lower panel of Fig. 3. For the Lipkin-Nogami method, we show $`\mathrm{\Delta }+\lambda _2`$, which is to be compared with experimental data . The closed shell nucleus <sup>16</sup>O and its neighbour nuclei <sup>15,17</sup>O have a zero pairing gap in the BCS approximation, and the Lipkin-Nogami method does not work well for these nuclei, as can be casted in Fig. 2. On the other hand, the RPA approach reproduces the correct $`A`$ dependence of the binding energy. Evidently, the RPA formula provides a better method to compute correlation energies than the Lipkin-Nogami method, especially for shell closures.
## IV Summary
Returning to our initial motivation, we seek a computationally tractable way to include pairing effects in a global model of nuclear binding energies, going beyond the BCS theory. Two attractive possibilities are the Lipkin-Nogami method and the RPA, in particular if the pairing interaction has a separable form. In this paper, we used a solvable two-level pairing Hamiltonian to show that the RPA formula for the correlation energy reproduces well the exact solutions both in a normal-fluid and a superfluid phases. On the contrary, the Lipkin-Nogami method is considerably less accurate for weak pairing, and therefore is not suitable in transition regions. As a consequence, the Lipkin-Nogami method fails to reproduce the correct mass number dependence of the binding energy around a shell closure. The correlation energy associated with the number fluctuation for the neutron mode in O isotopes was shown to be order of 2 MeV and has a relatively strong mass dependence. Although the correlation energy is small compared with the absolute value of the ground state energy, this suggests that including the correlation energy in the RPA provides a promising way to develop a better microscopic systematic theory for nuclear binding energies.
## Acknowledgments
The authors thank P.-G. Reinhard and W. Nazarewicz for useful discussions. This work was supported by the U.S. Department of Energy under Grant DOE-ER-40561.
## A Richardson solution
For a system where the number of pairs is $`N_{pair}`$, i.e., a 2$`N_{pair}`$-fermion system, the ground state energy of the Hamiltonian (2) is given by
$$E_{gs}=\underset{\lambda =1}{\overset{N_{pair}}{}}z_\lambda ,$$
(A1)
where $`z_\lambda `$ are $`N_{pair}`$ solutions of $`N_{pair}`$ coupled equations given by
$$\underset{j}{}\frac{\mathrm{\Omega }_j}{2ϵ_jz_\lambda }\stackrel{}{}_\lambda ^{}\frac{2}{z_\lambda ^{}z_\lambda }=\frac{1}{G}.$$
(A2)
The prime in the summation of $`\lambda ^{}`$ means to take only those $`\lambda ^{}`$ which are different from $`\lambda `$, and $`\mathrm{\Omega }_j=(2j+1)/2`$ is the pair degeneracy of the $`j`$-shell. In general, $`z_\lambda `$ are complex.
For a two pair system ($`N_{pair}`$=2) in a single $`j`$-shell, setting $`ϵ_j=0`$ and $`\mathrm{\Omega }_j=\mathrm{\Omega }`$, the Richardson equation reads
$`{\displaystyle \frac{\mathrm{\Omega }}{z_1}}{\displaystyle \frac{2}{z_2z_1}}`$ $`=`$ $`{\displaystyle \frac{1}{G}}`$ (A3)
$`{\displaystyle \frac{\mathrm{\Omega }}{z_2}}{\displaystyle \frac{2}{z_1z_2}}`$ $`=`$ $`{\displaystyle \frac{1}{G}}.`$ (A4)
The solutions of these equations are found to be
$$z=G(\mathrm{\Omega }1)\pm i\sqrt{G^2(\mathrm{\Omega }1)},$$
(A5)
from which the ground state energy reads $`E_{gs}=2G(\mathrm{\Omega }1)`$. This result coincides with the solution obtained using the seniority scheme,
$$E(N_{pair})=G\left\{\frac{\mathrm{\Omega }}{2}\left(\frac{\mathrm{\Omega }}{2}+1\right)\frac{1}{4}(2N_{pair}\mathrm{\Omega })^2+\frac{1}{2}(2N_{pair}\mathrm{\Omega })\right\}.$$
(A6)
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# Will the Population of Humanity in the Future be Stabilized?
## I Introduction
The problem of mankind population growth is the one of the global problems concerning the development of the mankind and its future. Will the demographic explosion existing now at the mankind population of the whole world stop as already it has stopped in the most developed countries? Will the population of the Earth will be stabilized as it follows from S.Kapitza’s theory (see )near 12 billions, or it will continue its growth at slower rate? What are the driving parameters that govern the development of mankind (such as presents in non-linear open system (see)? Are these parameters genetically predetermined or can they be changed and controlled by means of human activity? Stabilization of the mankind population of earth in the future is a sad prospect for mankind, because the absence of the numerical increasing of any biological population almost always leads, early or late, to cessation of any development (the examples are many species of insects, e.g. termites, frozen in development for millions years). Hence, the appearance of more active biological species becomes in that case quite probable. These active form of life will dominate the mankind and may force it out from its present ecological niche. The aim of the present paper is to introduce a new parameter in the phenomenological theories of humankind growth describing the development of mankind population if regard humankind as a large non-linear multifractal system. Namely, the fractal dimension of the whole mankind number in arbitrary moment of time - the local fractional dimension $`d(t)`$ (the fractional dimension for whole population of mankind). The introducing of this parameter allows to receive, as a special cases, the results of theories (see ),-), and several new scenarios of the mankind’s future as well. Alongside with probable increasing of the mankind population, the scenario of ruin of the present civilization - diminution of its number down to zero - is shown to be possible. We note, that correct analysis of dynamics of multifractal sets (see) requires the introduction of the mathematical concept of fractional derivatives (see -), which allow to take into account partly the memory of system about the past (in this case it is the memory, that includes genetic memory of mankind about its past development).
## II Fractional Dimension Quantity of Mankind on Time axis
We assume that the dynamics of the human population can be described within the framework of fractal geometry concepts and mathematical formalism of fractional differential equations (see -). For this purpose let us consider the set of all people as a multifractal set $`N(t)`$ $`(0t\mathrm{})`$ consisting of $`N(t)`$ of elements at the given time $`t`$. Following the Kapitza theory (see ), assume for a certain interval of time the existence of a self-similarity of this set characterized by its fractal dimension and introduce a local fractional dimension (LFD) $`d(t)`$ which describes the fractal properties of the set at the time $`t`$. This local fractional dimension is determined by those variables and their functions (these variables will be defined later) , that are treated as the driving (control) parameters of the human community development. Among such parameters there can be parameters of genetic origin (for example, probably, density of the population in cities per unit urban area, etc.) and ”external” parameters (e.g., a possibility of supplying the mankind population of earth by necessary amounts of food, water and energy or quick climate changes or pollution of the environment and the atmosphere, etc.). We shall characterize an alteration of the mankind population $`N(t)`$ over a small time interval by the generalized fractional Riemann - Liouville derivative (which coincides with the usual Riemann - Liouville derivative if $`d`$=const.)
$$D_{+,t}^{1+\nu (t)}N(t)=\frac{^\alpha }{t^\alpha }\underset{0}{\overset{t}{}}\frac{N(t^{})dt^{}}{\mathrm{\Gamma }(\alpha d(t^{}))(tt^{})^{d(t^{})\alpha +1}},$$
(1)
$`d(t)=1+\nu (t)>0`$ (2)
$$\nu (t)\nu (X_1(t),X_2(t)\mathrm{}X_i(t)),i=1,2,\mathrm{},\alpha =\{d\}+1$$
(3)
In (1) $`\nu (t)`$ is the fractional quantity and defines the differences between the derivatives of integer order and fractional derivatives (1) thus being the driving parameter for the growth of mankind as a whole, $`\{d\}`$ is equals to the integer part of $`d(t)>0`$ $`(\alpha 1d(t)<\alpha )`$ , $`\alpha =0`$ for $`d<0`$ , and the set $`X_i`$ are the control parameters determining all external and internal influences on the mankind population growth .
The explicit information about the function $`\nu `$ and, hence, about LFD can be obtained only after a careful investigations and processing the statistical data of different events, circumstances and situations impact on the development of a human population. The fractional derivative using, as defined in (1), instead of the integer first derivative allows to introduce and take into account an obvious thing - a certain kind of mankind’s memory of the past and memory about the development rates over the past years (integration with a weight function over all times till $`t`$ beginning from a fixed moment). It gives a way of considering of different parameters $`X_i`$ that influence at the mankind’s development by means of LFD’s dependence upon them. The present theory, as well as (), is a phenomenological theory and the exact definition of function $`\nu (t)`$ form is beyond its scope. We note that it has sense to consider $`\nu (t)1`$, apparently, only for times greater than $`T_1^{}`$ because before a certain time, introduced in (, -) (the time of demographic transition $`T=T_1^{}`$), the theories mentioned describe the empirical data about the number and progress of the mankind population quite well. Moreover, we shall restrict ourselves in this paper to analyzing growth of the mankind population for three special cases of $`\nu (t):\nu (t)=0,\nu (t)=1`$ and $`\nu (t)=1`$.
## III Foerster,Hoerner, Kapitza Theories
It is known, that growth of the population of earth from ancient times untill now is well approximated by an empirical relation suggested by Foerster Von H. (see ) and improved by H. von Xoerner (see )
$$N(t)=\frac{C}{T_1T^{}},$$
(4)
$$\frac{d(N(t))}{dt}=\frac{C}{(T_1T)^2}$$
(5)
$`C=2^.10^{11},T_1=2025`$
The reasonable generalization (4),(5) for future time (suitable for $`T>T_1^{}`$ that not gives as result an infinite value $`N(t)`$ for values $`T=T_1`$) was given in the phenomenological theory of the population S.Kapitza (see ) with the help of introducing of the mean people’s lifetime $`\tau `$ ($`\tau `$ =42 of years). In this theory the relation for $`N(T)`$ (with C’=1,86$`{}_{}{}^{.}10_{}^{11},T_1^{}=2007`$) takes place
$$\frac{N(t)}{dt}=\frac{C^{}}{(T_1^{}T)^2+\tau ^2}$$
(6)
From (6) the restriction of the mankind population of earth by quantity $`14^.10^9`$ follows. Unfortunately, the theory of the population S.Kapitza does not take into account neither exterior nor interior control parameters (even in simple form) basing only at the self-similarity of growth of a population of the people.
## IV Generalization of Foerster, Hoerner and Kapitza Theories for Multifractal Set of a Quantity of Mankind $`N(t)`$
Let us assume the hypothesis about a fractal nature of set $`N(t)`$ (maintaining assumption about selfsimilar the sets $`N(t)`$). In that case derivative on time $`\frac{}{t}`$ in equation (6) should be substitute for fractional derivative $`D_{+,t}^{1+\nu (t)}`$. This operation take into account the memory of mankind about the past development. The right part (6) must be changed too for including in (6) an influence of the fractal dimension. So, instead of (6), obtain an equation
$$D_{+,t}^{1+\nu (t)}N(t)=\frac{C^{}}{|T_1^{}T|^{2+\nu (t)}+\frac{2+\nu (t)}{2}\tau ^{2+\nu (t)}}$$
(7)
The equation (7) can be considered as the basic equation of the phenomenological theory of development mankind’s population offered in this paper. The selection of different functions of fractional corrections for $`\nu (t)`$ allows to estimate character of changing $`N(t)`$ as functions of time.
## V Forecasts of the Future Development of Mankind Population for Special Cases of the Fractional Dimension Choice
Some simple special cases forecasting of growth of the mankind’s population of earth are considered below at the basis of the equation (7). The meanings a fractal dimension, for simplicity , chosen as integer.
a. Case $`\nu (t)=0`$
At $`\nu (t)=0`$ the fractional derivative $`D`$ coincides with $`\frac{}{t}`$ derivative and the equation (7) coincides with equation (6) (hence, (7) includes S.Kapitza theory (see ) as a special case). It corresponds, probably, a compensation the positive and the negative control parameters $`X_i`$ (including exterior and ”interior” parameters), driving population of mankind.
b. Case $`\nu (t)1`$
The selection $`\nu (t)1`$ corresponds to a dominance of the negative tendencies in the future development of mankind and it is stipulated, for example, occurrence of irreversible changes in molecules DNA (owing to epidemic AIDS, etc.), irreversible cosmic cataclysmic (for example, drop on earth of meteorites of huge mass), impossibility for mankind to cope with negative factors of biogeozinos results in by effects of development of a technical civilization. In this case the equation (7) reduce to equation
$$N(t)=\frac{2C^{}}{2|T_1^{}T|+\tau }$$
(8)
From (8) the presence of the maximum of number of mankind follows at $`T=T^{}`$ (i.e. in 2007) and it is equal $`8,86^.10^9`$ (if $`\tau `$ is not changed). After transition through a maximum $`N(t)`$ the number of mankind decreases (if the scenario will not vary) and in 2107 year $`N(t)`$ will be equal $`1,54^.10^9`$. By the year 3007 the number of mankind will decrease down to $`182^.10^6`$, i.e. the mankind population of all earth will be equal to number of the people is occupying a dozen of large modern cities. The complete degenerating of mankind, as a result of decreasing of mankind population in the considered scenario and subsequent leaving mankind from biological arena may be expected through millions years ( not large time from biological point of view). By then population the mankind of earth will decrease up to several thousand.
c. Case $`\nu (t)=2.`$
The equation (7) for this case is transformed into an integral equation at maintenance of the general tendency to accumulation of the negative factors resulting in to negative value $`d(t)`$ . At negative values $`d(t)`$ the integral equation gives the prompt diminution of mankind population and extinction of existence of mankind as a species is follows. So, at $`d(t)1`$ equations (7) transforms to a
$$\underset{T^{}}{\overset{T}{}}N(t)𝑑t=C^{}$$
(9)
supposing the absence the mankind at the earth (in that case equation (9) has no solution for $`N(t)0`$ thou for $`\nu >2`$ solution of equation (7) is exist (so $`N(t)0`$ for $`\nu 2`$). The time interval (necessary for disappearance mankind) is determined in this case by time for which $`d(t)`$ will transfer from value equal to unity (status of mankind now) to value equal to negative unity. This time interval can be very short: from several years (cosmic cataclysmic) to about several centuries (virus pandemia with lethal change of a heredity, increasing of the mean temperature of earth on some degrees owing to throw out of carbonic gas etc.). It is necessary to pay attention on possibility of change of the negative scenario of mankind development stipulated by appearance and dominance of the positive control factors $`X_i`$ (including those factors that due to conscious activity of mankind). In that case the inevitability of diminution of mankind population and destruction of mankind are not inevitable.
d. case $`\nu =1.`$
We shall consider the optimistic scenario of change of mankind population. It relevant to a dominance of positive driving parameters: $`d(t)>1.`$ So, e.g. for $`d(t)=2`$ $`(\nu =1)`$ from (7) the equation follows
$$\frac{^2}{t^2}N(t)=\frac{C^{}}{|T^{}T|^3+1,5\tau ^3}$$
(10)
The precise solution of an initial value problem of this equation is unwieldy, so we shall note , that at $`T>>\tau `$ the quantity $`N(t)`$ will increase faster then first degree of time ($`N(t)(TT_1)`$). The mankind population is increase and it characterizes in this case by following form (if there will not be includes appearance of a factors of conscious mankind activity which change the scenario and restricts unlimited growth of population)
$$N(t)|^t\mathrm{}\frac{C^{}(TT_1^{})}{2,29\tau ^2}\mathrm{ln}\frac{(TT_1^{})}{\tau }$$
(11)
So, at 3000 years, at the rate of increasing of the population defined by (11) (with the allowances that the corrections to $`N(t)`$ of value $`ln[\frac{(TT_1^{})}{\tau }]`$ are dropped) population of mankind will increase up to 150 billions. That is improbable large, though, but it is may be not unreasonable because of future technical possibilities of mankind (probably, this number is the upper number for existence mankind population occupying the earth). For fractional values LFD $`d(t)`$, increasing or decreasing of a population of the mankind will be characterized by intermediate dependencies between the received for the whole values $`d(t)`$. In case of a periodic dependencies $`d(t)`$ from time the population of the world will change periodically depending on a concrete choice of $`d`$ and $`N(t)`$ and will not be a monotonous function of time.
The examples are considered allow to determine the interval of change of fractional dimension $`d(t)`$ in reason boundary for number of quantity mankind in future: $`1<d<2`$. The boundary values $`(d=1,d=2)`$ are result in or to ruin of mankind, or to so large mankind population that Earth may not endure. The last case must result in to change of scenario and it consist in the change in correlation between positive and negative components of control parameters $`X_i`$ towards increasing of influence of negative parameters.
Let the basic equation (7) is replaced by generalization of the S.Kapitza equation (6). Basic equation in that case reads
$$\frac{^{1+\nu (t)}}{t^{1+\nu (t)}}n(t)=K\mathrm{sin}^2\frac{n(t)}{K}=\frac{1}{K},$$
(12)
$`n={\displaystyle \frac{N(t)}{K}},t={\displaystyle \frac{TT_1^{}}{\tau }},K=\sqrt{C^{}\tau ^1}`$ (13)
It gives for $`N(t)`$ the qualitative effects analogous to results obtained from the equation (7). So, in this connection, a selection for describing the future increasing the population of mankind by equation (7), as the basis, or equation (12), containing, as well as the equation (7), driving parameters $`X_i`$ in fractional dimension $`d(t)`$, is not simple. As one of advantages of use of the equations (7) or (12) for describing demographic problems (with some of them the mankind already has confronted now) we shall stress an opportunity of insert and account in the theory many factors (such as , incurable illnesses, natural cataclysmic etc.) defining future of mankind as a result of influences the control parameters $`X_i`$ (included in the dimension $`d(t)`$ of fractal set for number’s distribution of the people in the time axis).
## VI Conclusions
1.The main purpose of this paper was to analyze possibility of introducing mankind’s population driving parameters $`X_i`$ in the phenomenological theories of the mankind population of the earth (, , ) by method of attributing to set of the people $`N(t)`$ the fractional dimension $`d(t)`$, depending from these parameters. At a choice $`d(t)=1`$ the numbered theories are a special case of this phenomenological theory. The consideration of examples with integer meanings of $`d(t)`$ is caused only by their mathematical simplicity and gives a reasonable meaning of fractal dimensions ($`1<d<2`$) for describing the time dependence of population of mankind.
2. In case when the interpretation the fractal properties of set of the homosapience given in this paper corresponds to a reality, ( more real cases correspond to fractional meanings $`d(t)`$) the future of mankind is not so mournful as in the case of the S.Kapitza theory (see ) and the exposition of change of number of mankind within the framework of the phenomenological theories of the population can be reduced to a selection of control parameters and filling them by the concrete contents.
3. The change of number of mankind (described in the framework a phenomenological theories of the population) can be adjusted by such choice of control parameters (and filling the fractal dimension $`d(t)`$ by the concrete contents of dependencies of them) at which the population of mankind will grow so slowly, that overpopulation and the problems connected with it will do not arise in the foreseen future. Last will allow the theory be more realistic for predicting and menaging the future growth population of mankind as one of biological species occupying our world.
4. The growth of mankind population regulation (included in the parameters $`X_i`$) will allow to avoid degenerating of mankind and to keep as much as long time the main ecological niche at Earth occupied by mankind. Last will give time for more realistic forecasting of the future mankind as one of biological kinds occupying our world.
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# THE BEST UNBIASED ESTIMATOR FOR THE CMB ANGULAR BISPECTRUM
## References
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# Emergence of intense jets and Jupiter Great Red Spot as maximum entropy structures.
## 1 Introduction
Atmospheric and oceanic flows are often organized into narrow jets. They can zonally flow around the planet like the jet streams in the earth stratosphere, or the eastward jet at 24 in the northern hemisphere of Jupiter ( Maxworthy 1984). Jets can alternatively organize into rings, forming vortices, like the rings shed by the meandering of the Gulf-Stream in the western Atlantic Ocean. The flow field in Jupiter most famous feature, the Great Red Spot, is an oval-shaped jet, rotating in the anticyclonic direction and surrounding an interior area with a weak mean flow ( Dowling and Ingersoll 1989), see figure 1(a). Robust cyclonic vortices have a similar jet structure ( Hatzes et al 1981 ), see figure 1(b).
Such jets and vortices are in a turbulent surrounding, and the persistence of their strength and concentration in the presence of eddy mixing is intriguing. The explanation proposed in this paper is based on a statistical mechanical approach: the narrow jet or vortex appears as the most probable state of the flow after a turbulent mixing of potential vorticity, taking into account constraints due to the quantities conserved by the dynamics, especially energy. Such a statistical theory has been first proposed for the two-dimensional Euler equations by Kuz’min (1982), Robert (1990) , Robert and Sommeria (1991), Miller (1990). See Brandt et al (1999) for a recent review and discussion. This theory predicts an organization of two-dimensional turbulence into a steady flow (with fine scale, ’microscopic’ vorticity fluctuations). Complete vorticity mixing is prevented by the conservation of the energy, which can be expressed as a constraint in the accessible vorticity fields. A similar, but quantitatively different, organization had been previously obtained with statistical mechanics of singular point vortices with the mean field approximation, instead of continuous vorticity fields (Onsager 1949, Joyce and Montgomery 1973).
Extension to the quasi-geostrophic (QG) model has been discussed by Sommeria et al (1991), Michel & Robert (1994), Kazantsev Sommeria and Verron (1998). This model describes a shallow water system with a weak vorticity in comparison with the planetary vorticity (small Rossby number), such that the pressure is in geostrophic balance, and the corresponding free surface deformation is supposed small in comparison with the layer thickness. For Jupiter the free surface would be rather at the bottom of the active atmospheric layer, floating on a denser fluid, as discussed by Dowling and Ingersoll (1989), see Dowling (1995) for a review. The gradient of planetary vorticity is accounted by a beta-effect. An additional beta-effect, depending on the latitude coordinate $`y`$, is introduced to represent the influence on the active atmospheric layer of a steady zonal flow in the deep interior, as discussed by Dowling and Ingersoll (1989).
The free surface deformability, representing the strength of the density stratification, is controlled by the Rossby radius of deformation $`R^{}`$. The two-dimensional Euler equation is recovered in the limit of very strong stratification for which $`R^{}\mathrm{}`$. We shall consider in this paper the opposite limit of weak stratification for which $`R^{}`$ is much smaller than the scale of the system $`L`$. This limit is appropriate for large scale oceanic currents, as the radius of deformation is typically 10-100 km. For Jupiter, $`R^{}`$ is estimated to be in the range 500-2500 km, while the Great Red Spot extends over 20,000 km in longitude, and 10,000 km in latitude, so the limit $`R^{}/L0`$ seems relevant. We show that in this limit the statistical equilibrium is made of quiescent zones with well mixed uniform potential vorticity, bounded by jets with thickness of order $`R^{}`$. This provides therefore a general justification of jet persistence. Some of the ideas used have been already sketched in Sommeria et al (1991), but we here provide a systematic derivation and thorough analysis.
In principle, the Quasi Geostrophic approximation breaks down for scales much larger than the radius of deformation, so that the limit $`R^{}/L0`$ seems inconsistent with the QG approximation. However the relevant scale is the jet width, which remains of order $`R^{}`$, so that the Quasi Geostrophic approximation remains valid in this limit. This point has been discussed by Marcus (1993) for the Great Red Spot, which he supposes to be a uniform Potential Vorticity ( PV ) spot surrounded by a uniform Potential Vorticity background ( we here justify this structure as the result of Potential Vorticity mixing with constraints on the conserved quantities ). Analyzing wind data in the Great Red Spot, Dowling & Ingersoll (1989) concluded that the QG approximation is good up within typically 30% error, which is reasonable to a first approximation. Statistical mechanics of the more general shallow water system (to be published), predicts a similar jet structure. The present Quasi Geostrophic results therefore provide a good description as a first approximation.
We first consider the case without beta-effect in section 2. We furthermore assume periodic boundary conditions (along both coordinates) in this section to avoid consideration of boundary effects. Starting from some initial condition with patches of uniform PV, we find that these patches mix with uniform density (probability) in two sub-domains, with strong density gradient at the interface, corresponding to a free jet. The coexistence of the two sub-domains can be interpreted as an equilibrium between two thermodynamic phases. We find that the interface has a free energy per unit of length, and its minimization leads to a minimum length at equilibrium. This results in a constant radius of curvature, in analogy with surface tension effects in thermodynamics, leading to spherical bubbles or droplets. The range of the vortex interaction is of the order $`R^{}`$, therefore becoming very small in the limit of small radius of deformation, so the statistical equilibrium indeed behaves like in usual thermodynamics with short range molecular interactions. This contrasts with the case of Euler equation, with long range vortex interactions, analogous to gravitational effects (Chavanis Sommeria and Robert 1996, Chavanis 1998).
Figure 7 summarizes the calculated equilibrium states, depending on the total energy and a parameter $`B`$ representing an asymmetry between the initial Potential Vorticity patch areas, before the mixing process. We obtain straight jets for a weak asymmetry and circular jets for higher asymmetry. Such circular a jet reduces to an axisymmetric vortex, with radius of order $`R^{}`$, in the limit of low energy.
We discuss the influence of the beta-effect or the deep zonal flow in section 3. The channel geometry, representing a zonal band periodic in the longitude $`x`$ is appropriate for that study. With the usual beta-effect $`\beta y`$, linear in the transverse coordinate $`y`$, statistical equilibrium is, depending on the initial parameters, a zonal flow, or a meandering eastward jet, or a uniform velocity $`v_m=R^2\beta `$ whose induced free surface slope cancels the beta-effect ( uniformization of Potential Vorticity ) on which circular vortices can coexist.
For more general beta-effects, due to the deep zonal flow, we find that the jet curvature depends on latitude $`y`$. In particular a quadratic beta effect $`ay^2`$ leads to oval-shape jets, similar to the Great Red Spot. Using the determination of the sublayer flow from Voyager data by Dowling and Ingersoll (1989), we show in section 4, that such a quadratic effective beta-effect is indeed a realistic model for Jupiter atmosphere in the latitude range of the Great Red Spot and the White Ovals, the other major coherent vortices on Jupiter. Using these data on beta-effect, as well as the shear in the zonal flow at the latitude of the Great Red Spot, the jet width and its maximum velocity, we deduce all the parameters of our model.
## 2 The case with periodic boundary conditions
### 2.1 The dynamical system
We start from the barotropic Quasi Geostrophic (QG) equation :
$$\frac{q}{t}+𝐯q=0$$
(1)
$$q=\mathrm{\Delta }\psi +\frac{\psi }{R^2}h(y)$$
(2)
$$𝐯=\widehat{𝐳}\psi $$
(3)
where $`q`$ is the potential vorticity ( PV ), advected by the non-divergent velocity $`𝐯`$, $`\psi `$ is the stream function, $`R`$ is the internal Rossby deformation radius between the layer of fluid under consideration and a deep layer unaffected by the dynamics. $`x`$ and $`y`$ are respectively the zonal and meridional coordinates ( $`x`$ is directed eastward and $`y`$ poleward ). The term $`h(y)`$ represents the combined effect of the planetary vorticity gradient and of a given stationary zonal flow in the deep layer, with stream function $`\psi _d(y)`$: $`h(y)=\beta y+\psi _d/R^2`$. This deep flow induces a constant deformation of the free surface, acting like a topography on the active layer. We shall therefore call $`h(y)`$ the ’topography’, and study its influence in section 3. Let us assume $`h(y)=0`$ in this section. We define the QG equations (1,2) in the non-dimensional square $`D=[\frac{1}{2},\frac{1}{2}]^2`$. $`R`$ is then the ratio of the internal Rossby deformation radius $`R^{}`$ to the physical scale of the domain $`L`$.
Let $`f_Dfd^2𝐫`$ be the average of $`f`$ on $`D`$ for any function $`f`$. Physically, as the stream function $`\psi `$ is related to the geostrophic pressure, $`\psi `$ is proportional to the mean height at the interface between the fluid layer and the bottom layer, and due to the mass conservation it must be constant (Pedlosky 1987). We make the choice
$$\psi =0$$
(4)
without loss of generality.
The total circulation is $`q=\mathrm{\Delta }\psi +\psi /R^2=\psi /R^2`$ due to the periodic boundary conditions. Therefore
$$q=0$$
(5)
We note that the Dirichlet problem (2) on $`D`$ with periodic boundary conditions has a unique solution $`\psi `$ for a given PV field.
Due to the periodic conditions for $`\psi `$, the linear momentum is also equal to 0,
$$𝐯=0$$
(6)
The energy
$$E=\frac{1}{2}_Dq\psi d^2𝐫=\frac{1}{2}_D[(\psi )^2+\frac{\psi ^2}{R^2}]d^2𝐫$$
(7)
is conserved ( we note that the first term in the right hand side of (7) is the kinetic energy whereas the second one is the gravitational available potential energy ).
The integrals
$$C_f(q)=_Df(q)d^2𝐫$$
(8)
for any continuous function $`f`$ are also conserved, in particular the different moments of the PV. In the case of an initial condition made of a finite number of PV levels, the areas initially occupied by each of these levels is conserved, and this is equivalent to the conservation of all the constants of motion (8).
### 2.2 The statistical mechanics on a two PV levels configuration.
#### 2.2.1 The macroscopic description.
The QG equations (1) (2) are known to develop very complex vorticity filaments. Because of the rapidly increasing amount of information it would require, as time goes on, a deterministic description of the flow for long time is both unrealistic and meaningless. The statistical theory adopts a probabilistic description for the vorticity field. The statistical equilibrium depends on the energy and of the global probability distribution of PV levels. Various previous studies (Sommeria & al 1991), (Kazantsev & al 1998) indicate that a model with only two PV levels provides a good approximation in many cases. The determination of the statistical equilibrium is then simplified as it depends only on the energy, on the two PV levels, denoted $`q=a_1`$ and $`q=a_1`$ and on their respective areas $`A`$ and $`(1A)`$ in $`D`$. The number of free parameters can be further reduced by appropriate scaling. Indeed a change in the time unit permits to define the PV levels up to a multiplicative constant, and we choose for the sake of simplicity :
$$\frac{a_1a_1}{2}=1$$
(9)
and define the non-dimensional parameter $`B`$ as :
$$B\frac{a_1+a_1}{2}$$
(10)
The condition (5) of zero mean PV imposes that $`a_1A+a_1(1A)=0`$. This means that $`a_1`$ and $`a_1`$ must be of opposite sign and, using (9) and (10), $`A=(1B)/2`$. The distribution of PV levels is therefore fully characterized by the single asymmetry parameter $`B`$, which takes values between -1 and +1. The symmetric case of two PV patches with equal area $`A=1/2`$ corresponds to $`B=0`$, while the case of a patch with small area (but high PV, such that $`q=0`$) corresponds to $`B1`$. Note that we can restrict the discussion to $`B1`$ as the QG system is symmetric by a change of sign of the PV.
The two PV levels mix due to turbulent effects, and the resulting state is locally described by the local probability (local area proportion) $`p(𝐫)`$ to find the first level at the location $`𝐫`$. The probability to find the complementary PV level $`a_1`$ is $`1p`$, and the locally averaged PV at each point is then
$$\overline{q}(𝐫)=a_1p(𝐫)+a_1(1p(𝐫))=2\left(p\frac{1}{2}\right)+B$$
(11)
where the second relation is obtained by using (9) and (10).
As we consider the evolution of two PV patches, the conservation of all invariants (8) is equivalent to the conservation of the area $`A`$ of the patch with PV value $`a_1`$ ( the area of the other PV level $`a_1`$ being $`1A`$ ). The integral of $`p`$ over the domain must be therefore equal to the initial area $`A`$ ( the patch with PV level $`a_1`$ is mixed but globally conserved ),
$$A\frac{1B}{2}=_Dp(𝐫)d^2𝐫$$
(12)
As the effect of local PV fluctuations is filtered out by integration, the stream function and the velocity field are fully determined by the locally averaged PV $`\overline{q}`$ as the solution of
$$\overline{q}=\mathrm{\Delta }\psi +\frac{\psi }{R^2};\psi periodic$$
(13)
$$and𝐯=\widehat{𝐳}\psi $$
Therefore the energy is also expressed in terms of the field $`\overline{q}`$ :
$$E=\frac{1}{2}_D\left[(\psi )^2+\frac{\psi ^2}{R^2}\right]d^2𝐫=\frac{1}{2}_D\overline{q}\psi d^2𝐫$$
(14)
Here the energy of the ’microscopic’ PV fluctuations has been neglected (replacing $`q`$ by $`\overline{q}`$), as justified in the case of Euler equation by Robert and Sommeria (1991). Indeed, considering a ’cutoff’ for the microscopic fluctuations much smaller than $`R`$, the small scale dynamics coincides with the Euler case.
The central result of the statistical mechanics of the QG equations (1,2) is that, under an ergodic hypothesis, we expect the long time dynamics to converge towards the Gibbs states defined by maximizing the mixing entropy
$$S=_D[p(𝐫)\mathrm{ln}p(𝐫)+(1p(𝐫))\mathrm{ln}(1p(𝐫))]d^2𝐫$$
(15)
under the constraints of the global PV distribution (12) and energy (14). It can be shown that the microscopic states satisfying the constraints given by the conservation laws are overwhelmly concentrated near the Gibbs state, which is therefore likely to be reached after a complex flow evolution. A good justification of this statement is obtained by the construction of converging sequences of approximations of the QG equation (1,2), in finite dimensional vector spaces, for which a Liouville theorem holds. This is a straightforward translation of the work of Robert (1999) for 2D Euler equations. The sequence of such Liouville measures has then the desired concentration properties as (1,2) enters in the context considered in Michel & Robert (1994).
#### 2.2.2 The Gibbs states
Following Robert & Sommeria (1991), we seek maxima of the entropy (15) under the constraints (12) and (14). To account for these constraints, we introduce two corresponding Lagrange multipliers, which we denote $`2\alpha `$ and $`C/R^2`$ for convenience in future calculations. Then the first variation of the functionals satisfies :
$$\delta S2\alpha \delta A+\frac{C}{R^2}\delta E=0$$
for all variations $`\delta p`$ of the probability field $`p`$. After straightforward differentiation we obtain:
$$\delta S=_D[\mathrm{ln}p\mathrm{ln}(1p)]\delta pd^2𝐫,\delta A=_D\delta pd^2𝐫$$
$$\delta E=_D\psi \delta \overline{q}d^2𝐫=_D2\psi \delta pd^2𝐫$$
(16)
where the expression of $`\delta E`$ has been obtained by integrating by part and expressing $`\overline{q}`$ by (11). Then we can write the first variation under the form $`_D[\mathrm{ln}p+\mathrm{ln}(1p)2\alpha +2C\psi /R^2]\delta pd^2𝐫`$ which must vanish for any small variation $`\delta p`$. This implies that the integrand must vanish, and yields the equation for the optimum state:
$$p=\frac{1\mathrm{tanh}(\alpha \frac{C\psi }{R^2})}{2},$$
(17)
and using (11) and (13), the partial differential equation
$$q=\mathrm{\Delta }\psi +\frac{\psi }{R^2}=B\mathrm{tanh}\left(\alpha \frac{C\psi }{R^2}\right)$$
(18)
determining the Gibbs states (statistical equilibrium). From now on we forget the $`q`$ over-line for the locally averaged PV and refer to it as the PV.
Therefore, we have shown that for any solution of the variational problem, two constants $`\alpha `$ and $`C`$ exist such that $`\psi `$ satisfies (18). Conversely it can be proved that for any such two constants, a solution to equation (18), in general not unique, always exists. Then $`p`$ associated with one of these solutions by (17) is a critical point of the ’free energy’ $`S(p)+2\alpha A(p)\frac{C}{R^2}E(p)`$ (i.e. its first variation vanishes). Then the Lagrange multipliers are not given but have to be calculated by prescribing the constraints (12) and (14) corresponding to the two parameters $`B`$ and $`E`$ respectively, given by the initial condition. Furthermore, among the states of given energy $`E`$ and asymmetry parameter $`B`$, we shall select the actual maxima.
Finally, let us find a lower bound for the parameter $`C`$ of the Gibbs states with non-zero energy (i.e. $`\psi `$ is not constant on $`D`$). Multiplying (18) by $`\mathrm{\Delta }\psi `$, integrating by part and defining $`f(C\psi )B\mathrm{tanh}(\alpha \frac{C\psi }{R^2})`$, we obtain :
$$C=\frac{_D\left((\mathrm{\Delta }\psi )^2+\frac{1}{R^2}(\psi )^2\right)d^2𝐫}{_Df^{}(C\psi )(\psi )^2d^2𝐫}$$
From which, using $`0<f^{}(C\psi )\frac{1}{R^2}`$ it follows that when $`\psi `$ is not constant :
$$C>1$$
(19)
### 2.3 The limit of small Rossby deformation radius
As suggested by oceanographic or Jovian parameters, we seek solutions for the Gibbs states equation in the limit of a small ratio between the Rossby deformation radius and the length scale of the domain : $`R<<1`$ with our non-dimensional coordinates <sup>1</sup><sup>1</sup>1 Modica (1987) considered the minimization of the functional $`E_ϵ(u)=_\mathrm{\Omega }[ϵ\left(u\right)^2+W_0(u)]𝑑𝐱`$ with the constraint $`_\mathrm{\Omega }u(𝐱)𝑑𝐱=m`$ in the limit $`ϵ0^+`$ where $`W_0`$ is a real function with two relative minima. He proved, in a mathematical framework, working with bounded variations functions, that if $`(u_ϵ)`$ are solutions of this variational problem, for any subsequence of $`(u_ϵ)`$ converging in $`L^1(\mathrm{\Omega })`$ as $`ϵ0`$, this subsequence converge to a function $`u_0`$ which takes only the values where $`W_0`$ reaches its minima ; with the interface between the corresponding subdomains having minimal area ( See Modica (1987) for a precise statement ). We note that the Euler equation of this variational problem may be the same as the Gibbs States equation (18) for a convenient choice of $`W_0`$. However as the variational problem itself is different this beautiful result cannot be used in our context.
#### 2.3.1 The uniform subdomains
Then we expect that the Laplacian in the Gibbs states equation (18) can be neglected with respect to $`\psi /R^2`$, except possibly in transition regions of small area. This transforms (18) into the algebraic equation :
$$q=\frac{\psi }{R^2}=B\mathrm{tanh}\left(\alpha \frac{C\psi }{R^2}\right)$$
(20)
Depending on the parameters, this equation has either one, two or three solutions, denoted $`\psi _1,\psi _0`$ and $`\psi _1`$ in increasing order (see figure 2 ). The case with a single solution would correspond to a uniform $`\psi `$, which should be equal to 0 due to the condition $`\psi =0`$. This is only possible for $`E=0`$. Otherwise, we have therefore two or three solutions, with different solutions occurring in subdomains. This condition of multiple solutions requires that the maximum slope for the right hand side of (20) must be greater than $`1/R^2`$ ; this is always realized due to the inequality (19). Furthermore $`\alpha `$ must be in an interval centered in $`CB`$ ( $`\alpha =CB`$ in the symmetric case of figure 2 ).
At the interface between two constant stream function subdomains, a strong gradient of $`\psi `$ necessarily occurs, corresponding to a jet along the interface. These jets give first order contributions to the entropy and energy, but let us first describe the zero order problem. Suppose that $`\psi `$ takes the value $`\psi _1`$ ( resp $`\psi _1`$ ) in subdomains of total area $`A_1`$ ( resp $`A_1`$ ). The reason why we do not select the value $`\psi _0`$ will soon become clear. Using (11) we conclude that the probability $`p`$ takes two constants values $`p_{\pm 1}`$ in their respective subdomains. The two areas $`A_{\pm 1}`$ ( measured from the middle of the jet ) are complementary such that :
$$A_1+A_1=1$$
(21)
Furthermore the constraint (4) of zero domain average for $`\psi `$ implies at zero order,
$$\psi _1A_1+\psi _1A_1=0,$$
(22)
or equivalently, using $`q_{\pm 1}=\psi _{\pm 1}/R^2`$, (11) and (21) :
$$2A_1\left(p_1\frac{1}{2}\right)+2(1A_1)\left(p_1\frac{1}{2}\right)=B$$
(23)
This can be obtained as well from the constraint on the (microscopic) PV patch area (12). The energy inside the subdomains reduces to the potential term $`\psi ^2/2R^2`$, since velocity vanishes. This area energy $`E_A`$ can be computed in terms of $`p_{\pm 1}`$ using $`q_{\pm 1}=\psi _{\pm 1}/R^2`$ and (11) :
$$\begin{array}{cc}\hfill E_A& \frac{\psi _1^2A_1+\psi _1^2A_1}{2R^2}=A_1e(p_1)+(1A_1)e(p_1)\hfill \\ \hfill \text{with}e(p)& R^2\left(2\left(p\frac{1}{2}\right)^2+2B\left(p\frac{1}{2}\right)+\frac{B^2}{2}\right)\hfill \end{array}$$
(24)
There is also an energy in the jet at the interface of subdomains, but it is small with respect to $`E_A`$. Indeed the velocity in the jet, of width $`R`$, is of order $`(\psi _1\psi _1)/R\psi /R`$, and the corresponding integrated kinetic energy is of order $`\psi ^2/R`$. This is small in comparison with the area energy $`E_A`$ (mostly potential) of order $`\psi ^2/R^2`$. A precise calculation will confirm this estimate in next sub-section.
We need to determine three unknown, the area $`A_1`$ and the probabilities $`p_{\pm 1}`$ of the PV level $`a_1`$ in each subdomain, while the constraints (23) and (24) provide two relations. An additional relation will be given by entropy maximization. As we neglect the jet area, the entropy reduces at order zero to the area entropy :
$$\begin{array}{cc}\hfill S_A& A_1s(p_1)+(1A_1)s(p_1)\hfill \\ \hfill \text{with}s(p)& p\mathrm{log}p(1p)\mathrm{log}(1p)\hfill \end{array}$$
(25)
Thus the zero order problem corresponds to the maximization of the area entropy (25) with respect to the 3 parameters $`p_{\pm 1}`$ and $`A_1`$, under the 2 constraints (23) and (24). A necessary condition for a solution of this variational problem is the existence of two Lagrange parameters $`\alpha _0`$ and $`C_0`$, associated respectively with the circulation constraint (23) and with the energy constraint (24), such that the first variations of the total free energy,
$$F_AS_A\frac{C_0}{R^2}E_A+\alpha _0\frac{\psi }{R^2},$$
(26)
vanish. Let us calculate $`F_A`$ using (23) and (24):
$$\begin{array}{cc}\hfill F_A& =A_1f(p_1)+(1A_1)f(p_1),\hfill \\ \hfill \text{with}f(p)& s(p)2C_0\left(p\frac{1}{2}\right)^22\left(C_0B\alpha _0\right)\left(p\frac{1}{2}\right)\frac{C_0B^2}{2}\alpha _0B.\hfill \end{array}$$
(27)
The vanishing of the variations with respect to $`p_1`$ and $`p_1`$ gives that $`f(p_{\pm 1})`$ are local minima of the free energy $`f(p)`$. It is easily proven that the function $`f`$ has two local minima and one local maximum ( for $`C_0>1`$ and $`(C_0B\alpha _0)`$ small enough ) ( see figure 3). The local maximum is achieved for $`p_0`$ corresponding to the value $`\psi _0`$. It is the reason why it has not been taken into account in this analysis. In addition, the vanishing of the first variations with respect to the area $`A_1`$ imposes the free energies $`f(p_{\pm 1})`$ in the two subdomains to be equal. This is like the condition of thermodynamic equilibrium for a chemical species shared by two coexisting phases.
In the expression (27) of $`f(p)`$, the entropy term $`s(p)`$ is symmetric with respect to $`p=\frac{1}{2}`$, as well as the quadratic term. Therefore if the linear term in $`\left(p\frac{1}{2}\right)`$ vanishes the two maxima are equal, with $`p_{\pm 1}`$ symmetric with respect to $`\frac{1}{2}`$. The addition of a linear term obviously breaks this condition of two equal maxima, so the coefficient of the linear term must vanish, thus :
$$\alpha _0=C_0B.$$
(28)
Since $`p_{\pm 1}`$ are symmetric with respect to $`\frac{1}{2}`$, we introduce the parameter $`u`$ by :
$$p_{\pm 1}=\frac{1}{2}(1\pm u).$$
(29)
Using (11),(23) we deduce:
$$\psi _{\pm 1}=R^2(B\pm u)$$
(30)
From (22) we state that the two constant stream function (30) have to be of opposite sign, so that $`u>|B|`$. Introducing (29) in the circulation constraint (23), and using (21), we get :
$$A_{\pm 1}=\frac{1}{2}\left(1\frac{B}{u}\right).$$
(31)
Using these results, the energy (24) becomes
$$EE_A=\frac{R^2}{2}(u^2B^2)$$
(32)
This relates the parameter $`u`$ to the given energy $`E`$ and asymmetry parameter $`B`$. Finally the condition that $`f(p_{\pm 1})`$ are maxima of $`f`$ leads to :
$$u=\mathrm{tanh}(C_0u),$$
(33)
which determines the ’temperature’ parameter $`C_0`$, as represented in figure 4. Therefore all the quantities are determined from the asymmetry parameter $`B`$ and from the parameter $`u`$, related to the energy by (32).
In the limit of low energy, $`u|B|`$, when for instance $`B>0`$, then $`A_1`$ goes to zero, so that $`\psi _1`$ tends to occupy the whole domain. This state is the most mixed one compatible with the constraint of a given value of $`B`$ (or equivalently a given initial patch area $`A=(1B)/2`$). In the opposite limit $`u1`$, we see from (30) that in the two subdomains $`q=\psi /R^2`$ tends to the two initial PV levels $`a_1=1+B`$ and $`a_1=1+B`$. Thus, this state is an unmixed state. It achieves the maximum possible energy $`E=\frac{R^2}{2}(1B^2)`$ under the constraint of a given patch area. We conclude that the parameter $`u`$, or the related ’temperature’ $`C_0`$, linked with the difference between the energy and the maximum accessible energy for the two given initial PV levels, characterizes the mixing of these two PV levels. We shall call $`u`$ the segregation parameter, as it quantifies the segregation of the PV level $`a_1`$ ( or its complementary $`a_1`$ ) between the two phases.
Let us now study the interface between the subdomains.
#### 2.3.2 Interior Jets
At the interface between two constant stream function subdomains, a strong gradient of $`\psi `$ necessarily occurs, corresponding to a jet along the interface. To study these jets, we come back to the Gibbs state equation (18). We expect the Lagrange parameters $`\alpha `$ and $`C`$ to be close to the zero order parameters $`\alpha _0`$ and $`C_0`$, computed in the previous sub-section, so we use $`\alpha =\alpha _0`$ and $`C=C_0`$ to calculate the jet structure. In such a jet, we cannot neglect the Laplacian term in (18), but a boundary layer type approximation can be used: we neglect the tangential derivative with respect to the derivative along the coordinate normal to the interface, $`\zeta `$. Accordingly, we neglect the inverse of the curvature radius of the jet with respect to $`1/R`$.
Thus, from the Gibbs states equation (18), using (28), we deduce the jet equation :
$$\frac{d^2\psi }{d\zeta ^2}+\frac{\psi }{R^2}=B\mathrm{tanh}\left(C_0\left(B\frac{\psi }{R^2}\right)\right)$$
(34)
As the stream function depends only on the normal coordinate $`\zeta `$, the velocity is tangent to the interface, forming a jet with a typical width scaling like $`R`$. We thus make the change of variables defined by :
$$\tau \frac{\zeta }{R};\varphi B+\frac{\psi }{R^2},$$
(35)
leading to the rescaled jet equation:
$$\frac{d^2\varphi }{d\tau ^2}=\mathrm{tanh}(C_0\varphi )+\varphi $$
(36)
The jet equation (36) is similar to a one dimensional equation of motion for a particle ( with the position $`\varphi `$ depending on a time $`\tau `$) under the action of a force $`dU/d\varphi `$ deriving from the potential,
$$U(\varphi )=\frac{\mathrm{ln}(\mathrm{cosh}(C_0\varphi ))}{C_0}\frac{\varphi ^2}{2},$$
(37)
represented in figure 2(b). In its trajectory the particle energy is conserved :
$$\frac{1}{2}\left(\frac{d\varphi }{d\tau }\right)^2+U(\varphi )=Cst$$
(38)
Let $`\varphi _i\psi _i/R^2B`$, $`i=1,0,1`$, corresponding to the solutions $`\psi _i`$ of the algebraic equation (20). From (30), we have $`\varphi _{\pm 1}=\pm u`$. Note that the stationary limit of (36), which must be reached for $`lim_{\tau \pm \mathrm{}}`$, yields again (33). Moreover, the particle energy conservation (38) imposes the integrability condition,
$$U(\varphi _1)=U(\varphi _1),$$
(39)
which is indeed satisfied due to the symmetry of the potential $`U`$. We note that the Lagrange parameter determination (28) and the symmetry of the probabilities (29) with respect to $`\frac{1}{2}`$ could have been deduced from this integrability condition (39) instead of minimizing the free energy (27) (we shall proceed in this way in section 3 to take into account the beta-effect).
The jet equation (36) has been numerically integrated. Figure 5 shows a typical stream function and velocity profile in the jets. Figure 6(a) shows how the jet width depends on the segregation parameter $`u`$. We note that the width of the jet is an increasing function of the mixing and therefore a decreasing function of the energy. Figure 6(b) shows the dependence in $`u`$ of the total non-dimensional energy $`e(u)=\frac{1}{2}_{\mathrm{}}^+\mathrm{}\left(d\varphi /d\tau \right)^2𝑑\tau `$ and of the maximum non dimensional jet velocity $`\left(d\varphi /d\tau \right)_{max}`$.
As the jet structure (36) does not depend on the coordinate tangent to the jet, we can define the jet entropy ( respectively energy, free energy ) per unit length $`S_{Jet}`$ ( respectively $`E_{Jet}`$, $`F_{Jet}`$ ). Multiplied by the jet length, these quantities are the first order corrections to the entropy ( respectively energy, free energy ). Using the change of variables (35), we calculate the jet entropy per unit length :
$$S_{Jet}=R_{\mathrm{}}^+\mathrm{}[s(p(\tau ))s(p_{\pm 1})]𝑑\tau $$
where $`s`$ is defined in (25), and $`p_{\pm 1}`$ are defined in (29). Using the probability equation (17) and (35) we obtain :
$$S_{Jet}=R_{\mathrm{}}^+\mathrm{}[\stackrel{~}{s}(\varphi ))\stackrel{~}{s}(\varphi _{\pm 1}))]d\tau $$
(40)
involving the function $`\stackrel{~}{s}(\varphi )\mathrm{ln}(\mathrm{cosh}(C_0\varphi ))C_0\varphi \mathrm{tanh}(C_0\varphi )`$. Similarly we straightforwardly calculate the potential and kinetic energy per unit length for the jet :
$$E_{Jet}^P=\frac{R^3}{2}_{\mathrm{}}^+\mathrm{}(\varphi ^2\varphi _1^2)𝑑\tau E_{Jet}^K=\frac{R^3}{2}_{\mathrm{}}^+\mathrm{}\left(\frac{d\varphi }{d\tau }\right)^2𝑑\tau $$
We use the integral (38) to calculate $`d\varphi /d\tau `$ and conclude :
$$E_{Jet}=R^3_{\mathrm{}}^+\mathrm{}[\stackrel{~}{g}(\varphi )\stackrel{~}{g}(\varphi _{\pm 1})]𝑑\tau ,$$
(41)
with $`\stackrel{~}{g}(\varphi )=\frac{\mathrm{ln}(\mathrm{cosh}(C_0\varphi )}{C_0}+\varphi ^2`$. Due to the symmetry of the jets, the jets provide no perturbation to the zero-order circulation, so there is no circulation term in the jet free energy expression : $`F_{Jet}=S_{Jet}C_0/R^2E_{Jet}`$. Then
$$F_{Jet}=C_0R_{\mathrm{}}^+\mathrm{}[\stackrel{~}{h}(\varphi )\stackrel{~}{h}(\varphi _{\pm 1})]𝑑\tau $$
(42)
with $`\stackrel{~}{h}(\varphi )=\varphi (\varphi \mathrm{tanh}(C_0\varphi ))`$.
Let us study the sign of $`F_{Jet}`$. As $`\varphi _1`$ verifies $`\varphi _1=\mathrm{tanh}(C_0\varphi _1)`$ and as $`\varphi (\tau )`$ is an increasing function of $`\tau `$ with $`lim_{\tau +\mathrm{}}\varphi (\tau )=\varphi _1`$ we conclude that $`\stackrel{~}{h}(\varphi )\stackrel{~}{h}(\varphi _1)>0`$ for any $`\tau >0`$. Thus $`F_{Jet}>0`$. Using the analogy with usual thermodynamics, the ‘surface tension’ is positive. This favors large ‘bubbles’ which minimize the interfacial length and therefore the corresponding free energy (27). Our initial hypothesis of well separated domains with uniform $`\psi `$ is thus supported, as discussed more precisely in next subsection.
#### 2.3.3 Selection of the sub-domain shape
The above analysis has permitted us to determine the areas of subdomains on which the stream function $`\psi `$ takes the constant values $`\psi _{\pm 1}`$, but the subdomains shape is still to be selected. There is an analogy with two phases coexisting in thermodynamic equilibrium, for instance a gas bubble in a liquid medium, for which a classical thermodynamic argument explains the spherical shape of the bubble by minimizing its free energy, proportional to the bubble area. Our system is isolated rather than in a thermal bath, but the jet energy is small (of order $`R`$) with respect to the total energy. Therefore the subdomain interior behaves like a thermal bath with respect to the jet, so the usual argument on free energy minimization applies. We shall now show this more precisely by directly maximizing the total entropy with constraints, taking into account the jet contribution.
The jet with length $`L`$ has an entropy $`S_{jet}L`$ and energy $`E_{jet}L`$. Since the total energy $`E=E_A(C)+LE_{Jet}`$ is given, the jet has also an indirect influence in the area energy $`E_A`$. This small energy change $`\delta E_A`$ results in a corresponding change in the area entropy $`\delta S_A=(C_0/R^2)\delta E_A`$, from the condition (2.2.2) of zero first variations. Note that there is no area change $`\delta A`$ since the jet is symmetric and has therefore no influence in the condition (12) of a given integral of $`p`$ (the difference in $`p`$ with the case of two uniform patches with boundary at the jet center has zero integral). Therefore, adding the direct and indirect contribution of the jet entropy leads to the total entropy
$$S=S_A(C_0)+\left(\frac{C_0}{R^2}E_{Jet}+S_{Jet}\right)L=S_A(C_0)F_{Jet}L$$
(43)
where $`S_A(C_0)`$ is the zero order, area entropy, obtained in the limit of vanishing jet width <sup>2</sup><sup>2</sup>2 This reasoning to obtain the first order entropy can be precised by evaluating explicitly the first order modification of the Lagrange parameter $`C`$ ( let say $`C_1CC_0`$ ) due to the jet energy, and the first order modification of the Lagrange parameter $`\alpha `$ ( let say $`\alpha _1\alpha \alpha _0`$ ) due to the jet curvature and computing the first order entropy from its definition (15). We have calculated the first order entropy in this way to actually obtain (43).
We deduce from (43) that the maximization of the entropy is achieved by minimizing the total free energy $`F_{Jet}L`$, which we have proved to be positive at the end of previous sub-section. Thus we conclude that the maximum entropy state minimizes the jet length, with a given area of the subdomains (31). The subdomains shape will therefore be a circle or a stripe. More precisely if $`A_1<1/\pi `$ the jet forms a circle enclosing the positive constant stream function domain ( the jet bounds a cyclonic vortex ), if $`1/\pi <A_1<11/\pi `$ two straight lines jets form a stripe and if $`A_1>11/\pi `$ the jet form a circle enclosing the negative constant stream function domain ( the jet bounds anti-cyclonic vortex ).
The different types of solutions can be summarized in a (E,B) diagram : figure 7. The outer parabola is the maximum energy achievable for a fixed B : $`E=R^2(1B^2)/2`$. The frontier lines between the straight jets and the circular jets corresponds to $`A_1=1/\pi `$ or $`A_1=1/\pi `$. It has been calculated using (31) and (32) : $`E=R^2B^2(2\pi 2)/(\pi 2)^2`$. Note that the maximum accessible energy is in $`R^2`$, but it has been scaled by the normalization condition (9) on PV levels, so the real energy is not bounded.
All this analysis assumes that the vortex size is much larger than the jet width $`\lambda `$, given by figure 6. In other words, the area $`A_1`$ or $`A_1`$ (31) must be larger than $`(2\lambda )^2`$. This is satisfied on the right side of the dashed line represented in figure 7. The dashed line itself corresponds to the equality, and the condition of large vortex is clearly not satisfied on its left side, for low energy. The position of the dashed line depends on the numerical value of $`R`$ (it has been here numerically computed for $`R=0.03`$), and it gets closer to the origin as $`R0`$. We shall now determine the statistical equilibrium in this case of low energy.
#### 2.3.4 Axisymmetric vortices
We have noted in subsection (2.3.1) that, when for instance $`B>0`$, in the limit of small energy ($`E0`$ or equivalently $`u|B|`$, for fixed $`B`$ and $`R`$), the area $`A_1`$ occupied by $`\psi _1`$ tends to 1, the whole domain. Therefore, in this limit, the complementary area $`A_1`$ tends to 0 and the vortex becomes smaller than the deformation radius, so we can no more neglect the curvature radius of the jet.
In this limit $`u|B|`$, as the vortex has a small area with respect to the total domain, it is not affected by the boundary conditions, so it can be supposed axisymmetric. From the general Gibbs states equation (18), we deduce the axisymmetric vortex equation :
$$\frac{d^2\psi }{d\zeta ^2}\frac{1}{\zeta }\frac{d\psi }{d\zeta }+\frac{\psi }{R^2}=B\mathrm{tanh}\left(\alpha \frac{C\psi }{R^2}\right)$$
(44)
As $`R`$ will be a typical scale of the vortex, we make the change of variable,
$$\zeta =Rr;\varphi =\frac{\alpha }{C}+\frac{\psi }{R^2},$$
(45)
leading to the rescaled axisymmetric vortex equation :
$$\frac{d^2\varphi }{dr^2}\frac{1}{r}\frac{d\varphi }{dr}+\varphi +\frac{\alpha }{C}=B+\mathrm{tanh}(C\varphi ))$$
(46)
From now on, we shall consider the case $`B>0`$ (the case $`B<0`$ is just the symmetric case of a negative vortex).
For this equation to describe a localized vortex, we impose $`lim_r\mathrm{}\varphi (r)=\varphi _1\alpha /C+\psi _1/R^2`$, where $`\psi _1`$ is the positive solution of the algebraic equation (20). Since nearly the whole fluid domain is covered by the asymptotic stream function $`\psi _1`$ outside the vortex, the condition of zero total circulation $`q=\psi /R^2=0`$ imposes that $`\psi _10`$ (it is of order $`R`$), so that $`\varphi _1=\alpha /C`$, and the algebraic equation (20) then leads to:
$$\alpha =\mathrm{arg}\mathrm{tanh}(B)$$
(47)
We can thus eliminate $`\alpha `$ in (46), leading to an equation depending on two parameters, $`B`$ and $`C`$,
$$\begin{array}{c}\hfill \frac{d^2\varphi }{dr^2}=\frac{1}{r}\frac{d\varphi }{dr}\mathrm{tanh}(C\varphi )+\varphi B+\frac{\mathrm{arg}\mathrm{tanh}B}{C}\\ \hfill \frac{d\varphi }{dr}(r=0)=0and\underset{r\mathrm{}}{lim}\varphi (r)=\frac{\mathrm{arg}\mathrm{tanh}B}{C}\end{array}$$
(48)
where the regularity condition at $`r=0`$ has been included. Let us consider, as in section (2.3.2), the analogy of equation (48) with a one particle motion with ’position’ $`\varphi `$ and ’time’ $`r`$.
The last four terms on the right-hand side of (46) can be written as the derivative $`dU/d\varphi `$ of the potential,
$$U(\varphi )=\frac{\mathrm{ln}(\mathrm{cosh}(C\varphi ))}{C}\frac{\varphi ^2}{2}+\left(B\frac{\mathrm{arg}\mathrm{tanh}B}{C}\right)\varphi ,$$
(49)
(represented in figure 8), while the first term can be interpreted as a friction effect. Indeed, an integration of (48) leads to:
$$U(\varphi _1)U(\varphi (r=0))=_0^+\mathrm{}\frac{1}{r}\left(\frac{d\varphi }{dr}\right)^2𝑑r<0$$
(50)
Thus, in figure 8(a), the hatched area on the right side must be greater than the one on the left (since $`(U(\varphi _1)U(\varphi _1))>U(\varphi (r=0)U(\varphi _1)>0`$). It is clear from figure 8 that this is possible only if $`\varphi _0<0`$ and $`\alpha /C<B`$, or, using (47), $`C>\alpha /B=\mathrm{arg}\mathrm{tanh}B/B`$. The value $`C=\alpha /B`$ corresponds to the integrability condition (39) when the effect of jet curvature is neglected. This effect is now taken into account by the departure of $`C`$ from this value, which we shall denote $`\mathrm{\Delta }CC\mathrm{arg}\mathrm{tanh}B/B`$. Then $`\mathrm{\Delta }C>0`$ and we expect to recover the results of section 2.3.2 in the limit $`\mathrm{\Delta }C0`$. Moreover, we must reach a uniform stream function at large distance, solution of the algebraic equation(20), so it must have three solutions. We see in figure 8 that the corresponding $`\mathrm{\Delta }C`$ must not exceed a maximal value, denoted $`\mathrm{\Delta }C_{max}`$ .
We can prove that for any $`B>0`$ and $`\mathrm{arg}\mathrm{tanh}B/B<C<\mathrm{arg}\mathrm{tanh}B/B+\mathrm{\Delta }C_{max}`$, equation (48) has a unique solution. Such solutions have been numerically obtained for $`B=0.75`$ and $`0<\mathrm{\Delta }C<\mathrm{\Delta }C_{max}`$. Corresponding stream function profiles are shown in figure 9.
As $`\mathrm{\Delta }C`$ is decreased from $`\mathrm{\Delta }C_{max}`$ to zero, two stages can be seen in figure 9. First the maximum value for the stream function is increased while the mean width of the vortex remains of the order of $`R`$. In a second stage, when $`\mathrm{\Delta }C`$ goes to zero, as we are closer to the integrability condition for big vortices (39), $`\varphi `$ remains longer in the vicinity of $`\varphi _1`$ so the vortex size increases. Note that the energy monotically increases as $`\mathrm{\Delta }C`$ is decreased, first by an increase in the vortex maximum stream function and then by an increase in size. Finally the case of a jet with negligible curvature studied in subsection 2.3.2 is reached when $`\mathrm{\Delta }C0`$.
In conclusion, we have shown that in the limit of small energy, with fixed $`B`$ and $`R`$, the Gibbs states are approximated by axisymmetric vortices, whose radial structure depends on the parameter $`\mathrm{\Delta }C`$, which monotonically decreases from $`\mathrm{\Delta }C_{max}`$ to 0 as energy is increased.
#### 2.3.5 The linear approximation for the Gibbs states
The previous discussion of axisymmetric vortices was concerned with the limit of small energy with fixed $`B`$ and (small) fixed $`R`$. We consider now the limit of small $`E`$ and $`B`$, i.e. the neighborhood of the origin in the phase diagram of figure 7. Then from (32), $`u|B|<<1`$. Figure 6 shows that for $`|u|<<1`$, the jet width diverges and therefore the jet tends to develop on the scale of the whole domain, so the approximation of a localized jet, or an isolated axisymmetric vortex, falls down.
In this limit of small $`E`$ and $`B`$, we can however linearize the Gibbs states equation (18), following the work of Chavanis and Sommeria (1996) for Euler equation. After linearization, solutions are expressed in terms of the eigenmodes of the Laplacian, and only the first eigenmodes can be entropy maxima.
These results are unchanged by the linear deformation term $`\psi /R^2`$, so the work of Chavanis and Sommeria (1996) directly applies here. With the periodic boundary conditions, the first eigenmode of the Laplacian, a sine function of one of the coordinates, for instance $`y`$, is thus selected. This corresponds to the low energy limit of the two jet configuration shown in figure 7. The next eigenmode, in $`\mathrm{sin}(\pi x)\mathrm{sin}(\pi y)`$, has the topology of the vortex states. A competition between these two modes is expected in the neighborhood of the origin for small $`E`$ and $`B`$. Note that the range of validity of the linear approximation is limited to a smaller range of parameters than in the Euler case, and this range of validity gets smaller and smaller as $`R0`$. The dominant solution with uniform subdomains and interfacial jets relies by contrast on the tanh like relation between PV and stream function, and it is genuinely non-linear.
## 3 The channel case
We now consider the channel geometry, which represents a zonal band around a given latitude. It is then natural to introduce a beta-effect, or a mean sublayer zonal flow ( topography ), with the term $`h(y)`$ in (1). We shall study the two cases of a linear $`h(y)`$ ( beta effect and/or uniform velocity for the sublayer flow ) or a quadratic $`h(y)`$. We shall follow the presentation for the periodic boundary conditions, stressing only the new features.
### 3.1 The dynamical system
Let us consider the barotropic QG equations (1, 2 and 3) in a channel $`D=[\frac{1}{2},\frac{1}{2}]^2`$ with the velocity $`𝐯`$ tangent to the boundary for $`y=\pm \frac{1}{2}`$ and 1-periodicity in the zonal direction. Thus we choose for the boundary conditions a constant $`\psi `$, denoted $`\psi _b`$, the same on the two boundaries $`y=\pm \frac{1}{2}`$. We note that, due to these conditions, the physical momentum (6) is equal to zero. It is always possible to satisfy this condition by a change of reference frame with a zonal velocity $`V`$ such that it moves with the center of mass of the fluid layer, and a corresponding change of the deep flow, resulting in an additional beta effect $`hh+\frac{Vy}{R^2}`$.
As in section (2), we need to specify the gage constant in the stream function $`\psi `$, and we generalize the integral condition (4) as,
$$\frac{\psi }{R^2}h(y)=0.$$
(51)
The total mass $`\psi `$ is then constant in time (but not the boundary value $`\psi _b`$ in general). With these conditions, the Dirichlet problem (2) has a unique solution $`\psi `$ for a given PV field $`q`$. We note that the scale unit is chosen such that the area of $`D`$ is equal to 1.
The integral of any functions of the potential vorticity (8) is still conserved. Let in particular $`\mathrm{\Gamma }`$ be the global PV, or circulation :
$$\mathrm{\Gamma }q=_D\mathrm{\Delta }\psi d^2𝐫=_D𝐯.d𝐥$$
(52)
By contrast with the doubly periodic boundary conditions, the circulation $`\mathrm{\Gamma }`$ is not necessarily equal to zero. The expression of the energy in terms of the PV (see equation (7)) is therefore modified (due to the boundary term in the integration by parts) :
$$E=\frac{1}{2}_D\left[(\psi )^2+\frac{\psi ^2}{R^2}\right]d^2𝐫=\frac{1}{2}_D(q+h(y))\psi d^2𝐫\frac{1}{2}\mathrm{\Gamma }\psi _b$$
(53)
Due to the invariance under zonal translation of the system, another conserved quantity exists :
$$M=_Dyqd^2𝐫$$
(54)
This constant moment fixes the ‘center of mass’ latitude for the PV field.
### 3.2 General form of the Gibbs states
et us consider the statistical mechanics on a two PV level configuration : the initial states is made of patches with two levels of potential vorticity, $`q=a_1`$ and $`q=a_1`$, occupying respectively the areas $`A`$ and $`(1A)`$ in $`D`$. We keep the normalization (9) and definition (10) for $`B`$. Now, since the circulation $`\mathrm{\Gamma }`$ is non-zero, the area $`A`$ is related to $`B`$ by $`A=(1B)/2+\mathrm{\Gamma }/2`$. The boundary term in the expression of the energy (53) leads to an obvious change in the energy variation (16). Let $`\gamma `$ be the Lagrange multiplier associated with the conservation (54) of the momentum $`M`$.
Adapting the periodic case computations, we then calculate the probability equation and the Gibbs state equation :
$$p=\frac{1\mathrm{tanh}\left(\alpha ^{}\frac{C\psi }{R^2}+\gamma y\right)}{2}$$
(55)
$$q=\mathrm{\Delta }\psi +\frac{\psi }{R^2}h(y)=B\mathrm{tanh}\left(\alpha ^{}\frac{C\psi }{R^2}+\gamma y\right)$$
(56)
with $`\alpha ^{}\alpha +C\psi _b/R^2`$. These results generalize (17) and (18) of section 2.
In the case of a Gibbs state depending on $`x`$, the Lagrange parameter $`\gamma `$ is related to a zonal propagation of the equilibrium structure. The statistical theory only predicts a set of equilibria shifted in $`x`$, but introducing the result back in the dynamical equation yields the propagation. Indeed the Gibbs state equation (56) is of the form $`q=f(\psi ,y)`$, which can be inverted (as it is monotonous in $`\psi `$, see Robert and Sommeria (1991)) to yield:
$$\psi =g(q)+R^2\frac{\gamma y}{C}$$
(57)
where $`g`$ is a function of the potential vorticity. From this relation we calculate the velocity using (3) : $`𝐯=R^2\frac{\gamma }{C}\widehat{𝐱}g^{}(q)\widehat{𝐳}q`$. As PV is advected ( equation (1) ) we obtain :
$$\frac{q}{t}+R^2\frac{\gamma }{C}\frac{q}{x}=0$$
(58)
Thus the PV field is invariant in a frame propagating with the zonal speed $`V_{sr}=R^2\frac{\gamma }{C}`$.
### 3.3 The limit of small Rossby deformation radius.
In this sub-section we propose to analyze the Gibbs state equation (56) in the limit of small deformation radius ( $`R<<1`$ ). The main difference with the periodic case resides in the latudinally depending topography $`h(y)`$, resulting in two effects. Firstly the subdomains of uniform PV are no more strictly uniform, and contain a weak zonal flow. Secondly the jet curvature is no longer constant in general, but depends on the local topography. Using the same boundary layer approximation as in the periodic case, the Laplacian term in the Gibbs state equation (56) will be neglected, except possibly in an interior jet and in the vicinity of the boundaries $`y=\pm \frac{1}{2}`$ ( boundary jets ).
Outside such jets, (56) reduces to an algebraic equation, like previously :
$$\frac{\psi }{R^2}h(y)=B\mathrm{tanh}\left(\alpha ^{}C\frac{\psi }{R^2}+\gamma y\right).$$
(59)
This is like (20), replacing the constants $`B`$ by $`B+h(y)`$ and $`\alpha `$ by $`\alpha ^{}+\gamma y`$. The three solutions can be still visualized by figure 8, but the position of the straight line with respect to the tanh curve now depends on $`y`$, due to the terms $`h(y)`$ and $`\gamma y`$. We assume that this dependence is linear in $`y`$ or varies on scales much larger than $`R`$ so that the Laplacian term remains indeed negligible. The zero order Lagrange parameters $`\alpha ^{}`$, $`C`$, $`\gamma `$, involved in this expression, can be obtained by directly maximizing the entropy by the same method as in section (2.3.1). A relation between the jet curvature and topography is then obtained at first order. This approach is developed in appendix (A).
However it is more simple to proceed differently : we start from the jet equation and show that its integrability condition provides the relation between the jet curvature and topography. To catch this effect, we take into account the radius of curvature of the jet, denoted $`r`$, like in section 2.3.4, but state that $`r`$ is constant across the jet, assumed thin. From the Gibbs state equation (56), using the boundary layer approximation, we thus obtain the jet equation in term of the transverse coordinate $`\zeta `$ :
$$\frac{d^2\psi }{d\zeta ^2}ϵ\frac{1}{r}\frac{d\psi }{d\zeta }+\frac{\psi }{R^2}h(y)=B\mathrm{tanh}\left(\alpha ^{}C\frac{\psi }{R^2}+\gamma y\right).$$
(60)
We have introduced $`ϵ=\pm 1`$ to account for the direction of curvature (keeping $`r>0`$). We define $`ϵ=1`$ (respectively -1) if the curvature of the jet is such that $`\varphi _1`$ ( respectively $`\varphi _1`$ ) is in the inner part of the jet. Note that, in the case of a vortex, as in our notations $`\psi `$ is proportional to the opposite of the pressure the case $`ϵ=1`$ ( resp $`ϵ=1`$ ) corresponds to a cyclone ( resp an anticyclone ).
The algebraic equation (59) depends on three Lagrange parameters, instead of two for the periodic case of previous section, but we have three additional constraints, the condition that $`\psi `$ has the same value at the two boundaries, the circulation constraint (52) and the momentum constraint (54). This will be achieved in general by boundary jets. Let us first study the interior jets.
To study the interior jet, we make the change of variables :
$$\tau \frac{\zeta }{R};\varphi \frac{\alpha ^{}+\gamma y}{C}+\frac{\psi }{R^2}$$
(61)
We assume the variations of $`y`$ in the jet width are negligible ( $`R<<`$ scale of variation of $`h(y)`$ ), so that $`y`$ is treated as a constant. Then we obtain the jet equation :
$$\frac{d^2\varphi }{d\tau ^2}ϵ\frac{R}{r}\frac{d\varphi }{d\tau }+\varphi +\frac{\alpha ^{}+\gamma y}{C}=B+h(y)+\mathrm{tanh}(C\varphi )$$
(62)
with $`\varphi \varphi _{\pm 1}`$ for $`\tau \pm \mathrm{}`$, where again $`\varphi _{\pm 1}`$ corresponds to the solutions of the algebraic equation (59), rescaled as
$$\varphi +\frac{\alpha ^{}}{C}B+\frac{\gamma y}{C}h(y)=\mathrm{tanh}(C\varphi ).$$
(63)
Let us consider, as in section (2.3.4), the analogy of the equation (62) with the motion equation of a particle in the potential :
$$U(\varphi )=\frac{\mathrm{ln}\mathrm{cosh}(C\varphi )}{C}\frac{\varphi ^2}{2}+\left(B+h(y)\frac{\alpha ^{}+\gamma y}{C}\right)\varphi .$$
(64)
Like in section 2.3.4., integration of (62) from $`\mathrm{}`$ to $`+\mathrm{}`$ imposes the integrability condition :
$$U(\varphi _1)U(\varphi _1)=ϵ\frac{R}{r}_{\mathrm{}}^+\mathrm{}\left(\frac{d\varphi }{d\tau }\right)^2𝑑\tau $$
(65)
The second term of the l.h.s. of equation (62) can be interpreted as a friction term: if $`ϵ=1`$, the ’particle’ starting from rest at $`\varphi _1`$ can reach a state of rest at $`\varphi _1`$ only if the difference of ’potential’ corresponds to the energy loss (65) by friction (if $`ϵ=1`$ the same is true in the reversed direction).
As in the periodic case we have made the thin jet assumption $`R<<r`$, so that the friction term ( rhs of (65)) is a correction of order $`R/r`$ : $`U(\varphi _1)U(\varphi _1)=O(R/r)`$. We first neglect it to get the zero order results, so we write $`U(\varphi )=U(\varphi _1)`$. Therefore the two hatched areas in figure 8 must be equal, like in figure 2. Due to the symmetry of the tanh function, this clearly implies that the central solution of the rescaled algebraic equation (63) must be $`\varphi _0=0`$, so that $`\alpha _0^{}/C_0B+\gamma y/C_0h(y)=0`$ (denoting the zero order Lagrange parameters by the index 0). This is possible at different latitudes $`y`$ only if $`\gamma y/C_0h(y)=0`$, or is of order $`R`$ (so that it can be neglected at zero order). Then the integrability condition becomes
$$\alpha _0^{}=C_0B.$$
(66)
Furthermore $`\varphi _{\pm 1}`$ are symmetric with respect to 0, of the form $`\varphi _{\pm 1}=\pm u`$, determined by equation (33), like in section 2.3. This parameter $`u`$ is again related to the energy by (32). Finally, the terms $`\gamma y/C_0h(y)=0`$ and the curvature term disappears in the jet equation (62), which therefore reduces to (36), discussed in section 2.3.2.
The first order solution outside the jet is obtained as a small correction $`\delta \varphi _{\pm 1}`$ to the zero order solutions $`\pm u`$, with also a small correction $`\alpha _1`$ and $`C_1`$ to the parameters $`\alpha _0^{}/C_0`$ and $`C_0`$,
$$\varphi _{\pm 1}=\pm u+\delta \varphi _{\pm 1}(y),\frac{\alpha ^{}}{C}=\frac{\alpha _0^{}}{C_0}+\alpha _1,C=C_0+C_1$$
(67)
From (65), we deduce that $`U(\varphi _1)U(\varphi _1)`$ has the sign of $`ϵ`$. As may be seen in figure 2, when $`\alpha _1`$ is positive, the line $`\varphi +\frac{\alpha ^{}}{C}`$ moves upward, so that $`U(\varphi _1)<U(\varphi _1)`$. Thus $`\alpha _1`$ has the sign opposite to the sign of $`U(\varphi _1)U(\varphi _1)`$ ; and we conclude that $`ϵ\alpha _1`$ is always positive. Introducing this expansion (67) in the algebraic equation (63), using the zero order results (33) and (66), we obtain :
$$C_0\delta \varphi _{\pm 1}(y)=\frac{\alpha _1C_0+C_1[B\pm C_0u(1u^2)]\gamma y+C_0h(y)}{1C_0(1u^2)}$$
(68)
Coming back to the stream function $`\psi `$, using (61), we deduce the corresponding velocity $`𝐯`$ by differentiation with respect to $`y`$. This velocity outside the jet is zonal, along the unit vector $`\widehat{𝐱}`$, and verifies:
$$𝐯=R^2\left(\frac{\frac{dh(y)}{dy}\gamma (1u^2)}{1C_0(1u^2)}\right)\widehat{𝐱}$$
(69)
It is therefore a constant plus a term proportional to the local beta-effect $`\frac{dh}{dy}`$. Notice that the corresponding shear $`\frac{d𝐯}{dy}`$ is stronger than the deep shear $`d_d^\psi /dy^2=R^2d^2h/dy^2`$ by the factor $`[1C_0(1u^2)]^1>1`$. The integrability condition (65) now provides the curvature of the jet. We can approximate the r.h.s. of this relation, of order $`R`$, by the zero order jet profile (36), denoting :
$$e(u)\frac{1}{2}_{\mathrm{}}^+\mathrm{}\left(\frac{d\varphi }{d\tau }\right)^2𝑑\tau .$$
(70)
( see figure 6(b) ). The l.h.s. can be expanded, using (67). We first expand the expression (64) of the potential, $`U(\varphi )=U_0(\varphi )+C_1/C_0\left[\varphi \mathrm{tanh}(C_0\varphi )\mathrm{log}\mathrm{cosh}(C_0\varphi )/C_0\right]+\left[h(y)\gamma y/C_0\alpha _1\right]\varphi `$. We can approximate $`\varphi \pm u`$ in the correction terms, and expand $`U_0(\varphi )=U_0(\pm u)+dU_0/d\varphi (\pm u)`$. The zero order equilibrium condition imposes that $`dU_0/d\varphi (\pm u)=0`$, so that (65) becomes
$$ϵu\left(h(y)\frac{\gamma y}{C_0}\alpha _1\right)=e(u)\frac{R}{r}$$
(71)
This equation (71), expresses the dependence in latitude $`y`$ of the curvature radius $`r`$ of the curve on which the jet is centered, thus defining the shape of the sub-domain interface as a function of the topography.
Without topography and for $`\gamma =0`$, we get a constant jet curvature. The same result was obtained in sub-section 2.3.3 by a different argument of free energy minimization. The parameter $`u`$, related to the energy by (32) and to $`C_0`$ by (33), quantifies the strength of the jet. By contrast, the vortex area is determined by the constraint on PV patch area (parameter $`B`$), but it is also related to the jet curvature, proportional by (71) to the small shift $`\alpha _1`$ in chemical potential and temperature. Likewise the equilibrium temperature at a liquid-gas interface slightly depends on the bubble curvature, due to capillary effects.
As explained in the end of section (3.2) the parameter $`\gamma `$ is linked to the zonal propagation speed of the structure. The term $`\gamma y`$ in (71), combined with a usual beta effect (linear topography term $`h(y)`$), leads to an oscillation with latitude $`y`$ of the jet curvature $`1/r`$, i.e. a meandering jet. Another possibility is an exact compensation of the beta-effect by the $`\gamma y`$ term, leading to a propagating circular vortex, and the selection between these two alternatives is discussed in next sub-section. An oval shaped, zonally elongated vortex, such as on Jupiter, is obtained when this compensation occurs, but with an additional quadratic topography $`h(y)`$. Indeed, to get a zonally elongated vortex, supposed latitudinally centered in zero, the radius of curvature of the jet must decrease for $`y>0`$ and increase for $`y<0`$. As a consequence, we deduce from (71) that the topography must be extremal at the latitude on which the vortex is centered (it actually admits a maximum in the cyclonic case and a minimum in the anticyclonic case ). Moreover, we deduce from (72) that the surrounding flow must have a zero velocity at the latitude on which the vortex is centered and that the shear is cyclonic when the vortex is a cyclone and anticyclonic when the vortex is an anticyclone. More generally, the curvature can be related to the zonal velocity outside the jet, eliminating the topography between (72) and (71),
$$𝐯=R^2\left(\frac{\gamma }{C_0}+\frac{ϵe(u)}{u\left(1C_0(1u^2)\right)}\frac{d}{dy}\left(\frac{R}{r}\right)\right)\widehat{𝐱}$$
(72)
Let us recall our approximations. Writing down the jet equation (60), when making the boundary layer approximation, we have assumed $`R<<r`$. We also assumed that in the jet width the topography can be considered as a constant. If $`1/\sqrt{a}`$ is a typical length scale for the topography variation this gives $`aR^2<<1`$. Moreover we have assumed that the effective topography effect $`h(y)\gamma y/C_0`$ remains small along the jet. If $`L_V`$ denotes the jet extension ( for example a vortex latitudinal size ) this approximation is valid as long as $`aL_V^2<<1`$.
#### 3.3.1 Beta effect or linear topography.
Let $`h(y)\beta y`$ in the following of this section. $`\beta `$ may mimic the beta-effect or a uniform velocity in the sublayer ( but we will refer it as the beta-effect ).
A first class of equilibrium states corresponds to a single solution $`\psi (y)`$ of the algebraic equation (59). This determines a smooth zonal flow, with possibly intense jets at the boundaries $`y=\pm \frac{1}{2}`$. The solution depends on the unknown parameters $`C`$, $`\alpha ^{}`$ and $`\gamma `$, which are indirectly determined by the energy $`E`$, momentum $`M`$, and the condition $`\psi =0`$. The limit of small energy corresponds to $`C\mathrm{}`$, for which we can neglect the term $`\psi /R^2`$ on the left hand side of (59), which then reduces to $`\psi =R^2/C[\mathrm{arg}\mathrm{tanh}(\beta yB)+\gamma y+\alpha ^{}]`$. This corresponds indeed to arbitrarily small values of $`\psi `$ (small energy) as $`C\mathrm{}`$.
When the particular energy value $`E=R^2\beta ^2/24`$ is reached, a uniform PV is possible, with $`\psi /R^2=\beta y`$. Then PV mixing is complete, which clearly maximizes the mixing entropy. In this case, $`\gamma =C\beta `$, so that $`\gamma y`$ cancels the term $`C\psi /R^2`$ in (59). Physically, the uniform westward zonal velocity,
$$v_m=R^2\beta $$
(73)
tilts the free surface with uniform slope by the geostrophic balance, and the corresponding topographic beta-effect exactly balances the imposed beta-effect.
For a still higher energy, a first possibility is that again
$$\gamma =C\beta $$
(74)
so that the beta effect exactly balanced by the $`\gamma y`$ term in the jet equation of previous subsection. This cancellation is directly obtained in the general Gibbs state equations (55) and (56). Indeed the modified stream function $`\psi ^{}=\psi +R^2\beta y`$ satisfies the same equations as in the doubly-periodic case. Therefore in the limit of small $`R`$, the Gibbs states are made of subdomains with uniform $`\psi ^{}`$ (uniform PV), separated by straight zonal jets or circular vortices.
However jets or a vortex persist in this sea of uniform PV due to the constraint of energy conservation. The vortex moves westward at the same velocity $`v_m`$, according to (58) , so they are just entrained by the background flow, without relative propagation (this can be physically understood by the cancellation of the beta-effect).
The selection of the subdomain areas and PV values is given like in the periodic case of section 2, just replacing $`\psi `$ by $`\psi ^{}=\psi \beta R^2y`$. Therefore we get again probabilities $`p_{\pm 1}=(1\pm u)`$ in the two subdomains with respective areas $`A_{\pm 1}`$ given by (31), and stream function,
$$\psi _{\pm 1}=R^2(B\pm u)R^2\beta y$$
(75)
From this relation, we can calculate the energy $`E=\frac{1}{2}(\psi _1^2A_1+\psi _1^2A_1)/R^2`$, so the energy condition (32) then becomes
$$E=\frac{R^2}{2}(u^2B^2)+\frac{R^2\beta ^2}{24}$$
(76)
Therefore these solutions with canceled beta effect can be obtained only beyond a minimum energy $`R^2\beta ^2/24`$, corresponding to the potential energy of the surface tilting associated with the drift velocity $`v_m`$. Then the excess energy will control the organization in two uniform PV areas.
The shape of these subdomains can be obtained again by minimization of the jet free energy. However, unlike in the periodic case, jets occur at the boundaries $`y=\pm \frac{1}{2}`$ as well as at subdomain interfaces. Indeed, such boundary jets are in general necessary to satisfy $`𝐯=\mathrm{𝟎}`$, or equivalently that the stream function $`\psi _b`$ must be equal at the two boundaries $`y=\pm \frac{1}{2}`$. In particular, the solutions (75) necessarily involve a stream function difference (or mass flux) $`R^2\beta `$ associated with the drift velocity $`v_m`$. This stream function difference must be compensated by boundary eastward jets with opposite total mass flux. We show in Appendix B that for two PV levels with similar initial areas a single eastward jet, separating two regions of uniform PV and weak westward drift, is the selected state (instead of two opposite jets in the periodic case). In the case of a strong PV level with a small initial area, the system organizes in a circular vortex like in the periodic case. In the limit $`u|B|`$, as one of the areas $`A_{\pm 1}`$ goes to zero, the jet approximation falls down. The corresponding analysis of axisymmetric vortices and of the linear approximation for the Gibbs states, as performed in section 2, is still valid here.
Up to now we have ignored the constraint of the momentum $`M`$ (54). This constraint imposes the latitude $`y_0`$ of the equilibrium structure ( a circular patch or a zonal band with uniform PV ). For instance in the case $`B>0`$, for which $`A_1>A_1`$ (as seen from (31)), we define $`y_0_{A_1}yd^2𝐫/A_1`$. Then
$$M_Dyqd^2𝐫=_{A_1}y(B+u)d^2𝐫+_{A_1}y(Bu)d^2𝐫=2uy_0A_1$$
(77)
We thus deduce the latitudinal position of the equilibrium structure :
$$y_0=\frac{M}{2uA_1}$$
(78)
In the case of a single eastward jet, the subdomain position has been already fixed by the area ($`y_0=A_1/2`$). Then the only possibility to satisfy a moment $`M`$ different from $`uA_1^2`$ is that the jet oscillates in latitude with some amplitude $`\mathrm{\Lambda }`$ (then $`MuA_1^2u\mathrm{\Lambda }^2`$ ). This is possible if $`\gamma \beta C`$ according to (71), which becomes
$$\frac{1}{r}=b(yy_0)$$
(79)
where $`b(u(\gamma +\beta C))/(2Ce(u)R)<0`$ and $`y_0\alpha _1/b`$. This equation clearly leads to a jet oscillating around the mean latitude $`y_0`$ ( as the curvature $`r`$ is positive for $`y<y_0`$ and negative for $`y>y_0`$ ; recall that the curvature is by definition positive when positive PV is in the inner part of the jet ). Note that this oscillation propagates eastward at speed $`R^2\frac{\gamma }{C}`$ ( given by (58) ). Since $`b<0`$, $`\frac{\gamma }{C}>\beta `$, this speed is eastward with respect to the background drift $`v_m`$ (73).
#### 3.3.2 Quadratic sublayer topography
As explained in section (3.3), in the limit of small Rossby deformation radius $`R`$, the Gibbs state equation has solutions consisting of a vortex bounded by a strong jet on the scale of $`R`$. This corresponds to the case of an initial patch with strong PV and small area ( the asymmetry parameter $`B`$ is sufficiently large ) with an energy sufficiently strong to get a structure of closed jet ( see figure 7). In the presence of a moderate topography $`h(y)`$, this internal jet is no more circular but its radius of curvature $`r<<R`$ depends on $`y`$ according to (71). We have seen in previous subsection that a linear topography $`h(y)=\beta y`$ leads to jets oscillating ( or to circular jets when $`\gamma =C\beta `$ ). We shall discuss here how a quadratic term in $`h(y)`$ modifies the shape of closed jets. We therefore assume a topography $`h(y)`$ of the form :
$$h(y)ay^2+by$$
(80)
This corresponds to a uniform deep zonal shear, with velocity $`v_d=R^2d(h\beta y)/dy=2aR^2y+b\beta `$. We focus our attention on vortex solutions, seeking close curves solutions of equation (71). The vortices will be typically oval shaped as the ones seen on Jupiter. We then study how this shape ( for instance the ratio of the great axis of the oval to the small one ) depends on the topography ( sublayer flow ) and on the jet parameters. Application of these results to Great Red Spot observations will be discussed in next section.
To make equation (71) more explicit, let $`s`$ be a curvilinear parameterization of our curve, $`𝐓(s)`$ the tangent unit vector to the curve and $`\theta (s)`$ the angular function of the curve defined by $`𝐓(s)=(\mathrm{cos}\theta (s),\mathrm{sin}\theta (s))`$ for any $`s`$. Then the radius of curvature $`r`$ of the curve is linked to $`\theta (s)`$ by $`1/r=d\theta /ds`$ and (71) yields the differential equations :
$`\{\begin{array}{cc}\frac{d\theta }{dS}\hfill & =dY^2+1\hfill \\ \frac{dY}{dS}\hfill & =\mathrm{sin}\theta (S)\hfill \end{array}`$ (81)
and $`{\displaystyle \frac{dX}{dS}}=\mathrm{cos}\theta (S)`$
with $`c^{}X=x`$, $`c^{}Y=yy_0`$, $`c^{}S=s`$ ; where $`d=(e^2(u)R^2a)/(ϵ\alpha _1^3u^2)`$, $`y_0e(u)R(C_0b\gamma )/(ϵ\alpha _1C_0u)`$. The space coordinates $`X`$, $`Y`$ and $`S`$ are here non dimensional and have been obtained by dividing the real coordinates by the scale $`c^{}e(u)R/(ϵ\alpha _1u)`$. Note that as explained in section (3.3) $`ϵ\alpha _1>0`$, so that $`c^{}>0`$. We further assume that $`a>0`$, so that $`d>0`$.
We first note that the two variables $`\theta `$ and $`Y`$ are independent of $`X`$. We will therefore consider the system formed by the two first differential equations (81). It is easily verified that this system is Hamiltonian, with $`\theta `$ and $`Y`$ the two conjugate variables and
$$H\mathrm{cos}\theta d\frac{Y^3}{3}+Y$$
(82)
the Hamiltonian. Thus $`H`$ is constant on the solution curves. We look for vortex solutions of our problem (81 and 81). Thus we require $`\theta `$ to be a monotonic function of $`S`$. Moreover the curves must close, that is $`X`$ and $`Y`$ must be periodic. For symmetry reasons, it is easily verified that the solutions of (81,81) with initial conditions $`\theta (0)=\frac{\pi }{2}`$, $`Y(0)=0`$ ( $`H=0`$ ) and some $`X(0)`$ are periodic. We prove in appendix (C) that these initial conditions are the only ones leading to closed curves. We also prove that the solutions of (81 and 81) when $`d>d_{max}\frac{4}{9}`$ does not define $`\theta `$ as a monotonic function of $`S`$. They contain double points and thus are not possible solution for our problem. Once given these initial conditions, we can easily prove that the structure has both a zonal symmetry axis and a latitudinal one passing through $`y_0`$.
To study the shape of the jets, we numerically solve equations (81,81) with initial conditions : $`\theta (0)=\frac{\pi }{2},y(0)=0`$. We obtain closed curves with oval shapes, as shown in figure 10. In figure 11 we have represented the width, the length and the aspect ratio of these vortices versus the parameter $`d`$. When $`d`$ tends to $`\frac{4}{9}`$, the vortex width tends to a maximum value : $`w_{max}=\frac{3}{2}`$ whereas the length diverges. In this limit, the vortices are thus very elongated.
## 4 Application to the Jupiter’s Great Red Spot and Oval BC
In previous sections, we have found maximum entropy states with the following properties :
* The fluid domain is partitioned in two subdomains with weak velocity, separated by jets whose width scales as the Rossby deformation radius. A strong initial PV level occupying a small area mixes in a subdomain with the form of a vortex bounded by an annular jet.
* In the presence of a parabolic topography $`h(y)`$ (due to a sublayer zonal flow with uniform shear), outside the jet, exists a zonal flow (69) with uniform shear. The velocity at the latitude of the vortex center vanishes in the reference frame of the vortex.
* The curvature of the jet is linked to the topography by (71). For the parabolic topography, solutions are oval shaped vortices, symmetric in latitude and longitude.
These properties of our solutions are the main qualitative properties of the Jovian vortices. Moreover, would this description be correct, it would predict that the topography has an extremum at the center of the vortex.
Dowling and Ingersoll (1989) derive the bottom topography, using the GRS and Oval BC velocity fields obtained from cloud motion. They analyze the results in the frame of a 1-1/2 shallow water model (SW), with an active shallow layer floating on a much deeper layer. This deep layer is in steady zonal motion which acts like a topography $`h_2`$. The SW topography $`h_2`$ is defined by $`fv_d=1/R_l(\lambda )(gh_2)/\lambda `$ where $`v_d`$ is the deep layer flow and $`R_l(\lambda )`$ is the latitudinal radius of curvature of Jupiter. Dowling and Ingersoll (1989) have deduced this SW topography $`h_2`$ by assuming the conservation of the shallow water potential vorticity $`(\omega +f)/h_1`$ ( $`h_1`$ is the upper layer thickness ) along the streamlines of the observed steady vortex flow. The vorticity $`\omega `$ is deduced from the measured velocity field, and the planetary vorticity $`f`$ is known, so that the variation of $`h`$ along each streamline is deduced. The pressure field is then obtained from the Bernoulli relation and the hydrostatic balance, leading to the field $`h_2`$. The result depends on the radius of deformation of Rossby ( $`R^{}=gh_0/f_0`$, where $`h_0`$ is the mean upper layer height and $`f_0`$ the mean Coriolis parameter ), a free parameter in this analysis. Three test values have been chosen, $`R_1^{}=1700,R_2^{}=2200`$ and $`R_3^{}=2600`$ km for the GRS and $`R_1^{}=1100,R_2^{}=1600`$ and $`R_3^{}=2000`$ km for the Oval BC ( we denote by a star superscript the physical parameters, to distinguish them from the non-dimensional quantities used earlier ). The height $`h_2`$ under each vortex has been found to depend only on latitude, and has been fitted as a quartic the planetographic latitude $`\lambda `$:
$$gh_2=A_0+A_1\lambda +A_2\lambda ^2+A_3\lambda ^3+A_4\lambda ^4.$$
(83)
The values obtained for the coefficients $`A_i`$ in the vortex reference frame, for each of the vortices and for each of the values $`R_1^{},R_2^{},R_3^{}`$ are reported in table 1 of Dowling and Ingersoll (1989).
Our model is the QG limit of this shallow water system. Starting from (SW) equations, we can derive the QG equations (1) by assuming the geostrophic balance and weak free surface deformation in comparison with the mean layer thickness $`h_0`$. The validity of this QG approximation has been discussed by Dowling and Ingersoll (1989) and was found reasonably good as a first approach, although not accurate. We furthermore use the beta-plane approximation, linearizing the planetary vorticity around a reference latitude $`\lambda _0`$ ( $`\lambda _0`$ is taken to be $`23^0`$ for the GRS and $`33.5^0`$ for the Oval BC). Therefore we write $`f=f_0+\beta y`$, with $`f_0=2\mathrm{\Omega }sin\lambda _0`$ and $`\beta 2\mathrm{\Omega }\mathrm{cos}\lambda _0/r_z(\lambda _0)`$ ( $`\mathrm{\Omega }`$ is the planetary angular speed of rotation : $`2\pi /\mathrm{\Omega }=9`$ h 55 mn 29.7 s and $`r_z(\lambda _0)`$ is the zonal planetary radius, which slightly depends on the latitude $`\lambda _0`$, due to the ellipsoidal planetary shape, see formula (4) of Dowling and Ingersoll (1989) ) . We then obtain the QG potential vorticity (2) with the QG topography $`h^{}(y^{})`$ linked to the SW topography (83) by:
$$h^{}(y^{})=\frac{gh_2}{f_0R^2}+\beta y^{}$$
(84)
We have computed the QG topography (84) using results of Dowling and Ingersoll (1989) for the SW topography (83) for the three values of the Rossby deformation radius $`R_1^{},R_2^{}`$ and $`R_3^{}`$, for the GRS and for the Oval BC. The result in figure 12 shows that, for both the GRS and the Oval BC, the QG topography has an extremum at a latitude which is nearly the center of the vortex. As far as we know, this fact has not been noticed in the literature. This result is in agreement with the predictions of our model. We note moreover that the two extrema of the topography are minima, thus our model predicts anticyclonic shears around the GRS and the Oval BC, as observed. Figure 84 shows a comparison of the QG topography derived from Dowling results with a quadratic approximation, in the case $`R^{}=R_2^{}=2200`$ km. This shows that the quadratic approximation $`h^{}(y^{})=a^{}y^2`$ is a good approximation on the latitudinal extension of the GRS. This also provides values of the parameter $`a^{}`$ (80) : $`a^{}=\mathrm{9.2\; 10}^{13}`$, $`a^{}=\mathrm{7.2\; 10}^{13}`$ and $`a^{}=\mathrm{6.4\; 10}^{13}`$ km<sup>-2</sup> s<sup>-1</sup> for $`R^{}=R_1^{},R_2^{}`$ and $`R_3^{}`$ respectively.
Let us deduce the corresponding non-dimensional parameters. First, the PV levels were normalized by (9), so our time unit $`T^{}`$ will depend on the real PV level difference : $`(a_1^{}a_1^{})/2\frac{1}{T^{}}`$. The other parameters are the Rossby deformation radius $`R^{}`$, the segregation parameter $`u`$ and the topography coefficient $`a^{}`$ (80).
We will consider $`R^{}`$ as a free parameter and use the following data from GRS observation :
* The jet width $`l^{}`$. Let us define the width of the jet $`l^{}`$ as the width on which the jet velocity is greater than one half of the maximum velocity. We use velocity measurement within the GRS of Mitchell et al (1981). They have used small clouds as tracers to measure velocities, and have observed that the velocity is nearly tangential to ellipses. Using a grid of concentric ellipses of constant eccentricity, the velocities have been plotted with respect to the semi-major axe $`𝒜`$ of the ellipse on which the measurement point lies. Results have been fitted by a quartic in $`𝒜`$ and may be seen in figure 13. Using these results, assuming that the velocity profile in the jet is symmetric with respect to its maximum, we choose as jet width $`l^{}=\mathrm{5.6\; 10}^3`$ km. In our model, the normalized jet width $`l=l^{}/R^{}`$ can be computed from the jet equation (34) as shown in figure 6(a). This determines the parameter $`u`$ from the parameter $`R^{}`$. The corresponding theoretical jet velocity can be compared to observations in figure 13, for $`R^{}=2500`$ km (the shape is not very sensitive to this parameter). The computed results for $`u`$ versus $`R^{}`$ is shown in figure 15.
* The maximum jet velocity $`v_{max}^{}`$. We will use the value $`v_{max}^{}=110`$ ms<sup>-1</sup> ( Mitchell et al 1981). Using (35) and giving real dimension gives :
$$v_{max}^{}=\frac{R^{}}{T^{}}\frac{d\varphi }{d\tau }|_{max}(u).$$
(85)
$`d\varphi /d\tau |_{max}(u)`$ has been obtained by solving the non-dimensional jet equation (34), and is shown in figure 6(b). This now determine $`T^{}`$ from $`R^{}`$.
* The velocity shear surrounding the vortex. The ambient zonal shear measured at the latitude of the GRS from Limaye et al. (1986) is $`\sigma ^{}=1.5e5`$ s<sup>-1</sup>. Using (69) in its dimensional form, for a quadratic topography gives :
$$a^{}=\frac{\sigma ^{}}{2R^{}}\left(1\frac{C_0}{\mathrm{cosh}^2(C_0u)}\right)$$
(86)
This permits to compute $`a^{}`$ as a function of $`R^{}`$ (since $`u`$ has been determined, as well as $`C_0`$, related to $`u`$ by (33) ).
The computed results for $`a^{}`$ and $`T^{}`$ versus $`R^{}`$ are shown in figure 14. This shows that our determination of the topography is in agreement with the QG topography deduced from the shallow water model of Dowling and Ingersoll (1989) within a factor of two. The corresponding PV level difference $`a_1^{}a_1^{}`$ is comparable to the planetary vorticity $`f_0`$ at the latitude of the center of the GRS when $`R2400`$ km. For this value of $`R^{}`$, figure 15 shows that $`u`$ is very close to 1. Furthermore, as the GRS area is very small compared to the global area of a latitudinal band centered around the GRS, the non dimensional area occupied by the positive PV is very close to 1. Using this area expression (31) we conclude that $`B`$ is very close to 1. Using the definition of $`B`$ (10), we conclude that $`a_1^{}0`$ and $`a_1^{}f_0`$. As discussed below a forcing mechanism by convective plumes incoming from the sublayer is expected to yield this result.
The shape of the jet depends on the parameter $`d=\frac{e^2(u)}{ϵ\alpha _1^3u^2}\frac{a^{}R^^2}{a_1^{}a_1^{}}`$ and on the length scale $`c^{{}_{}{}^{}}\frac{e(u)R^{}}{ϵ\alpha _1u}`$. We can determine these two parameters from the previously determined values of $`a^{}`$,$`T^{}`$ and $`u`$, and from the observed half width of the Great Red Spot : $`y_{max}^{}=4900`$km. This permits to calculate $`c^{}`$, $`\alpha _1`$ and $`d`$ versus $`R^{}`$. Figure 16 show $`d`$ versus $`R^{}`$. The dot line represents the critical value $`d=\frac{4}{9}`$ beside which a vortex solution exists. The ratio of the length to the width of the GRS is approximatively $`2`$, which would correspond to $`d=0.441`$ ( figure 11 ) ; this is very close to the critical value $`\frac{4}{9}`$. From figure 16, our model predicts that the Rossby deformation radius is $`R^{}=1800`$ km. Figure 10 shows the actual shape of the vortex for $`d=0.441`$.
However, in the jet shape analyze, to obtain (81,81) we have supposed $`a^{}L_V^^2<<(a_1^{}a_1^{})/2`$ where $`L_V^{}`$ is the maximal latitudinal extension of the vortex ( the topography part of the PV remains negligible with respect to the PV ). For the value of $`R^{}`$ calculated above, we find $`a^{}y_{max}^{}{}_{}{}^{2}=3.910^5s^1`$ whereas $`a_1^{}a_1^{}=1.310^4s^1`$. We are thus at the limit of validity of our assumption.
Now we can reverse the procedure and propose a predictive model of the Great Red Spot. Assume a steady deep zonal flow with uniform shear, $`v_d=2a^{}R^2(y^{}y_I^{})`$, vanishing at an origin $`y_I`$ depending on our reference frame (which we shall choose in order to cancel the vortex drift). This flow is presumably generated by deep thermal convection but we are concerned here only with the dynamics of the upper layer, assumed stably stratified (due to cooling by radiative effects). We model this stratified upper layer as a shallow layer with radius of deformation $`R^{}2400`$ km. This layer is submitted to a total beta effect or ’topography’ $`h(y)=ay^2(\beta +2ay_I)y`$.
Assume that PV spots with value $`f_0`$, occupying small area proportion are randomly generated in this layer. This would be the result of intense incoming thermal plumes, as recently discussed by Ingersoll et al (2000), : conservation of the absolute angular momentum during the radial expansion leads to strong decrease of the local absolute vorticity, which comes close to zero. This means that in the planetary reference frame, a local vorticity patch with value $`f_0`$ is created. The opposite vorticity is globally created by the subducting flow, but it is close to 0 due the much larger area. This gives our time unit $`(2f_0)^1`$ and $`B=12A`$.
The outcome of random PV mixing with the constraint of the conservation laws is then a zonal velocity (69) in the observed upper layer and a vortex with area (31) with velocity profile shown by the dot curve in fig 14. The vortex moves with the upper velocity at $`y=0`$, it drifts with respect to the deep layer at velocity (73) so that the beta effect is suppressed. The shape is an oval symmetric in $`x`$ and $`y`$, with aspect ratio computed from the shape parameter $`d`$ ( figure 11 ).
Note that a slightly larger area , or stronger energy could lead to very elongated vortices. Then argument of free energy minimization show that this would lead to a single eastward zonal jet. This may explain the jet observed in the Northern hemisphere of Jupiter at the same latitude as the GRS.
For smaller Jovian vortices such as the Oval BC, the size of the vortices is comparable with the Rossby deformation radius. Thus such vortices may be described as done for axisymmetric vortices ( section 2.3.4 ). This explain why such vortices do not have a quiescent core as the GRS.
Let us describe dark brown cyclonic spots ( ’Barges’ ) at $`14`$ N on Jupiter. Their first interest is to stress that cyclonic vortices embedded in cyclonic shear exist on Jupiter. Let us go further. The greater of these barges is studied from Voyager observations by Hatzes at al (1981). The meridional velocities measured at the latitude of the center of the barge ( Hatzes et al (1981), figure 7 ) show a boundary jet organization around the perimeter of the barge ( $`v_{max}=25`$ ms<sup>-1</sup> ), see figure 1(b). The surrounding shear is such that the shear velocity at the maximum latitude of the barge is the same as the maximum jet velocity. Thus our approximation $`aL_V^2<<1`$ is not good. We can however explain the elongated shape, similar to figure 10(b) obtained for $`d`$ very close to $`d_{max}`$.
We conclude that the Gibbs state equation (56) derived from maximization of entropy of the QG model (1) is in the limit of small Rossby deformation radius a model that explains the main qualitative features of Jovian vortices. The statistical mechanics itself explains the organization of a turbulent flow in coherent structures.
## 5 Conclusion
Our first result is to provide a general explanation for the emergence and robustness of intense jets in atmospheric or oceanic turbulent flows. In the absence of topography or beta-effect turbulence mixes potential vorticity in subdomains, and such jets occur at the interface of these subdomains, with a width of the order of the deformation radius. From a thermodynamic point of view, this is like coexistence of two phases. Indeed the vortex interaction becomes short ranged in the limit of small deformation radius, and statistical mechanics leads to a thermodynamic equilibrium between two 'phases', with different concentrations of the 2 Potential Vorticity levels. Another approach leading to the same result is to consider the general partial differential equation (18) characterizing the equilibrium states. This equation reduces to the algebraic equation (20) in the limit of small deformation radius. The two uniform subdomains correspond to two solutions $`\psi _1`$ and $`\psi _1`$ of this equation. At the interface of these subdomains, the general pde reduces to the equation (36), whose solution determines the jet profile. In addition, a solvability condition of this equation confirms the relation of equilibrium between the two 'phases', which was obtained in the thermodynamic approach.
All our results have been obtained for a 2 potential vorticity level case, but cases with more levels would lead to qualitatively similar results, although the quantitative analysis would be more involved due to the additional parameters.
In the presence of beta effect or topography, for low energy, purely zonal flows with gentle variation in latitude are obtained. A critical energy is the energy of the state where the zonal flow just compensates the beta-effect. For this state the PV is strictly uniform in the whole domain. For greater energy, two well mixed domain separated by jets appears, as in the without topography case. However the PV is no more strictly uniform in the well mixed subdomain : a zonal flow exists. Also the jet curvature depends on latitude. With ordinary beta effect this yields an intense eastward jet, purely zonal or wavy depending on the constraint on the momentum $`M`$. With the quadratic beta-effect generated by a deep shear, this can produce an oval shaped vortex. This vortex then drift so as to compensate the beta-effect. In other words, in the vortex reference frame the equivalent topography $`h(y)`$ admits an extremum, and this is in agreement with the data of Dowling and Ingersoll (1989).
Our quasi analytical approach therefore explains most of the basic features of the Great Red Spot and other Jovian vortices. It can be developed into a more accurate predictive model along the following lines. First the approximation $`R<<r`$ of thin jet is convenient for a qualitative understanding but is only marginally satisfied. However this limitation can be overcomed by numerical determination of the equilibrium state equation (18) by methods like used by Turkington and Whitaker (1995) or using relaxation equations toward the maximum entropy state as described by Robert et Sommeria (1992). Furthermore, extension to the more general shallow water model is desirable, as the Rossby number ( $`0.36`$ where it a maximal ) is not very small. This can be formally achieved ( in preparation ).
Finally the results rely on an assumption of ergodicity, or complete potential vorticity mixing consistent with the constraints on the conservation laws. Various numerical and laboratory experiments in the case of Euler equations ( see e.g. Brands et al. 1989 ) indicate that mixing may not be global but be restricted to active regions. Organization into local vortices, rather that at the scale of the whole domain, is more likely with a small radius of deformation, as vortex interactions leading to coalescence are then screened. This is observed for instance in the numerical computations of Kukharkin Orszag and Yakhot (1995). By contrast, the zonal shear here promote vortex encounters ( as observed in Voyager data ) and we expect a much better relaxation toward the global statistical equilibrium, which involves always a single vortex in a given shear zone ( as it minimizes the interfacial free energy ).
## 6 Acknowledgments
The authors thank R. Robert for collaboration on statistical mechanics approach and for useful comments on the present work.
## Appendix A Determination of the Gibbs state by direct entropy maximization, in the presence of a topography ( beta-effect ) .
In section (3.3), we studied the limit of small Rossby deformation radius in the Gibbs state equation (56) by considering the jet equation (60) and its integrability condition (65). We deduced that the Gibbs states are composed of subdomains in which $`\psi `$ verifies the algebraic equation (59) separated by an interfacial jet whose curvature verifies (71). The aim of this annex is to prove that these results can be obtained by directly maximizing the entropy, adapting the method used in section (2.3.1).
Let us make the following assumptions :
1. In the limit of small Rossby deformation radius, the probability $`p`$ of finding the PV level $`a_1`$ takes two values $`p_{\pm 1}(y)`$, depending only on $`y`$. We are looking for vortex solutions. The vortex shape is described by the length $`l(y)`$ on which the probability $`p`$ takes the value $`p_1(y)`$ ( See figure 17 ).
2. The two subdomains where $`p`$ take the two values $`p_{\pm 1}(y)`$ are separated by a jet. The probabilities $`p_{\pm 1}(y)`$ are supposed to be close to the their values without topography $`p_{\pm 1}=\pm u`$, such that the free energy per unit length of the jet is well approximated by the one calculated without topography (42). If $`L_V`$ denotes the vortex size, $`1/\sqrt{a}`$ a typical length on which topography varies, we will show this approximation to be valid as soon as $`aL_V^2<<1`$.
3. The boundary conditions can be relaxed ; that is no boundary term appears in the variation of the free energy at the order considered here. ( See discussion concerning boundary jets in section (3.3).
Given these hypothesis, the Gibbs states is described by the 3 functions $`p_{\pm 1}(y)`$ and $`l(y)`$. We will determine them by maximizing the entropy $`S`$ (15) under the 3 constraints : energy (53), mass conservation (51) and momentum (54). A necessary condition for a solution to this variational problem is the existence of 3 Lagrange parameters $`C`$,$`\alpha `$ and $`\gamma `$ such that the first variations of the free energy :
$$FS\frac{C_0}{R^2}E+\alpha _0\frac{\psi }{R^2}+\gamma M$$
(87)
vanish. Using (15),(53),(51) and (54) ; (2) and (11) and the above hypothesis ( see figure 17 ), a direct calculation shows that the free energy (87) is up to a constant :
$$\begin{array}{cc}\hfill F& =_{y_{min}}^{y_{max}}\left[f(p_1(y),y)(1l(y))+f(p_1(y),y)l(y)\right]𝑑y+_{\frac{1}{2}}^{y_{min}}f(p_1(y),y)𝑑y+\hfill \\ & _{y_{max}}^{\frac{1}{2}}f(p_1(y),y)𝑑y+LF_{Jet}(u)\hfill \\ \hfill \text{with}f(p,y)& (p\mathrm{log}p+(1p)\mathrm{log}(1p))2C_0\left(p\frac{1}{2}\right)^22\left(C_0B\alpha _0\right)\left(p\frac{1}{2}\right)\hfill \\ & 2\left(C_0h(y)\gamma y\right)\left(p\frac{1}{2}\right)\hfill \end{array}$$
(88)
and where $`L`$ is the jet length, $`F_{Jet}(u)`$ is the jet free energy per unit length (42), calculated without topography.
Considering first variations of the free energy (88) under variations of $`p_1(y)`$ ( resp $`p_1(y)`$ ) proves that $`f/p(p_1(y),y)=0`$ and that $`f/p(p_1(y),y)=0`$. A direct calculation shows that :
$$2\left(p_{\pm 1}\frac{1}{2}\right)=\mathrm{tanh}\left(2C_0\left(p_{\pm 1}\frac{1}{2}\right)+C_0h(y)\gamma y+C_0B\alpha \right)$$
(89)
Using (11) and (2) ; recalling that we neglect the Laplacian term, a straight calculation shows that (89) is equivalent to the algebraic equation (59).
Let us consider now first variations of the free energy (88) under small variations $`\delta l(y)`$ of $`l(y)`$. Using that the length of the jet is given by $`L=2_{y_{min}}^{y_{max}}\sqrt{1+\frac{1}{4}\left(dl/dy\right)^2}𝑑y`$, a straightforward calculation shows that $`\delta L=_{y_{min}}^{y_{max}}\delta l(y)/r𝑑y`$ where $`r`$ is the radius of curvature of the jet. We thus deduce from first variations of the free energy (88) :
$$\frac{F_{Jet}(u)}{r}=f(p_1(y),y)f(p_1(y),y)$$
(90)
Hypothesis (2) : $`aL_V^2<<1`$ permitted us to consider $`F_{Jet}(u)`$ as independent of $`y`$. In accordance with this hypothesis we evaluate $`f(p_1(y),y)f(p_1(y),y)`$ at order zero, with $`p_{\pm 1}=\frac{1}{2}(1\pm u)`$ at this order. We obtain : $`f(p_1(y),y)f(p_1(y),y)=2u(\alpha 1+C_0h(y)\gamma y)`$. Moreover using the free energy per unit length expression (42), (37) and (38), one can show that $`F_{Jet}(u)=2e(u)C_0R`$ where $`e(u)`$ is defined by (70). These two last results show that (90) is equivalent to (71) the expression for the radius of curvature $`r`$ found by the integrability condition for the jet.
## Appendix B Boundary jets in the channel case. Beta-effect or linear sublayer topography.
Let us derive the boundary jet properties in the case of a beta-effect (or linear sublayer topography). For the sake of simplicity, we will treat the case of zero circulation $`\mathrm{\Gamma }=0`$. Let $`\varphi _b^+`$ and $`\varphi _b^{}`$ be the values of $`\varphi `$ on the $`y=\frac{1}{2}`$ and $`y=\frac{1}{2}`$ boundary respectively. The boundary jet satisfies the jet equation (36), but with boundary conditions $`\varphi (\tau =0)=\varphi _b^\pm `$ and $`\varphi (\tau +\mathrm{})=\varphi _1or\varphi _1`$. Thus, using (38), we deduce that :
$$\frac{1}{2}\left(\frac{d\varphi }{d\tau }\right)^2\left(y=\pm \frac{1}{2}\right)+U(\varphi _b^\pm )=U(\varphi _1)=U(\varphi _1)$$
(91)
(where the last equality comes from the integrability condition for the interfacial jet).
We can relate $`d\varphi /d\tau `$ to the derivative $`d\psi /d\zeta `$ normal to the boundary ( denoting the normal coordinate $`\zeta =\pm \frac{1}{2}y`$ ), using (61),
$$\frac{d\varphi }{d\tau }\left(y=\pm \frac{1}{2}\right)=\frac{1}{R}\frac{d\psi }{d\zeta }\left(y=\pm \frac{1}{2}\right)R\frac{\gamma }{C}$$
(92)
Furthermore the condition $`\mathrm{\Gamma }=0`$ imposes
$$\mathrm{\Gamma }=_{y=\frac{1}{2}}\frac{d\psi }{dy}𝑑x+_{y=\frac{1}{2}}\frac{d\psi }{dy}𝑑x=\frac{d\psi }{d\zeta }\left(y=\frac{1}{2}\right)+\frac{d\psi }{d\zeta }\left(y=\frac{1}{2}\right)=0,$$
so that $`\frac{d\psi }{d\zeta }(y=\frac{1}{2})=\frac{d\psi }{d\zeta }(y=\frac{1}{2})`$. Then (91) becomes :
$$U(\varphi _b^\pm )=U(\varphi _1)\frac{1}{2}\left(\frac{1}{R}\frac{d\psi }{d\zeta }\left(y=\frac{1}{2}\right)R\frac{\gamma }{C}\right)^2$$
(93)
so that
$$U(\varphi _b^+)=U(\varphi _b^{})$$
(94)
Note that the two values $`\varphi _b^\pm `$ cannot be equal. They must indeed satisfy, from (61) and the condition of zero mass flux ($`\psi ^+=\psi `$),
$$\varphi _b^+\varphi _b^{}=\gamma /C$$
(95)
To solve these two equations (94) and (95), let us have a look at the potential $`U`$ in figure 19. We have to compare $`\varphi _b^+\varphi _b^{}`$ and $`\varphi _1\varphi _1`$. We thus distinguish two cases. Using (95) and $`\varphi _{\pm 1}=\pm u`$, we calculate : $`(\varphi _b^+\varphi _b^{})/(\varphi _1\varphi _1)=\gamma /(2Cu)`$. We recall that $`u`$ is of order unity.
* Low beta-effect case : $`\gamma /C<2u`$.
For an x-independent statistical equilibrium, one possibility is a zonal band with $`\varphi _1`$ inside the band and $`\varphi _1`$ near both boundaries, with $`\varphi _{b1}^+`$ and $`\varphi _{b1}^{}`$ at the boundaries. The symmetric solution, with $`\varphi _{b1}^+`$ and $`\varphi _{b1}^{}`$ at the boundaries and $`\varphi _1`$ near the boundaries, has the same free energy ( due to the symmetry of the potential $`U`$ ). However the solution which maximizes the total free energy of the jets corresponds to $`\varphi _1`$ in the lower part of the domain, and $`\varphi _1`$ in the upper part, with $`\varphi _{b0}^+`$ at the upper boundary ( $`y=\frac{1}{2}`$ ) and $`\varphi _{b0}^{}`$ at the lower boundary ( $`y=\frac{1}{2}`$ ), and a single eastward interior jet.
* High beta-effect case : $`\beta >2u`$
In this case, we note that $`\varphi _b^+\varphi _b^{}>\varphi _1\varphi _1`$. Then the two equations (95) and (94) determine a unique solution for $`(\varphi _b^+,\varphi _b^{})`$ ( see figure 19 ). We are then necessarily in the case with $`\varphi =\varphi _1`$ near the boundaries $`y=\frac{1}{2}`$ and $`\varphi =\varphi _1`$ near the other one. It again involves a single eastward internal jet. This jet can oscillate in latitude due to the momentum constraint ( according to (79) ).
Once $`\varphi _b^\pm `$ are fixed, we use equation (61) on the two boundaries $`y=\pm \frac{1}{2}`$ to conclude that $`\alpha =\frac{1}{2}(\varphi _b^++\varphi _b^{})`$. Then using the integrability condition (66) we calculate :
$$\psi _b=R^2\left(B+\frac{\varphi _b^++\varphi _b^{}}{2}\right)$$
(96)
This closes the determination of our parameters.
## Appendix C Boundary condition for the equation (81,81) of the oval-shaped vortices boundary
We want to see whether equations (81,81) with $`d>0`$, defining the curve formed by the jet, admit periodic solutions both in $`x`$ and $`y`$ corresponding to vortices, or not. As stressed in section (3.3.2), equations (81) derive from the Hamiltonian (82) ; $`y`$ and $`\theta `$ are the two conjugated variables. Let us study the phase portrait of $`H`$. For $`\theta `$ in $`[0,2\pi [`$, there are 4 critical points : $`P_1=(0,1/d),P_2=(0,1/d),P_3=(\pi ,1/d),P_4=(\pi ,1/d)`$. By linearization around these fixed points, one easily prove that $`P_1`$ and $`P_3`$ are stable fixed points whereas $`P_2`$ and $`P_4`$ are hyperbolic fixed points. This permits to draw the phase portrait : figure 18. Using $`H`$ (82), we obtain that the unstable manifolds are given by $`12/(3\sqrt{d})=H`$ and $`1+2/(3\sqrt{d})=H`$ respectively. The parameter $`d`$ governs a transition of the phase space structure. This transition occurs when the two unstable manifolds merge ; this permits to compute the transition value for $`d`$ : $`d=\frac{4}{9}`$. We are looking for vortices solutions of (81,81). We recall that $`\theta `$ is the angle with the $`x`$ axe. We thus impose $`\theta `$ to be a monotonic function of $`s`$ on a trajectory. Thus areas $`c`$ of the phase portrait (figure 18 ) are forbidden ; they would correspond to $`\theta `$ varying in a finite interval. We won't address their analyze there, but oscillating jet solutions can be found in these areas. Area $`b2`$ : it can be shown that, as on such trajectories, as $`\theta `$ is not strictly increasing, double points exists, thus forbidding area $`b2`$. Conversely areas $`a`$ and $`b1`$ are admissible trajectories, giving $`y`$ as periodic functions of $`\theta `$.We have to had the condition that $`x`$ (81) must also be a periodic function of $`\theta `$ for the curve to close. Let us denote $`\mathrm{\Delta }x`$ the $`x`$ variation when $`\theta `$ runs in $`[0,2\pi ]`$. We thus impose the condition $`\mathrm{\Delta }x=0`$. Using (81) and (81) we calculate :
$$\mathrm{\Delta }x=_0^L\mathrm{cos}\theta ds=_0^{2\pi }\frac{\mathrm{cos}\theta d\theta }{dy^2(\theta )+1}=_{\frac{\pi }{2}}^{\frac{\pi }{2}}\mathrm{cos}\theta \frac{d[y^2(\theta )y^2(\theta \pi )]}{(dy^2(\theta )+1)(dy^2(\theta \pi )+1)}d\theta =0$$
(97)
( The last expression is obtained rewriting the integral as a sum on $`[\frac{\pi }{2},\frac{\pi }{2}]`$ plus a sum on $`[\frac{\pi }{2},\frac{3\pi }{2}]`$ ; and performing a variable change ). Let us study the sign of $`y^2(\theta )y^2(\theta \pi )`$. Using (82) we deduce that $`d\frac{y^3(\theta )}{3}+y(\theta )=H\mathrm{cos}\theta `$ and $`d\frac{y^3(\theta \pi )}{3}+y(\theta \pi )=H+\mathrm{cos}\theta `$. From these two relations we conclude that $`y(\theta \pi )=y(\theta )`$ implies $`\mathrm{cos}(\theta )=0`$ and that $`y(\theta \pi )=y(\theta )`$ implies $`H=0`$. Thus if $`H0`$ ; $`y^2(\theta )y^2(\theta \pi )`$ does not change sign on $`[\frac{\pi }{2},\frac{\pi }{2}]`$. Moreover, on the areas a) and b1) $`(dy^2(\theta )+1)`$ does not change sign. Thus if $`H0`$, the argument of the last integral of (97) does not change sign and $`\mathrm{\Delta }x`$ can not be zero. We thus conclude that the only solution where $`x`$ is a periodic function of $`\theta `$ is the solution corresponding to $`H=0`$. This solution is the one obtained from (81,81) with initial conditions $`y(0)=0`$ and $`\theta (0)=\frac{\pi }{2}`$. As we have previously excluded the area $`b2`$ ( when $`d>\frac{4}{9}`$ ) we conclude that no vortex solution exists when $`d>\frac{4}{9}`$. We conclude that equations (81,81) with $`d>0`$, defining the curve formed by the jet, admit one periodic solution both in $`x`$ and $`y`$ corresponding to vortices, only when $`d<d_{max}=\frac{4}{9}`$. This solution corresponds to $`H=0`$ (82). The vortex then admits a latitudinal and a zonal axes of symmetry.
Figure Captions.
Figure 1 : Annular jets observed in the atmosphere of Jupiter. a) Velocity field in the Great Red Spot of Jupiter (20<sup>0</sup> South), from Dowling and Ingersoll (1989). b) Velocity field in the cyclonic Barge of Jupiter (14<sup>0</sup> North) from Hatzes et al (1981)
Figure 2 : (a) Graphical representation of the algebraic equation (20), with the rescaled variable $`\varphi \alpha /C+\psi /R^2`$. The three solutions are at the intersection of the curve (left-hand side) and straight line (right-hand side). Here the integrability condition $`\alpha =C_0B`$ for the differential equation (36)
Figure 3 : The free energy density $`f(p)`$ (27) versus the probability $`p`$. For $`C_0>1`$ and $`(C_0B\alpha _0)`$ small enough $`f(p)`$ has two local minima and one local maximum, allowing to obtain two values $`p_{\pm 1}`$ in the maximization of entropy under constraints.
Figure 4 : The parameter $`u`$ versus the Lagrange parameter $`C_0`$, as the solution of (33).
Figure 5 : Typical stream function profile in a jet ($`u`$ = 0.75 ) versus the transverse coordinate $`\tau =\zeta /R`$ ( Top ) and corresponding velocity profile ( Bottom ).
Figure 6 : Jet properties versus the segregation parameter $`u`$. a) Jet width, defined as the width of the region with velocity greater than half the maximum jet velocity. b) Maximum velocity $`\left(d\varphi /d\tau \right)_{max}`$ and jet kinetic energy $`e(u)`$ ( dotted line).
Figure 7 : Phase diagram of the Gibbs states versus the energy $`E`$ and the asymmetry parameter B. The outer line is the maximum energy achievable for a fixed B : $`E=\frac{R^2}{2}(1B^2)`$. The frontiers line between the straight jets and the circular jets corresponds to $`A_1=1/\pi `$ or $`A_1=1/\pi `$. it as been calculated using (31) and (32) : $`E=R^2B^2(2\pi 2)/(\pi 2)^2`$. The dot line represents the frontiers between axisymmetric vortices and the circular jets. We define it as the energy value for which the circular vortex area $`A_1`$ or $`A_1`$ (31) is equal to $`(2l)^2`$, where $`l`$ is the typical jets width ( figure 6 ). Such a line depends on the numerical value of R the ratio of the Rossby deformation radius to the domain scale. It has been here numerically calculated for R = 0.03.
Figure 8 : (a) Graphical representation of the algebraic equation (20), with the rescaled variable $`\varphi \alpha /C+\psi /R^2`$, like in figure 2, but in the case of a Gibbs state with an axisymmetric vortex ($`\mathrm{\Delta }C>0`$). Then the rhs hatched area is greater than the lhs one.
(b) The corresponding potential $`U(\varphi )`$, given by (49), is asymmetric to compensate for the friction term in equation (50).
Figure 9 : Various axisymmetric stream-function profiles for decreasing $`\mathrm{\Delta }C`$ ( $`\mathrm{\Delta }C=[\mathrm{0.90.60.30.10.050.030.01}]`$ and $`B=0.75`$.
Figure 10 : a ) Typical sub-domain shape with a topography $`h(y)=ay^2`$. The parameter $`d`$ has been chosen such that the ratio of the length on the width be 2 ; as on Jupiter's GRS. b ) Typical sub-domain shape with a topography $`h(y)=ay^2`$ when the parameter $`d`$ id very close to its maximum value $`d=\frac{4}{9}`$. The shape is then very elongated, with latitudinal boundaries quasi parallel, as for instance the Jovian cyclonic vortices ( 'Barges' ) described by ( Hatzes et al 1981 ).
Figure 11 : a ) Sub-domain non dimensional length and width versus the parameter $`d`$ ( topography $`h(y)=ay^2`$ ). b ) Sub-domain aspect ratio versus the parameter $`d`$.
Figure 12 : QG topography ( units $`s^1`$ ) versus latitude computed from data of Dowling and Ingersoll (1989) : a) under the GRS ; b) under the Oval BC.
Figure 13 : Velocity profile within the GRS from ( Mitchell et al 1981 ). They have observed that the velocity is nearly tangential to ellipses. Using a grid of concentric ellipses of constant eccentricity, the velocities have been plotted with respect to the semi-major axis $`𝒜`$ of the ellipse on which the measurement point lies. Results have been fitted by a quartic in $`𝒜`$. $`l^{}`$ is the jet width, defined as the width on which the jet velocity is greater than half the maximum velocity.
Figure 14 : A) Coefficient $`a`$ for a quadratic topography $`h(y)=ay^2`$ versus $`R^{}`$ computed from our QG model (86). The three cross show the coefficient $`a`$ computed from QG topography deduced from Dowling SW observed results. B) Difference of PV levels $`a_1a_1`$ versus $`R^{}`$ computed from our QG model (85). The dot line represents the planetary vorticity $`f_0`$ at the latitude of the center of the GRS. This show that one of the PV levels may be interpreted as vorticity generated by convection-plumes from the sublayer.
Figure 15 : The segregation parameter $`u`$ versus the Rossby deformation Radius $`R^{}`$ for the GRS. $`u`$ has been computed using the actual jet maximum velocity and width ( see section (4) ).
Figure 16 : The non dimensional parameter giving the shape of the curve : $`d`$ ( see (81) ) with respect to $`R^{}`$ in our model of the GRS. The dot line represents the critical value $`d=\frac{4}{9}`$ below which a vortex solution exists. The ratio of the length to the width of the GRS is approximately $`2`$. From figure 11 we conclude that this correspond to $`d`$ very close to the critical value $`\frac{4}{9}`$. From this figure, our model predicts that the Rossby deformation radius is $`R^{}=1800`$ km ( see section (4) for comments ).
Figure 17 : Definition of $`l(y)`$.
Figure 18 : Phase portraits of the Hamiltonian H (82) for $`y_0=0`$, governing the jet shape via differential equations (81) ( two periods in $`\theta `$ ). For vortices, we are looking for periodic solutions in $`y`$. Thus only trajectories of areas a) and b) are under interest. Conversely trajectories of area c) could correspond to oscillating jets. The parameters $`d`$ governs a transition between two type of phase portraits. A) For $`d<\frac{4}{9}`$ ( here $`d=0.075`$ ), trajectories of area a) can define $`y`$ as a function of $`\theta `$ corresponding to convex vortices. B ) For $`d>\frac{4}{9}`$ ( here $`d=0.075`$ ), for trajectories of area a), the curve $`y(\theta )`$ admits double points. Thus they can not define vortex boundaries.
Figure 19 : Resolution of the equations (93 and 95). The long-doted and the doted lines represent the two cases discussed in appendix (B).
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# Quantum Computation by Geometrical Means
## 1. Prologue
Abelian \[Sha\] and non-abelian \[Wil\] geometrical phases in quantum theory have been considered as a deep and fascinating subject. They provide a natural connection between the evolution of a physical system with degenerate structure and differential geometry. Here we shall present a model where these concepts can be explicitly applied for quantum computation \[Zan\].
The physical setup consists of an energy degenerate quantum system on which we perform an adiabatic isospectral evolution described by closed paths in the parametric space of external variables. The corresponding evolution operators acting on the code-state in the degenerate eigenspace are given in terms of holonomies and we can use them as quantum logical gates. This is a generalization of the Berry phase or geometrical phase, to the non-abelian case, where a non-abelian adiabatic connection, $`A`$, is produced from the geometrical structure of the degenerate spaces. In particular, on each point of the manifold of the external parameters there is a code-state attached and a transformation between these bundles of codes is dictated by the connection $`A`$.
In order to apply this theoretical construction to a concrete example we employ a model with $`\mathrm{𝐂𝐏}^2`$ geometry, that is a complex projective manifold with two complex coordinates. This is interpreted as a qubit \[Pac\]. A further generalization with the tensor product of $`m`$ $`\mathrm{𝐂𝐏}^2`$ models and additional interaction terms parametrized by the Grassmannian manifold, $`𝐆(4,2)`$, is interpreted as a model of quantum computer.
The initial code-state is written on the degenerate eigenspace of the system. The geometrical evolution operator is a unitary acting on it and it is interpreted as a logical gate. Due to adiabaticity the geometrical part of the evolution operator has a dimensionality equal to the degree of degeneracy of the eigenspace. Specific logical gates given by holonomies are constructed for a system with a tensor product structure resulting in universality while at the end a quantum optical application is sketched.
## 2. Coset Space Geometry
A transformation $`U(n)`$ between the states $`|\alpha `$, $`\alpha =1,\mathrm{},n`$ can be realized by all possible sub-$`U(2)`$ transformations between any two of those states, $`|i`$ and $`|j`$. A coset space can be produced as the factor with respect to some particular $`U(2)`$ symmetries of these transformations.
Examples of such constructions are given in the following:
* the $`\mathrm{𝐂𝐏}^2`$ projective space:
$$\mathrm{𝐂𝐏}^2\frac{U(3)}{U(2)\times U(1)}$$
The lines denote $`U(2)`$ transformations between the states represented here by “holes”. The $`U(3)`$ group could be interpreted by three lines connecting all the holes together. The distinction between “filled” and “unfilled” holes is due to the coset structure, which factors out the symmetry transformations, between $`|1`$ and $`|2`$, and denotes explicitly the non-symmetric ones between $`|1`$ or $`|2`$ and $`|\stackrel{~}{2}`$.
* the $`\left(\mathrm{𝐂𝐏}^2\right)^{\times m}\times \left(𝐆(4,2)_{int}\right)^{\times (m1)}`$ product space:
$$\mathrm{}\frac{U(3)}{U(2)\times U(1)}\times \frac{U(3)}{U(2)\times U(1)},\frac{U(4)}{U(2)\times U(2)}|_{int}\mathrm{}$$
Transformations can be performed between the states $`\{|1,|2\}`$ and $`\{|3,|4\}`$ due to their connections with the states $`\{|\stackrel{~}{2},|\stackrel{~}{4}\}`$, while an interaction in the tensor product space between $`|24`$ and $`|\stackrel{~}{2}\stackrel{~}{4}`$ gives transformations between those two sets. The transformations between only two states may be performed by linear operations with respect to $`U(2)`$ generators, while the combined transformations between the two qubits can be produced by bilinear generators which act simultaneously on the states of both of the $`\mathrm{𝐂𝐏}^2`$ models. The latter is denoted in the previous figure by the dashed lines, where the connection between the black dots indicates the simultanious action.
## 3. Degeneracy, Adiabaticity and Holonomies
Let us introduce the degenerate Hamiltonians $`H_0^1`$ and $`H_0^m`$ as follows
$$H_0^1=\left[\begin{array}{ccccc}0& 0& 0& & \\ 0& 0& 0& & \\ 0& 0& 1& & \end{array}\right],H_0^m=\underset{m}{}H_0^1.$$
The orbit, that is the parametric manifold of the unitary transformations which preserves the degenerate spectrum of $`H_0^1`$ is given by $`\mathrm{𝐂𝐏}^2`$. A sub-manifold of the orbit of $`H_0^m`$ in which we are interested in is given by the $`\left(\mathrm{𝐂𝐏}^2\right)^{\times m}\times \left(𝐆(4,2)_{int}\right)^{\times (m1)}`$ product manifold.
A general transformation parametrized by the $`\mathrm{𝐂𝐏}^2`$ space is given by $`𝐔(𝐳):=U_1(z_1)U_2(z_2)`$, with $`U_\alpha (z_\alpha )=\mathrm{exp}G_\alpha (z_\alpha )=\mathrm{exp}(z_\alpha |\alpha \stackrel{~}{2}|\overline{z}_\alpha |\stackrel{~}{2}\alpha |)`$. The complex parameter $`z_\alpha `$ may be decomposed as $`z_\alpha =\theta _\alpha \mathrm{exp}i\varphi _\alpha `$. Due to the $`2\times 2`$ sub-form of $`G_\alpha (z_\alpha )`$ we can rewrite $`U_\alpha (z_\alpha )`$ as
$$U_\alpha (z_\alpha )=\mathrm{𝟏}_\alpha ^{}+\mathrm{cos}\theta _\alpha \mathrm{𝟏}_\alpha +\frac{\mathrm{sin}\theta _\alpha }{\theta _\alpha }G_\alpha (z_\alpha ),$$
where $`\mathrm{𝟏}_\alpha ^{}=\mathrm{𝟏}\mathrm{𝟏}_\alpha `$ and $`\mathrm{𝟏}_\alpha =|\alpha \alpha |+|\stackrel{~}{2}\stackrel{~}{2}|`$. For the Grassmannian manifold $`𝐆(4,2)`$ we have, for example, the $`U(2)`$ rotation in the tensor product basis of two qubits, between the states $`|24`$ and $`|\stackrel{~}{2}\stackrel{~}{4}`$, given by $`𝐔(z)=\mathrm{exp}(z|24\stackrel{~}{2}\stackrel{~}{4}|\overline{z}|\stackrel{~}{2}\stackrel{~}{4}24|)`$, with $`z=\theta \mathrm{exp}i\varphi `$. The coordinates $`\{\lambda ^a\}=`$ {$`𝜽`$, $`\mathit{\varphi }`$} provide the parametric space which the experimenters control.
In the four dimensional manifold $`\mathrm{𝐂𝐏}^2`$ with coordinates $`\{\lambda ^a\}`$ a closed path, $`C`$, is drawn on a two-submanifold. Consider this evolution to be adiabatic as well as isospectral which is provided by the formula $`H(\lambda (t))=𝐔(\lambda (t))H_0𝐔^{}(\lambda (t))`$. As a result the state of the system, $`|\psi (t)`$, stays on the same eigenvalue of the Hamiltonian, taken in our example to be $`E_0=0`$, without level-crossing.
At the end of the loop $`C`$, spanned in time $`T=N\mathrm{\Delta }t`$, when divided in $`N`$ equal time intervals, we obtain
$`|\psi (T)`$ $`=𝐓e^{i_0^T𝐔H_0𝐔^{}𝑑t}|\psi (0)`$
$`=𝐓\underset{N\mathrm{}}{lim}{\displaystyle \underset{i=1}{\overset{N}{}}}𝐔_ie^{iH_0\mathrm{\Delta }t}𝐔_i^{}|\psi (0)`$
$`=𝐏\underset{N\mathrm{}}{lim}\left(\mathrm{𝟏}+{\displaystyle \underset{i=1}{\overset{N}{}}}A_i\mathrm{\Delta }\lambda _i\right)|\psi (0)`$
with $`A_i=𝐔_i^{}{\displaystyle \frac{\mathrm{\Delta }𝐔_i}{\mathrm{\Delta }\lambda _i}}\text{and}𝐔_i=𝐔(\lambda (t_i)).`$
Hence, the state $`|\psi (0)`$ acquires a geometrical unitary operator given by the holonomy of a connection $`A`$ as
$$\mathrm{\Gamma }_A(C):=𝐏\mathrm{exp}_CA,\text{where}(A^{\lambda ^a})_{\alpha \beta }:=\alpha |𝐔^{}(\lambda )\frac{}{\lambda ^a}𝐔(\lambda )|\beta .$$
The states $`|\alpha `$ and $`|\beta `$ belong to the same degenerate eigenspace of $`H_0`$ and $`\lambda ^a`$’s are the real control parameters.
The produced unitary operator is an effect of the non-commutativity of the control transformations which produce effectively a curvature. In the case of the Berry phase produced for example in front of the spin states of an electron when placed in a magnetic field, the non-commutativity is between the $`U(2)`$ control unitaries which change the direction of the magnetic field in the three dimensional space. What is presented here is the generalization of the Berry phase to the non-abelian case.
The $`\mathrm{\Gamma }_A(C)`$’s produced by $`\mathrm{𝐂𝐏}^2`$ for various loops $`C`$, generate the whole $`U(2)`$,
$$\{\mathrm{\Gamma }_A(C);C\mathrm{𝐂𝐏}^2\}U(2).$$
In the case of $`m`$ qubits with their proper interactions, the produced group is $`U(2^m)`$,
$$\{\mathrm{\Gamma }_A(C);C\left(\mathrm{𝐂𝐏}^2\right)^{\times m}\times \left(𝐆(4,2)_{int}\right)^{\times (m1)}\}U(2^m).$$
## 4. Quantum Computation
In order to perform quantum computation by using the above constructions we consider the following identifications:
QUANTUM CODE $``$ Degenerate States, $`|\psi (0)`$
LOGICAL GATES $``$ Holonomies, $`\mathrm{\Gamma }_A(C)`$
Let us first investigate the $`\mathrm{𝐂𝐏}^2`$ case. The basic question is how we can generate a general $`U(2)`$ element by moving along a closed path, $`C`$. Or in other words, for a specific $`UU(2)`$ which loop $`C`$ is such that $`\mathrm{\Gamma }_A(C)=U`$. In general, $`𝐠u(2)\text{loop}C\mathrm{𝐂𝐏}^2`$ manifold, such that $`\mathrm{\Gamma }_A(C)=\mathrm{exp}𝐠`$, which is the statement of irreducibility of the connection $`A`$ \[Zan\]. To answer the above question we perform the following analysis. The loop integral
$$_CA=_CA_{\lambda ^1}𝑑\lambda ^1+A_{\lambda ^2}d\lambda ^2+\mathrm{}$$
is the main ingredient of the holonomy. Due to the path ordering symbol it is not possible to just calculate it and evaluate its exponential, as in general the connection components do not commute with each other. Still it is possible to consider the following restrictions in the position of the loop. Choose $`C`$ such that:
* it belongs to one plane $`(\lambda ^i,\lambda ^j)=(\theta _i,\varphi _j)`$ or $`(\theta _i,\theta _j)`$, hence only two components of $`A`$ are involved,
* the position of the plane is such that the connection, $`A`$, restricted on it become, $`A|_{(\lambda ^i,\lambda ^j)}=(A^{\lambda ^i}=0,A^{\lambda ^j})`$, that is these two components commute with each other. Still it is important that their related field strength component, $`F_{ij}=_iA_j_jA_i+[A_i,A_j]`$, is non-vanishing in order to obtain a non-trivial holonomy. Such a requirement is possible for the $`\mathrm{𝐂𝐏}^2`$ model and for a wide class of other models.
On the planes where those conditions are satisfied the evaluation of the holonomy is trivially given by just exponentiating the loop integral of the connection without worrying about the path ordering symbol. Hence,
$$\mathrm{\Gamma }_A(C)=𝐏\mathrm{exp}_CA=\mathrm{exp}(\mathrm{\Sigma }𝐠)=\mathrm{𝟏}_{2\times 2}\mathrm{cos}\mathrm{\Sigma }+𝐠\mathrm{sin}\mathrm{\Sigma },$$
where $`\mathrm{\Sigma }`$ represents the area enclosed by the loop $`C`$ projected on the sphere associated with the compactified $`\mathrm{𝐂𝐏}^2`$ manifold. This area may be varied desirably. Furthermore, we are able to obtain a complete set of generators $`𝐠`$ by choosing $`C`$ to lie on different planes. In detail we may obtain for $`𝐠`$ the following forms
$`i|\alpha \alpha |`$ $`:=`$ $`i\sigma _\alpha ^3,\alpha =1,\mathrm{\hspace{0.17em}\hspace{0.17em}2}`$
$`i(i|12|+i|21|)`$ $`:=`$ $`i\sigma ^2,`$
$`i(|12|+|21|)`$ $`:=`$ $`i\sigma ^1.`$
The $`\sigma _\alpha ^3`$ generators are similar to a Berry phase and they are produced by paths $`C_1`$ on the $`(\theta _\alpha ,\varphi _\alpha )`$ planes. The corresponding holonomy is the exponential of this generator multiplied by the area, $`\mathrm{\Sigma }_1`$, of the surface the path $`C_1`$ encloses when projected on a sphere $`S^2(2\theta _\alpha ,\varphi _\alpha )`$ with spherical coordinates $`2\theta _\alpha `$ and $`\varphi _\alpha `$,
$$i\sigma _\alpha ^3:C_1(\theta _\alpha ,\varphi _\alpha )\mathrm{\Gamma }_A(C_1)=\mathrm{exp}i\mathrm{\Sigma }_1\sigma _\alpha ^3.$$
The $`\sigma ^2`$ generator is produced by a path $`C_2`$ along the plane $`(\theta _1,\theta _2)`$ positioned at $`\varphi _1=\varphi _2=0`$, while $`\sigma ^1`$ is produced by a path, $`C_3`$, along a parallel plane positioned at $`\varphi _1=\frac{\pi }{2}`$ and $`\varphi _2=0`$. Their corresponding areas are $`\mathrm{\Sigma }_2`$ and $`\mathrm{\Sigma }_3`$. For example
$$i\sigma ^2:C_2(\theta _1,\theta _2)|_{\varphi _1=0,\varphi _2=0}\mathrm{\Gamma }_A(C_2)=\mathrm{exp}i\mathrm{\Sigma }_2\sigma ^2.$$
Altogether we have $`2^2`$ independent generators spanning the Lie algebra of $`U(2)`$.
For the case of the two qubit interaction the corresponding connection components are given by
$$A_\theta =\text{diag}(0,0,0,0),A_\varphi =\text{diag}(0,0,0,i\mathrm{sin}^2\theta )$$
which are written in the basis $`\{|13,|14,|23,|24\}`$. A loop $`C`$ on the $`(\theta ,\varphi )`$ plane will produce the following holonomy
$$\mathrm{\Gamma }_A(C)=\text{diag}(1,1,1,e^{i\mathrm{\Sigma }}),\mathrm{\Sigma }=_{D(C)}𝑑\theta 𝑑\varphi \mathrm{sin}2\theta ,$$
where $`\mathrm{\Sigma }`$ can also be interpreted as an area on the sphere $`S^2(2\theta ,\varphi )`$.
## 5. One and Two Qubit Logical Gates
By performing appropriate loops we can obtain one qubit phase rotations as well as two qubit gates such as a controlled phase rotation $`U_{CPH}`$.
Analytically, by spanning the indicated areas we may obtain
$$U_1=\mathrm{exp}\left[\begin{array}{ccc}i\mathrm{\Sigma }_1& 0& \\ 0& 0& \end{array}\right],U_2=\mathrm{exp}\left[\begin{array}{ccc}0& 0& \\ 0& i\mathrm{\Sigma }_1& \end{array}\right],U_3=\mathrm{exp}\left[\begin{array}{ccc}0& \mathrm{\Sigma }_2& \\ \mathrm{\Sigma }_2& 0& \end{array}\right].$$
The combinations
$$U_1U_2=\mathrm{exp}(i\sigma _3\mathrm{\Sigma }_1),U_3=\mathrm{exp}(i\sigma _2\mathrm{\Sigma }_2)$$
can give any $`U(2)`$ transformation and hence any one qubit rotation.
For the two qubit gates we can construct easily the controlled rotation $`U_{CPH}=\text{diag}(1,1,1,\mathrm{exp}i\mathrm{\Sigma })`$ between any pair of qubits. It is generated by a loop $`C`$ on the $`(\theta ,\varphi )`$ plane. Together with the one qubit rotations they provide a universal set of gates.
## 6. Epilogue
Apart from the intriguing theoretical formulation of holonomic computation there are several aspects of it, which have appealing technical advantages. Without overlooking the difficulties posed to an experimenter for performing continuous control over a system in order to span a loop, there are several unique characteristics of it, which await for exploitation. For example, robustness of the control procedure in terms of the spanned area, according to errors in the actual form of the performed loop, as well as the isolation of the degenerate states as a calculational space may prove to be advantages worth exploring.
In quantum optics displacing devices, squeezing devices and interferometers acting on laser beams can provide the control parameters for the holonomic computation. Each laser beam is placed in a non-linear Kerr medium with degenerate Hamiltonian $`H_0=n(n1)`$, where $`n`$ is the photon numbering operator. The degenerate states $`|0`$ and $`|1`$ are the basis for encoding one qubit which is manipulated by displacing and squeezing devices. Any two qubit interactions can be implemented by interferometers \[Cho\].
It is challenging for the experimenters to produced the desired closed paths.
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# Untitled Document
Gerbes of chiral differential operators. II
Vassily Gorbounov, Fyodor Malikov, Vadim Schechtman
Contents
Introduction
§0. Recollections and notation
§1. Vertex algebroids
§2. From vertex algebras to vertex algebroids
§3. Category of vertex algebroids
§4. Cofibered structure
§5. Chern-Simons
§6. Atiyah
§7. Gerbes of vertex algebroids
§8. Vertex envelope of a conformal algebra
§9. Enveloping algebra of a vertex algebroid
References
Introdution
The aim of this paper is a study of certain class of chiral (vertex) algebras which we call chiral algebras of differential operators.
I.0. To explain this notion, let us take up a ”naive” approach to vertex algebras. Let us fix (in this Introduction) a ground field $`k`$ of characteristic $`0`$. Recall (cf. \[K\]) that a $`_0`$-graded vertex algebra is a $`_0`$-graded $`k`$-vector space $`V=_{i0}V_i`$ equipped with a distinguished vacuum vector $`\text{1}V_0`$ and a family of bilinear operations
$${}_{(n)}{}^{}:V_kVV,xyx_{(n)}y,$$
$`(n)`$, the operation <sub>(n)</sub> having the degree $`n1`$, that is, $`V_{i(n)}V_jV_{i+jn1}`$ where by definition $`V_i=0`$ for $`i<0`$. These operations must satisfy a family of quadratic Borcherds identities. Throughout this paper we shall deal only with $`_0`$-graded vertex algebras, so we shall omit the words ”$`_0`$-graded” when speaking about them.
Note that in particular the operation <sub>(-1)</sub> has degree $`0`$. It is neither commutative nor associative in general.
We are interested in a question: what structure on the subspace $`V_1:=V_0V_1`$ is induced by the structure of a vertex algebra on $`V`$?
The claims in italics below are easily proved using the Borcherds identities, see the main body of the paper.
First of all, $`V_0`$ is a commutative associative $`k`$-algebra with unit 1 with respect to <sub>(-1)</sub>. Let us denote this algebra by $`A`$, and the operation $`a_{(1)}b`$ on it by $`ab`$. Note that for each $`i`$ the operation <sub>(-1)</sub> gives a mapping
$$AV_iV_i,$$
$`(I1)`$
to be denoted $`ax`$, but this mapping does not make $`V_i`$ an $`A`$-module since the associativity $`abx=(ab)x`$ is not satisfied in general.
We have a mapping $`:V_iV_{i+1}`$ defined by $`x=x_{(2)}\text{1}`$. Let us denote by $`\mathrm{\Omega }V_1`$ a $`k`$-subspace generated by all elements $`ab,a,bA`$. The operation (I1) makes $`\mathrm{\Omega }`$ an $`A`$-module (i.e. the associativity holds true on this subspace) and $`:A\mathrm{\Omega }`$ is a derivation.
Let $`T`$ be the quotient space $`V_1/\mathrm{\Omega }`$. The operation (I1) induces the structure of an $`A`$-module on $`T`$ and operation <sub>(0)</sub> induces a structure of a Lie algebra on $`T`$ and an action of $`T`$ on $`A`$ by derivations. These structures are compatible: they make $`T`$ a Lie algebroid (see 0.2) over $`A`$.
The operation <sub>(0)</sub> induces an action of $`T`$ on $`\mathrm{\Omega }`$ which makes $`\mathrm{\Omega }`$ a module over the Lie algebra $`T`$, and $``$ is a morphism of $`T`$-modules.
The operation <sub>(1)</sub> induces an $`A`$-bilinear pairing $`T\mathrm{\Omega }A`$ which is a morphism of $`T`$-modules.
Let us choose a splitting
$$s:TV_1$$
$`(I2)`$
of the canonical projection $`\pi :V_1T`$. Define a mapping $`\gamma :A\times T\mathrm{\Omega }`$ by $`\gamma (a,\tau )=s(a\tau )as(\tau )`$. Define a symmetric mapping $`,:T\times TA`$ by $`\tau ,\tau ^{}=s(\tau )_{(1)}s(\tau ^{})`$. Finally, define a skew symmetric operation $`c:T\times T\mathrm{\Omega }`$ by $`c(\tau ,\tau ^{})=s([\tau ,\tau ^{}])[s(\tau ),s(\tau ^{})]`$ where we set $`[x,y]:=\frac{1}{2}(x_{(0)}yy_{(0)}x)`$.
These three operations, $`\gamma ,,`$ and $`c`$, satisfy certain compatibilities, listed below in 1.3, (A1) — (A5).
Thus, we have assigned to a vertex algebra $`V`$ (with a splitting (I2)) a collection of data $`𝒜=(A,T,\mathrm{\Omega },,\gamma ,,,c)`$ as above satisfying properties (A1) — (A5). We call such collection of data a vertex algebroid which is a central hero of our story.
The above procedure gives a functor form the category of vertex algebras to the (appropriately defined) category of vertex algebroids. This functor admits a left adjoint $`U`$, called vertex envelope. A vertex algebra of the form $`U𝒜`$ is called a chiral algebra of differential operators.
I.2. Let $`X`$ be a smooth algebraic variety over $`k`$, $`UX`$ an affine Zariski open subset, $`A=𝒪(U)`$. Let us consider the class of vertex algebroids of the form $`𝒜=(A,T,\mathrm{\Omega },,\mathrm{})`$ with $`T=Der_k(A),\mathrm{\Omega }=\mathrm{\Omega }_k^1(A),=d_{DR}`$ — the usual de Rham differential, and with the standard action of $`T`$ on $`\mathrm{\Omega }`$ by the Lie derivative.
Such vertex algebroids (with suitably defined morphisms) form a groupoid (i.e. a category where all morphisms are invertible) $`𝒜lg(A)`$. This category is nonempty if there exists an $`A`$-base of $`T`$ consisting of commuting vector fields.
This groupoid has a remarkable additional structure. To describe this structure, let us define another groupoid $`\mathrm{\Omega }^{[2,3}(A)`$ as follows. By definition, its objects will be all closed $`3`$-forms on $`A`$. If $`\omega ,\omega ^{}`$ are two such $`3`$-forms, the set of morphisms between them is defined as
$$Hom_{\mathrm{\Omega }^{[2,3}(A)}(\omega ,\omega ^{})=\{\eta \mathrm{\Omega }_{A/k}^2|d_{DR}\eta =\omega ^{}\omega \}$$
The compostion is defined in the obvious way. The addition of $`3`$-forms induces a functor
$$\mathrm{\Omega }^{[2,3}(A)\times \mathrm{\Omega }^{[2,3}(A)\mathrm{\Omega }^{[2,3}(A)$$
which makes $`\mathrm{\Omega }^{[2,3}(A)`$ an Abelian group in categories.
The first main result of this paper is
Theorem I1. The groupoid $`𝒜lg(A)`$ admits a canonical structure of a $`\mathrm{\Omega }^{[2,3}(A)`$-Torseur.
This is Theorem 7.2 of the paper. It means that one has a functor
$$\stackrel{}{+}:𝒜lg(A)\times \mathrm{\Omega }^{[2,3}(A)𝒜lg(A),(𝒜,\omega )𝒜\stackrel{}{+}\omega $$
which defines an Action of the Abelian group $`\mathrm{\Omega }^{[2,3}(A)`$ on $`𝒜lg(A)`$, such that for each $`𝒜𝒪b𝒜lg(A)`$ the induced functor
$$𝒜\stackrel{}{+}\mathrm{?}:\mathrm{\Omega }^{[2,3}(A)𝒜lg(A)$$
is an equivalence of categories.
I.3. Sheafifying the previous construction one gets a sheaf of groupoids, i.e. gerbe $`𝒜lg_X`$ over $`X`$, with the lien $`\mathrm{\Omega }_X^{[2,3}:=\mathrm{\Omega }_X^2\mathrm{\Omega }^{3,closed}`$ (a short complex of sheaves over $`X`$, with $`\mathrm{\Omega }_X^2`$ sitting in degree $`0`$).
By a simple general homological construction (see 7.3), this gerbe gives rise to a characteristic class
$$c(𝒜lg_X)H^2(X;\mathrm{\Omega }_X^{[2,3})$$
with the following property.
The groupoid of global sections $`𝒜lg_X(X)`$ is nonempty iff $`c(𝒜lg_X)=0`$. If this is the case then its set of connected components is a nonempty $`H^1(X;\mathrm{\Omega }_X^{[2,3})`$-torseur and the group of automorphisms of its object is isomorphic to $`H^0(X;\mathrm{\Omega }_X^{[2,3})`$.
Our next aim is to calculate the class $`c(𝒜lg_X)`$. One shows (see Theorem 7.10) that it essentially coincides with the second component of the Chern character of the tangent bundle $`𝒯_X`$. Let us describe an ”explicit formula” for it.
Let $`E`$ be an arbitrary finite dimensional vector bundle over $`X`$, given by a Cech cocycle $`(g_{ij})Z^1(𝒰;GL_r(𝒪_X))`$ on a Zariski open covering $`𝒰=\{U_i\}`$ of $`X`$. Define a Cech two-cocycle $`ch_2(E)Z^2(𝒰;\mathrm{\Omega }_X^{[2,3})`$ by
$$ch_2(E)=(\frac{1}{2}tr(g_{ij}^1g_{jk}^1dg_{jk}dg_{ij}),\frac{1}{6}tr((g_{ij}^1dg_{ij})^3))$$
$`(ACS)`$
This class may be called the ”Atiyah-Chern-Simons class” of $`E`$. Its first component, $`\alpha (E)`$, is an element of $`Z^2(𝒰;\mathrm{\Omega }_X^2)`$ which is the usual ”Atiyah” representative of the degree $`2`$ part of the Chern character of $`E`$ ”style Hodge” living in $`H^2(X;\mathrm{\Omega }_X^2)`$, while the second component, $`\beta (E)`$ (”Chern-Simons”), is a Cech $`1`$-cochain with coefficients in $`\mathrm{\Omega }_X^{3,closed}`$. The de Rham differential of $`\alpha (E)`$ is equal to the Cech coboundary of $`\beta (E)`$. One checks that the cohomology class of (ACS) does not depend on the choice of trivialization of $`E`$.
Our second main result (Theorem 7.5) is
Theorem I.2. The class $`c(𝒜lg)`$ is equal to (the cohomology class of) $`2ch_2(𝒯_X)`$.
I.4. Our third main topic is an explicit construction of the enveloping algebra $`U𝒜`$ of a vertex algebroid $`𝒜`$ and ”Poincaré-Birkhoff-Witt” type theorem for it. Let us formulate the last theorem (see Theorem 9.18).
Theorem I.3. Each sheaf of algebras of chiral differential operators $`𝒟=U𝒜,𝒜𝒜lg_X(X)`$ admits a canonical filtration $`F`$, compatible with the conformal grading and finite on each component of fixed conformal weight, such that the corresponding graded algebra (which is a sheaf of $`_0`$-graded commutative vertex algebras) is
$$gr_F(𝒟)=Sym_{𝒪_X}\left\{(_{i1}𝒯^{(i)})(_{i1}\mathrm{\Omega }^{(i)})\right\}$$
where $`𝒯^{(i)}`$ (resp. $`\mathrm{\Omega }^{(i)}`$) is a copy of $`𝒯_X`$ (resp. of $`\mathrm{\Omega }_X^1`$) living in conformal weight $`i`$.
As a preparation to this theorem, we study in Section 8 the vertex envelope $`UC`$ of an arbitrary conformal algebra $`C`$. Although its construction is more or less contained in Kac’s book \[K\], we present here in a sense more direct construction of $`C`$, maybe of independent interest. Note in particular an amusing Lie bracket (8.21.1) which is defined on an arbitrary conformal algebra.
In the hope of possible arithmetical applications, we adopted in this paper a ”characteristic free” approach to vertex algebras. Contrary to this Introduction (and to the tradition), in the main body of this work the ground ring will be an arbitrary commutative ring containing $`1/2`$ (with one exception: the PBW theorem 9.18). This generalization is achieved without great effort. Note that the original Borcherds’ definition of vertex algebras was characteristic-free.
I.4. This paper may be regarded as an ”algebraic” version of our last note \[GMS\] where we worked in the analytical category. However, the approach adopted here is quite different from op. cit. (cf. also \[MSV\] and \[MS\]).
During the preparation of this work we have enourmously benefited from the discussions and correspondence with Sasha Beilinson. He was the first to suggest (by analogy with the classical picture, \[BB\]) that in the algebraic category the situation is more subtle than in the analytical one. He has sent us a note \[BD1\] where similar questions are treated from a different point of view. In a sense, a great part of this work is a result of the attempts to understand this note (which still remains mysterious for us). For the details on Beilinson-Drinfeld approach, see \[BD2\], 3.8.
Another special gratitude goes to Hélène Esnault. She has found a mistake in the previous formulation of our main result whose correction lead to the discovery of the class (ACS), and greatly helped in some computations. Theorem 7.10 was obtained in collaboraion with her.
This work was mostly done while V.S. visited IHES. He is grateful to this institute for the support and excellent working conditions.
§ 0. Recollections and Notation
0.0. Throughout this paper, a commutative ring (algebra) will mean a commutative associative ring (algebra) with unit.
$`k`$ will denote a fixed ground commutative ring. We will often assume that $`1/2k`$; we will indicate this assumption when necessary.
Algebra will mean a $`k`$-algebra; $``$ will mean the tensor product over $`k`$.
In a nonassociatve algebra, $`abc\mathrm{}`$ will mean $`a(b(c(\mathrm{}))\mathrm{})`$.
$`_0`$ will denote the set of nonnegative integers. The binomial coefficients are defined by
$$\left(\genfrac{}{}{0pt}{}{a}{b}\right)=\frac{a(a1)\mathrm{}(ab+1)}{b!},a,b_0,$$
$`(\mathrm{0.0.1})`$
and $`0`$ if $`b<0`$. We have
$$\left(\genfrac{}{}{0pt}{}{a1}{b}\right)=(1)^b\left(\genfrac{}{}{0pt}{}{a+b}{b}\right)$$
$`(\mathrm{0.0.2})`$
for all $`a0`$.
0.1. Let $`A`$ be a commutative algebra. The Lie algebra of $`k`$-derivations $`T_A:=Der_k(A,A)`$ acts on the $`A`$-module of Kähler $`1`$-differentials $`\mathrm{\Omega }_A^1:=\mathrm{\Omega }_{A/k}^1`$ according to the usual rule
$$\tau (adb)=\tau (a)db+ad\tau (b)$$
$`(\mathrm{0.1.1})`$
The de Rham differential $`d:A\mathrm{\Omega }_A^1`$ commutes with the $`T_A`$-action.
We have the canonical $`A`$-bilinear pairing
$$,:T_A\times \mathrm{\Omega }_A^1A,\tau ,adb=a\tau (b)$$
$`(\mathrm{0.1.2})`$
We have
$$\tau (a\omega )=\tau (a)\omega +a\tau (\omega )$$
$`(\mathrm{0.1.3})`$
$$(a\tau )(\omega )=a\tau (\omega )+\tau ,\omega da$$
$`(\mathrm{0.1.4})`$
$$\tau (\nu ,\omega )=[\tau ,\nu ],\omega +\nu ,\tau (\omega )$$
$`(\mathrm{0.1.5})`$
$$\tau ,da=\tau (a)$$
$`(\mathrm{0.1.6})`$
$`(aA,\tau ,\nu T_A,\omega \mathrm{\Omega }_A^1)`$. (Of course the last formula is a particular case of (0.1.2).)
0.2. A Lie $`A`$-algebroid is a Lie algebra $`T`$ acting by derivations on $`A`$ and equipped with a structure of an $`A`$-module, such that
$$[\tau ,a\nu ]=a[\tau ,\nu ]+\tau (a)\nu $$
$`(\mathrm{0.2.1})`$
and
$$(a\tau )(b)=a\tau (b)$$
$`(\mathrm{0.2.2})`$
for all $`\tau ,\nu T;a,bA`$.
0.2.1. Direct image (pushout). Let $`B`$ be a commutative $`A`$-algebra, $`i:AB`$ be the structure morphism and $`T`$ be an $`A`$-algebroid Lie. Assume that $`T`$, as a Lie algebra, acts on $`B`$ by derivations in such a way that $`i`$ is a morphism of $`T`$-modules and
$$(a\tau )(b)=a\tau (b)$$
$`(aA,bB,\tau T)`$.
Then the $`B`$-module $`T_B:=B_AT`$ admits a canonical structure of a $`B`$-algebroid Lie. Namely, the Lie bracket on $`T_B`$ is given by
$$[b_1\tau _1,b_2\tau _2]=b_1b_2[\tau _1,\tau _2]+\tau _1(b_2)b_1\tau _2\tau _2(b_1)b_2\tau _1$$
$`(\mathrm{0.2.1.1})`$
and the action of $`T_B`$ on $`B`$ is defined by
$$(b_1\tau )(b_2)=b_1\tau (b_2)$$
$`(\mathrm{0.2.1.2})`$
Vertex algebras
0.3. In his original paper, \[B\], Borcherds defined a vertex algebra over an arbitrary commutative ring. However, later most people preferred to work over the complex numbers. It is fairly obvious that all the general theorems of Kac’s book \[K\] are true over an arbitrary field of characteristic $`0`$. A little less obvious, but true, is that with a minor modification of the definitions, they remain true over an arbitrary commutative ring. Below we recall the definitions and results to be used in the sequel, and explain these modifications.
Throughout this work, we will deal only with $`_0`$-graded vertex algebras.
0.4. Definition. A $`_0`$-graded conformal algebra is a $`_0`$-graded $`k`$-module $`C=C_i`$, together with a family of endomorphisms
$$^{(j)}:CC,\text{of degree}j,j_0,$$
such that
$$^{(i)}^{(j)}=\left(\genfrac{}{}{0pt}{}{i+j}{i}\right)^{(i+j)};^{(0)}=Id,$$
$`(\mathrm{0.4.1})`$
and a family of bilinear operations
$${}_{(n)}{}^{}:C\times CC,(a,b)a_{(n)}b,\text{of degree }n1,n_0,$$
such that
$$(^{(i)}a)_{(n)}b=(1)^i\left(\genfrac{}{}{0pt}{}{n}{i}\right)a_{(ni)}b$$
$`(\mathrm{0.4.2})`$
$$a_{(n)}b=(1)^{n+1}\underset{j=0}{\overset{\mathrm{}}{}}(1)^j^{(j)}(b_{(n+j)}a)$$
$`(\mathrm{0.4.3})`$
$$a_{(m)}b_{(n)}c=b_{(n)}a_{(m)}c+\underset{j=0}{\overset{m}{}}\left(\genfrac{}{}{0pt}{}{m}{j}\right)(a_{(j)}b)_{(m+nj)}c$$
$`(\mathrm{0.4.4})`$
for all $`a,b,cC;m,n,i_0`$.
Cf. \[K\], Definition 2.7b. In each conformal algebra we have the following identities:
$$(a_{(m)}b)_{(n)}c=\underset{j=0}{\overset{m}{}}(1)^j\left(\genfrac{}{}{0pt}{}{m}{j}\right)\{a_{(mj)}b_{(n+j)}c(1)^mb_{(m+nj)}a_{(j)}c\}$$
$`(\mathrm{0.4.5})`$
and
$$^{(j)}(a_{(n)}b)=\underset{p=0}{\overset{j}{}}^{(p)}a_{(n)}^{(jp)}b$$
$`(\mathrm{0.4.6})`$
We leave their direct proof to the reader. The proof of (0.4.5) uses only the axiom (0.4.4) (and (0.4.1)), and the proof of (0.4.6) uses the axioms (0.4.2) and (0.4.3).
We also have
$$a_{(n)}^{(j)}b=\underset{p=0}{\overset{j}{}}\left(\genfrac{}{}{0pt}{}{n}{p}\right)^{(jp)}(a_{(np)}b)$$
$`(\mathrm{0.4.7})`$
This is proven by induction on $`j`$.
0.5. First Definition of a Vertex Algebra. A $`_0`$-graded vertex algebra is a $`_0`$-graded $`k`$-module $`V=V_i`$, equipped with a distinguished vector $`\text{1}V_0`$ (vacuum vector) and a family of bilinear operations
$${}_{(n)}{}^{}:V\times VV,(a,b)a_{(n)}b,\text{ of degree }n1,n,$$
$`(\mathrm{0.5.1})`$
such that
$$\text{1}_{(n)}a=\delta _{n,1}a;a_{(1)}\text{1}=a;a_{(n)}\text{1}=0\text{ if }n0,$$
$`(\mathrm{0.5.2})`$
and
$$\underset{j=0}{\overset{\mathrm{}}{}}\left(\genfrac{}{}{0pt}{}{m}{j}\right)(a_{(n+j)}b)_{(m+lj)}c=$$
$$=\underset{j=0}{\overset{\mathrm{}}{}}(1)^j\left(\genfrac{}{}{0pt}{}{n}{j}\right)\left\{a_{(m+nj)}b_{(l+j)}c(1)^nb_{(n+lj)}a_{(m+j)}c\right\}$$
$`(\mathrm{0.5.3})`$
for all $`a,b,cV,m,n,l`$.
Cf. \[B\], Section 4, \[K\], Prop. 4.8. (b). The important particular case of (0.5.3) corresponds to $`m=0`$:
$$(a_{(n)}b)_{(l)}c=\underset{j=0}{\overset{\mathrm{}}{}}(1)^j\left(\genfrac{}{}{0pt}{}{n}{j}\right)\left\{a_{(nj)}b_{(l+j)}c(1)^nb_{(n+lj)}a_{(j)}c\right\}$$
$`(\mathrm{0.5.4})`$
cf. (0.4.5).
Set
$$^{(j)}a:=a_{(1j)}\text{1},j_0$$
$`(\mathrm{0.5.5})`$
This way we get endomorphisms $`^{(j)}`$ of $`V`$ of degree $`j`$. It follows from (0.5.2) that
$$^{(j)}\text{1}=\delta _{j,0}\text{1}$$
$`(\mathrm{0.5.6})`$
and
$$^{(0)}=Id$$
$`(\mathrm{0.5.7})`$
and (0.5.4) applied to $`b=c=\text{1}`$ gives
$$^{(i)}^{(j)}=\left(\genfrac{}{}{0pt}{}{i+j}{i}\right)^{(i+j)}$$
$`(\mathrm{0.5.8})`$
We have the commutativity formula
$$a_{(n)}b=(1)^{n+1}\underset{j=0}{\overset{\mathrm{}}{}}(1)^j^{(j)}(b_{(n+j)}a)$$
$`(\mathrm{0.5.9})`$
for all $`a,bV,n`$. The proof will be given below, see the last paragraph of the next subsection.
One deduces from (0.5.4) that
$$(^{(j)}a)_{(n)}b=(1)^j\left(\genfrac{}{}{0pt}{}{n}{j}\right)a_{(nj)}b$$
$`(\mathrm{0.5.10})`$
and
$$^{(j)}(a_{(n)}b)=\underset{p=0}{\overset{j}{}}(^{(p)}a)_{(n)}^{(jp)}b$$
$`(\mathrm{0.5.11})`$
for all $`n`$. We have the Operator Product Expansion (OPE) formula,
$$[x_{(m)},y_{(n)}]=\underset{j0}{}\left(\genfrac{}{}{0pt}{}{m}{j}\right)(x_{(j)}y)_{(m+nj)}(m,n)$$
$`(\mathrm{0.5.12})`$
cf. \[K\], (4.6.7) and the end of the next Subsection.
0.6. Following \[K\], let us give an equivalent definition of a vertex algebra. If $`M`$ is a $`k`$-module, $`M[[z,z^1]]`$ will denote as usually the module of formal power series $`_{n=\mathrm{}}^{\mathrm{}}a_nz^n,a_nM`$. Let us define endomorphisms $`^{(j)},j_0`$, of $`M[[z,z^1]]`$ by
$$^{(j)}(\underset{n}{}a_nz^n)=\underset{n}{}\left(\genfrac{}{}{0pt}{}{n}{j}\right)a_nz^{nj}$$
$`(\mathrm{0.6.1})`$
Second Definition of a Vertex Algebra. A $`_0`$-graded vertex algebra is a $`_0`$-graded $`k`$-module $`V=_iV_i`$ equipped with
— a distingushed vector $`\text{1}V_0`$ (vacuum vector);
— a family of endomorphisms $`^{(j)}:VV`$ of degree $`j`$, $`j_0`$, such that $`^{(0)}=Id`$ and $`^{(i)}^{(j)}=\left(\genfrac{}{}{0pt}{}{i+j}{i}\right)^{(i+j)}`$;
— a linear mapping $`VEnd(V)[[z,z^1]],aa(z)=_na_{(n)}z^{n1}`$ such that $`dega_{(n)}=degan1`$.
These data should satisfy the following axioms.
Translation invariance. For all $`j_0`$, $`[^{(j)},a(z)]=^{(j)}a(z)`$.
Vacuum. $`^{(j)}\text{1}=\delta _{j,0}\text{1};\text{1}(z)=\text{1};a_{(n)}\text{1}=0`$ for $`n0`$ and $`a_{(1)}\text{1}=a`$.
Locality. $`(zw)^N[a(z),b(w)]=0`$ for $`N>>0`$.
This is a modification of the definition given in \[K\], 4.1, which works over an arbitrary base commutative ring. Proposition 4.1 of op. cit. remains true if we understand by $`e^z`$ the expression
$$e^z:=\underset{j=0}{\overset{\mathrm{}}{}}z^j^{(j)}$$
$`(\mathrm{0.6.2})`$
(we cannot use Lemma 4.1 of op.cit. anymore!). The argument of op. cit. 4.2 shows that the commutativity formula (0.5.9) is true for the vertex algebras in the second definition. Finally, the proof of op. cit., Proposition 4.8 (b) shows that the first and the second definitions are equivalent. The proof from \[K\], 4.6 works to give the proof of the OPE formula (0.5.12).
0.7. Theorem. Let $`V=V_i`$ be a $`_0`$-graded $`k`$-module equipped with a distinguished vector $`\text{1}V_0`$ and a family of endomorphisms $`^{(j)}`$ of degree $`j`$, $`j_0`$ such that $`^{(0)}=Id,^{(i)}^{(j)}=\left(\genfrac{}{}{0pt}{}{i+j}{i}\right)^{(i+j)}`$ and $`^{(j)}\text{1}=\delta _{j,0}\text{1}`$.
Assume that we are given a family of homogeneous vectors $`\{a^\alpha \}V`$ and a family of formal power series (”distributions”)
$$\{a^\alpha (z)=\underset{n}{}a_{(n)}^\alpha z^{n1}\}End(V)[[z,z^1]],dega_{(n)}^\alpha =dega^\alpha n1,$$
satisfying the following conditions
(t) $`[^{(j)},a^\alpha (z)]=^{(j)}a^\alpha (z)`$;
(v) $`a_{(n)}^\alpha \text{1}=0`$ for $`n0`$; $`a_1^\alpha \text{1}=a^\alpha `$;
(l) The distributions $`a^\alpha (z)`$ are mutually local, i.e. for all $`\alpha ,\beta `$, $`(zw)^N[a^\alpha (z),a^\beta (w)]=0`$ for sufficiently large $`N`$.
Let $`S`$ denote the set of all vectors of the form
$$a_{j_1\mathrm{}j_N}^{\alpha _1\mathrm{}\alpha _N}:=a_{(1j_1)}^{\alpha _1}\mathrm{}a_{(1j_N)}^{\alpha _N}\text{1},N,j_i_0,$$
$`(\mathrm{0.7.1})`$
Assume that we are given a map
$$VEnd(V)[[z,z^1]],aa(z)=a_{(n)}z^{n1}$$
$`()`$
having the following property:
(P) There exists a subset $`S^{}S`$ which generates $`V`$ as a $`k`$-module such that for each $`a_{j_1\mathrm{}j_N}^{\alpha _1\mathrm{}\alpha _N}S^{}`$, we have
$$a_{j_1\mathrm{}j_N}^{\alpha _1\mathrm{}\alpha _N}(z)=:^{j_1}\alpha ^{\alpha _1}(z)\mathrm{}^{j_N}a^{\alpha _N}(z):$$
$`(\mathrm{0.7.2})`$
Then the mapping (\*) defines the structure of a $`_0`$-graded vertex algebra on $`V`$, and (0.7.2) holds true for all $`a_{j_1\mathrm{}j_N}^{\alpha _1\mathrm{}\alpha _N}S`$.
This is a version of ”Existence Theorem”, \[K\], Theorem 4.5 and Corollary 4.5, whose proof goes through.
0.8. Let $`𝒱ert,𝒞onf`$ denote the categories of $`_0`$-graded vertex and conformal algebras respectively. We have a functor
$$c:𝒱ert𝒞onf$$
$`(\mathrm{0.8.1})`$
which assigns to a vertex algebra $`V`$ the same space $`V`$, with the operations <sub>(n)</sub>, $`n<0,`$ forgotten. The axioms of a conformal algebra are satisfied due to (0.5.10), (0.5.9) and (0.5.12). This functor admits a left adjoint, to be constructed in Section 8 below.
In the sequel we will call $`_0`$-graded vertex (conformal) algebras simply vertex (conformal) algebras.
0.9. A vertex algebra $`V`$ is called commutative if $`a_{(n)}b=0`$ for all $`a,bV,n0`$ (Kac uses the term ”holomorphic”).
Let $`V`$ be a commutative vertex algebra. Then, with respect to the operation $`ab:=a_{(1)}b`$, $`V`$ becomes a commutative associative algebra with the unity 1. The operations with negative indices are recovered from the formula
$$a_{(1j)}b=(^{(j)}a)b$$
$`(\mathrm{0.9.1})`$
This way we get an equivalence of categories
$$𝒱ert𝒜b\stackrel{}{}𝒜lg$$
$`(\mathrm{0.9.2})`$
Here $`𝒱ert𝒜b`$ denotes the category of commutative vertex algebras and $`𝒜lg`$ denotes the category whose objects are $`_0`$-graded vector spaces $`V=V_i`$ equipped with a structure of a commutative algebra such that $`V_iV_jV_{i+j}`$ and $`\text{1}V_0`$, and with a family of endomorphisms $`^{(j)}`$ of degree $`j`$, $`j_0`$, such that
$$^{(i)}^{(j)}=\left(\genfrac{}{}{0pt}{}{i+j}{i}\right)^{(i+j)},^{(0)}=Id$$
$`(\mathrm{0.9.3})`$
and
$$^{(j)}(ab)=\underset{p=0}{\overset{j}{}}^{(p)}a^{(jp)}b$$
$`(\mathrm{0.9.4})`$
Cf. \[B\]. Objects of $`𝒜lg`$ will be called $``$-algebras.
§ 1. Vertex Algebroids
1.1. Let us define an extended Lie algebroid to be a quintuple $`𝒯=(A,T,\mathrm{\Omega },,,)`$ where $`A`$ is a commutative $`k`$-algebra, $`T`$ is a Lie $`A`$-algebroid, $`\mathrm{\Omega }`$ is an $`A`$-module equipped with a structure of a module over the Lie algebra $`T`$, $`:A\mathrm{\Omega }`$ is an $`A`$-derivation and a morphism of $`T`$-modules, $`,:T\times \mathrm{\Omega }A`$ is an $`A`$-bilinear pairing.
These data must satisfy the following properties $`(aA,\tau ,\nu T,\omega \mathrm{\Omega })`$:
$$\tau ,a=\tau (a)$$
$`(\mathrm{1.1.1})`$
$$\tau (a\omega )=\tau (a)\omega +a\tau (\omega )$$
$`(\mathrm{1.1.2})`$
$$(a\tau )(\omega )=a\tau (\omega )+\tau ,\omega a$$
$`(\mathrm{1.1.3})`$
$$\tau (\nu ,\omega )=[\tau ,\nu ],\omega +\nu ,\tau (\omega )$$
$`(\mathrm{1.1.4})`$
Cf. 0.1. We will also say that $`𝒯=(A,T,\mathrm{})`$ is an extended Lie $`A`$-algebroid.
Let define a morphism between two extended Lie algebroids $`𝒯=(A,T,\mathrm{})`$ and $`𝒯^{}=(A^{},T^{},\mathrm{})`$ to be a triple $`g=(g_A,g_T,g_\mathrm{\Omega })`$ where $`g_A:AA^{}`$ is a morphism of $`k`$-algebras, $`g_T:TT^{}`$ is a morphism of Lie algebras and $`A`$-modules, $`g_\mathrm{\Omega }:\mathrm{\Omega }\mathrm{\Omega }^{}`$ is a morphism of $`A`$-modules. We require that the following properties should hold:
$$g_A(\tau (a))=g_T(\tau )(g_A(a))$$
$`(\mathrm{1.1.5})`$
$$g_\mathrm{\Omega }(a)=g_A(a)$$
$`(\mathrm{1.1.6})`$
$$g_A(\tau ,\omega =g_T(\tau ),g_\mathrm{\Omega }(\omega )$$
$`(\mathrm{1.1.7})`$
Composition of morphisms is defined in the obvious way. This way we get a category $`ieAlg`$ of extended Lie algebroids.
1.2. Let $`T`$ be a Lie $`A`$-algebroid. Set $`\mathrm{\Omega }:=Hom_A(T,A)`$. Define $`:A\mathrm{\Omega }`$ by (1.1.1); let $`,`$ be the evident pairing. Action of $`T`$ on $`\mathrm{\Omega }`$ is defined by (1.1.4). This way we get an extended Lie $`A`$-algebroid $`𝒯_T`$.
Let us call an extended Lie algebroid $`𝒯=(A,T,\mathrm{\Omega },,)`$ perfect if the pairing $`,`$ induces an isomorphism $`\mathrm{\Omega }\stackrel{}{}Hom_A(T,A)`$.
The correspondence $`T𝒯_T`$ provides an equivalence of the category of Lie algebroids with the full subcategory $`ieAlg^{perf}ieAlg`$ of perfect extended Lie algebroids.
1.3. De Rham - Chevalley complex. Let $`𝒯=(A,T,\mathrm{\Omega },\mathrm{})`$ be an extended Lie $`A`$-algebroid. Let us define $`A`$-modules $`\mathrm{\Omega }^i=\mathrm{\Omega }^i(𝒯),i_0`$, as follows. Set $`\mathrm{\Omega }^0=A,\mathrm{\Omega }^1=\mathrm{\Omega }`$. For $`i2`$, $`\mathrm{\Omega }^i`$ is the submodule of the module of $`A`$-polylinear homomorphisms $`h`$ from $`T^{i1}`$ to $`\mathrm{\Omega }`$ such that the function $`\tau _1,h(\tau _2,\mathrm{},\tau _i)`$ is skew symmetric with respect to all permutations of $`(\tau _1,\mathrm{},\tau _i)`$.
For example, if $`𝒯`$ is as in the previous example, then $`\mathrm{\Omega }^i=Hom_A(\mathrm{\Lambda }_A^iT,A)`$.
Let us define the maps $`d_{DR}=d_{DR}^i:\mathrm{\Omega }^i\mathrm{\Omega }^{i+1}`$ as follows. For $`i=0`$ we set $`d_{DR}a=a`$. For $`i1`$ we set
$$d_{DR}h(\tau _1,\mathrm{},\tau _i)=d_{Lie}h(\tau _1,\mathrm{},\tau _i)\tau _1,h(\tau _2,\mathrm{},\tau _i)$$
$`(\mathrm{1.3.1})`$
where
$$d_{Lie}h(\tau _1,\mathrm{},\tau _i)=\underset{p=1}{\overset{i}{}}(1)^{p+1}\tau _p(h(\tau _1,\mathrm{},\widehat{\tau }_p,\mathrm{},\tau _i))+$$
$$+\underset{1p<qi}{}(1)^{p+q}h([\tau _p,\tau _q],\tau _1,\mathrm{},\widehat{\tau }_p,\mathrm{},\widehat{\tau }_q,\mathrm{},\tau _i)$$
$`(\mathrm{1.3.2})`$
For example,
$$d_{DR}\omega (\tau )=\tau (\omega )\tau ,\omega ,$$
$`(\mathrm{1.3.3})`$
for $`\omega \mathrm{\Omega }^1=\mathrm{\Omega }`$; and
$$d_{DR}h(\tau _1,\tau _2)=h([\tau _1,\tau _2])+\tau _1(h(\tau _2))\tau _2(h(\tau _1))\tau _1,h(\tau _2),$$
$`(\mathrm{1.3.4})`$
for $`h\mathrm{\Omega }^2`$.
Let us introduce the action of the Lie algebra $`T`$ on the modules $`\mathrm{\Omega }^i`$ by
$$\tau (h)(\tau _1,\mathrm{},\tau _{i1})=\tau (h(\tau _1,\mathrm{},\tau _{i1}))\underset{p=1}{\overset{i1}{}}h(\tau _1,\mathrm{},[\tau ,\tau _p],\mathrm{},\tau _i)$$
$`(\mathrm{1.3.5})`$
Let us define the convolution operators $`\tau ,:\mathrm{\Omega }^i\mathrm{\Omega }^{i1}`$ by
$$\tau ,h(\tau _1,\mathrm{},\tau _{i2}=h(\tau ,\tau _1,\mathrm{},\tau _{i2})$$
$`(\mathrm{1.3.6})`$
The maps $`\{d_{DR}^i\}`$ may be characterized as a unique collection of maps such that $`d_{DR}^0=`$ and the Cartan formula
$$\tau (h)=\tau ,d_{DR}h+d_{DR}\tau ,h$$
$`(\mathrm{1.3.7})`$
holds true.
The maps $`d_{DR}`$ commute with the action of $`T`$. One checks that $`d_{DR}^2=0`$, so we get a complex $`(\mathrm{\Omega }^{}(𝒯),d_{DR})`$ called the de Rham-Chevalley complex of $`𝒯`$.
1.4. In the definition below we assume that $`1/2k`$.
A vertex algebroid is a septuple $`𝒜=(A,T,\mathrm{\Omega },,\gamma ,,,c)`$ where $`A`$ is a commutative $`k`$-algebra, $`T`$ is a Lie $`A`$-algebroid, $`\mathrm{\Omega }`$ is an $`A`$-module equipped with an action of the Lie algebra $`T`$, $`:A\mathrm{\Omega }`$ is a derivation commuting with the $`T`$-action,
$$,:(T\mathrm{\Omega })\times (T\mathrm{\Omega })A$$
is a symmetric $`k`$-bilinear pairing equal to zero on $`\mathrm{\Omega }\times \mathrm{\Omega }`$ and such that $`𝒯_𝒜=(A,T,\mathrm{\Omega },,,|_{T\times \mathrm{\Omega }})`$ is an extended Lie $`A`$-algebroid; $`c:T\times T\mathrm{\Omega }`$ is a skew symmetric $`k`$-bilinear pairing and $`\gamma :A\times T\mathrm{\Omega }`$ is a $`k`$-bilinear map.
The following axioms must hold $`(a,bA;\tau ,\tau _iT)`$:
$$\gamma (a,b\tau )=\gamma (ab,\tau )a\gamma (b,\tau )\tau (a)b\tau (b)a$$
$`(A1)`$
$$a\tau _1,\tau _2=a\tau _1,\tau _2+\gamma (a,\tau _1),\tau _2\tau _1\tau _2(a)$$
$`(A2)`$
$$c(a\tau _1,\tau _2)=ac(\tau _1,\tau _2)+\gamma (a,[\tau _1,\tau _2])\gamma (\tau _2(a),\tau _1)+\tau _2(\gamma (a,\tau _1))$$
$$\frac{1}{2}\tau _1,\tau _2a+\frac{1}{2}\tau _1\tau _2(a)\frac{1}{2}\tau _2,\gamma (a,\tau _1)$$
$`(A3)`$
$$[\tau _1,\tau _2],\tau _3+\tau _2,[\tau _1,\tau _3]=\tau _1(\tau _2,\tau _3)\frac{1}{2}\tau _2(\tau _1,\tau _3)\frac{1}{2}\tau _3(\tau _1,\tau _2)+$$
$$+\tau _2,c(\tau _1,\tau _3)+\tau _3,c(\tau _1,\tau _2)$$
$`(A4)`$
$$d_{Lie}c(\tau _1,\tau _2,\tau _3)=\frac{1}{2}\{[\tau _1,\tau _2],\tau _3+[\tau _1,\tau _3],\tau _2[\tau _2,\tau _3],\tau _1$$
$$\tau _1(\tau _2,\tau _3)+\tau _2(\tau _1,\tau _3)2\tau _3,c(\tau _1,\tau _2)\}$$
$`(A5)`$
where $`d_{Lie}`$ is defined by (1.3.2).
To stress the dependence on $`A`$, we shall sometimes say that $`𝒜`$ is a vertex $`A`$-algebroid.
1.5. Multiplying the identity (A4) by $`l`$ and subtracting the $`j`$-th multiple of (A4) corresponding to the triple $`(\tau _2,\tau _1,\tau _3)`$ we get an equivalent form of (A4):
for every $`l,jk`$,
$$j\tau _1,[\tau _2,\tau _3]+l\tau _2,[\tau _1,\tau _3]+(l+j)\tau _3,[\tau _1,\tau _2]$$
$$(l+\frac{1}{2}j)\tau _1(\tau _2,\tau _3)+(\frac{1}{2}l+j)\tau _2(\tau _1,\tau _3)+\frac{1}{2}(lj)\tau _3(\tau _1,\tau _2)$$
$`(A4)_{l,j}^{bis}`$
$$+j\tau _1,c(\tau _2,\tau _3)l\tau _2,c(\tau _1,\tau _3)(l+j)\tau _3,c(\tau _1,\tau _2)=0$$
1.6. The left hand side of (A5) is skew symmetric with respect to all permutations of $`(\tau _1,\tau _2,\tau _3)`$. The right hand side is manifestly symmetric only with respect to the transposition of $`(\tau _1,\tau _2)`$. However, if we replace $`\tau _1(\tau _2,\tau _3)`$ in the right hand side by its expression from (A4) we get the following equivalent form of (A5):
$$d_{Lie}c(\tau _1,\tau _2,\tau _3)=\frac{1}{2}\{\tau _1,[\tau _2,\tau _3]+\frac{1}{2}\tau _2(\tau _1,\tau _3)\frac{1}{2}\tau _3(\tau _1,\tau _2)+$$
$$+\tau _2,c(\tau _1,\tau _3)\tau _3,c(\tau _1,\tau _2)\}$$
$`(A5)^{bis}`$
The rhs of $`(A5)^{bis}`$ is skew symmetric with respect to the transposition of $`(\tau _2,\tau _3)`$. Consequently, the rhs of (A5) is completely skew symmetric since the symmetric group $`S_3`$ is generated by two transpositions $`(12)`$ and $`(23)`$.
Let us replace in $`(A5)^{bis}`$ the triple $`(\tau _1,\tau _2,\tau _3)`$ by the triples $`(\tau _2,\tau _3,\tau _1)`$ and $`(\tau _3,\tau _1,\tau _2)`$ and sum up the three identities. We will get another equivalent for of (A5):
$$3d_{Lie}c(\tau _1,\tau _2,\tau _3)=\{\tau _1,\frac{1}{2}[\tau _2,\tau _3]+c(\tau _2,\tau _3)+$$
$$+\tau _2,\frac{1}{2}[\tau _3,\tau _1]+c(\tau _3,\tau _1)+\tau _3,\frac{1}{2}[\tau _1,\tau _2]+c(\tau _1,\tau _2)\}$$
$`(A5)^{ter}`$
1.7. Example. Let $`𝔤`$ be a Lie algebra over $`k`$ equipped with a symmetric invariant form $`,:𝔤\times 𝔤k`$. Then $`𝒜_{𝔤,,}=(k,𝔤,0,0,0,,,0)`$ is a vertex $`k`$-algebroid.
1.8. Let $`A`$ be a commutative $`k`$-algebra and $`𝒯=(A,T,\mathrm{\Omega },,,)`$ be an extended Lie $`A`$-algebroid. Let $`T_0T`$ be a $`k`$-submodule which generates $`T`$ as an $`A`$-module, and assume that we are given $`k`$-bilinear mappings
$$\gamma :A\times T_0\mathrm{\Omega };,:T_0\times T_0A,\text{symmetric};$$
$$c:T_0\times T_0\mathrm{\Omega },\text{skew symmetric}$$
$`(\mathrm{1.8.1})`$
It is clear that there exists not more than one extension of the mappings to mappings
$$\gamma :A\times T\mathrm{\Omega };,:T\times TA,c:T\times T\mathrm{\Omega },$$
$`(\mathrm{1.8.2})`$
satisfying (A1), (A2), (A3). This extension, if it exists, must be given by the formulas $`(a,bA,\tau ,\tau _iT_0)`$:
$$\gamma (a,b\tau )=\gamma (ab,\tau )a\gamma (b,\tau )\tau (a)b\tau (b)a$$
$`(\mathrm{1.8.3})_\gamma `$
$$a\tau _1,b\tau _2=ab\tau _1,\tau _2+a\tau _1,\gamma (b,\tau _2)+b\tau _2,\gamma (a,\tau _1)$$
$$a\tau _2\tau _1(b)b\tau _1\tau _2(a)\tau _1(b)\tau _2(a)$$
$`(\mathrm{1.8.3})_,`$
$$c(a\tau _1,b\tau _2)=abc(\tau _1,\tau _2)+\gamma (ab,[\tau _1,\tau _2])+$$
$$+\gamma (a\tau _1(b),\tau _2)\gamma (b\tau _2(a),\tau _1)a\tau _1(\gamma (b,\tau _2))+b\tau _2(\gamma (a,\tau _1))+$$
$$+\frac{1}{2}\tau _1,\tau _2\left\{abba\right\}[\tau _1,\tau _2](a)b[\tau _1,\tau _2](b)a+\frac{1}{2}\left\{\tau _1(b)\tau _2(a)\tau _2(a)\tau _1(b)\right\}+$$
$$+\frac{1}{2}\left\{\tau _1,\gamma (b,\tau _2)a+\tau _2,\gamma (a,\tau _1)b+a\tau _1,\gamma (b,\tau _2)b\tau _2,\gamma (a,\tau _1)\right\}+$$
$$+\frac{1}{2}\left\{b\tau _1\tau _2(a)a\tau _2\tau _1(b)\right\}$$
$`(\mathrm{1.8.3})_c`$
When deducing the last formula, one should take into account (1.1.3).
1.9. Theorem. (Extension of Identities) Assume that (A4) and (A5) hold true for all $`\tau _iT_0`$ and that the formulas (1.8.3) provide well defined mappings (1.8.2).
Then the axioms (A1) — (A5) hold true for all $`\tau ,\tau _iT`$, i.e. $`𝒜=(A,T,\mathrm{})`$ is a vertex $`A`$-algebroid.
Proof. Let us check the axiom (A1). It is enough to show that if (A1) holds true for some $`\tau `$ and all $`a,b`$ then it holds true for $`c\tau (cA)`$ and all $`a,b`$. Thus, we have to check the identity
$$\gamma (ab,c\tau )=\gamma (a,bc\tau )+a\gamma (b,c\tau )+c\tau (a)b+c\tau (b)a$$
$`(\mathrm{?})`$
The left hand side of it equal to
$$\gamma (abc,\tau )ab\gamma (c,\tau )\tau (ab)c\tau (c)(ab),$$
cf. $`(\mathrm{1.8.3})_\gamma `$. On the other hand, the first two terms in the right hand side are equal to
$$\gamma (a,bc\tau )=\gamma (abc,\tau )a\gamma (bc,\tau )\tau (a)(bc)\tau (bc)a$$
and
$$a\gamma (b,c\tau )=a\gamma (bc,\tau )ab\gamma (c,\tau )a\tau (b)ca\tau (c)b,$$
again by $`(\mathrm{1.8.3})_\gamma `$. Comparing, we get the desired identity (?).
The other axioms are checked in a similar way. This is a tiresome but straightforward calculation. When checking (A5) one should use the identity $`(A4)_{1,1}^{bis}`$. $``$
1.10. Pushout. Let $`𝒜=(A,T,\mathrm{\Omega },\mathrm{})`$ be a vertex $`A`$-algebroid. Let $`B`$ be a commutative $`A`$-algebra, $`i:AB`$ the structure morphism. Set $`\mathrm{\Omega }_B:=B_A\mathrm{\Omega },T_B:=BT`$. The $`A`$-derivation $`:A\mathrm{\Omega }`$ induces a $`B`$-derivation $`_B:B\mathrm{\Omega }_B`$. The $`A`$-bilinear pairing $`,:T\times \mathrm{\Omega }A`$ uniquely extends to a $`B`$-bilinear pairing $`,_B:T_B\times \mathrm{\Omega }_BB`$.
Assume that the Lie algebra $`T`$ acts on $`B`$ by derivations in such a way that $`\tau (i(a))=i(\tau (a))`$ and $`(a\tau )(b)=a\tau (b)(aA,bB,\tau T)`$. Then $`T_B`$ acquires a canonical structure of a Lie $`B`$-algebroid, cf. 0.2.1, and $`(T_B,\mathrm{\Omega }_B,_B,,_B)`$ becomes an extended Lie $`B`$-algebroid.
1.10.1. Theorem. Assume that we are given a $`k`$-bilinear mapping $`\gamma :B\times T\mathrm{\Omega }_B`$ such that $`\gamma (i(a),\tau )=1\gamma (a,\tau )`$ and that (A1) holds true for all $`\tau T,aB,bA`$.
Then there exists a unique extension of $`\gamma `$ to a $`k`$-bilinear mapping $`\gamma _B:B\times T_B\mathrm{\Omega }_B`$ satisfying (A1) for all $`a,bB,\tau T_B`$;
there exists a unique extension of the pairing $`,:T\times TA`$ to a pairing $`,_B:T_B\times T_BB`$ satisfying (A2) for all $`aB,\tau _iT_B`$;
there exists a unique extension of the pairing $`c:T\times T\mathrm{\Omega }`$ to a pairing $`c_B:T_B\times T_B\mathrm{\Omega }_B`$ satisfying (A3) for all $`aB,\tau _iT_B`$.
The septuple $`𝒜_B=(B,T_B,\mathrm{\Omega }_B,_B,\gamma _B,,_B,c_B)`$ is a vertex $`B`$-algebroid.
Proof. Apply 1.9 to $`T_0:=Im(TT_B)`$. $``$
1.11. Example. Let $`A`$ be a commutative $`k`$-algebra, $`\mathrm{\Omega }=\mathrm{\Omega }_k^1`$ — the $`A`$-module of Kählerian $`1`$-differentials, $`:A\mathrm{\Omega }`$ the canonical $`A`$-derivation. Let $`𝔤`$ be a Lie $`k`$-algebra acting on $`A`$ by derivations and equipped with an invariant bilinear form $`,:𝔤\times 𝔤k`$.
Due to the morphism $`𝔤Der_k(A)`$, $`\mathrm{\Omega }`$ is equipped with a canonical structure of a $`𝔤`$-module such that $``$ is a morphism of $`𝔤`$-modules and there is a canonical pairing $`𝔤\times \mathrm{\Omega }A`$, cf. 0.1.
Set $`T=A𝔤`$. Then $`T`$ is canonically a Lie $`A`$-algebroid, there exists a unique extension of the zero map $`A\times 𝔤A`$ (resp., of $`,`$ and of the zero map $`𝔤\times 𝔤\mathrm{\Omega }`$) to the map $`\gamma :A\times T\mathrm{\Omega }`$ (resp., to the pairing $`,:T\times TA`$ and to the map $`c:T\times T\mathrm{\Omega }`$) which satisfies (A1) (resp., (A2) and (A3)).
This way we get a vertex $`A`$-algebroid $`𝒜=(A,T,\mathrm{\Omega },,\gamma ,,,c)`$.
For the proof, apply Theorem 1.9 to $`T_0=𝔤T`$.
1.12. Example. Let $`A`$ be a commutative $`k`$ algebra, $`𝔤T=Der_k(A)`$ a Lie $`k`$-subalgebra such that $`T=A𝔤`$, $`\mathrm{\Omega }=\mathrm{\Omega }_k^1`$.
For example, $`A`$ is the coordinate ring of an algebraic group, $`𝔤`$ is the Lie algebra of left invariant vector fields.
Assume that $`𝔤`$ is equipped with an invariant bilinear form $`,`$. Then we are in the situation of 1.11 and get a vertex $`A`$-algebroid $`𝒜=𝒜_{A,𝔤,,}`$.
1.13. Example. Let $`A`$ be a commutative $`k`$-algebra equipped with an étale morphism $`\text{f}=(f_1,\mathrm{},f_n):A_0=k[x_1,\mathrm{},x_n]A`$. Such an f is called an étale coordinate system on $`A`$; it exists Zariski locally for every $`A`$ smooth over $`k`$.
The commuting vector fields $`/x_iDer_k(A_0)`$ admit unique liftings to vector fields $`\tau _iT=Der_k(A)`$. The $`k`$-submodule $`𝔤T`$ spanned by $`\tau _i`$ is an abelian Lie subalgebra; let us equip it with the trivial bilinear form. Then we are in the situation of 1.11 and we get a vertex $`A`$-algebroid $`𝒜=𝒜_{A;\text{f}}`$.
§2. From Vertex Algebras to Vertex Algebroids
2.1. Let $`V`$ be a vertex algebra over $`k`$, cf. 0.5. In this section we assume that $`1/2k`$; recall that we always deal with $`_0`$-graded vertex algebras. We denote $`:=^{(1)}`$.
The identity (0.5.4) at $`n=0`$ gives
$$a_{(0)}b_{(n)}c=(a_{(0)}b)_{(n)}c+b_{(n)}a_{(0)}c$$
$`(\mathrm{2.1.1})`$
for all $`n`$, i.e. the operation <sub>(0)</sub> is a derivation with respect to all operations <sub>(n)</sub>.
We shall denote the operation $`a_{(1)}b`$ by $`ab`$ or $`ab`$.
Set $`A=V_0`$. In the sequel $`a,b,c`$ will denote elements of $`A`$. Since $`a_{(n)}b=0`$ for all $`n0`$, (0.5.4) and (0.5.9) imply that $`A`$ is a commutative associative $`k`$-algebra with the unity 1 with respect to the operation $`ab`$.
Elements of $`V_1`$ will be denoted $`x,y,z`$. It follows from (0.5.4), taking into account that $`\left(\genfrac{}{}{0pt}{}{1}{j}\right)=(1)^j`$, that
$$(ab)x=abx+ab_{(0)}x+ba_{(0)}x$$
$`(\mathrm{2.1.2})`$
The identity (0.5.9) implies
$$a_{(0)}x=x_{(0)}a$$
$`(\mathrm{2.1.3})`$
and
$$xa=ax+(x_{(0)}a)$$
$`(\mathrm{2.1.4})`$
Let $`\mathrm{\Omega }V_1`$ be the $`k`$-submodule generated by all elements $`ab`$; $`\omega `$ will denote an element of $`\mathrm{\Omega }`$.
From (2.1.2) we have
$$(ab)c=abc+ab_{(0)}c+ba_{(0)}c$$
Note that
$$u_{(0)}v=0$$
$`(\mathrm{2.1.5})`$
for all $`u,vV`$, due to (0.5.10). Hence $`b_{(0)}c=c_{(0)}b=0`$; therefore $`(ab)c=abc`$. On the other hand, $`\text{1}c=c\text{1}c_{(0)}\text{1}=c\text{1}=c`$. It follows that the operation $`a\omega `$ gives a structure of an $`A`$-module on $`\mathrm{\Omega }`$. (It is not true for $`V_1`$ since $`ax`$ is not associative.)
Note that by (0.5.4) $`(ab)_{(0)}c=ab_{(0)}c=0`$, in other words
$$\omega _{(0)}a=0$$
$`(\mathrm{2.1.6})`$
It follows from (0.5.9) that
$$a\omega =\omega a$$
$`(\mathrm{2.1.7})`$
We have the map $`:A\mathrm{\Omega }`$. By (0.5.11) and (2.1.7) $``$ is an $`A`$-derivation.
Let $`T`$ denote the quotient module $`V_1/\mathrm{\Omega }`$; let $`\pi :V_1T`$ be the canonical projection. Elements of $`T`$ will be denoted $`\tau `$. The operation $`ax`$ induces an operation $`A\times TT,(a,\tau )a\tau `$. By (2.1.2) $`(ab)\tau =ab\tau `$. On the other hand,
$$\text{1}x=x\text{1}(x_{(0)}\text{1})=x\text{1}=x$$
$`(\mathrm{2.1.8})`$
i.e. the operation $`a\tau `$ provides $`T`$ with a structure of an $`A`$-module.
By (0.5.4)
$$(ab)_{(0)}x=ab_{(0)}xa_{(2)}b_{(1)}x+ba_{(0)}x$$
Note that by (0.5.10) and (2.1.7)
$$a_{(2)}b=ab=ba$$
$`(\mathrm{2.1.9})`$
Therefore
$$(ab)_{(0)}x=(b_{(0)}x)a+(a_{(0}x)b$$
whence
$$\mathrm{\Omega }_{(0)}V_1\mathrm{\Omega }$$
$`(\mathrm{2.1.10})`$
On the other hand,
$$x_{(0)}ab=(x_{(0)}a)b+ax_{(0)}b$$
by (2.1.1), and
$$x_{(0)}b=b_{(0)}x+(b_{(1)}x)=(b_{(0)}x)=(x_{(0)}b),$$
i.e.
$$x_{(0)}(ab)=(x_{(0)}a)b+a(x_{(0)}b)$$
$`(\mathrm{2.1.11})`$
It follows that
$$V_{1(0)}\mathrm{\Omega }\mathrm{\Omega }$$
$`(\mathrm{2.1.12})`$
Therefore the operation <sub>(0)</sub> induces an operation $`T\times TT`$, to be denoted $`[,]`$. By (0.5.9)
$$x_{(0)}y=y_{(0)}x+(y_{(1)}x)$$
$`(\mathrm{2.1.13})`$
hence $`[,]`$ is skew symmetric. (2.1.1) implies that $`[,]`$ satisfies the Jacobi identity, hence it provides a structure of a Lie algebra (over $`k`$) on $`T`$.
By (0.5.4) and (2.1.5) $`(ab)_{(0)}c=ab_{(0)}c=0`$, i.e.
$$\mathrm{\Omega }_{(0)}A=0$$
$`(\mathrm{2.1.14})`$
It follows that the operation $`x_{(0)}a`$ induces a pairing $`T\times AA`$, to be denoted $`\tau (a)`$. By (2.1.1), $`\tau _1\tau _2(a)=[\tau _1,\tau _2](a)+\tau _2\tau _1(a)`$ and $`\tau (ab)=\tau (a)b+a\tau (b)`$. By (0.5.4), $`(a\tau )(b)=a\tau (b)`$. Thus, $`T`$ is equipped with a structure of a Lie $`A`$-algebroid.
By (2.1.11) and (2.1.14)
$$\mathrm{\Omega }_{(0)}\mathrm{\Omega }=0$$
$`(\mathrm{2.1.15})`$
It follows that the operation $`x_{(0)}\omega `$ induces an operation $`T\times \mathrm{\Omega }\mathrm{\Omega }`$, to be denoted $`\tau (\omega )`$, which makes $`\mathrm{\Omega }`$ a module over the Lie algebra $`T`$, by (2.1.1). Again by (2.1.1), $`\tau (a\omega )=\tau (a)\omega +a\tau (\omega )`$.
It follows from (2.1.11) that
$$\tau (ab)=\tau (a)b+a\tau (b)$$
$`(\mathrm{2.1.16})`$
In particular $`\tau (b)=\tau (b)`$, i.e. $``$ is a morphism of $`T`$-modules.
By (0.5.9)
$$x_{(1)}y=y_{(1)}x$$
$`(\mathrm{2.1.17})`$
We have
$$(ab)_{(1)}x=ab_{(1)}x=ab_{(0)}x=ax_{(0)}b$$
$`(\mathrm{2.1.18})`$
by (0.5.4), (2.1.5), (0.5.10) and (2.1.3). It follows that
$$\mathrm{\Omega }_{(1)}\mathrm{\Omega }=0$$
$`(\mathrm{2.1.19})`$
by (2.1.14). Therefore the operation $`x_{(1)}y`$ induces an operation $`T\times \mathrm{\Omega }A`$, to be denoted $`\tau ,\omega `$.
By (2.1.17) and (2.1.18)
$$\tau ,ab=a\tau (b)$$
$`(\mathrm{2.1.20})`$
From (0.5.4), $`a\tau ,\omega =a\tau ,\omega `$. From (2.1.1),
$$\tau (\tau ^{},\omega )=[\tau ,\tau ^{}],\omega +\tau ^{},\tau (\omega )$$
Finally, again from (0.5.4), $`(ax)_{(0)}\omega =ax_{(0)}\omega +a_{(2)}x_{(1)}\omega `$ which implies $`(a\tau )(\omega )=a\tau (\omega )+\tau ,\omega a`$, by (2.1.9).
Therefore we have canonically associated with our vertex algebra an extended Lie $`A`$-algebroid $`𝒯(V)=(T,\mathrm{\Omega },,,)`$.
2.2. Let us assume that the projection $`\pi :V_1T`$ admits a splitting, i.e. there exists a morphism of $`k`$-modules $`s:TV_1`$ such that $`\pi s=Id_T`$. Let us fix such a splitting. A vertex algebra with a chosen $`s`$ will be called split.
Define a skew symmetric operation $`[,]:V_1\times V_1V_1`$ by
$$[x,y]:=\frac{1}{2}(x_{(0)}yy_{(0)}x)$$
$`(\mathrm{2.2.1})`$
We have $`\pi ([x,y])=[\pi (x),\pi (y)]`$.
Let us use the notation $`x,y`$ for $`x_{(1)}y`$. It follows form (0.5.9) that
$$[x,y]=x_{(0)}y\frac{1}{2}x,y$$
$`(\mathrm{2.2.2})`$
Set
$$\gamma (a,\tau )=s(a\tau )as(\tau )$$
$`(\mathrm{2.2.3})_\gamma `$
$$\tau _1,\tau _2=s(\tau _1),s(\tau _2)$$
$`(\mathrm{2.2.3})_,`$
$$c(\tau _1,\tau _2)=s([\tau _1,\tau _2])[s(\tau _1),s(\tau _2)]$$
$`(\mathrm{2.2.3})_c`$
2.3. Theorem. The septuple $`𝒜=(A,T,\mathrm{\Omega },,\gamma ,,,c)`$ is a vertex algebroid.
Proof. We have to check the axioms (A1) — (A5) from 1.4.
Check of (A1). We have $`\gamma (ab,\tau )=s(ab\tau )(ab)s(\tau )=s(ab\tau )as(b\tau )+as(b\tau )(ab)s(\tau )`$. By (2.1.2) and (2.1.3),
$$(ab)s(\tau )=abs(\tau )+ab_{(0)}s(\tau )+ba_{(0)}s(\tau )=abs(\tau )\tau (b)a\tau (a)b$$
which implies (A1).
Check of (A2). We have $`a\tau _1,\tau _2=s(a\tau _1),s(\tau _2)=as(\tau _1)+\gamma (a,\tau _1),s(\tau _2)`$. By (0.5.4),
$$ax,y=ax,yx_{(0)}y_{(0)}a$$
$`(\mathrm{2.3.1})`$
which implies (A2).
Check of (A3). We have $`c(a\tau _1,\tau _2)=s([a\tau _1,\tau _2])[s(a\tau _1),s(\tau _2)]`$;
$$s([a\tau _1,\tau _2])=s(a[\tau _1,\tau _2]\tau _2(a)\tau _1)=as([\tau _1,\tau _2])+\gamma (a,[\tau _1,\tau _2])\tau _2(a)s(\tau _1)\gamma (\tau _2(a),\tau _1);$$
$`[s(a\tau _1),s(\tau _2)]=[as(\tau _1)+\gamma (a,\tau _1),s(\tau _2)]`$.
We have $`[ax,y]=(ax)_{(0)}y\frac{1}{2}ax,y`$ (see (2.2.2)). By (0.5.4),
$$(ax)_{(0)}y=ax_{(0)}y+a_{(2)}x_{(1)}y+xa_{(0)}y=$$
$$=ax_{(0)}y+x,ya\pi (y)(a)x\pi (x)\pi (y)(a)$$
$`(\mathrm{2.3.2})`$
By (2.3.1)
$$\frac{1}{2}ax,y=\frac{1}{2}\left\{ax,yx_{(0)}y_{(0)}a\right\}=\frac{1}{2}x,ya\frac{1}{2}ax,y+\frac{1}{2}\pi (x)\pi (y)(a)$$
Therefore
$$[ax,y]=a[x,y]+\frac{1}{2}x,ya\frac{1}{2}\pi (x)\pi (y)(a)\pi (y)(a)x$$
On the other hand,
$$[\omega ,x]=[x,\omega ]=\pi (x)\omega +\frac{1}{2}x,\omega $$
It follows that
$$[s(a\tau _1),s(\tau _2)]=a[s(\tau _1),s(\tau _2)]+\frac{1}{2}\tau _1,\tau _2a\frac{1}{2}\tau _1\tau _2(a)\tau _2(a)s(\tau _1)\tau _2(\gamma (a,\tau _1))+$$
$`+\frac{1}{2}\gamma (a,\tau _1),\tau _2`$. This implies (A3).
Check of (A4). We shall make use of the formula
$$s(\tau _1)_{(0)}s(\tau _2),s(\tau _3)=s(\tau _1)_{(0)}s(\tau _2),s(\tau _3)+s(\tau _2),s(\tau _1)_{(0)}s(\tau _3)$$
$`(\mathrm{2.3.3})`$
which is a particular case of (2.1.1). We have $`s(\tau _1)_{(0)}s(\tau _2),s(\tau _3)=\tau _1(\tau _2,\tau _3)`$. On the other hand,
$$s(\tau _1)_{(0)}s(\tau _2)=[s(\tau _1),s(\tau _2)]+\frac{1}{2}\tau _1,\tau _2=s([\tau _1,\tau _2])c(\tau _1,\tau _2)+\frac{1}{2}\tau _1,\tau _2$$
$`(\mathrm{2.3.4})`$
Therefore
$$s(\tau _1)_{(0)}s(\tau _2),s(\tau _3)=[\tau _1,\tau _2],\tau _3c(\tau _1,\tau _2),\tau _3+\frac{1}{2}\tau _3(\tau _1,\tau _2)$$
Interchanging $`\tau _2`$ with $`\tau _3`$ and plugging these expressions into (2.3.3) we get (A4).
Check of (A5). We shall use the formulas
$$s(\tau _1)_{(0)}s(\tau _2)_{(0)}s(\tau _3)=(s(\tau _1)_{(0)}s(\tau _2))_{(0)}s(\tau _3)+s(\tau _2)_{(0)}s(\tau _1)_{(0)}s(\tau _3)$$
and (2.3.4). We get
$$s(\tau _1)_{(0)}s(\tau _2)_{(0)}s(\tau _3)=s(\tau _1)_{(0)}\left\{s([\tau _2,\tau _3])c(\tau _2,\tau _3)+\frac{1}{2}\tau _2,\tau _3\right\}=$$
$$=s([\tau _1,[\tau _2,\tau _3])c(\tau _1,[\tau _2,\tau _3])+\frac{1}{2}\tau _1,[\tau _2,\tau _3]\tau _1(c(\tau _2,\tau _3))+\frac{1}{2}\tau _1(\tau _2,\tau _3)$$
Interchanging $`\tau _1`$ with $`\tau _2`$ we get
$$s(\tau _2)_{(0)}s(\tau _1)_{(0)}s(\tau _3)=s([\tau _2,[\tau _1,\tau _3])$$
$$c(\tau _2,[\tau _1,\tau _3])+\frac{1}{2}\tau _2,[\tau _1,\tau _3]\tau _2(c(\tau _1,\tau _3))+\frac{1}{2}\tau _2(\tau _1,\tau _3)$$
Similarly,
$$(s(\tau _1)_{(0)}s(\tau _2))_{(0)}s(\tau _3)=\left\{s([\tau _1,\tau _2])c(\tau _1,\tau _2)+\frac{1}{2}\tau _1,\tau _2\right\}_{(0)}s(\tau _3)=$$
$$=s([\tau _1,\tau _2],\tau _3])c([\tau _1,\tau _2],\tau _3)+\frac{1}{2}[\tau _1,\tau _2],\tau _3+\tau _3(c(\tau _1,\tau _2))\tau _3,c(\tau _1,\tau _2)$$
where we have used the formula
$$\omega _{(0)}x=\pi (x)(\omega )+\pi (x),\omega $$
$`(\mathrm{2.3.5})`$
which is a consequence of (0.5.9). The axiom (A5) follows. This completes the proof of the theorem. $``$
2.4. Thus, to a pair $`(V,s)`$, where $`V`$ is a vertex algebra and $`s:TV_1`$ is a splitting on $`V`$, we have assigned a vertex algebroid $`𝒜(V,s)=(A,T,\mathrm{\Omega },,\gamma ,,,c)`$. Note that by definition $`𝒜(V,s)`$ has the property
(Sur) The module $`\mathrm{\Omega }`$ is generated as an $`A`$-module by the subspace $`A`$. Equivalently, the morphism $`\mathrm{\Omega }_{A/k}^1\mathrm{\Omega }`$ induced by $``$ is epimorphic.
§3. Category of Vertex Algebroids
3.1. Let us define a $`1`$-truncated vertex algebra to be a septuple $`v=(V_0,V_1,\text{1},,_{(1)},_{(0)},_{(1)})`$ where $`V_0,V_1`$ are two $`k`$-modules, 1 an element of $`V_0`$ (vacuum vector), $`:V_0V_1`$ a morphism of $`k`$-modules,
$`{}_{(i)}{}^{}:(V_0V_1)\times (V_0V_1)V_0V_1(i=1,0,1)`$
are $`k`$-bilinear operations of degree $`i1`$. Elements of $`V_0`$ (resp., $`V_1`$) will be denoted $`a,b,c`$ (resp., $`x,y,z`$). So, we have $`7`$ operations: $`a_{(1)}b,a_{(1)}x,x_{(1)}a,a_{(0)}x`$, $`x_{(0)}a,x_{(0)}y`$ and $`x_{(1)}y`$. The following axioms must be satisfied:
(Vacuum)
$$a_{(1)}\text{1}=a;x_{(1)}\text{1}=x;x_{(0)}\text{1}=0$$
$`(Vac)`$
(Derivation)
$$(a)_{(0)}b=0;(a)_{(0)}x=0;(a)_{(1)}x=a_{(0)}x$$
$`(Der)_1`$
$$(a_{(1)}b)=(a)_{(1)}b+a_{(1)}b;(x_{(0)}a)=x_{(0)}a$$
$`(Der)_2`$
(Commutativity)
$$a_{(1)}b=b_{(1)}a;a_{(1)}x=x_{(1)}a(x_{(0)}a)$$
$`(Com)_1`$
$$x_{(0)}a=a_{(0)}x;x_{(0)}y=y_{(0)}x+(y_{(1)}x)$$
$`(Com)_0`$
$$x_{(1)}y=y_{(1)}x$$
$`(Com)_1`$
(Associativity)
$$(a_{(1)}b)_{(1)}c=a_{(1)}b_{(1)}c$$
$`(Ass)_1`$
Operation <sub>(0)</sub> is a derivation with respect to all operations <sub>(i)</sub>, i.e.
$$\alpha _{(0)}\beta _{(i)}\gamma =(\alpha _{(0)}\beta )_{(i)}\gamma +\beta _{(i)}\alpha _{(0)}\gamma ,(\alpha ,\beta ,\gamma V_0V_1)$$
$`(Ass)_0`$
whenever the both sides are defined.
$$(a_{(1)}x)_{(0)}b=a_{(1)}x_{(0)}b$$
$`(Ass)_1`$
$$(a_{(1)}b)_{(1)}x=a_{(1)}b_{(1)}x+(a)_{(1)}b_{(0)}x+(b)_{(1)}a_{(0)}x$$
$`(Ass)_2`$
$$(a_{(1)}x)_{(1)}y=a_{(1)}x_{(1)}yx_{(0)}y_{(0)}a$$
$`(Ass)_3`$
A morphism between two $`1`$-truncated vertex algebras $`f:v=(V_0,V_1,\mathrm{})v^{}=(V_0^{},V_1^{},\mathrm{})`$ is a pair of maps of $`k`$-modules $`f=(f_0,f_1),f_i:V_iV_i^{}`$ such that $`f_0(\text{1})=\text{1}^{},f_1(a)=f_0(a)`$ and $`f(\alpha _{(i)}\beta )=f(\alpha )_{(i)}f(\beta )`$, whenever both sides are defined.
This way we get a category $`𝒱ert_1`$ of $`1`$-truncated vertex algebras. We have an obvious truncation functor
$$t:𝒱ert𝒱ert_1$$
$`(\mathrm{3.1.1})`$
which assignes to a vertex algebra $`V=V_i`$ the truncated algebra $`tV:=(V_0,V_1,\mathrm{})`$. In Section 9 below we shall construct a left adjoint to $`t`$.
3.2. Let $`v=(V_0,V_1,\mathrm{})`$ be a $`1`$-truncated vertex algebra. Let $`\mathrm{\Omega }_vV_1`$ be the $`k`$-submodule generated by all elements $`a_{(1)}b`$. Set $`T_v=V_1/\mathrm{\Omega }_v`$; let $`\pi :V_1T_v`$ be the canonical projection.
Let us call $`v`$ splittable if $`\pi `$ admits a $`k`$-linear splitting $`s:TV_1`$, cf. 2.2. This is of course a weak condition; it holds true for example if $`T`$ is a projective $`k`$-module. Let $`𝒱ert_1^{}𝒱ert_1`$ denote the full subcategory of splittable $`1`$-truncated vertex algebras.
A $`1`$-truncated vertex algebra with a chosen splitting $`s`$ will be called split. The argument of Section 2 assigns a vertex algebroid $`𝒜(v;s)`$ to every split $`1`$-truncated vertex algebra $`(v;s)`$.
We have by definition
$$𝒜(V;s)=𝒜(tV;s)$$
$`(\mathrm{3.2.2})`$
for every split vertex algebra $`(V;s)`$.
3.3. Conversely, let $`𝒜=(A,T,\mathrm{\Omega },\mathrm{})`$ be a vertex algebroid. We want to assign to it a split $`1`$-truncated vertex algebra. To this end one needs simply to invert the construction of the previous Section.
Namely $`V_0=A,V_1=T\mathrm{\Omega }`$; let $`:V_0V_1`$ be the composition of $`:A\mathrm{\Omega }`$ with the obvious embedding $`\mathrm{\Omega }V_1`$; let $`s:TV_1`$ be the obvious embedding.
Let us define the operations <sub>(i)</sub> as follows:
$$a_{(1)}b=ab;a_{(1)}\omega =a\omega ;a_{(1)}\tau =a\tau \gamma (a,\tau )$$
$`(\mathrm{3.3.1})`$
$$a_{(0)}b=a_{(0)}\omega =\omega _{(0)}\omega ^{}=0;\tau _{(0)}a=\tau (a);\tau _{(0)}\omega =\tau (\omega )$$
$`(\mathrm{3.3.2})`$
$$\tau _{(0)}\tau ^{}=[\tau ,\tau ^{}]c(\tau ,\tau ^{})+\frac{1}{2}\tau ,\tau ^{}$$
$`(\mathrm{3.3.3})`$
$$x_{(1)}y=x,y$$
$`(\mathrm{3.3.4})`$
By inverting the argument of the previous Section, one sees easily that we get a split $`1`$-truncated vertex algebra $`(V_0,V_1,\mathrm{})`$, to be denoted by $`u𝒜`$.
Let $`v𝒱ert_1^{}`$. For each splitting $`s`$ of $`v`$ we have by construction a canonical isomorphism
$$v\stackrel{}{}u𝒜(v;s)$$
$`(\mathrm{3.3.5})`$
For two vertex algebroids $`𝒜,𝒜^{}`$ define the set of morphisms between $`𝒜`$ and $`𝒜^{}`$ to be the subset of $`Hom_{𝒱ert_1}(u𝒜,u𝒜^{})`$ which consists of all morphisms $`g:u𝒜u𝒜^{}`$ such that $`g(\mathrm{\Omega })\mathrm{\Omega }^{}`$. This way we get a category of vertex algebroids, to be denoted $`𝒜lg`$. The mapping $`u`$ induces a functor
$$u:𝒜lg𝒱ert_1^{}$$
$`(\mathrm{3.3.6})`$
which is in fact an equivalence of categories, due to (3.3.5).
3.4. Let $`𝒜=(A,T,\mathrm{\Omega },\mathrm{})`$ and $`𝒜^{}=(A^{},T^{},\mathrm{\Omega }^{},\mathrm{})`$ be two vertex algebroids. We want to describe explicitly the set of morphisms $`Hom_{𝒜lg}(𝒜,𝒜^{})Hom_{𝒱ert_1}(u𝒜,u𝒜^{})`$.
Let $`g:u𝒜u𝒜^{}`$ be a morphism belonging to $`Hom_{𝒜lg}(𝒜,𝒜^{})`$. Let $`g_A:AA^{}`$ be its component of weight $`0`$. Since $`g`$ preserves the operation <sub>(-1)</sub> and the vacuum, $`g_A`$ is a morphism of commutative $`k`$-algebras. Let $`g_1:T\mathrm{\Omega }T^{}\mathrm{\Omega }^{}`$ be its component of weight $`1`$. We have $`g=g`$; therefore $`g_1(A)A^{}`$. By definition, $`g_1(\mathrm{\Omega })\mathrm{\Omega }^{}`$. Let us denote by $`g_\mathrm{\Omega }:\mathrm{\Omega }\mathrm{\Omega }^{}`$ and $`g_T:T=U𝒜/\mathrm{\Omega }U𝒜^{}/\mathrm{\Omega }^{}`$ the morphisms induced by $`g_1`$. Thus, in components $`g_1`$ has the form
$$g_1(\tau ,\omega )=(g_T(\tau ),g_\mathrm{\Omega }(\omega )+h(\tau ))$$
$`(\mathrm{3.4.1})`$
where $`h:T\mathrm{\Omega }^{}`$ is some $`k`$-linear mapping.
3.5. Theorem. The correspondence described in 3.4 provides a canonical isomorphism of the set $`Hom_{𝒜lg}(𝒜,𝒜^{})`$ with the set of all quadruples $`(g_A,g_T,g_\mathrm{\Omega },h)`$ where
$`(3.5A)`$ $`g_A:AA^{}`$ is a morphism of $`k`$-algebras;
$`(3.5\mathrm{\Omega })`$ $`g_\mathrm{\Omega }:\mathrm{\Omega }\mathrm{\Omega }^{}`$ is a morphism of $`k`$-modules such that $`g_\mathrm{\Omega }(a)=g_A(a)`$ and $`g_\mathrm{\Omega }(a\omega )=g_A(a)g_\mathrm{\Omega }(\omega )`$;
$`(3.5T)`$ $`g_T:TT^{}`$ is a morphism of Lie $`k`$-algebras such that $`g_T(a\tau )=g_A(a)g_T(\tau )`$$`g_A(\tau (a))=g_T(\tau )(g_A(a))`$$`g_\mathrm{\Omega }(\tau (\omega ))=g_T(\tau )(g_\mathrm{\Omega }(\omega ))`$ and $`g_A(\tau ,\omega )=g_T(\tau ),g_\mathrm{\Omega }(\omega )`$;
$`(3.5\gamma )`$ $`h:T\mathrm{\Omega }^{}`$ is a morphism of $`k`$-modules such that
$$h(a\tau )=g_A(a)h(\tau )\gamma ^{}(g_A(a),g_T(\tau ))+g_\mathrm{\Omega }(\gamma (a,\tau ));$$
$`(3.5,)`$ $`g_A(\tau _1,\tau _2)=g_T(\tau _1),g_T(\tau _2)^{}+g_T(\tau _1),h(\tau _2)+g_T(\tau _2),h(\tau _1)`$;
$`(3.5c)`$ $`g_\mathrm{\Omega }(c(\tau _1,\tau _2))=c^{}(g_T(\tau _1),g_T(\tau _2))+\frac{1}{2}g_T(\tau _1),h(\tau _2)\frac{1}{2}g_T(\tau _2),h(\tau _1)`$
$`g_T(\tau _1)(h(\tau _2))+g_T(\tau _2)(h(\tau _1))+h([\tau _1,\tau _2])`$.
3.6. Proof. By definition, a quadruple $`g=(g_A,g_T,g_\mathrm{\Omega },h)`$ defines a morphism $`u𝒜u𝒜^{}`$ iff it satisfies the identities (3.6.1) through (3.6.6) below $`(a,bA,x,yu𝒜_1)`$:
$$g_0(\text{1})=\text{1};g_0(a_{(1)}b)=g_0(a)_{(1)}g_0(b)$$
$`(\mathrm{3.6.1})`$
$$g_1(a)=g_0(a)$$
$`(\mathrm{3.6.2})`$
$$g_1(a_{(1)}x)=g_0(a)_{(1)}g_0(x)$$
$`(\mathrm{3.6.3})`$
$$g_1(x_{(0)}a)=g_1(x)_{(0)}g_0(a)$$
$`(\mathrm{3.6.4})`$
$$g_1(x_{(0)}y)=g_1(x)_{(0)}g_1(y)$$
$`(\mathrm{3.6.5})`$
$$g_0(x_{(1)}y)=g_1(x)_{(1)}g_1(y)$$
$`(\mathrm{3.6.6})`$
The condition (3.6.1) is equivalent to (3.5A); (3.6.2) and (3.6.3) for $`x=\omega `$ is equivalent to $`(3.5\mathrm{\Omega })`$. Let us write down (3.6.3) for $`x=\tau `$. We have
$$a_{(1)}\tau =(a\tau ,\gamma (a,\tau ))$$
$`(\mathrm{3.6.7})`$
Therefore $`g_1(a_{(1)}\tau )=(g(a\tau ),h(a\tau )g(\gamma (a,\tau ))`$ and
$$g_1(a)_{(1)}g_1(\tau )=g_1(a)_{(1)}(g(\tau ),h(\tau ))=(g(a)g(\tau ),\gamma (g(a),g(\tau ))+g(a)h(\tau ))$$
It follows that (3.6.3) for $`x=\tau `$ is equivalent to $`g(a\tau )=g(a)g(\tau )`$ and $`(3.5\gamma )`$. (3.6.4) for $`x=\omega `$ is vacuous. (3.6.4) for $`x=\tau `$ is quivalent to $`g(a\tau )=g(a)g(\tau )`$.
(3.6.6) for $`x,y\mathrm{\Omega }`$ is vacuous. (3.6.6) for $`x=\omega ,y=\tau `$ (or vice versa) is equivalent to $`g(\tau ,\omega )=g(\tau ),g(\omega )`$. We have
$$g_1(\tau _1)_{(1)}g_1(\tau _2)=g(\tau )_1,g(\tau _2)^{}+g(\tau _1),h(\tau _2)+g(\tau _2),h(\tau _1)$$
Therefore (3.6.6) for $`x,yT`$ is equivalent to $`(3.5,)`$.
(3.6.5) for $`x=\tau ,y=\omega `$ (or vice versa) is equivalent to $`g(\tau (\omega ))=g(\tau )(g(\omega ))`$. Finally, let us write down (3.6.5) for $`x,yT`$. We have
$$\tau _{1(0)}\tau _2=([\tau _1,\tau _2],c(\tau _1,\tau _2)+\frac{1}{2}\tau _1,\tau _2)$$
$`(\mathrm{3.6.8})`$
cf. (2.2.2) and $`(\mathrm{2.2.3})_c`$. Therefore
$$g_1(\tau _{1(0)}\tau _2)=(g([\tau _1,\tau _2],h([\tau _1,\tau _2])g(c(\tau _1,\tau _2)+\frac{1}{2}g(\tau _1),g(\tau _2)+$$
$$+\frac{1}{2}g(\tau _1),h(\tau _2)+\frac{1}{2}g(\tau _2),h(\tau _1))$$
where we have used $`(3.5,)`$. On the other hand,
$$g_1(\tau _1)_{(0)}g_1(\tau _2)=([g(\tau _1),g(\tau _2)],c^{}(g(\tau _1),g(\tau _2))+\frac{1}{2}g(\tau _1),g(\tau _2)+$$
$$+g(\tau _1)(h(\tau _2))g(\tau _2)(h(\tau _1))+h(\tau _1),g(\tau _2))$$
where we have used the formula
$$\omega _{(0)}\tau =\tau (\omega )+\tau ,\omega $$
$`(\mathrm{3.6.9})`$
following from (0.5.9) (cf. (2.3.5)). Therefore (3.6.5) for $`x,yT`$ is equivalent to the requirement that $`g_T`$ is a morphism of Lie algebras, and to $`(3.5c)`$. The theorem is proved. $``$
3.7. Let $`(A,T,\mathrm{})\stackrel{g}{}(A^{},T^{},\mathrm{})\stackrel{g^{}}{}(A^{\prime \prime },T^{\prime \prime },\mathrm{})`$ be two morphisms of vertex algebroids. Then their composition is obviously equal to
$$g^{}g=(g_A^{}g_A,g_T^{}g_T,g_\mathrm{\Omega }^{}g_\mathrm{\Omega },g_\mathrm{\Omega }^{}h+h^{}g_T)$$
$`(\mathrm{3.7.1})`$
cf. (3.4.1). The identity morphisms are
$$Id_𝒜=(Id_A,Id_T,Id_\mathrm{\Omega },0)$$
$`(\mathrm{3.7.2})`$
3.8. Theorem. (Extension of Mappings.) In the conditions of (3.5), let $`T_0T`$ be a $`k`$-submodule which generates $`T`$ as an $`A`$-module. Let $`g=(g_A,g_T,g_\mathrm{\Omega },h)`$ be such a quadruple that $`(3.5A),(3.5\mathrm{\Omega })`$ and $`(3.5T)`$ are fulfilled and such that $`(3.5\gamma ),(3.5,)`$ and $`(3.5c)`$ are fulfilled for all $`\tau ,\tau _iT_0`$.
Then $`(3.5\gamma ),(3.5,)`$ and $`(3.5c)`$ hold true for all $`\tau ,\tau _iT`$, i.e. $`g`$ defines a morphism $`𝒜𝒜^{}`$.
Cf. Theorem 1.9.
Proof. Let us prove that if $`(3.5\gamma )`$ is true $`(a,\tau )`$ with a fixed $`\tau `$ and all $`a`$, then it is true for all couples $`(a,b\tau )`$. We have to prove that
$$h(ab\tau )=g(a)h(b\tau )\gamma (g(a),g(b\tau ))+g(\gamma (a,b\tau ))$$
$`(\mathrm{?})`$
The left hand side is equal to
$$h(ab\tau )=g(ab)h(\tau )\gamma (g(ab),g(\tau ))+g(\gamma (ab,\tau ))=$$
$$=g(ab)h(\tau )\gamma (g(a),g(b)g(\tau ))g(a)\gamma (g(b),g(\tau ))g(\tau )(g(a))g(b)g(\tau )(g(b))g(a)+$$
$$+g(\gamma (a,b\tau ))+g(a\gamma (b,\tau ))+g(\tau (a)b)+g(\tau (b))a)$$
where we have used (A1). On the other hand, the first summand in the rhs is equal to
$$g(a)h(b\tau )=g(a)g(b)h(\tau )g(a)\gamma (g(b),g(\tau ))+g(a)g(\gamma (b,\tau )),$$
and we easily see that we have the required identity indeed.
The similar claims connected with the equations $`(3.5,)`$ and $`(3.5c)`$ are proved analogously, and we leave them to the reader. $``$
§4. Cofibered Structure
4.1. Let us introduce objects which lie ”in between” extended Lie algebroids and vertex algebroid. Let us define a vertex prealgebroid to be a sextuple $`=(A,T,\mathrm{\Omega },,\gamma ,,)`$ where $`A,T,\mathrm{\Omega },,\gamma ,,`$ are as in 1.4.
The data $`(A,T,\mathrm{\Omega },,|_{T\times \mathrm{\Omega }})`$ should form an extended Lie algebroids and the axioms 1.4 (A1) and (A2) must be satisfied.
Let us define a morphism between two vertex prealgebroids $`=(A,T,\mathrm{})`$ and $`^{}=(A^{},T^{},\mathrm{})`$ to be a quadruple $`g=(g_A,g_T,g_\mathrm{\Omega },h)`$ as in Theorem 3.5, satisfying the properties $`(3.5A)`$$`(3.5,)`$. The composition of morphisms and the identity morphisms are defined by the rules (3.7.1), (3.7.2). This way we get a category $`𝒫reAlg`$ of vertex prealgebroids.
We have obvious forgetful functors
$$𝒜lg\stackrel{P}{}𝒫reAlg\stackrel{Q}{}ieAlg$$
$`(\mathrm{4.1.1})`$
Note that $`P`$ is injective on $`Hom`$’s. The composition $`QP:𝒜lgieAlg`$ will be denoted by $`\mathrm{\Theta }`$.
We will use the standard notation for fibers. For example, if $`𝒫reAlg`$ then $`𝒜lg_{}`$ will denote the category whose objects are $`𝒜𝒜lg`$ such that $`P(𝒜)=`$ and morphisms are morphisms in $`𝒜`$ which go to $`Id_{}`$ after applying $`P`$, and similarly with the other fibers.
4.2. Let $`=(A,T,\mathrm{\Omega },,\gamma ,,)𝒫reAlg`$; let $`𝒯=Q()=(A,T,\mathrm{\Omega },,|_{T\times \mathrm{\Omega }})`$. Let $`\stackrel{~}{𝒜lg}_{}`$ be a full subcategory of $`𝒜lg_𝒯`$ with $`Ob\stackrel{~}{𝒜lg}_{}=Ob𝒜lg_{}`$. In other words, $`Hom_{\stackrel{~}{𝒜lg}_{}}(𝒜,𝒜^{})`$ consists of quadruples of the form $`g=(Id_A,Id_T,Id_\mathrm{\Omega },h)`$.
Note that all such $`g`$ are invertible, namely $`g^1=(Id_A,Id_T,\text{Id}_\mathrm{\Omega },h)`$, cf. (3.7.1), (3.7.2). In other words, $`\stackrel{~}{𝒜lg}_{}`$ is a groupoid.
Consider the de Rham-Chevalley complex of $`𝒯`$:
$$\mathrm{\Omega }^{}(𝒯):0A\mathrm{\Omega }\stackrel{d_{DR}}{}\mathrm{\Omega }^2\mathrm{}$$
cf. 1.3. Let $`𝒜_i=(\mathrm{},c_i)Ob\stackrel{~}{𝒜lg}_{},i=1,2`$. Set $`\alpha =c_1c_2`$. It follows from (A4) that
$$\tau _2,\alpha (\tau _1,\tau _3)+\tau _3,\alpha (\tau _1,\tau _2)=0,$$
i.e. $`\alpha \mathrm{\Omega }^3`$. It follows from (A5) that $`\alpha \mathrm{\Omega }^{3,cl}:=Ker(\mathrm{\Omega }^3\mathrm{\Omega }^4)`$, cf. (1.3.2).
Conversely, if $`𝒜=(\mathrm{},c)Ob\stackrel{~}{𝒜lg}_{}`$ and $`\alpha \mathrm{\Omega }^{3,cl}`$ then $`𝒜\stackrel{}{+}\alpha :=(\mathrm{},c+\alpha )Ob𝒜lg_{}`$. We have proven that
4.2.1. The set $`Ob\stackrel{~}{𝒜lg}_{}`$ is canonically an $`\mathrm{\Omega }^{3,cl}`$-torseur.
Let $`𝒜Ob\stackrel{~}{𝒜lg}_{}`$, $`𝒜^{}=𝒜\stackrel{}{+}\alpha ,\alpha \mathrm{\Omega }^{3,cl}`$, $`g=(Id_A,Id_T,Id_\mathrm{\Omega },h)Hom_{\stackrel{~}{𝒜lg}_{}}(𝒜,𝒜^{})`$. Then, due to $`(3.5\gamma )`$, $`h`$ must be $`A`$-linear; by $`(3.5,)`$, $`h\mathrm{\Omega }^2`$ and by $`(3.5c)`$, $`d_{DR}h=\alpha `$. Therefore,
4.2.2. the set $`Hom_{\stackrel{~}{𝒜lg}_{}}(𝒜,𝒜\stackrel{}{+}\alpha )`$ may be canonically identified with the set of $`h\mathrm{\Omega }^2`$ such that $`d_{DR}h=\alpha `$.
Consequently, the set of isomorphism classes of objects of the groupoid $`\stackrel{~}{𝒜lg}_{}`$ is equal to the third de Rham cohomology $`H^3(\mathrm{\Omega }^{}(𝒯))`$.
4.3. Let $`𝒯=(A,T,\mathrm{}),𝒯^{}=(A^{},T^{},\mathrm{})ieAlg`$, $`g=(g_A,g_T,g_\mathrm{\Omega })Hom_{ieAlg}(𝒯^{},𝒯)`$. Let $`𝒫reAlg_𝒯,^{}𝒫reAlg_𝒯^{}`$.
Let us consider the set $`Hom_g(^{},)Hom_{𝒫reAlg}(^{},)`$ consisting of all morphisms $`\stackrel{~}{g}`$ such that $`Q(\stackrel{~}{g})=g`$. Let $`\stackrel{~}{g}=(g_A,g_T,g_\mathrm{\Omega },h),\stackrel{~}{g}^{}=(g_A,g_T,g_\mathrm{\Omega },h^{})Hom_g(^{},)`$. Set $`\beta :=hh^{}:T\mathrm{\Omega }^{}`$. Due to $`(3.5\gamma )`$ and $`(3.5,)`$, $`\beta `$ will satisfy the properties
$$\beta (a\tau )=g_A(a)\beta (\tau )$$
$`(\mathrm{4.3.1})`$
and
$$g_T(\tau _1),\beta (\tau _2)+g_T(\tau _2),\beta (\tau _1)=0$$
$`(\mathrm{4.3.2})`$
Let us denote by $`\mathrm{\Omega }_g^2`$ the set of all maps $`\beta :T^{}\mathrm{\Omega }`$ satisfying (4.3.1) and (4.3.2). It is an $`A`$-module in the obvious way, and consequently an $`A^{}`$-module, by restriction of scalars. In particular, $`\mathrm{\Omega }_{Id_𝒯}^2=\mathrm{\Omega }^2(𝒯)`$.
Vice versa, for each $`\beta \mathrm{\Omega }_g^2`$, the map $`\stackrel{~}{g}^{}\stackrel{}{+}\beta :=(g_A,g_T,g_\mathrm{\Omega },h+\beta )`$ belongs to $`Hom_g(^{},)`$. Thus we have proven that
4.3.1. the set $`Hom_g(^{},)`$ is canonically an $`\mathrm{\Omega }_g^2`$-torseur. In particular, if $`𝒯=𝒯^{}`$ then $`Hom_{𝒫reAlg_𝒯}(^{},)`$ is an $`\mathrm{\Omega }^2(𝒯)`$-torseur.
Similarly, let $`𝒜𝒜lg_𝒯,𝒜^{}𝒜lg_𝒯^{}`$. Let $`Hom_g(𝒜^{},𝒜)`$ be the subset of $`Hom_{𝒜lg}(𝒜^{},𝒜)`$ consisting of all morphisms $`\stackrel{~}{g}`$ such that $`\mathrm{\Theta }(\stackrel{~}{g})=g`$. Let $`\mathrm{\Omega }_g^{2,cl}`$ be the $`k`$-module of all maps $`\beta :T^{}\mathrm{\Omega }`$ satisfying (4.3.1), (4.3.2) and
$$\beta ([\tau _1,\tau _2])+g_T(\tau _1)(h(\tau _2))g_T(\tau _2)(h(\tau _1))g_T(\tau _1),h(\tau _2)=0$$
$`(\mathrm{4.3.3})`$
In particular, $`\mathrm{\Omega }_{Id_𝒯}^{2,cl}=\mathrm{\Omega }^{2,cl}(𝒯)`$, cf. (1.3.4). As above, due to (3.5c), we get
4.3.2. the set $`Hom_g(𝒜^{},𝒜)`$ is canonically an $`\mathrm{\Omega }_g^{2,cl}`$-torseur. In particular, if $`𝒯=𝒯^{}`$ then $`Hom_{𝒜lg_𝒯}(𝒜^{},𝒜)`$ is an $`\mathrm{\Omega }^{2,cl}(𝒯)`$-torseur.
If $`g^{}:𝒯^{\prime \prime }𝒯^{}`$ is another morphism of extended Lie algebroids, we have a composition map
$$\nu _{g,g^{}}:\mathrm{\Omega }_g^2\times \mathrm{\Omega }_g^{}^2\mathrm{\Omega }_{gg^{}}^2,\nu _{g,g^{}}(\beta ,\beta ^{}):=g_\mathrm{\Omega }\beta ^{}+\beta g_T^{}$$
$`(\mathrm{4.3.4})`$
It is a morphism of $`A^{}`$-modules. We have $`\nu _{g,g^{}}(\mathrm{\Omega }_g^{2,cl}\times \mathrm{\Omega }_g^{}^{2,cl})\mathrm{\Omega }_{gg^{}}^{2,cl}`$ (cf. Remark 4.4 below). The maps (4.3.4) are associative in the obvious sense (with respect to triples of composable morphisms in $`ieAlg`$).
On the other hand, given $`^{\prime \prime }𝒫reAlg_{𝒯^{\prime \prime }}`$, we have the composition
$$Hom_g(^{},)\times Hom_g^{}(^{\prime \prime },^{})Hom_{gg^{}}(^{\prime \prime },)$$
$`(\mathrm{4.3.5})`$
4.3.3. The map (4.3.5) is compatible with (4.3.4). Therefore, we have a canonical isomorphism of $`\mathrm{\Omega }_{gg^{}}^2`$-torseurs
$$\rho _{g,g^{}}:Hom_{gg^{}}(^{\prime \prime },)\stackrel{}{}\nu _{g,g^{}}(Hom_g(^{},)\times Hom_g^{}(^{\prime \prime },))$$
$`(\mathrm{4.3.6})`$
The isomorphisms $`\rho _{g,g^{}}`$ satisfy obvious $`2`$-cocycle equations connected with triples of composable morphisms $`(g,g^{},g^{\prime \prime })`$.
This follows from (3.7.1).
4.4. Remark. Generalizing 1.3, one can define the de Rham-Chevalley complexes $`\mathrm{\Omega }_g^{}`$ such that the modules $`\mathrm{\Omega }_g^2`$ and $`\mathrm{\Omega }_g^{2,cl}`$ from the previous no. become really the module of two-forms and closed ones respectively. These complexes come equipped with the composition maps $`\mathrm{\Omega }_g^{}\times \mathrm{\Omega }_g^{}^{}\mathrm{\Omega }_{gg^{}}^{}`$ which are compatible with the de Rham differentials and satisfy associativity.
We leave the necessary definitions as an exercise to the reader.
4.5. Theorem. Let $`g=(g_A,g_T,g_\mathrm{\Omega },h):^{}`$ be a morphism between vertex prealgebroids such that $`g_T`$ is an isomorphism. Let $`𝒜𝒜lg_{}`$. There exists a unique pair $`(g_{}𝒜,\stackrel{~}{g})`$ where $`g_{}𝒜𝒜lg_{^{}}`$ and $`\stackrel{~}{g}:𝒜g_{}𝒜`$ is a morphism of vertex algebroids such that $`P(\stackrel{~}{g})=\stackrel{~}{g}`$.
Proof. Uniqueness. Let $`𝒜=(A,T,\mathrm{\Omega },,\gamma ,,,c)`$;
let $`𝒜^{}=(A^{},T^{},\mathrm{\Omega }^{},^{},,^{},\gamma ^{},c^{})𝒜lg_{^{}}`$. A morphism $`\stackrel{~}{g}:𝒜𝒜^{}`$ such that $`P(\stackrel{~}{g})=g`$ must of course be represented by the same quadruple as $`g`$, i.e. $`\stackrel{~}{g}=(g_A,g_T,g_\mathrm{\Omega },h)`$. It is a morphism of vertex algebroids iff the condition $`(3.5c)`$ is fulfilled, i.e. iff
$$c^{}(g_T(\tau _1),g_T(\tau _2))=g_\mathrm{\Omega }(c(\tau _1,\tau _2))\frac{1}{2}^{}g_T(\tau _1),h(\tau _2)^{}+\frac{1}{2}^{}g_T(\tau _2),h(\tau _1)^{}+$$
$$+g_T(\tau _1)(h(\tau _2))g_T(\tau _2)(h(\tau _1))h([\tau _1,\tau _2])$$
$`(\mathrm{4.5.1})`$
This equation defines $`c^{}`$ uniquely since $`g_T`$ is an isomorphism by assumption.
Existence. We have to check that the map $`c^{}`$ defined by (4.5.1) satisfies axioms (A3), (A4), (A5). Let us check (A4) for example.
To unburden the notation, let us assume that $`Q()=Q(^{})`$ and $`g=(Id_A,Id_T,Id_\mathrm{\Omega },h)`$ (the general case being treated by an identical computation). We have to prove that
$$\tau _2,c^{}(\tau _1,\tau _3)+\tau _3,c^{}(\tau _1,\tau _2)=$$
$$=[\tau _1,\tau _2],\tau _3^{}+[\tau _1,\tau _3],\tau _2^{}\tau _1(\tau _2,\tau _3^{})+\frac{1}{2}\tau _2(\tau _1,\tau _3)^{})+\frac{1}{2}\tau _3(\tau _1,\tau _2^{})$$
$`(\mathrm{?})`$
We have
$$\tau _2,c^{}(\tau _1,\tau _3)=\tau _2,c(\tau _1,\tau _3)\frac{1}{2}\tau _1,h(\tau _3)\tau _3)+\frac{1}{2}\tau _3,h(\tau _1)+\tau _1(h(\tau _3))$$
$$\tau _3(h(\tau _1))h([\tau _1,\tau _3])=\tau _2,c(\tau _1,\tau _3)\frac{1}{2}\tau _2(\tau _1,h(\tau _3))+\frac{1}{2}\tau _2(\tau _3,h(\tau _1))+$$
$$+\tau _1(\tau _2,h(\tau _3))[\tau _1,\tau _2],h(\tau _3)\tau _3(\tau _2,h(\tau _1))[\tau _3,\tau _2],h(\tau _1)\tau _2,h([\tau _1,\tau _3])$$
Similarly, interchanging $`\tau _2`$ with $`\tau _3`$,
$$\tau _3,c^{}(\tau _1,\tau _2)=\tau _3,c(\tau _1,\tau _2)\frac{1}{2}\tau _3(\tau _1,h(\tau _2))+\frac{1}{2}\tau _3(\tau _2,h(\tau _1))+$$
$$+\tau _1(\tau _3,h(\tau _2))[\tau _1,\tau _3],h(\tau _2)\tau _2(\tau _3,h(\tau _1))[\tau _2,\tau _3],h(\tau _1)\tau _3,h([\tau _1,\tau _2])$$
By assumption,
$$\tau _2,c(\tau _1,\tau _3)+\tau _3,c(\tau _1,\tau _2)=$$
$$=[\tau _1,\tau _2],\tau _3+[\tau _1,\tau _3],\tau _2\tau _1(\tau _2,\tau _3)+\frac{1}{2}\tau _2(\tau _1,\tau _3))+\frac{1}{2}\tau _3(\tau _1,\tau _2)$$
Using this and the axiom $`(3.5,)`$ which takes the form
$$\tau _1,\tau _2^{}=\tau _1,\tau _2\tau _1,h(\tau _2)\tau _2,h(\tau _1)$$
we get the required identity (?).
The other two axioms, (A3) and (A5), are checked in a similar manner, and we leave them to the reader. It is convenient to check (A5) in its equivalent form $`(A5)^{ter}`$, see 1.6. $``$
§5. Chern-Simons
5.1. Let us define a frame of a Lie $`A`$-algebroid $`T`$ to be a $`k`$-submodule $`𝔤T`$ such that $`A_k𝔤=T`$. For example, if $`T`$ is a free $`A`$-module and $`\{\tau _i\}`$ is some $`A`$-base of $`T`$ then $`𝔤_{\{\tau _i\}}:=k\tau _i`$ is a frame in $`T`$.
A frame of a vertex $`A`$-algebroid (resp. prealgebroid, extended Lie algebroid) is by definition a frame of the underlying Lie algebroid $`T`$. A vertex algebroid (resp. prealgebroid, …) equipped with a frame will be called framed.
We will call an extended Lie algebroid $`𝒯`$ quasiregular if it is perfect (see 1.2) and admits a frame.
5.2. Let us consider the situation of 4.3, and assume that $`𝒯^{}`$ is equipped with a frame $`𝔤^{}`$ and $`𝒯`$ is perfect. Let us assume that $`g_T`$ is an isomorphism. Then $`𝔤:=g_T(𝔤^{})`$ is a frame in $`T`$. Let us define a $`k`$-linear map $`h_𝔤^{}:𝔤^{}\mathrm{\Omega }`$ by the condition
$$g_T(\tau _1),h_𝔤^{}(\tau _2)=\frac{1}{2}\{g_A(\tau _1,\tau _2^{})g_T(\tau _1),g_T(\tau _2)\}(\tau _i𝔤^{})$$
$`(\mathrm{5.2.1})`$
Since $`g_T`$ is an isomorphism, (5.2.1) defines $`h_𝔤^{}`$ uniquely. Then the condition $`(3.5,)`$ is obviously satisfied for all $`\tau _i𝔤^{}`$. There is a unique extension of $`h_𝔤^{}`$ to a map $`h_𝔤^{}:T^{}\mathrm{\Omega }`$ satisfying $`(3.5\gamma )`$ for all $`\tau 𝔤^{}`$.
By Theorem 3.8, the conditions $`(3.5\gamma )`$ and $`(3.5,)`$ are then fulfilled for all $`\tau ,\tau _iT`$, that is, $`g_𝔤^{}:=(g,h_𝔤^{})`$ is a morphism of vertex prealgebroids $`^{}`$.
In other words,
5.2.1. the $`\mathrm{\Omega }^2(𝒯)`$-torseur $`Hom_g(^{},)`$ is equipped with a trivialization $`g_𝔤^{}`$.
5.3. Let $`(𝒯,𝔤)`$ be a framed extended Lie algebroid. We can define a vertex algebroid $`𝒜_{𝒯,𝔤}`$ by setting $`\gamma (a,\tau )=0,\tau ,\tau ^{}=c(\tau ,\tau ^{})=0`$ for all $`aA,\tau ,\tau ^{}𝔤`$ and then extending the operations $`\gamma ,,,c`$ to the whole of $`T`$ using the axioms (A1) — (A3). By Theorem 1.9 we get a vertex algebroid, to be denoted $`𝒜_{𝒯;𝔤}`$.
If $`𝔤`$ is a Lie subalgebra of $`T`$ then $`𝒜_{𝒯;𝔤}`$ may be defined as follows. Take the vertex $`k`$-algebroid $`𝒜_𝔤=(k,𝔤,𝔤^{},0,0,,,0)`$ where $`𝔤`$ acts on $`𝔤^{}`$ in the coadjoint way, $`,|_{𝔤\times 𝔤^{}}`$ is the obvious pairing, the other components of $`,`$ being zero. Then apply to $`𝒜_𝔤`$ the pushout with respect to the structure morphism $`kA`$, where we set $`\gamma (a,\tau )=0`$ for $`\tau 𝔤`$ and the Lie algebra acts on $`A`$ due to the embedding $`𝔤T`$, cf. 1.10. We get a vertex $`A`$-algebroid which is equal to $`𝒜_{𝒯;𝔤}`$.
We set $`_{𝒯;𝔤}:=P(𝒜_{𝒯;𝔤})`$.
Let us call the frame $`𝔤`$ abelian if $`𝔤`$ is an abelian Lie subalgebra of $`T`$, i.e. $`[\tau ,\tau ^{}]=0`$ for all $`\tau ,\tau ^{}𝔤`$. We call an extended Lie algebroid $`𝒯`$ regular if it is perfect and admits an abelian frame.
5.4. Let $`𝒯=(A,T,\mathrm{\Omega },\mathrm{})`$ be a regular extended Lie algebroid. Let $`𝔤,𝔤^{}`$ be two abelian frames of $`𝒯`$. Let us assume that $`𝔤`$ and $`𝔤^{}`$ are free $`k`$-modules of finite rank $`n`$. In the sequel we will need some formulas pertaining to this situation.
Choose some $`k`$-bases $`\{\tau _i\},\{\tau _i^{}\},i=1,\mathrm{},n,`$ of $`𝔤`$ and $`𝔤^{}`$ respectively; let $`\{\omega _i\}𝔤^{},\{\omega _i^{}\}\mathrm{\Omega }^{}`$ be the dual bases. Note that
$$\tau _i(\omega _j)=0\text{for all }i,j,$$
$`(\mathrm{5.4.1})`$
since for all $`p`$, $`\tau _p,\tau _i(\omega _j)=\tau _i(\tau _p,\omega _j)[\tau _i,\tau _p],\omega _j=0`$ because $`𝔤`$ is abelian and $`\tau _p,\omega _j=\delta _{pj}`$.
Recall that
$$\tau ,ab=a\tau (b)$$
$`(\mathrm{5.4.2})`$
and
$$(a\tau )(\omega )=a\tau (\omega )+\tau ,\omega a$$
$`(\mathrm{5.4.3})`$
Since $`\tau _i,a=\tau _i(a)`$,
$$a=\tau _i(a)\omega _i$$
$`(\mathrm{5.4.4})`$
where we always imply the summation over repeating indices.
Define the matrices $`\varphi =(\varphi ^{ij}),\rho =(\rho ^{ij})GL_n(A)`$ by $`\tau _i^{}=\varphi ^{ij}\tau _j,\omega _i^{}=\rho ^{ij}\omega _j`$. Since $`\tau _i,\omega _j=\tau _i^{},\omega _j^{}=\delta _{ij}`$, we have
$$\varphi \rho ^t=1$$
$`(\mathrm{5.4.5})`$
where $`(\rho ^t)^{ij}=\rho ^{ij}`$. Since $`[\tau _i,\tau _j]=[\tau _i^{},\tau _j^{}]=0`$,
$$\varphi ^{ip}\tau _p(\varphi ^{jq})=\varphi ^{jp}\tau _p(\varphi ^{iq})$$
$`(\mathrm{5.4.6})`$
for all $`i,j,q`$. Applying $`\tau _r`$ to (5.4.6) we get
$$\varphi ^{ip}\tau _r\tau _p(\varphi ^{jq})\varphi ^{jp}\tau _r\tau _p(\varphi ^{iq})=\tau _r(\varphi ^{jp})\tau _p(\varphi ^{iq})\tau _r(\varphi ^{ip})\tau _p(\varphi ^{jq})$$
$`(\mathrm{5.4.7})`$
for all $`i,j,q,r`$. Setting here $`r=q`$ and summing up by $`q`$ we get
$$\varphi ^{ip}\tau _q\tau _p(\varphi ^{jq})=\varphi ^{jq}\tau _p\tau _q(\varphi ^{ip})$$
$`(\mathrm{5.4.8})`$
for all $`i,j`$. Applying $`\tau _r`$ we get
$$\varphi ^{ip}\tau _r\tau _q\tau _p(\varphi ^{jq})\varphi ^{jq}\tau _r\tau _p\tau _q(\varphi ^{ip})=\tau _r(\varphi ^{jq})\tau _p\tau _q(\varphi ^{ip})\tau _r(\varphi ^{ip})\tau _q\tau _p(\varphi ^{jq})$$
$`(\mathrm{5.4.9})`$
It follows from (5.4.5) that
$$\tau _r(\rho )=\rho \tau _r(\varphi ^t)\rho ;\tau _r(\varphi )=\varphi \tau _r(\rho ^t)\varphi $$
$`(\mathrm{5.4.10})`$
for all $`r`$. Multiplying (5.4.6) by $`\rho ^{ia}`$ and summing up by $`i`$, we get
$$\tau _a(\varphi ^{jq})=\varphi ^{jp}\rho ^{ia}\tau _p(\varphi ^{iq})$$
$`(\mathrm{5.4.11})`$
whence
$$\rho ^{ib}\tau _a(\varphi ^{iq})=\rho ^{ia}\tau _b(\varphi ^{iq})$$
$`(\mathrm{5.4.12})`$
By (5.4.10)
$$\tau _c(\rho ^{ub})=\rho ^{uk}\tau _c(\varphi ^{lk})\rho ^{lb}\stackrel{\left(\mathrm{5.4.11}\right)}{=}$$
$$=\rho ^{uk}\varphi ^{lp}\rho ^{ic}\tau _p(\varphi ^{ik})\rho ^{lb}=\rho ^{uk}\tau _b(\varphi ^{ik})\rho ^{ic}=\tau _b(\rho ^{uc})$$
Thus,
$$\tau _a(\rho ^{bc})=\tau _c(\rho ^{ba})$$
$`(\mathrm{5.4.13})`$
5.5. In the situation 5.4, consider the vertex prealgebroids $`_{𝒯;𝔤}=(A,T,\mathrm{\Omega },,\gamma ,,)`$ and $`_{𝒯;𝔤^{}}=(A,T,\mathrm{\Omega },,\gamma ^{},,^{})`$. According to 5.2, we have an isomorphism
$$g_{𝔤,𝔤^{}}=(Id_A,Id_T,Id_\mathrm{\Omega },h_{𝔤,𝔤^{}}):_{𝒯;𝔤^{}}\stackrel{}{}_{𝒯;𝔤}$$
$`(\mathrm{5.5.1})`$
where $`h=h_{𝔤;𝔤^{}}`$ is defined by
$$\tau _i^{},h(\tau _j^{})=\frac{1}{2}\tau _i^{},\tau _j^{}$$
$`(\mathrm{5.5.2})`$
which is (5.2.1) in our situation.
Now consider the vertex algebroids $`𝒜_{𝒯;𝔤}=(\mathrm{},c),𝒜_{𝒯;𝔤^{}}=(\mathrm{},c^{})`$. We have $`g_{}𝒜_{𝒯;𝔤^{}},𝒜_{𝒯;𝔤}𝒜lg_{_{𝒯;𝔤}}`$. Therefore, by 4.2.1 $`𝒜_{𝒯;𝔤}=g_{}𝒜_{𝒯;𝔤^{}}\stackrel{}{+}\beta `$ for some closed $`3`$-form $`\beta `$.
Let us define a form $`\beta =\beta _{𝔤,𝔤^{}}\mathrm{\Omega }^3(𝒯)`$ by
$$\beta (\tau _i^{},\tau _j^{})=\beta ^{ijr}\omega _r^{}=\frac{1}{2}\left\{\tau _u(\varphi ^{jp})\tau _p(\varphi ^{iq})\tau _q(\varphi ^{ru})\tau _u(\varphi ^{ip})\tau _p(\varphi ^{jq})\tau _q(\varphi ^{ru})\right\}\omega _r^{}$$
$`(\mathrm{5.5.3})`$
5.6. Magic Lemma. The form $`\beta _{𝔤,𝔤^{}}`$ is closed and we have $`𝒜_{𝒯;𝔤}=g_{}𝒜_{𝒯;𝔤^{}}\stackrel{}{+}\beta _{𝔤,𝔤^{}}`$.
5.7. Proof. Let us write down the things explicitly. The following formulas hold true in the algebroid $`𝒜_{𝒯;𝔤}`$ (they are (1.8.3) in our situation):
$$\gamma (a,b\tau _i)=\tau _i(a)b\tau _i(b)a$$
$`(\mathrm{5.7.1})_\gamma `$
$$a\tau _i,b\tau _j=b\tau _i\tau _j(a)a\tau _j\tau _i(b)\tau _i(b)\tau _j(a)$$
$`(\mathrm{5.7.1})_,`$
$$c(a\tau _i,b\tau _j)=\frac{1}{2}\{\tau _i(b)\tau _j(a)\tau _j(a)\tau _i(b)\}+\frac{1}{2}\{b\tau _i\tau _j(a)a\tau _j\tau _i(b)\}$$
$`(\mathrm{5.7.1})_c`$
Define the matrix $`(h^{ij})Mat_n(A)`$ by $`h(\tau _i^{})=h^{ij}\omega _j`$. The left hand side of (5.5.2) is equal to
$$\varphi ^{ip}\tau _p,h^{jq}\omega _q=\varphi ^{ip}h^{jp}$$
By $`(\mathrm{5.7.1})_,`$, the right hand side is equal to
$$\frac{1}{2}\{\varphi ^{jq}\tau _p\tau _q(\varphi ^{ip})+\varphi ^{ip}\tau _q\tau _p(\varphi ^{jq})+\tau _p(\varphi ^{jq})\tau _q(\varphi ^{ip})\}\stackrel{\left(\mathrm{5.4.8}\right)}{=}$$
$$=\varphi ^{ip}\tau _q\tau _p(\varphi ^{jq})+\frac{1}{2}\tau _p(\varphi ^{jq})\tau _q(\varphi ^{ip})$$
Thus, the equation (5.5.2) takes the form
$$\varphi ^{ip}h^{jp}=\varphi ^{ip}\tau _q\tau _p(\varphi ^{jq})+\frac{1}{2}\tau _p(\varphi ^{jq})\tau _q(\varphi ^{ip})$$
wherefrom, applying (5.4.5),
$$h^{ij}=\tau _p\tau _j(\varphi ^{ip})+\frac{1}{2}\tau _q(\varphi ^{ip})\tau _p(\varphi ^{rq})\rho ^{rj}$$
$`(\mathrm{5.7.2})`$
We have to prove that
$$c(\tau _i^{},\tau _j^{})=g_{}c^{}(\tau _i^{},\tau _j^{})+\beta (\tau _i^{},\tau _j^{})$$
where $`g_{}c^{}`$ is defined by
$$g_{}c^{}(\tau _i^{},\tau _j^{})=c^{}(\tau _i^{},\tau _j^{})+\tau _i^{}(h(\tau _j^{}))\tau _j^{}(h(\tau _i^{})),$$
by (4.5.1). Thus, we have to prove that
$$c(\tau _i^{},\tau _j^{})c^{}(\tau _i^{},\tau _j^{})\tau _i^{}(h(\tau _j^{}))+\tau _j^{}(h(\tau _i^{}))=\beta (\tau _i^{},\tau _j^{})$$
$`(\mathrm{5.7.3})`$
By definition, $`c^{}(\tau _i^{},\tau _j^{})=0`$. By $`(\mathrm{5.7.1})_c`$ and (5.4.8) we have
$$c(\tau _i^{},\tau _j^{})=c(\varphi ^{ip}\tau _p,\varphi ^{jq}\tau _q)=\frac{1}{2}\left\{\tau _p(\varphi ^{jq})\tau _q(\varphi ^{ip})\tau _q(\varphi ^{ip})\tau _p(\varphi ^{jq})\right\}$$
$`(\mathrm{5.7.4})`$
On the other hand, by (5.4.1) and (5.4.3),
$$\tau _i^{}(h(\tau _j^{}))=(\varphi ^{ip}\tau _p)(h^{jq}\omega _q)=\varphi ^{ip}\tau _p(h^{jq}\omega _q)+\tau _p,h^{jq}\omega _q\varphi ^{ip}=\varphi ^{ip}\tau _p(h^{jq})\omega _q+h^{jp}\varphi ^{ip}$$
where we have used (5.4.1). Thus, (5.7.3) takes the form
$$\frac{1}{2}\left\{\tau _p(\varphi ^{jq})\tau _q(\varphi ^{ip})\tau _q(\varphi ^{ip})\tau _p(\varphi ^{jq})\right\}$$
$$\varphi ^{ip}\tau _p(h^{jq})\omega _qh^{jp}\varphi ^{ip}+\varphi ^{jp}\tau _p(h^{iq})\omega _q+h^{ip}\varphi ^{jp}=\beta ^{ijr}\omega _r^{}$$
$`(\mathrm{5.7.5})`$
We have to prove that the matrix $`(h^{ij})`$ defined by (5.7.2) satisfies the differential equation (5.7.5). Using (5.4.4), rewrite (5.7.5) in the form
$$\frac{1}{2}\left\{\tau _p(\varphi ^{jq})\tau _l\tau _q(\varphi ^{ip})\tau _q(\varphi ^{ip})\tau _l\tau _p(\varphi ^{jq})\right\}$$
$$\varphi ^{ip}\tau _p(h^{jl})h^{jp}\tau _l(\varphi ^{ip})+\varphi ^{jp}\tau _p(h^{il})+h^{ip}\tau _l(\varphi ^{jp})=\beta ^{ijr}\rho ^{rl}\omega _l$$
$`(\mathrm{5.7.6})`$
Denote
$$A=\frac{1}{2}\left\{\tau _p(\varphi ^{jq})\tau _l\tau _q(\varphi ^{ip})\tau _q(\varphi ^{ip})\tau _l\tau _p(\varphi ^{jq})\right\},$$
$$B=B^{ij}=\varphi ^{ip}\tau _p(h^{jl})h^{jp}\tau _l(\varphi ^{ip}),$$
and
$$C=B^{ji}=\varphi ^{jp}\tau _p(h^{il})+h^{ip}\tau _l(\varphi ^{jp})$$
We have to prove that $`A+B+C=\beta ^{ijr}\rho ^{rl}\omega _l`$ where $`h^{ij}`$ is given by (5.7.2) and $`\beta ^{ijr}`$ is given by (5.5.3). Thus, we have
$$B=\varphi ^{ip}\tau _p\left\{\tau _q\tau _l(\varphi ^{jq})+\frac{1}{2}\tau _u(\varphi ^{jq})\tau _q(\varphi ^{ru})\rho ^{rl}\right\}\tau _l(\varphi ^{ip})\left\{\tau _q\tau _p(\varphi ^{jq})+\frac{1}{2}\tau _u(\varphi ^{jq})\tau _q(\varphi ^{ru})\rho ^{rp}\right\}$$
and
$$C=\varphi ^{jp}\tau _p\left\{\tau _q\tau _l(\varphi ^{iq})+\frac{1}{2}\tau _u(\varphi ^{iq})\tau _q(\varphi ^{ru})\rho ^{rl}\right\}+\tau _l(\varphi ^{jp})\left\{\tau _q\tau _p(\varphi ^{iq})+\frac{1}{2}\tau _u(\varphi ^{iq})\tau _q(\varphi ^{ru})\rho ^{rp}\right\}$$
Let us denote the $`n`$-th summand in an expression $`X`$ by $`Xn`$ (where we open the brackets). We have
$$B1+C1=\varphi ^{ip}\tau _p\tau _q\tau _l(\varphi ^{jq})+\varphi ^{jp}\tau _p\tau _q\tau _l(\varphi ^{iq})\stackrel{\left(\mathrm{5.4.9}\right)}{=}$$
$$=\tau _l(\varphi ^{jq})\tau _p\tau _q(\varphi ^{ip})+\tau _l(\varphi ^{iq})\tau _p\tau _q(\varphi ^{jp})=C3B3$$
Next,
$$B2=\frac{1}{2}\varphi ^{ip}\left\{\tau _p\tau _u(\varphi ^{jq})\tau _q(\varphi ^{ru})\rho ^{rl}+\tau _u(\varphi ^{jq})\tau _p\tau _q(\varphi ^{ru})\rho ^{rl}+\tau _u(\varphi ^{jq})\tau _q(\varphi ^{ru})\tau _p(\rho ^{rl})\right\}$$
We have
$$B21=\frac{1}{2}\varphi ^{ip}\tau _p\tau _u(\varphi ^{jq})\tau _q(\varphi ^{ru})\rho ^{rl}\stackrel{\left(\mathrm{5.4.7}\right)}{=}$$
$$=\frac{1}{2}\left\{\varphi ^{jp}\tau _p\tau _u(\varphi ^{iq})+\tau _u(\varphi ^{jp})\tau _p(\varphi ^{iq})\tau _u(\varphi ^{ip})\tau _p(\varphi ^{jq})\right\}\tau _q(\varphi ^{ru})\rho ^{rl}$$
Next,
$$B22=\frac{1}{2}\varphi ^{ip}\tau _p\tau _q(\varphi ^{ru})\tau _u(\varphi ^{jq})\rho ^{rl}\stackrel{\left(\mathrm{5.4.7}\right)}{=}$$
$$=\frac{1}{2}\tau _u(\varphi ^{jq})\rho ^{rl}\left\{\varphi ^{rp}\tau _p\tau _q(\varphi ^{iu})+\tau _q(\varphi ^{rp})\tau _p(\varphi ^{iu})\tau _q(\varphi ^{ip})\tau _p(\varphi ^{ru})\right\}=$$
$$=\frac{1}{2}\tau _u(\varphi ^{jq})\tau _l\tau _q(\varphi ^{iu})\frac{1}{2}\tau _u(\varphi ^{jq})\tau _q(\varphi ^{rp})\tau _p(\varphi ^{iu})\rho ^{rl}+\frac{1}{2}\tau _u(\varphi ^{jq})\tau _q(\varphi ^{ip})\tau _p(\varphi ^{ru})\rho ^{rl}$$
We see that $`B221=A1,B222=B213`$ and $`B223=B212`$. Similarly, $`A2=C221`$. We compute $`B23`$ using (5.4.10):
$$B23=\frac{1}{2}\varphi ^{ip}\tau _u(\varphi ^{jq})\tau _q(\varphi ^{ru})\rho ^{ra}\tau _p(\varphi ^{ba})\rho ^{bl}\stackrel{\left(\mathrm{5.4.6}\right)}{=}$$
$$=\frac{1}{2}\varphi ^{bp}\tau _p(\varphi ^{ia})\tau _u(\varphi ^{jq})\tau _q(\varphi ^{ru})\rho ^{ra}\rho ^{bl}=\frac{1}{2}\tau _l(\varphi ^{ia})\tau _u(\varphi ^{jq})\tau _q(\varphi ^{ru})\rho ^{ra}=B4$$
Finally, $`B211+C211=\beta ^{ijr}\rho ^{rl}\omega _l`$. Everything except these last terms cancels out, and this proves the Lemma. $``$
§6. Atiyah
6.1. Let $`𝒯=(A,T,\mathrm{\Omega },)`$ be a perfect extended Lie algebroid. Let $`𝔤,𝔤^{},𝔤^{\prime \prime }`$ be three frames in $`T`$. According to 5.2, we have the morphisms of the corresponding vertex prealgebroids
$$_{𝒯;𝔤^{\prime \prime }}\stackrel{g_{𝔤^{},𝔤^{\prime \prime }}}{}_{𝒯;𝔤^{}}\stackrel{g_{𝔤,𝔤^{}}}{}_{𝒯;𝔤}$$
$`(\mathrm{6.1.1})`$
as well as the morphism $`g_{𝔤;𝔤^{\prime \prime }}:_{𝒯;𝔤^{\prime \prime }}_{𝒯;𝔤}`$, all of them over $`Id_𝒯`$. Recall that $`Hom_{𝒫re𝒜lg_𝒯}(_{𝒯;𝔤^{\prime \prime }},_{𝒯;𝔤^{}})`$ is an $`\mathrm{\Omega }^2(𝒯)`$-torseur, cf. 4.3.1. We are aiming to compute the discrepancy
$$\alpha _{𝔤,𝔤^{},𝔤^{\prime \prime }}:=g_{𝔤;𝔤^{}}g_{𝔤^{},𝔤^{\prime \prime }}g_{𝔤,𝔤^{\prime \prime }}\mathrm{\Omega }^2(𝒯)$$
$`(\mathrm{6.1.2})`$
We have the functions $`h_{𝔤,𝔤^{}}`$, etc., acting from $`T`$ to $`\mathrm{\Omega }`$ (not $`A`$-linear!), as in the previous Section, which define our morphisms. By (3.7.1) the composition $`g_{𝔤;𝔤^{}}g_{𝔤^{},𝔤^{\prime \prime }}`$ is defined by the function $`h_{𝔤;𝔤^{}}+h_{𝔤^{};𝔤^{\prime \prime }}`$; therefore the discrepancy (6.1.2) is defined by the $`A`$-linear function
$$\alpha _{𝔤,𝔤^{},𝔤^{\prime \prime }}=h_{𝔤,𝔤^{}}+h_{𝔤^{},𝔤^{\prime \prime }}h_{𝔤,𝔤^{\prime \prime }}\mathrm{\Omega }^2(𝒯)Hom_A(T,\mathrm{\Omega })$$
$`(\mathrm{6.1.3})`$
which by definition coincides with (6.1.2).
Note that the functions (6.1.3) obviously satisfy the $`2`$-cocycle condition
$$\alpha _{𝔤^{},𝔤^{\prime \prime },𝔤^{\prime \prime \prime }}\alpha _{𝔤,𝔤^{\prime \prime },𝔤^{\prime \prime \prime }}+\alpha _{𝔤,𝔤^{},𝔤^{\prime \prime \prime }}\alpha _{𝔤,𝔤^{},𝔤^{\prime \prime }}=0$$
$`(\mathrm{6.1.4})`$
6.2. Choose some bases $`\{\tau _i\},\{\tau _i^{}\}`$ and $`\{\tau _i^{\prime \prime }\}`$ of our frames $`𝔤,𝔤^{},𝔤^{\prime \prime }`$; let $`\{\omega _i\},\{\omega _i^{}\},\{\omega _i^{\prime \prime }\}`$ be the dual bases in $`\mathrm{\Omega }`$. Define two matrices $`\varphi ,\psi GL_n(A)`$ by $`\tau _i^{}=\varphi ^{ij}\tau _j,\tau _i^{\prime \prime }=\psi ^{ij}\tau _j^{}`$ and set $`\rho :=\varphi ^{1t},\sigma :=\psi ^{1t}`$.
The maps $`h_{𝔤,𝔤^{}}`$, etc., are defined by the matrices $`(h_{𝔤,𝔤^{}}^{ij})`$, etc., where $`h_{𝔤,𝔤^{}}(\tau _i^{})=h_{𝔤,𝔤^{}}^{ij}\omega _j`$, etc.
By $`(3.5\gamma )`$ and $`(\mathrm{5.7.1})_\gamma `$ we have
$$h_{𝔤,𝔤^{}}(a\tau _i^{})=ah_{𝔤,𝔤^{}}(\tau _i^{})\gamma (a,\tau _i^{})=ah_{𝔤,𝔤^{}}(\tau _i^{})\gamma (a,\varphi ^{ip}\tau _p)=$$
$$=ah_{𝔤,𝔤^{}}(\tau _i^{})+\tau _p(a)\varphi ^{ip}+\tau _p(\varphi ^{ip})a$$
Therefore,
$$A:=h_{𝔤,𝔤^{}}(\tau _i^{\prime \prime })=h_{𝔤,𝔤^{}}(\psi ^{ij}\tau _j^{})=\psi ^{ij}h_{𝔤,𝔤^{}}(\tau _j^{})+\tau _p(\psi ^{ij})\varphi ^{jp}+\tau _p(\varphi ^{jp})\psi ^{ij};$$
$$B:=h_{𝔤^{},𝔤^{\prime \prime }}(\tau _i^{\prime \prime })=h_{𝔤^{},𝔤^{\prime \prime }}^{ia}\omega _a^{}=h_{𝔤^{},𝔤^{\prime \prime }}^{ia}\rho ^{al}\omega _l=$$
by (5.7.2)
$$=\left\{\tau _p^{}\tau _a^{}(\psi ^{ip})+\frac{1}{2}\tau _q^{}(\psi ^{ip})\tau _p^{}(\psi ^{rq})\sigma ^{ra}\right\}\rho ^{al}\omega _l$$
and
$$C=h_{𝔤,𝔤^{\prime \prime }}(\tau _i^{\prime \prime })=h_{𝔤,𝔤^{\prime \prime }}^{il}\omega _l=\left\{\tau _p\tau _l((\psi \varphi )^{ip})+\frac{1}{2}\tau _q((\psi \varphi )^{ip})\tau _p((\psi \varphi )^{rq})(\sigma \rho )^{rl}\right\}\omega _l$$
We have to calculate $`\alpha _{𝔤,𝔤^{},𝔤^{\prime \prime }}(\tau _i^{\prime \prime })=A+B+C`$. In the computation we shall use the same convention for the notation of various terms in $`A,B,C,\mathrm{}`$ as in 5.7.
So, we have
$$A1=\psi ^{ij}h_{𝔤,𝔤^{}}^{jl}\omega _l=\psi ^{ij}\left\{\tau _p\tau _l(\varphi ^{jp})+\frac{1}{2}\tau _q(\varphi ^{jp})\tau _p(\varphi ^{rq})\rho ^{rl}\right\}\omega _l$$
By (5.4.4),
$$A2=\tau _p(\psi ^{ij})\tau _l(\varphi ^{jp})\omega _l$$
and
$$A3=\tau _p(\varphi ^{jp})\tau _l(\psi ^{ij})\omega _l$$
Next,
$$B1=\varphi ^{pu}\tau _u\varphi ^{av}\tau _v(\psi ^{ip})\rho ^{al}\omega _l=\varphi ^{pu}\tau _u\tau _l(\psi ^{ip})\omega _l+\varphi ^{pu}\tau _u(\varphi ^{av})\tau _v(\psi ^{ip})\rho ^{al}\omega _l$$
and
$$B2=\frac{1}{2}\varphi ^{qu}\varphi ^{pv}\tau _u(\psi ^{ip})\tau _v(\psi ^{rq})\sigma ^{ra}\rho ^{al}\omega _l$$
Finally,
$$C1=\tau _p\tau _l(\psi ^{iu}\varphi ^{up})\omega _l=\varphi ^{up}\tau _p\tau _l(\psi ^{iu})\omega _l\psi ^{iu}\tau _p\tau _l(\varphi ^{up})\omega _l$$
$$\tau _p(\psi ^{iu})\tau _l(\varphi ^{up})\omega _l\tau _l(\psi ^{iu})\tau _p(\varphi ^{up})\omega _l$$
and
$$C2=\frac{1}{2}\left\{\varphi ^{up}\tau _q(\psi ^{iu})+\psi ^{iu}\tau _q(\varphi ^{up})\right\}\left\{\varphi ^{vq}\tau _p(\psi ^{rv})+\psi ^{rv}\tau _p(\varphi ^{vq})\right\}\sigma ^{rs}\rho ^{sl}\omega _l=$$
$$=\frac{1}{2}\{\varphi ^{up}\varphi ^{vq}\tau _q(\psi ^{iu})\tau _p(\psi ^{rv})\sigma ^{rs}\rho ^{sl}+\psi ^{iu}\varphi ^{vq}\tau _p(\psi ^{rv})\tau _q(\varphi ^{up})\sigma ^{rs}\rho ^{sl}+$$
$$+\varphi ^{up}\tau _q(\psi ^{iu})\tau _p(\varphi ^{sq})\rho ^{sl}+\psi ^{iu}\tau _q(\varphi ^{up})\tau _p(\varphi ^{sq})\rho ^{sl}\}\omega _l$$
We see first of all the terms of the second order cancel out, as they should: $`A11=C12`$ and $`B11=C11`$. Most of the other terms also cancel out, and in $`A+B+C`$ we are left only with
$$B12+C23=\frac{1}{2}\varphi ^{up}\tau _p(\varphi ^{sq})\tau _q(\psi ^{iu})\rho ^{sl}\omega _l\stackrel{\left(\mathrm{5.4.6}\right)}{=}$$
$$=\frac{1}{2}\varphi ^{sp}\tau _p(\varphi ^{uq})\tau _q(\psi ^{iu})\rho ^{sl}\omega _l=\frac{1}{2}\tau _l(\varphi ^{uq})\tau _q(\psi ^{iu})\omega _l$$
and
$$C22=\frac{1}{2}\psi ^{iu}\varphi ^{vq}\tau _q(\varphi ^{up})\tau _p(\psi ^{rv})\sigma ^{rs}\rho ^{sl}\omega _l\stackrel{\left(\mathrm{5.4.6}\right)}{=}\frac{1}{2}\psi ^{iu}\varphi ^{uq}\tau _q(\varphi ^{vp})\tau _p(\psi ^{rv})\sigma ^{rs}\rho ^{sl}\omega _l$$
Thus, we have
$$\alpha _{𝔤,𝔤^{},𝔤^{\prime \prime }}(\tau _i^{\prime \prime })=\frac{1}{2}\tau _l(\varphi ^{uq})\tau _q(\psi ^{iu})\omega _l\frac{1}{2}\psi ^{iu}\varphi ^{uq}\tau _q(\varphi ^{vp})\tau _p(\psi ^{rv})\sigma ^{rs}\rho ^{sl}\omega _l$$
Rewriting the right hand side in the base $`\{\omega _i^{\prime \prime }\}`$ we get
$$\alpha (\tau _i^{\prime \prime }):=\alpha _{𝔤,𝔤^{},𝔤^{\prime \prime }}(\tau _i^{\prime \prime })=\alpha ^{ir}\omega _r^{\prime \prime }=\alpha _1^{ir}\omega _r^{\prime \prime }+\alpha _2^{ir}\omega _r^{\prime \prime }$$
$`(\mathrm{6.2.1})`$
where
$$\alpha _1^{ir}=\frac{1}{2}\psi ^{ra}\varphi ^{al}\tau _l(\varphi ^{uq})\tau _q(\psi ^{iu})$$
$`(\mathrm{6.2.2})`$
and
$$\alpha _2^{ir}=\frac{1}{2}\psi ^{iu}\varphi ^{uq}\tau _q(\varphi ^{vp})\tau _p(\psi ^{rv})$$
$`(\mathrm{6.2.3})`$
So we see that $`\alpha _1^{ir}=\alpha _2^{ri}`$, i.e. the matrix $`\alpha ^{ir}`$ is skew symmetric, that is, $`\alpha \mathrm{\Omega }^2(𝒯)`$ as it should be.
6.3. Let us rewrite the expression (6.2.3) in terms of vector fields $`\tau _i^{\prime \prime }`$:
$$\alpha _2^{ir}=\frac{1}{2}\tau _i^{\prime \prime }(\varphi ^{vp})\varphi ^{1pa}\psi ^{1ab}\tau _b^{\prime \prime }(\psi ^{rv})=$$
(using the identity $`\tau (\varphi )=\varphi \tau (\varphi ^1)\varphi `$)
$$=\frac{1}{2}\varphi ^{vs}\tau _i^{\prime \prime }(\varphi ^{1sa})\psi ^{1ab}\psi ^{ru}\tau _b^{\prime \prime }(\psi ^{1uc})\psi ^{cv}$$
Using (5.4.13), we have
$$\psi ^{1ab}\tau _b^{\prime \prime }(\psi ^{1uc})=\tau _a^{}(\psi ^{1uc})=\tau _u^{}(\psi ^{1ac})=\psi ^{1ub}\tau _b^{\prime \prime }(\psi ^{1ac})$$
whence
$$\alpha _2(\tau _i^{\prime \prime }):=\alpha _2^{ir}\omega _r^{\prime \prime }=\frac{1}{2}\varphi ^{vs}\tau _i^{\prime \prime }(\varphi ^{1sa})\psi ^{ru}\psi ^{1ub}\tau _b^{\prime \prime }(\psi ^{1ac})\psi ^{cv}\omega _r^{\prime \prime }=$$
$$=\frac{1}{2}\varphi ^{vs}\tau _i^{\prime \prime }(\varphi ^{1sa})\tau _r^{\prime \prime }(\psi ^{1ac})\psi ^{cv}\omega _r^{\prime \prime }=\frac{1}{2}tr\left\{\varphi \tau _i^{\prime \prime }(\varphi ^1)\tau _r^{\prime \prime }(\psi ^1)\psi \right\}\omega _r^{\prime \prime }$$
Skew symmetrizing, we arrive at the first part of
6.4. Theorem. (a) The cocycle $`\alpha _{𝔤,𝔤^{},𝔤^{\prime \prime }}`$, (6.1.3) is given in coordinates by the expression
$$\alpha _{𝔤,𝔤^{},𝔤^{\prime \prime }}(\tau _i^{\prime \prime })=\alpha (\psi ,\varphi )(\tau _i^{\prime \prime })=$$
$$=\frac{1}{2}tr\left\{\tau _i^{\prime \prime }(\psi ^1)\psi \varphi \tau _r^{\prime \prime }(\varphi ^1)\tau _r^{\prime \prime }(\psi ^1)\psi \varphi \tau _i^{\prime \prime }(\varphi ^1)\right\}\omega _r^{\prime \prime }$$
$`(\mathrm{6.4.1})`$
(b) The $`3`$-form $`\beta _{𝔤,𝔤^{}}=\beta (\varphi )`$ (5.5.3) is equal to
$$\beta _{𝔤,𝔤^{}}(\tau _i^{},\tau _j^{})=\beta (\varphi )(\tau _i^{},\tau _j^{})=$$
$$=\frac{1}{2}tr\left\{\tau _i^{}(\varphi ^1)\varphi \tau _j^{}(\varphi ^1)\varphi \tau _r^{}(\varphi ^1)\varphi \tau _j^{}(\varphi ^1)\varphi \tau _i^{}(\varphi ^1)\varphi \tau _r^{}(\varphi ^1)\varphi \right\}\omega _r^{}$$
$`(\mathrm{6.4.2})`$
Part (b) is proven by the same argument as in 6.3, and we leave it to the reader.
From the expression (6.4.2) we see directly that the form $`\beta (\varphi )`$ is closed, and from (6.4.1) one checks that the form $`\alpha (\varphi ,\psi )`$ satisfies the group (or Cech) cocycle condition
$$\alpha (\psi ,\chi )\alpha (\varphi \psi ,\chi )+\alpha (\varphi ,\psi \chi )\alpha (\varphi ,\psi )=0$$
$`(\mathrm{6.4.3})`$
for all $`\varphi ,\psi ,\chi GL_n(A)`$, which is the same as (6.1.4), and
$$d\alpha (\psi ,\varphi )=\beta (\varphi )+\beta (\psi )\beta (\psi \varphi )$$
$`(\mathrm{6.4.4})`$
Thus, a couple $`c(𝒯)=(\alpha ,\beta )`$ represents a cohomology class in $`H^2(GL_n(A),\mathrm{\Omega }^{[2,3}(𝒯))`$ of the group $`GL_n(A)`$ with coefficients in the complex
$$\mathrm{\Omega }^{[2,3}(𝒯):=\mathrm{\Omega }^2(𝒯)\mathrm{\Omega }^{3,cl}(𝒯)$$
$`(\mathrm{6.4.5})`$
where $`\mathrm{\Omega }^2(𝒯)`$ sits in degree $`0`$ (the action of $`GL_n(A)`$ on $`\mathrm{\Omega }^{[2,3}`$ being trivial).
The cocycle $`\alpha `$ is classical, and essentially goes back to Atiyah; it is written down explicitly by Harris, \[H\], p. 280. The class $`\beta `$ resembles ”Chern-Simons” form.
The whole cocycle $`c(𝒯)`$ may be thought of as an integration of a cocycle $`\stackrel{~}{c}`$ from \[MSV\], (5.16), (5.17).
§7. Gerbes of Vertex Algebroids
7.1. Let us reformulate the results of the last three Sections in language of Torseurs. This reformulation was inspired by \[BD1\].
Let $`𝒯=(A,T,\mathrm{})`$ be a quasiregular extended Lie algebroid (see 5.1). Let us define a groupoid $`\mathrm{\Omega }_𝒯^{[2,3}`$ is follows. We set $`Ob\mathrm{\Omega }_𝒯^{[2,3}=\mathrm{\Omega }^{3,cl}(𝒯)`$; for $`\omega _1,\omega _2\mathrm{\Omega }^{3,cl}(𝒯)`$ a morphism $`\omega _1\omega _2`$ is by definition a two-form $`\eta \mathrm{\Omega }^2(𝒯)`$ such that $`d_{DR}(\eta )=\omega _2\omega _1`$. The composition of morphisms is defined in an obvious manner. Note that $`\mathrm{\Omega }_𝒯^{[2,3}`$ is in fact an abelian group in categories.
Similarly, let $`\mathrm{\Omega }_𝒯^{2,cl}`$ denote a groupoid with the unique object and the set of morphisms equal to $`\mathrm{\Omega }^{2,cl}(𝒯)`$. It is also an abelian group in categories. We have a fully faithful monoidal functor
$$\mathrm{\Omega }_𝒯^{2,cl}\mathrm{\Omega }_𝒯^{[2,3}$$
$`(\mathrm{7.1.1})`$
sending the unique object in $`\mathrm{\Omega }_𝒯^{2,cl}`$ to $`0\mathrm{\Omega }^{3,cl}(𝒯)`$.
Consider the groupoid $`𝒜lg_𝒯`$. According to 5.3, it is nonempty.
We can define an Action of $`\mathrm{\Omega }_𝒯^{[2,3}`$ on $`𝒜lg_𝒯`$, i.e. a monoidal functor
$$\stackrel{}{+}:𝒜lg_𝒯\times \mathrm{\Omega }_𝒯^{[2,3}𝒜lg_𝒯$$
$`(\mathrm{7.1.2})`$
as follows. For $`𝒜𝒜lg,\omega \mathrm{\Omega }^{3,cl}(𝒯)`$, a couple $`(𝒜,\omega )`$ goes to the vertex algebroid $`𝒜\stackrel{}{+}\omega `$ defined in 4.2. If $`\eta \mathrm{\Omega }^2(𝒯)`$, $`d_{DR}(\eta )=\omega ^{}\omega `$ then a morphism
$`\stackrel{}{+}(\eta ):𝒜\stackrel{}{+}\omega 𝒜\stackrel{}{+}\omega ^{}`$ is defined according to 4.2.2.
Let us fix $`𝒜`$ and consider the functor
$$\mathrm{\Omega }_𝒯^{[2,3}𝒜lg_𝒯,\omega 𝒜\stackrel{}{+}\omega $$
$`(\mathrm{7.1.3})`$
induced by (7.1.2). By 4.2.2 this functor is fully faithful.
Let $`𝒜^{}𝒜lg_𝒯`$ be another object. By 5.2, there exists a morphism of vertex prealgebroids $`g:P(𝒜)P(𝒜^{})`$ lying over $`Id_𝒯`$; it is necessarily an isomorphism. Consider the vertex algebroid $`g_{}𝒜`$ constructed in Theorem 4.5. By definition, $`g`$ is lifted to an isomorphism $`\stackrel{~}{g}:𝒜\stackrel{}{}g_{}𝒜`$. Since $`g_{}𝒜𝒜lg_{P(𝒜^{})}`$, by 4.2.1 $`𝒜^{}=g_{}𝒜\stackrel{}{+}\omega `$ for some $`\omega \mathrm{\Omega }^{3,cl}(𝒯)`$. Therefore $`\stackrel{~}{g}`$ induces an isomorphism $`\stackrel{~}{g}\stackrel{}{+}\omega :𝒜\stackrel{}{+}\omega \stackrel{}{}g_{}𝒜\stackrel{}{+}\omega =𝒜^{}`$.
In other words, we have checked that (7.1.3) is surjective on isomorphism classes of objects, hence it is an equivalence of categories. This proves
7.2. Theorem. Let $`𝒯`$ be a quasiregular extended Lie algebroid. Then the Action (7.1.2) makes the groupoid $`𝒜lg_𝒯`$ a nonempty $`\mathrm{\Omega }_𝒯^{[2,3}`$-Torseur. $``$
7.3. Let $`(X,𝒪_X)`$ be a topological space ringed by a sheaf of commutative $`k`$-algebras $`𝒪_X`$. Let us call an extended Lie $`𝒪_X`$ Lie algebroid $`𝒯=(𝒪_X,T,\mathrm{})`$ (quasi)regular if there exists an open covering $`X=U_i`$ of $`X`$ such that all $`𝒯(U_i)`$ are (quasi)regular.
For example, if $`T`$ is a Lie $`𝒪_X`$-algebroid which is locally free as an $`𝒪_X`$-module then the corresponding extended algebroid $`𝒯_T`$ is quasiregular. If $`X`$ is a smooth $`k`$-scheme of finite type and $`T=T_{X/k}`$ is the sheaf of vector fields then $`𝒯_T`$ is regular.
We can sheafify the constructions of the previous Subsections and obtain the sheaves (champs) of groupoids (i.e. gerbes) $`𝒜lg_𝒯`$, etc.
Let $`𝒯`$ be quasiregular. Consider the gerbe $`𝒜lg_𝒯`$. According to 7.1 and 7.2, it is locally nonempty but not locally connected in general; its sheaf of connected components is an $`H_{DR}^3(𝒯)`$-torseur.
By the general procedure this gerbe defines a characteristic class
$$c(𝒯):=c(𝒜lg_𝒯)H^2(X;\mathrm{\Omega }^{[2,3}(𝒯))$$
$`(\mathrm{7.3.1})`$
Here in the right hand side we consider the hypercohomology with coefficients in the complex $`\mathrm{\Omega }^{[2,3}(𝒯)`$.
Let us explain how to define the class (7.3.1). Choose an open covering $`𝒰=\{U_i\}`$ of $`X`$ such that all groupoids $`𝒜lg_{𝒯(U_i)}`$ are nonempty. Choose an object $`𝒜_i`$ in each $`𝒜lg_{𝒯(U_i)}`$. Over double intersections $`U_{ij}:=U_iU_j`$ we get two objects $`𝒜_i|_{U_{ij}},𝒜_j|_{U_{ij}}𝒜lg_{𝒯(U_{ij})}`$. Choose $`3`$-forms $`\omega _{ij}\mathrm{\Omega }^{3,cl}(𝒯(U_{ij})`$ such that there exist isomorphisms
$$h_{ij}:𝒜_j|_{U_{ij}}\stackrel{}{}𝒜_i|_{U_{ij}}\stackrel{}{+}\omega _{ij}$$
$`(\mathrm{7.3.2})`$
Choose some isomorphisms (7.3.2). Then on triple intersections we get isomorphisms
$$𝒜_i|_{U_{ijk}}\stackrel{}{+}\omega _{ij}\stackrel{}{+}\omega _{ik}\stackrel{}{}𝒜_i|_{U_{ijk}}\stackrel{}{+}\omega _{ik}$$
$`(\mathrm{7.3.3})`$
The isomorphisms (7.5.3) must be given by the $`2`$-forms $`\eta _{ijk}\mathrm{\Omega }^2(𝒯)`$ such that $`d_{DR}(\eta _{ijk})=c_{ij}c_{ik}+c_{jk}`$. Then $`(\omega _{ij},\eta _{ijk})`$ is a $`2`$-cocycle in the Cech complex $`C^{}(𝒰;\mathrm{\Omega }^{[2,3}(𝒯))`$ representing the class (7.3.1).
7.4. Now let us assume that the groupoid $`𝒯`$ is regular. Theorem 6.4 calculates the class (7.3.1).
Namely, define the ”Atiyah-Chern-Simons” class $`ch_2(𝒯)H^2(X;\mathrm{\Omega }^{[2,3}(𝒯))`$ by the following procedure. Choose some bases of local sections $`\tau ^{(i)}=\{\tau _\alpha ^{(i)}\}T(U_i)`$ over some open covering. Let $`\varphi _{ij}GL_n(𝒪(U_{ij}))`$ be the transition matrix from $`\tau ^{(i)}`$ to $`\tau ^{(j)}`$ over $`U_{ij}`$. By definition, $`ch_2(𝒯)`$ is represented by the Cech $`2`$-cocycle
$$2ch_2(𝒯):=(\alpha (𝒯),\beta (𝒯))$$
$`(\mathrm{7.4.1})`$
where
$$\alpha (𝒯)=(\alpha (𝒯)_{ijk})=tr(\varphi _{ij}^1\varphi _{ij}\varphi _{jk}\varphi _{jk}^1)\}Z^2(𝒰;\mathrm{\Omega }^2(𝒯))$$
$`(\mathrm{7.4.2})`$
and
$$\beta (𝒯)=(\beta (𝒯)_{ij})=\{\frac{1}{3}tr(\varphi _{ij}\varphi _{ij}^1\varphi _{ij}\varphi _{ij}^1\varphi _{ij}\varphi _{ij}^1)\}C^1(𝒰;\mathrm{\Omega }^{3,cl}(𝒯))$$
$`(\mathrm{7.4.3})`$
Theorem 6.4 implies
7.5. Theorem. Let $`𝒯`$ be a regular extended Lie $`𝒪_X`$-algebroid. Then $`c(𝒯)=2ch_2(𝒯)`$.
Therefore, the gerbe $`𝒜lg_𝒯`$ admits a global section iff $`ch_2(𝒯)=0`$. If so, then the groupoid of global sections $`𝒜lg_𝒯(X)`$ is equivalent to the groupoid of $`\mathrm{\Omega }^{[2,3}(𝒯)`$-torseurs, whence $`\pi _0(𝒜lg_𝒯(X))\stackrel{}{=}H^1(X;\mathrm{\Omega }^{[2,3}(𝒯))`$ and the automorphism group of an object of this groupoid is isomorphic to $`H^0(X,\mathrm{\Omega }^{[2,3}(𝒯))`$.
$``$
It is instructive to compare the previous discussion with \[BB\], 2.1. In the chiral situation the degree of cohomology goes one step up.
7.6. Let us identify the class $`c(𝒯)`$ when $`X`$ is a smooth $`k`$-scheme and $`𝒯=𝒯_X`$ is the tangent bundle. If $`E`$ is an arbitrary vector bundle over $`X`$ given by a Cech $`1`$-cocycle
$$g=(g_{ij})Z^1(𝒰;GL_r(𝒪_X))$$
$`(\mathrm{7.6.1})`$
on some open covering $`𝒰`$ then (7.4.1) - (7.4.3) define a cocycle
$$c(g)=(\alpha (g),\beta (g))Z^2(𝒰;\mathrm{\Omega }_X^{[2,3})$$
$`(\mathrm{7.6.2})`$
If $`g_{ij}=\varphi _ih_{ij}\varphi _j^1`$ for some $`\varphi =(\varphi _i)C^1(𝒰;GL(𝒪_X)`$ then one checks by a direct computation that
7.6.1. Claim. (a)
$$\alpha (g)\alpha (h)=d_{Cech}\eta $$
$`(\mathrm{7.6.3})`$
where $`\eta =\eta (h,\varphi )=(\eta _{ij})C^2(𝒰;\mathrm{\Omega }_X^2)`$ is given by
$$\eta _{ij}=tr\left\{h_{ij}^1dh_{ij}\varphi _j^1d\varphi _j\varphi _i^1d\varphi _idh_{ij}h_{ij}^1+h_{ij}^1\varphi _i^1d\varphi _ih_{ij}\varphi _j^1d\varphi _j\right\}$$
$`(\mathrm{7.6.4})`$
(b)
$$d_{DR}\eta =\beta (g)\beta (h)d_{Cech}\gamma $$
$`(\mathrm{7.6.5})`$
where $`\gamma =\gamma (h,\varphi )=(\gamma _i)C^0(𝒰;\mathrm{\Omega }^{3,cl})`$ is defined by
$$\gamma _i=\frac{1}{3}tr((\varphi _i^1d\varphi _i)^3)$$
$`(\mathrm{7.6.6})`$
$``$
7.6.2. Corollary.
$$c(g)c(h)=de$$
$`(\mathrm{7.6.7})`$
where
$$e=(\eta ,\gamma )C^1(𝒰;\mathrm{\Omega }_X^{[2,3})$$
$`(\mathrm{7.6.8})`$
$``$
Therefore (7.6.2) gives rise to a well defined characteristic class $`c(E)H^2(X,\mathrm{\Omega }_X^{[2,3})`$. The following Lemma is obvious.
7.7. Lemma. (a) If $`f:YX`$ is an arbitrary morphism from another smooth scheme then $`c(f^{}E)=f^{}c(E)`$.
(b) If $`0E^{}EE^{\prime \prime }0`$ is a short exact sequence of vector bundles over $`X`$ then $`c(E)=c(E^{})+c(E^{\prime \prime })`$.
(c) If $`L`$ is a line bundle then $`c(L)`$ is equal to the image of $`c_1(L)c_1(L)H^1(X;𝒪_X^{})^2`$ under the composition
$$H^1(X;𝒪_X^{})^2\stackrel{dlog^2}{}H^1(X;\mathrm{\Omega }_X^{1,cl})^2H^2(X;\mathrm{\Omega }_X^{2,cl})H^2(X;\mathrm{\Omega }_X^{[2,3})$$
7.8. Recall (cf. \[S\]) that for an arbitrary $`E`$ we have a characteristic class
$$2ch_2^{(K)}(E):=c_1^{(K)2}(E)2c_2^{(K)}H^2(X;𝒦_{2,X})$$
from which we can get a class $`2ch_2(E)H^2(X;\mathrm{\Omega }_X^{[2,3})`$ using the $`dlog`$ map
$$H^2(X;𝒦_2)H^2(X;\mathrm{\Omega }_X^{2,cl})H^2(X,\mathrm{\Omega }_X^{[2,3})$$
The class $`2ch_2(E)`$ also satisfies properties 7.7 (a) - (c).
It is obvious that the natural map
$$H^2(X;\mathrm{\Omega }_X^{[2,3})H^2(X;\mathrm{\Omega }_X^{[2})$$
is injective, where $`\mathrm{\Omega }_X^{[2}:=\mathrm{\Omega }_X^2\mathrm{\Omega }_X^3\mathrm{}`$ is the stupidly truncated (and shifted, so that $`\mathrm{\Omega }_X^2`$ sits in degree $`0`$) de Rham complex.
The proof of the following lemma was provided to us by H. Esnault.
7.9. Lemma. The inverse image map
$$H^2(X;\mathrm{\Omega }_X^{[2})H^2((E);\mathrm{\Omega }_{(E)}^{[2})$$
is injective. $``$
In fact, more is true. One can consider the cohomology theory which assigns to a smooth $`X`$ a collection of cohomology groups $`\{H^i(X;\mathrm{\Omega }^{[i})\}`$. This theory has the standard Grothendieck’s properties needed to define the Chern classes, cf. \[Gr\]. This fact was communicated to us by A. Beilinson.
Anyway, the inverse image map $`H^2(X;\mathrm{\Omega }_X^{[2,3})H^2((E);\mathrm{\Omega }_{(E)}^{[2,3})`$ is also injective, hence by splitting principle we get
7.10. Theorem. For all vector bundles $`E`$, $`c(E)=2ch_2(E)`$. $``$
This theorem was obtained in collaboration with H. Esnault.
7.11. Corollary. In the situation of 7.6 the class described in Theorem 7.5 is equal to $`2ch_2(𝒯_X)`$. $``$
§8. Vertex Envelope of a Conformal Algebra
8.1. Our aim in this Section will be to construct a left adjoint $`U`$ to the forgetful functor (0.8.1) and to prove the ”Poincaré-Birkhoff-Witt” theorem for algebras $`UC,C𝒞onf`$.
Let $`V`$ be a vertex algebra. We have the following two particular cases of the OPE formula (0.5.12). The first one corresponds to $`m=n=1`$:
$$[x_{(1)},y_{(1)}]=\underset{j0}{}(1)^j^{(j+1)}(x_{(j)}y)_{(1)}$$
$`(\mathrm{8.1.1})`$
where we have used (0.5.10). The second one corresponds to $`m0,n=1`$:
$$[x_{(m)},y_{(1)}]=\underset{j=0}{\overset{m}{}}\left(\genfrac{}{}{0pt}{}{m}{j}\right)(x_{(j)}y)_{(mj1)}$$
$`(\mathrm{8.1.2})`$
8.2. Let $`C=C_i`$ be a conformal algebra. Let $`TC=_{j_0}T^jC,T^jC:=C^j`$ be the tensor algebra of $`C`$ over $`k`$. The multiplication in $`TC`$ will be denoted $`xy`$ or $`xy`$. The $`_0`$-grading of $`C`$ induces the $`_0`$-grading of $`TC`$ such that $`TC`$ becomes a $`_0`$-graded associative algebra with the unit $`\text{1}T^0C=k`$. We have a canonical embedding of $`k`$-modules $`C=T^1CTC`$.
There is a unique extension of the operators $`^{(j)}`$ on $`C`$ to the whole space $`TC`$ satisfying
$$^{(j)}(xy)=\underset{p=0}{\overset{j}{}}^{(p)}x^{(jp)}y$$
$`(\mathrm{8.2.0})`$
cf. (0.4.6). These operators will satisfy (0.4.1).
Let $`RTC`$ be a two-sided ideal generated by all elements
$$r(x,y):=xyyx\underset{j0}{}(1)^j^{(j+1)}(x_{(j)}y),x,yC,$$
$`(\mathrm{8.2.1})`$
cf. (8.1.1). Set $`UC=TC/R`$. We have a canonical morphism $`i:CUC`$ equal to the restriction of the projection $`p:TCUC`$ to $`C`$.
Since the relations (8.2.1) are homogeneous, the algebra $`UC`$ inherits a $`_0`$-grading from $`TC`$. Using (0.4.6) and an obvious identity $`^{(i)}^{(j)}=^{(j)}^{(i)}`$ one sees easily that the operators $`^{(i)}`$ respect the ideal $`R`$. Hence they induce the operators $`^{(i)}`$ on $`UC`$ of degree $`i`$ which satisfy (0.4.1).
8.3. Theorem. There is a unique structure of a vertex algebra on the $`k`$-module $`UC`$ such that for all $`xC,zTC`$,
$$p(xz)=i(x)_{(1)}p(z)$$
$`(\mathrm{8.3.1})`$
The corresponcence $`CUC`$ defines a functor $`U:𝒞onf𝒱ert`$ left adjoint to the forgetful functor.
The algebra $`UC`$ will be called a vertex envelope of a conformal algebra $`C`$.
This Theorem will be proven in 8.20, after some preparation.
8.4. Let us define $`k`$-linear operators $`x_{(j)},xC,j`$ of degree $`j1`$ acting on the module $`TC`$. If $`j<0,j=n1`$, we set
$$x_{(n1)}z=(^{(n)}x)z$$
$`(\mathrm{8.4.1})`$
In particular, $`x_{(1)}z=xz`$.
Each element of $`TC`$ is a linear combination of monomials $`z=z_1z_2\mathrm{}z_n,z_iC`$. We define $`x_{(j)}z,j_0`$ by induction on the length $`n`$ of the monomial. If $`n=1`$, i.e. $`zC`$, then we already have $`x_{(j)}z`$ due to the structure of a conformal algebra on $`C`$.
If $`z=yu,yC`$, we set
$$x_{(j)}yu=yx_{(j)}u+\underset{p=0}{\overset{j}{}}\left(\genfrac{}{}{0pt}{}{j}{p}\right)(x_{(p)}y)_{(jp1)}u$$
$`(\mathrm{8.4.2})`$
cf. (8.1.2).
We leave to the reader an easy proof of the lemma below.
8.4.1. Lemma. For all $`i_0,n`$,
$$(^{(i)}x)_{(n)}=(1)^i\left(\genfrac{}{}{0pt}{}{n}{i}\right)x_{(ni)}$$
$`(\mathrm{8.4.3})`$
8.5. Lemma. The operators $`x_{(j)}`$ introduced above respect the ideal $`R`$.
The proof will be given in 8.6 — 8.8 below.
8.6. Since $`R`$ is a left ideal, it is respected by all operators $`x_{(n)}`$ with $`n<0`$.
It follows from the commutation formula (8.4.2) that if all the operators $`x_{(n)},n0,`$ respect a subset $`STC`$ then they respect the left ideal generated by $`S`$. Therefore we need to prove that
$$u_{(n)}xyzu_{(n)}yxz\underset{j0}{}(1)^ju_{(n)}^{(j+1)}(x_{(j)}y)zR$$
$`(\mathrm{8.6.1})`$
for all $`u,x,yC,zTC,n0`$. Let us denote the summands in (8.6.1) by $`A,B`$ and $`C`$. We have
$$A=u_{(n)}xyz=(u_{(n)}x)yz+xu_{(n)}yz+\underset{p=o}{\overset{n1}{}}\left(\genfrac{}{}{0pt}{}{n}{p}\right)(u_{(p)}x)_{(np1)}yz$$
$`(\mathrm{8.6.2})`$
We shall use the same agreement as in 5.7: in an expression $`X`$, the $`n`$-th summand will be denoted by $`Xn`$. We have
$$A2=x(u_{(n)}y)z+xyu_{(n)}z+\underset{q=0}{\overset{n1}{}}\left(\genfrac{}{}{0pt}{}{n}{q}\right)x(u_{(q)}y)_{(nq1)}z$$
$`(\mathrm{8.6.3})`$
$$A3=\underset{p=0}{\overset{n1}{}}\left(\genfrac{}{}{0pt}{}{n}{p}\right)\{\left((u_{(p)}x)_{(np1)}y\right)z+y(u_{(p)}x)_{(np1)}z+$$
$$+\underset{q=0}{\overset{np2}{}}\left(\genfrac{}{}{0pt}{}{np1}{q}\right)((u_{(p)}x)_{(q)}y)_{(npq2)}z\}$$
$`(\mathrm{8.6.4})`$
Similarly,
$$B=u_{(n)}yxz=(u_{(n)}y)xzyu_{(n)}xz\underset{p=0}{\overset{n1}{}}\left(\genfrac{}{}{0pt}{}{n}{p}\right)(u_{(p)}y)_{(np1)}xz$$
$`(\mathrm{8.6.5})`$
$$B2=y(u_{(n)}x)zyxu_{(n)}z\underset{p=0}{\overset{n1}{}}\left(\genfrac{}{}{0pt}{}{n}{p}\right)y(u_{(p)}x)_{(np1)}z$$
$`(\mathrm{8.6.6})`$
and
$$B3=\underset{p=0}{\overset{n1}{}}\left(\genfrac{}{}{0pt}{}{n}{p}\right)\{\left((u_{(p)}y)_{(np1)}\right)z+x(u_{(p)}y)_{(np1)}z+$$
$$+\underset{q=0}{\overset{np2}{}}\left(\genfrac{}{}{0pt}{}{np1}{q}\right)\left((u_{(p)})_{(q)}x\right)_{(npq2)}z\}$$
$`(\mathrm{8.6.7})`$
Next,
$$C=\underset{j0}{}(1)^j\{\left(u_{(n)}^{(j+1)}(x_{(j)}y)\right)z+^{(j+1)}(x_{(j)}y)u_{(n)}z+$$
$$+\underset{p=0}{\overset{n1}{}}\left(\genfrac{}{}{0pt}{}{n}{p}\right)\left(u_{(p)}^{(j+1)}(x_{(j)}y)\right)_{(np1)}z\}$$
$`(\mathrm{8.6.8})`$
Due to (0.4.7),
$$C1=\underset{j0}{}(1)^j\left\{^{(j+1)}(u_{(n)}x_{(j)}y)+\underset{p=1}{\overset{min(n,j+1)}{}}\left(\genfrac{}{}{0pt}{}{n}{p}\right)^{(j+1p)}(u_{(np)}x_{(j)}y)\right\}z$$
$`(\mathrm{8.6.9})`$
Next,
$$C11=\underset{j0}{}(1)^j^{(j+1)}\left\{(u_{(n)}x)_{(j)}y+x_{(j)}u_{(n)}y+\underset{p=0}{\overset{n1}{}}\left(\genfrac{}{}{0pt}{}{n}{p}\right)(u_{(p)}x)_{(n+jp)}y\right\}z$$
$`(\mathrm{8.6.10})`$
and
$$C12=\underset{j0}{}(1)^j\underset{p=1}{\overset{min(n,j+1)}{}}\left(\genfrac{}{}{0pt}{}{n}{p}\right)^{(j+1p)}\{(u_{(np)}x)_{(j)}y+x_{(j)}u_{(np)}y+$$
$$+\underset{q=0}{\overset{np1}{}}\left(\genfrac{}{}{0pt}{}{np}{q}\right)(u_{(q)}x)_{(j+npq)}y\}z$$
$`(\mathrm{8.6.11})`$
We have $`A21+B21+C2=r(x,y)u_{(n)}zR`$. Next, $`A23=B32`$ and $`A32=B23`$. Next,
$$A21+B1=x(u_{(n)}y)z(u_{(n)}y)xz\stackrel{R}{}\underset{j}{}(1)^j^{(j+1)}(x_{(j)}u_{(n)}y)z=C112$$
Similarly,
$$A1+B21=(u_{(n)}x)yzy(u_{(n)}x)z\stackrel{R}{}\underset{j}{}(1)^j^{(j+1)}\left((u_{(n)}x)_{(j)}y\right)z=C111$$
8.7. Below we shall often use the identity
$$\left(\genfrac{}{}{0pt}{}{n}{p}\right)\left(\genfrac{}{}{0pt}{}{p}{r}\right)=\left(\genfrac{}{}{0pt}{}{n}{r}\right)\left(\genfrac{}{}{0pt}{}{nr}{np}\right)$$
$`(\mathrm{8.7.1})`$
We claim that
$$C113+C121+C123=A31$$
$`(\mathrm{8.7.2})`$
Indeed, we have
$$C113=\underset{j0}{}\underset{q=0}{\overset{n1}{}}(1)^j\left(\genfrac{}{}{0pt}{}{n}{q}\right)^{(j+1)}\left\{(u_{(q)}x)_{(n+jq)}\right\}z$$
$`(\mathrm{8.7.3})`$
and
$$C121+C123=\underset{j0}{}\underset{p=1}{\overset{min(n,j+1)}{}}$$
$$\underset{q=0}{\overset{np1}{}}(1)^j\left(\genfrac{}{}{0pt}{}{n}{p}\right)\left(\genfrac{}{}{0pt}{}{np}{q}\right)^{(j+1p)}\left\{(u_{(q)}x)_{(j+npq)}y\right\}z$$
$`(\mathrm{8.7.4})`$
Using (8.7.1), we have
$$\left(\genfrac{}{}{0pt}{}{n}{p}\right)\left(\genfrac{}{}{0pt}{}{np}{q}\right)=\left(\genfrac{}{}{0pt}{}{n}{np}\right)\left(\genfrac{}{}{0pt}{}{np}{q}\right)=\left(\genfrac{}{}{0pt}{}{n}{q}\right)\left(\genfrac{}{}{0pt}{}{nq}{p}\right)$$
Consider the part of (8.7.4) at a fixed $`q`$, $`0qn1`$:
$$(\mathrm{8.7.4})_q=\left(\genfrac{}{}{0pt}{}{n}{q}\right)\underset{j0}{}\underset{p=1}{\overset{min(nq,j+1)}{}}$$
$$(1)^j\left(\genfrac{}{}{0pt}{}{nq}{p}\right)^{(j+1p)}\left\{(u_{(q)}x)_{(j+npq)}y\right\}z=$$
$$=\left(\genfrac{}{}{0pt}{}{n}{q}\right)\underset{p=1}{\overset{nq}{}}\underset{jp1}{}(1)^j\left(\genfrac{}{}{0pt}{}{nq}{p}\right)^{(j+1p)}\left\{(u_{(q)}x)_{(j+npq)}y\right\}z$$
$`(\mathrm{8.7.5})`$
It is easy to see that the part of (8.7.5) corresponding to $`j=p1`$ is equal to $`A31_q`$ and the part corresponding to $`jp`$ equals $`C113_q`$. This proves (8.7.2).
Using the commutativity formula (0.4.3), one sees that $`B31=C122`$.
8.8. We claim that
$$A33+B33+C3=0$$
$`(\mathrm{8.8.1})`$
We have
$$C3=\underset{j0}{}(1)^j\underset{p=0}{\overset{n1}{}}\left(\genfrac{}{}{0pt}{}{n}{p}\right)\underset{q=0}{\overset{min(p,j+1)}{}}\left(\genfrac{}{}{0pt}{}{p}{q}\right)^{(j+1q)}(u_{(pq)}x_{(j)}y)_{(np1)}z=$$
$$=\underset{j0}{}(1)^j\underset{p=0}{\overset{n1}{}}\left(\genfrac{}{}{0pt}{}{n}{p}\right)\underset{q=0}{\overset{min(p,j+1)}{}}\left(\genfrac{}{}{0pt}{}{p}{q}\right)(1)^{j+1q}\left(\genfrac{}{}{0pt}{}{np1}{jq+1}\right)(u_{(pq)}x_{(j)}y)_{(njp+q2)}z=$$
(we set $`r:=pq`$)
$$=\underset{j0}{}\underset{p=0}{\overset{n1}{}}\left(\genfrac{}{}{0pt}{}{n}{p}\right)\underset{r=pj1}{\overset{p}{}}(1)^{pr}\left(\genfrac{}{}{0pt}{}{p}{r}\right)\left(\genfrac{}{}{0pt}{}{np1}{jp+r+1}\right)(u_{(r)}x_{(j)}y)_{(njr2)}z=$$
$$=\underset{j0}{}\underset{p=0}{\overset{n1}{}}\left(\genfrac{}{}{0pt}{}{n}{p}\right)\underset{r=pj1}{\overset{p}{}}(1)^{pr}\left(\genfrac{}{}{0pt}{}{p}{r}\right)\left(\genfrac{}{}{0pt}{}{np1}{jp+r+1}\right)$$
$$\left\{(x_{(j)}u_{(r)}y)_{(njr2)}z+\underset{s=0}{\overset{r}{}}\left(\genfrac{}{}{0pt}{}{r}{s}\right)\left((u_{(s)}x)_{(r+js)}y\right)_{(njr2)}z\right\}$$
$`(\mathrm{8.8.2})`$
So, we have written $`C3`$ as a sum of two terms. Now we claim that
$$C32=A33$$
$`(\mathrm{8.8.3})`$
Set $`l:=r+js`$. We have
$$C32=\underset{l0,s0,l+sn2}{}\left(\genfrac{}{}{0pt}{}{n}{s}\right)\underset{p=0}{\overset{l+s+1}{}}\underset{r=s}{\overset{min(p,l+s)}{}}$$
$$(1)^{pr}\left(\genfrac{}{}{0pt}{}{ns}{np}\right)\left(\genfrac{}{}{0pt}{}{ps}{rs}\right)\left(\genfrac{}{}{0pt}{}{np1}{l+sp+1}\right)\left((u_{(s)}x)_{(l)}y\right)_{(nls2)}z$$
where we have used that
$$\left(\genfrac{}{}{0pt}{}{n}{p}\right)\left(\genfrac{}{}{0pt}{}{p}{r}\right)\left(\genfrac{}{}{0pt}{}{r}{s}\right)=\left(\genfrac{}{}{0pt}{}{n}{r}\right)\left(\genfrac{}{}{0pt}{}{r}{s}\right)\left(\genfrac{}{}{0pt}{}{nr}{np}\right)=\left(\genfrac{}{}{0pt}{}{n}{s}\right)\left(\genfrac{}{}{0pt}{}{ns}{nr}\right)\left(\genfrac{}{}{0pt}{}{nr}{np}\right)=\left(\genfrac{}{}{0pt}{}{n}{s}\right)\left(\genfrac{}{}{0pt}{}{ns}{np}\right)\left(\genfrac{}{}{0pt}{}{ps}{rs}\right)$$
Thus, at fixed $`l,s`$,
$$C32_{l,s}=\left(\genfrac{}{}{0pt}{}{n}{s}\right)\underset{p=0}{\overset{l+s+1}{}}(1)^p\left(\genfrac{}{}{0pt}{}{ns}{np}\right)\left(\genfrac{}{}{0pt}{}{np1}{l+sp+1}\right)\underset{r=s}{\overset{min(p,l+s)}{}}(1)^r\left(\genfrac{}{}{0pt}{}{ps}{rs}\right)((u_{s)}x)_{(l)}y)_{(nls2)}z$$
The last sum is non-zero in two cases:
(a) $`p=s`$, so that $`r=s=p`$, and we get the coefficient $`\left(\genfrac{}{}{0pt}{}{n}{s}\right)\left(\genfrac{}{}{0pt}{}{ns1}{l+1}\right)`$,
and
(b) $`p=l+s+1`$ where we get $`\left(\genfrac{}{}{0pt}{}{n}{s}\right)\left(\genfrac{}{}{0pt}{}{ns}{l+1}\right)`$.
Since
$$\left(\genfrac{}{}{0pt}{}{ns1}{l+1}\right)\left(\genfrac{}{}{0pt}{}{ns}{l+1}\right)=\left(\genfrac{}{}{0pt}{}{ns1}{l}\right)$$
$`(\mathrm{8.8.4})`$
we get
$$C32_{l,s}=\left(\genfrac{}{}{0pt}{}{n}{s}\right)\left(\genfrac{}{}{0pt}{}{ns1}{l}\right)=A33_{l,s}$$
which proves (8.8.3).
Next, we claim that
$$C31=B33$$
$`(\mathrm{8.8.5})`$
Indeed,
$$C31=\underset{j0}{}\underset{p=0}{\overset{n1}{}}(1)^{pr}\left(\genfrac{}{}{0pt}{}{n}{p}\right)\underset{r=max(0,pj1)}{\overset{p}{}}\left(\genfrac{}{}{0pt}{}{np1}{jp+r+1}\right)\left(\genfrac{}{}{0pt}{}{p}{r}\right)(x_{(j)}u_{(r)}y)_{(njr2)}z=$$
(using (0.4.3))
$$=\underset{j0}{}\underset{p=0}{\overset{n1}{}}(1)^{pr}\left(\genfrac{}{}{0pt}{}{n}{p}\right)\underset{r=max(0,pj1)}{\overset{p}{}}\left(\genfrac{}{}{0pt}{}{np1}{jp+r+1}\right)\left(\genfrac{}{}{0pt}{}{p}{r}\right)$$
$$(1)^{j+1}\underset{s0}{}(1)^s^{(s)}\left((u_{(r)}y)_{j+s}x\right)_{(njr2)}z=$$
$$=\underset{j0}{}\underset{p=0}{\overset{n1}{}}(1)^{pr}\left(\genfrac{}{}{0pt}{}{n}{p}\right)\underset{r=max(0,pj1)}{\overset{p}{}}\left(\genfrac{}{}{0pt}{}{np1}{jp+r+1}\right)\left(\genfrac{}{}{0pt}{}{p}{r}\right)(1)^{j+1}$$
$$\underset{s0}{}\left(\genfrac{}{}{0pt}{}{njr2}{s}\right)\left((u_{(r)}y)_{(j+s)}x\right)_{(njsr2)}z$$
Set $`l:=j+s`$. We get
$$C31=\underset{l0,r0,l+rn2}{}\underset{s0,pr,p+sr+l+1}{}$$
$$(1)^{pr+ls+1}\left(\genfrac{}{}{0pt}{}{n}{p}\right)\left(\genfrac{}{}{0pt}{}{p}{r}\right)\left(\genfrac{}{}{0pt}{}{np1}{l+rsp+1}\right)\left(\genfrac{}{}{0pt}{}{nlr+s2}{s}\right)\left((u_{(r)}y)_{(l)}x\right)_{(nlr2)}z$$
Using the identities $`\left(\genfrac{}{}{0pt}{}{n}{p}\right)\left(\genfrac{}{}{0pt}{}{p}{r}\right)=\left(\genfrac{}{}{0pt}{}{nr}{np}\right)`$ and
$$\left(\genfrac{}{}{0pt}{}{np1}{l+rsp+1}\right)\left(\genfrac{}{}{0pt}{}{nlr+s2}{s}\right)=\left(\genfrac{}{}{0pt}{}{np1}{nl+sr2}\right)\left(\genfrac{}{}{0pt}{}{nlr+s2}{s}\right)=$$
$$=\left(\genfrac{}{}{0pt}{}{np1}{s}\right)\left(\genfrac{}{}{0pt}{}{nps1}{l+rsp1}\right)=\left(\genfrac{}{}{0pt}{}{np1}{nps1}\right)\left(\genfrac{}{}{0pt}{}{nps1}{nlr2}\right)=$$
$$=\left(\genfrac{}{}{0pt}{}{np1}{nlr2}\right)\left(\genfrac{}{}{0pt}{}{l+rp+1}{s}\right)$$
we get at fixed $`l,r`$
$$C31_{l,r}=(1)^{l+r+1}\left(\genfrac{}{}{0pt}{}{n}{r}\right)\underset{p=r}{\overset{r+l+1}{}}(1)^p\left(\genfrac{}{}{0pt}{}{nr}{np}\right)\left(\genfrac{}{}{0pt}{}{np1}{nlr2}\right)$$
$$\underset{s=0}{\overset{min(r+lp+1,l)}{}}(1)^s\left(\genfrac{}{}{0pt}{}{l+rp+1}{s}\right)\left((u_{(r)}y)_{(l)}x\right)_{(nlr2)}z$$
The last sum is non-zero in two cases:
(a) $`p=l+r+1`$, so that $`s=0`$, and we get the coefficient $`\left(\genfrac{}{}{0pt}{}{n}{r}\right)\left(\genfrac{}{}{0pt}{}{nr}{nlr1}\right)=\left(\genfrac{}{}{0pt}{}{n}{r}\right)\left(\genfrac{}{}{0pt}{}{nr}{l+1}\right)`$,
and
(b) $`p=r`$, so that we get the coefficient $`\left(\genfrac{}{}{0pt}{}{n}{r}\right)\left(\genfrac{}{}{0pt}{}{nr1}{nlr2}\right)=\left(\genfrac{}{}{0pt}{}{n}{r}\right)\left(\genfrac{}{}{0pt}{}{nr1}{l+1}\right)`$.
It follows that
$$C31_{r,l}=\left(\genfrac{}{}{0pt}{}{n}{r}\right)\left(\genfrac{}{}{0pt}{}{nr1}{l}\right)\left((u_{(r)}y)_{(l)}x\right)_{(nlr2)}z=B33_{r,l}$$
which proves (8.8.5). The identities (8.8.3) and (8.8.5) together imply (8.8.1).
After bookkeeping, the computations of 8.6 — 8.8 show that the left hand side of (8.6.1) belongs to the ideal $`R`$, which finishes the proof of Lemma 8.5. $``$
8.9. By Lemma 8.5, the operators $`x_{(j)}`$ induce operators
$$x_{(j)}:UCUC,xC,j$$
$`(\mathrm{8.9.1})`$
8.10. Lemma. The operators (8.9.1) satisfy the OPE formula (0.5.12).
This lemma will be proven in 8.11 — 8.19 below.
8.11. We have to prove the identity between the operators acting on $`UC`$,
$$[x_{(m)},y_{(n)}]=\underset{j0}{}\left(\genfrac{}{}{0pt}{}{m}{j}\right)(x_{(j)}y)_{(m+nj)}$$
$`(\mathrm{8.11.1})`$
for all $`m,n`$.
Let us discuss separately three cases.
8.12. Case A. $`m<0`$ and $`n<0`$. First of all, note that for $`m=n=1`$ the relation (8.11.1) is nothing but (8.2.1) and holds true by definition, since $`R`$ is a left ideal.
In general, set $`m=1a,n=1b`$ for $`a,b0`$. We have by definition (8.4.1)
$$[x_{(1a)},y_{(1b)}]=[^{(a)}x,^{(b)}y]=$$
$$=\underset{j0}{}(1)^j^{(j+1)}((^{(a)}x)_{(j)}^{(b)}y)$$
$`(\mathrm{8.12.1})`$
The right hand side of (8.11.1) is equal to
$$\underset{j0}{}\left(\genfrac{}{}{0pt}{}{1a}{j}\right)(x_{(j)}y)_{(2abj)}=\underset{j0}{}(1)^j\left(\genfrac{}{}{0pt}{}{a+j}{j}\right)^{(a+b+j+1)}(x_{(j)}y)$$
$`(\mathrm{8.12.2})`$
where we have used that
$$\left(\genfrac{}{}{0pt}{}{1a}{j}\right)=(1)^j\left(\genfrac{}{}{0pt}{}{a+j}{j}\right)$$
$`(\mathrm{8.12.3})`$
for $`a,j0`$.
On the other hand,
$$(^{(a)}x)_{(j)}^{(b)}y=(1)^a\left(\genfrac{}{}{0pt}{}{j}{a}\right)x_{(ja)}^{(b)}y=(1)^a\left(\genfrac{}{}{0pt}{}{j}{a}\right)\underset{p=0}{\overset{b}{}}\left(\genfrac{}{}{0pt}{}{ja}{p}\right)^{(bp)}(x_{(jap)}y),$$
by (0.4.2) and (0.4.7). Therefore, the rhs of (8.12.1) is equal to
$$\underset{j0}{}(1)^{j+a}\left(\genfrac{}{}{0pt}{}{j}{a}\right)\underset{p0}{}\left(\genfrac{}{}{0pt}{}{ja}{p}\right)^{(j+1)}^{(bp)}(x_{(jap)}y)=$$
$$=\underset{j0}{}(1)^{j+a}\left(\genfrac{}{}{0pt}{}{j}{a}\right)\underset{p=0}{\overset{min(b,ja)}{}}\left(\genfrac{}{}{0pt}{}{ja}{p}\right)\left(\genfrac{}{}{0pt}{}{j+1+bp}{j+1}\right)^{(j+1+bp)}(x_{(jap)}y)=$$
(substituting $`k:=jap`$)
$$=\underset{k0}{}\underset{p=0}{\overset{b}{}}(1)^{p+k}\left(\genfrac{}{}{0pt}{}{k+a+p}{a}\right)\left(\genfrac{}{}{0pt}{}{k+p}{p}\right)\left(\genfrac{}{}{0pt}{}{k+a+b+1}{k+a+p+1}\right)^{(k+a+b+1)}(x_{(k)}y)=$$
(using that $`\left(\genfrac{}{}{0pt}{}{k+a+p}{a}\right)\left(\genfrac{}{}{0pt}{}{k+p}{p}\right)=\left(\genfrac{}{}{0pt}{}{k+a+p}{k+p}\right)\left(\genfrac{}{}{0pt}{}{k+p}{p}\right)=\left(\genfrac{}{}{0pt}{}{k+a+p}{p}\right)\left(\genfrac{}{}{0pt}{}{k+a}{a}\right)`$)
$$=\underset{k0}{}(1)^k\left(\genfrac{}{}{0pt}{}{k+a}{a}\right)\left\{\underset{p=0}{\overset{b}{}}(1)^p\left(\genfrac{}{}{0pt}{}{k+a+p}{p}\right)\left(\genfrac{}{}{0pt}{}{k+a+b+1}{k+a+p+1}\right)\right\}^{(k+a+b+1)}(x_{(k)}y)$$
$`(\mathrm{8.12.4})`$
8.13. Lemma. For all $`b,q_0`$,
$$\underset{p=0}{\overset{b}{}}(1)^p\left(\genfrac{}{}{0pt}{}{p+q}{p}\right)\left(\genfrac{}{}{0pt}{}{b+q+1}{p+q+1}\right)=1$$
This is easily proved by induction on $`b`$. $``$
It follows from this lemma that (8.12.4) is equal to (8.12.2), which completes the check of Case A.
8.14. Case B. $`m0`$ and $`n<0`$. If $`n=1`$ then (8.11.1) is the same as the definition (8.4.2). Now let $`n=1a,a0`$. The lhs of (8.11.1) is equal to
$$[x_{(m)},y_{(1a)}]=[x_{(m)},^{(a)}y]=\underset{j=0}{\overset{m1}{}}\left(\genfrac{}{}{0pt}{}{m}{j}\right)(x_{(j)}^{(a)}y)_{(mj1)}+x_{(m)}^{(a)}y$$
$`(\mathrm{8.14.1})`$
The rhs of (8.11.1) is equal to
$$\underset{k=0}{\overset{m}{}}\left(\genfrac{}{}{0pt}{}{m}{k}\right)(x_{(k)}y)_{(m1ak)}=$$
$$=\underset{k=0}{\overset{m1a}{}}\left(\genfrac{}{}{0pt}{}{m}{k}\right)(x_{(k)}y)_{(m1ak)}+\underset{kmax(ma,0)}{}\left(\genfrac{}{}{0pt}{}{m}{k}\right)(x_{(k)}y)_{(m1ak)}$$
$`(\mathrm{8.14.2})`$
Consider the first sum in the rhs of (8.14.1). It is equal to
$$\underset{j=0}{\overset{m1}{}}\left(\genfrac{}{}{0pt}{}{m}{j}\right)\underset{p=o}{\overset{min(j,a)}{}}\left(\genfrac{}{}{0pt}{}{j}{p}\right)^{(ap)}(x_{(jp)}y)_{(mj1)}=$$
$$=\underset{j=0}{\overset{m1}{}}\left(\genfrac{}{}{0pt}{}{m}{j}\right)\underset{p=max(0,am+j+1)}{\overset{min(j,a)}{}}\left(\genfrac{}{}{0pt}{}{j}{p}\right)(1)^{ap}\left(\genfrac{}{}{0pt}{}{mj1}{ap}\right)(x_{(jp)}y)_{(mja+p1)}=$$
(making the substitution $`k=jp`$)
$$=\underset{k=0}{\overset{ma1}{}}\underset{p=0}{\overset{a}{}}(1)^{ap}\left(\genfrac{}{}{0pt}{}{m}{k+p}\right)\left(\genfrac{}{}{0pt}{}{k+p}{p}\right)\left(\genfrac{}{}{0pt}{}{mpk1}{ap}\right)(x_{(k)}y)_{(mka1)}=$$
(using that $`\left(\genfrac{}{}{0pt}{}{m}{k+p}\right)\left(\genfrac{}{}{0pt}{}{k+p}{p}\right)=\left(\genfrac{}{}{0pt}{}{m}{k}\right)\left(\genfrac{}{}{0pt}{}{mk}{mpk}\right)`$)
$$=\underset{k=0}{\overset{ma1}{}}\left(\genfrac{}{}{0pt}{}{m}{k}\right)\left\{\underset{p=0}{\overset{a}{}}(1)^{ap}\left(\genfrac{}{}{0pt}{}{mk}{mpk}\right)\left(\genfrac{}{}{0pt}{}{mpk1}{ap}\right)\right\}(x_{(k)}y)_{(mka1)}=$$
(substituting $`r=pa,s=mak10`$)
$$=\underset{k=0}{\overset{ma1}{}}\left(\genfrac{}{}{0pt}{}{m}{k}\right)\left\{\underset{r=0}{\overset{a}{}}(1)^r\left(\genfrac{}{}{0pt}{}{a+s+1}{r+s+1}\right)\left(\genfrac{}{}{0pt}{}{r+s}{r}\right)\right\}(x_{(k)}y)_{(mka1)}$$
By Lemma 8.13, this is equal to the first sum in the rhs of (8.14.2).
The second term in the rhs of (8.14.1) is equal to
$$x_{(m)}^{(a)}y=\underset{p=0}{\overset{min(m,a)}{}}\left(\genfrac{}{}{0pt}{}{m}{p}\right)^{(ap)}(x_{(mp)}y)=$$
(substituting $`k=mp`$)
$$=\underset{kmax(ma,0)}{}\left(\genfrac{}{}{0pt}{}{m}{mk}\right)^{(am+k)}(x_{(k)}y)=\underset{kmax(ma,0)}{}\left(\genfrac{}{}{0pt}{}{m}{k}\right)(x_{(k)}y)_{(1a+mk)},$$
which is the same as the second sum in (8.14.2).
Therefore, (8.14.1)$`=`$(8.14.2), which finishes the proof of (8.11.1) in the Case B.
8.15. Case C. $`m<0`$ and $`n0`$. First let us treat the case $`m=1`$. We need to prove that
$$[x_{(1)},y_{(n)}]=\underset{j0}{}(1)^j(x_{(j)}y)_{(n1j)}$$
$`(\mathrm{8.15.1})`$
The lhs of (8.15.1) is equal to
$$[y_{(n)},x_{(1)}]=\underset{p0}{}\left(\genfrac{}{}{0pt}{}{n}{p}\right)(y_{(p)}x)_{(np1)}$$
$`(\mathrm{8.15.2})`$
On the other hand, the rhs of (8.15.1) equals (we use (0.4.3))
$$\underset{j0}{}\underset{q0}{}(1)^q^{(q)}(y_{(j+q)}x)_{(n1j)}\stackrel{\left(\mathrm{8.4.3}\right)}{=}\underset{j0,q0}{}\left(\genfrac{}{}{0pt}{}{n1j}{q}\right)(y_{(j+q)}x)_{(n1jq)}=$$
(substituting $`p=j+q`$)
$$=\underset{p0}{}(y_{(p)}x)_{(n1p)}\left\{\underset{q=0}{\overset{p}{}}\left(\genfrac{}{}{0pt}{}{n1p+q}{q}\right)\right\}$$
Therefore, (8.15.1) follows from the identity below, which is easily checked by induction on $`n`$:
$$\underset{q=0}{\overset{p}{}}\left(\genfrac{}{}{0pt}{}{n1p+q}{q}\right)=\left(\genfrac{}{}{0pt}{}{n}{p}\right)$$
$`(\mathrm{8.15.3})`$
for all $`n,p0`$.
Now assume that $`m=1a,a0`$. We need to prove that
$$[x_{(1a)},y_{(n)}]=\underset{p0}{}\left(\genfrac{}{}{0pt}{}{1a}{p}\right)(x_{(p)}y)_{(1a+np)}$$
$`(\mathrm{8.15.4})`$
The lhs is equal to
$$[^{(a)}x_{(1)},y_{(n)}]=\underset{j0}{}(1)^j(^{(a)}x_{(j)}y)_{(n1j)}=\underset{ja}{}(1)^{j+a}\left(\genfrac{}{}{0pt}{}{j}{a}\right)(x_{(ja)}y)_{(n1j)}=$$
(substituting $`p=ja`$)
$$=\underset{p0}{}(1)^p\left(\genfrac{}{}{0pt}{}{p+a}{a}\right)(x_{(p)}y)_{(n1pa)}$$
so (8.15.4) follows from (0.0.2). The completes the proof of Case C.
8.16. Case D. $`m,n0`$. We have to prove that
$$[x_{(m)},y_{(n)}]v=\underset{j=0}{\overset{m}{}}\left(\genfrac{}{}{0pt}{}{m}{j}\right)(x_{(j)}y)_{(m+nj)}v$$
$`(\mathrm{8.16.1})`$
for all $`x,yC,vUC`$. In fact, we shall prove a stronger statement, which we prefer to formulate as a separate
8.17. Lemma. If $`m`$ and $`n`$ are nonnegative then the OPE identity (8.11.1) holds true on $`TC`$. In other words, (8.16.1) holds true for all $`m,n0,x,yC`$ and $`vTC`$.
We shall prove this by induction on the length of $`v`$. If $`vC`$ then the desired identity holds true by definition of a conformal algebra, (0.4.4). Now, let $`v=zu`$ where $`zC,uTC`$. So, we have to prove that
$$[x_{(m)},y_{(n)}]zu=\underset{j=0}{\overset{m}{}}\left(\genfrac{}{}{0pt}{}{m}{j}\right)(x_{(j)}y)_{(m+nj)}zu$$
$`(\mathrm{8.17.1})`$
Consider the lhs first. We have
$$A:=x_{(m)}y_{(n)}zu=x_{(m)}\left\{zy_{(n)}u+(y_{(n)}z)u+\underset{p=0}{\overset{n1}{}}\left(\genfrac{}{}{0pt}{}{n}{p}\right)(y_{(p)}z)_{(np1)}u\right\}=$$
$$=zx_{(m)}y_{(n)}u+(x_{(m)}z)y_{(n)}u+\underset{q=0}{\overset{m1}{}}\left(\genfrac{}{}{0pt}{}{m}{q}\right)(x_{(q)}z)_{(mq1)}y_{(n)}u+$$
$$+(y_{(n)}z)x_{(m)}u+(x_{(m)}y_{(n)}z)u+\underset{q=0}{\overset{m1}{}}\left(\genfrac{}{}{0pt}{}{m}{q}\right)(x_{(q)}y_{(n)}z)_{(mq1)}u+\underset{p=0}{\overset{n1}{}}\left(\genfrac{}{}{0pt}{}{n}{p}\right)x_{(m)}(y_{(p)}z)_{(np1)}u$$
Similarly,
$$B:=y_{(n)}x_{(m)}zu=zy_{(n)}x_{(m)}u(y_{(n)}z)x_{(m)}u\underset{p=0}{\overset{n1}{}}\left(\genfrac{}{}{0pt}{}{n}{p}\right)(y_{(p)}z)_{(np1)}x_{(m)}u$$
$$(x_{(m)}z)y_{(n)}u(y_{(n)}x_{(m)}z)u\underset{p=0}{\overset{n1}{}}\left(\genfrac{}{}{0pt}{}{n}{p}\right)(y_{(p)}x_{(m)}z)_{(np1)}u\underset{q=0}{\overset{m1}{}}\left(\genfrac{}{}{0pt}{}{m}{q}\right)y_{(n)}(x_{(q)}z)_{(mq1)}u$$
Obviously, the lhs of (8.17.1) is equal to $`A+B`$, and we should compare this with the rhs:
$$R:=\underset{j=0}{\overset{m}{}}\left(\genfrac{}{}{0pt}{}{m}{j}\right)(x_{(j)}y)_{(m+nj)}zu=\underset{j=0}{\overset{m}{}}\left(\genfrac{}{}{0pt}{}{m}{j}\right)\{z(x_{(j)}y)_{(m+nj)}u+$$
$$+((x_{(j)}y)_{(m+nj)}z)u+\underset{r=0}{\overset{m+nj1}{}}\left(\genfrac{}{}{0pt}{}{m+nj}{r}\right)((x_{(j)}y)_{(r)}z)_{(m+njr1)}u\}$$
First of all, $`A1+B1=R1`$ by induction hypothesis. Next, $`A2=B4`$ and $`A4=B2`$. Further, $`A5+B5=R2`$ by the axiom of a conformal algebra. Next,
$$C:=A3+B7=\underset{q=0}{\overset{m1}{}}\left(\genfrac{}{}{0pt}{}{m}{q}\right)\left\{(x_{(q)}z)_{(mq1)}y_{(n)}uy_{(n)}(x_{(q)}z)_{(mq1)}u\right\}=$$
(by induction hypothesis)
$$=\underset{q=0}{\overset{m1}{}}\left(\genfrac{}{}{0pt}{}{m}{q}\right)\underset{s=0}{\overset{n}{}}\left(\genfrac{}{}{0pt}{}{n}{s}\right)(y_{(s)}x_{(q)}z)_{(m+nqs1)}u$$
Similarly,
$$D:=A7+B3=\underset{p=0}{\overset{n1}{}}\left(\genfrac{}{}{0pt}{}{n}{p}\right)\underset{r=0}{\overset{m}{}}\left(\genfrac{}{}{0pt}{}{m}{r}\right)(x_{(r)}y_{(p)}z)_{(m+npr1)}u$$
We have
$$E:=A6+C_{s=n}=\underset{q=0}{\overset{m1}{}}\left(\genfrac{}{}{0pt}{}{m}{q}\right)\left\{(x_{(q)}y_{(n)}z)_{(mq1)}(y_{(n)}x_{(q)}z)_{(mq1)}\right\}u=$$
$$=\underset{q=0}{\overset{m1}{}}\underset{s=0}{\overset{q}{}}\left(\genfrac{}{}{0pt}{}{m}{q}\right)\left(\genfrac{}{}{0pt}{}{q}{s}\right)((x_{(s)}y)_{(q+ns)}z)_{(mq1)}u=$$
$$=\underset{q=0}{\overset{m1}{}}\underset{s=0}{\overset{q}{}}\left(\genfrac{}{}{0pt}{}{m}{s}\right)\left(\genfrac{}{}{0pt}{}{ms}{mq}\right)((x_{(s)}y)_{(q+ns)}z)_{(mq1)}u$$
Similarly,
$$F:=D_{r=m}+B6=\underset{p=0}{\overset{n1}{}}\left(\genfrac{}{}{0pt}{}{n}{p}\right)\left\{(x_{(m)}y_{(p)}z)_{(np1)}(y_{(p)}x_{(m)}z)_{(np1)}\right\}u=$$
$$=\underset{p=0}{\overset{n1}{}}\underset{r=0}{\overset{m}{}}\left(\genfrac{}{}{0pt}{}{n}{p}\right)\left(\genfrac{}{}{0pt}{}{m}{r}\right)((x_{(r)}y)_{(m+pr)}z)_{(np1)}u$$
Finally,
$$G:=C_{s<n}+D_{r<m}=\underset{p=0}{\overset{n1}{}}\underset{q=0}{\overset{m1}{}}\left(\genfrac{}{}{0pt}{}{n}{p}\right)\left(\genfrac{}{}{0pt}{}{m}{q}\right)\left\{(x_{(q)}y_{(p)}z)_{(m+npq1)}(y_{(p)}x_{(q)}z)_{(m+npq1)}\right\}u=$$
$$=\underset{p=0}{\overset{n1}{}}\underset{q=0}{\overset{m1}{}}\left(\genfrac{}{}{0pt}{}{n}{p}\right)\left(\genfrac{}{}{0pt}{}{m}{q}\right)\underset{s=0}{\overset{q}{}}\left(\genfrac{}{}{0pt}{}{q}{s}\right)\left((x_{(s)}y)_{(p+qs)}z\right)_{(m+npq1)}u=$$
$$=\underset{p=0}{\overset{n1}{}}\underset{q=0}{\overset{m1}{}}\left(\genfrac{}{}{0pt}{}{n}{p}\right)\underset{s=0}{\overset{q}{}}\left(\genfrac{}{}{0pt}{}{m}{s}\right)\left(\genfrac{}{}{0pt}{}{ms}{mq}\right)\left((x_{(s)}y)_{(p+qs)}z\right)_{(m+npq1)}u$$
A careful examination of the sums involved, together with a formula
$$\underset{p}{}\left(\genfrac{}{}{0pt}{}{a}{rp}\right)\left(\genfrac{}{}{0pt}{}{n}{p}\right)=\left(\genfrac{}{}{0pt}{}{a+n}{r}\right)$$
$`(\mathrm{8.17.2})`$
with $`a=mj`$, which is obvious from the combinatorial definition of the binomial coefficients, shows that $`E+F+G=R3`$.
This completes the proof of Lemma. $``$
8.18. Remark. Lemmas 8.17 and 8.4.1 mean that $`TC`$ is canonically a module over conformal algebra $`C`$.
8.19. Lemma 8.17 implies Case D, which completes the proof of Lemma 8.10. $``$
8.20. Now we can finish the proof of Theorem 8.3. According to Lemmas 8.5 and 8.10, we have the collection of mutually local fields $`a(z)=_na_{(n)}z^{n1},aC`$, acting on the space $`UC`$. These fields satisfy the conditions of Theorem 0.7, which gives the desired structure of a vertex algebra on $`UC`$.
The other claims of Theorem 8.3 are obvious. $``$
8.21. Let $`C`$ be a conformal algebra. Let us introduce an operation $`[x,y]`$ on the space $`C`$ by
$$[x,y]=\underset{j0}{}(1)^j^{(j+1)}(x_{(j)}y)$$
$`(\mathrm{8.21.1})`$
8.22. Theorem. The operation (8.21.1) is a Lie bracket on $`C`$.
The space $`C`$ with the Lie algebra structure given by (8.21.1) will be denoted $`C^{Lie}`$.
The proof is given in 8.23 — 8.26 below.
8.23. Skew symmetry. Using (0.4.3), we have
$$[x,y]=\underset{j0}{}(1)^j^{(j+1)}\left\{(1)^{j+1}\underset{p0}{}(1)^p^{(p)}(y_{(j+p)}x)\right\}=$$
$$=\underset{j,p0}{}(1)^p\left(\genfrac{}{}{0pt}{}{j+1+p}{p}\right)^{(j+p+1)}(y_{(j+p)}x)=$$
$$=\underset{n0}{}\left\{\underset{p=0}{\overset{n}{}}(1)^p\left(\genfrac{}{}{0pt}{}{n+1}{p}\right)\right\}^{(n+1)}(y_{(n)}x)=[y,x],QED$$
8.24. Jacobi identity. We need to prove that
$$[x,[y,z]]=[[x,y],z]+[y,[x,z]]$$
This is an obvious consequence of two lemmas below.
8.25. Lemma. For all $`x,y,zC`$
$$[x,[y,z]]=\underset{a,b0}{}(1)^{a+b}^{(a+b+2)}(x_{(a)}y_{(b)}z)$$
$`(\mathrm{8.25.1})`$
Indeed, we have
$$[x,[y,z]]=\underset{j0}{}(1)^j[x,^{(j+1)}(y_{(j)}z)]=\underset{i,j0}{}(1)^{i+j}^{(i+1)}\left(x_{(i)}^{(j+1)}(y_{(j)}z)\right)=$$
(using (0.4.6))
$$=\underset{i,j0}{}(1)^{i+j}^{(i+1)}\left\{\underset{p0}{}\left(\genfrac{}{}{0pt}{}{i}{p}\right)^{(j+1p)}(x_{(ip)}y_{(j)}z)\right\}=$$
$$=\underset{i,j0}{}(1)^{i+j}\underset{p0}{}\left(\genfrac{}{}{0pt}{}{i+j+2p}{i+1}\right)\left(\genfrac{}{}{0pt}{}{i}{p}\right)^{(i+j+2p)}(x_{(ip)}y_{(j)}z)=$$
(substituting $`a:=ip,b:=j`$)
$$=\underset{a,b0}{}(1)^{a+b}^{(a+b+2)}(x_{(a)}y_{(b)}z)\left\{\underset{p0}{}(1)^p\left(\genfrac{}{}{0pt}{}{a+p}{p}\right)\left(\genfrac{}{}{0pt}{}{a+b+2}{a+p+1}\right)\right\},$$
and we conclude by Lemma 8.13.
8.26. Lemma. For all $`x,y,zC`$
$$[[x,y],z]=\underset{a,b0}{}(1)^{a+b}^{(a+b+2)}\left\{x_{(a)}y_{(b)}zy_{(b)}x_{(a)}z\right\}$$
$`(\mathrm{8.26.1})`$
Indeed, we have
$$[[x,y],z]=\underset{j0}{}(1)^j\underset{i0}{}(1)^i^{(i+1)}\left\{^{(j+1)}(x_{(j)}y)_{(i)}z\right\}=$$
$$=\underset{i,j0}{}(1)^{i+1}\left(\genfrac{}{}{0pt}{}{i}{j+1}\right)^{(i+1)}((x_{(j)}y)_{(ij1)}z)=$$
(by (0.4.5))
$$=\underset{i,j0}{}(1)^{i+1}\left(\genfrac{}{}{0pt}{}{i}{j+1}\right)\underset{q=0}{\overset{j}{}}(1)^q\left(\genfrac{}{}{0pt}{}{j}{q}\right)^{(i+1)}\left\{x_{(jq)}y_{(ij+q1)}z(1)^jy_{(iq1)}x_{(q)}z\right\},$$
and we again conclude after an easy rearrangement of the indices and using Lemma 8.13.
This completes the proof of lemma and of the theorem. $``$
8.27. Corollary. The vertex enevlope $`UC`$ of a conformal algebra $`C`$ is equal to the associative enveloping algebra $`UC^{Lie}`$ of the Lie algebra $`C^{Lie}`$.
If $`i:CUC`$ denotes the canonical morphism,
$$i([x,y])=i(x)_{(1)}i(y)i(y)_{(1)}i(x)$$
$`(\mathrm{8.27.1})`$
for all $`x,yC`$.
This follows immediately from the definition in 8.2.
Note that in an arbitrary vertex algebra, although the operation $`x_{(1)}y`$ is not associative, the commutator $`x_{(1)}yy_{(1)}x`$ is a Lie bracket, according to a remark by A.Radul cited in \[K\], 3.1, page 82.
This corollary supports the point of view advocated by Beilinson and Drinfeld, namely, that conformal algebras ($`Lie^{}`$-algebras in the language of BD) are analogs of Lie algebras, while vertex algebras (chiral algebras) are analogs of associative algebras.
8.28. Let us define a canonical increasing filtration $`F_{}UC`$ on $`UC`$ by induction: $`F_1UC=0,F_0UC=k\text{1},F_{i+1}UC=F_iUC+C_{(1)}F_iUC`$.
Obviously, the operation $`x_{(1)}y`$ induces on the associated graded space $`gr_F(UC):=_{i0}gr_F^i(UC):=_iF_iUC/F_{i1}UC`$ a structure of a commutative associative unitary $`k`$-algebra.
We have an evident surjective morphism of commutative algebras
$$Sym_k(C)gr_F(UC)$$
$`(\mathrm{8.28.1})`$
8.29. Theorem. If $`C`$ is a free $`k`$-module then (8.28.1) is an isomorphism.
This is immediate consequence of Corollary 8.27 and the usual Poincaré-Birkhoff-Witt theorem for Lie algebras, cf. \[Se\], Part I, Chapter III, Theorem 4.3.
§9. Enveloping Algebra of a Vertex Algebroid
Abelian case
9.1. Let us define a $``$-module (over $`k`$) to be a $`_0`$-graded $`k`$-module $`C=_{i0}C_i`$ equipped with a family of endomorphisms $`^{(j)}:CC,j0,`$ where $`^{(j)}`$ has degree $`j`$, $`^{(0)}=Id`$ and
$$^{(i)}^{(j)}=\left(\genfrac{}{}{0pt}{}{i+j}{i}\right)^{(i+j)}$$
Obviously, a $``$-module is the same as an abelian conformal algebra, i.e. a conformal algebra in which all operations <sub>(n)</sub> are trivial. These objects, with obvious morphisms, form a category $`od`$.
Let us define a $`1`$-truncated $``$-module to be a triple $`𝒳=(C_0,C_1,)`$ where $`C_i`$ are $`k`$-modules and $`:C_0C_1`$ is a map of $`k`$-modules. These objects form a category $`od_1`$. We have an obvious truncation functor
$$odod_1$$
$`(\mathrm{9.1.1})`$
It is easy to construct the left adjoint to (9.1.1). Namely, given $`𝒳=(C_0,C_1,)od_1`$, consider the direct sum
$$M=C_0C_1(_{i2}^{(i)}C_0)(_{i1}^{(i)}C_1)$$
$`(\mathrm{9.1.2})`$
where $`^{(i)}X`$ denotes a copy of a $`k`$-module $`X`$, whose elements are denoted $`^{(i)}x,xX`$. Define a $``$-module $`𝒳^{}=_{i0}C_i`$ as the quotient of $`M`$ modulo the submodule generated by the following elements:
$$^{(i+1)}x(i+1)^{(i)}(x),i1,xC_0$$
$`(\mathrm{9.1.3})`$
The grading and operators $`^{(j)}:C_iC_{i+j}`$ are defined in the obvious manner.
Note that if the ground ring $`k`$ contains the field of rational numbers $``$ then $`𝒳`$ is equal simply to
$$𝒳^{}=C_0C_1(_{i1}^{(i)}C_1)$$
$`(\mathrm{9.1.4})`$
9.2. Recall that in 0.9 one introduced the category $`𝒜lg`$ of $``$-algebras which are identified with commutative vertex algebras. A $``$-ideal $`J`$ in a $``$-algebra $`B`$ is a $`_0`$-graded ideal stable under all endomorphisms $`^{(i)}`$; this is the same as a vertex ideal if $`B`$ is understood as a vertex algebra.
We have an obvious forgetful functor
$$𝒜lgod$$
$`(\mathrm{9.2.1})`$
This functor admits a left adjoint
$$Sym:od𝒜lg$$
$`(\mathrm{9.2.2})`$
that assigns to a $``$-module $`C`$ the symmetric algebra (over $`k`$) $`Sym(C)`$. It is nothing but the restriction of the functor defined in 8.3 to the full subcategory of abelian conformal algebras.
9.3. Jet algebra. Let us call a vertex algebroid $`𝒜=(A,T,\mathrm{\Omega },\mathrm{})`$ abelian if $`T=0`$. Thus, an abelian vertex algebroid is simply a triple $`𝒜=(A,\mathrm{\Omega },)`$ where $`A`$ is a commutative algebra, $`\mathrm{\Omega }`$ is an $`A`$-module and $`:A\mathrm{\Omega }`$ is a derivation. Thus, abelian vertex algebroids are the same as ”1-truncated $``$-algebras”. The category of abelian vertex algebroids will be denoted $`𝒜lg𝒜b`$; it is a full subcategory of $`𝒜lg`$.
We have an obvious forgetful functor
$$o:𝒜lg𝒜bod$$
$`(\mathrm{9.3.1})`$
and the truncation functor
$$𝒜lg𝒜lg𝒜b$$
$`(\mathrm{9.3.2})`$
Let us construct a left adjoint to (9.3.2). Given $`𝒜=(A,\mathrm{\Omega },)𝒜lg𝒜b`$, consider the $``$-module $`(o𝒜)^{}`$ and its symmetric algebra $`S(𝒜):=Sym((o𝒜)^{})𝒜lg`$. Note that the grading on $`S(𝒜)`$ is defined in such a way that $`AS(𝒜)`$ has grading $`0`$ and $`\mathrm{\Omega }S(𝒜)`$ has grading $`1`$.
Let $`RS(𝒜)`$ be a $``$-ideal generated by all elements
$$1_A1_{S(𝒜)};abab;a\omega a\omega ,$$
$`(\mathrm{9.3.3})`$
$`a,bA,\omega \mathrm{\Omega }`$ (note that these elements are indeed homogeneous). We denote the quotient $``$-algebra $`S(𝒜)/R`$ by $`J(𝒜)`$, and call it the jet algebra of $`𝒜`$. The assignment $`𝒜J(𝒜)`$ defines a functor
$$J:𝒜lg𝒜b𝒜lg$$
$`(\mathrm{9.3.4})`$
left adjoint to (9.3.2).
The composition $`A(o𝒜)^{}=Sym^1((o𝒜)^{})S(𝒜)J(𝒜)`$ defines a map
$$i_A:AJ(𝒜)$$
$`(\mathrm{9.3.5})`$
which is a map of $`k`$-algebras, due to the relations (9.3.3), which identifies $`A`$ with $`J(A)_0`$. This map makes $`J(𝒜)`$ an $`A`$-algebra.
A similar composition defines a map
$$i_\mathrm{\Omega }:\mathrm{\Omega }J(𝒜)$$
$`(\mathrm{9.3.6})`$
which is a map of left $`A`$-modules, again due to the relations (9.3.3). We have the compatibility: $`(i_A(a))=i_\mathrm{\Omega }((a))`$, and the triple $`(J(𝒜),i_A,i_\mathrm{\Omega })`$ has a
9.3.1. Universal Property. Given a triple $`(B,i_A^{},i_\mathrm{\Omega }^{})`$ where $`B`$ is a $``$-algebra, $`i_A^{}:AB`$ a morphism of $`k`$-algebras such that $`i_A^{}(A)B_0`$, $`i_\mathrm{\Omega }^{}:\mathrm{\Omega }B`$ a morphism of $`A`$-modules such that $`i_\mathrm{\Omega }^{}(\mathrm{\Omega })B_1`$ and $`i_\mathrm{\Omega }^{}((a))=(i_A^{}(a))`$, there exists a unique map of $``$-algebras $`f:J(A)B`$ such that $`fi_A=i_A^{}`$ and $`fi_\mathrm{\Omega }=i_\mathrm{\Omega }^{}`$.
9.4. Let $`J^+=_{i1}J(𝒜)_i`$ be the augmentation ideal. Consider the associated graded algebra with respect to the $`J^+`$-adic filtration
$$grJ(𝒜)=_{i0}J^{+i}/J^{+i+1}$$
$`(\mathrm{9.4.1})`$
It inherits a $`_0`$-grading from $`J(𝒜)`$. Let $`\mathrm{\Omega }^{(i)}`$ denote the image of the composition
$$^{(i)}A^{(i1)}\mathrm{\Omega }J^+J^+/J^{+2}$$
Note that $`\mathrm{\Omega }^{(i)}`$ is an $`A`$-submodule of $`J^+/J^{+2}`$, and is contained inside the homogeneous component $`(J^+/J^{+2})_i`$.
Adding up, we get a map of $`A`$-modules
$$_{i1}\mathrm{\Omega }^{(i)}J^+/J^{+2}$$
$`(\mathrm{9.4.2})`$
and hence a morphism of $`A`$-algebras
$$Sym_A(_{i1}\mathrm{\Omega }^{(i)})Sym_A(J^+/J^{+2})grJ(𝒜)$$
$`(\mathrm{9.4.3})`$
the second arrow being the evident canonical map.
9.5. Theorem. Assume that $`k`$ and $`\mathrm{\Omega }`$ is a free $`A`$-module. Then
(i) the maps $`\mathrm{\Omega }\mathrm{\Omega }^{(i)}`$ sending $`\omega `$ to (the image of) $`^{(i1)}\omega ,i1,`$ are isomorphisms of $`A`$-modules;
(ii) both maps in (9.4.3) are isomorphisms of $`A`$-algebras.
Indeed, choose an $`A`$-base $`\{\omega _r\}`$ of $`\mathrm{\Omega }`$. Let $`\stackrel{~}{\mathrm{\Omega }}^{(i)}`$ denote a free $`A`$-module with the base $`\{^{(i1)}\omega _r\}`$. Let $`\stackrel{~}{J}`$ be the symmetric algebra $`Sym_A\{A(_{i1}\stackrel{~}{\mathrm{\Omega }}^{(i)})\}`$. It has an obvious structure of a $``$-algebra, and satisfies the universal property 9.3.1. Therefore, we have canonical isomorphism
$$J(𝒜)\stackrel{}{}\stackrel{~}{J}$$
$`(\mathrm{9.5.1})`$
We can apply the construction of 9.4 to the algebra $`\stackrel{~}{J}`$ as well, and for it the claims of Theorem 9.5 are clear. On the other hand, the maps $`\mathrm{\Omega }^{(i)}\stackrel{~}{\mathrm{\Omega }}^{(i)}`$ induced by (9.5.1) are isomorphisms. This implies the theorem. $``$
Conformal algebra associated with a vertex algebroid
9.6. Following the pattern of 3.1, let us define a $`1`$-truncated conformal algebra to be a quintuple $`c=(C_0,C_1,,_{(0)},_{(1)})`$ where $`C_i`$ are $`k`$-modules, $`:C_0C_1`$ is a $`k`$-module map and $`{}_{(i)}{}^{}:(C_0C_1)^2(C_0C_1)`$ are bilinear operations of degree $`i1`$.
Elements of $`C_0`$ (resp. of $`C_1`$) will be denoted $`a,b,c`$ (resp. $`x,y,z`$).
These data must satisfy the following axioms.
(Derivation)
$$(a)_{(0)}=0;(a)_{(1)}=a_{(0)};(x_{(0)}a)=x_{(0)}a$$
$`(Der)`$
(Commutativity)
$$x_{(0)}a=a_{(0)}x;x_{(0)}y=y_{(0)}x+(y_{(1)}x);x_{(1)}y=y_{(1)}x$$
$`(Com)`$
(Associativity) For all $`\alpha ,\beta ,\gamma C_0C_1`$ and $`i=0,1`$,
$$\alpha _{(0)}\beta _{(i)}\gamma =(\alpha _{(0)}\beta )_{(i)}\gamma +\beta _{(i)}\alpha _{(0)}\gamma $$
$`(Ass)`$
$`1`$-truncated conformal algebras, with obvious morphisms, form a category $`𝒞onf_1`$.
We have an evident forgetful functor
$$𝒱ert_1𝒞onf_1$$
$`(\mathrm{9.6.1})`$
9.7. We have an obvious truncation functor
$$𝒞onf𝒞onf_1$$
$`(\mathrm{9.7.1})`$
which assigns to a conformal algebra $`C`$ its part $`C_1`$ of degree $`1`$. Let us construct a left adjoint to this functor.
Given $`c=(C_0,C_1,\mathrm{})𝒞onf_1`$, consider $`c`$ as a $`1`$-truncated $``$-module (forgetting the operations), and consider the corresponding $``$-module $`C:=c^{}`$, cf. 9.1.
9.8. Theorem. There is a unique structure of a conformal algebra on $`C`$ such that operations <sub>(i)</sub> and $``$ on the subspace $`C_1`$ coincide with the ones given by the structure of a $`1`$-truncated conformal algebra on $`c`$.
The assignement $`cC`$ gives a functor
$$𝒞onf_1𝒞onf$$
$`(\mathrm{9.8.1})`$
left adjoint to (9.7.1).
Indeed, uniqueness is clear, due to the axioms (0.4.2) and (0.4.3) of a conformal algebra. We leave the proof of existence to the reader. The claim about adjointness is evident.
9.9. Let $`𝒜=(A,T,\mathrm{\Omega },\mathrm{})`$ be a vertex algebroid. We assign to $`𝒜`$ a $`1`$-truncated conformal algebra $`c𝒜=(C_0,C_1,\mathrm{})`$ by setting $`C_0=A,C_1=T\mathrm{\Omega },:C_0C_1`$ to be the composition of $`:A\mathrm{\Omega }`$ with the embedding $`\mathrm{\Omega }C_1`$, and defining the operations <sub>(i)</sub> by
$$a_{(0)}b=a_{(0)}\omega =\omega _{(0)}\omega ^{}=0;\tau _{(0)}a=\tau (a);\tau _{(0)}\omega =\tau (\omega )$$
$`(\mathrm{9.9.1})`$
$$\tau _{(0)}\tau ^{}=[\tau ,\tau ^{}]c(\tau ,\tau ^{})+\frac{1}{2}\tau ,\tau ^{}$$
$`(\mathrm{9.9.2})`$
$$x_{(1)}y=x,y$$
$`(\mathrm{9.9.3})`$
In other words, $`c𝒜`$ is the result of application of the forgetful functor (9.6.1) to the $`1`$-truncated vertex algebra $`u𝒜`$ defined in 3.3.
9.10. We may apply the functor (9.8.1) to the $`1`$-truncated conformal algebra $`c𝒜`$, and get a conformal algebra, to be denoted $`C𝒜`$.
Assume for simplicity that $`k`$. Then, as a $`k`$-module,
$$C𝒜=A(_{i0}^{(i)}T)(_{i0}^{(i)}\mathrm{\Omega })$$
$`(\mathrm{9.10.1})`$
Here are the explicit formulas for the operations <sub>(n)</sub> (below we agree that $`^{(i)}\alpha =0`$ if $`i<0`$). First of all, $`\alpha _{(n)}\beta =0`$ for all $`\alpha ,\beta A(_{i0}^{(i)}\mathrm{\Omega })`$ and $`n0`$. Next,
$$a_{(n)}^{(i)}\tau =^{(in)}\tau (a);\tau _{(n)}^{(i)}a=^{(in)}\tau (a)$$
$`(\mathrm{9.10.2})`$
$$\omega _{(n)}^{(i)}\tau =^{(in)}\tau (\omega )+(i+1)^{(in+1)}\omega ,\tau $$
$`(\mathrm{9.10.3})`$
$$\tau _{(n)}^{(i)}\omega =^{(in)}\tau (\omega )+n^{(in+1)}\tau ,\omega $$
$`(\mathrm{9.10.4})`$
Finally,
$$\tau _{(n)}^{(i)}\tau ^{}=^{(in)}\left\{[\tau ,\tau ^{}]c(\tau ,\tau ^{})\right\}+\frac{1}{2}(i+n+1)^{(in+1)}\tau ,\tau ^{}$$
$`(\mathrm{9.10.5})`$
Construction of envelope
9.11. Now we can construct the vertex envelope of a vertex algebroid $`𝒜=(A,T,\mathrm{\Omega },\mathrm{})`$. Consider the conformal algebra $`C𝒜`$; our vertex envelope $`U𝒜`$ will be a quotient of the vertex envelope $`UC𝒜`$ defined in the previous Section by certain vertex ideal.
To define this ideal, let us return to the tensor algebra $`T:=TC𝒜`$, as in 8.2. Let us define a left $``$-ideal $`RT`$ generated by the following elements:
$$r_0(u):=(1_A1_T)u,uT$$
$`(\mathrm{9.11.1})`$
$$r(a,x):=axax,aA,xA\mathrm{\Omega };r(a,\tau )=a\tau a\tau +\gamma (a,\tau )$$
$`(\mathrm{9.11.2})`$
Here by $``$-ideal we mean an ideal closed under all operators $`^{(i)}`$.
9.12. Lemma. The ideal $`R`$ is respected by all operations $`y_{(n)},yC𝒜,n`$.
Indeed, one has to check this for $`yAT\mathrm{\Omega }`$, $`n0`$, and only on generators of $`R`$. This is done by an easy case-by-case computation that we will leave to the reader, restricting ourselves by an example.
Let us check that $`\tau _{(0)}r(a,\tau ^{})S`$. We have
$$\tau _{(0)}r(a,\tau ^{})=\tau _{(0)}\left\{a\tau ^{}a\tau ^{}+\gamma (a,\tau ^{})\right\};$$
$$\tau _{(0)}(a\tau ^{})=\tau _{(0)}a\tau ^{}+a\tau _{(0)}\tau ^{}=\tau (a)\tau ^{}+a([\tau ,\tau ^{}]c(\tau ,\tau ^{})+\frac{1}{2}\tau ,\tau ^{});$$
$$\tau _{(0)}(a\tau ^{})=[\tau ,a\tau ^{}]+c(\tau ,a\tau ^{})\frac{1}{2}\tau ,a\tau ^{}$$
We have
$$[\tau ,a\tau ^{}]=a[\tau ,\tau ^{}]\tau (a)\tau ^{};$$
by axiom (A3) of a vertex algebroid,
$$c(\tau ,a\tau ^{})=c(a\tau ^{},\tau )=ac(\tau ^{},\tau )\gamma (a,[\tau ^{},\tau ])+\gamma (\tau (a),\tau ^{})\tau (\gamma (a,\tau ^{}))+$$
$$+\frac{1}{2}\tau ^{},\tau a\frac{1}{2}\tau ^{}\tau (a)+\frac{1}{2}\tau ,\gamma (a,\tau ^{})$$
By axiom (A2),
$$\frac{1}{2}\tau ,a\tau ^{}=\frac{1}{2}\left\{a\tau ^{},\tau +\gamma (a,\tau ^{}),\tau \tau ^{}\tau (a)\right\}=$$
$$=\frac{1}{2}\tau ^{},\tau a\frac{1}{2}a\tau ^{},\tau \frac{1}{2}\gamma (a,\tau ^{}),\tau +\frac{1}{2}\tau ^{}\tau (a)$$
Finally,
$$\tau _{(0)}\gamma (a,\tau ^{})=\tau (\gamma (a,\tau ^{}))$$
After bookkeeping, we see that
$$\tau _{(0)}r(a,\tau ^{})=r(a,[\tau ,\tau ^{}])+r(\tau (a),\tau ^{})r(a,c(\tau ,\tau ^{}))+\frac{1}{2}r(a,\tau ,\tau ^{})S,QED$$
9.13. Let $`\overline{R}UC𝒜`$ denote the image of $`R`$ under the canonnical projection $`TC𝒜UC𝒜`$. The previous lemma means that $`\overline{R}`$ is a vertex ideal in $`UC𝒜`$. Let us denote by $`U𝒜`$ the vertex algebra $`UC𝒜/\overline{R}`$. It is equipped with an obvious splitting $`TT\mathrm{\Omega }=U𝒜_1`$.
Conversely, given a splittable vertex algebra $`V`$, choose a splitting $`s:TV_1`$, and consider the vertex algebroid $`𝒜(V;s)`$, cf. §2. Let $`𝒱ert^{}𝒱ert`$ denote the full subcategory of splittable vertex algebras. We get a functor
$$𝒜:𝒱ert^{}𝒜lg$$
$`(\mathrm{9.13.1})`$
which is in fact the composition of the truncation functor (3.1.1) (restricted to $`𝒱ert^{}`$) and of the functor quasiinverse to (3.3.6).
9.14. Theorem. The assignement $`𝒜U𝒜`$ provides a functor
$$U:𝒜lg𝒱ert^{}$$
$`(\mathrm{9.14.1})`$
which is left adjoint to (9.13.1).
Indeed, this is evident from the construction.
Poincaré-Birkhoff-Witt
9.15. Let $`𝒜=(A,T,\mathrm{\Omega },\mathrm{})`$ be a vertex algebroid. The enveloping algebra $`U𝒜`$ is generated as a $`k`$-module by the monomials of the form
$$x^px^{p1}\mathrm{}x^1$$
$`(\mathrm{9.15.1})`$
where $`x^i`$ has the form $`^{(j)}y,y=a,\omega `$ or $`\tau `$. Here we have denoted for brevity by $`xx^{}`$ the operation $`x_{(1)}x^{}`$.
Let us introduce a canonical increasing exhaustive filtration $`F_0U𝒜F_1U𝒜\mathrm{}U𝒜`$ by setting $`F_iU𝒜`$ to be equal the $`k`$-submodule of $`U𝒜`$ generated by all monomials (9.15.1) where there are $`i`$ letters $`x_j`$ of the form $`^{(a)}\tau `$.
Obviously all submodules $`F_iU𝒜`$ are stable under $`^{(j)}`$ and $`F_iU𝒜F_jU𝒜F_{i+j}U𝒜`$. Consider the associated graded module $`gr_FU𝒜=_{i0}F_iU𝒜/F_{i1}U𝒜`$.
It is easy to see that the operation $`xy`$ induces a commutative and associative multiplication on $`gr_FU𝒜`$, i.e. $`gr_FU𝒜`$ becomes a $``$-algebra (an abelian vertex algebra).
Let $`𝒜^{ab}`$ denote an abelian vertex algebroid $`(A,\mathrm{\Omega }T,)`$. We have
$$𝒜(gr_FU𝒜)=𝒜^{ab}$$
$`(\mathrm{9.15.2})`$
hence by adjunction we get a canonical map of $``$-algebras
$$J(𝒜^{ab})gr_FU𝒜$$
$`(\mathrm{9.15.3})`$
It is clear the this map is surjective.
Recall that by 9.4 we have a canonical filtration on $`J(𝒜^{ab})`$ (let us denote it by $`G`$ here) and a canonical surjective map of $``$-algebras
$$Sym_A\left\{(_{i1}T^{(i)})(_{i1}\mathrm{\Omega }^{(i)})\right\}gr_GJ(𝒜^{ab})$$
$`(\mathrm{9.15.4})`$
where $`X^{(i)}(X=T`$ or $`\mathrm{\Omega }`$) denotes a copy of $`X`$ sitting in weight $`i`$.
By refining the filtration $`F`$ using the filtration $`G`$, we get a canonical filtration $`H`$ on $`U𝒜`$ together with a surjective map (of $``$-algebras over $`A`$)
$$Sym_A\left\{(_{i1}T^{(i)})(_{i1}\mathrm{\Omega }^{(i)})\right\}gr_HU𝒜$$
$`(\mathrm{9.15.5})`$
The multiplication in the right hand side is induced by the operation <sub>(-1)</sub> on $`U𝒜`$.
9.16. Theorem. Assume that $`k`$ and both $`\mathrm{\Omega }`$ and $`T`$ are free $`A`$-modules. Then the maps (9.15.3) and (9.15.5) are isomorphisms.
Filtration $`H`$ is compatible with the conformal grading, finite on each component with the fixed conformal weight and is canonical in following sense: for any morphism of vertex algebroids $`𝒜𝒜^{}`$, the corresponding morphism of vertex envelopes $`U𝒜U𝒜^{}`$ respects these filtrations.
To prove this, we use the strategy of Serre’s proof of the usual PBW theorem, cf. \[Se\], pp. 14-16. Let us choose well ordered $`A`$-bases $`\{\tau _s\},\{\omega _r\}`$ of $`T`$ and $`\mathrm{\Omega }`$. They give rise to a well ordered $`A`$-base
$$\{x_u\}_{uI}=\{^{(a)}\tau _s,^{(b)}\omega _r\}$$
of $`(_{i1}T^{(i)})(_{i1}\mathrm{\Omega }^{(i)})`$. For a sequence $`M=(u_1,\mathrm{},u_m),u_jI,u_1u_2\mathrm{}u_m`$ define a monomial $`x_MU𝒜`$ by
$$x_M=x_{u_1}\mathrm{}x_{u_m}$$
Similarly, $`M`$ defines monomials $`\overline{x}_Mgr_HU𝒜`$ and
$$\stackrel{~}{x}_MS:=Sym_A\left\{(_{i1}T^{(i)})(_{i1}\mathrm{\Omega }^{(i)})\right\}$$
Obviously, the monomials $`\{\stackrel{~}{x}_M\}`$ form an $`A`$-base of $`S`$. On the other hand, it is easy to see that each element of $`U𝒜`$ may by written as $`a_Mx_M`$ with some $`a_MA`$. To prove our claim it is enough to show that a relation
$$a_Mx_M=0$$
$`(\mathrm{9.16.1})`$
implies that all $`a_M=0`$ (cf. loc. cit., Lemma 4.5).
Proceeding in a manner similar to 8.4, we define an action of the Lie algebra $`C𝒜^{Lie}`$, and hence of its envelope $`UC𝒜^{Lie}=UC𝒜`$, on $`S`$. Next one checks the relations (9.11.1) and (9.11.2) and therefore gets an action of $`U𝒜`$ on $`S`$.
One sees immediately from the definitions that $`(ax_M)1_S=a\stackrel{~}{x}_M`$. Hence (9.16.1) implies the relation $`a_M\stackrel{~}{x}_M=0`$ in $`S`$ and therefore $`a_M=0`$, which proves the theorem. We leave the details to the reader. $``$
9.17. Let $`X`$ be a smooth $`k`$-scheme, $`UX`$ a Zariski open in $`X`$. A vertex algebra of the form $`U𝒜`$ where $`𝒜𝒜lg_{𝒯_X}(U)`$ is a section of the gerbe $`𝒜lg_{𝒯_X}`$ discussed in Section 7, is called an algebra of chiral differential operators on $`U`$.
These algebras form a gerbe $`𝒟iff_X`$, by definition isomorphic to the gerbe $`𝒜lg_X`$.
Note that all isomorphisms of algebras of chiral differential operators respect the canonical filtrations on them, and there is no obstruction to the gluing of associated graded algebras.
In particular, we have
9.18. Theorem. Assume that $`k`$. Each algebra of chiral do $`𝒟_X𝒟iff_X(X)`$ admits a canonical filtration whose graded algebra is isomorphic to
$$gr(𝒟_X)=Sym_{𝒪_X}\left\{(_{i1}𝒯_X^{(i)})(_{i1}\mathrm{\Omega }_X^{1(i)})\right\}$$
$`(\mathrm{9.18.1})`$
where $`𝒯_X^{(i)}`$ (resp. $`\mathrm{\Omega }_X^{1(i)}`$) denotes a copy of the tangent bundle (resp. of the bundle of $`1`$-forms) sitting in conformal degree $`i`$. $``$
References
\[BB\] A. Beilinson, J. Bernstein, A proof of Jantzen conjectures, I.M.Gelfand Seminar, Adv. in Soviet Math. 16, Part 1, AMS, Providence, RI, 1993, pp. 1-50.
\[BD1\] A. Beilinson, V. Drinfeld, Excerpt from ”Chiral Algebras II” on chiral tdo and Tate structures, Preprint, 1999.
\[BD2\] A. Beilinson, V. Drinfeld, Chiral algebras, Preprint, Version 2000.
\[B\] R. Borcherds, Vertex algebras, Kac-Moody Lie algebras, and the Monster, Proc. Nat. Acad. Sci. USA, 83 (1986), no. 10, 3068-3071.
\[GMS\] V. Gorbounov, F. Malikov, V. Schechtman, Gerbes of chiral differential operators, math.AG/9906117; Math. Research Letters, 7 (2000), 55-66.
\[Gr\] A. Grothendieck, Classes de Chern et répresentations linéaires des groupes discrets, pp. 215-305 in: Dix exposés sur la cohomologie des schémas, Masson & Cie, Paris — North Holland, Amsterdam, 1968.
\[H\] B. Harris, Group cohomology classes with differential form coefficients,
Algebraic $`K`$-theory, Evanston 1976, Lect. Notes in Math. 551, Springer, 1976; pp. 278-282.
\[K\] V. Kac, Vertex algebras for beginners, Second Edition, University Lecture Series, 10, American Mathematical Society, Providence, Rhode Island, 1998.
\[MS\] F. Malikov, V. Schechtman, Chiral de Rham complex. II, D.B.Fuchs’ 60-th Anniversary volume, Differential topology, infinite-dimensional Lie algebras, and applications, 149-188, AMS Transl. Ser. 2, 194, AMS, Providence, RI, 1999.
\[MSV\] F. Malikov, V. Schechtman, A. Vaintrob, Chiral de Rham complex, Comm. Math. Phys., 204 (1999), 439-473.
\[S\] V. Schechtman, Riemann-Roch theorem and the Atiyah-Hirzebruch spectral sequence, Usp. Mat. Nauk, 35, no. 6 (1980), 179-180 (Russian).
\[Se\] J.-P. Serre, Lie algebras and Lie groups, 1964 lectures given at Harvard University, Second Edition, Lect. Notes in Math. 1500, Springer-Verlag, Berlin, 1992.
V.G.: Department of Mathematics, University of Kentucky, Lexington, KY 40506, USA; vgorbms.uky.edu
F.M.: Department of Mathematics, University of Southern California, Los Angeles, CA 90089, USA; fmalikovmathj.usc.edu
V.S.: IHES, 35 Route de Chartres, 91440 Bures-sur-Yvette, France; vadikihes.fr
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# Abstract
## Abstract
We show that and how the Coulomb potential $`V(x)=Ze^2/x`$ can be regularized and solved exactly at the imaginary coupling $`Ze^2`$. The new spectrum of energies is real and bounded as expected, but its explicit form proves totally different from the usual real-coupling case.
## 1 Introduction
Quantum mechanics often works with the exactly solvable simplified models. For the precise fits of data or for some more subtle quantitative analyses, unfortunately, the number of solvable models is too limited. In $`D`$ dimensions, the only really useful and easily tractable interactions are harmonic oscillators and/or the central Coulomb well $`V^{(Z)}(\stackrel{}{r})=Ze^2/|\stackrel{}{r}|`$. A new way out of this deadlock emerges within the framework of the alternative, so called “$`𝒫𝒯`$ symmetric” quantum mechanics. With its complex Hamiltonians $`H`$ breaking both the parity $`𝒫`$ and the time-reflection symmetry $`𝒯`$ and commuting only with their product $`𝒫𝒯`$, this formalism was proposed by Bessis and by Bender et al as a possible way towards weakening of the standard requirements of Hermiticity.
Several new exactly solvable $`𝒫𝒯`$ symmetric models have been proposed recently . This is a promising development with possible applications ranging from field theories to supersymmetric models and from quasi-classical methods to perturbation theory .
Even the solvable harmonic oscillator itself acquires a richer spectrum after its consequent $`𝒫𝒯`$ symmetric regularization in $`D`$ dimensions . The detailed structure of spectrum of this prominent example does not in fact offer any really serious surprise. A manifest violation of the parity $`𝒫`$ is compensated by an emergence of the so called quasi-parity $`q=\pm 1`$ tractable as a signature of two equidistant subspectra. The new quantum number $`q`$ degenerates back to the eigenvalue of parity after a return to the standard Hermitean and one-dimensional oscillator.
No immediate surprise emerges also for the quartic anharmonic oscillator . The situation only becomes less clear after one moves towards the asymptotically vanishing models. They exhibit several counterintuitive properties and open new mathematical challenges . In particular, the popular Coulomb potential did not not even seem particularly suitable for any immediate $`𝒫𝒯`$ symmetric regularization . A psychological barrier has been created by the numerical and semiclassical studies of the general power-law forces $`V(x)(ix)^\delta `$. They may be well defined everywhere near the harmonic exponents $`\delta =2`$, $`\delta =6`$ etc . At the same time, the related analyses hinted that it is apparently difficult to move beyond the Herbst’s singularity located at $`\delta =1`$ .
In what follows, we intend to employ a slightly different strategy and try to study the Coulomb problem directly, via its well known correspondence to the harmonic oscillator. This correspondence is based on an elementary change of variables. Its background dates back to the nineteenth century mathematics and, in particular, to the work of Liouville . The Newton’s monograph cites also Fivel as a newer source of the idea. In the contemporary literature (cf., e.g., for further references) people usually speak about the Kustaanheimo - Steifel (KS) transformation . In all the implementations of this idea the parameters appearing in the Coulombic and oscillator problems are interrelated of course. Details will be mentioned below. Preliminarily, let us only warn the reader that all the KS-type mappings can also change the dimensions and angular momenta and that the energies of one problem are related to the coupling constants of the other one and vice versa. Within the ”normal” quantum mechanics, all this has already been thoroughly discussed elsewhere: In ref. for $`D=3`$ and in ref. for the continuous transformation between Coulomb problems and harmonic oscillators in various dimensions.
## 2 Liouvillean changes of variables
The change-of-variable approach to the Coulombic bound-state problem enables us to start directly from the harmonic oscillator potential $`W(r)=r^2`$ or, in the present less traditional context, from its $`𝒫𝒯`$ symmetric radial Schrödinger equation
$$\left[\frac{d^2}{dr^2}+\frac{l(l+1)}{r^2}+W(r)\right]\chi (r)=\epsilon ^2\chi (r)$$
(1)
of ref. , using the complex coordinate $`r=xic`$ with real $`x(\mathrm{},\mathrm{})`$ and with, say, positive $`c>0`$. This means that the integration path has been shifted down from the position where it would cross the strong centrifugal singularity. Such a regularization preserves the asymptotic decrease of the normalizable solutions. Only in the limit $`c0`$ and beyond the trivial one-dimensional case one has to omit all the so called irregular solutions defined by their “physically unnacceptable” $`\chi (r)r^l`$ behaviour near the origin.
Within the framework of the general Liouville method the change of variables mediates a transition to the different potential $`V(t)`$. It is easy to show that once we forget about boundary conditions one merely has to demand the existence of an invertible function $`r=r(t)`$ and its few derivatives $`r^{}(t),r^{\prime \prime }(t),\mathrm{}`$. Then, the explicit correspondence between the two bound state problems may be explicitly given by the elementary formulae. From our original eq. (1) (i.e., in our case, harmonic oscillator) one obtains the new (i.e., in our case, Coulombic) Schrödinger equation
$$\left[\frac{d^2}{dt^2}+\frac{L(L+1)}{t^2}+V(t)\right]\mathrm{\Psi }(t)=E\mathrm{\Psi }(t)$$
(2)
with the new wave functions
$$\mathrm{\Psi }(t)=\chi [r(t)]/\sqrt{r^{}(t)}$$
(3)
and with the new interaction and the new energies ,
$$\frac{L(L+1)}{t^2}+V(t)E=\left[r^{}(t)\right]^2\left\{\frac{l(l+1)}{r^2(t)}+W[r(t)]\epsilon ^2\right\}+\frac{3}{4}\left[\frac{r^{\prime \prime }(t)}{r^{}(t)}\right]^2\frac{1}{2}\left[\frac{r^{\prime \prime \prime }(t)}{r^{}(t)}\right].$$
Thus, it only remains for us to re-analyse the boundary conditions.
## 3 $`𝒫𝒯`$ symmetric KS transformation
Without any serious formal difficulties let us extend the scope of the present considerations to all the singular forces $`\widehat{W}(r)=W(r)+f/r^2`$ and and/or $`\widehat{V}(t)=V(t)+F/t^2`$. Both these central forces may act in the respective $`d`$ and $`D`$ dimensions. This means that
$$l(l+1)=\left\{[j+(d3)/2][j+(d1)/2]+f\right\}$$
or rather
$$(l+1/2)^2(=\alpha ^2)=[j+(d2)/2]]^2+f$$
(with Re $`\alpha >0`$) and
$$(L+1/2)^2(=A^2)=[J+(D2)/2]]^2+F$$
(with Re $`A>0`$) where the partial waves are numbered by the respective integers $`j=0,1,\mathrm{}`$ and $`J=0,1,\mathrm{}`$.
An important simplification of our effort is provided by our knowledge of the complete harmonic oscillator solution as derived in ref. . Its two equidistant subsets of energies
$$\epsilon ^2=\epsilon _{(n,q)}^2=4n+22q\alpha ,q=\pm 1,n=0,1,\mathrm{}$$
correspond to the two families of the Laguerre-polynomial wave functions
$$\chi _{(n,q)}(r)=𝒩r^{1/2q\alpha }e^{r^2/2}L_n^{(q\alpha )}(r^2).$$
Integration path $`r=xic`$ lies in the lower half of the complex plane and does not change after the subsequent $`𝒫`$ and $`𝒯`$ transformations $`r=r(x)r`$ and $`rr^{}=r(x)`$.
In the spirit of the above-mentioned KS mapping of harmonic oscillators on Coulombic bound states we now have to define a complex variable $`t`$ as a re-scaled square of $`r(x)`$ such that the resulting path $`t(x)`$ remains $`𝒫𝒯`$ invariant. Our requirement implies that the complex plane of $`r`$ will cover twice the complex plane of $`t`$. In such an arrangement, our lower half plane of $`r`$ should cover the Riemann sheet given as a whole plane of $`t`$ which is cut upwards from the origin. In the polar representation, one has $`r\mathrm{exp}(i\phi )`$ mapped upon $`t\mathrm{exp}(2i\phi )`$ with $`\phi (0,\pi )`$. Once we introduce a suitable free parameter $`\kappa >0`$ we can put, say, $`r^2=2\kappa ^2z`$ and then rotate the $`z`$plane (which is cut, by construction, along the real and positive semi-axis) by the angle $`\pi /2`$ giving $`t=iz`$. Our final recipe
$$r^2=2i\kappa ^2t$$
(4)
maps the above-mentioned straight line $`r(x)=xic`$ upon a curve $`t(x)=u+iv`$. Its real part $`u=u(x)=xc/\kappa ^2`$ and imaginary part $`v=v(x)=(x^2c^2)/2\kappa ^2`$ form, as required, a $`𝒫𝒯`$ symmetric and upwards-oriented parabola $`v=c^2/2\kappa ^2+(\kappa ^2/2c^2)u^2`$ in complex plane. Obviously, a small asymptotic deformation of our original curve $`r(x)`$ with the modified shifts $`c=c(x)1/x^{1+\eta }`$ would transform the parabola to a pair of lines which are parallel to the cut in the asymptotic domain of $`|x|1`$.
Having achieved a $`𝒫𝒯`$ symmetry in the complex plane of $`t`$, we may move to the (trivial) insertions and conclude that all the above-mentioned harmonic oscillator bound-state solutions are in a one-to-one correspondence to the solutions of the Coulombic Schrödinger equation (2),
$$\left[\frac{d^2}{dt^2}+\frac{L(L+1)}{t^2}+i\frac{Ze^2}{t}\right]\mathrm{\Psi }(t)=E\mathrm{\Psi }(t),t=u(x)+iv(x),xIR.$$
(5)
The underlying assignment of constants is such that $`\alpha =2A`$ while $`\kappa `$ itself becomes $`n`$dependent, $`\kappa ^2=2Ze^2/\epsilon ^2=Ze^2/(2n+12qA)`$. In full detail one gets the new, Laguerre-polynomial wave functions
$$\mathrm{\Psi }_{(n,q)}(t)=t^{1/2qA}e^{i\kappa ^2t}L_n^{(2qA)}(2i\kappa ^2t)$$
(6)
and their energy spectrum specified by the elementary formula
$$E_{(n,q)}=\kappa _{(n,q)}^4=\frac{Z^2e^4}{(2n+12qA)^2}q=\pm 1,n=0,1,\mathrm{}.$$
(7)
This is our main result.
## 4 Discussion
### 4.1 Consequences of the curvature of our integration path
The latter two formulae exhibit several unusual features. The first concerns the asymptotics of the wave functions which are determined by the decreasing exponential $`\mathrm{exp}(i\kappa ^2t)`$. Its form re-confirms the correctness of the above, slightly counter-intuitive KS-dictated choice of our $`𝒫𝒯`$ symmetric integration path. Asymptotically, it encircles more or less closely the positive imaginary axis in $`t`$ plane. This clarifies the apparent paradox.
The second unexpected result is the positivity and unusual $`n`$dependence of the energies. This can be related to the choice of the KS integration path $`t(x)`$ again. In the very vicinity of the origin, one can visualize this path as a circle with radius $`\sigma `$,
$$u^2(x)+v^2(x)=\sigma ^2,|x|1.$$
From the appropriate definitions we get the formula
$$\sigma =c^2(0)/2\kappa _{(n,q)}^2+𝒪(x^2)$$
and see that this radius is $`n`$dependent and increases with the growth of this principal quantum number, $`\sigma nc^2(0)/Ze^2`$, $`n1`$. As a consequence, an “effective charge” of our $`𝒫𝒯`$ symmetric Coulomb potential appears to decrease with $`n`$ since
$$\left|\frac{iZe^2}{t}\right|\frac{Ze^2}{\sigma }+𝒪(t)=𝒪(1/n).$$
This offers a “rule-of-thumb” guide to the unusual and certainly counterintuitive $`n`$dependence of the energy levels (7). Of course, in practice, a preferred integration path will be $`n`$independent. In such a case, the $`n`$dependence re-appears in the small-$`x`$ deformation of the initial harmonic-oscillator path with $`c=c_n(x)=𝒪(1/n)`$. Such a flexible transfer of the excitation-dependence throws also a new light on the complexified KS transformation itself.
### 4.2 “Flown-away” energies and unavoided crossings
Let us return in more detail to the $`A`$dependence of our energies (7). Firstly, we notice their power-law dependence on $`n`$ and $`A`$ (with exponent = -2) as somewhat similar to the spectra in $`𝒫𝒯`$ symmetric oscillator well (with exponent = +1) and in the Morse potential of ref. (with exponent = +2).
In the present case, obviously, we have to distinguish between the two separate families $`E_{(n,q)}`$ with $`q=+1`$ (cf. Figure 1) and $`q=1`$ (cf. Figure 2). The latter set is, up to its sign, analogous to the ordinary Coulombic spectrum. By far not so the former one. Its energies enrich and dominate the spectrum. The $`n_{div}`$th energy “flies away” and disappears from the spectrum at $`A_{div}=n_{div}+1/2`$. Moreover, at all the positive integers and half-integers $`A`$ one encounters the unavoided level crossings. In contrast to the harmonic case, they appear at both the opposite and identical (viz., positive) quasiparities $`q`$. The former case takes place at $`A=A_{crit}=(nn^{})/2`$ while in the latter case we must fulfill the condition $`A=A_{crit}=(n+n^{}+1)/2`$. A sample of this phenomenon is given in Figure 3.
Formally the unavoided crossings generate certain identities which connect different Laguerre polynomials (cf. their sample in ref. ). In applications, these “critical” cases are not exceptional at all. For $`F=0`$ forces without a spike, the critical integer or half-integer coordinates $`A_{crit}=J1+D/2`$ correspond precisely to the physical (namely, integer) dimensions $`D`$ and partial waves $`J`$.
In the conclusion, let us not forget about many open questions. Pars pro toto, one could mention a not yet clear possibility of re-interpretation of our bound states, say, in the limit $`treal`$, i.e., beyond the mathematical and apparently natural boundaries of our present approach. Moreover, we must keep in mind that via our complexification of the coordinates we of course broke their immediate connection to any standard $`D`$dimensional problem. In this sense, our $`F=0`$ and $`F0`$ Hamiltonians differ just in an inessential way. One could even prefer the latter, Kratzer-like option as a model which is formally simpler, due to the generic absence of the puzzling unavoided crossings. At $`F0`$ the structure of the spectrum of our present Coulomb model becomes also richer and, in this sense, more interesting.
## Acknowledgements
M. Z. thanks for the hospitality of the Institute of Nuclear Research of the Hungarian Academy of Sciences, in Debrecen, and appreciates the support by the grant Nr. A 1048004 of GA AS CR. G. L. acknowledges the OTKA grant no. T031945.
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