id
stringlengths 27
33
| source
stringclasses 1
value | format
stringclasses 1
value | text
stringlengths 13
1.81M
|
|---|---|---|---|
warning/0003/hep-th0003168.html
|
ar5iv
|
text
|
# 1 Introduction and conclusions
## 1 Introduction and conclusions
M-theory compactified on $`O_7X_6\times S^1/𝐙_2`$, where $`X_6`$ is a Calabi-Yau (CY) three-fold, leads to a four-dimensional theory with $`N=1`$ local supersymmetry. In the low-energy limit, M-theory information can be organized as an expansion in powers of the eleven-dimensional gravitational constant $`\kappa _{11}`$ . The lowest order $`\kappa _{11}^2`$ is eleven-dimensional supergravity . In a compactification on $`S^1/𝐙_2`$ only, the next orders are known to include orbifold plane contributions as well as gauge and gravitational anomaly-cancelling terms . Similarly, the effective four-dimensional supergravity can be formulated as an expansion in the four-dimensional gravitational constant $`\kappa `$, even if string theory rather suggests to use the dilaton as expansion parameter. The lowest order $`\kappa ^2`$ is the $`S^1/𝐙_2`$ truncation of eleven-dimensional supergravity on a CY three-fold. The next order includes super-Yang-Mills (SYM) and charged matter kinetic and superpotential contributions. Sigma-model anomaly-cancelling terms modifying in particular the gauge thresholds are then also involved. These first corrections to the low-energy limit of M-theory compactifications on $`O_7`$ are identical to those obtained from heterotic compactifications on CY. The literature gives a detailed description of these results, with particular attention paid to the ‘strong-coupling heterotic limit’ in which the size of the CY space is smaller than the orbifold length, supersymmetry breaking by gaugino condensation and non-standard embeddings .
In this note, we give the structure of the four-dimensional $`N=1`$ wilsonnian effective supergravity describing the universal massless sector of M-theory compactified on $`O_7`$. We begin by writing the theory corresponding to the reduction of the bulk eleven-dimensional supergravity directly in terms of four-dimensional ‘M-theory supermultiplets’. The supersymmetrized Bianchi identities for the components of the M-theory tensor field strength are promoted to equations of motion using ‘Lagrange multiplets’. Within this ‘off-shell’ approach, we can then introduce ‘source multiplets’ to take into account the contributions of the $`S^1/𝐙_2`$ planes which appear as modifications of the Bianchi identities. This formulation is also particularly appropriate for the inclusion of non-perturbative states (M-theory five-branes, condensates, etc.).
The material presented here is detailed in ref. and a forthcoming publication will contain a direct application of our approach (the coupling of five-brane moduli to the background).
## 2 The bulk Lagrangian
In this section, we establish our basic procedure by considering the well-known ‘bulk dynamics’, which follows from $`O_7`$ compactification of eleven-dimensional supergravity. The resulting Lagrangian is the lowest order in the $`\kappa `$-expansion and describes Kaluza-Klein (KK) massless modes of eleven-dimensional supergravity.
We will precisely describe two aspects which may be of importance in M-theory compactifications. Firstly, we will introduce chiral, linear or vector supermultiplets with constraints in order to obtain a supersymmetric version of the Bianchi identities satisfied by antisymmetric tensors. Secondly, we will use superconformal supergravity in which we can keep open the choice of gravity frame.
### 2.1 Superconformal formalism
We use the superconformal formulation of $`N=1`$ supergravity with a chiral compensating multiplet $`S_0`$ (with conformal and chiral weights $`w=1`$ and $`n=1`$) to generate Poincaré theories by gauge fixing. In this formalism, a change of frame corresponds to a different Poincaré gauge condition applied on the modulus of the scalar compensator $`z_0`$, which fixes dilatation symmetry. Up to terms with more than two derivatives and up to terms which would contribute to kinetic terms in a fermionic background only , the most general supergravity Lagrangian reads<sup>1</sup><sup>1</sup>1 Except otherwise mentioned, our notations for superconformal expressions are as in refs. , from where the original literature can also be traced back. The appendix of ref. displays the conventions we follow through this note.
$$=\left[S_0\overline{S}_0\mathrm{\Phi }\right]_D+\left[S_0^3W\right]_F+\frac{1}{4}\left[f_{ab}𝒲^a𝒲^b\right]_F.$$
(2.1)
The symbols $`[\mathrm{}]_D`$ and $`[\mathrm{}]_F`$ denote the invariant $`D`$\- and $`F`$-density formulas given by (all fermion contributions are omitted)
$$[𝒱]_D=e(d+\frac{1}{3}cR)\text{ and }[𝒮]_F=e(f+\overline{f}),$$
(2.2)
where $`𝒱`$ is a vector multiplet with components $`(c,\chi ,m,n,b_\mu ,\lambda ,d)`$ and $`𝒮`$ a chiral multiplet with components $`(z,\psi ,f)`$. The real vector multiplet $`\mathrm{\Phi }`$ (zero weights) is a function (in the sense of tensor calculus) of the multiplets present in the theory, including in general the compensating multiplet. The holomorphic function $`W`$ of the chiral multiplets is the superpotential. The chiral multiplet $`𝒲`$ is the gauge field strength for the gauge multiplets and $`f_{ab}`$ is the holomorphic gauge kinetic function of the chiral multiplets. Besides $`S_0`$ and $`𝒲`$, we will use chiral multiplets with zero weights and neither $`W`$ nor $`f_{ab}`$ will depend on the compensator.
Using a $`U(1)`$/Kähler gauge fixing the supergravity Lagrangian (2.1) can also take the form
$$=\left[S_0\overline{S}_0\mathrm{\Phi }\right]_D+c\left[S_0^3\right]_F+\frac{1}{4}\left[f_{ab}𝒲^a𝒲^b\right]_F,$$
(2.3)
with an arbitrary constant $`c`$ as superpotential and two arbitrary functions $`\mathrm{\Phi }`$ and $`f_{ab}`$.
### 2.2 Supermultiplets with constraints
The Lagrangian of eleven-dimensional supergravity can be written as
$$\begin{array}{ccc}\hfill e^1_{\mathrm{CJS}}& =& \frac{1}{2\kappa _{11}^2}R\frac{1}{4\kappa _{11}^2}\frac{1}{4!}G_{M_1M_2M_3M_4}G^{M_1M_2M_3M_4}\hfill \\ \multicolumn{3}{c}{}\\ & & \frac{1}{12\kappa _{11}^2}\frac{1}{4!4!3!}e^1ϵ^{M_1\mathrm{}M_{11}}G_{M_1M_2M_3M_4}G_{M_5M_6M_7M_8}C_{M_9M_{10}M_{11}}\hfill \\ \multicolumn{3}{c}{}\\ & & +\mathrm{fermionic}\mathrm{terms}.\hfill \end{array}$$
(2.4)
Omitting all fields related to the detailed geometry of the CY manifold, the particle content of the four-dimensional theory is the $`N=1`$ supergravity multiplet, with metric tensor $`g_{\mu \nu }`$, and matter multiplets including on-shell four bosons and four fermions. Two bosons are scalars and correspond to the dilaton and the ‘universal modulus’ of the CY space, the massless volume mode. Two bosons are KK modes of the field strength $`G`$, with Bianchi identity $`dG=0`$. Explicitly, these two last fields and their Bianchi identities read<sup>2</sup><sup>2</sup>2 In our notations, $`x^4`$ is the orbifold coordinate.
$$\begin{array}{ccc}G_{\mu \nu \rho 4},\hfill & & _{[\mu }G_{\nu \rho \sigma 4]}=0,\hfill \\ \multicolumn{3}{c}{}\\ G_{\mu j\overline{k}4}=iT_\mu \delta _{j\overline{k}},\hfill & & _{[\mu }T_{\nu ]}=0.\hfill \end{array}$$
(2.5)
It will prove useful to identify these fields with the vector components of two real vector multiplets $`V`$ ($`w=2`$, $`n=0`$) and $`V_T`$ ($`w=n=0`$), and to impose the Bianchi identities as field equations using a chiral multiplet $`S`$ ($`w=n=0`$) and a real linear multiplet $`L_T`$ ($`w=2`$, $`n=0`$) as Lagrange multipliers. The bulk supergravity Lagrangian takes then the form
$$_\mathrm{B}=\left[(S_0\overline{S}_0V_T)^{3/2}(2V)^{1/2}(S+\overline{S})V+L_TV_T\right]_D.$$
(2.6)
The various superconformal multiplets appearing in this Lagrangian have the following components expressions<sup>3</sup><sup>3</sup>3 We only explicitly consider the bosonic sector of the theory and omit all fermions in the $`N=1`$ supermultiplets. We gauge-fix the superconformal symmetries not contained in $`N=1`$ Poincaré supersymmetry, except dilatation symmetry. Notice also that our component expansion of vector multiplets differs in its highest component from refs. .
$$\begin{array}{ccc}\hfill V& =& (C,0,H,K,v_\mu ,0,d\mathrm{}C\frac{1}{3}CR),\hfill \\ \multicolumn{3}{c}{}\\ \hfill V_T& =& (C_T,0,H_T,K_T,T_\mu ,0,d_T\mathrm{}C_T),\hfill \\ \multicolumn{3}{c}{}\\ \hfill S& =& (s,0,f,if,i_\mu s,0,0),\hfill \\ \multicolumn{3}{c}{}\\ \hfill L_T& =& (\mathrm{}_T,0,0,0,t_\mu ,0,\mathrm{}\mathrm{}_T\frac{1}{3}\mathrm{}_TR),\hfill \\ \multicolumn{3}{c}{}\\ \hfill S_0& =& (z_0,0,f_0,if_0,iD_\mu ^cz_0,0,0).\hfill \end{array}$$
(2.7)
The role of the Lagrange multipliers $`S`$ and $`L_T`$ follows from the two relations
$$\begin{array}{ccc}\hfill e^1[(S+\overline{S})V]_D& =& 2Ims^\mu v_\mu +2dResf(HiK)\overline{f}(H+iK)\hfill \\ \multicolumn{3}{c}{}\\ & & +\mathrm{derivative},\hfill \\ \multicolumn{3}{c}{}\\ \hfill e^1[L_TV_T]_D& =& \mathrm{}_T(d_T\mathrm{}C_T)\frac{e}{2}ϵ_{\mu \nu \rho \sigma }(^\mu T^\nu )t^{\rho \sigma }+\mathrm{derivative}.\hfill \end{array}$$
(2.8)
In the last equality, we have used the constraint imposed to the linear multiplet $`L_T`$, $`^\mu t_\mu =0`$, to write $`t_\mu =\frac{e}{2}ϵ_{\mu \nu \rho \sigma }^\nu t^{\rho \sigma }`$. Solving for the components of $`S`$ leads to $`^\mu v_\mu =d=H=K=0`$, and $`V`$ is a linear multiplet $`L`$ ($`w=2`$, $`n=0`$). Solving for the components of $`L_T`$ leads to $`d_T\mathrm{}C_T=_{[\mu }T_{\nu ]}=0`$, and $`V_T`$ can be written as $`T+\overline{T}`$, with a chiral weightless multiplet $`T`$<sup>4</sup><sup>4</sup>4 With components: $`C_T=2ReT`$, $`T_\mu =2_\mu ImT`$, $`H_T=2Ref_T`$, $`K_T=2Imf_T`$.. Since one can always write $`v_\mu =\frac{e}{6}ϵ_{\mu \nu \rho \sigma }v^{\nu \rho \sigma }`$, we have generated with $`Ims`$ and $`t_{\mu \nu }`$ the Bianchi identities $`_{[\mu }v_{\nu \rho \sigma ]}=_{[\mu }T_{\nu ]}=0`$. A modification of these Bianchi identities, as induced by $`S^1/𝐙_2`$ compactification or by five-brane couplings will then be phrased as a modification of the supermultiplets appearing multiplied by $`S+\overline{S}`$ or $`L_T`$ in Eqs. (2.8).
The structure of the Lagrangian (2.6) reflects the familiar duality relating scalars and antisymmetric tensors or, for superfields, chiral and linear multiplets.
Solving in Eq. (2.6) for the Lagrange multipliers $`S`$ and $`L_T`$ leads to the ‘standard form’ of the bulk four-dimensional Lagrangian
$$_{\mathrm{B},\mathrm{l}}=\left[\left(S_0\overline{S}_0e^{\widehat{K}/3}\right)^{3/2}(2L)^{1/2}\right]_D,$$
(2.9)
with the Kähler potential $`\widehat{K}=3\mathrm{log}(T+\overline{T})`$ for the volume modulus $`T`$. We will see again below that this standard form is naturally obtained by direct reduction of the Cremmer, Julia and Scherk version of eleven-dimensional supergravity on $`O_7`$. Clearly, theory (2.9) is also the CY truncation of ten-dimensional $`N=1`$ pure supergravity .
Solving for $`V`$ and $`L_T`$ in Eq. (2.6) leads to the familiar chiral form
$$_{\mathrm{B},\mathrm{c}}=\frac{3}{2}\left[S_0\overline{S}_0e^{K/3}\right]_D,$$
(2.10)
with $`K=\mathrm{log}(S+\overline{S})+\widehat{K}`$.
#### 2.2.1 Choice of Poincaré frame
According to the component expression for the $`D`$-density and the tensor calculus of superconformal multiplets , the Einstein term included in the bulk Lagrangian (2.6) is
$$_\mathrm{E}=\frac{1}{2}eR\left[\left(z_0\overline{z}_0C_T\right)^{3/2}\left(2C\right)^{1/2}\right].$$
(2.11)
As they should, the terms introduced to impose Bianchi identities do not contribute. We then select the Einstein frame, in which the gravitational Lagrangian is $`\frac{1}{2\kappa ^2}eR`$, by the dilatation gauge condition
$$\kappa ^2=(z_0\overline{z}_0C_T)^{3/2}(2C)^{1/2}.$$
(2.12)
It will be convenient to introduce the (composite) real vector multiplet
$$\mathrm{{\rm Y}}=(S_0\overline{S}_0V_T)^{3/2}(2V)^{1/2},$$
(2.13)
with conformal weight two. In the Poincaré theory and in the Einstein frame, its lowest component is equal to $`\kappa ^2`$.
#### 2.2.2 Identification of the components
Choosing the Einstein frame, $`\mathrm{{\rm Y}}=\kappa ^2`$, and solving for the components of $`S`$ and $`L_T`$, the complete bosonic expansion of the four-dimensional supergravity (2.6) is
$$\begin{array}{ccc}\hfill e^1_\mathrm{B}& =& \frac{1}{2\kappa ^2}R\frac{1}{4\kappa ^2}C^2[(_\mu C)(^\mu C)v_\mu v^\mu ]\hfill \\ \multicolumn{3}{c}{}\\ & & \frac{3}{4\kappa ^2}C_T^2[(_\mu C_T)(^\mu C_T)+T_\mu T^\mu ],\hfill \end{array}$$
(2.14)
with $`v_\mu =\frac{e}{2}ϵ_{\mu \nu \rho \sigma }^\nu b^{\rho \sigma }`$ since $`V`$ is a linear multiplet, $`C_T=2ReT`$ and $`T_\mu =2_\mu ImT`$ since $`V_T=T+\overline{T}`$.
This Lagrangian is to be compared with the one we obtain from the reduction of eleven-dimensional supergravity (2.4). The $`𝐙_2`$ orbifold projection eliminates all states which are odd under $`x^4x^4`$, and the reduction of the eleven-dimensional space-time metric is
$$g_{MN}=\left(\begin{array}{ccc}e^\gamma e^{2\sigma }g_{\mu \nu }& 0& 0\\ 0& e^{2\gamma }e^{2\sigma }& 0\\ 0& 0& e^\sigma \delta _{i\overline{j}}\end{array}\right).$$
(2.15)
The surviving components of the field strength $`G_{MNPQ}`$ are only $`G_{\mu \nu \rho 4}`$ and $`G_{\mu i\overline{j}4}`$, with
$$G_{\mu \nu \rho 4}=3_{[\mu }C_{\nu \rho ]4},G_{\mu i\overline{j}4}=_\mu C_{i\overline{j}4},C_{i\overline{j}4}=ia(x)\delta _{i\overline{j}}.$$
(2.16)
The resulting four-dimensional Lagrangian is
$$\begin{array}{ccc}\hfill e^1_{\mathrm{CJS}}& =& \frac{1}{2\kappa ^2}R\frac{1}{4\kappa ^2}\left[9(_\mu \sigma )(^\mu \sigma )+\frac{1}{6}e^{6\sigma }G_{\mu \nu \rho 4}G^{\mu \nu \rho 4}\right]\hfill \\ \multicolumn{3}{c}{}\\ & & \frac{3}{4\kappa ^2}\left[(_\mu \gamma )(^\mu \gamma )+e^{2\gamma }(_\mu a)(^\mu a)\right].\hfill \end{array}$$
(2.17)
In this expression, $`\kappa `$ is the four-dimensional gravitational coupling with $`\kappa ^2=\kappa _{11}^2/V_7`$, $`V_7=V_1V_6`$ being the volume of the compact space $`S^1\times X_6`$.
At this stage, the identification of the bosonic components $`C`$, $`b_{\mu \nu }`$, $`C_T`$ and $`T_\mu `$ with the bulk fields $`\sigma `$, $`C_{\mu \nu 4}`$, $`\gamma `$ and $`a`$ can only be determined up to two proportionality constants (one for each ‘M-theory multiplet’ $`V`$ and $`V_T`$). These constants can however be determined from the couplings of $`C`$ and $`C_T`$ to charged matter and gauge fields . The result is
$$\begin{array}{cccccc}\hfill 4\kappa ^2C& =& \frac{\lambda ^2}{V_6}e^{3\sigma },\hfill & \hfill 4\kappa ^2b_{\mu \nu }& =& \frac{\lambda ^2}{V_6}C_{\mu \nu 4},\hfill \\ \multicolumn{6}{c}{}\\ \hfill C_T& =& 2\frac{\lambda ^2}{V_6}e^\gamma ,\hfill & \hfill T_\mu & =& 2\frac{\lambda ^2}{V_6}_\mu a.\hfill \end{array}$$
(2.18)
The quantity $`\lambda `$ is the gauge coupling constant on the $`𝐙_2`$ fixed planes. The dimensionless number $`\lambda ^2/V_6`$ actually never appears in the four-dimensional effective theory.
#### 2.2.3 Addition of a superpotential
The standard reduction of eleven-dimensional supergravity with unbroken $`N=1`$ supersymmetry does not generate a superpotential. This fact is however not a direct consequence of the eleven-dimensional Bianchi identity or of the CY and $`S^1/𝐙_2`$ symmetries. In principle, the Bianchi identity $`_{[M}G_{NPQR]}=0`$ allows a solution
$$G_{ijk4}=2i\kappa ^1hϵ_{ijk},G_{\overline{ijk}4}=2i\kappa ^1hϵ_{\overline{ijk}},$$
(2.19)
where $`h`$ is a real constant and $`ϵ_{ijk}`$ is the $`SU(3)`$-invariant CY tensor. The second term in the Lagrangian (2.4) generates then an extra contribution in the effective supergravity which corresponds to the addition of a superpotential term $`[ihS_0^3]_F`$ to the bulk Lagrangian. This contribution however breaks supersymmetry . Since we have insisted in writing Lagrangians in which all Bianchi identities are field equations, we prefer to consider
$$[U(W+\overline{W})]_D+[S_0^3W]_F.$$
(2.20)
In this way, the fact that the chiral multiplet $`W`$ ($`w=n=0`$) is an arbitrary imaginary constant is imposed by the field equation of the vector multiplet $`U`$ ($`w=2`$, $`n=0`$).
With the addition of a superpotential, the bulk Lagrangian takes its final ‘off-shell’ form
$$_\mathrm{B}=\left[\mathrm{{\rm Y}}(S+\overline{S})V+L_TV_T+U(W+\overline{W})\right]_D+[S_0^3W]_F,$$
(2.21)
in which the Bianchi identities of eleven-dimensional supergravity are translated into field equations of the Lagrange multipliers $`S`$, $`L_T`$ and $`U`$.
### 2.3 Modified Bianchi identities and $`\kappa `$-expansion
Compactification of M-theory on $`S^1/𝐙_2`$ is usually discussed in an expansion in powers of $`\kappa _{11}`$. Compactification on $`O_7`$ can similarly be formulated with $`\kappa `$ as expansion parameter. In the upstairs version, Bianchi identities are modified at the ten-dimensional planes fixed by $`S^1/𝐙_2`$. Suppose now that we modify the four-dimensional supersymmetric Bianchi identities of the bulk Lagrangian in the following way (we set $`h=0`$):
$$_\mathrm{B}=\left[\mathrm{{\rm Y}}(S+\overline{S})(V+\mathrm{\Delta }_V)+L_T(V_T+\mathrm{\Delta }_T)\right]_D,$$
(2.22)
with two composite vector multiplets $`\mathrm{\Delta }_V`$ ($`w=2`$, $`n=0`$) and $`\mathrm{\Delta }_T`$ ($`w=n=0`$). Solving for the Lagrange multipliers now leads to
$$V=L\mathrm{\Delta }_V,V_T=T+\overline{T}\mathrm{\Delta }_T.$$
The Lagrangian to first order in these modifications is then
$$=_\mathrm{B}\left[\frac{\mathrm{{\rm Y}}}{2V}\mathrm{\Delta }_V\frac{3}{2V_T}(\mathrm{{\rm Y}}\mathrm{\Delta }_T)\right]_D,$$
(2.23)
with $`V`$ and $`V_T`$ respectively replaced by $`L`$ and $`T+\overline{T}`$. The multiplets $`\mathrm{\Delta }_V`$ and $`\mathrm{{\rm Y}}\mathrm{\Delta }_T`$, with ‘canonical’ dimension $`w=2`$, appear at order $`\mathrm{{\rm Y}}^0\kappa ^0`$, in comparison with bulk terms of order $`\mathrm{{\rm Y}}\kappa ^2`$. This is the relation with the expansion in powers of $`\kappa _{11}`$ of M-theory in the low-energy limit. In M-theory compactification, the multiplets $`\mathrm{\Delta }_V`$ and $`\mathrm{\Delta }_T`$ can thus be obtained either by considering the modified Bianchi identities on $`O_7`$, formulated as in Eq. (2.22), or from corrections to the Lagrangian of eleven-dimensional supergravity on $`O_7`$, as in expression (2.23).
## 3 Gauge and matter contributions from the two $`𝐙_2`$ fixed planes
In this section, we show that the introduction of the next to lowest order corrections (gauge multiplets and charged matter contributions) is controlled by a simple modification of the four-dimensional Bianchi identities, in analogy with the appearance of $`𝐙_2`$ fixed planes contributions in the M-theory Bianchi identities.
We start by considering the well-known dependence on charged matter (in chiral multiplets collectively denoted by $`M`$, with $`w=n=0`$) and gauge multiplets (vector multiplet $`A`$, in the adjoint representation, with $`w=n=0`$) of the effective $`N=1`$ four-dimensional supergravity for CY compactifications of heterotic strings . The Lagrangian in the chiral formulation (2.10) becomes
$$_\mathrm{c}=\frac{3}{2}\left[S_0\overline{S}_0e^{K/3}\right]_D+\left[S_0^3W\right]_F+\frac{1}{4}\left[S𝒲𝒲\right]_F,$$
(3.1)
with
$$K=\mathrm{log}(S+\overline{S})3\mathrm{log}(T+\overline{T}2\overline{M}e^AM)$$
(3.2)
and $`W=\alpha M^3`$. The superpotential should be understood as a gauge invariant trilinear interaction with coupling constant $`\alpha `$ defined as an integral over the CY space. The chiral multiplet $`𝒲`$ ($`w=n=3/2`$) is the gauge field-strength for $`A`$. The gauge group is in general not simple, and
$$𝒲𝒲=\underset{a}{}c^a𝒲^a𝒲^a,$$
(3.3)
with a real coefficient $`c^a`$ for each simple or abelian factor. In the linear equivalent version of the theory, the Lagrangian (2.9) reads now
$$_\mathrm{l}=\left[(S_0\overline{S}_0e^{\widehat{K}/3})^{3/2}(2\widehat{L})^{1/2}\right]_D+[\alpha S_0^3M^3]_F,$$
(3.4)
where the new modulus and matter Kähler potential is
$$\widehat{K}=3\mathrm{log}(T+\overline{T}2\overline{M}e^AM).$$
(3.5)
The linear multiplet $`L`$ is replaced by
$$\widehat{L}=L2\mathrm{\Omega },$$
(3.6)
with the Chern-Simons vector multiplet $`\mathrm{\Omega }`$ ($`w=2`$, $`n=0`$) defined by<sup>5</sup><sup>5</sup>5 In global Poincaré supersymmetry, $`\mathrm{\Sigma }(\mathrm{\Omega })=\frac{1}{4}\overline{DD}\mathrm{\Omega }`$.
$$\mathrm{\Omega }=\underset{a}{}c^a\mathrm{\Omega }^a,\mathrm{\Sigma }(\mathrm{\Omega }^a)=\frac{1}{16}𝒲^a𝒲^a.$$
(3.7)
Insisting as before on Bianchi identities, both forms (3.1) and (3.4) are equivalent to
$$\begin{array}{ccc}\hfill & =& \left[\mathrm{{\rm Y}}(S+\overline{S})(V+2\mathrm{\Omega })+L_T(V_T+2\overline{M}e^AM)\right]_D\hfill \\ \multicolumn{3}{c}{}\\ & & +[U(W\alpha M^3)+\mathrm{c}.\mathrm{c}.]_D+\left[S_0^3W\right]_F\hfill \\ \multicolumn{3}{c}{}\\ & =& \left[\mathrm{{\rm Y}}(S+\overline{S})(V+2\mathrm{\Omega })+L_T(V_T+2\overline{M}e^AM)\right]_D\hfill \\ \multicolumn{3}{c}{}\\ & & +\left[S_0^3(ih+\alpha M^3)\right]_F.\hfill \end{array}$$
(3.8)
Supersymmetric vacua have $`h=0`$. As before, solving for $`S`$ and $`L_T`$ imposes respectively $`V=L2\mathrm{\Omega }=\widehat{L}`$ and $`V_T=T+\overline{T}2\overline{M}e^AM`$, leading to Eq. (3.4). Alternatively, with the tensor calculus identity (and up to an irrelevant total derivative)
$$2[(S+\overline{S})\mathrm{\Omega }]_D=\frac{1}{4}\underset{a}{}c^a[S𝒲^a𝒲^a]_F,$$
(3.9)
the resolution for $`V`$ and $`L_T`$ leads back to the chiral form (3.1).
This reformulation of the gauge invariant Lagrangian suggests some remarks. Firstly, it enhances the importance of gauge and matter Chern-Simons multiplets in superstring effective actions. Secondly, the Chern-Simons vector multiplet $`\mathrm{\Omega }(A)`$ is not gauge invariant: its variation is a linear multiplet. The variation of $`[(S+\overline{S})2\mathrm{\Omega }]_D`$ is then a derivative and $`V`$ remains gauge invariant. When solving for $`S`$, it simply follows that $`\widehat{L}`$ is gauge invariant and that the linear multiplet transforms as $`\delta L=2\delta \mathrm{\Omega }`$. Finally, expression (3.8) shows that all gauge and chiral matter contributions can be viewed as the supersymmetrization of modified Bianchi identities imposed by $`S`$, $`L_T`$ and $`U`$. This observation provides the link to the approach based on M-theory on $`O_7`$, in which the $`𝐙_2`$ fixed planes carrying the Yang-Mills fields induce because of supersymmetry modifications to the Bianchi identity of the four-form field strength of eleven-dimensional supergravity.
In the effective supergravity of M-theory on $`O_7`$ (‘upstairs formulation’), the various components of the Lagrangian (3.8) have the following origin. The first term is the bulk supergravity contribution. The second term, $`[(S+\overline{S})(V+2\mathrm{\Omega })]_D`$, is the supersymmetrization of the Bianchi identity verified by the component $`G_{\mu \nu \rho 4}`$ of the field $`G`$, modified by gauge contributions on the fixed planes. Similarly, the two last terms, $`[L_T(V_T+2\overline{M}e^AM)]_D`$ and $`[U(W\alpha M^3)+\mathrm{c}.\mathrm{c}.]_D`$, are respectively the supersymmetric extensions of the Bianchi identities of $`G_{\mu j\overline{k}4}`$ and $`G_{ijk4}`$. All the fixed plane contributions are then given at this order by the supersymmetrization of Bianchi identities, as obtained by direct $`O_7`$ truncation of the eleven-dimensional identities .
At this point, the gauge coupling constant for each simple or abelian factor $`a`$ in the gauge group appears to be
$$\frac{1}{g_a^2}=c^aRes.$$
(3.10)
At this order, $`g_a`$ is the tree-level wilsonnian and physical<sup>6</sup><sup>6</sup>6The coefficient of $`\frac{1}{4}F_{\mu \nu }^aF^{a\mu \nu }`$ in the generating functional of one-particle irreducible Green’s functions. gauge coupling.
It is clear, as already observed , that as far as the structure of the four-dimensional effective supergravity is concerned, the same information follows from $`O_7`$ compactification of M-theory at the next to lowest order in the $`\kappa `$-expansion and from CY compactifications of the heterotic strings, at zero string loop order.
## 4 Anomaly-cancelling terms
In the ten-dimensional heterotic string, cancellation of gauge and gravitational anomalies is a one-loop effect in string or effective supergravity perturbation theory. In four space-time dimensions, the nature of the cancelled anomalies is known from studies of $`(2,2)`$ compactifications of heterotic strings in the Yang-Mills sector : target-space duality of the modulus $`T`$ has a one-loop anomaly which is cancelled by a counterterm in the one-loop Wilson Lagrangian $`_W^{(1)}`$<sup>7</sup><sup>7</sup>7 The expressions given in the previous sections were for $`_W^{(0)}`$, or for the tree-level standard effective Lagrangian $`_\mathrm{\Gamma }`$., in a generalization to sigma-model anomalies of the Green-Schwarz mechanism . The derivation of the complete counterterm requires a calculation to all orders in the modulus $`T`$ . However, at the present stage of understanding, the M-theory approach should be regarded as a large-$`T`$ limit in which T-duality reduces to a shift symmetry in the imaginary part of $`T`$.
In the large-$`T`$ limit, the $`T`$-dependent corrections to gauge kinetic terms are of the form (see ref. and citations therein)
$$\frac{1}{4}\underset{a}{}\beta ^a\left[T𝒲^a𝒲^a\right]_F,$$
(4.1)
where the coefficients $`\beta ^a`$ are in principle calculable in heterotic strings. Taking also into account the $`D`$-density $`\left[L_T(V_T+2\overline{M}e^AM)\right]_D`$ present in Lagrangian (3.8), we can rewrite expression (4.1) in terms of the ‘M-theory multiplets’:
$$\left[(L_T2\underset{a}{}\beta ^a\mathrm{\Omega }^a)(V_T+2\overline{M}e^AM)\right]_D.$$
(4.2)
The correction (4.1) to the SYM Lagrangian is independent of the matter fields and can be seen as a correction to the holomorphic gauge kinetic function $`f_{ab}`$. A possible matter-dependent contribution to gauge kinetic terms is the gauge invariant real density<sup>8</sup><sup>8</sup>8 For simplicity, we consider the standard embedding with a gauge group $`E_6\times E_8`$, with the notation $`\mathrm{\Omega }=\mathrm{\Omega }^1+\mathrm{\Omega }^2`$, and with a matter multiplet $`M`$ transforming as (27,1) of $`E_6\times E_8`$. :
$$2\delta \left[\overline{M}e^AM(L2\underset{a=1}{\overset{2}{}}\mathrm{\Omega }^a)\right]_D,$$
(4.3)
or
$$2\delta \left[\overline{M}e^AMV\right]_D,$$
(4.4)
using the ‘M-theory multiplet’ $`V`$.
The M-theory anomaly-cancelling terms generate a further contribution of the form
$$ϵ\left[V|\alpha M^3|^2\right]_D.$$
(4.5)
In summary, the Wilson Lagrangian up to string one-loop order is expected to become
$$\begin{array}{ccc}\hfill & =& \left[\mathrm{{\rm Y}}(S+\overline{S})(V+2\mathrm{\Omega })+(L_T2_{a=1}^2\beta ^a\mathrm{\Omega }^a)(V_T+2\overline{M}e^AM)\right]_D\hfill \\ \multicolumn{3}{c}{}\\ & & +[U(W\alpha M^3)+\mathrm{c}.\mathrm{c}.]_D+\left[S_0^3W\right]_F\hfill \\ \multicolumn{3}{c}{}\\ & & +\left[V(ϵ|\alpha M^3|^22\delta \overline{M}e^AM)\right]_D.\hfill \end{array}$$
(4.6)
Each of the one-loop corrections, with coefficients $`\beta ^1`$, $`\beta ^2`$, $`ϵ`$ and $`\delta `$ is related to a well-defined counterterm which can be easily identified in the KK reduction of the ten-dimensional Green-Schwarz counterterms arising from M-theory on $`S^1/𝐙_2`$ . An explicit computation predicts in particular the relations $`\beta ^1=\beta ^2=\delta `$ .
From the general expression (4.6), we can derive various equivalent forms. For instance, solving for $`S`$, $`L_T`$ and $`U`$ gives the version of the effective supergravity in which the dilaton is described by a linear multiplet:
$$\begin{array}{ccc}\hfill _\mathrm{l}& =& \left[(S_0\overline{S}_0)^{3/2}\left(T+\overline{T}2\overline{M}e^AM\right)^{3/2}(2\widehat{L})^{1/2}\right]_D+\left[S_0^3(ih+\alpha M^3)\right]_F\hfill \\ \multicolumn{3}{c}{}\\ & & +\frac{1}{4}\left[T_{a=1}^2\beta ^a𝒲^a𝒲^a\right]_F+\left[\widehat{L}\left(ϵ|\alpha M^3|^22\delta \overline{M}e^AM\right)\right]_D.\hfill \end{array}$$
(4.7)
The threshold corrections are the holomorphic $`T`$-dependent terms controlled by $`\beta ^1`$ and $`\beta ^2`$.
We can also solve for $`V`$, $`L_T`$ and $`U`$ in Eq. (4.6) to get the version with a chiral dilaton multiplet:
$$_\mathrm{c}=\frac{3}{2}\left[S_0\overline{S}_0e^{K/3}\right]_D+\left[S_0^3(ih+\alpha M^3)\right]_F+\frac{1}{4}\underset{a=1}{\overset{2}{}}\left[(S+\beta ^aT)𝒲^a𝒲^a\right]_F,$$
(4.8)
with the Kähler potential
$$K=\mathrm{log}\left(S+\overline{S}+2\delta \overline{M}e^AMϵ|\alpha M^3|^2\right)3\mathrm{log}\left(T+\overline{T}2\overline{M}e^AM\right),$$
(4.9)
and the gauge kinetic functions $`f^a=S+\beta ^aT`$. The term with coefficient $`\delta `$ has been obtained in direct CY reductions of M-theory on $`S^1/𝐙_2`$ (see for instance ). The charged matter contribution with coefficient $`ϵ`$ was not included in these analyses. Observe however that an ambiguity exists because of the possibility to perform a holomorphic redefinition of the two chiral multiplets $`S`$ and $`T`$. To remove this ambiguity, we can use information from M-theory compactification , or choose the unequivocal linear version.
The gauge contributions appearing in Eq. (4.6) read
$$2\underset{a=1}{\overset{2}{}}\left[\left(S+\overline{S}+\beta ^a(V_T+2\overline{M}e^AM)\right)\mathrm{\Omega }^a\right]_D,$$
so that the gauge coupling constants are given by
$$\frac{1}{g_a^2}=Res+\frac{1}{2}\beta ^a\left(C_T+2\overline{M}M\right).$$
(4.10)
This expression becomes harmonic once the Bianchi identity imposing $`C_T+2\overline{M}M=2ReT`$ has been used.
Acknowledgments
This research has been supported in part by the European Union under the TMR contract ERBFMRX-CT96-0045, the Swiss National Science Foundation and the Swiss Office for Education and Science.
|
warning/0003/hep-th0003210.html
|
ar5iv
|
text
|
# Introduction
## Introduction
Although the Wess-Zumino-Novikov-Witten (WZNW) model was first formulated in terms of a (multivalued) action , it was originally solved by using axiomatic conformal field theory methods. The two dimensional (2D) Euclidean Green functions have been expressed as sums of products of analytic and antianalytic conformal blocks. Their operator interpretation exhibits some puzzling features: the presence of noninteger (”quantum”) statistical dimensions (that appear as positive real solutions of the fusion rules ) contrasted with the local (”Bose”) commutation relations (CR) of the corresponding 2D fields. The gradual understanding of both the factorization property and the hidden braid group statistics (signaled by the quantum dimensions) only begins with the development of the canonical approach to the model (for a sample of references, see ) and the associated splitting of the basic group valued field $`g:𝕊^1\times G`$ into chiral parts. The resulting zero mode extended phase space displays a new type of quantum group gauge symmetry: on one hand, it is expressed in terms of the quantum universal enveloping algebra $`U_q(𝒢),`$ a deformation of the finite dimensional Lie algebra $`𝒢`$ of $`G`$ – much like a gauge symmetry of the first kind; on the other, it requires the introduction of an extended, indefinite metric state space, a typical feature of a (local) gauge theory of the second kind.
Chiral fields admit an expansion into chiral vertex operators (CVO) which diagonalize the monodromy and are expressed in terms of the currents’ degrees of freedom with ”zero mode” coefficients that are independent of the world sheet coordinate . Such a type of quantum theory has been studied in the framework of lattice current algebras (see and references therein). Its accurate formulation in the continuum limit has only been attempted in the case of $`G=SU(2)`$ (see ). The identification (in ) of the zero mode ($`U_q`$)vertex operators $`a_\alpha ^i`$ (the ”$`U_q`$-oscillators” of the $`SU(2)`$ case ) with the generators of a quantum matrix algebra defined by a pair of (dynamical) $`R`$-matrices allows to extend this approach to the case of $`G=SU(n)`$.
The basic group valued chiral field $`u_\alpha ^A(x)`$ is thus expanded in CVO $`u_i^A(x,p)`$ which interpolate between chiral current algebra modules of weight $`p=p_jv^{(j)}`$ and $`p+v^{(i)},i=1,\mathrm{},n`$ (in the notation of to be recapitulated in Section 1 below). The operator valued coefficients $`a_\alpha ^i`$ of the resulting expansion intertwine finite dimensional irreducible representations (IR) of $`U_qU_q(sl_n)`$ that are labeled by the same weights. For generic $`q`$ ($`q`$ not a root of unity) they generate, acting on a suitably defined vacuum vector, a Fock-like space $``$ that contains every (finite dimensional) IR of $`U_q`$ with multiplicity $`1,`$ thus providing a model for $`U_q`$ in the sense of . This result (established in Section 3.1) appears to be novel even in the undeformed case ($`q=1`$) giving rise to a new (for $`n>2`$) model of $`SU(n).`$ In the important case of $`q`$ an even root of unity ($`q^h=1`$) we have prepared the ground (in Sections 3.2 and 3.3) for a (co)homological study of the two dimensional (left and right movers’) zero mode problem .
It should be emphasized that displaying the quantum group’s degrees of freedom requires an extension of the phase space of the models under consideration. Much interesting work on both physical and mathematical aspects of $`2D`$ conformal field theory has been performed without going to such an extension – see e.g. . The concept of a quantum group, on the other hand, has emerged in the study of closely related integrable systems and its uncovering in conformal field theory models has fascinated researchers from the outset – see e.g. . (For a historical survey of an early stage of this development see . Significant later developments in different directions – beyond the scope of the present paper – can be found e.g. in .)
Even within the scope of this paper there remain unresolved problems. We have, for instance, no operator realization of the extended chiral WZNW model, involving indecomposable highest weight modules of the Kac-Moody current algebra.
The paper is organized as follows. Section 1 provides an updated summary of recent work \- on the $`SU(n)`$ WZNW model. A new point here is the accurate treatment of the path dependence of the exchange relations in both the $`x`$ and the $`z=e^{ix}`$ pictures (Proposition 1.3). In Section 2 we carry out the factorization of the chiral field $`u(x)`$ into CVO and $`U_q`$ vertex operators and review relevant results of computing, in particular, the determinant of the quantum matrix $`a`$ as a function of the $`U_q(sl_n)`$ weights. The discussion of the interrelation between the braiding properties of four-point blocks and the exchange relations among zero modes presented in Section 2.2 is new; so are some technical results like Proposition 2.3 used in the sequel. Section 3.1 introduces the Fock space ($``$) representation of the zero mode algebra $`𝒜`$ for generic $`q`$; the main result is summed up in Proposition 3.3. In Section 3.2 we compute inner products for the canonical bases in the $`U_q`$ modules $`_p`$ for $`n=2,3.`$ In Section 3.3 we study the kernel of the inner product in $``$ for $`q`$ an even root of unity,
$$q=e^{i\frac{\pi }{h}}(h=k+nn).$$
(0.1)
It is presented in the form $`\stackrel{~}{}_h`$ where $`\stackrel{~}{}_h`$ is an ideal in $`𝒜.`$ We select a smaller ideal $`_h\stackrel{~}{}_h`$ (introduced in ) such that the factor algebra $`𝒜_h=𝒜/_h`$ is still finite-dimensional but contains along with each physical weight $`p`$ (with $`p_{1n}<h`$) a weight $`\stackrel{~}{p}`$ corresponding to the first singular vector of the associated Kac-Moody module (cf. Remark 2.1).
## 1 Monodromy extended $`SU(n)`$ WZNW model: a synopsis
### 1.1 Exchange relations; path dependent monomials of chiral fields
The WZNW action for a group valued field on a cylindric space-time $`^1\times 𝕊^1`$ is written as
$$S=\frac{k}{4\pi }\{\mathrm{tr}(g^1_+g)(g^1_{}g)dx^+dx^{}+s^{}\omega (g)\},x^\pm =x\pm t$$
(1.1a)
where $`s^{}\omega `$ is the pullback ($`s^{}g^1dg=g^1_+gdx^++g^1_{}gdx^{}`$) of a two-form $`\omega `$ on $`G`$ satisfying
$$d\omega (g)=\frac{1}{3}\mathrm{tr}(g^1dg)^3.$$
(1.1b)
The general, $`G=SU(n)`$ valued (periodic) solution, $`g(t,x+2\pi )=g(t,x),`$ of the resulting equations of motion factorizes into a product of group valued chiral fields
$$g_B^A(t,x)=u_\alpha ^A(x+t)(\overline{u}^1)_B^\alpha (xt)(\mathrm{classically},g,u,\overline{u}SU(n)),$$
(1.2a)
where $`u`$ and $`\overline{u}`$ satisfy a twisted periodicity condition
$$u(x+2\pi )=u(x)M,\overline{u}(x+2\pi )=\overline{u}(x)\overline{M}$$
(1.2b)
with equal monodromies, $`\overline{M}=M.`$ The symplectic form of the $`2D`$ model is expressed as a sum of two chiral $`2`$-forms involving the monodromy:
$$\mathrm{\Omega }^{(2)}=\mathrm{\Omega }(u,M)\mathrm{\Omega }(\overline{u},M),$$
(1.3)
$$\mathrm{\Omega }(u,M)=\frac{k}{4\pi }\left(\mathrm{tr}\left(\underset{\pi }{\overset{\pi }{}}(u^1du)u^1dudxb^1dbdMM^1\right)+\rho (M)\right).$$
Here $`b=u(\pi )`$ and the $`2`$-form $`\rho (M)`$ is restricted by the requirement that $`\mathrm{\Omega }(u,M)`$ is closed, $`d\mathrm{\Omega }(u,M)=\mathrm{\hspace{0.17em}0}`$ which is equivalent to
$$d\rho (M)=\frac{1}{3}\mathrm{tr}(dMM^1)^3$$
(1.4a)
(in other words, $`\rho `$ satisfies the same equation (1.1b) as $`\omega `$).
Such a $`\rho `$ can only be defined locally – in an open dense neighbourhood of the identity of the complexification of $`SU(n)`$ to $`SL(n,)`$. An example is given by
$$\rho (M)=\mathrm{tr}(M_+^1dM_+M_{}^1dM_{})$$
(1.4b)
where $`M_\pm `$ are the Gauss components of $`M`$ (which are well defined for $`M_{nn}0\mathrm{det}\left(\begin{array}{cc}M_{n1n1}& M_{n1n}\\ M_{nn1}& M_{nn}\end{array}\right)`$ etc.):
$$M=M_+M_{}^1,M_+=N_+D,M_{}^1=N_{}D,$$
(1.5)
$$N_+=\left(\begin{array}{cccc}1& f_1& f_{12}& \mathrm{}\\ 0& 1& f_2& \mathrm{}\\ 0& 0& 1& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\end{array}\right),N_{}=\left(\begin{array}{cccc}1& 0& 0& \mathrm{}\\ e_1& 1& 0& \mathrm{}\\ e_{21}& e_2& 1& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\end{array}\right),D=(d_\alpha \delta _\beta ^\alpha ),$$
and the common diagonal matrix $`D`$ has unit determinant: $`d_1d_2\mathrm{}d_n=1.`$ Different solutions $`\rho `$ of (1.4a) correspond to different non-degenerate solutions of the classical Yang-Baxter equation .
The closed $`2`$-form (1.3) on the space of chiral variables $`u,\overline{u},M`$ is degenerate. This fact is related to the non-uniqueness of the decomposition (1.2a): $`g(t,x)`$ does not change under constant right shifts of the chiral components, $`uuh,\overline{u}\overline{u}h`$, $`hG`$. Under such shifts the monodromy changes as $`Mh^1Mh`$ (see also the discussion of this point in ). We restore non-degeneracy by further extending the phase space, assuming that the monodromies $`M`$ and $`\overline{M}`$ of $`u`$ and $`\overline{u}`$ are independent so that the left and the right sector completely decouple. As a result monodromy invariance in the extended phase space is lost since $`M`$ and $`\overline{M}`$ satisfy Poisson bracket relations of opposite sign (due to (1.3)) and hence cannot be identified. Singlevaluedness of $`g(t,x)`$ can only be recovered in a weak sense, when $`g`$ is applied to a suitable subspace of ”physical states” in the quantum theory .
We require that quantization respects all symmetries of the classical chiral theory. Apart from conformal invariance and invariance under periodic left shifts the $`(u,M)`$ system admits a Poisson–Lie symmetry under constant right shifts which gives rise to a quantum group symmetry in the quantized theory. The quantum exchange relations so obtained ,
$$u_2(y)u_1(x)=u_1(x)u_2(y)R(xy),\overline{u}_1(x)\overline{u}_2(y)=\overline{u}_2(y)\overline{u}_1(x)R(xy)$$
(1.6)
(for $`0<|xy|<2\pi `$) can be also written as braid relations:
$`Pu_1(y)u_2(x)=u_1(x)u_2(y)\widehat{R}(xy)`$
$`\overline{u}_1^1(y)\overline{u}_2^1(x)P=\widehat{R}^1(xy)\overline{u}_1^1(x)\overline{u}_2^1(y)`$
$``$ $`\overline{u}_1(x)\overline{u}_2(y)P=\overline{u}_2(y)\overline{u}_1(x)\widehat{R}(xy).`$ (1.7b)
Here $`R(x)`$ is related to the (constant, Jimbo) $`SL(n)R`$-matrix by
$`R(x)`$ $`=`$ $`R\theta (x)+PR^1P\theta (x)`$ (1.8a)
$`R_{\beta _1\beta _2}^{\alpha _1\alpha _2}`$ $`=`$ $`\overline{q}^{\frac{1}{n}}\left(\delta _{\beta _1\beta _2}^{\alpha _1\alpha _2}q^{\delta _{\alpha _1\alpha _2}}+(q\overline{q})\delta _{\beta _2\beta _1}^{\alpha _1\alpha _2}\theta _{\alpha _1\alpha _2}\right)`$ (1.8b)
where
$$\delta _{\beta _1\beta _2}^{\alpha _1\alpha _2}=\delta _{\beta _1}^{\alpha _1}\delta _{\beta _2}^{\alpha _2},\theta _{\alpha \beta }=\{\begin{array}{ccc}1& \mathrm{if}& \alpha >\beta \\ 0& \mathrm{if}& \alpha \beta \end{array},\overline{q}:=q^1,$$
(1.8c)
$`P`$ stands for permutation of factors in $`VV,V=^n,`$ while $`\widehat{R}`$ is the corresponding braid operator:
$$\widehat{R}=RP,P\left(X_{|1𝒾}Y_{|2𝒾}\right)=X_{|2𝒾}Y_{|1𝒾}$$
(1.9a)
$$\widehat{R}(x)=R(x)P=\{\begin{array}{ccc}\widehat{R}& \mathrm{for}& x>0\\ \widehat{R}^1& \mathrm{for}& x<0\end{array}.$$
(1.9b)
We are using throughout the tensor product notation of Faddeev et al. : $`u_1=u1\mathrm{I},u_2=1\mathrm{I}u`$ are thus defined as operators in $`VV.`$
Restoring all indices we can write Eq.(1.7a) as
$$u_\alpha ^B(y)u_\beta ^A(x)=u_\sigma ^A(x)u_\tau ^B(y)\widehat{R}(xy)_{\alpha \beta }^{\sigma \tau }.$$
$`(1.7c)`$
Whenever dealing with a tensor product of $`3`$ or more copies of $`V`$ we shall write $`R_{ij}`$ to indicate that $`R`$ acts non-trivially on the $`i`$-th and $`j`$-th factors (and reduces to the identity operator on all others).
Remark 1.1 The operator $`\widehat{R}`$ (1.9a) coincides with $`\widehat{R}_{21}=P\widehat{R}_{12}P(P=P_{12})`$ in the notation of and . We note that if $`\widehat{R}_{ii+1}`$ satisfy the Artin braid relations then so do $`\widehat{R}_{i+1i}`$; we have, in particular,
$$\widehat{R}_{12}\widehat{R}_{23}\widehat{R}_{12}=\widehat{R}_{23}\widehat{R}_{12}\widehat{R}_{23}\widehat{R}_{32}\widehat{R}_{21}\widehat{R}_{32}=\widehat{R}_{21}\widehat{R}_{32}\widehat{R}_{21}.$$
(1.10)
Indeed, the two relations are obtained from one another by acting from left and right on both sides with the permutation operator $`P_{13}=P_{12}P_{23}P_{12}=P_{23}P_{12}P_{23}(=P_{31})`$ and taking into account the identities
$$P_{13}\widehat{R}_{12}P_{13}=\widehat{R}_{32},P_{13}\widehat{R}_{23}P_{13}=\widehat{R}_{21}.$$
(1.11)
Here we shall stick, following Refs. , to the form (1.7), (1.9) of the basic exchange relations. Note however that (1.1a) involves a change of sign in the WZ term (as compared to ) which yields the exchange of the $`x^+`$ and $`x^{}`$ factors in (1.2a) and is responsible for the sign change in the phase of $`q`$ (0.1).
The multivaluedness of chiral fields requires a more precise formulation of (1.7). To give an unambiguous meaning to such exchange relations we shall proceed as follows.
Energy positivity implies that for any $`l>0`$ the vector valued function
$$\mathrm{\Psi }(\zeta _1,\mathrm{},\zeta _l)=u_1(\zeta _1)\mathrm{}u_l(\zeta _l)|0𝒾$$
is (single valued) analytic on a simply connected open subset
$$\{\zeta _j=x_j+iy_j;|x_j|<\pi ,j=1,\mathrm{},l;y_j<y_{j+1},j=1,\mathrm{},l1\},$$
($`x_{jk}:=x_jx_k`$) of the manifold $`^l\mathrm{Diag}`$ where Diag is defined as the partial diagonal set in $`^l:\mathrm{Diag}=\{(\zeta _1,\mathrm{},\zeta _l),\zeta _j=\zeta _k\mathrm{for}\mathrm{some}jk\}.`$
Introduce (exploiting reparametrization invariance – cf. ) the analytic ($`z`$-) picture fundamental chiral field
$$\phi (z)=e^{i\mathrm{\Delta }\zeta }u(\zeta ),z=e^{i\zeta },\mathrm{\Delta }=\frac{n^21}{2hn},$$
(1.12)
$`\mathrm{\Delta }`$ standing for the conformal dimension of $`u,`$ and note that the variables $`z_j`$ are radially ordered in the domain $`𝒪_l`$:
$$𝒪_l=\{z_j=e^{y_j+ix_j};|z_j|>|z_{j+1}|,j=1,\mathrm{},l1;|\mathrm{arg}z_j|<\pi ,j=1,\mathrm{},l\}.$$
(1.13)
Remark 1.2 The time evolution law
$$e^{itL_0}u(x)e^{itL_0}=u(x+t)$$
(1.14a)
for the ”real compact picture” field $`u(x)`$ implies
$$e^{itL_0}\phi (z)e^{itL_0}=e^{it\mathrm{\Delta }}\phi (ze^{it}).$$
(1.14b)
Energy positivity, combined with the prefactor in (1.12), guarantees that the state vector $`\phi (z)|0𝒾`$ is a single valued analytic function of $`z`$ in the neighbourhood of the origin (in fact, for a suitably defined inner product, its Taylor expansion around $`z=0`$ is norm convergent for $`|z|<1`$ – see ).
The vector valued functions
$$\mathrm{\Phi }(z_1,\mathrm{},z_l)=\phi _1(z_1)\mathrm{}\phi _l(z_l)|0𝒾$$
(1.15a)
and
$$\mathrm{\Psi }(\zeta _1,\mathrm{},\zeta _l)=u_1(\zeta _1)\mathrm{}u_l(\zeta _l)|0𝒾=\underset{j}{}e^{i\mathrm{\Delta }\zeta _j}\mathrm{\Phi }(e^{i\zeta _1},\mathrm{},e^{i\zeta _l})$$
(1.15b)
are both analytic in their respective domains (cf. (1.13)) and are real analytic (and still single valued) on the parts
$$\{\zeta _j=x_j(z_j=e^{ix_j}),x_1>x_2>\mathrm{}>x_l,x_{1l}<\pi \}$$
of their physical boundaries.
The following Proposition allows to continue these boundary values through the domain $`𝒪_l`$ to any other ordered set of $`x_j`$ (the result will be a path dependent multivalued function for $`\{z_1,\mathrm{},z_l\}^l\mathrm{Diag}`$).
Proposition 1.3 Let $`z_1=e^{ix_1},z_2=e^{ix_2},0<x_{12}<2\pi `$; the path exchanging $`x_1`$ and $`x_2`$ (and hence, $`z_1`$ and $`z_2`$),
$$C_{12}:\zeta _{1,2}(t)=e^{i\frac{\pi }{2}t}(x_{1,2}\mathrm{cos}\frac{\pi }{2}t+ix_{2,1}\mathrm{sin}\frac{\pi }{2}t),0t1$$
(1.16a)
turns clockwise around the middle of the segment $`(x_1,x_2):`$
$$\zeta _1(t)+\zeta _2(t)=x_1+x_2,\zeta _{12}(t):=\zeta _1(t)\zeta _2(t)=x_{12}e^{i\pi t}.$$
(1.16b)
Furthermore, if $`z_a(t)=e^{i\zeta _a(t)},a=1,2,`$ then
$$|z_1(t)|^2=e^{x_{12}\mathrm{sin}\pi t}=|z_2(t)|^2>1\mathrm{for}0<t<1$$
(1.16c)
so that the pair $`(z_1(t),z_2(t))`$ satisfies the requirement (1.13) for two consecutive arguments in the analyticity domain $`𝒪_l.`$ For $`0<x_{21}<2\pi `$ one has to change the sign of $`t`$ (and thus the orientation of the path (1.16)) in order to preserve the inequality $`|z_1(t)|>|z_2(t)|.`$
Proof All assertions are verified by a direct computation; in particular, (1.16b) implies
$$2\mathrm{Im}\zeta _2(t)=x_{12}\mathrm{sin}\pi t=2\mathrm{Im}\zeta _1(t)$$
(1.16d)
which yields (1.16c). a
We note that for $`\zeta _{1,2}`$ given by (1.16a) one has $`|\zeta _1(t)|^2+|\zeta _2(t)|^2=x_1^2+x_2^2.`$ Proposition 1.3 supplements (1.12) in describing the relationship (the essential equivalence) between the real compact and the analytic picture allowing us to use each time the one better adapted to the problem under consideration.
We are now prepared to give an unambiguous formulation of the exchange relations (1.7).
Let $`\mathrm{\Pi }_{12}u_1(x_2)u_2(x_1)\left(\mathrm{\Pi }_{12}\phi _1(z_2)\phi _2(z_1)\right)`$ be the analytic continuation of $`u_1(x_1)u_2(x_2)`$ (respectively, $`\phi _1(z_1)\phi _2(z_2)`$ ) along a path in the homotopy class of $`C_{12}`$ (1.16). Then Eq.(1.7a) should be substituted by
$$P\mathrm{\Pi }_{12}u_1(x_2)u_2(x_1)=u_1(x_1)u_2(x_2)\widehat{R},$$
$`(1.7d)`$
$$P\mathrm{\Pi }_{12}\phi _1(z_2)\phi _2(z_1)=\phi _1(z_1)\phi _2(z_2)\widehat{R}$$
for $`z_j=e^{ix_j},\mathrm{\hspace{0.17em}0}<x_{12}<2\pi .`$ For $`0<x_{21}<2\pi `$ and a positively oriented path one should replace $`\widehat{R}`$ by $`\widehat{R}^1.`$
We recall (see ) that the quantized $`u`$ (and $`g`$) cannot be treated as group elements. We can just assert that the operator product expansion of $`u`$ with its conjugate only involves fields of the family (or, rather, the Verma module) of the unit operator. The relation
$$u(x+2\pi )=e^{2\pi iL_0}u(x)e^{2\pi iL_0}=u(x)M,$$
(1.17)
on the other hand, gives (by (1.14) for $`\mathrm{\Delta }`$ given by (1.12))
$$\left(M_\beta ^\alpha q^{\frac{1}{n}n}\delta _\beta ^\alpha \right)|0>=\mathrm{\hspace{0.17em}0};$$
(1.18a)
hence, in order to preserve the condition $`d_1\mathrm{}d_n=1`$ for the product of diagonal elements of $`M_+`$ and $`M_{}^1`$ we should substitute (1.5) by its quantum version
$$M=q^{\frac{1}{n}n}M_+M_{}^1.$$
(1.18b)
The tensor products of Gauss components, $`M_{2\pm }M_{1\pm }`$, of the monodromy matrix commute with the braid operator,
$`[\widehat{R},M_{2\pm }M_{1\pm }]=`$ $`0`$ $`=[\widehat{R},\overline{M}_{1\pm }\overline{M}_{2\pm }],`$ (1.19a)
(and hence, with its inverse) but
$`\widehat{R}M_2M_{1+}=M_{2+}M_1\widehat{R}`$ , $`\widehat{R}\overline{M}_{1+}\overline{M}_2=\overline{M}_1\overline{M}_{2+}\widehat{R}`$ (1.19b)
while the exchange relations between $`u`$ and $`M_\pm `$ can be written in the form (cf. \- )
$$M_{1\pm }Pu_1(x)=u_2(x)\widehat{R}^1M_{2\pm },\overline{M}_{2\pm }P\overline{u}_2(x)=\overline{u}_1(x)\widehat{R}^\pm \overline{M}_{1\pm }.$$
(1.20)
The left and right sectors decouple completely as a consequence of the separation of variables in the classical extended phase space,
$$[M_1,\overline{u}_2]=[u_1,\overline{u}_2]=[M_1,\overline{M}_2]=[u_1,\overline{M}_2]=0.$$
(1.21)
The above relations for the left sector variables $`(u,M)`$ are invariant under the left coaction of $`SL_q(n),`$
$$u_\alpha ^A(x)(T^1)_\alpha ^\beta u_\beta ^A(u^A(x)T^1)_\alpha ,M_\beta ^\alpha T_\gamma ^\alpha (T^1)_\beta ^\delta M_\delta ^\gamma (TMT^1)_\beta ^\alpha $$
(1.22a)
while the right sector is invariant under its right coaction,
$$\overline{u}_A^\alpha (x)\overline{u}_A^\beta (x)(\overline{T})_\beta ^\alpha =(\overline{T}\overline{u}_A(x))^\alpha ,\overline{M}_\beta ^\alpha \overline{M}_\delta ^\gamma \overline{T}_\gamma ^\alpha (\overline{T}^1)_\beta ^\delta (\overline{T}\overline{M}\overline{T}^1)_\beta ^\alpha $$
(1.22b)
provided
$$\widehat{R}T_2T_1=T_2T_1\widehat{R},\widehat{R}\overline{T}_1\overline{T}_2=\overline{T}_1\overline{T}_2\widehat{R}$$
(1.23)
where we have used concise notation in the right hand side of (1.22). The elements $`T^\alpha _\beta `$ of $`T`$ commute with $`u`$, $`M`$, $`\overline{u}`$ and $`\overline{M}`$. The fact that the maps (1.22a) and (1.22b) are respectively left and right coactions can be proven by checking the comodule axioms, see e.g. . There are corresponding transformations of the elements of $`M_\pm `$ and $`\overline{M}_\pm `$.
Thus the Latin and Greek indices of $`u`$ and $`\overline{u}`$ in (1.2a) transform differently: A, B correspond to the (undeformed) $`SU(n)`$ action while $`\alpha `$ is a quantum group index.
It is known, on the other hand, that the first equations in (1.19a) and (1.19b) for the matrices $`M_\pm `$ are equivalent to the defining relations of the (”simply connected” ) quantum universal enveloping algebra (QUEA) $`U_q(sl_n)`$ that is paired by duality to $`Fun(SL_q(n))`$ (see ). The Chevalley generators of $`U_q`$ are related to the elements $`d_i,e_i,f_i`$ of the matrices (1.5) by (; see also )
$`d_i=q^{\mathrm{\Lambda }_{i1}\mathrm{\Lambda }_i}(i=1,\mathrm{},n,\mathrm{\Lambda }_0=0=\mathrm{\Lambda }_n),`$
$`e_i=(\overline{q}q)E_i,f_i=(\overline{q}q)F_i`$ (1.24a)
$`(\overline{q}q)f_{12}=f_2f_1qf_1f_2=(\overline{q}q)^2(F_2F_1qF_1F_2)etc.,`$
$`(\overline{q}q)e_{21}=e_1e_2qe_2e_1=(\overline{q}q)^2(E_1E_2qE_2E_1)etc..`$ (1.24b)
Here $`\mathrm{\Lambda }_i`$ are the fundamental co-weights of $`sl(n)`$ (related to the co-roots $`H_i`$ by $`H_i=2\mathrm{\Lambda }_i\mathrm{\Lambda }_{i1}\mathrm{\Lambda }_{i+1}`$); $`E_i`$ and $`F_i`$ are the raising and lowering operators satisfying
$$[E_i,F_j]=[H_i]\delta _{ij}\left([H]:=\frac{q^H\overline{q}^H}{q\overline{q}}\right),$$
$$[E_i,E_j]=0=[F_i,F_j]for|ji|2,$$
(1.25a)
$$q^{\mathrm{\Lambda }_i}E_j=E_jq^{\mathrm{\Lambda }_i+\delta _{ij}},q^{\mathrm{\Lambda }_i}F_j=F_jq^{\mathrm{\Lambda }_i\delta _{ij}},$$
$$\left[2\right]X_iX_{i\pm 1}X_i=X_{i\pm 1}X_i^2+X_i^2X_{i\pm 1}forX=E,F.$$
(1.25b)
We note that the invariance under the coaction of $`SL_q(n)`$ (1.22a) is, in effect, equivalent to the covariance relations
$$\begin{array}{c}q^{H_i}u_\alpha (x)\overline{q}^{H_i}=q^{\delta _\alpha ^i\delta _\alpha ^{i+1}}u_\alpha (x),[E_i,u_\alpha ]=\delta _\alpha ^{i+1}u_{\alpha 1}(x)q^{H_i},\\ \\ F_iu_\alpha (x)q^{\delta _\alpha ^{i+1}\delta _\alpha ^i}u_\alpha (x)F_i=\delta _\alpha ^iu_{\alpha +1}(x).\end{array}$$
(1.26)
### 1.2 $`R`$-matrix realizations of the Hecke algebra; quantum antisymmetrizers
The $`R`$-matrix for the quantum deformation of any (simple) Lie algebra can be obtained as a representation of Drinfeld’s universal $`R`$-matrix . In the case of the defining representation of $`SU(n)`$ the braid operator (1.9) gives rise, in addition, to a representation of the Hecke algebra. This fact, exploited in , is important for our understanding of the dynamical $`R`$-matrix. We recall the basic definitions.
For any integer $`k2`$ let $`H_k(q)`$ be an associative algebra with generators $`1,g_1,\mathrm{},g_{k1}`$, depending on a non-zero complex parameter $`q`$, with defining relations
$`g_ig_{i+1}g_i=g_{i+1}g_ig_{i+1}`$ $`\mathrm{for}1ik2(\mathrm{if}k3),`$ (1.27a)
$`g_ig_j=g_jg_i`$ $`\mathrm{for}|ij|1,1i,jk1,`$ (1.27b)
$`g_i^2=1+(q\overline{q})g_i`$ $`\mathrm{for}1ik1,\overline{q}:=q^1.`$ (1.27c)
The $`SL(n)`$ braid operator $`\widehat{R}`$ (see (1.8b,c), (1.9a)) generates a representation $`\rho _n:H_k(q)End(V^k)`$, $`V=^n`$ for any $`k2,`$
$$\rho _n(g_i)=q^{\frac{1}{n}}\widehat{R}_{ii+1}\mathrm{or}[\rho _n(g_i)]^{\pm 1}=q^{\pm 1}1\mathrm{I}A_i,$$
(1.28a)
where $`A`$ is the $`q`$-antisymmetrizer
$$A_{\beta _1\beta _2}^{\alpha _1\alpha _2}=q^{ϵ_{\alpha _2\alpha _1}}\delta _{\beta _1\beta _2}^{\alpha _1\alpha _2}\delta _{\beta _2\beta _1}^{\alpha _1\alpha _2},q^{ϵ_{\alpha _2\alpha _1}}=\{\begin{array}{ccc}\overline{q}& \mathrm{for}& \alpha _1>\alpha _2\\ 1& \mathrm{for}& \alpha _1=\alpha _2\\ q& \mathrm{for}& \alpha _1<\alpha _2\end{array}$$
(1.28b)
($`A_i=[2]A^{(i+1,i)}`$ in the – suitably extended – notation of ; note that for $`q^2=1,[2]=0`$ the ”normalized antisymmetrizer” $`A^{(i+1,i)}`$ is ill defined while $`A_i`$ still makes sense). Eqs.(1.27) are equivalent to the following relations for the antisymmetrizers $`A_i`$:
$`A_iA_{i+1}A_iA_i=A_{i+1}A_iA_{i+1}A_{i+1}`$ (1.29a)
$`A_iA_j=A_jA_i\mathrm{for}|ij|1`$ (1.29b)
$`A_i^2=[2]A_i.`$ (1.29c)
Remark 1.3 We can define (see, e.g. ) the higher antisymmetrizers $`A_{ij},i<j`$ inductively, setting
$`A_{ij+1}:=A_{ij}(q^{ji+1}q^{ji}\rho _n(g_j)+\mathrm{}+(1)^{ji+1}\rho _n(g_jg_{j1}\mathrm{}g_i))=`$
$`=A_{i+1j+1}(q^{ji+1}q^{ji}\rho _n(g_i)+\mathrm{}+(1)^{ji+1}\rho _n(g_ig_{i+1}\mathrm{}g_j)).`$ (1.30a)
They can be also expressed in terms of antisymmetrizers only:
$`A_{ii+1}=A_i,A_{ij+1}={\displaystyle \frac{1}{[ji]!}}\left(A_{ij}A_jA_{ij}[ji][ji]!A_{ij}\right)=`$
$`={\displaystyle \frac{1}{[ji]!}}\left(A_{i+1j+1}A_iA_{i+1j+1}[ji][ji]!A_{i+1j+1}\right).`$ (1.30b)
The term ”$`q`$-antisymmetrizer” is justified by the relation
$$\left(\rho _n(g_i)+\overline{q}\right)A_{1j}=\mathrm{\hspace{0.17em}0}=A_{1j}\left(\rho _n(g_i)+\overline{q}\right)$$
(1.31a)
or
$$A_{1i}A_{1j}=[i]!A_{1j}\mathrm{for}1<ij.$$
(1.31b)
The dependence of the representation $`\rho _n`$ on $`n`$ (for $`G=SU(n)`$) is manifest in the relations
$$A_{1n+1}=0,\mathrm{rank}A_{1n}=1.$$
(1.32)
$`A_{1n}`$ can be written as a (tensor) product of two Levi-Civita tensors:
$$A_{1n}=^{|1\mathrm{}n𝒾}_{𝒽1\mathrm{}n|},_{𝒽1\mathrm{}n|}^{|1\mathrm{}n𝒾}=[n]!,$$
(1.33)
the second equation implying summation in all $`n`$ repeated indices. We can (and shall) choose the covariant and the contravariant $``$-tensors equal,
$$^{\alpha _1\alpha _2\mathrm{}\alpha _n}=_{\alpha _1\alpha _2\mathrm{}\alpha _n}=\overline{q}^{n(n1)/4}(q)^{\mathrm{}(\sigma )}\mathrm{for}\sigma =\left(\genfrac{}{}{0pt}{}{n_{},\mathrm{},1_{}}{\alpha _1,\mathrm{},\alpha _n}\right)\left(𝒮_n\right);$$
(1.34)
here $`\mathrm{}(\sigma )`$ is the length of the permutation $`\sigma `$. (Note the difference between (1.34) and the expression (2.5) of for $``$ which can be traced back to our present choice (1.28) for $`\rho _n(g_i)`$ – our $`\widehat{R}_{ii+1}`$ corresponding to $`\widehat{R}_{i+1i}`$ of – cf. Remark 1.1.)
We recall for further reference that the first equation (1.32) is equivalent to either of the following two relations:
$`A_{1n}A_{2n+1}A_{1n}`$ $`=`$ $`([n1]!)^2A_{1n}`$ (1.35a)
$`A_{2n+1}A_{1n}A_{2n+1}`$ $`=`$ $`([n1]!)^2A_{2n+1}`$ (1.35b)
(see Lemma 1.1 of ); this agrees with (1.33), (1.34) since
$`_{𝒽2\mathrm{}n+1|}^{|1\mathrm{}n𝒾}`$ $`=`$ $`(1)^{n1}[n1]!\delta _{𝒽n+1|}^{|1𝒾}`$ (1.36a)
$`_{𝒽1\mathrm{}n|}^{|2\mathrm{}n+1𝒾}`$ $`=`$ $`(1)^{n1}[n1]!\delta _{𝒽1|}^{|n+1𝒾}.`$ (1.36b)
We shall encounter in Section 2 below another, ”dynamical”, Hecke algebraic representation of the braid group which has the same form (1.28a) but with a ”dynamical antisymmetrizer”, i.e., $`A_i=A_i(p)`$, a (rational) function of the $`q`$-weights $`(q^{p_1},\mathrm{},q^{p_n})`$ which satisfies a finite difference (”dynamical”) version of (1.29a).
### 1.3 Barycentric basis, shifted $`su(n)`$ weights; conformal dimensions
Let $`\{v^{(i)},i=1,\mathrm{},n\}`$ be a symmetric ”barycentric basis” of (linearly dependent) real traceless diagonal matrices (thus $`\{v^{(j)}\}`$ span a real Cartan subalgebra $`𝔥sl(n)`$):
$$(v^{(i)})_k^j=\left(\delta _{ij}\frac{1}{n}\right)\delta _k^j\underset{i=1}{\overset{n}{}}v^{(i)}=0,(v^{(i)}|v^{(j)})=\delta _{ij}\frac{1}{n}.$$
(1.37)
(The inner product of two matrices is given by the trace of their product.) Analogously, the $`n`$ ”barycentric” components $`p_i`$ of a vector in the $`n1`$ dimensional dual space $`𝔥^{}`$ are determined up to a common additive constant and can be fixed by requiring $`_{i=1}^np_i=0.`$ Specifying thus the bases, we can make correspond to any such vector in $`𝔥^{}`$ a unique diagonal matrix $`p𝔥,`$
$$p=\underset{i=1}{\overset{n}{}}p_iv^{(i)}=\left(\begin{array}{cccc}p_1& 0& \mathrm{}& 0\\ 0& p_2& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& \mathrm{}& p_n\end{array}\right),\underset{i=1}{\overset{n}{}}p_i=0,$$
(1.38)
and vice versa. In particular, the simple $`sl(n)`$ roots $`\alpha _i`$ and the fundamental $`sl(n)`$ weights $`\mathrm{\Lambda }^{(j)},i,j=1,\mathrm{},n1,`$ satisfying $`(\mathrm{\Lambda }^{(j)}|\alpha _i)=\delta _i^j,`$ correspond to the following diagonal matrices (denoted by the same symbols),
$$\alpha _i=v^{(i)}v^{(i+1)},\mathrm{\Lambda }^{(j)}=\underset{\mathrm{}=1}{\overset{j}{}}v^{(\mathrm{})}(1\frac{j}{n})\underset{\mathrm{}=1}{\overset{j}{}}v^{(\mathrm{})}\frac{j}{n}\underset{\mathrm{}=j+1}{\overset{n}{}}v^{(\mathrm{})},$$
(1.39)
respectively. Expanding $`p`$ (1.38) in the basis of fundamental weights,
$$p=\underset{i=1}{\overset{n}{}}p_iv^{(i)}=\underset{j=1}{\overset{n1}{}}p_{jj+1}\mathrm{\Lambda }^{(j)},p_{ij}:=p_ip_j,$$
one can characterize a shifted dominant weight
$$p=\mathrm{\Lambda }+\rho ,\mathrm{\Lambda }=\underset{i=1}{\overset{n1}{}}\lambda _i\mathrm{\Lambda }^{(i)},\lambda _i_+,\rho =\underset{i=1}{\overset{n1}{}}\mathrm{\Lambda }^{(i)}=\frac{1}{2}\underset{\alpha >0}{}\alpha $$
(1.40)
($`\rho `$ is the $`sl(n)`$ Weyl vector) by the relations
$$p_{ii+1}=\lambda _i+1,i=1,2,\mathrm{},n1.$$
(1.41)
The non-negative integers $`\lambda _i=p_{ii+1}1`$ count the number of columns of length $`i`$ in the Young tableau that corresponds to the IR of highest weight $`p`$ of $`SU(n)`$ – see, e.g., . Conversely, $`p_i`$ satisfying (1.38) can be expressed in terms of the integer valued differences $`p_{ij}`$ as $`p_i=\frac{1}{n}_{j=1}^np_{ij}.`$
Dominant weights $`p`$ also label highest weight representations of $`U_q.`$ For integer heights $`h(n)`$ and $`q`$ satisfying (0.1) these are (unitary) irreducible if $`(n1)p_{1n}h`$. The quantum dimension of such an IR is given by (see, e.g., )
$$d_q(p)=\underset{i=1}{\overset{n1}{}}\{\frac{1}{[i]!}\underset{j=i+1}{\overset{n}{}}[p_{ij}]\}(0\mathrm{for}p_{1n}=p_1p_nh).$$
(1.42)
For $`q1(h\mathrm{}),[m]m`$ we recover the usual (integral) dimension of the IR under consideration.
The chiral observable algebra of the $`SU(n)`$ WZNW model is generated by a local current $`j(x)su(n)`$ of height $`h`$. In contrast to gauge dependent charged fields like $`u(x)`$, it is periodic, $`j(x+2\pi )=j(x)`$. The quantum version of the classical field-current relation $`ij(x)=ku^{}(x)u^1(x)`$ is the operator Knizhnik-Zamolodchikov equation in which the level $`k`$ gets a quantum correction (equal to the dual Coxeter number $`n`$ of $`su(n)`$):
$$hu^{}(x)=i:j(x)u(x):,h=k+n.$$
(1.43)
Here the normal product is defined in terms of the current’s frequency parts
$$:ju:=j_{(+)}u+uj_{()},j_{(+)}(x)=\underset{\nu =1}{\overset{\mathrm{}}{}}J_\nu e^{i\nu x},j_{()}(x)=\underset{\nu =0}{\overset{\mathrm{}}{}}J_\nu e^{i\nu x}.$$
(1.44)
The canonical chiral stress energy tensor and the conformal energy $`L_0`$ are expressed in terms of $`j`$ and its modes by the Sugawara formula:
$$𝒯(x)=\frac{1}{2h}\mathrm{tr}:j^2:(x)L_0=_\pi ^\pi 𝒯(x)\frac{dx}{2\pi }=\frac{1}{2h}\mathrm{tr}\left(J_0^2+2\underset{\nu =1}{\overset{\mathrm{}}{}}J_\nu J_\nu \right).$$
(1.45)
Energy positivity implies that the state space of the chiral quantum WZNW theory is a direct sum of (height $`h`$) ground state modules $`_p`$ of the Kac-Moody algebra $`\widehat{su}(n)`$ each entering with a finite multiplicity:
$$=\underset{p}{}_p_p,\mathrm{dim}_p<\mathrm{}.$$
(1.46)
We are not fixing at this point the structure of the internal spaces $`_p.`$ In the simpler but unrealistic case of generic $`q`$ explored in Section 3.1 each $`_p`$ is an irreducible $`U_q`$ module and the direct sum $`_p_p`$ carries a Fock type representation of the intertwining quantum matrix algebra $`𝒜`$ introduced below. The irreducibility property fails, in general, for $`q`$ a root of unity (as discussed in Section 3.3). It is conceivable that in this (realistic) case the label $`p`$ should be substituted by the set of eigenvalues of the $`U_q`$ Casimir operators which are symmetric polynomials in $`q^{p_i}.`$ (In the case of $`U_q(sl_2)`$ the single Casimir invariant depends on $`q^p+\overline{q}^p,pp_{12},`$ which suggests that $`p`$ and $`2hp`$ should refer to the same internal space.)
Each $`_p`$ in the direct sum (1.46) is a positive energy graded vector space,
$$_p=_{\nu =0}^{\mathrm{}}_p^\nu ,\left(L_0\mathrm{\Delta }(p)\nu \right)_p^\nu =\mathrm{\hspace{0.17em}0},\mathrm{dim}_p^\nu <\mathrm{}.$$
(1.47)
It follows from here and from the current algebra and Virasoro CR
$$[J_\nu ,L_0]=\nu J_\nu ,[L_\nu ,L_0]=\nu L_\nu (\nu )$$
(1.48)
that $`J_\nu _p^0=0=L_\nu _p^0\mathrm{for}\nu =1,2,\mathrm{}.`$ Furthermore, $`_p^0`$ spans an IR of $`su(n)`$ of (shifted) highest weight $`p`$ and dimension $`d_1(p)`$ (the $`q1`$ limit of the quantum dimension (1.42)). The conformal dimension (or conformal weight) $`\mathrm{\Delta }(p)`$ is proportional to the ($`su(n)`$-) second order Casimir operator $`|p|^2|\rho |^2:`$
$$2h\mathrm{\Delta }(p)=|p|^2|\rho |^2=\frac{1}{n}\underset{1i<jn}{}p_{ij}^2\frac{n(n^21)}{12}.$$
(1.49)
Note that the conformal dimension $`\mathrm{\Delta }(p^{(0)})`$ of the trivial representation
$$p^{(0)}=\{p;p_{ii+1}=1,1in1\}$$
(1.50)
is zero. This follows from the identity
$`n|p^{(0)}|^2={\displaystyle \underset{i=1}{\overset{n1}{}}}{\displaystyle \underset{j=i+1}{\overset{n}{}}}(p_{ij}^{(0)})^2=`$
$`={\displaystyle \underset{i=1}{\overset{n1}{}}}{\displaystyle \underset{j=i+1}{\overset{n}{}}}(ji)^2={\displaystyle \underset{i=1}{\overset{n1}{}}}{\displaystyle \frac{ni}{6}}(2n2i+1)(ni+1)=`$
$`={\displaystyle \frac{n^2(n^21)}{12}}=n|\rho |^2|p^{(0)}|^2|\rho |^2=0.`$ (1.51)
The eigenvalues of the braid operator $`\widehat{R}`$ (1.9), (1.28) are expressed as products of exponents of conformal dimensions. Let indeed $`p^{(1)}`$ be the weight of the defining $`n`$-dimensional IR of $`su(n)`$:
$$p_{12}^{(1)}=\mathrm{\hspace{0.17em}2},p_{ii+1}^{(1)}=\mathrm{\hspace{0.17em}1}\mathrm{for}i2$$
(1.52)
while $`p^{(s)}`$ and $`p^{(a)}`$ be the weights of the symmetric and the antisymmetric squares of $`p^{(1)}`$, respectively,
$`p_{12}^{(a)}=\mathrm{\hspace{0.17em}1}(=p_{ii+1}^{(a)}\mathrm{for}i3),p_{23}^{(a)}=\mathrm{\hspace{0.17em}2}(\mathrm{for}n3)`$
$`p_{12}^{(s)}=\mathrm{\hspace{0.17em}3},p_{ii+1}^{(s)}=\mathrm{\hspace{0.17em}1}\mathrm{for}n1i2.`$ (1.53)
The corresponding conformal dimensions $`\mathrm{\Delta }_{}=\mathrm{\Delta }(p^{(1)}),\mathrm{\Delta }_a=\mathrm{\Delta }(p^{(a)})`$ and $`\mathrm{\Delta }_s=\mathrm{\Delta }(p^{(s)})`$ are computed from (1.49):
$`2h\mathrm{\Delta }`$ $`=`$ $`|p^{(1)}|^2|\rho |^2={\displaystyle \frac{n^21}{n}},`$
$`2h\mathrm{\Delta }_a`$ $`=`$ $`|p^{(a)}|^2|\rho |^2=\mathrm{\hspace{0.17em}2}{\displaystyle \frac{n+1}{n}}(n2),`$ (1.54)
$`2h\mathrm{\Delta }_s`$ $`=`$ $`|p^{(s)}|^2|\rho |^2=\mathrm{\hspace{0.17em}2}{\displaystyle \frac{n1}{n}}(n+2).`$
The two eigenvalues of $`\widehat{R}`$ (evaluated from the nonvanishing $`3`$-point functions that involve two fields $`u(x)`$ (or $`\phi (z)`$ – see (1.12)) of conformal weight $`\mathrm{\Delta },`$
$$e^{i\pi (2\mathrm{\Delta }\mathrm{\Delta }_s)}=q^{\frac{n1}{n}},e^{i\pi (2\mathrm{\Delta }\mathrm{\Delta }_a)}=\overline{q}^{\frac{n+1}{n}},$$
(1.55)
appear with multiplicities
$$d_s=\left(\genfrac{}{}{0pt}{}{n+1}{2}\right),d_a=\left(\genfrac{}{}{0pt}{}{n}{2}\right)(d_a+d_s=n^2),$$
(1.56)
respectively. The deformation parameter
$$q^{\frac{1}{n}}=e^{i\frac{\pi }{nh}},$$
(1.57)
computed from here, satisfies (0.1) as anticipated. For $`q\mathrm{\hspace{0.17em}1}`$ the eigenvalues (1.55) of $`\widehat{R}`$ go into the corresponding eigenvalues $`\pm 1`$ of the permutation matrix $`P`$; furthermore,
$$\mathrm{det}\widehat{R}=\mathrm{det}P=(1)^{d_a}=\{\begin{array}{ccc}1& \mathrm{for}& n=2,3\mathrm{mod}\mathrm{\hspace{0.17em}4}\\ 1& \mathrm{for}& n=0,1\mathrm{mod}\mathrm{\hspace{0.17em}4}\end{array}.$$
(1.58)
Eq.(1.55) illustrates the early observation (see, e.g., ) that the quantum group is determined by basic characteristics (critical exponents) of the underlying conformal field theory.
## 2 CVO and $`U_q`$-vertex operators: monodromy and braiding
### 2.1 Monodromy eigenvalues and $`_p`$ intertwiners
The labels $`p`$ of the two factors in each term of the expansion (1.46) have different nature. While $`_p`$ is a ground state current algebra module for which $`p`$ stands for the shifted weight $`p^{KM}`$ (such that $`(p_i^{KM}p_i)_p=0`$) of the ground state representation of the $`\widehat{su}(n)`$ current algebra of minimal conformal dimension (or energy) $`\mathrm{\Delta }(p),_p`$ is a $`U_q`$ module (the quantum group commuting with the currents). We introduce, accordingly, the field of rational functions of the commuting operators $`q^{\widehat{p}_i}`$ (giving rise to an abelian group and) such that
$$\underset{i=1}{\overset{n}{}}q^{\widehat{p}_i}=\mathrm{\hspace{0.17em}1},\left(q^{\widehat{p}_i}q^{p_i}\right)_p=\mathrm{\hspace{0.17em}0},[q^{\widehat{p}_i},j(x)]=\mathrm{\hspace{0.17em}0}$$
(2.1)
with $`p_{ij}`$ obeying the condition (1.41) for dominant weights.
We shall split the $`SU(n)\times U_q(sl_n)`$ covariant field $`u(x)=\left(u_\alpha ^A(x)\right)`$ into factors which intertwine separately different $`_p`$ and $`_p`$ spaces.
A CVO $`u_j(x,p)`$ (for $`pp^{KM}`$) is defined as an intertwining map between $`_p`$ and $`_{p+v^{(j)}}`$ (for each $`p`$ in the sum (1.46)). Noting that $`_p`$ is an eigenspace of $`e^{2\pi iL_0}`$,
$$\mathrm{Spec}L_0|__p\mathrm{\Delta }_h(p)+_+\{e^{2\pi iL_0}e^{2\pi i\mathrm{\Delta }_h(p)}\}_p=\mathrm{\hspace{0.17em}0},$$
(2.2)
we deduce that $`u_j(x,p)`$ is an eigenvector of the monodromy automorphism:
$$u_j(x+2\pi ,p)=e^{2\pi iL_0}u_j(x,p)e^{2\pi iL_0}=u_j(x,p)\mu _j(p)$$
(2.3a)
where, using (1.49) and the relation $`(p|v^{(j)})=p_j,`$ we find
$$\mu _j(p):=e^{2\pi i\{\mathrm{\Delta }_h(p+v^{(j)})\mathrm{\Delta }_h(p)\}}=q^{\frac{1}{n}2p_j1}.$$
(2.3b)
The monodromy matrix (1.5) is diagonalizable whenever its eigenvalues (2.3b) are all different. In particular, for the ”physical IR”, characterized by $`p_{1n}<h,M`$ is diagonalizable. The exceptional points are those $`p`$ for which there exists a pair of indices $`1i<jn`$ such that $`q^{2p_{ij}}=\mathrm{\hspace{0.17em}1}`$, since we have
$$\frac{\mu _i(p)}{\mu _j(p)}=q^{2p_{ij}}.$$
(2.3c)
According to (1.42) all such ”exceptional” $`_p`$ have zero quantum dimension ($`[p_{ij}]=0`$).
Remark 2.1 The simplest example of a non-diagonalizable $`M`$ appears for $`n=2,p(p_{12})=h`$ when $`\mu _1(h)=\overline{q}^{\frac{1}{2}}=\mu _2(h).`$ In fact any $`\widehat{su}(2)`$ module $`_h\mathrm{},0\mathrm{}h1`$ contains a singular (invariant) subspace isomorphic to $`_{h+\mathrm{}}`$ ; note that, for $`p=h\mathrm{},\stackrel{~}{p}=h+\mathrm{},`$
$$\mu _1(p)=q^{\mathrm{}\frac{1}{2}}=\mu _2(\stackrel{~}{p}),\mu _2(p)=\overline{q}^{\mathrm{}+\frac{1}{2}}=\mu _1(\stackrel{~}{p})$$
(2.3d)
(cf. (2.3b)). It turns out that, in general,
$$\mu _1(p)=\mu _n(\stackrel{~}{p}),\mu _n(p)=\mu _1(\stackrel{~}{p}),\mu _i(p)=\mu _i(\stackrel{~}{p}),i=2,3,\mathrm{},n1.$$
(2.3e)
Indeed, if $`|\mathrm{hwv}𝒾_p`$ is the highest weight vector in the minimal energy subspace $`_p^0`$ of the $`\widehat{su}(n)`$ module $`_p`$ and $`\theta =\alpha _1+\mathrm{}\alpha _{n1}`$ is the $`su(n)`$ highest root, then the corresponding singular vector can be written in the form
$`\left(E_1^\theta \right)^{hp_{1n}}|\mathrm{hwv}𝒾_p|\mathrm{hwv}𝒾_{\stackrel{~}{p}},`$
$`\stackrel{~}{p}_1=h+p_n,\stackrel{~}{p}_n=h+p_1,\stackrel{~}{p}_i=p_i,i=2,3,\mathrm{},n1,`$ (2.4)
$`\mathrm{\Delta }(\stackrel{~}{p})\mathrm{\Delta }(p)=hp_{1n}.`$
To prove (2.3e) one just has to insert $`\stackrel{~}{p}_i`$ into (2.3b). One concludes that the monodromy $`M`$ always has coinciding eigenvalues on $`_p_{\stackrel{~}{p}}`$ (suggesting the inclusion of this direct sum into a single indecomposable $`U_q`$ module – cf. ). The nondiagonalizability of the monodromy matrix in the extended state space may require a modification of the splitting (2.7) below for $`p_{1n}h.`$ The study of this question is, however, beyond the scope of the present paper.
The CVO $`u_j`$ only acts on the factor $`_p`$ of the tensor product $`_p_p`$ in (1.46); hence, it shifts the Kac-Moody operators $`p_i^{KM}`$ but not the quantum group ones:
$$[p_i^{KM},u_j]=(\delta _{ij}\frac{1}{n})u_j,[q^{\widehat{p}_i},u_j]=\mathrm{\hspace{0.17em}0}\left((q^{p_i^{KM}}q^{\widehat{p}_i})=0\right).$$
(2.5)
We shall skip from now on (as we did in Eqs.(2.3),(2.1)) the superscript $`KM`$ of the argument $`p`$ of the CVO $`u_j`$ as well as the hat on the quantum group operators $`q^{p_i}`$ since the distinction between the two labels $`p`$ should be obvious from the context (and does not matter when acting on the diagonal chiral state space (1.46)).
The intertwining quantum matrix algebra $`𝒜`$ is generated by $`q^{p_i}`$ and by the elements of the matrix $`a=\left(a_\alpha ^j\right),j,\alpha =1,\mathrm{},n`$ which shift $`p`$ (commuting with $`p^{KM}`$),
$$a_\alpha ^j:_p_{p+v^{(j)}},q^{p_i}a_\alpha ^j=a_\alpha ^jq^{p_i+\delta _i^j\frac{1}{n}}.$$
(2.6)
The $`SU(n)\times U_q(sl_n)`$ covariant field $`u_\alpha ^A(x)`$ is related to the CVO $`u_j^A(x,p)`$ by the so called vertex-IRF (interaction-round-a-face) transformation
$$u_\alpha ^A(x)=u_i^A(x,p)a_\alpha ^i.$$
(2.7)
According to (2.7), $`a_\alpha ^j`$ act on the second factor of (1.46) only and hence commute with the currents and with the Virasoro generators.
It will be proven in Section 3.1 that the Fock space representation of $`𝒜`$ in $`_p_p`$ for generic $`q`$ provides a model of $`U_q.`$ The exchange relations of $`a_\alpha ^i`$ (displayed in Section 2.3 below), combined with (2.7) and with the defining property $`a_\alpha ^j|0𝒾=0`$ for $`j>1`$ of the vacuum vector $`|0𝒾`$ (the unique normalized state in $`_{p^{(0)}}`$ for $`p^{(0)}`$ given by (1.50)) yield, in particular, the relation
$$a^j_p=0\mathrm{for}j>1\mathrm{and}p_{j1}=p_j+1.$$
(2.8)
The meaning of (2.6), (2.8) can be visualized as follows. To each finite dimensional representation of $`U_q`$ with (dominant) highest weight $`p`$ we associate, as usual, a Young tableau $`Y_{[\lambda _1,\mathrm{},\lambda _{n1}]}`$ with $`\lambda _i(=p_{ii+1}1)`$ columns of height $`i(=1,2,\mathrm{},n1).`$ Then $`a^j`$ adds a box to the $`j^{\mathrm{th}}`$ row of the Young tableau of $`p`$ (provided $`p_{j1}>p_j+1`$ for $`j=2,\mathrm{},n).`$ Here are some examples for $`n=4:`$
The exchange relations of $`a_\alpha ^i`$ with the Gauss components (1.5) of the monodromy are dictated by (1.20),
$$M_{1\pm }Pa_1=a_2\widehat{R}^1M_{2\pm }$$
(2.9)
and reflect, in view of (1.1), the $`U_q`$ covariance of $`a`$:
$`[E_a,a_\alpha ^i]=\delta _{a\alpha 1}a_{\alpha 1}^iq^{H_a},{}_{a}{}^{}=1,\mathrm{},n1,`$ (2.10a)
$`[q^{H_a}F_a,a_\alpha ^i]=\delta _{a\alpha }q^{H_a}a_{\alpha +1}^i,`$ (2.10b)
$`q^{H_a}a_\alpha ^i=a_\alpha ^iq^{H_a+\delta _{a\alpha }\delta _{a\alpha 1}}.`$ (2.10c)
The transformation law (2.10) expresses the coadjoint action of $`U_q`$. Comparing (1.2b), (2.3) and (2.7) we deduce that the zero mode matrix $`a`$ diagonalizes the monodromy (whenever the quantum dimension (1.42) does not vanish); setting
$$aM=M_pa$$
(2.11a)
we find (from the above analysis of Eq.(2.3)) the implication
$$d_q(p)\mathrm{\hspace{0.17em}0}\left(M_p\right)_j^i=\delta _j^i\mu _j(pv^{(j)}),\mu _j(pv^{(j)})=q^{2p_j+1\frac{1}{n}}.$$
(2.11b)
It follows from (2.11) that the subalgebra of $`𝒜`$ generated by the matrix elements of $`M`$ commutes with all $`q^{p_i}.`$ As recalled in (1.5), (1.1) and (1.24), the Gauss components of $`M`$ are expressed in terms of the $`U_q`$ generators. We can thus state that the centralizer of $`q^{p_i}`$ in $`𝒜`$ is compounded by $`U_q`$ and $`q^{p_i}.`$
### 2.2 Exchange relations among zero modes from braiding properties of $`4`$-point blocks
The exchange relations (1.1) for $`u`$ given by (2.7) can be translated into quadratic exchange relations for the ”$`U_q`$ vertex operators” $`a_\alpha ^i`$ provided that the CVO $`u(x,p)`$ satisfy standard braid relations: if $`0<xy<2\pi ,`$ then
$$\mathrm{\Pi }_{xy}u_i^B(y,p+v^{(j)})u_j^A(x,p)=u_k^A(x,p+v^{(l)})u_l^B(y,p)\widehat{R}(p)_{ij}^{kl};$$
(2.12)
if $`0<yx<2\pi ,`$ then $`\widehat{R}(p)`$ in the right hand side should be substituted by the inverse matrix (cf. (1.9b)). Indeed, consistency of (2.12) with (1.7d) on the diagonal state space $``$ (1.46) requires that
$$\widehat{R}(p)^{\pm 1}a_1a_2=a_1a_2\widehat{R}^{\pm 1},$$
(2.13)
where $`p`$ in $`\widehat{R}(p)`$ should be understood as an operator, see (2.1) and (2.6).
It has been proven in that Eq. (2.12) is in fact a consequence of the properties of the chiral 4-point function
$$\begin{array}{c}w_{p^{}p}^{(4)}=𝒽0|\varphi _p_{}^{}{}_{}{}^{}(z_1)\phi (z_2)\phi (z_3)\varphi _p(z_4)|0𝒾=\\ \\ =\underset{i,j}{}𝒮^{ij}(p)s_{ij}(z_1,\mathrm{},z_4;p)\delta _{p^{},p+v^{(i)}+v^{(j)}}\end{array}$$
(2.14)
(we assume that the vacuum vector is given by the tensor product of the vacuum vectors for the affine and quantum matrix algebras). Here $`\varphi _p(z)`$ and $`\varphi _p_{}^{}{}_{}{}^{}(z)`$ are general $`z`$-picture primary chiral fields of weights $`p`$ and $`p_{}^{}{}_{}{}^{},`$ respectively, where $`p^{}`$ is the weight conjugate to $`p,`$
$$pp^{}=\{p_i^{}=p_{n+1i}\}p_{ii+1}^{}=p_{nin+1i},$$
(2.15)
$`\phi (z)`$ is the ”step operator” (1.12) (of weight $`p^{(1)}`$, see (1.52), i.e., $`\phi (z)\varphi _{p^{(1)}}(z)`$), $`𝒮^{ij}(p)`$ is the zero mode correlator
$$𝒮^{ij}(p):=𝒽p+v^{(i)}+v^{(j)}|a^ia^j|p𝒾,$$
(2.16)
while $`s_{ij}`$ is the conformal block expressed in terms of a function of the cross ratio $`\eta `$:
$$\begin{array}{c}s_{ij}(z_1,z_2,z_3,z_4;p):=\\ \\ =𝒽0|\varphi _p_{}^{}{}_{}{}^{}^{p^{(0)}}(z_1,p^{})\phi _i(z_2,p+v^{(j)})\phi _j(z_3,p)\varphi _p^p(z_4,p^{(0)})|0𝒾=\\ \\ =D_{ij}(z_1,z_2,z_3,z_4;p)f_{ij}(\eta ,p).\end{array}$$
(2.17)
Here we use the standard notation
$$\varphi _p^{p_2}(z,p_1):_{p_1}\stackrel{\varphi _p}{}_{p_2}$$
for a CVO of weight $`p`$ (so that $`\phi _{\mathrm{}}(z,p)\varphi _{p^{(1)}}^{p+v^{(\mathrm{})}}(z,p)`$ is the $`z`$-picture counterpart of $`u_{\mathrm{}}(x,p)`$),
$`D_{ij}(z_1,z_2,z_3,z_4;p)=\left({\displaystyle \frac{z_{24}}{z_{12}z_{14}}}\right)^{\mathrm{\Delta }(p^{})}\left({\displaystyle \frac{z_{13}}{z_{14}z_{34}}}\right)^{\mathrm{\Delta }(p)}z_{23}^{2\mathrm{\Delta }}\eta ^{\mathrm{\Delta }_j\mathrm{\Delta }}(1\eta )^{\mathrm{\Delta }_a},`$
$`\eta ={\displaystyle \frac{z_{12}z_{34}}{z_{13}z_{24}}},p^{}=p+v^{(i)}+v^{(j)};`$
$`\mathrm{\Delta }(p)`$ is given by (1.49) and $`\mathrm{\Delta }=\mathrm{\Delta }(p^{(1)})=\frac{n^21}{2hn},\mathrm{\Delta }_j=\mathrm{\Delta }(p+v^{(j)}),\mathrm{\Delta }_a=\frac{(n+1)(n2)}{2hn}`$ ($`\mathrm{\Delta }_a`$ is the dimension (1.3) of the antisymmetric tensor representation of weight $`p^{(a)}`$ in (1.53)). We are omitting here both $`SU(n)`$ and $`SL_q(n)`$ indices: $`s_{ij}`$ (2.17) (and hence $`f_{ij}`$) is an $`SU(n)`$ invariant tensor in the tensor product of four IRs, while $`𝒮^{ij}(p)`$ (2.16) is an $`SL_q(n)`$ invariant tensor. Only terms for which both $`p+v^{(j)}`$ and $`p+v^{(i)}+v^{(j)}`$ are dominant weights contribute to the sum (2.14). The prefactor $`D_{ij}(z_1,z_2,z_3,z_4;p)`$ in the right hand side of (2.17) is fixed, up to a multiplicative function of $`\eta ,`$ by the Möbius invariance of $`f_{ij}.`$ The choice of the powers of $`\eta `$ and $`1\eta `$ corresponds to extracting the leading singularities (in both $`s`$\- and $`u`$-channels) so that $`f_{ij}(\eta ,p)`$ should be finite (and nonzero) at $`\eta =0`$ and $`\eta =1.`$
We shall sketch the proof of (2.12); the reader could find the details in (see also ).
The conformal block $`s_{ij}`$ (2.17) is determined as the $`SU(n)`$ invariant solution of the Knizhnik-Zamolodchikov equation
$$\left(h\frac{}{z_2}+\frac{𝒞_{12}}{z_{12}}\frac{𝒞_{23}}{z_{23}}\frac{𝒞_{24}}{z_{24}}\right)s_{ij}(z_1,z_2,z_3,z_4;p)=\mathrm{\hspace{0.17em}0}$$
(2.19a)
satisfying the above boundary conditions. Inserting the expression (2.17) for $`s_{ij},`$ one gets a system of ordinary differential equations for the Möbius invariant amplitudes $`f_{ij}`$:
$$\left(h\frac{d}{d\eta }\frac{\mathrm{\Omega }_{12}}{\eta }+\frac{\mathrm{\Omega }_{23}}{1\eta }\right)f_{ij}(\eta ;p)=\mathrm{\hspace{0.17em}0}.$$
(2.19b)
Here $`𝒞_{ab}=\stackrel{}{t_a}.\stackrel{}{t_b},1a<b4`$ is the Casimir invariant of the corresponding tensor product of IRs of $`SU(n);`$ in our case $`t_a,a=1,2,3,4`$ generate the IRs of weights $`p_{}^{}{}_{}{}^{},p^{(1)},p^{(1)}`$ and $`p,`$ respectively. The prefactor $`D_{ij}`$ being an $`SU(n)`$ scalar, $`SU(n)`$ invariance of $`s_{ij}`$ implies
$$\left(𝒞_{12}+𝒞_{23}+𝒞_{24}+\frac{n^21}{n}\right)f_{ij}=\mathrm{\hspace{0.17em}0},$$
(2.19c)
so that
$$\mathrm{\Omega }_{12}=𝒞_{12}+p_m+\delta _{ij}+\frac{n^2+n4}{2n},\mathrm{\Omega }_{23}=𝒞_{23}+\frac{n+1}{n},$$
(2.19d)
where $`m=\mathrm{min}(i,j).`$
Our objective is to study the braiding properties of the solution $`f_{ij}`$ of (2.19b) that is analytic in $`\eta `$ (and non-zero) around $`\eta =0.`$
It is important to observe that the space of invariant $`SU(n)`$ tensors is in the case at hand at most two dimensional; this allows us to find a convenient realization of the operators $`\mathrm{\Omega }_{12},\mathrm{\Omega }_{23}`$ . (In the $`n=2`$ case this can be done even for four general isospins, due to the simple rules for tensor multiplication in the $`SU(2)`$ representation ring.)
The existence of a solution of (2.19) is guaranteed whenever the quantum dimension (1.42) for each weight encountered in (2.14) is positive,
$$n1p_{1n},(p+v^{(j)})_{1n},p_{}^{}{}_{1n}{}^{}<h,p^{}p+v^{(i)}+v^{(j)}.$$
(2.20)
In fact, for fixed $`p`$ and $`p^{}`$ in (2.17) and $`ij`$ the $`2\times 2`$ matrix system Eq. (2.19b) gives rise to a hypergeometric equation. Assume, in addition, that $`p+v^{(i)}`$ is also a dominant weight. Then both $`s_{ij}`$ and $`s_{ji}`$ will satisfy Eq.(2.19a) and provide a basis of independent solutions of that equation (note that the sum in (2.14) reduces to two terms with permuted $`i`$ and $`j`$). More precisely, let $`P_{23}\mathrm{\Pi }_{23}s_{ij}(z_1,z_3,z_2,z_4;p)`$ be the analytic continuation of $`s_{ij}`$ along a path $`C_{23}`$ obtained from $`C_{12}`$ (1.16a) by the substitution $`12,\mathrm{\hspace{0.17em}2}3`$ (that is, $`C_{23}=\{z_a(t)=e^{i\zeta _a(t)},a=2,3;\zeta _2(t)+\zeta _3(t)=x_2+x_3,\zeta _{23}(t)=e^{i\pi t}x_{23},0t1\}`$ ) with permuted $`SU(n)`$ indices $`2`$ and $`3`$. It satisfies again Eq.(2.19a) and hence is a linear combination of $`s_{kl}(z_1,z_2,z_3,z_4;p)`$ with $`(k,l)=(i,j)`$ and $`(k,l)=(j,i)`$:
$$P_{23}\mathrm{\Pi }_{23}s_{ij}(z_1,z_3,z_2,z_4;p)=s_{kl}(z_1,z_2,z_3,z_4;p)\widehat{R}_{ij}^{kl}(p).$$
(2.21)
Here $`\widehat{R}(p)`$ satisfies the ice condition: its components $`\widehat{R}_{ij}^{kl}(p)`$ do not vanish only if the unordered pairs $`i,j`$ and $`k,l`$ coincide – i.e.,
$$\widehat{R}_{ij}^{kl}(p)=a^{kl}(p)\delta _j^k\delta _i^l+b^{kl}(p)\delta _i^k\delta _j^l.$$
(2.22a)
Eq.(2.21) is nothing but a matrix element version of (2.12); hence, it yields the exchange relation (2.13) for $`ij`$ ($`kl`$).
For $`i=j`$ the analytic continuation in the left hand side of (2.21) reduces to a multiplication by a phase factor. In this case the space of $`SU(n)`$ invariant tensors is $`1`$-dimensional (since the skewsymmetric invariant vanishes), and so is the space of $`U_q(s\mathrm{}_n)`$ invariants. The resulting equation for $`f_{ii}(\eta ;p)`$ is first order:
$$\left(h\frac{d}{d\eta }+\frac{2}{1\eta }\right)f_{ii}(\eta ;p)=0,$$
and is solved by $`f_{ii}(\eta ;p)=c_{ii}(p)(1\eta )^{\frac{2}{h}}.`$ Substituting
$$z_{23}e^{i\pi }z_{23}1\eta e^{i\pi }\frac{1\eta }{\eta },D_{ii}\overline{q}^{\frac{n+1}{n}}\eta ^{\frac{2}{h}}D_{ii},$$
we get
$$s_{ii}\stackrel{}{}q^{1\frac{1}{n}}s_{ii}.$$
Explicitly, the ($`4\times 4`$) $`(i,j)`$-block of $`\widehat{R}(p)`$ has the form
$$\widehat{R}^{(i,j)}(p_{ij})=\overline{q}^{\frac{1}{n}}\left(\begin{array}{cccc}q& 0& 0& 0\\ 0& \frac{q^{p_{ij}}}{[p_{ij}]}& \frac{[p_{ij}1]}{[p_{ij}]}\alpha (p_{ij})& 0\\ 0& \frac{[p_{ij}+1]}{[p_{ij}]}\alpha (p_{ij})& \frac{\overline{q}^{p_{ij}}}{[p_{ij}]}& 0\\ 0& 0& 0& q\end{array}\right),\alpha (p)\alpha (p)=1,$$
(2.22b)
i.e. (cf. (2.22a))
$$q^{\frac{1}{n}}a^{kl}(p_{kl})=\alpha (p_{kl})\frac{[p_{kl}1]}{[p_{kl}]},q^{\frac{1}{n}}b^{kl}(p_{kl})=\frac{q^{p_{kl}}}{[p_{kl}]}\mathrm{for}kl.$$
The arbitrariness reflected by $`\alpha (p)`$ is related to the freedom of choosing the normalization of the two independent solutions of the hypergeometric equation.
The matrix (2.22b) coincides with the one, obtained independently in by imposing consistency conditions on the intertwining quantum matrix algebra of $`SL(n)`$ type. We shall display the ensuing properties of $`\widehat{R}(p)`$ in the following subsection.
### 2.3 The intertwining quantum matrix algebra
Among the various points of view on the $`U_q(sl_2)`$ intertwiners (or ”$`U_q`$ vertex operators”) $`a_\alpha ^i`$ the one which yields an appropriate generalization to $`U_q(sl_n)`$ is the so called ”quantum $`6j`$ symbol” $`\widehat{R}(p)`$-matrix formulation of . The $`n^2\times n^2`$ matrix $`\widehat{R}(p)`$ satisfies the dynamical Yang-Baxter equation (DYBE) first studied in whose general solution obeying the ice condition was found in .
The associativity of triple tensor products of quantum matrices together with Eq.(1.10) for $`\widehat{R}`$ yields the DYBE for $`\widehat{R}(p)`$:
$$\widehat{R}_{12}(p)\widehat{R}_{23}(pv_1)\widehat{R}_{12}(p)=\widehat{R}_{23}(pv_1)\widehat{R}_{12}(p)\widehat{R}_{23}(pv_1)$$
(2.23)
where we use again the succinct notation of Faddeev et al. (cf. Section 1):
$$\left(\widehat{R}_{23}(pv_1)\right)_{j_1j_2j_3}^{i_1i_2i_3}=\delta _{j_1}^{i_1}\widehat{R}(pv^{(i_1)})_{j_2j_3}^{i_2i_3}.$$
(2.24)
In deriving (2.23) from (2.13) we use (2.6). (The DYBE (2.23) is only sufficient for the consistency of the quadratic matrix algebra relations (2.13); it would be also necessary if the matrix $`a`$ were invertible – i.e., if $`d_q(p)0`$.)
The property of the operators $`\widehat{R}_{ii+1}(p)`$ to generate a representation of the braid group is ensured by the additional requirement (reflecting (1.27b))
$$\widehat{R}_{12}(p+v_1+v_2)=\widehat{R}_{12}(p)\widehat{R}_{kl}^{ij}(p)a_\alpha ^ka_\beta ^l=a_\alpha ^ka_\beta ^l\widehat{R}_{kl}^{ij}(p).$$
(2.25)
The Hecke algebra condition (1.27c) for the rescaled matrices $`\rho _n(g_i)`$ (1.28a) also fits our analysis of braiding properties of conformal blocks displayed in the previous subsection.
It is not surprising that the direct inspection of the braiding properties of the conformal blocks, from one side, and the common solution of the DYBE, (2.25) and the Hecke algebra conditions , from the other, lead to the same result. The solution (2.22b) can be presented in a form similar to (1.28):
$$q^{\frac{1}{n}}\widehat{R}(p)=q\text{1I}A(p),A(p)_{kl}^{ij}=\frac{[p_{ij}1]}{[p_{ij}]}\left(\delta _k^i\delta _l^j\delta _l^i\delta _k^j\right).$$
(2.26)
It is straightforward to verify the relations (1.29) for $`A_i(p):=q\text{1I}_{ii+1}q^{\frac{1}{n}}\widehat{R}_{ii+1}(p)`$; in particular,
$$[p_{ij}1]+[p_{ij}+1]=[2][p_{ij}]A^2(p)=[2]A(p).$$
(2.27)
According to the general $`SL(n)`$-type dynamical $`R`$-matrix can be obtained from (2.26) by either an analog of Drinfeld’s twist (see Lemma 3.1 of ) or by a canonical transformation $`p_ip_i+c_i`$ where $`c_i`$ are constants (numbers) such that $`_{i=1}^nc_i=0`$. The interpretation of the eigenvalues $`p_i`$ of $`\widehat{p}_i`$ as (shifted) weights (of the corresponding representations of $`U_q`$) allows to dispose of the second freedom.
Inserting (2.26) into the exchange relations (2.13) allows to present the latter in the following explicit form:
$$[a_\alpha ^i,a_\alpha ^j]=0,a_\alpha ^ia_\beta ^i=q^{ϵ_{\alpha \beta }}a_\beta ^ia_\alpha ^i$$
(2.28)
$$[p_{ij}1]a_\alpha ^ja_\beta ^i=[p_{ij}]a_\beta ^ia_\alpha ^jq^{ϵ_{\beta \alpha }p_{ij}}a_\alpha ^ia_\beta ^j\mathrm{for}\alpha \beta \mathrm{and}ij,$$
(2.29)
where $`q^{ϵ_{\alpha \beta }}`$ is defined in (1.28b).
There is, finally, a relation of order $`n`$ for $`a_\alpha ^i`$, derived from the following basic property of the quantum determinant:
$$det(a)=\frac{1}{[n]!}\epsilon _{𝒽1\mathrm{}n|}a_1\mathrm{}a_n^{|1\mathrm{}n𝒾}\frac{1}{[n]!}\epsilon _{i_1\mathrm{}i_n}a_{\alpha _1}^{i_1}\mathrm{}a_{\alpha _n}^{i_n}^{\alpha _1\mathrm{}\alpha _n}$$
(2.30)
where $`^{\alpha _1\mathrm{}\alpha _n}`$ is given by (1.34) while $`\epsilon _{i_1\mathrm{}i_n}`$ is the dynamical Levi-Civita tensor with lower indices (which can be consistently chosen to be equal to the undeformed one , a convention which we assume throughout this paper), normalized by $`\epsilon _{n\mathrm{}1}=1`$. The ratio $`det(a)\left(_{i<j}[p_{ij}]\right)^1`$ belongs to the centre of the quantum matrix algebra $`𝒜=𝒜(\widehat{R}(p),\widehat{R})`$ (see Corollary 5.1 of Proposition 5.2 of ). It is, therefore, legitimate to normalize the quantum determinant setting
$$det(a)=\underset{i<j}{}[p_{ij}]𝒟(p).$$
(2.31)
It is proportional (with a positive $`p`$-independent factor) to the quantum dimension (1.42).
Remark 2.2 The results of this section are clearly applicable if the determinant $`𝒟(p)`$ does not vanish (i.e., either for generic $`q`$ or, if $`q`$ is given by (0.1), for $`p_{1n}<h`$). As noted in the introduction, the notion of a CVO and the splitting (2.7) may well require a modification if this condition is violated.
To sum up: the intertwining quantum matrix algebra $`𝒜`$ is generated by the $`n^2`$ elements $`a_\alpha ^i`$ and the field $`(q,q^{p_i})`$ of rational functions of the commuting variables $`q^{p_i}`$ whose product is $`1`$, subject to the exchange relations (2.6) and (2.13) and the determinant condition (2.31).
The centralizer of $`q^{p_i}`$ in $`𝒜`$ (i.e., the maximal subalgebra of $`𝒜`$ commuting with all $`q^{p_i}`$) is spanned by the QUEA $`U_q`$ over the field $`(q,q^{p_i})`$ and $`a_\alpha ^i`$ obey the $`U_q`$ covariance relations (2.10). The expressions for the $`U_q`$ generators in terms of $`n`$-linear products of $`a_\alpha ^i`$ are worked out for $`n=2`$ and $`n=3`$ in Appendix A.
We shall use in what follows the intertwining properties of the product $`a_1\mathrm{}a_n`$ (see Proposition 5.1 of ):
$$\epsilon _{𝒽1\mathrm{}n|}a_1\mathrm{}a_n=𝒟(p)_{𝒽1\mathrm{}n|}$$
(2.32a)
or, in components,
$$\epsilon _{i_1\mathrm{}i_n}a_{\alpha _1}^{i_1}\mathrm{}a_{\alpha _n}^{i_n}=𝒟(p)_{\alpha _1\mathrm{}\alpha _n};$$
(2.32b)
$$a_1\mathrm{}a_n^{|1\mathrm{}n𝒾}=\epsilon ^{|1\mathrm{}n𝒾}(p)𝒟(p).$$
(2.33)
Here $`\epsilon (p)`$ is the dynamical Levi-Civita tensor with upper indices given by
$$\epsilon ^{i_1\mathrm{}i_n}(p)=(1)^{\mathrm{}(\sigma )}\underset{1\mu <\nu n}{}\frac{[p_{i_\mu i_\nu }1]}{[p_{i_\mu i_\nu }]},$$
(2.34)
$`\mathrm{}(\sigma )`$ standing again for the length of the permutation $`\sigma =\left(\genfrac{}{}{0pt}{}{n_{},\mathrm{},1_{}}{i_1,\mathrm{},i_n}\right)`$.
Remark 2.3 Selfconsistency of (1.17) requires that $`\mathrm{det}(a)=\mathrm{det}(aM).`$ Indeed, the non-commutativity of $`q^{p_j}`$ and $`a^i,`$ see Eq.(2.6), exactly compensates the factors $`q^{1\frac{1}{n}}`$ when computing the determinant of $`aM`$ (cf. (2.11a), (2.11b)); we have
$$q^{2p_n1+\frac{1}{n}}a_{\alpha _1}^nq^{2p_{n1}1+\frac{1}{n}}a_{\alpha _2}^{n1}\mathrm{}q^{2p_11+\frac{1}{n}}a_{\alpha _n}^1=a_{\alpha _1}^na_{\alpha _2}^{n1}\mathrm{}a_{\alpha _n}^1$$
(2.35)
since
$$q^{\frac{2}{n}(1+2+\mathrm{}+n1)n+1}=\mathrm{\hspace{0.17em}1}.$$
(2.36)
An important consequence of the ice property (2.22a) (valid for both $`\widehat{R}`$ and $`\widehat{R}(p)`$) is the existence of subalgebras of $`𝒜`$ with similar properties.
Let
$$I=\{i_1,i_2,\mathrm{},i_m\},1i_1<i_2<\mathrm{}<i_mn$$
and
$$\mathrm{\Gamma }=\{\alpha _1,\alpha _2,\mathrm{},\alpha _m\},1\alpha _1<\alpha _2<\mathrm{}<\alpha _mn$$
be two ordered sets of $`m`$ integers ($`1mn`$). Let $`A_{1m}|_\mathrm{\Gamma }`$ be the restriction of the $`q`$-antisymmetrizer $`(A_{1m})_{\beta _1\beta _2\mathrm{}\beta _m}^{\alpha _1\alpha _2\mathrm{}\alpha _m},\alpha _k,\beta _k\{1,2,\mathrm{},n\}`$ (for its definition see (1.30)) to a subset of indices $`\alpha _k,\beta _k\mathrm{\Gamma }.`$ Then $`rankA_{1m}|_\mathrm{\Gamma }=1`$ and one can define the corresponding restricted Levi-Civita tensors satisfying
$$A_{1m}|_\mathrm{\Gamma }=|_{\mathrm{\Gamma }}^{}{}_{}{}^{|1\mathrm{}m𝒾}|_{\mathrm{\Gamma }}^{}{}_{𝒽1\mathrm{}m|}{}^{},|_{\mathrm{\Gamma }}^{}{}_{𝒽1\mathrm{}m|}{}^{}|_{\mathrm{\Gamma }}^{}{}_{}{}^{|1\mathrm{}m𝒾}=[m]!.$$
(2.37)
In the same way one defines restricted dynamical Levi-Civita tensors
$$\epsilon |_{I}^{}{}_{}{}^{|1\mathrm{}m𝒾}(p)\mathrm{and}\epsilon |_{I}^{}{}_{𝒽1\mathrm{}m|}{}^{}(p)$$
for the subset $`I\{1,2,\mathrm{},n\}`$ (the last one of these does not actually depend on $`p`$ and coincides with the classical Levi-Civita tensor).
Consider the subalgebra $`𝒜(I,\mathrm{\Gamma })𝒜`$ generated by $`(q,q^{p_{ij}})`$, $`i,jI`$ and the elements of the submatrix $`a|_{I,\mathrm{\Gamma }}:=a_{\alpha \mathrm{\Gamma }}^{iI}`$ of the quantum matrix $`a`$.
Proposition 2.4 The normalized minor
$$\mathrm{\Delta }_{I,\mathrm{\Gamma }}(a):=\frac{det(a|_{I,\mathrm{\Gamma }})}{𝒟_I(p)}:=\frac{1}{[m]!𝒟_I(p)}\epsilon |_{I}^{}{}_{𝒽1\mathrm{}m|}{}^{}(a_1a_2\mathrm{}a_m)|_{I,\mathrm{\Gamma }}|_{\mathrm{\Gamma }}^{}{}_{}{}^{|1\mathrm{}m𝒾},$$
(2.38)
where
$$𝒟_I(p):=\underset{i<j;i,jI}{}[p_{ij}]$$
(2.39)
belongs to the centre of $`𝒜(I,\mathrm{\Gamma })`$.
The statement follows from the observation that relations (2.32a2.33) and (2.34) are valid for the restricted quantities $`|_\mathrm{\Gamma }`$, $`\epsilon |_I`$, $`𝒟_I(p)`$ and $`a|_{I,\mathrm{\Gamma }}`$. a
Using restricted analogs of the relations (2.33), (2.34), we can derive alternative expressions for the normalized minors:
$$\mathrm{\Delta }_{I,\mathrm{\Gamma }}(a)=\frac{1}{𝒟_I^+(p)}a_{\alpha _1}^{i_1}a_{\alpha _2}^{i_2}\mathrm{}a_{\alpha _m}^{i_m}|_\mathrm{\Gamma }^{\alpha _1\mathrm{}\alpha _m},$$
(2.40)
where
$$𝒟_I^+(p):=\underset{i<j;i,jI}{}[p_{ij}+1],$$
(2.41)
the indices $`i_kI`$ are in descendant order, $`i_1>i_2>\mathrm{}>i_m,`$ and the indices $`\alpha _k\mathrm{\Gamma }`$ are summed over.
## 3 The Fock space representation of $`𝒜`$. The ideal $`_h`$ for $`q^h=1`$.
### 3.1 The Fock space $`(𝒜)`$ (the case of generic $`q`$).
The ”Fock space” representation of the quantum matrix algebra $`𝒜`$ was anticipated in Eq.(2.7) and the subsequent discussion of Young tableaux. We define $``$ and its dual $`^{}`$ as cyclic $`𝒜`$ modules with one dimensional $`U_q`$-invariant subspaces of multiples of (non-zero) bra and ket vacuum vectors $`𝒽0|`$ and $`|0𝒾`$ such that $`𝒽0|𝒜=^{},𝒜|0𝒾=`$ satisfying
$`a_\alpha ^i|0𝒾=0\mathrm{for}i>1,𝒽0|a_\alpha ^j=0\mathrm{for}j<n,`$ (3.1a)
$`q^{p_{ij}}|0𝒾=q^{ji}|0𝒾,𝒽0|q^{p_{ij}}=q^{ji}𝒽0|,`$ (3.1b)
$`(X\epsilon (X))|0𝒾=0=𝒽0|(X\epsilon (X))`$ (3.1c)
for any $`XU_q`$ (with $`\epsilon (X)`$ the counit). The duality between $``$ and $`^{}`$ is established by a bilinear pairing $`𝒽.|.𝒾`$ such that
$$𝒽0|0𝒾=1,𝒽\mathrm{\Phi }|A|\mathrm{\Psi }𝒾=𝒽\mathrm{\Psi }|A^{}|\mathrm{\Phi }𝒾$$
(3.2)
where $`AA^{}`$ is a linear antiinvolution (transposition) of $`𝒜`$ defined for generic $`q`$ by
$$𝒟_i(p)(a_\alpha ^i)^{}=\stackrel{~}{a}_i^\alpha :=\frac{1}{[n1]!}^{\alpha \alpha _1\mathrm{}\alpha _{n1}}\epsilon _{ii_1\mathrm{}i_{n1}}a_{\alpha _1}^{i_1}\mathrm{}a_{\alpha _{n1}}^{i_{n1}},(q^{p_i})^{}=q^{p_i},$$
(3.3)
$`𝒟_i(p)`$ standing for the product
$$𝒟_i(p)=\underset{j<l,jil}{}[p_{jl}]([𝒟_i(p),a_\alpha ^i]=0=[𝒟_i(p),\stackrel{~}{a}_i^\alpha ]).$$
(3.4)
We verify in Appendix B the involutivity property, $`A^{\prime \prime }=A,`$ of (3.3) for $`n=3.`$ Eq. (3.3) implies the following formulae for the transposed of the Chevalley generators of $`U_q`$:
$$E_i^{}=F_iq^{H_i1},F_i^{}=q^{1H_i}E_i,(q^{H_i})^{}=q^{H_i}.$$
(3.5)
The main result of this section is the proof of the statement that for generic $`q`$ ($`q`$ not a root of unity) $``$ is a model space for $`U_q`$: each finite dimensional IR of $`U_q`$ is encountered in $``$ with multiplicity one.
Lemma 3.1 For generic $`q`$ the space $``$ is spanned by antinormal ordered polynomials applied to the vacuum vector:
$$P_{m_{n1}}(a_\alpha ^{n1})\mathrm{}P_{m_1}(a_\alpha ^1)|0𝒾\mathrm{with}m_1m_2\mathrm{}m_{n1}.$$
(3.6)
Here $`P_{m_i}(a_\alpha ^i)`$ is a homogeneous polynomial of degree $`m_i`$ in $`a_1^i,\mathrm{},a_n^i.`$
Proof We shall first prove the weaker statement that $``$ is spanned by vectors of the type $`P_{m_n}(a_\alpha ^n)\mathrm{}P_{m_1}(a_\alpha ^1)|0𝒾`$ (without restrictions on the nonnegative integers $`m_1,\mathrm{},m_n`$). It follows from the exchange relations (2.29) for $`i>j`$ and from the observation that $`[p_{jl}+1]0`$ for generic $`q`$ and $`j<l`$ in view of (3.1b).
Next we note that if $`m_{j1}=0`$ but $`m_j>0`$ for some $`j>1,`$ the resulting vector vanishes. Indeed, we can use in this case repeatedly (2.29) for $`i<j1`$ to move an $`a_\alpha ^j`$ until it hits the vacuum giving zero according to (3.1a).
If all $`m_i>0,i=1,\mathrm{},n,`$ we move a factor $`a_{\alpha _i}^i`$ of each monomial to the right to get rid of an $`n`$-tuple of $`a_{\alpha _i}^i`$ since
$$a_{\alpha _n}^n\mathrm{}a_{\alpha _1}^1|0𝒾=[n1]!_{\alpha _n\mathrm{}\alpha _1}|0𝒾;$$
(3.7)
here we have used once more (3.1a), and also (2.32) and (3.1b). Repeating this procedure $`m_n`$ times, we obtain an expression of the type (3.6) (or zero, if $`m_n>\mathrm{min}(m_1,\mathrm{},m_{n1})`$ ).
To prove the inequalities $`m_im_{i+1}`$ we can reduce the problem (by the same procedure of moving whenever possible $`a_\alpha ^i`$ to the right) to the statement that any expression of the type $`a_{\beta _1}^{i+1}a_{\beta _2}^{i+1}a_{\alpha _i}^i\mathrm{}a_{\alpha _1}^1|0𝒾`$ vanishes. We shall display the argument for a special case proving that
$$a_\alpha ^3a_\beta ^3a_2^2a_1^1|0𝒾=0\mathrm{for}n3.$$
(3.8)
This is a simple consequence of (2.28), (2.29) and (3.1a) if either $`\alpha `$ or $`\beta `$ is $`1`$ or $`2`$. We can hence write, using (2.10b),
$$0=F_2a_2^3a_3^3a_2^2a_1^1|0𝒾=(a_3^3)^2a_2^2a_1^1|0𝒾+a_2^3a_3^3a_3^2a_1^1|0𝒾=(a_3^3)^2a_2^2a_1^1|0𝒾.$$
(3.9)
By repeated application of $`F_i`$ (with $`i3`$ for $`n4`$) exploiting the $`U_q`$ invariance of the vacuum (3.1c), we thus complete the proof of (3.8) and hence, of Lemma 3.1. a
Corollary It follows from Lemma 3.1 that the space $``$ splits into a direct sum of weight spaces $`_p`$ spanned by vectors of type (3.6) with fixed degrees of homogeneity $`m_1,\mathrm{},m_{n1}`$:
$$=_p_p,p_{ij}=m_im_j+ji(ji\mathrm{for}i<j),$$
(3.10)
each subspace $`_p`$ being characterized by (2.1).
In order to exhibit the $`U_q`$ properties of $`_p`$ we shall introduce the highest and lowest weight vectors (HWV and LWV)
$$|\lambda _1\mathrm{}\lambda _{n1}𝒾\mathrm{and}|\lambda _{n1}\mathrm{}\lambda _1𝒾,$$
obeying
$$(q^{H_i}q^{\lambda _i})|\lambda _1\mathrm{}\lambda _{n1}𝒾=0=(q^{H_i}q^{\lambda _{ni}})|\lambda _{n1}\mathrm{}\lambda _1𝒾$$
(3.11)
for $`\lambda _i=m_im_{i+1}=p_{ii+1}1,1in1.`$
Lemma 3.2 Each $`_p`$ contains a unique (up to normalization) HWV and a unique LWV satisfying (3.11). They can be written in either of the following three equivalent forms:
$`|\lambda _1\mathrm{}\lambda _{n1}𝒾=`$ (3.12)
$`=(\mathrm{\Delta }_{n\mathrm{1\hspace{0.33em}1}}^{n\mathrm{1\hspace{0.33em}1}})^{\lambda _{n1}}(\mathrm{\Delta }_{n\mathrm{2\hspace{0.33em}1}}^{n\mathrm{2\hspace{0.33em}1}})^{\lambda _{n2}}\mathrm{}(\mathrm{\Delta }_{\mathrm{2\hspace{0.33em}1}}^{\mathrm{2\hspace{0.33em}1}})^{\lambda _2}(a_1^1)^{\lambda _1}|0𝒾=`$
$`=(a_1^1)^{\lambda _1}(\mathrm{\Delta }_{\mathrm{2\hspace{0.33em}1}}^{\mathrm{2\hspace{0.33em}1}})^{\lambda _2}\mathrm{}(\mathrm{\Delta }_{n\mathrm{1\hspace{0.33em}1}}^{n\mathrm{1\hspace{0.33em}1}})^{\lambda _{n1}}|0𝒾`$
$`(a_{n1}^{n1})^{\lambda _{n1}}(a_{n2}^{n2})^{\lambda _{n2}+\lambda _{n1}}\mathrm{}(a_1^1)^{\lambda _1+\mathrm{}+\lambda _{n1}}|0𝒾,`$
$`|\lambda _{n1}\mathrm{}\lambda _1𝒾=`$ (3.13)
$`=(\mathrm{\Delta }_{n\mathrm{\hspace{0.33em}2}}^{n\mathrm{1\hspace{0.33em}1}})^{\lambda _{n1}}(\mathrm{\Delta }_{n\mathrm{\hspace{0.33em}3}}^{n\mathrm{2\hspace{0.33em}1}})^{\lambda _{n2}}\mathrm{}(\mathrm{\Delta }_{nn1}^{\mathrm{2\hspace{0.33em}1}})^{\lambda _2}(a_n^1)^{\lambda _1}|0𝒾=`$
$`=(a_n^1)^{\lambda _1}(\mathrm{\Delta }_{nn1}^{\mathrm{2\hspace{0.33em}1}})^{\lambda _2}\mathrm{}(\mathrm{\Delta }_{n\mathrm{\hspace{0.33em}2}}^{n\mathrm{1\hspace{0.33em}1}})^{\lambda _{n1}}|0𝒾`$
$`(a_2^{n1})^{\lambda _{n1}}(a_3^{n2})^{\lambda _{n2}+\lambda _{n1}}\mathrm{}(a_n^1)^{\lambda _1+\mathrm{}+\lambda _{n1}}|0𝒾.;`$
here $`\mathrm{\Delta }_{i\mathrm{\hspace{0.33em}1}}^{i\mathrm{\hspace{0.33em}1}}`$ and $`\mathrm{\Delta }_{nni+1}^{i\mathrm{\hspace{0.33em}1}}`$ are normalized minors of the type (2.40):
$$\mathrm{\Delta }_{i\mathrm{\hspace{0.33em}1}}^{i\mathrm{\hspace{0.33em}1}}=\mathrm{\Delta }_{I_i,\mathrm{\Gamma }_i}(a)=\frac{1}{𝒟_{I_i}^+(p)}a_{\alpha _1}^i\mathrm{}a_{\alpha _i}^1|_{\mathrm{\Gamma }_i}^{\alpha _1\mathrm{}\alpha _i}$$
(3.14)
for $`I_i:=\{1,2,\mathrm{},i\}=:\mathrm{\Gamma }_i,`$ and
$$\mathrm{\Delta }_{nni+1}^{i\mathrm{\hspace{0.33em}1}}\mathrm{\Delta }_{I_i,\mathrm{\Gamma }_n^i}(a)=\frac{1}{𝒟_{I_i}^+(p)}a_{\alpha _1}^i\mathrm{}a_{\alpha _i}^1|_{\mathrm{\Gamma }_n^i}^{\alpha _1\mathrm{}\alpha _i}$$
(3.15)
where $`\mathrm{\Gamma }_n^i:=\{ni+1,ni+2,\mathrm{},n\}.`$
Proof We shall prove the uniqueness of the HWV by reducing an arbitrary eigenvector of $`q^{H_i}`$ of eigenvalue $`q^{\lambda _i},1in1,`$ to the form of the second equation (3.12). To this end we again apply the argument in the proof of Lemma 3.1. Let $`k(n1)`$ be the maximal numeral for which $`\lambda _k>0.`$ By repeated application of the exchange relations (2.29) we can arrange each $`k`$-tuple $`a_{\alpha _1}^k\mathrm{}a_{\alpha _k}^1`$ to hit a vector $`|v𝒾`$ such that $`(p_{ii+1}1)|v𝒾=0\mathrm{for}i<k`$. (Observe that all vectors of the type $`|v_k𝒾=(\mathrm{\Delta }_{k\mathrm{\hspace{0.33em}1}}^{k\mathrm{\hspace{0.33em}1}})^{\lambda _k}\mathrm{}(\mathrm{\Delta }_{n\mathrm{1\hspace{0.33em}1}}^{n\mathrm{1\hspace{0.33em}1}})^{\lambda _{n1}}|0𝒾,`$ for various choices of the non-negative integers $`\lambda _k,\mathrm{},\lambda _{n1},`$ have this property.) Noting then that $`a_\alpha ^{i+1}|v𝒾=0`$ whenever $`(p_{ii+1}1)|v𝒾=0`$ and using once more Eq.(2.29) we find
$$(p_{ii+1}1)|v𝒾=0(a_\beta ^{i+1}a_\alpha ^i+q^{ϵ_{\alpha \beta }}a_\alpha ^{i+1}a_\beta ^i)|v𝒾=0$$
(3.16)
which implies that we can substitute the product $`a_\alpha ^{i+1}a_\beta ^i`$ (acting on such a vector) by its antisymmetrized expression:
$$a_\alpha ^{i+1}a_\beta ^i|v𝒾=\frac{1}{[2]}(q^{ϵ_{\beta \alpha }}a_\alpha ^{i+1}a_\beta ^ia_\beta ^{i+1}a_\alpha ^i)|v𝒾(\mathrm{for}(p_{ii+1}1)|v𝒾=0).$$
(3.17)
Such successive antisymmetrizations will give rise to the minor $`\mathrm{\Delta }_{k1}^{k1}`$ yielding eventually the second expression (3.12) for the HWV.
To complete the proof of Lemma 3.2, it remains to prove the first equalities in (3.12) and (3.13). The commutativity of all factors $`\mathrm{\Delta }_{i\mathrm{\hspace{0.33em}1}}^{i\mathrm{\hspace{0.33em}1}},1in1(\mathrm{\Delta }_{\mathrm{1\hspace{0.33em}1}}^{\mathrm{1\hspace{0.33em}1}}a_1^1)`$ follows from Proposition 2.4 which implies
$$[a_\alpha ^i,\mathrm{\Delta }_{k\mathrm{\hspace{0.33em}1}}^{k\mathrm{\hspace{0.33em}1}}]=0\mathrm{for}1\alpha ,ik.$$
(3.18)
In order to compute the proportionality factors between the second and the third expressions in (3.12), (3.13) one may use the general exchange relation
$$[p_{ij}1](a_\alpha ^j)^ma_\beta ^i=[p_{ij}+m1]a_\beta ^i(a_\alpha ^j)^mq^{ϵ_{\beta \alpha }(p_{ij}+m1)}[m](a_\alpha ^j)^{m1}a_\alpha ^ia_\beta ^j$$
(3.19)
(valid for $`ij`$ and $`\alpha \beta `$) which follows from (2.29).
Lemmas 3.1 and 3.2 yield the main result of this section.
Proposition 3.3 The space $``$ is (for generic $`q`$) a model space of $`U_q.`$ a
We proceed to defining the $`U_q`$ symmetry of a Young tableau $`Y`$. A $`U_q`$ tensor $`T_{\alpha _1\mathrm{}\alpha _s}`$ is said to be $`q`$-symmetric if for any pair of adjacent indices $`\alpha \beta `$ we have
$$T_{\mathrm{}\alpha ^{}\beta ^{}\mathrm{}}A_{\alpha \beta }^{\alpha ^{}\beta ^{}}=0T_{\mathrm{}\alpha \beta \mathrm{}}=q^{ϵ_{\alpha \beta }}T_{\mathrm{}\beta \alpha \mathrm{}}$$
(3.20)
where $`q^{ϵ_{\alpha \beta }}`$ is defined in (1.28b). A tensor $`F_{\alpha _1\mathrm{}\alpha _s}`$ is $`q`$-skewsymmetric if it is an eigenvector of the antisymmetrizer (1.28b):
$$F_{\mathrm{}\alpha ^{}\beta ^{}\mathrm{}}A_{\alpha \beta }^{\alpha ^{}\beta ^{}}=[2]F_{\mathrm{}\alpha \beta \mathrm{}}F_{\mathrm{}\alpha \beta \mathrm{}}=q^{ϵ_{\beta \alpha }}F_{\mathrm{}\beta \alpha \mathrm{}}.$$
(3.21)
A $`U_q`$ tensor of $`\lambda _1+2\lambda _2+\mathrm{}+(n1)\lambda _{n1}`$ indices has the $`q`$-symmetry of a Young tableau $`Y=Y_{[\lambda _1,\mathrm{},\lambda _{n1}]}`$ (where $`\lambda _i`$ stands for the number of columns of height $`i`$) if it is first $`q`$-symmetrized in the indices of each row and then $`q`$-antisymmetrized along the columns.
The $`q`$-symmetry of a tensor associated with a Young tableau allows to choose as independent components an ordered set of values of the indices $`\alpha ,\beta `$ that monotonically increase along rows and strictly increase down the columns (as in the undeformed case - see ). Counting such labeled tableaux of a fixed type $`Y`$ allows to reproduce the dimension $`d_1(p)`$ of the space $`_p.`$
### 3.2 Canonical basis. Inner product.
We shall introduce a canonical basis in the $`U_q`$ modules $`_p`$ in the simplest cases of $`n=2,3`$ preparing the ground for the computation of inner products in $`_p`$ for such low values of $`n`$.
We shall follow Lusztig for a general definition of a canonical basis. It is, to begin with, a basis of weight vectors, a property which determines it (up to normalization) for $`n=2.`$ We shall set in this case
$$|p,m𝒾=(a_1^1)^m(a_2^1)^{p1m}|0𝒾,0mp1(pp_{12}).$$
(3.22)
Introducing (following ) the operators
$$E^{[m]}=\frac{1}{[m]!}E^m,F^{[m]}=\frac{1}{[m]!}F^m$$
(3.23)
we can relate $`|p,m𝒾`$ to the HWV and LWV in $`_p`$:
$$F^{[p1m]}|p,p1𝒾=\left[\genfrac{}{}{0pt}{}{p1}{m}\right]|p,m𝒾=E^{[m]}|p,0𝒾.$$
(3.24)
The situation for $`n=3`$ can still be handled more or less explicitly. A basis in $`_p`$ is constructed in that case by applying Lusztig’s canonical basis in either of the two conjugate Hopf subalgebras of raising or lowering operators
$$X_1^{[m]}X_2^{[\mathrm{}]}X_1^{[k]}\mathrm{and}X_2^{[k]}X_1^{[\mathrm{}]}X_2^{[m]}\mathrm{for}X=E\mathrm{or}F,\mathrm{}k+m,$$
(3.25)
the $`U_q`$ Serre relations implying
$$X_1^{[m]}X_2^{[k+m]}X_1^{[k]}=X_2^{[k]}X_1^{[k+m]}X_2^{[m]},$$
(3.26)
to the lowest or to the highest weight vector, respectively,
$$E_1^{[m]}E_2^{[\mathrm{}]}E_1^{[k]}|\lambda _2\lambda _1𝒾,E_2^{[k]}E_1^{[\mathrm{}]}E_2^{[m]}|\lambda _2\lambda _1𝒾,$$
(3.27)
$$F_1^{[m]}F_2^{[\mathrm{}]}F_1^{[k]}|\lambda _1\lambda _2𝒾,F_2^{[k]}F_1^{[\mathrm{}]}F_2^{[m]}|\lambda _1\lambda _2𝒾,0k+m\mathrm{}\lambda _1+\lambda _2$$
(3.28)
where we are setting
$`|\lambda _1\lambda _2𝒾=(a_1^1)^{\lambda _1}(qa_{3}^{3}{}_{}{}^{})^{\lambda _2}|0𝒾,`$ (3.29)
$`|\lambda _2\lambda _1𝒾=(a_3^1)^{\lambda _1}(\overline{q}a_{1}^{3}{}_{}{}^{})^{\lambda _2}|0𝒾.`$ (3.30)
(These expressions differ by an overall power of $`q`$ from (3.12) and (3.13).)
Lemma 3.4 The action of $`F_i^{[m]}(E_i^{[m]}),m`$ on a HWV (LWV) is given by
$`F_1^{[m]}|\lambda _1\lambda _2𝒾`$ $`=`$ $`\left[{\displaystyle \genfrac{}{}{0pt}{}{\lambda _1}{m}}\right](a_1^1)^{\lambda _1m}(a_2^1)^m(qa_{3}^{3}{}_{}{}^{})^{\lambda _2}|0𝒾,`$ (3.31a)
$`F_2^{[m]}|\lambda _1\lambda _2𝒾`$ $`=`$ $`\left[{\displaystyle \genfrac{}{}{0pt}{}{\lambda _2}{m}}\right](a_1^1)^{\lambda _1}(qa_{3}^{3}{}_{}{}^{})^{\lambda _2m}(a_{2}^{3}{}_{}{}^{})^m|0𝒾,`$
$`E_1^{[m]}|\lambda _2\lambda _1𝒾`$ $`=`$ $`\left[{\displaystyle \genfrac{}{}{0pt}{}{\lambda _2}{m}}\right](a_3^1)^{\lambda _1}(a_{2}^{3}{}_{}{}^{})^m(\overline{q}a_{1}^{3}{}_{}{}^{})^{\lambda _2m}|0𝒾,`$ (3.31b)
$`E_2^{[m]}|\lambda _2\lambda _1𝒾`$ $`=`$ $`\left[{\displaystyle \genfrac{}{}{0pt}{}{\lambda _1}{m}}\right](a_2^1)^m(a_3^1)^{\lambda _1m}(\overline{q}a_{1}^{3}{}_{}{}^{})^{\lambda _2}|0𝒾.`$
The proof uses (2.10), (2.28) and the relations
$$a_{2}^{3}{}_{}{}^{}a_{3}^{3}{}_{}{}^{}=qa_{3}^{3}{}_{}{}^{}a_{2}^{3}{}_{}{}^{},a_{1}^{3}{}_{}{}^{}a_{2}^{3}{}_{}{}^{}=qa_{2}^{3}{}_{}{}^{}a_{1}^{3}{}_{}{}^{}$$
(3.32)
obtained by transposing the second equality in (2.28) for $`i=3`$.
We shall turn now to the computation of the $`U_q`$ invariant form.
Conjecture 3.5 The scalar square of the HWV (3.12) and the LWV (3.13) of $`U_q`$ is given by
$$𝒽\lambda _1\mathrm{}\lambda _{n1}|\lambda _1\mathrm{}\lambda _{n1}𝒾=\underset{i<j}{}[p_{ij}1]!=𝒽\lambda _{n1}\mathrm{}\lambda _1|\lambda _{n1}\mathrm{}\lambda _1𝒾.$$
(3.33)
For $`n=2`$ the result is a straightforward consequence of Eqs.(3.22) and (Appendix A. Monodromy matrix and identification of $`U_q(sl_n)`$ generators for $`n=2`$ and $`n=3`$) (of Appendix A). For $`n=3`$ Eq.(3.33) reads
$$𝒽\lambda _1\lambda _2|\lambda _1\lambda _2𝒾=[\lambda _1]![\lambda _2]![\lambda _1+\lambda _2+1]!=𝒽\lambda _2\lambda _1|\lambda _2\lambda _1𝒾$$
(3.34)
which is proven in Appendix C. We conjecture that the argument can be extended to prove (3.33) for any $`n2.`$
For $`n=2`$ we can also write the inner products of any two vectors of the canonical basis :
$$𝒽p,m|p^{},m^{}𝒾=\delta _{pp^{}}\delta _{mm^{}}\overline{q}^{m(p1m)}[m]![p1m]!.$$
(3.35)
### 3.3 The case of $`q`$ a root of unity. Subspace of zero norm vectors. Ideals in $`𝒜`$.
In order to extend our results to the study of a WZNW model, we have to describe the structure of the $`U_q`$ modules $`_p`$ for $`q`$ a root of unity, (0.1). Here $`_p`$ is, by definition, the space spanned by vectors of type (3.6) (albeit the proof of Lemma 3.1 does not apply to this case). We start by recalling the situation for $`n=2`$ (see ).
The relations
$$E|p,m𝒾=[pm1]|p,m+1𝒾,F|p,m𝒾=[m]|p,m1𝒾$$
(3.36)
show that for $`ph`$ the $`U_q`$ module $`_p`$ admits a single HWV and LWV and is, hence, irreducible. For $`p>h`$ the situation changes.
Proposition 3.6 For $`h<p<2h`$ and $`q`$ given by (0.1) the module $`_p`$ is indecomposable. It has two $`U_q(sl_2)`$ invariant subspaces with no invariant complement:
$`_{p,h}^+`$ $`=`$ $`\mathrm{Span}\{|p,m𝒾,hmp1\},`$
$`_{p,h}^{}`$ $`=`$ $`\mathrm{Span}\{|p,m𝒾,0mp1h\}.`$
It contains a second pair of singular vectors: the LWV $`|p,h𝒾`$ and the HWV $`|p,p1h𝒾.`$ The vector $`|p,ph𝒾`$ is cosingular, – i.e., it cannot be written in the form $`E|v𝒾`$ with $`v(p);`$ similarly, the vector $`|p,h1𝒾`$ cannot be presented as $`F|v𝒾.`$
The statement follows from (3.36) and from the fact that
$$F|p,ph𝒾=[ph]|p,ph1𝒾0E|p,h1𝒾$$
(3.38)
so that the invariant subspace $`_{p,h}^+_{p,h}^{}`$ indeed has no invariant complement in $`_p`$.
The factor space
$$\stackrel{~}{}_p=_p/(_{p,h}^+_{p,h}^{})(h<p<2h)$$
(3.39)
carries an IR of $`U_q(sl_2)`$ of weight $`\stackrel{~}{p}=2hp`$ (cf. (2.1)).
The inner product (3.35) vanishes for vectors of the form (3.22) with $`p>h`$ and either $`mh`$ or $`mp1h.`$ Writing similar conditions for the bra vectors we end up with the following proposition: all null vectors belong to the set $`_h|0𝒾`$ or $`𝒽0|_h`$ where $`_h`$ is the ideal generated by $`[hp],[hH],q^{hp}+q^{hH}`$, and by $`h`$-th powers of the $`a_\alpha ^i`$ or, equivalently, by the $`h`$-th powers of $`\stackrel{~}{a}_i^\alpha `$. The factor algebra $`𝒜_h=𝒜/_h`$ is spanned by monomials of the type
$$q^{\mu p}q^{\nu H}(a_1^1)^{m_1}(a_2^1)^{m_2}(a_1^2)^{n_1}(a_2^2)^{n_2},h<\mu h,0\nu <h,0m_i,n_i<h$$
(3.40)
and is, hence, (not more than) $`2h^6`$ dimensional.
The definition of the ideal $`_h`$ can be generalized for any $`n2`$ assuming that it includes the $`h`$-th powers of all minors of the quantum matrix $`(a_\alpha ^i)`$ (for $`n=3,`$ equivalently, the $`h`$-th powers of $`a_\alpha ^i`$ and $`\stackrel{~}{a}_i^\alpha `$). It follows from Eq.(3.19), taking into account the vanishing of $`[h],`$ and from (2.28) that
$$(a_\alpha ^i)^ha_\beta ^j+(1)^{\delta _{\alpha \beta }}a_\beta ^j(a_\alpha ^i)^h=\mathrm{\hspace{0.17em}0}(=[[p_{ij}],(a_\alpha ^i)^h]_+)$$
(3.41)
implying also
$$(a_\alpha ^i)^h\stackrel{~}{a}_j^\beta +(1)^{\delta _\alpha ^\beta }\stackrel{~}{a}_j^\beta (a_\alpha ^i)^h=\mathrm{\hspace{0.17em}0}forn=3.$$
(3.42)
Similar relations are obtained (by transposition of (3.41), (3.42)) for $`(\stackrel{~}{a}_i^\alpha )^h`$ thus proving that the ideal $`_h`$ is indeed nontrivial, $`_h𝒜.`$
One can analyze on the basis of Lemma 3.4 the structure of indecomposable $`U_q(sl_3)`$ modules for, say, $`h<p_{13}<3h,`$ thus extending the result of Proposition 3.6. For example, as a corollary of (3.31), for $`q`$ given by (0.1) (a $`2h`$-th root of $`1`$) a HWV (a LWV) is annihilated by $`F_i(E_i)`$ if $`\lambda _i=0\mathrm{mod}h(\lambda _{\overline{i}}=0\mathrm{mod}h)`$ where $`\lambda _{\overline{1}}=\lambda _2,\lambda _{\overline{2}}=\lambda _1`$. If, in particular, both $`\lambda _i`$ are multiples of $`h,`$ then the corresponding weight vector spans a one-dimensional IR of $`U_q(sl_3).`$
For $`n>2,`$ however, the subspace $`_h|0𝒾`$ does not exhaust the set of null vectors in $`.`$ Indeed, for $`n=3`$ it follows from (3.34) and from the non-degeneracy of the highest and the lowest weight eigenvalues of the Cartan generators that the HWV and the LWV are null vectors for $`p_{13}>h`$:
$$𝒽|\lambda _1\lambda _2𝒾=0=𝒽|\lambda _2\lambda _1𝒾for\lambda _1+\lambda _2+1=p_{13}1h.$$
(3.43)
(If the conjecture (3.33) is satisfied then the HWV and the LWV for any $`n`$ are null vectors for $`p_{1n}h+1`$.) Since the representation of highest weight $`(\lambda _1,\lambda _2)`$ is irreducible for $`\lambda _ih1`$ (cf. (3.31)), the subspace $`𝒩`$ of null vectors contains $`_p`$ for $`p_{12}=\lambda _1+1h,p_{23}=\lambda _2+1h,p_{13}=p_{12}+p_{23}>h:`$
$$𝒫_{\lambda _1\lambda _2}(a_\alpha ^1;a_{\beta }^{3}{}_{}{}^{})|0𝒾𝒩for\lambda _ih1,\lambda _1+\lambda _2h1$$
(3.44)
for $`𝒫_{\lambda _1\lambda _2}(\rho _1a_\alpha ^1;\rho _2a_{\beta }^{3}{}_{}{}^{})=\rho _1^{\lambda _1}\rho _2^{\lambda _2}𝒫_{\lambda _1\lambda _2}(a_\alpha ^1;a_{\beta }^{3}{}_{}{}^{}),`$ – i.e., for any homogeneous polynomial $`𝒫_{\lambda _1\lambda _2}`$ of degree $`\lambda _1`$ in the first three variables, $`a_\alpha ^1,`$ and of degree $`\lambda _2h\lambda _11`$ in $`a_{\beta }^{3}{}_{}{}^{}`$. It follows that $`𝒩`$ contains all $`U_q`$ modules $`_{\stackrel{~}{p}}`$ of weights (2.1) corresponding to the first Kac-Moody singular vector for $`p_{13}<h(\stackrel{~}{p}_{13}=2hp_{13}>h`$ – see Remark 2.1). Hence, the factor space $`/𝒩`$ would be too small to accommodate the gauge theory treatment of the zero mode counterpart of such singular vectors.
We can write the null space $`𝒩`$ in the form $`𝒩=\stackrel{~}{}_h|0𝒾`$ where $`\stackrel{~}{}_h𝒜`$ is the ideal containing all $`𝒫_{\lambda _1\lambda _2}`$ appearing in (3.44) and closed under transposition, which contains $`_h`$ as a proper subideal. (We note that the transposition (3.3) is ill defined for $`q`$ a root of unity whenever $`𝒟_i(p)`$ vanishes.) The above discussion induces us to define the factor algebra
$$𝒜_h=𝒜/_h$$
(3.45)
(rather than $`𝒜/\stackrel{~}{}_h`$) as the restricted zero-mode algebra for $`q`$ a root of unity. It is easily verified (following the pattern of the $`n=2`$ case) that $`𝒜_h`$ is again a finite dimensional algebra. Its Fock space $`^h`$ includes vectors of the form
$$(a_1^2)^{m_1}(a_2^2)^{m_2}(a_3^2)^{m_3}(a_1^1)^{n_1}(a_2^1)^{n_2}(a_3^1)^{n_3}|0𝒾$$
(3.46)
for $`m_i,n_i<h(_im_i_in_i)`$ thus allowing for weights
$$p_{13}=n_1+n_2+n_3+23h1.$$
(3.47)
This justifies the problem of studying indecomposable $`U_q(sl_3)`$ modules for $`p_{13}<3h`$.
To sum up, the intertwining quantum matrix algebra $`𝒜`$ introduced in is an appropriate tool for studying the WZNW chiral zero modes. Its Fock space representation provides the first known model of $`U_q`$ for generic $`q`$. For exceptional $`q`$ (satisfying (0.1)) it gives room – by the results of this section – to the ”physical $`U_q`$ modules” coupled to the integrable (height $`h`$) representations of the $`\widehat{su}(n)`$ Kac-Moody algebra. This is a prerequisite for a BRS treatment of the zero mode problem of the two dimensional WZNW model (carried out, for $`n=2,`$ in ).
## Acknowledgements
The final version of this paper has been completed during visits of O.V.O. at ESI, Vienna, of L.K.H. at ESI, Vienna and at ICTP and INFN, Trieste, and of I.T.T. at ESI, Vienna, at SISSA, Trieste and at IHES, Bures-sur-Yvette. The authors thank these institutions for hospitality. P.F. acknowledges the support of the Italian Ministry of Education, University and Research (MIUR). This work is supported in part by the Bulgarian National Council for Scientific Research under contract F-828, and by CNRS and RFBR grants PICS-608, RFBR 98-01-22033. The work of A.P.I. and P.N.P. is supported in part by RFBR grant No. 00-01-00299.
## Appendix A. Monodromy matrix and identification of $`U_q(sl_n)`$ generators for $`n=2`$ and $`n=3`$
Eq.(2.11) rewritten as
$$M=a^1M_pa\mathrm{or}M_\beta ^\alpha =\underset{i=1}{\overset{n}{}}(a^1)_i^\alpha a_\beta ^iq^{2p_i1+\frac{1}{n}},$$
(A.1)
together with the Gauss decomposition (1.5) of the monodromy allows to express by (1.1) the Chevalley generators of $`U_q`$ as well as the operators $`E_iE_{i+1}qE_{i+1}E_i,F_{i+1}F_iqF_iF_{i+1}`$ etc. as linear combinations of products $`a_{\alpha _1}^1\mathrm{}a_{\alpha _n}^n`$ (with coefficients that depend on $`q^{p_i}`$ and the Cartan elements $`q^{\pm H_i}`$). Indeed, in view of (2.30), (2.31), we can express the elements of the inverse quantum matrix in terms of the (noncommutative) algebraic complement $`\stackrel{~}{a}_i^\alpha `$ of $`a_\alpha ^i:`$
$$𝒟(p)(a^1)_i^\alpha =\stackrel{~}{a}_i^\alpha =\frac{(1)^{n1}}{[n1]!}\epsilon _{i_1\mathrm{}i_{n1}i}^{\alpha \alpha _1\mathrm{}\alpha _{n1}}a_{\alpha _1}^{i_1}\mathrm{}a_{\alpha _{n1}}^{i_{n1}}.$$
(A.2)
Eq.(A.2) is equivalent to (3.3) since for the constant $`\epsilon `$-tensor used here we have $`(1)^{n1}\epsilon _{i_1\mathrm{}i_{n1}i}=\epsilon _{ii_1\mathrm{}i_{n1}}.`$ Thus we can recast (2.30-33) and (A.1) in the form
$$\stackrel{~}{a}_i^\alpha a_\beta ^i=𝒟(p)\delta _\beta ^\alpha ,\underset{i=1}{\overset{n}{}}\stackrel{~}{a}_i^\alpha a_\beta ^iq^{2p_i1+\frac{1}{n}}=𝒟(p)M_\beta ^\alpha .$$
(A.3)
Using Eqs.(4.10-12) of we can also write
$$\frac{1}{𝒟(p)}a_\alpha ^i\stackrel{~}{a}_j^\alpha =N_j^i(p)=\delta _j^i\underset{k<i}{}\frac{[p_{ki}+1]}{[p_{ki}]}\underset{i<l}{}\frac{[p_{il}1]}{[p_{il}]}.$$
(A.4)
We can express the $`U_q`$ generators in terms of products $`\stackrel{~}{a}_i^\alpha a_\beta ^i`$ (no summation over $`i`$). To this end we use (1.5), (1.1) to write
$$M_\beta ^\alpha =q^{\frac{1}{n}n}\underset{\sigma =\mathrm{max}(\alpha ,\beta )}{\overset{n}{}}f_{\alpha \sigma 1}d_\sigma e_{\sigma 1\beta }d_\beta $$
(A.5)
with $`f_{\alpha \alpha }=f_\alpha ,e_{\alpha \alpha }=e_\alpha ;f_{\alpha \alpha 1}=1=e_{\alpha 1\alpha }`$ (see (1.1)). It is thus simpler to start the identification of the elements with $`M_\beta ^n`$ and $`M_n^\alpha .`$ Using (1.1), we find, in particular,
$$d_{n}^{}{}_{}{}^{2}=q^{2\mathrm{\Lambda }_{n1}}=\frac{1}{𝒟(p)}\underset{i=1}{\overset{n}{}}\stackrel{~}{a}_i^na_n^iq^{n12p_i}.$$
We shall spell out the full set of resulting relations for $`n=2`$ and $`n=3.`$
The general relation between Cartan generators and $`sl_n`$ weights
$$H_i\underset{j}{}c_{ij}\mathrm{\Lambda }_j=2\mathrm{\Lambda }_i\mathrm{\Lambda }_{i1}\mathrm{\Lambda }_{i+1}(\mathrm{\Lambda }_0=\mathrm{\Lambda }_n=0)$$
(A.6)
tells us, for $`n=2,`$ that $`2\mathrm{\Lambda }_1=H.`$ This allows (using (1.28b) and (1.1)) to write the relations (A.3) in the form
$$\stackrel{~}{a}_i^\alpha a_\beta ^i=[p]\delta _\beta ^\alpha ,$$
(A.7)
$$\overline{q}^p\stackrel{~}{a}_1^\alpha a_\beta ^1+q^p\stackrel{~}{a}_2^\alpha a_\beta ^2=\overline{q}[p](M_+M_{}^1)_\beta ^\alpha .$$
(A.8)
Inserting for $`M_+M_{}^1`$ (1.5), (1.1) and (1.24) we find for $`n=2`$
$`\overline{q}M_+M_{}^1=\overline{q}\left(\begin{array}{cc}\overline{q}^{\frac{H}{2}}& (\overline{q}q)Fq^{\frac{H}{2}}\\ 0& q^{\frac{H}{2}}\end{array}\right)\left(\begin{array}{cc}\overline{q}^{\frac{H}{2}}& 0\\ (\overline{q}q)E\overline{q}^{\frac{H}{2}}& q^{\frac{H}{2}}\end{array}\right)=`$
$`=\left(\begin{array}{cc}q^p+\overline{q}^pq^{H+1}& (\overline{q}q)E^{}\\ (\overline{q}q)E& q^{H1}\end{array}\right),E^{}=Fq^{H1}.`$ (A.9)
As a result we obtain
$`\overline{q}^p\stackrel{~}{a}_1^1a_1^1+q^p\stackrel{~}{a}_2^1a_1^2`$ $`=`$ $`[p](q^p+\overline{q}^pq^{H+1}),`$
$`\overline{q}^p\stackrel{~}{a}_1^2a_2^1+q^p\stackrel{~}{a}_2^2a_2^2`$ $`=`$ $`[p]q^{H1},`$ (A.10)
$`\overline{q}^p\stackrel{~}{a}_1^1a_2^1+q^p\stackrel{~}{a}_2^1a_2^2`$ $`=`$ $`[p](\overline{q}q)E^{},`$
$`\overline{q}^p\stackrel{~}{a}_1^2a_1^1+q^p\stackrel{~}{a}_2^2a_1^2`$ $`=`$ $`[p](\overline{q}q)E.`$
Together with (A.7) this gives $`8`$ equations for the $`8`$ products $`\stackrel{~}{a}_i^\alpha a_\beta ^i`$ which can be solved with the result
$`\stackrel{~}{a}_1^1a_1^1={\displaystyle \frac{q^{H+1}\overline{q}^p}{q\overline{q}}}(=qa_2^2\stackrel{~}{a}_2^2)`$ , $`\stackrel{~}{a}_2^1a_1^2={\displaystyle \frac{q^pq^{H+1}}{q\overline{q}}}(=qa_2^1\stackrel{~}{a}_1^2)`$
$`\stackrel{~}{a}_2^2a_2^2={\displaystyle \frac{q^{H1}\overline{q}^p}{q\overline{q}}}(=\overline{q}a_1^1\stackrel{~}{a}_1^1)`$ , $`\stackrel{~}{a}_1^2a_2^1={\displaystyle \frac{q^pq^{H1}}{q\overline{q}}}(=\overline{q}a_1^2\stackrel{~}{a}_2^1);`$ (A.11)
$`\stackrel{~}{a}_1^2a_1^1=E=\stackrel{~}{a}_2^2a_1^2(=a_1^1\stackrel{~}{a}_1^2)`$ , $`\stackrel{~}{a}_1^1a_2^1=E^{}=\stackrel{~}{a}_2^1a_2^2(=a_2^1\stackrel{~}{a}_1^1)`$
further implying
$`a_2^2\stackrel{~}{a}_2^2a_1^1\stackrel{~}{a}_1^1=`$ $`\overline{q}^p`$ $`=\overline{q}\stackrel{~}{a}_1^1a_1^1q\stackrel{~}{a}_2^2a_2^2,`$
$`a_1^2\stackrel{~}{a}_2^1a_2^1\stackrel{~}{a}_1^2=`$ $`q^p`$ $`=q\stackrel{~}{a}_1^2a_2^1\overline{q}\stackrel{~}{a}_2^1a_1^2,`$ (A.12)
$`\stackrel{~}{a}_1^1a_1^1\stackrel{~}{a}_2^2a_2^2=`$ $`q^H`$ $`=\stackrel{~}{a}_1^2a_2^1\stackrel{~}{a}_2^1a_1^2.`$
In deriving the relations including products of the type $`a_\alpha ^i\stackrel{~}{a}_i^\beta `$ (appearing in parentheses in (Appendix A. Monodromy matrix and identification of $`U_q(sl_n)`$ generators for $`n=2`$ and $`n=3`$)), we have used (A.2) and (2.28).
These relations agree with (2.28), (2.29) for $`\stackrel{~}{a}_i^\alpha `$ given by (A.2) which becomes
$$\stackrel{~}{a}_i^\alpha =^{\alpha \beta }\epsilon _{ij}a_\beta ^j,\mathrm{i}.\mathrm{e}.,\stackrel{~}{a}_1^1=q^{1/2}a_2^2,\stackrel{~}{a}_2^1=q^{1/2}a_2^1,\stackrel{~}{a}_1^2=\overline{q}^{1/2}a_1^2,\stackrel{~}{a}_2^2=\overline{q}^{1/2}a_1^1,$$
(A.13)
implying
$`\stackrel{~}{a}_1^1a_1^1=qa_2^2\stackrel{~}{a}_2^2,\stackrel{~}{a}_2^2a_2^2=\overline{q}a_1^1\stackrel{~}{a}_1^1,`$
(A.14)
$`\stackrel{~}{a}_1^2a_2^1=\overline{q}a_1^2\stackrel{~}{a}_2^1,\stackrel{~}{a}_2^1a_1^2=qa_2^1\stackrel{~}{a}_1^2.`$
In the case of $`n=3`$ we make (A.3) and (A.4) explicit by noting the identities
$$3p_1=p_{12}+p_{13},3p_2=p_{23}p_{12},3p_3=p_{13}p_{23},$$
(A.15)
$$𝒟(p)=𝒟(p_1,p_2,p_3)=[p_{12}][p_{23}][p_{13}],$$
(A.16)
$$𝒟(p)N_j^i(p)=$$
(A.17)
$$=\mathrm{diag}([p_{23}][p_{12}1][p_{13}1],[p_{13}][p_{12}+1][p_{23}1],[p_{12}][p_{13}+1][p_{23}+1])$$
$`(N_i^i(p)=[3])`$. We find, in particular,
$$𝒟(p)q^{2\mathrm{\Lambda }_22}=\stackrel{~}{a}_1^3a_3^1\overline{q}^{\frac{2}{3}(p_{12}+p_{13})}+\stackrel{~}{a}_2^3a_3^2q^{\frac{2}{3}(p_{12}p_{23})}+\stackrel{~}{a}_3^3a_3^3q^{\frac{2}{3}(p_{13}+p_{23})},$$
(A.18)
$$𝒟(p)(\overline{q}^21)q^{\mathrm{\Lambda }_1}E_2=\stackrel{~}{a}_1^3a_2^1\overline{q}^{\frac{2}{3}(p_{12}+p_{13})}+\stackrel{~}{a}_2^3a_2^2q^{\frac{2}{3}(p_{12}p_{23})}+\stackrel{~}{a}_3^3a_2^3q^{\frac{2}{3}(p_{13}+p_{23})},$$
etc.
## Appendix B. Transposition in $`𝒜`$ for $`n=3`$
The involutivity of the transposition (3.3) is easily verified for $`n=2`$. Here we shall verify it for $`n=3`$ which is indicative for the general case.
Proposition A.1 The (linear) antihomomorphism of $`𝒜`$ defined by (3.3) is involutive: $`a_{\alpha }^{i}{}_{}{}^{\prime \prime }=a_\alpha ^i.`$
Proof Starting with the relation (3.3) for
$$\stackrel{~}{a}_i^\alpha =\frac{1}{[2]}\epsilon _{ijk}^{\alpha \beta \gamma }a_\beta ^ja_\gamma ^k$$
(B.1)
we shall prove, say, for $`i=1,`$ that
$`[2][p_{23}]a_{\alpha }^{1}{}_{}{}^{\prime \prime }=_{\alpha \beta \gamma }(a_\beta ^3a_\gamma ^2a_\beta ^2a_\gamma ^3)^{}=`$ (B.2)
$`={\displaystyle \frac{1}{[2]}}_{\alpha \beta \gamma }^{\gamma \rho \sigma }\{{\displaystyle \frac{1}{[p_{13}]}}(a_\rho ^1a_\sigma ^3a_\rho ^3a_\sigma ^1){\displaystyle \frac{\stackrel{~}{a}_3^\beta }{[p_{12}]}}{\displaystyle \frac{1}{[p_{12}]}}(a_\rho ^2a_\sigma ^1a_\rho ^1a_\sigma ^2){\displaystyle \frac{\stackrel{~}{a}_2^\beta }{[p_{13}]}}\}.`$
Noting the relation between the contraction of two Levi-Civita tensors and the $`q`$-antisymmetrizer (1.28b),
$$_{\alpha \beta \gamma }^{\gamma \rho \sigma }=A_{\alpha \beta }^{\rho \sigma }=\overline{q}^{ϵ_{\alpha \beta }}\delta _{\alpha \beta }^{\rho \sigma }\delta _{\beta \alpha }^{\rho \sigma },$$
(B.3)
we can rewrite (B.2) as
$`[2]^2𝒟(p)a_{\alpha }^{1}{}_{}{}^{\prime \prime }={\displaystyle \frac{[p_{12}]}{[p_{12}1]}}\{\overline{q}^{ϵ_{\alpha \beta }}(a_\alpha ^1a_\beta ^3a_\alpha ^3a_\beta ^1)a_\beta ^1a_\alpha ^3+a_\beta ^3a_\alpha ^1)\}\stackrel{~}{a}^\beta _3+`$
$`+{\displaystyle \frac{[p_{13}]}{[p_{13}1]}}\{\overline{q}^{ϵ_{\alpha \beta }}(a_\alpha ^1a_\beta ^2a_\alpha ^2a_\beta ^1)a_\beta ^1a_\alpha ^2+a_\beta ^2a_\alpha ^1)\}\stackrel{~}{a}^\beta _2.`$ (B.4)
Applying four times Eq.(2.29) in the form
$`a_\beta ^1a_\alpha ^i={\displaystyle \frac{[p_{1i}1]}{[p_{1i}]}}a_\alpha ^ia_\beta ^1+{\displaystyle \frac{\overline{q}^{ϵ_{\alpha \beta }}}{[p_{1i}]}}a_\alpha ^1a_\beta ^i,`$
$`a_\beta ^ia_\alpha ^1={\displaystyle \frac{[p_{1i}+1]}{[p_{1i}]}}a_\alpha ^1a_\beta ^i{\displaystyle \frac{q^{ϵ_{\alpha \beta }}}{[p_{1i}]}}a_\alpha ^ia_\beta ^1,`$ (B.5)
for $`i=2,3,\alpha \beta ,`$ and using (A.4), (A.16) and the identities
$$\overline{q}^ϵ[p]+[p+1]\overline{q}^{ϵp}=[2][p]=\overline{q}^ϵ[p]+[p1]+q^{ϵp}$$
(B.6)
for $`ϵ=\pm 1`$, we find that (B.2) is equivalent to
$`[2]𝒟(p)a_{\alpha }^{1}{}_{}{}^{\prime \prime }=`$
$`={\displaystyle \frac{[p_{12}]}{[p_{12}1]}}a_\alpha ^1[p_{12}][p_{13}+1][p_{23}+1]+{\displaystyle \frac{[p_{13}]}{[p_{13}1]}}a_\alpha ^1[p_{13}][p_{12}+1][p_{23}1]=`$
$`=[p_{12}][p_{13}]([p_{23}+1]+[p_{23}1])a_\alpha ^1=[2]𝒟(p)a_\alpha ^1.`$ (B.7)
The last equality is satisfied due to the CR (2.6) and the ”$`q`$-formula”
$$[p1]+[p+1]=[2][p].$$
## Appendix C. Computation of the scalar square of highest and lowest weight vectors in the $`n=3`$ case
According to the general definition (3.12), the scalar square of the HWV in the $`U_q(sl_3)`$ module $`_p,`$
$$𝒽HWV(p)|HWV(p)𝒾=𝒽\lambda _1\lambda _2|\lambda _1\lambda _2𝒾(p_{12}=\lambda _1+1,p_{23}=\lambda _2+1)$$
is given by
$`𝒽\lambda _1\lambda _2|\lambda _1\lambda _2𝒾=𝒽0|(qa_3^3)^{\lambda _2}(a_{1}^{1}{}_{}{}^{})^{\lambda _1}(a_1^1)^{\lambda _1}(qa_{3}^{3}{}_{}{}^{})^{\lambda _2}|0𝒾=`$
$`=q^{2\lambda _2}𝒽0|(a_{1}^{1}{}_{}{}^{})^{\lambda _1}(a_3^3)^{\lambda _2}(a_{3}^{3}{}_{}{}^{})^{\lambda _2}(a_1^1)^{\lambda _1}|0𝒾`$ (C.1)
where
$$q[p_{12}+1]a_{3}^{3}{}_{}{}^{}=\overline{q}^{\frac{1}{2}}a_2^2a_1^1q^{\frac{1}{2}}a_1^2a_2^1,$$
(C.2)
$$\overline{q}[p_{23}+1]a_{1}^{1}{}_{}{}^{}=\overline{q}^{\frac{1}{2}}a_3^3a_2^2q^{\frac{1}{2}}a_2^3a_3^2.$$
(C.3)
We shall prove (3.34) in four steps.
Step 1 The exchange relation
$$a_3^3a_{3}^{3}{}_{}{}^{}=\frac{[p_{23}][p_{13}]}{[p_{23}1][p_{13}1]}a_{3}^{3}{}_{}{}^{}a_3^3+B_1a_1^3+B_2a_2^3$$
(C.4)
where
$`q^{\frac{3}{2}}[p_{12}+1]B_1=\overline{q}^{p_{23}}[p_{13}+1]a_3^2a_2^1\overline{q}^{p_{13}}a_2^2a_3^1,`$
(C.5)
$`q^{\frac{1}{2}}[p_{12}+1]B_2=\overline{q}^{p_{13}}a_1^2a_3^1\overline{q}^{p_{23}}[p_{13}+1]a_3^2a_1^1,`$
obtained by repeated application of (2.29), implies
$$a_3^3(a_{3}^{3}{}_{}{}^{})^{\lambda _2}(a_1^1)^{\lambda _1}|0𝒾=\frac{[\lambda _2][\lambda _1+\lambda _2+1]}{[\lambda _1+2]}(a_{3}^{3}{}_{}{}^{})^{\lambda _21}a_3^3(a_1^1)^{\lambda _1}a_{3}^{3}{}_{}{}^{}|0𝒾.$$
(C.6)
Proof The last two terms in (C.4), proportional to $`a_1^3`$ and $`a_2^3,`$ do not contribute to (C.6) since, when moved to the right, they yield expressions proportional to $`a_\alpha ^3a_{3}^{3}{}_{}{}^{}|0𝒾(=0\mathrm{for}\alpha =1,2`$). Repeated application of (C.4) in which only the first term in the right hand side is kept gives (C.6).
Step 2 The exchange relation
$$[p_{ij}m]a_\alpha ^j(a_\beta ^i)^m=[p_{ij}](a_\beta ^i)^ma_\alpha ^jq^{ϵ_{\beta \alpha }(p_{ij}m+1)}[m](a_\beta ^i)^{m1}a_\alpha ^ia_\beta ^j$$
(C.7)
which is a consequence of (2.29), implies
$$a_3^3(a_1^1)^{\lambda _1}a_{3}^{3}{}_{}{}^{}|0𝒾=\frac{[p_{13}]}{[p_{13}\lambda _1]}(a_1^1)^{\lambda _1}a_3^3a_{3}^{3}{}_{}{}^{}|0𝒾=\overline{q}^2[\lambda _1+2](a_1^1)^{\lambda _1}|0𝒾.$$
(C.8)
Proof Eq.(C.7) is established by induction in $`m.`$ Eq.(C.8) then follows from the identity $`q^2a_3^3a_{3}^{3}{}_{}{}^{}|0𝒾=[2]|0𝒾.`$
Step 3 Applying $`\lambda _2`$ times steps $`1`$ and $`2`$ one gets
$$q^{2\lambda _2}(a_3^3)^{\lambda _2}(a_{3}^{3}{}_{}{}^{})^{\lambda _2}(a_1^1)^{\lambda _1}|0𝒾=\frac{[\lambda _2]![\lambda _1+\lambda _2+1]!}{[\lambda _1+1]!}(a_1^1)^{\lambda _1}|0𝒾.$$
(C.9)
Step 4 Eqs. (3.19), (C.7) and (2.29) imply
$$(a_{1}^{1}{}_{}{}^{})(a_1^1)^{\lambda _1}|0𝒾=[\lambda _1][\lambda _1+1](a_1^1)^{\lambda _11}|0𝒾;$$
(C.10)
as a result,
$$𝒽\lambda _10|\lambda _10𝒾=𝒽0|(a_{1}^{1}{}_{}{}^{})^{\lambda _1}(a_1^1)^{\lambda _1}|0𝒾=[\lambda _1]![\lambda _1+1]!.$$
(C.11)
The last two steps are obvious.
An analogous computation gives the same result (3.34) for the scalar square of the LWV $`𝒽\lambda _2\lambda _1|\lambda _2\lambda _1𝒾`$.
|
warning/0003/math0003147.html
|
ar5iv
|
text
|
# Characteristic Classes for 𝐺𝑂(2𝑛,𝐶𝐼)
## 1 Introduction
The complex Lie group $`GO(n)`$ is by definition the closed subgroup of $`GL(n)`$ consisting of all matrices $`g`$ such that $`{}_{}{}^{t}gg`$ is a scalar matrix $`\lambda I`$ for some $`\lambda CI^{}`$. (We write simply $`GL(n)`$, $`SO(n)`$, $`O(n)`$, etc. for the complex Lie groups $`GL(n,CI)`$, $`SO(n,CI)`$, $`O(n,CI)`$, etc.)
The group $`GO(1)`$ is just $`CI^{}`$, with classifying space $`BCI^{}=𝐏_{CI}^{\mathrm{}}`$, whose cohomology ring $`H^{}(BCI^{};ZZ/(2))`$ is the polynomial ring $`ZZ/(2)[\lambda ]`$ where $`\lambda H^2(BCI^{};ZZ/(2))`$ is the Euler class. For an odd number $`2n+13`$, the group $`GO(2n+1)`$ is isomorphic to the direct product $`CI^{}\times SO(2n+1)`$. Hence $`BGO(2n+1)`$ is homotopic to the direct product $`BCI^{}\times BSO(2n+1)`$, with cohomology ring the polynomial ring $`ZZ/(2)[\lambda ,w_2,w_3,\mathrm{},w_{2n+1}]`$ where for $`2i2n+1`$, the elements $`w_iH^i(BSO(2n+1);ZZ/(2))`$ are the Stiefel-Whitney classes.
In this paper, we consider the even case $`GO(2n)`$. The main result is an explicit determination, in terms of generators and relations, of the singular cohomology ring $`H^{}(BGO(2n);ZZ/(2))`$. This is the Theorem 4.3.5 below.
An outline of the argument is as follows.
To each action of the group $`CI^{}`$ on any space $`X`$, we functorially associate a certain derivation $`s:H^{}(X)H^{}(X)`$ on the cohomology ring of $`X`$, which is graded of degree $`1`$, with square zero (see Lemma 3.1.2 below). In terms of the action $`\mu :CI^{}\times XX`$ and the projection $`p:CI^{}\times XX`$, it is given by the formula
$$\mu ^{}p^{}=\eta s$$
where $`\eta `$ is the positive generator of $`H^1(CI^{})`$.
When $`X`$, together with its given $`CI^{}`$-action, is the total space of a principal $`CI^{}`$-bundle $`\pi :XY`$ over some base $`Y`$, recall that we have a long exact Gysin sequence
$$\mathrm{}\stackrel{\lambda }{}H^i(Y)\stackrel{\pi ^{}}{}H^i(X)\stackrel{d}{}H^{i1}(Y)\stackrel{\lambda }{}H^{i+1}(Y)\stackrel{\pi ^{}}{}\mathrm{}$$
We show (Lemma 3.2.1) that in this case, the derivation $`s`$ on $`H^{}(X)`$ equals the composite
$$s=\pi ^{}d$$
As $`O(2n)GO(2n)`$ is normal with quotient $`CI^{}`$, $`BO(2n)`$ can be regarded as the total space of a principal $`CI^{}`$-bundle $`BO(2n)BGO(2n)`$. The resulting action of $`CI^{}`$ on $`BO(2n)`$ gives a derivation $`s`$ on the ring $`H^{}(BO(2n))=ZZ/(2)[w_1,\mathrm{},w_{2n}]`$ with $`\mu ^{}p^{}=\eta s`$. The last equality enables us to write the following expression for $`s`$ (see Lemma 4.2.1 below)
$$s=\underset{i=1}{\overset{n}{}}w_{2i1}\frac{}{w_{2i}}$$
Determining the invariant subring $`ker(s)`$ of $`H^{}(BO(2n))`$ thus becomes a purely commutative algebraic problem, which we solve by using the technique of regular sequences and Koszul complex, and express the invariant subring $`B=ker(s)`$ in terms of generators and relations (see Theorem 2.0.1 below).
As explained above, in terms of the Gysin sequence
$$\mathrm{}\stackrel{\lambda }{}H^i\left(BGO\left(2n\right)\right)\stackrel{\pi ^{}}{}H^i\left(BO\left(2n\right)\right)\stackrel{d}{}H^{i1}\left(BGO\left(2n\right)\right)\stackrel{\lambda }{}H^{i+1}\left(BGO\left(2n\right)\right)\stackrel{\pi ^{}}{}\mathrm{}$$
the derivation $`s`$ equals the composite $`\pi ^{}d`$. This fact, together with the explicit knowledge of the kernel and image of $`s`$ (in terms of the ring $`B`$), allows us to ‘solve’ the above long exact Gysin sequence, to obtain the main Theorem 4.3.5, which gives generators and relations for the ring $`H^{}(BGO(2n))`$.
Note that $`H^{}(BGL(2n))`$ is the polynomial ring $`ZZ/(2)[\overline{c}_1,\mathrm{},\overline{c}_{2n}]`$ in the variables $`\overline{c}_i`$ which are the images (mod $`2`$) of the Chern classes $`c_iH^{2i}(BGL(2n);ZZ)`$. The natural inclusion $`GO(2n)GL(2n)`$ induces a ring homomorphism $`H^{}(BGL(2n))H^{}(BGO(2n))`$, and we determine the images of the $`\overline{c}_i`$ in $`H^{}(BGO(2n))`$ in terms of our generators of $`H^{}(BGO(2n))`$ (see Proposition 5.0.1 below).
The paper is arranged as follows. The purely algebraic problem of determining the kernel ring of the derivation $`s`$ is solved in Section 2. In Section 3, we describe the derivation on the cohomology of a space associated to a $`CI^{}`$-action, and connect this with the Gysin sequence for principal $`CI^{}`$-bundles. The above material is applied to determine the cohomology ring of $`BGO(2n)`$ in Section 4. In Section 5, we determine the ring homomorphism $`H^{}(BGL(2n))H^{}(BGO(2n))`$ which is induced by the natural inclusion $`GO(2n)GL(2n)`$. (This section is new in this version, and does not appear in \[H-N\].)
## 2 Derivation on a polynomial algebra
Let $`k`$ denote the field $`ZZ/(2)`$. For any integer $`n1`$, consider the polynomial ring $`C=k[w_1,\mathrm{},w_{2n}]`$ where $`w_1,\mathrm{},w_{2n}`$ are independent variables (the ring $`C`$ will stand for $`H^{}(BO(2n))`$, with $`w_i`$ the Stiefel-Whitney classes, when we apply this to topology).
In what follows, we will treat the odd variables $`w_{2i1}`$ differently from the even variables $`w_{2i}`$, so for the sake of clarity we introduce the notation
$$u_i=w_{2i1}\text{and }v_i=w_{2i}\text{for }1in.$$
On the ring $`C=k[u_1,\mathrm{},u_n,v_1,\mathrm{},v_n]`$, we introduce the derivation
$$s=\underset{i=1}{\overset{n}{}}u_i\frac{}{v_i}:CC$$
Note that $`s`$ is determined by the conditions $`s(u_i)=0`$ and $`s(v_i)=u_i`$ for $`1in`$. As $`2=0`$ in $`k`$, it follows from the above that
$$ss=0$$
We are interested in the subring $`B=ker(s)C`$. Let $`ABC`$ be the subring $`A=k[u_1,\mathrm{},u_n,v_1^2,\mathrm{},v_n^2]`$. For any subset $`T=\{p_1,\mathrm{},p_r\}\{1,\mathrm{},n\}`$ of cardinality $`r1`$, consider the monomial $`v_T=v_{p_1}\mathrm{}v_{p_r}`$. We put $`v_T=1`$ when $`T`$ is empty. It is clear that $`C`$ is a free $`A`$-module of rank $`2^n`$ over $`A`$, with $`A`$-module basis formed by the $`v_T`$. Hence as the ring $`A`$ is noetherian, the submodule $`BC`$ is also finite over $`A`$. In particular, this proves that $`B`$ is a finite type $`k`$-algebra.
We are now ready to state the main result of this section, which is a description of the ring $`B`$ by finitely many generators and relations.
###### Theorem 2.0.1
(I) (Generators) Let the derivation $`s:CC`$ be defined by $`s=_{i=1}^nu_i\frac{}{v_i}`$. Consider the subring $`B=ker(s)C`$. We have $`s^2=0`$, and $`B=ker(s)=im(s)+A`$. Consequently, $`B`$ is generated as an $`A`$-module by $`1`$ together with the $`2^nn1`$ other elements $`s(v_T)`$, where $`T`$ ranges over all subsets $`T\{1,\mathrm{},n\}`$ of cardinality $`2`$, and $`v_T=_{iT}v_i`$.
(II) (Relations) Let $`A[c_T]`$ be the polynomial ring over $`A`$ in the $`2^nn1`$ algebraically independent variables $`c_T`$, indexed by all subsets $`T\{1,\mathrm{},n\}`$ of cardinality $`2`$. Let $`N`$ be the ideal in $`A[c_T]`$ generated by
(1) the $`2^nn(n1)/2n1`$ elements of the form
$$\underset{iT}{}u_ic_{T\{i\}}$$
where $`T\{1,\mathrm{},n\}`$ is a subset of cardinality $`3`$, and
(2) the $`n(n1)/2`$ elements of the form
$$(c_{\{i,j\}})^2+u_i^2v_j^2+u_j^2v_i^2$$
where $`\{i,j\}\{1,\mathrm{},n\}`$ is a subset of cardinality $`2`$, and
(3) the $`(2^nn1)^2n(n1)/2`$ elements (not necessarily distinct) of the form
$$c_Tc_U+\underset{pT}{}\underset{qTU\{p\}}{}u_pv_q^2c_{(T\{p\})\mathrm{\Delta }U}$$
where $`TU`$ in case both $`T`$ and $`U`$ have cardinality $`2`$, and $`\mathrm{\Delta }`$ denotes the symmetric difference of sets $`X\mathrm{\Delta }Y=(XY)(YX)`$.
Then we have an isomorphism of $`A`$-algebras $`A[c_T]/NB`$, mapping $`c_Ts(v_T)`$.
Proof of (I) (Generators) : As already seen, $`ss=0`$, hence we have $`im(s)ker(s)=B`$. As $`AB`$, we have the inclusion $`A+im(s)B`$. We next prove that in fact $`B=A+im(s)`$. The problem translates into the exactness of a certain Koszul resolution, as follows. An exposition of the elementary commutative algebra used below can be found, for example, in §16 of Matsumura \[M\].
Let $`M`$ be a free $`A`$-module of rank $`n`$, with basis $`e_1,\mathrm{},e_n`$. Consider the $`A`$-linear map $`u:MA:e_iu_i`$. For any subset $`T=\{p_1,\mathrm{},p_r\}\{1,\mathrm{},n\}`$ of cardinality $`r`$ where $`0rn`$, let $`e_T=e_{p_1}\mathrm{}e_{p_r}^rM`$. These $`e_T`$ form a free $`A`$-module basis of $`^rM`$. Recall that the Koszul complex for $`u`$ is the complex
$$0\stackrel{n}{}M\stackrel{d_n}{}\stackrel{n1}{}M\stackrel{d_{n1}}{}\mathrm{}\stackrel{d_3}{}\stackrel{2}{}M\stackrel{d_2}{}M\stackrel{d_1}{}A0$$
where for $`1rn`$, the differential $`d_r:^rM^{r1}M`$ is defined by putting
$`d_r(e_T)={\displaystyle \underset{jT}{}}u_je_{T\{j\}}`$ (1)
(note that $`T`$ is non-empty as $`r1`$). As $`u_1,\mathrm{},u_n`$ is a regular sequence in $`A`$, the Koszul complex is exact except in degree $`0`$ (in fact the Koszul complex gives rise to a projective resolution of the $`A`$-module $`A/im(d_1)=A/(u_1,\mathrm{},u_n)`$, but we do not need this).
Let $`(M)=_{0rn}^rM`$, which we regard as a free $`A`$-module of rank $`2^n`$ with basis $`e_T`$. (The algebra structure of $`(M)`$ is not relevant to us.) The graded module $`(M)`$ comes with an $`A`$-linear endomorphism $`d`$ of degree $`1`$, where by definition $`d`$ is $`d_r`$ on the graded piece $`^rM`$. Hence we have the equality $`ker(d)=im(d)+A`$ as submodules of $`(M)`$.
Now consider the $`A`$-linear isomorphism $`\psi :(M)C`$ of free $`A`$-modules under which $`e_Tv_T`$. From the expression $`s=_iu_i(/v_i)`$, it follows that
$`s(v_T)={\displaystyle \underset{jT}{}}u_jv_{T\{j\}}`$ (2)
when $`T`$ is non-empty. By comparing (1) and (2) we see that $`\psi :(M)C`$ takes the Koszul differential $`d`$ to the differential $`s`$, hence $`ker(s)=im(s)+A`$.
As $`im(s)`$ is generated by all the elements $`s(v_T)`$ where $`T`$ is non-empty, the $`A`$-module $`B`$ is generated by $`1`$ together with the $`2^n1`$ other elements $`s(v_T)`$, where $`T`$ ranges over all non-empty subsets $`T\{1,\mathrm{},n\}`$.
This completes the proof of part (I) (Generators) of the Theorem 2.0.1.
Proof of Theorem 2.0.1.(II) (Relations) : As $`ss=0`$, the equation (2) gives $`_{iT}u_is(v_{T\{i\}})=0`$. This proves that the relation II.(1) is satisfied by the assignment $`c_Ts(v_T)`$.
In particular, when $`T=\{i,j\}`$ has cardinality $`2`$, we get $`s(v_{\{i,j\}})=u_iv_j+u_jv_i`$, and hence $`s(v_{\{i,j\}})^2=u_i^2v_j^2+u_j^2v_i^2`$. This proves that the relation II.(2) is satisfied by the assignment $`c_Ts(v_T)`$.
Note that if $`X`$ and $`Y`$ are subsets of $`\{1,\mathrm{},n\}`$, then we have
$`v_Xv_Y={\displaystyle \underset{iXY}{}}v_i^2v_{X\mathrm{\Delta }Y}`$ (3)
where $`X\mathrm{\Delta }Y`$ denotes the symmetric difference $`(XY)(YX)`$.
From the equations (2) and (3) it follows by a straight-forward calculation that
$$s(v_T)s(v_U)=\underset{pT}{}\underset{qTU\{p\}}{}u_pv_q^2s(v_{(T\{p\})\mathrm{\Delta }U})$$
where $`TU`$ in case both $`T`$ and $`U`$ have cardinality $`2`$. This proves that the relation II.(3) is also satisfied by the assignment $`c_Ts(v_T)`$.
Hence the $`A`$-algebra homomorphism $`A[c_T]B`$ defined by sending $`c_Ts(v_T)`$ kills the ideal $`N`$, hence defines an $`A`$-algebra homomorphism $`\eta :A[c_T]/NB`$. By part (I) of the theorem, which we have already proved, the $`A`$-algebra $`B`$ is generated by the $`s(v_T)`$, hence $`\eta `$ is surjective.
To prove that $`\eta `$ is injective, we use the following lemma.
###### Lemma 2.0.2
Let $`FC`$ be the $`A`$-submodule generated by all $`v_T`$ where $`T\{1,\mathrm{},n\}`$ has cardinality $`1`$. Let $`EB`$ be the $`A`$-submodule generated by all $`s(v_T)`$ where $`T\{1,\mathrm{},n\}`$ has cardinality $`2`$. Then we have $`EF`$, and $`FA=0`$, and $`B=AE`$, as $`A`$-modules.
Proof Note that for each $`i`$, the $`v_i`$-degree of each non-zero monomial in any element of $`A`$ is even. On the other hand, if $`|T|1`$ (where $`|T|`$ denotes the cardinality of $`T`$), then the monomial $`v_T`$ has $`v_i`$-degree $`1`$ for for all $`iT`$. Hence in any $`A`$-linear combination $`w=_{|T|1}a_Tv_T`$, each monomial has odd $`v_i`$-degree for some $`i`$. It follows that $`AF=0`$.
From the equality $`s(v_T)=_{iT}u_iv_{T\{i\}}`$, it follows that if $`|T|2`$ then $`s(v_T)F`$, so $`EF`$. It follows that $`AE=0`$. By part (I) of Theorem 2.0.1, we already know that $`A+E=B`$, so the lemma follows. $`\mathrm{}`$
We now prove that $`\eta `$ is injective. For this, we define a homomorphism of $`A`$-modules $`\theta :CA[c_T]/N`$ by $`\theta (1)=0`$, $`\theta (v_i)=u_i`$, and $`\theta (v_T)=c_T`$ for $`|T|2`$. With this definition of $`\theta `$, note that the composite $`C\stackrel{\theta }{}A[c_T]/N\stackrel{\eta }{}BC`$ equals $`s:CC`$. Moreover, note that $`\theta (E)=0`$, where $`EC`$ is the $`A`$-submodule in the previous lemma. To see the above, we must show that $`\theta (s(v_T))=0`$ for all $`T`$. For $`|T|3`$, this follows from the relations (1) in part (II) of Theorem 2.0.1, while for $`|T|2`$, it follows from $`\theta (v_i)=u_i`$ and $`\theta (1)=0`$.
Now suppose $`xA[c_T]/N`$ with $`\eta (x)=0`$. From the generators (3) of $`N`$, it follows that any element of $`xA[c_T]/N`$ can be written in the form
$$x=a_0+\underset{|T|2}{}a_Tc_T\text{where }a_0,a_TA$$
Hence by definition of $`\eta `$, we have $`\eta (x)=a_0+_{|T|2}a_Ts(v_T)`$. As $`\eta (x)=0`$, and as $`AE=0`$ by the above lemma, we have
$$a_0=\underset{|T|2}{}a_Ts(v_T)=0$$
Hence we get
$$x=\underset{|T|2}{}a_Tc_T=\theta (x^{})\text{where }x^{}=\underset{|T|2}{}a_Tv_T$$
As $`s(x^{})=_{|T|2}a_Ts(v_T)=0`$, we have $`x^{}B`$. In the notation of the above lemma, $`x^{}`$ is in $`F`$, and as $`BF=E`$, we see that $`x^{}E`$. As $`\theta (E)=0`$, it finally follows that $`x=\theta (x^{})=0`$.
This proves $`\eta `$ is injective, and completes the proof of the Theorem 2.0.1. $`\mathrm{}`$
## 3 $`CI^{}`$-actions, derivations, and Gysin boundary
### 3.1 The derivation associated to a $`CI^{}`$-action
Given an action $`\mu :CI^{}\times XX`$ of the group $`CI^{}`$ on a space $`X`$, consider the pullback $`\mu ^{}:H^{}(X)H^{}(CI^{}\times X)`$, where the cohomology is with arbitrary coefficients. By Künneth formula, for each $`n`$ we get a homomorphism $`\mu ^{}:H^n(X)H^0(CI^{})H^n(X)H^1(CI^{})H^{n1}(Y)`$. As we have a section $`XCI^{}\times X:x(1,x)`$, it follows that for any $`\alpha `$ in $`H^n(X)`$, we have $`\mu ^{}(\alpha )=1\alpha +\eta \alpha ^{}`$ where $`\eta `$ is the positive generator of $`H^1(CI^{})`$, and the element $`\alpha ^{}H^{n1}`$ is uniquely determined by $`\alpha `$. On the other hand, under the projection $`p:CI^{}\times XX`$ we have $`p^{}(\alpha )=1\alpha `$.
We now define a linear operator $`s:H^n(X)H^{n1}(X)`$ on $`H^{}(X)`$ by the equality $`\mu ^{}(\alpha )=1\alpha +\eta s(\alpha )`$.
#### Example 3.1.1
Let $`X=CI^{}`$, and consider the action $`m:CI^{}\times CI^{}CI^{}:(x,y)xy`$. Then we have the basic equality
$$m^{}(\eta )=1\eta +\eta 1$$
It follows that the corresponding operator $`s`$ on $`H^{}(CI^{})=k[\eta ]/(\eta ^2)`$ is given by $`s(1)=0`$ and $`s(\eta )=1`$.
###### Lemma 3.1.2
Given an action $`\mu :CI^{}\times XX`$ of the group $`CI^{}`$ on a space $`X`$, the map $`s:H^{}(X)H^{}(X)`$ on the singular cohomology ring of $`X`$ with arbitrary coefficients, defined by the equality
$$\mu ^{}p^{}=\eta s$$
where $`p:CI^{}\times XX`$ is the projection and $`\eta `$ is the positive generator of $`H^1(CI^{})`$, is a graded anti-derivation of degree $`(1)`$ on the ring $`H^{}(X)`$, that is, for all $`\alpha H^i(X)`$ and $`\beta H^j(X)`$, we have
$$s(\alpha \beta )=s(\alpha )\beta +(1)^i\alpha s(\beta )H^{i+j1}(X)$$
Moreover, the composite $`ss=0`$.
Proof The verification of the derivation property is straight from the definition. We now prove the property $`ss=0`$. For this consider the commutative diagram
$$\begin{array}{ccc}CI^{}\times CI^{}\times X& \stackrel{(m,id_X)}{}& CI^{}\times X\\ (id_{CI^{}},\mu )& & \mu \\ CI^{}\times X& \stackrel{\mu }{}& X\end{array}$$
where $`m:CI^{}\times CI^{}CI^{}:(x,y)xy`$. As $`m^{}(\eta )=1\eta +\eta 1`$, it follows from the definition of $`s`$ that for any $`\alpha H^i(X)`$,
$`(m,id_X)^{}\mu ^{}(\alpha )`$ $`=`$ $`11\alpha +1\eta s(\alpha )+\eta 1s(\alpha ),\text{ and}`$
$`(id_{CI^{}},\mu )^{}\mu ^{}(\alpha )`$ $`=`$ $`11\alpha +1\eta s(\alpha )+\eta 1s(\alpha )+\eta \eta s^2(\alpha ).`$
Comparing, we get $`\eta \eta s^2(\alpha )=0`$, which means $`s^2(\alpha )=0`$. $`\mathrm{}`$
### 3.2 Derivation in the case of principal $`CI^{}`$-bundles
We next express the above derivation $`s`$ in the case when $`X`$ is the total space of a principal $`CI^{}`$-bundle $`\pi _X:XY`$ over some base $`Y`$, in terms of the Gysin boundary map $`d_X:H^{}(X)H^{}(Y)`$.
###### Lemma 3.2.1
Let $`\pi _X:XY`$ be a principal $`CI^{}`$-bundle. Then the graded anti-derivation $`s_X`$ on the singular cohomology ring $`H^{}(X)`$ with arbitrary coefficients, which is associated to the given $`CI^{}`$-action on $`X`$ by Lemma 3.1.2, equals the composite
$$s_X=\pi _X^{}d_X$$
where $`d_X:H^{}(X)H^{}(Y)`$ is the Gysin boundary map of degree $`(1)`$ and $`\pi _X:H^{}(Y)H^{}(X)`$ is induced by the projection $`\pi _X:XY`$.
Proof Let $`\mu _X:CI^{}\times XX`$ be the action, and $`p_X:CI^{}\times XX`$ be the projection. If the bundle $`X`$ is trivial, then the equality $`\mu _X^{}p_X^{}=\eta s_X`$ is obvious from the definitions. So it remains to prove this equality for a non-trivial $`X`$.
For any map $`f:BY`$, let $`\pi _M:MB`$ denote the pullback bundle. The two projections $`\pi _M:MB`$ and $`r:MX`$, and the Gysin boundary map $`d_M`$, satisfy
$$\pi _M^{}f^{}=r^{}\pi _X^{}\text{and }d_Mr^{}=f^{}d_X$$
Here the second equality follows from the fact that the Euler class of $`X`$ pulls back to the Euler class of $`M`$. Hence the following diagram commutes
$$\begin{array}{ccc}H^n(X)& \stackrel{\pi _X^{}d_X}{}& H^{n1}(X)\\ r^{}& & r^{}\\ H^n(M)& \stackrel{\pi _M^{}d_M}{}& H^{n1}(M)\end{array}$$
For the projection $`(id_{CI^{}},r):CI^{}\times MCI^{}\times X`$, we similarly have
$$p_M^{}r^{}=(id_{CI^{}},r)^{}p_X^{}\text{and }\mu _M^{}r^{}=(id_{CI^{}},r)^{}\mu _X^{}$$
Now take $`f:BY`$ to be $`\pi _X:XY`$. Then $`M`$ is trivial and so the lemma holds for $`s_M`$. The projection $`r:MX`$ has a tautological section $`\mathrm{\Delta }:XM`$ which is the diagonal of $`M=X\times _YX`$. (Under the canonical isomorphism $`(\mu _X,p_X):CI^{}\times XX\times _YX`$, the section $`\mathrm{\Delta }:XM`$ becomes the map $`u(1,u)`$.)
Hence given $`\alpha H^i(X)`$ we have
$`\eta \pi _X^{}d_X(\alpha )`$ $`=`$ $`\eta \mathrm{\Delta }^{}r^{}\pi _X^{}d_X(\alpha )\text{as }r\mathrm{\Delta }=id_X.`$
$`=`$ $`\eta \mathrm{\Delta }^{}\pi _M^{}d_Mr^{}(\alpha )\text{as }r^{}\pi _X^{}d_X=\pi _M^{}d_Mr^{}.`$
$`=`$ $`(id_{CI^{}},\mathrm{\Delta })^{}(\eta \pi _M^{}d_Mr^{}\alpha )`$
$`=`$ $`(id_{CI^{}},\mathrm{\Delta })^{}(\mu _M^{}p_M^{})(r^{}\alpha )\text{as }M\text{ is trivial.}`$
$`=`$ $`(id_{CI^{}},\mathrm{\Delta })^{}(id_{CI^{}},r)^{}(\mu _X^{}p_X^{})(\alpha )`$
$`=`$ $`(\mu _X^{}p_X^{})(\alpha )\text{as }(id_{CI^{}},r)(id_{CI^{}},\mathrm{\Delta })=id_{CI^{}\times X}.`$
$`=`$ $`\eta s_X(\alpha )\text{by definition of }s_X.`$
This proves the lemma in the general case. $`\mathrm{}`$
## 4 Cohomology of $`BGO(2n)`$
Note From now onwards, in the rest of this paper, singular cohomology will be with coefficients $`k=ZZ/(2)`$, unless otherwise indicated.
### 4.1 Principle $`GO(n)`$-bundles and reductions to $`O(n)`$
Principal $`GO(n)`$-bundles and triples $`(E,L,b)`$
Recall that principal $`O(n)`$-bundles $`Q`$ on a space $`X`$ are equivalent to pairs $`(E,q)`$ where $`E`$ is the rank $`n`$ complex bundle on $`X`$ associated to $`Q`$ via the defining representation of $`O(n)`$ on $`CI^n`$, and $`q:EE𝒪_X`$ is the everywhere non-degenerate symmetric bilinear form on $`E`$ with values in the trivial line bundle $`𝒪_X`$ on $`X`$, such that $`q`$ is induced by the standard quadratic form $`x_i^2`$ on $`CI^n`$. The converse direction of the equivalence is obtained via a Gram-Schmidt process, applied locally.
By a similar argument, principal $`GO(n)`$-bundles $`P`$ on $`X`$ are equivalent to triples $`(E,L,b)`$ where $`E`$ is the vector bundle associated to $`P`$ via the defining representation of $`GO(n)`$ on $`CI^n`$, and $`b:EEL`$ is the everywhere nondegenerate symmetric bilinear form on $`E`$ induced by the standard form $`x_i^2`$ on $`CI^n`$, which now takes values in the line bundle $`L`$ on $`X`$, associated to $`P`$ by the homomorphism $`\sigma :GO(n)CI^{}`$ defined by the equality $`{}_{}{}^{t}gg=\sigma (g)I`$.
#### Example 4.1.1
Let $`L`$ be any line bundle on $`X`$. On the rank $`2`$ vector bundle $`F=L𝒪_X`$, we define a nondegenerate symmetric bilinear form $`b`$ with values in $`L`$ by putting $`b:(x_1,x_2)(y_1,y_2)x_1y_2+y_1x_2`$. Now for any $`n1`$, let the triple $`(F^n,L,b)`$ be the orthogonal direct sum of $`n`$ copies of $`(F,L,b)`$. This shows that given any space $`X`$ and any even integer $`2n2`$, there exists some nondegenerate symmetric bilinear triple $`(E,L,b)`$ on $`X`$ of rank $`2n`$, where $`L`$ is a given line bundle.
Reductions to $`O(n)`$
Given a principle $`GO(n)`$-bundle $`P`$ on $`X`$, let $`(E,L,b)`$ be the corresponding triple, and let $`L_o=LX`$ (complement of zero section). As the sequence
$$1O(n)GO(n)\stackrel{\sigma }{}CI^{}1$$
is exact, $`L_oX`$ is the associated $`GO(n)/O(n)`$-bundle to $`P`$, and so reductions of structure group of $`P`$ from $`GO(n)`$ to $`O(n)`$ are the same as global sections (trivializations) of $`L`$. This can be described in purely linear terms, by saying that given an isomorphism $`v:L\stackrel{}{}𝒪_X`$, we get a pair $`(E,vb:EE𝒪_X)`$ from the triple $`(E,L,b)`$.
Let $`u:L\stackrel{}{}𝒪_X`$ and $`v:L\stackrel{}{}𝒪_X`$ be two such reductions. Then the two $`O(n)`$-bundles $`(E,ub)`$ and $`(E,vb)`$ are not necessarily isomorphic. In particular, the two sets of Stiefel-Whitney classes $`w_i(E,ub)`$ and $`w_i(E,vb)`$ need not coincide, but are related as follows.
Given any rank $`n`$ vector bundle $`E`$ together with an $`𝒪_X`$-valued nondegenerate bilinear form $`q:EE𝒪_X`$, consider the new bilinear form $`yq:EE𝒪_X`$ where $`y:XCI^{}`$ is a nowhere vanishing function. Let $`(y)H^1(X,ZZ/(2))`$ be the pull-back of the generator $`\eta H^1(CI^{},ZZ/(2))`$ (in other words, $`(y)`$ is the Kummer class of $`y`$). Note that $`(y)^2=0`$.
A simple calculation using the splitting principle shows the following.
###### Lemma 4.1.2
If $`w_i(E,q)`$ are the Stiefel-Whitney classes of the $`O(n)`$-bundle $`(E,q)`$, then the Stiefel-Whitney classes of the $`O(n)`$-bundle $`(E,yq)`$ are given by the formula
$$w_i(E,yq)=w_i(E,q)+(ni+1)(y)w_{i1}(E,q)$$
#### Remark 4.1.3
It can be seen that the homomorphism $`GO(2n)\{\pm 1\}:g\sigma (g)^n/det(g)`$ is surjective, with connected kernel $`GSO(2n)`$. This implies that $`\pi _0(GO(2n))=ZZ/(2)`$, and hence $`\pi _1(BGO(2n))=ZZ/(2)`$.
### 4.2 The natural derivation on the ring $`H^{}(BO(n))`$
Let $`G`$ be a Lie group and $`H`$ a closed subgroup. If $`PBG`$ is the universal bundle on the classifying space $`BG`$ of $`G`$, then the quotient $`P/H`$ can be taken to be $`BH`$, so that $`BHBG`$ is the bundle associated to $`P`$ by the action of $`G`$ on $`G/H`$. In particular, taking $`G=GO(n)`$ and $`H=O(n)`$, we see that $`BO(n)BGO(n)`$ is the fibration $`L_oBGO(n)`$, where $`L_o`$ is the complement of the zero section of the line bundle $`L`$ on $`BGO(n)`$ occuring in the universal triple $`(E,L,b)`$ on $`BGO(n)`$.
We denote by $`\pi :BO(n)BGO(n)`$ the projection and we denote by $`\lambda H^2(BGO(n))`$ the Euler class of $`L`$. This gives us the long exact Gysin sequence
$$\mathrm{}\stackrel{\lambda }{}H^r\left(BGO\left(n\right)\right)\stackrel{\pi ^{}}{}H^r\left(BO\left(n\right)\right)\stackrel{d}{}H^{r1}\left(BGO\left(n\right)\right)\stackrel{\lambda }{}H^{r+1}\left(BGO\left(n\right)\right)\stackrel{\pi ^{}}{}\mathrm{}$$
By Lemma 3.2.1, the composite maps $`s=\pi ^{}d:H^r(BO(n))H^{r1}(BO(n))`$ define a derivation $`s`$ on the graded ring $`H^{}(BO(n))`$. We now identify this derivation.
Recall that the singular cohomology ring $`H^{}(BO(n))`$ with $`k=ZZ/(2)`$ coefficients is the polynomial ring $`H^{}(BO(n))=k[w_1,\mathrm{},w_n]`$ in the Stiefel-Whitney classes $`w_i`$. We use the convention that $`w_0=1`$.
###### Proposition 4.2.1
Let $`s:H^{}(BO(n))H^{}(BO(n))`$ be the composite $`s=\pi ^{}d`$ as above. Then $`s`$ is a derivation of degree $`(1)`$ on the graded ring $`H^{}(BO(n))`$, with $`ss=0`$. In terms of the universal Stiefel-Whitney classes $`w_i`$, we have $`H^{}(BO(n))=k[w_1,\mathrm{},w_n]`$, and the derivation $`s`$ is given in terms of these generators by
$$s=\underset{i=1}{\overset{n}{}}(ni+1)w_{i1}\frac{}{w_i}:w_i(ni+1)w_{i1}$$
In particular, for $`n=2m`$ even, let $`u_i=w_{2i1}`$ and $`v_i=w_{2i}`$, where $`1im`$ be the generators of the polynomial ring $`H^{}(BO(2m))`$. Then the derivation $`s`$ on the polynomial ring $`k[u_1,\mathrm{}u_m,v_1,\mathrm{},v_m]`$ is given by $`s=u_i\frac{}{v_i}`$, with kernel ring explicitly given in terms of generators and relations by the Theorem 2.0.1.
Proof Let $`(E,L,b)`$ denote the universal triple on $`BGO(n)`$ and $`\pi :L_oBGO(n)`$ the projection, where $`L_o`$ is the complement of the zero section of $`L`$. The pullback $`\pi ^{}(L)`$ has a tautological trivialization $`\tau :\pi ^{}(L)\stackrel{}{}𝒪_{L_o}`$, which gives the universal pair $`(\pi ^{}E,\tau \pi ^{}q)`$ on $`L_o=BO(n)`$. Let $`f=\pi p:CI^{}\times L_oBGO(2n)`$ be the composite of the projection $`p:CI^{}\times L_oL_o`$ with $`\pi :L_oBGO(n)`$.
Under the projection $`p:CI^{}\times L_oL_o`$, the pair $`(\pi ^{}E,\tau \pi ^{}b)`$ pulls back to the pair $`(f^{}E,p^{}(\tau \pi ^{}q))`$ on $`CI^{}\times L_o`$. Under the scalar multiplication $`\mu :CI^{}\times L_oL_o`$, the pair $`(\pi ^{}E,\tau \pi ^{}q)`$ pulls back to the pair $`(f^{}E,yp^{}(\tau \pi ^{}b))`$ on $`CI^{}\times L_o`$, where $`y:CI^{}\times L_oCI^{}`$ is the projection. By Lemma 4.1.2, we have
$$w_i[\mu ^{}(\pi ^{}E,\tau \pi ^{}b)]=w_i[p^{}(\pi ^{}E,\tau \pi ^{}b)]+(ni+1)(y)w_{i1}[p^{}(\pi ^{}E,\tau \pi ^{}b)]$$
in the cohomology ring $`H^{}(CI^{}\times L_o)`$.
Note that the Stiefel-Whitney classes $`w_i(\pi ^{}E,\tau \pi ^{}b)`$ are simply the universal classes $`w_i`$. The class $`(y)`$ becomes $`\eta 1`$ under the Künneth isomorphism. Hence the above formula reads
$$\mu ^{}w_i=p^{}w_i+(ni+1)\eta w_{i1}$$
Hence the proposition follows by lemmas 3.1.2 and 3.2.1. $`\mathrm{}`$
### 4.3 Generators and relations for the ring $`H^{}(BGO(2n))`$
The elements $`\lambda `$, $`a_{2i1}`$, $`b_{4i}`$ and $`d_T`$ of $`H^{}(BGO(2n))`$
Let $`(E,L,b)`$ be the universal triple on $`BGO(2n)`$. Recall that we denote by $`\lambda H^2(BGO(2n))`$ the Euler class of $`L`$. By Example 4.1 applied to $`X=𝐏_{CI}^{\mathrm{}}`$ with $`L=𝒪_X(1)`$ the universal line bundle on $`X`$, together with the universal property of $`BGO(2n)`$, it follows that $`\lambda ^n0\text{for all }n1`$. For each $`1jn`$, we define elements
$$a_{2j1}=dw_{2j}H^{2j1}(BGO(2n)).$$
Note that we therefore have
$$\pi ^{}(a_{2j1})=s(w_{2j})=w_{2j1}H^{2j1}(BO(2n)).$$
More generally, for any subset $`T=\{i_1,\mathrm{},i_r\}\{1,\mathrm{},n\}`$ of cardinality $`r2`$, let $`v_T=w_{2i_1}\mathrm{}w_{2i_r}H^{2\mathrm{deg}(T)}(BO(2n))`$ where $`\mathrm{deg}(T)=i_1+\mathrm{}+i_r`$. We put
$$d_T=d(v_T)H^{2\mathrm{deg}(T)1}(BGO(2n)).$$
Next, let $`E`$ be the rank $`2n`$ complex vector bundle on $`BGO(2n)`$ which occurs in the universal triple $`(E,L,b)`$ on $`BGO(2n)`$. For $`1jn`$, let $`b_{4j}H^{4j}(BGO(2n))`$ be the element
$$b_{4j}=\overline{c}_{2j}(E)H^{4j}(BGO(2n))$$
which is the image of the $`2j`$ th Chern class $`c_{2j}(E)H^{2j}(BGO(2n);ZZ)`$ of $`E`$ under the change of coefficients from $`ZZ`$ to $`ZZ/(2)`$. For any $`m1`$, consider the composite map
$$H^{}(BGL(m);ZZ)H^{}(BGL(m);ZZ/(2))H^{}(BO(m);ZZ/(2))$$
where the first map is induced by the change of coefficients $`ZZZZ/(2)`$ and the second map is induced by the natural inclusion $`O(m)GL(m)`$. The ring $`H^{}(BGL(m);ZZ)`$ is the polynomial ring $`ZZ[c_1,\mathrm{},c_m]`$ where $`c_iH^{2i}(BGL(m);ZZ)`$ is the $`i`$ th universal Chern class. The image of $`c_i`$ in $`H^{}(BO(m);ZZ/(2))`$ is known to be $`w_i^2`$, for all $`1im`$. Hence the elements $`b_{4j}H^{4j}(BGO(2n))`$ satisfy the property that
$$\pi ^{}(b_{4j})=w_{2j}^2\text{for each }1jn.$$
In the ring $`H^{}(BGO(2n))`$, we have $`\lambda a_{2i1}=0`$ for all $`1in`$, and $`\lambda d_T=0`$ for every subset $`T\{1,\mathrm{},n\}`$ of cardinality $`|T|2`$, as follows from the fact that $`\lambda d=0`$ in the Gysin sequence.
The image of $`\pi ^{}`$
We now come to a crucial lemma, one which allows us to write down the image of the ring homomorphism $`\pi ^{}:H^{}(BGO(2n))H^{}(BO(2n))`$.
###### Lemma 4.3.1
We have the equality $`ker(d)=im(\pi ^{})=ker(s)H^{}(BO(2n))`$, in other words, the sequence $`H^i(BGO(2n))\stackrel{\pi ^{}}{}H^i(BO(2n))\stackrel{s}{}H^{i1}(BO(2n))`$ is exact.
Proof We have $`im(s)ker(d)=im(\pi ^{})ker(s)`$ by the exactness of the Gysin sequence. We have $`w_{2i1}=s(w_{2i})=\pi ^{}d(w_{2i})im(\pi ^{})`$ and as $`d(w_{2i}^2)=0`$, $`w_{2i}^2ker(d)=im(\pi ^{})`$. Hence we get the inclusion $`Aim(\pi ^{})`$, where $`A`$ is the polynomial ring in variables $`w_{2i1}`$ and $`w_{2i}^2`$. As already seen, $`im(s)im(\pi ^{})`$, so $`im(s)+Aim(\pi ^{})`$. This completes the proof, as $`ker(s)=im(s)+A`$ by Theorem 2.0.1. $`\mathrm{}`$
Generators for the ring $`H^{}(BGO(2n))`$
###### Lemma 4.3.2
The ring $`H^{}(BGO(2n))`$ is generated by $`\lambda ,(a_{2i1})_i,(b_{4i})_i,(d_T)_TH^{}(BGO(2n))`$.
Proof Let $`SH^{}(BGO(2n))`$ be the subring generated by the elements
$`\lambda ,(a_{2i1})_i,(b_{4i})_i,(d_T)_T`$. By Theorem 2.0.1, the ring $`B`$ is generated by the elements $`w_{2i1}=\pi ^{}(a_{2i1})`$ and $`w_{2i}^2=\pi ^{}(b_{4i})`$ where $`1in`$, together with the elements $`s(v_T)=\pi ^{}(d_T)`$ where $`T\{1,\mathrm{},n\}`$ with $`|T|2`$, which shows that $`\pi ^{}(S)=B`$. As $`1S`$, we have $`H^0(BGO(2n))S`$. Moreover, $`H^1(BGO(2n))=\{0,a_1\}S`$. We now proceed by induction. Suppose that $`H^j(BGO(2n))S`$ for all $`j<i`$. From $`\pi ^{}(S)=B`$, and the fact that $`\pi ^{}`$ is a graded homomorphism, it follows that given $`xH^i(BGO(2n))`$ there exists $`x^{}SH^i(BGO(2n))`$ such that $`\pi ^{}(x)=\pi ^{}(x^{})`$. Hence $`(xx^{})ker(\pi ^{})=im(\lambda )`$, so let $`x=x^{}+\lambda y`$ where $`yH^{i1}(BGO(2n))`$. By induction, $`yS`$, therefore $`xS`$. This proves the lemma. $`\mathrm{}`$
###### Lemma 4.3.3
In the $`k`$-algebra $`H^{}(BGO(2n))`$, the $`n+1`$ elements $`\lambda `$ and $`b_{4i}`$ (where $`1in`$) are algebraically independent over $`k`$.
Proof We recall the following commutative diagram in which the top row is exact.
$$\begin{array}{ccccccc}H^j(BGO(2n))& \stackrel{\pi ^{}}{}& H^j(BO(2n))& \stackrel{d}{}& H^{j1}(BGO(2n))& \stackrel{\lambda }{}& H^{j+1}(BGO(2n))\\ & & & s& \pi ^{}& & \\ & & & & H^{j1}(BO(2n))& & \end{array}$$
Let $`k[x_1,\mathrm{},x_n,y]`$ be a polynomial ring in the $`n+1`$ variables $`x_1,\mathrm{},x_n,y`$. Let $`fk[x_1,\mathrm{},x_n,y]`$ be a non-constant polynomial of the lowest possible total degree, with $`f(b_4,\mathrm{},b_{4n},\lambda )=0`$. Let $`f=f_0+yf_1`$ where $`f_0k[x_1,\mathrm{},x_n]`$. Then as $`\pi ^{}(\lambda )=0`$, we get
$$0=\pi ^{}f(b_4,\mathrm{},b_{4n},\lambda )=\pi ^{}f_0(b_4,\mathrm{},b_{4n})=f_0(w_2^2,\mathrm{},w_{2n}^2)$$
As $`w_{2i}^2`$ are algebraically independent elements of the $`k`$-algebra $`H^{}(BO(2n))`$, the equality $`\pi ^{}(b_{4i})=w_{2i}^2`$ implies that the $`n`$ elements $`b_{4i}H^{}(BGO(2n))`$ are algebraically independent over $`k`$. It follows that $`f_0=0k[x_1,\mathrm{},x_n]`$. Hence we have $`f=yf_1`$, and so $`\lambda f_1(b_4,\mathrm{},b_{4n},\lambda )=f(b_4,\mathrm{},b_{4n},\lambda )=0`$. Hence by exactness of the Gysin sequence, there exists $`zH^{}(BO(2n))`$ with $`f_1(b_4,\mathrm{},b_{4n},\lambda )=dz`$. Applying $`\pi ^{}`$ to both sides, this gives $`\pi ^{}f_1(b_4,\mathrm{},b_{4n},\lambda )=s(z)`$. Now let $`f_1=f_2+yf_3k[x_1,\mathrm{},x_n,y]`$, where $`f_2k[x_1,\mathrm{},x_n]`$. As $`\pi ^{}f_1(b_4,\mathrm{},b_{4n},\lambda )=\pi ^{}f_2(b_4,\mathrm{},b_{4n})`$, we get the equality $`\pi ^{}f_2(b_4,\mathrm{},b_{4n})=s(z)`$, that is,
$$f_2(w_2^2,\mathrm{},w_{2n}^2)im(s)$$
Now, from the formula $`s=_iu_i\frac{}{v_i}`$, it follows that $`im(s)`$ is contained in the ideal generated by the $`u_i=w_{2i1}`$. Hence the above means that $`f_2(w_2^2,\mathrm{},w_{2n}^2)=0`$. Hence as before, $`f_2=0k[x_1,\mathrm{},x_n]`$. Hence we get $`f=yf_1=y^2f_3`$.
As $`\lambda (\lambda f_3(b_4,\mathrm{},b_{4n},\lambda ))=\lambda ^2f_3(b_4,\mathrm{},b_{4n},\lambda )=0`$, by exactness of the Gysin sequence there is some $`qH^{}(BO(2n))`$ with $`d(q)=\lambda f_3(b_4,\mathrm{},b_{4n},\lambda )`$. Hence $`s(q)=\pi ^{}d(q)=\pi ^{}(\lambda f_3(b_4,\mathrm{},b_{4n},\lambda ))=0`$ as $`\pi ^{}(\lambda )=0`$. Hence $`qker(s)`$. The Lemma 4.3.1 showed that $`ker(s)=im(\pi ^{})`$, hence $`q=\pi ^{}(h)`$ for some $`hH^{}(BGO(2n))`$. But then we have
$$\lambda f_3(b_4,\mathrm{},b_{4n},\lambda )=d(q)=d\pi ^{}(h)=0$$
as $`d\pi ^{}=0`$ in the Gysin complex. Hence the polynomial $`g=yf_3k[x_1,\mathrm{},x_n,y]`$, which is non-constant with degree less than that of $`f=yg`$, has the property that $`g(b_4,\mathrm{},b_{4n},\lambda )=0`$. This contradicts the choice of $`f`$, proving the lemma. $`\mathrm{}`$
#### Remark 4.3.4
Let $`k[b_4,\mathrm{},b_{4n},\lambda ]`$ denote the polynomial ring in the variables $`b_4,\mathrm{},b_{4n},\lambda `$. We have a short exact sequence of $`k`$-modules
$$0k[b_4,\mathrm{},b_{4n},\lambda ]\stackrel{\lambda }{}H^{}(BGO(2n))\stackrel{\pi ^{}}{}B0$$
where injectivity of $`\lambda `$ is by Lemma 4.3.3, exactness in the middle is by exactness of the Gysin sequence, and surjectivity of $`\pi ^{}`$ is by Lemma 4.3.1.
Now we state and prove the main result.
Consider the $`2^n+n`$ algebraically independent indeterminates $`\lambda `$, $`a_{2i1}`$ and $`b_{4i}`$ where $`1in`$, and $`d_T`$ where $`T`$ varies over subsets of $`\{1,\mathrm{},n\}`$ of cardinality $`|T|2`$. Let $`k[\lambda ,(a_{2i1})_i,(b_{4i})_i,(d_T)_T]`$ be the polynomial rings in these variables.
###### Theorem 4.3.5
For any $`n1`$, the cohomology ring of $`BGO(2n)`$ with coefficients $`k=ZZ/(2)`$ is isomorphic to the quotient
$$H^{}(BGO(2n))=\frac{k[\lambda ,(a_{2i1})_i,(b_{4i})_i,(d_T)_T]}{I}$$
where $`I`$ is the ideal generated by
(1) the $`n`$ elements $`\lambda a_{2i1}`$ for $`1in`$, and
(2) the $`2^nn1`$ elements $`\lambda d_T`$ where $`T\{1,\mathrm{},n\}`$ is a subset of cardinality $`2`$, and
(3) the $`2^nn(n1)/2n1`$ elements of the form
$$\underset{iT}{}a_{2i1}d_{T\{i\}}$$
where $`T\{1,\mathrm{},n\}`$ is a subset of cardinality $`3`$, and
(4) the $`n(n1)/2`$ elements of the form
$$(d_{\{i,j\}})^2+a_{2i1}^2b_{4j}+a_{2j1}^2b_{4i}$$
where $`\{i,j\}\{1,\mathrm{},n\}`$ is a subset of cardinality $`2`$, and finally
(5) the $`(2^nn1)^2n(n1)/2`$ elements (not necessarily distinct) of the form
$$d_Td_U+\underset{pT}{}\underset{qTU\{p\}}{}a_{2p1}b_{4q}d_{(T\{p\})\mathrm{\Delta }U}$$
where $`TU`$ in case both $`T`$ and $`U`$ have cardinality $`2`$, and where $`\mathrm{\Delta }`$ denotes the symmetric difference of sets $`X\mathrm{\Delta }Y=(XY)(YX)`$.
Proof Let $`R=k[\lambda ,(a_{2i1})_i,(b_{4i})_i,(d_T)_T]`$ denote the polynomial ring. Consider the homomorphism $`RH^{}(BGO(2n))`$ which maps each of the variables in this polynomial ring to the corresponding element of $`H^{}(BGO(2n))`$. This map is surjective by Lemma 4.3.2.
From $`\pi ^{}(\lambda )=0`$ and from the description of $`B`$ in terms of generators and relations given in Theorem 2.0.1, it follows that all the generators of the ideal $`IR`$, which are listed above, map to $`0`$ under $`\pi ^{}:H^{}(BGO(2n))B`$. Hence we get an induced surjective homomorphism
$$\phi :R/IH^{}(BGO(2n))$$
From its definition, we see that the ideal $`IR`$ satisfies $`I\lambda k[\lambda ,b_4,\mathrm{},b_{4n}]=0`$. Hence we get an inclusion $`\lambda k[\lambda ,b_4,\mathrm{},b_{4n}]R/I`$. Next, consider the map $`RB`$ which sends $`\lambda 0`$, $`a_{2i1}w_{2i1}`$, $`b_{4i}w_{2i}^2`$, and $`d_Tc_T`$. From $`\pi ^{}(\lambda )=0`$ and from the description of $`B`$ in terms of generators and relations given in Theorem 2.0.1, it again follows that this map is surjective, and all the generators of the ideal $`IR`$ map to $`0`$, inducing a surjective homomorphism $`\psi :R/IB`$. Hence we get a short exact sequence $`0\lambda k[\lambda ,b_4,\mathrm{},b_{4n}]R/I\stackrel{\psi }{}B0`$.
The above short exact sequence and the short exact sequence of Remark 4.3 fit in the following commutative diagram.
$$\begin{array}{ccccccc}0& \lambda k[\lambda ,b_4,\mathrm{},b_{4n}]& & R/I& \stackrel{\psi }{}& B& 0\\ & & & \phi & & & \\ 0& \lambda k[\lambda ,b_4,\mathrm{},b_{4n}]& & H^{}(BGO(2n))& \stackrel{\pi ^{}}{}& B& 0\end{array}$$
Hence the theorem follows by five lemma. $`\mathrm{}`$
#### Example 4.3.6
In particular, the small dimensional cohomology vector spaces of the $`BGO(2n)`$ are as follows, in terms of linear bases over the coefficients $`k=ZZ/(2)`$.
$`H^0(BGO(2n))`$ $`=`$ $`<1>`$
$`H^1(BGO(2n))`$ $`=`$ $`<a_1>`$
$`H^2(BGO(2n))`$ $`=`$ $`<\lambda ,a_1^2>`$
$`H^3(BGO(2n))`$ $`=`$ $`\{\begin{array}{cc}<a_1^3>\hfill & \text{ for }n=1,\hfill \\ <a_1^3,a_3>\hfill & \text{ for }n2.\hfill \end{array}`$
$`H^4(BGO(2n))`$ $`=`$ $`\{\begin{array}{cc}<\lambda ^2,a_1^4,b_4>\hfill & \text{ for }n=1,\hfill \\ <\lambda ^2,a_1^4,a_1a_3,b_4>\hfill & \text{ for }n2.\hfill \end{array}`$
$`H^5(BGO(2n))`$ $`=`$ $`\{\begin{array}{cc}<a_1^5,a_1b_4>\hfill & \text{ for }n=1,\hfill \\ <a_1^5,a_1^2a_3,a_1b_4,d_{\{1,2\}}>\hfill & \text{ for }n=2,\hfill \\ <a_1^5,a_1^2a_3,a_1b_4,a_5,d_{\{1,2\}}>\hfill & \text{ for }n3.\hfill \end{array}`$
## 5 The map $`H^{}(BGL(2n))H^{}(BGO(2n))`$
The cohomology ring $`H^{}(BGL(2n))`$ is the polynomial ring $`ZZ/(2)[\overline{c}_1,\mathrm{},\overline{c}_{2n}]`$ where the $`2n`$ variables $`\overline{c_i}`$ are the Chern classes mod $`2`$. By definition of $`b_{4r}H^{}(BGO(2n))`$ in terms of the universal triple $`(E,L,b)`$, we have $`\overline{c}_{2i}(E)=b_{4i}`$ for $`1in`$. Hence under the ring homomorphism $`H^{}(BGL(2n))H^{}(BGO(2n))`$ induced by the inclusion $`GO(2n)GL(2n)`$, we have $`\overline{c}_{2i}b_{4i}`$. The following proposition gives the images of the odd Chern classes $`\overline{c}_{2i}`$ in $`H^{}(BGO(2n))`$, which completes the determination of the ring homomorphism $`H^{}(BGL(2n))H^{}(BGO(2n))`$.
###### Proposition 5.0.1
Consider the $`n\times n`$ matrix $`A`$ over $`ZZ/(2)[\lambda ^2]`$, with entries
$$A_{r,k}=\left(\genfrac{}{}{0pt}{}{nk}{2r2k}\right)\lambda ^{2r2k}$$
This is a lower triangular matrix, with all diagonal entries equal to $`1`$, so is invertible over the ring $`ZZ/(2)[\lambda ^2]`$. Let $`B=A^1`$ be its matrix inverse, which is again lower triangular, with all diagonal entries $`1`$. For $`1rn`$, let the polynomial $`f_r(\lambda ,b_4,\mathrm{},b_{4r4})`$ be defined by
$$f_r=\left(\genfrac{}{}{0pt}{}{n}{2r1}\right)\lambda ^{2r2}+\underset{1kr1}{}\left(\genfrac{}{}{0pt}{}{nk}{2r12k}\right)\lambda ^{2r22k}\left(\underset{1jk}{}B_{k,j}(b_{4j}\left(\genfrac{}{}{0pt}{}{n}{2j}\right)\lambda ^{2j})\right)$$
By definition, $`f_r(\lambda ,b_4,\mathrm{},b_{4r4})`$ is linear in $`b_4,\mathrm{},b_{4r4}`$, with the coefficient of $`b_{4r4}`$ equal to the constant $`nr+1ZZ/(2)`$. Then in the cohomology ring $`H^{}(BGO(2n))`$, we have the identity
$$\overline{c}_{2r1}(E)=a_{2r1}^2+\lambda f_r(\lambda ,b_4,\mathrm{},b_{4r4})$$
#### Proof
We divide the proof into two steps, (a) and (b).
(a) For each $`1rn`$, there exists a unique polynomial $`g_r(\lambda ,b_4,\mathrm{},b_{4r4})`$ such that $`\overline{c}_{2r1}(E)=a_{2r1}^2+\lambda g_r(\lambda ,b_4,\mathrm{},b_{4r4})`$.
(b) The polynomials $`g_r(\lambda ,b_4,\mathrm{},b_{4r4})`$ are the polynomials $`f_r(\lambda ,b_4,\mathrm{},b_{4r4})`$ defined in the statement of the proposition.
Proof of (a) The composite homomorphism $`O(2n)GO(2n)GL(2n)`$ induces $`H^{}(BGL(2n))H^{}(BGO(2n))H^{}(BO(2n))`$ under which $`\overline{c}_iw_i^2`$. As $`a_{2i1}^2w_{2i1}^2`$ under $`H^{}(BGO(2n))H^{}(BO(2n))`$, we must have
$$\overline{c}_{2i1}(E)=a_{2i1}^2+x_{4i2}$$
where $`x_{4i2}`$ lies in the kernel of $`\pi ^{}:H^{4i2}(BGO(2n))H^{4i2}(BO(2n))`$. By exactness of Gysin, we have $`x_{4i2}=\lambda y_{4i4}`$ for some $`y_{4i4}H^{4i4}(BGO(2n))`$. Now, by the structure of the ring $`H^{}(BGO(2n))`$, we know that $`\lambda `$ annihilates the $`a_{2j1}`$’s and the $`d_T`$’s. Hence we can replace $`y_{4i4}`$ by a polynomial $`g_i(\lambda ,b_{4j})`$ in $`\lambda `$ and the $`b_{4j}`$’s. As by Lemma 4.3.3 the variables $`\lambda `$ and the $`b_{4j}`$ are algebraically independent, $`g_i`$ is unique. By degree considerations, the highest $`b_{4j}`$ that occurs in $`g_i`$ can be at most $`b_{4i4}`$. This completes the proof of (a).
Proof of (b) We first determine the polynomial $`g_i(\lambda ,b_4,\mathrm{},b_{4i4})`$ in the following Example 5 of a triple $`𝒯=(E,L,b)`$, which generalizes the Example 4.1.
#### Example 5.0.2
Let $`𝒪(1)`$ be the universal line bundle on $`B(CI^{})=𝐏_{CI}^{\mathrm{}}`$. Let $`X=B(CI^{})\times \mathrm{}\times B(CI^{})`$ be the product of $`n+1`$ copies, and let $`p_i:XB(CI^{})`$ be the projections for $`0in`$. Let $`L=p_0^{}(𝒪(1))`$, and for $`1in`$ let $`K_i=p_i^{}(𝒪(1))`$. On $`X`$, we get the nondegenerate triples $`𝒯_i=((LK_i)K_i^1,L,b_i)`$ where $`b_i`$ is induced by the canonical isomorphism $`(LK_i)K_i^1\stackrel{}{}L`$. Let $`𝒯=(E,L,b)`$ be the triple $`𝒯_i`$. We now write down the characteristic classes of $`𝒯`$. Let $`\lambda =\overline{c}_1(L)`$. Then by definition, $`\lambda (𝒯)=\lambda `$. As the odd cohomologies of $`X`$ are zero, the classes $`a_{2i1}(𝒯)`$ and $`d_T(𝒯)`$ are zero, for all $`1in`$ and for all $`T\{1,\mathrm{},n\}`$ with $`|T|2`$.
Let $`x_i=\overline{c}_1(K_i)`$ for $`1in`$, and let $`s_k(\lambda x_i+x_i^2)`$ denote the $`k`$ th elementary symmetric polynomial in the variables $`\lambda x_1+x_1^2,\mathrm{},\lambda x_n+x_n^2`$. As $`E`$ is the direct sum of the $`(LK_i)K_i^1`$, its mod $`2`$ Chern classes $`\overline{c}_i`$ are given by
$`\overline{c}_{2r}(E)`$ $`=`$ $`\left({\displaystyle \genfrac{}{}{0pt}{}{n}{2r}}\right)\lambda ^{2r}+{\displaystyle \underset{1kr}{}}\left({\displaystyle \genfrac{}{}{0pt}{}{nk}{2r2k}}\right)\lambda ^{2r2k}s_k(\lambda x_i+x_i^2)`$
$`\overline{c}_{2r1}(E)`$ $`=`$ $`\left({\displaystyle \genfrac{}{}{0pt}{}{n}{2r1}}\right)\lambda ^{2r1}+{\displaystyle \underset{1kr1}{}}\left({\displaystyle \genfrac{}{}{0pt}{}{nk}{2r12k}}\right)\lambda ^{2r12k}s_k(\lambda x_i+x_i^2)`$
As $`b_{4r}(𝒯)=\overline{c}_{2r}(E)`$, we have the following equations for $`1rn`$.
$$b_{4r}\left(\genfrac{}{}{0pt}{}{n}{2r}\right)\lambda ^{2r}=\underset{1kr}{}\left(\genfrac{}{}{0pt}{}{nk}{2r2k}\right)\lambda ^{2r2k}s_k(\lambda x_i+x_i^2)$$
This is a system of linear equations with coefficients in $`ZZ/(2)[\lambda ^2]`$, for the $`b_{4r}\left(\genfrac{}{}{0pt}{}{n}{2r}\right)\lambda ^{2r}`$ in terms of the $`s_k(\lambda x_i+x_i^2)`$. It is given by the $`n\times n`$ matrix $`A`$ over $`ZZ/(2)[\lambda ^2]`$, with entries
$$A_{r,k}=\left(\genfrac{}{}{0pt}{}{nk}{2r2k}\right)\lambda ^{2r2k}$$
This is a lower triangular matrix, with all diagonal entries equal to $`1`$, so is invertible over the ring $`ZZ/(2)[\lambda ^2]`$. Let $`B=A^1`$ be its matrix inverse, which is again lower triangular, with diagonal entries $`1`$. Hence we get
$$s_r(\lambda x_i+x_i^2)=\underset{1kr}{}B_{r,k}(b_{4k}\left(\genfrac{}{}{0pt}{}{n}{2k}\right)\lambda ^{2k})$$
Substituting this in the equation for $`\overline{c}_{2r1}(E)`$, we get equations
$$\overline{c}_{2r1}(E)=\lambda f_r(\lambda ,b_4,\mathrm{},b_{4r4})$$
where
$$f_r=\left(\genfrac{}{}{0pt}{}{n}{2r1}\right)\lambda ^{2r2}+\underset{1kr1}{}\left(\genfrac{}{}{0pt}{}{nk}{2r12k}\right)\lambda ^{2r22k}\left(\underset{1jk}{}B_{k,j}(b_{4j}\left(\genfrac{}{}{0pt}{}{n}{2j}\right)\lambda ^{2j})\right)$$
Proof of (b) continued In the Example 5, we have the desired equality $`g_i=f_i`$. Note that the cohomology ring $`H^{}(X)`$ is the polynomial ring $`ZZ/(2)[\lambda ,x_1,\mathrm{},x_n]`$, in which the $`n+1`$ elements $`\lambda ,b_4,\mathrm{},b_{4n}`$ are algebraically independent, where $`b_{4r}=\left(\genfrac{}{}{0pt}{}{n}{2r}\right)\lambda ^{2r}+_{1kr}\left(\genfrac{}{}{0pt}{}{nk}{2r2k}\right)\lambda ^{2r2k}s_k(\lambda x_i+x_i^2)`$. Hence as $`g_i=f_i`$ in this example, we get $`g_i=f_i`$ universally.
This completes the proof of Proposition 5.0.1. $`\mathrm{}`$
#### Example 5.0.3
The above proposition in particular gives in $`H^{}(BGO(2n))`$ the identities
$`\overline{c}_1`$ $`=`$ $`a_1^2+n\lambda \text{for all }n1,\text{and}`$
$`\overline{c}_3`$ $`=`$ $`a_3^2+{\displaystyle \frac{n(n1)(2n1)}{6}}\lambda ^3+(n1)\lambda b_4\text{for all }n2.`$
Acknowledgement Y. I. Holla thanks ICTP Trieste, and N. Nitsure thanks the University of Essen, for support while this work was being completed.
## References
\[H-N\] Holla, Y. I. and Nitsure, N. : Characteristic Classes for $`BGO(2n)`$. (Earlier version of this paper, without the last section). To appear in Asian J. Mathematics.
\[M\] Matsumura, H. : Commutative ring theory, Cambridge Studies in Advanced Mathematics 8, Cambridge Univ. Press, 1986.
Address : School of Mathematics, Tata Institute of Fundamental Research, Homi Bhabha Road, Mumbai 400 005, India.
e-mails: yogi@math.tifr.res.in, nitsure@math.tifr.res.in
22-III-2000, revised on 22-V-2000
|
warning/0003/nucl-th0003004.html
|
ar5iv
|
text
|
# 1 Generic diagram of the pp→ppπ⁰. The empty blobs represent initial- and final-state interactions (ISI and FSI), and the shaded blob the two-nucleon irreducible transition operators. The momenta in the asymptotic region have bars attached to them.
## Acknowledgement
We thank M. Rho, K. Kubodera, F. Myhrer and the nuclear theory group of University of South Carolina for useful discussions. SA also thanks B.-Y. Park for discussions. This research was supported in part by the KOSEF Grant 985-0200-001-2 and KRF 1999-015-DI0023, and by the National Science Foundation, Grant No. PHY-9900756 and No. INT-9730847.
|
warning/0003/hep-ph0003280.html
|
ar5iv
|
text
|
# Electroproduction of pseudoscalar mesons on the deuteron.
## I Introduction
The reaction $`\gamma +dd+P^0`$, where $`P^0`$ is a neutral pseudoscalar meson ($`\pi ^0`$ or $`\eta `$), is the simplest coherent meson production process in $`\gamma d`$-collisions. The presence of a deuteron with zero isospin in the initial and final states leads to a specific isotopic structure for the corresponding amplitudes. Moreover, although the spin structure may be, in general, fairly complex, it is essentially simplified in the near theshold region making the $`\gamma +NN+P^0`$ (where $`N`$ denotes a nucleon) and $`\gamma +dd+P^0`$ reactions especially interesting for hadron electrodynamics studies.
The $`\gamma +dd+\pi ^0`$ reaction is important to test the predictions of low energy theorems (LET) for threshold $`\gamma +NN+\pi ^0`$ amplitudes. Multipole analyses of older $`\gamma +pp+\pi ^0`$ data were in serious discrepancy with the predictions of LET . Recent data obtained with tagged photons, combined with new theoretical developments have brought experiment and theory into agreement. These calculations show that the amplitudes for the $`\gamma +NN+\pi ^0`$ reaction near threshold have a complex isotopic structure. Calculations of the electric dipole $`E_{0+}`$ threshold amplitudes for $`\gamma +NN+\pi `$ processes in the framework of the dispersion relation method confirm this observation. Therefore the knowledge of the $`\gamma +nn+\pi ^0`$ reaction amplitude is very important and the $`\gamma +dd+\pi ^0`$ reaction appears the most suitable for that purpose. However the extraction of the $`\gamma +nn+\pi ^0`$ amplitude from $`\gamma d`$ experimental data requires a careful study of possible rescattering effects .
Pion-electropoduction $`e+de+d+\pi ^0`$ is even richer since it involves longitudinal as well as transverse photons. Experimental information about this process has been missing for a long time, but such an experiment can be performed at MAMI or at Jefferson Lab. In this case, with the experimental set-up which has been used to measure the tensor deuteron polarization in elastic $`ed`$-scattering , a sample of $`\pi ^0`$-electroproduction data were obtained during dedicated runs, at relatively large momentum transfer square ($`1.1÷1.6`$ (GeV/c)<sup>2</sup>) in the threshold and in the $`\mathrm{\Delta }`$-region.
The $`e+de+d+\pi ^0`$ reaction allows to ”scan” the isospin structure in the full resonance region and to separate isovector from isoscalar contributions. Moreover, experiments using a polarized deuteron target yield a different information compared to measurements of the polarization of the final deuteron.
Another interesting problem of near threshold meson photoproduction on deuterons concerns the isotopic structure of the $`\gamma +NS_{11}(1535)`$ transition. The results of different multipole analyses of the $`\gamma +NN+\pi `$ reactions have shown that the $`\gamma +NS_{11}(1535)`$ transition is essentially isovector , in agreement with predictions of quark models . Existing $`\gamma +pp+\eta `$ experimental data in the near threshold region indicate that the $`S_{11}(1535)`$ excitation is the main mechanism. On the other hand, the amplitude for $`\gamma +dd+\eta `$ in the near threshold region has to be isoscalar and, therefore, small in contradiction with earlier data which showed a large cross section. Recent $`d(\gamma ,\eta )`$ $`X`$ data have given an explanation by showing that this reaction is essentially inelastic.
The near-threshold region for $`\gamma +dd+\eta `$ is linked with the physics of the $`n+pd+\eta `$ reaction because both processes are connected via the unitarity condition. The cross section for $`n+pd+\eta `$, which was first measured at Saturne , was found to be very large : $`\sigma (npd\eta )=(100\pm 20)`$ $`\mu `$b. A recent experiment at CELSIUS has confirmed these data and has shown a steep decrease down to $`\sigma 40`$ $`\mu `$b up to $`q_{CM}=20`$ MeV, where $`q_{CM}`$ is the final kinetic energy in the reaction center of mass system (CMS). The shape of the energy dependence is reproduced by calculations taking into account the $`N^{}(1535)`$ resonance but more exotic explanations are not ruled out :
the existence of an isoscalar dibaryon resonance with zero isotopic spin and a small width, $`\mathrm{\Gamma }7`$ $`MeV`$, or
the existence of a quasi-bound $`\eta d`$-state (due to the strong $`\eta N`$ interaction), or
the possible presence of a nonperturbative $`s\overline{s}`$-component in the nucleon which could allow a strong $`\eta `$-production from spin singlet $`np`$ initial states .
The study of the processes $`\gamma +dn+p+\eta `$, $`e+de+d+\eta `$ and $`e+de+n+p+\eta `$ with particular emphasis on polarization observables, would help to identify the correct interpretation.
We derive here a general analysis for pseudoscalar meson electroproduction on deuterons, based on general symmetry properties of the hadron electromagnetic interaction. A similar analysis limited to pion photoproduction on deuterons has been published . Such a general analysis has to be considered as the first necessary step in the theoretical study of this process and is no substitute for dynamical model calculations . An adequate dynamical approach to pion electroproduction has to take into account all previous theoretical findings related to other electromagnetic processes on deuteron, such as elastic $`ed`$ scattering , $`\pi ^0`$-photoproduction, $`\gamma +dd+\pi ^0`$ , and deuteron photodisintegration $`\gamma +dn+p`$ . Similarly to these processes, the reaction $`e+de+d+\pi ^0`$ will face two main problems: the study of the deuteron structure and of the reaction mechanism, on one side, and the determination of the neutron elementary amplitude ($`\pi ^0`$-meson electroproduction on neutron, $`e^{}+ne^{}+n+\pi ^0`$), on another side.
Elastic $`ed`$-scattering, being the simplest process to access the deuteron structure, has been considered, for large momentum transfers, a good case to test different predictions of perturbative QCD, such as the scaling behavior of the deuteron electromagnetic form factors and the hypothesis of helicity conservation . The analysis of the scaling behavior should help in defining the kinematical region of the transition regime from the meson-nucleon degrees of freedom to the quark-gluon description of the deuteron structure. In this respect coherent $`\pi ^0`$-electroproduction of the deuteron opens new possibilities to study the scaling phenomena in different regions due to the more flexible kinematical conditions: it unifies the kinematics of elastic $`ed`$-scattering, with its single dynamical variable (the momentum transfer square, $`k^2`$) and the process of $`\pi ^0`$-photoproduction, with two independent dynamical variables (the total energy $`s`$ and the momentum transfer $`t`$ from the initial to the final deuteron). As a result, three kinematical variables drive the process $`e+de+d+\pi ^0`$. Different mechanisms have a leading role in different kinematical regions. In order to interpret the first experimental data for $`e+de+d+\pi ^0`$, with small excitation energy of the produced $`d\pi ^0`$-system (up to the $`\mathrm{\Delta }`$-resonance region) but at relatively large momentum transfer, $`k^2`$, the starting point of the theoretical analysis is naturally the impulse approximation (IA). Similarly to previous calculations of elastic $`ed`$\- scattering and $`\pi ^0`$-photoproduction processes, as a further step, contributions of meson exchange currents (MEC) have to be evaluated in the resonance region, while rescattering effects have to be taken into account in the near threshold region. Large disagreements exist, up to now, in a quantitative evaluation of these effects.
The present paper is organized as follows:
a) we first establish the spin structure of the matrix element for the $`\gamma ^{}+dd+P^0`$ reaction and give a formalism for the description of polarization observables. The dependence of the $`\stackrel{}{d}(e,e^{}P^0)d`$ differential cross section on the polarization characteristics of the deuteron target is derived in a general form, using a formalism of structure functions (SF), which is particularly adequate to describe, in the one-photon approximation, the polarization properties for any $`e+Ae+h+A^{}`$ process (where A is any nucleus and h is a single hadron or hadronic system). These structure functions are further expressed in terms of the scalar amplitudes which parametrize the spin structure of the corresponding electromagnetic current for the process $`\gamma ^{}+dd+P^0`$.
b) we study the isospin structure of these reactions,
c) using the IA , we give relations between the scalar amplitudes, describing the $`\gamma ^{}+dd+P^0`$ and the $`\gamma ^{}+NN+P^0`$ reactions,
d) we then examine the special kinematical conditions corresponding to threshold production,
e) finally, we calculate some observables for $`e+de+d+\pi ^0`$ in the framework of the IA in order to study its sensitivity to the isotopic structure of the $`\gamma ^{}+NN+\pi ^0`$ processes near threshold and in the region of $`\mathrm{\Delta }`$ excitation, at relatively large $`k^2`$.
The present analysis has been extended to the electroproduction of a ”scalar” deuteron ($`i.e.`$ $`np`$ pair with $`J^P=0^+`$) together with a pseudoscalar meson which would be much more difficult to investigate experimentally . The results are available on request to the authors.
## II General formalism for the description of $`e+de+d+P^0`$ processes
### A Derivation of the cross section
The general structure of the differential cross section for the $`e+de+d+P^0`$ reaction can be established in the framework of the one-photon mechanism (Fig. 1) by using only the most general symmetry properties of the hadronic electromagnetic interaction, such as gauge invariance (the conservation of hadronic and leptonic electromagnetic currents) and invariance upon mirror symmetry (parity invariance of the strong and electromagnetic interactions or, in short, $`P`$-invariance). The details of the reaction mechanism and the deuteron structure do not contribute at this step.
The transition matrix element can be written:
$`(ededP)`$ $`=`$ $`{\displaystyle \frac{e^2}{k^2}}\overline{u}(k_2)\gamma _\mu u(k_1)dP\left|\widehat{𝒥_\mu }\right|d{\displaystyle \frac{e^2}{k^2}}\mathrm{}_\mu 𝒥_{\mu ,}`$ (1)
$`\mathrm{}_\mu `$ $``$ $`\overline{u}(k_2)\gamma _\mu u(k_1),𝒥_\mu dP\left|\widehat{𝒥_\mu }\right|d,`$ (2)
where the notations of the particle four-momenta are explained in Fig. 1 and $`𝒥_\mu `$ is the electromagnetic current for the transition $`\gamma ^{}+dd+P^0`$. Using the conservation of leptonic and hadronic currents, ($`k𝒥=k\mathrm{}=0`$) one can rewrite the matrix element in terms of space-like components of currents only :
$`={\displaystyle \frac{e^2}{k^2}}\stackrel{}{e}\stackrel{}{𝒥},\stackrel{}{e}\stackrel{}{\mathrm{}}\stackrel{}{k}{\displaystyle \frac{\stackrel{}{k}\stackrel{}{\mathrm{}}}{k_{0}^{}{}_{}{}^{2}}},`$
where $`k=(k_0,\stackrel{}{k})`$, $`k_0`$ is the energy, $`\stackrel{}{k}`$ is the three-momentum of the virtual photon in the CMS of $`\gamma ^{}+dd+P^0`$. All observables will be determined by bilinear combinations of the components of the hadronic current $`\stackrel{}{𝒥}`$: $`H_{ab}=𝒥_a𝒥_b^{}`$. As a result, we obtain the following formula for the exclusive differential cross section in terms of the tensor components $`H_{ab}`$:
$$\frac{d^3\sigma }{dE_2d\mathrm{\Omega }_ed\mathrm{\Omega }_p}=\frac{\alpha ^2}{64\pi ^3}\frac{E_2}{E_1}\frac{\left|\stackrel{}{q}\right|}{M\sqrt{s}}\frac{1}{1\kappa }\frac{\text{X}}{(k^2)},$$
$$\text{X}=H_{xx}+H_{yy}+\kappa \mathrm{cos}2\phi \left(H_{xx}H_{yy}\right)2\kappa \frac{k^2}{k_{0}^{}{}_{}{}^{2}}H_{zz}$$
$$\sqrt{2\kappa (1+\kappa )\frac{(k^2)}{k_{0}^{}{}_{}{}^{2}}}\left[\mathrm{cos}\phi \left(H_{xz}+H_{zx}\right)\mathrm{sin}\phi \left(H_{yz}+H_{zy}\right)\right]$$
(3)
$$+\kappa \mathrm{sin}2\phi (H_{xy}+H_{yx})\lambda \sqrt{1\kappa }[\sqrt{1+\kappa }(H_{xy}H_{yx})\sqrt{2\kappa \frac{(k^2)}{k_{0}^{}{}_{}{}^{2}}}$$
$$(\mathrm{sin}\phi (H_{xz}H_{zx})\mathrm{cos}\phi (H_{yz}H_{zy})],$$
where $`\kappa ^1=12\stackrel{}{k}_L^2tg^2{\displaystyle \frac{\theta e}{2}}/k^2`$ is the polarization of the virtual photon. Here $`E_1(E_2)`$ is the energy of the initial (final) electron in the lab system; $`\theta _e`$ is the electron scattering angle in the lab; $`d\mathrm{\Omega }_e`$ is the solid angle of the scattered electron in the lab system ; $`d\mathrm{\Omega }_p`$ and $`\stackrel{}{q}`$ are respectively the solid angle and three-momentum of the produced $`P^0`$-meson in the $`CMS`$; $`M`$ is the target mass; $`\stackrel{}{k}_L`$ is the photon three-momentum in the lab system; $`\lambda =\pm 1`$ for the two possible initial electron helicities; $`\phi `$ is the azimuthal angle of the scattered electron with respect to the plane of the reaction $`\gamma ^{}+dd+P^0`$. The coordinate system is such that the $`z`$-axis is along $`\stackrel{}{k}`$ and the $`xz`$ plane is defined by $`\stackrel{}{k}`$ and $`\stackrel{}{q}`$.
The tensor structure of $`H_{ab}=\overline{𝒥_a𝒥_b^{}}`$ (where the line denotes the sum over the final deuteron polarizations) can be written in the following form:
$$H_{ab}=H_{ab}^{(0)}+H_{ab}^{(1)}+H_{ab}^{(2)},$$
(4)
where the indexes $`(0)`$, $`(1)`$ and $`(2)`$ correspond to unpolarized, vector and tensor polarized initial deuterons, respectively. The first term $`H_{ab}^{(0)}`$ can be parametrized as:
$$H_{ab}^{(0)}=\widehat{m}_a\widehat{m}_bh_1+\widehat{n}_a\widehat{n}_bh_2+\widehat{k}_a\widehat{k}_bh_3+\{\widehat{m},\widehat{k}\}_{ab}h_4+i[\widehat{m},\widehat{k}]_{ab}h_5,$$
(5)
with $`\{\widehat{m},\widehat{k}\}_{ab}=\widehat{m}_a\widehat{k}_b+\widehat{m}_b\widehat{k}_a,[\widehat{m},\widehat{k}]_{ab}=\widehat{m}_a\widehat{k}_b\widehat{m}_b\widehat{k}_a.`$ Here $`h_1`$ \- $`h_5`$ are the real SF’s, which depend on $`k^2`$, $`s`$ and $`t`$, $`\widehat{\stackrel{}{n}}=\stackrel{}{k}\times \stackrel{}{q}/\left|\stackrel{}{k}\times \stackrel{}{q}\right|`$, $`\widehat{\stackrel{}{m}}=\widehat{\stackrel{}{n}}\times \widehat{\stackrel{}{k}}`$, $`\widehat{\stackrel{}{k}}=\stackrel{}{k}/|\stackrel{}{k}|`$. The SF’s $`h_1`$ \- $`h_4`$ determine the cross section for the reaction $`e+de+d+P^0`$ with unpolarized particles. The SF $`h_5`$ (the so-called “fifth” structure function) determines the asymmetry of longitudinally polarized electrons scattered by an unpolarized target. This $`T`$-odd contribution is determined by the product of longitudinal and transverse components of the hadron electromagnetic current and it is nonzero only for noncoplanar kinematics, $`\phi 0`$. This contribution is very sensitive to the details of the final state interaction.
The tensor $`H_{ab}^{(1)}`$ is linear in the pseudovector $`\stackrel{}{S}`$ (vector polarization of the initial deuteron) and can be written in the following general form:
$$H_{ab}^{(1)}=\widehat{\stackrel{}{m}}\stackrel{}{S}(\{\widehat{m},\widehat{n}\}_{ab}h_6+\{\widehat{k},\widehat{n}\}_{ab}h_7+i[\widehat{m},\widehat{n}]_{ab}h_8+i[\widehat{k},\widehat{n}]_{ab}h_9)$$
$$+\widehat{\stackrel{}{n}}\stackrel{}{S}(\widehat{m}_a\widehat{m}_bh_{10}+\widehat{n}_a\widehat{n}_bh_{11}+\widehat{k}_a\widehat{k}_bh_{12}+\{\widehat{m},\widehat{k}\}_{ab}h_{13}+i[\widehat{m},\widehat{k}]_{ab}h_{14})$$
$$+\widehat{\stackrel{}{k}}\stackrel{}{S}(\{\widehat{m},\widehat{n}\}_{ab}h_{15}+\{\widehat{k},\widehat{n}\}_{ab}h_{16}+i[\widehat{m},\widehat{n}]_{ab}h_{17}+i[\widehat{k},\widehat{n}]_{ab}h_{18}).$$
(6)
So, 13 real SF’s $`h_6`$ \- $`h_{18}`$ describe the effects of the vector target polarization for the exclusive cross section in the one-photon approximation. The symmetric (antisymmetric) part of $`H_{ab}^{(1)}`$ determines the scattering of unpolarized (polarized) electrons by a vector-polarized target. In particular, it is the symmetric part of $`H_{ab}^{(1)}`$, which induces $`T`$-odd asymmetries in the $`\stackrel{}{d}(e,e^{^{}}P^0)d`$ reaction.
The integration of the tensor $`H_{ab}^{(1)}`$ over $`d\mathrm{\Omega }_p`$ can be done in the following way, typical for inclusive polarized electron-hadron collisions Note, that for an unpolarized deuteron target the following formula holds: $`H_{ab}^{(0)}𝑑\mathrm{\Omega }_p=\delta _{ab}w_1+\widehat{k_a}\widehat{k_b}w_2`$.:
$$H_{ab}^{(1)}d\mathrm{\Omega }{}_{p}{}^{}=i\epsilon _{abc}S_cw_3+i\epsilon _{abc}\widehat{k}_c\stackrel{}{S}\widehat{\stackrel{}{k}}w_4+(\widehat{k_a}[\widehat{\stackrel{}{k}}\times \stackrel{}{S}]_b+\widehat{k_b}[\widehat{\stackrel{}{k}}\times \stackrel{}{S}]_a)w_5.$$
For the inclusive structure functions $`w_3`$ \- $`w_5`$ one obtains the following expressions in terms of integrals of the linear combinations of SF’s $`h_i`$:
$`w_3`$ $`=`$ $`{\displaystyle (h_9h_{14}+h_{17})𝑑\mathrm{\Omega }_p},`$ (7)
$`w_3+w_4`$ $`=`$ $`{\displaystyle h_{17}𝑑\mathrm{\Omega }_p},`$ (8)
$`w_5`$ $`=`$ $`{\displaystyle (h_7h_{13})𝑑\mathrm{\Omega }_p},`$ (9)
$`i.e.`$ most of the exclusive $`SF^{^{}}s`$ $`h_6`$ \- $`h_{18}`$ do not contribute to the inclusive $`SF^{^{}}s`$ $`w_3`$ \- $`w_5`$.
Finally, for the tensor $`H_{ab}^{(2)}`$, characterizing the effects of the tensor target polarization, it is possible to write the following general expression :
$`H_{ab}^{(2)}`$ $`=`$ $`(Q_{cd}\widehat{m}_c\widehat{m}_d)(\widehat{m}_a\widehat{m}_bh_{19}+\widehat{n}_a\widehat{n}_bh_{20}+\widehat{k}_a\widehat{k}_bh_{21}+\{\widehat{m},\widehat{k}\}_{ab}h_{22}+i[\widehat{m},\widehat{k}]_{ab}h_{23})`$ (14)
$`+(Q_{cd}\widehat{n}_c\widehat{n}_d)(\widehat{m}_a\widehat{m}_bh_{24}+\widehat{n}_a\widehat{n}_bh_{25}+\widehat{k}_a\widehat{k}_bh_{26}+\{\widehat{m},\widehat{k}\}_{ab}h_{27}+i[\widehat{m},\widehat{k}]_{ab}h_{28})`$
$`+(Q_{cd}\widehat{m}_c\widehat{k}_d)(\widehat{m}_a\widehat{m}_bh_{29}+\widehat{n}_a\widehat{n}_bh_{30}+\widehat{k}_a\widehat{k}_bh_{31}+\{\widehat{m},\widehat{k}\}_{ab}h_{32}+i[\widehat{m},\widehat{k}]_{ab}h_{33})`$
$`+(Q_{cd}\widehat{m}_c\widehat{n}_d)(\{\widehat{m},\widehat{n}\}_{ab}h_{34}+\{\widehat{k},\widehat{n}\}_{ab}h_{35}+i[\widehat{m},\widehat{n}]_{ab}h_{36}+i[\widehat{k},\widehat{n}]_{ab}h_{37})`$
$`+(Q_{cd}\widehat{k}_c\widehat{n}_d)(\{\widehat{m},\widehat{n}\}_{ab}h_{38}+\{\widehat{k},\widehat{n}\}_{ab}h_{39}+i[\widehat{m},\widehat{n}]_{ab}h_{40}+i[\widehat{k},\widehat{n}]_{ab}h_{41}),`$
where $`Q_{ij}`$ is a tensor polarization component of the deuteron target, $`Q_{ii}=0`$, $`Q_{ij}=Q_{ji}`$, so the density matrix for the initial deuteron can be written as follows:
$$D_{1a}D_{1b}^{}=\frac{1}{3}\left(\delta _{ab}\frac{3}{2}i\epsilon _{abc}S_cQ_{ab}\right).$$
(15)
Therefore, for exclusive reactions like $`A(e,e)A^{}h`$, in the framework of the one-photon mechanism, the effects of the target tensor polarization are characterized by a set of 23 real SF’s, $`h_{19}`$ \- $`h_{41}`$. However the result of the integration of this tensor over the angle $`d\mathrm{\Omega }_p`$ of the $`P^0`$-meson reduces its dependence to 5 real structure functions only :
$`{\displaystyle H_{ab}^{(2}𝑑\mathrm{\Omega }_p}`$ $`=`$ $`\left(Q_{cd}\widehat{k}_c\widehat{k}_d\right)\left[w_6\left(\delta _{ab}\widehat{k}_a\widehat{k}_b\right)+w_7\widehat{k}_a\widehat{k}_b\right]`$
$`+Q_{ab}w_8+\left(Q_a\widehat{k}_b+Q_b\widehat{k}_a\right)w_9+i\left(Q_a\widehat{k}_bQ_b\widehat{k}_a\right)w_{10},Q_a=Q_{ab}\widehat{k}_b.`$
In summary, the exclusive differential cross section for unpolarized electron scattering in $`e^{}+de^{}+d+P`$ is determined by a set of 28 ($`4_0+8_1+16_2=28`$) SF’s, where the indexes $`0`$, $`1`$ and $`2`$ correspond to unpolarized target $`(0)`$, target with vector $`(1)`$ and tensor $`(2)`$ polarizations. For longitudinally polarized electron scattering there are additional $`1_0+5_1+7_2=13`$ SF’s. These 41 SF’s can be divided alternatively into 5 - describing electron scattering by an unpolarized deuteron target, 13 - describing the effect of the vector deuteron polarization and 23 - depending on the tensor deuteron polarization. Taking into account the T-invariance of the electromagnetic interaction of hadrons, we can classify the set of 41 SF’s in $`1_0+8_1+7_2=16`$ T-odd structures and $`4_0+5_1+16_2=25`$ T-even SF, as illustrated in Table 1.
For inclusive hadron electro-production, the number of SF’s reduces to two $`(w_1w_2)`$ for the unpolarized case, three $`(w_3w_5)`$, describing deuteron vector polarization effects and five $`(w_6w_{10})`$, depending on the tensor polarization.
This analysis takes into account the eventual vector and tensor polarizations of the target but not the polarization of the produced particles since a summation over the final polarization states has been done. It can be easily generalized to any other polarization observables such as the recoil deuteron polarization or the spin correlation coefficients.
### B Amplitude analysis
The next step in this analysis, is to establish the spin structure of the matrix element for the $`\gamma ^{}+dd+P^0`$ reaction without any constraint on the kinematical conditions.
This spin structure of the amplitude can be parametrized by different (and equivalent) methods, but for the analysis of polarization phenomena the choice of transverse amplitudes is sometimes preferable. Taking into account the $`P`$-invariance of the electromagnetic interaction of hadrons, the dependence of the amplitude of $`\gamma ^{}+dd+P^0`$ on the $`\gamma ^{}`$ polarization and polarization three-vectors $`\stackrel{}{D_1}`$ and $`\stackrel{}{D_2}`$ of the initial and final deuterons is given by:
$$F(\gamma ^{}ddP^0)=\stackrel{}{e}\widehat{\stackrel{}{m}}(g_1\widehat{\stackrel{}{m}}\stackrel{}{D_1}\widehat{\stackrel{}{n}}\stackrel{}{D_2^{}}+g_2\widehat{\stackrel{}{k}}\stackrel{}{D_1}\widehat{\stackrel{}{n}}\stackrel{}{D_2^{}}+g_3\widehat{\stackrel{}{n}}\stackrel{}{D_1}\widehat{\stackrel{}{m}}\stackrel{}{D_2^{}}+g_4\widehat{\stackrel{}{n}}\stackrel{}{D_1}\widehat{\stackrel{}{k}}\stackrel{}{D_2}^{})$$
$$+\stackrel{}{e}\widehat{\stackrel{}{n}}(g_5\widehat{\stackrel{}{m}}\stackrel{}{D_1}\widehat{\stackrel{}{m}}\stackrel{}{D_2^{}}+g_6\widehat{\stackrel{}{n}}\stackrel{}{D_1}\widehat{\stackrel{}{n}}\stackrel{}{D_2^{}}+g_7\widehat{\stackrel{}{k}}\stackrel{}{D_1}\widehat{\stackrel{}{k}}\stackrel{}{D_2^{}}+g_8\widehat{\stackrel{}{m}}\stackrel{}{D_1}\widehat{\stackrel{}{k}}\stackrel{}{D_2^{}}+g_9\widehat{\stackrel{}{k}}\stackrel{}{D_1}\widehat{\stackrel{}{m}}\stackrel{}{D_2^{}})$$
$$+\stackrel{}{e}\widehat{\stackrel{}{k}}(g_{10}\widehat{\stackrel{}{m}}\stackrel{}{D_1}\widehat{\stackrel{}{n}}\stackrel{}{D_2^{}}+g_{11}\widehat{\stackrel{}{k}}\stackrel{}{D_1}\widehat{\stackrel{}{n}}\stackrel{}{D_2^{}}+g_{12}\widehat{\stackrel{}{n}}\stackrel{}{D_1}\widehat{\stackrel{}{m}}\stackrel{}{D_2^{}}+g_{13}\widehat{\stackrel{}{n}}\stackrel{}{D_1}\widehat{\stackrel{}{k}}\stackrel{}{D_2^{}}),$$
(16)
The process $`\gamma ^{}+dd+P^0`$ is described in the general case, by a set of 9 amplitudes for the absorption of a virtual photon with transverse polarization and 4 amplitudes for the absorption of a virtual photon with longitudinal polarization. These numbers are dictated by the values of the spins of the particles and by the P-invariance of hadron electrodynamics. Taking into account the possible helicities for $`\gamma ^{}`$ and deuterons ( in the initial and final states) one can find $`3(\gamma ^{})\times 3(initialdeuteron)\times 3(finaldeuteron)=27`$ different transitions in $`\gamma ^{}+dd+P^0`$ and 27 corresponding helicity amplitudes $`f_{\lambda _\gamma ;\lambda _1;\lambda _2}`$, where $`\lambda _i`$ are the corresponding helicities. Not all these amplitudes are independent, due to the following relations: $`f_{\lambda _\gamma ;\lambda _1;\lambda _2}=(1)^{\lambda _\gamma \lambda _1\lambda _2}f_{\lambda _\gamma ;\lambda _1;\lambda _2}`$, which result from the P-invariance. It is then possible to find that $`f_{00,0}=0`$ and that it remains only 13 independent complex amplitudes. Therefore the complete experiment requires, at least, the measurement of 25 observables. Let us mention in this respect specific properties of polarization phenomena for inelastic electron-hadron scattering: in exclusive $`e+de+d+P^0`$ processes the virtual photon has a nonzero linear polarization, even for the scattering of unpolarized electrons by an unpolarized deuteron target. Therefore, the study of the $`\phi `$\- and $`\kappa `$-dependences of the $`d(e,eP^0)d`$ differential cross section - at a fixed values of the dynamical variables $`s`$, $`t`$ and $`k^2`$ \- allows, in principle, to find not a single, but 4 different quadratic combinations of scalar amplitudes simultaneously: $`h_1`$, $`h_2`$, $`h_3`$ and $`h_4`$. The relationships between the structure functions $`h_i,i=141`$, and the amplitudes $`g_k,k=113`$, are given in the Appendix.
## III The $`\gamma ^{}+dd+P^0`$ reaction at threshold
### A Derivation of the cross section
The threshold region is defined here as the $`\gamma ^{}`$ energy region in which $`P^0`$-meson production occurs in a S-state. This region may be wide as it happens in $`\gamma +NN+\eta `$ or very narrow as in $`\gamma +pp+\pi ^0`$. This region starts from $`s=\left(M+m_P\right)^2`$, where $`m_P`$ is the mass of the produced pseudoscalar meson, but the momentum transfer squared $`k^2`$ can take any value in the space-like region $`(k^20)`$.
For threshold $`P^0`$-meson production only one three-momentum, $`\stackrel{}{k}`$, is present (instead of two: $`\stackrel{}{k}`$ and $`\stackrel{}{q}`$, in the general case) and the full kinematics of the produced hadronic system is fixed by the kinematical conditions of the scattered electron only, similarly to elastic $`ed`$-scattering. For S-wave production any angular dependence in $`\gamma ^{}+dd+\pi ^0`$ disappears and the corresponding integration can be done trivially : $`X^{(t)}𝑑\mathrm{\Omega }_p=4\pi X^{(t)}.`$ Setting $`\phi =0`$ means that for inclusive electron scattering, the $`xz`$ plane is related to the electron scattering plane. The inclusive cross section is obtained by integrating the differential cross section (2):
$$\frac{d^2\sigma }{dE_2d\mathrm{\Omega }_e}=\frac{\alpha ^2}{16\pi ^2}\frac{E_2}{E_1}\frac{\left|\stackrel{}{q}\right|}{M\sqrt{s}}\frac{1}{1\kappa }\frac{X^{(t)}}{(k^2)},$$
$$X^{(t)}=H_{xx}^{(t)}+H_{yy}^{(t)}+\kappa \left(H_{xx}^{(t)}H_{yy}^{(t)}\right)2\kappa \frac{k^2}{k_{0}^{}{}_{}{}^{2}}H_{zz}^{(t)}\sqrt{2\kappa \left(1+\kappa \right)\frac{(k^2)}{k_{0}^{}{}_{}{}^{2}}}\left(H_{xz}^{(t)}+H_{zx}^{(t)}\right)$$
$$\lambda \sqrt{1\kappa }\left(\sqrt{1+\kappa }\left(H_{xy}^{(t)}H_{yx}^{(t)}\right)+\sqrt{2\kappa \frac{(k^2)}{k_{0}^{}{}_{}{}^{2}}}\left(H_{yz}^{(t)}H_{zy}^{(t)}\right)\right),$$
(17)
where the superscript $`(t)`$ stands for threshold.
The hadronic tensor $`H_{ab}^{(t)}`$, for the case of polarized deuteron target, can be written as :
$$H_{ab}^{(t)}=\left(\delta _{ab}\widehat{k_a}\widehat{k_b}\right)t_1\left(k^2\right)+\widehat{k_a}\widehat{k_b}t_2(k^2)+i\epsilon _{abc}S_ct_3(k^2)+i\epsilon _{abc}\widehat{k_c}\stackrel{}{S}\stackrel{}{k}t_4(k^2)$$
$$+\left[\widehat{k_a}\left(\stackrel{}{k}\times \stackrel{}{S}\right)_b+\widehat{k_b}\left(\stackrel{}{k}\times \stackrel{}{S}\right)_a\right]t_5(k^2)+\left(\stackrel{}{Q}\widehat{\stackrel{}{k}}\right)\left[\left(\delta _{ab}\widehat{k_a}\widehat{k_b}\right)t_6(k^2)+\widehat{k_a}\widehat{k_b}t_7(k^2)\right]$$
$$+Q_{ab}t_8(k^2)+\left(Q_a\widehat{k_b}+Q_b\widehat{k_a}\right)t_9\left(k^2\right)+i\left(Q_a\widehat{k_b}Q_b\widehat{k_a}\right)t_{10}(k^2).$$
(18)
The quantities $`t_i(k^2)`$, $`i=1`$ \- $`10`$, are real structure functions, which are bilinear combinations of threshold electromagnetic form factors which will be defined in the next section.
The symmetrical part of the tensor $`H_{ab}^{(t)}`$ determines the differential threshold cross section for the scattering of unpolarized electrons (by polarized and unpolarized deuterons), and the antisymmetrical part characterizes the scattering of longitudinally polarized electrons.
### B Amplitude analysis
Taking into account the $`P`$-invariance of the hadronic electromagnetic interaction, the following threshold multipole transitions for $`\gamma ^{}+dd+P^0`$ are allowed:
$$E1_{\mathrm{}},E1_t\text{ and }M2𝒥^P=1^{},$$
where $`𝒥`$ and $`P`$ are respectively the total angular momentum and parity of the $`\gamma ^{}d`$ system. Therefore, threshold $`P^0`$-electroproduction is characterized by two transitions with absorption of electric dipole virtual photons (with longitudinal $`\mathrm{}`$ and transverse $`t`$ polarizations) and one transition with absorption of magnetic quadrupole (transverse only) virtual photons.
The threshold amplitude of the process $`\gamma ^{}+dd+P^0`$ can be parametrized in the following way:
$`F_{th}`$ $`=`$ $`\left(\stackrel{}{e}\stackrel{}{D_1}\times \stackrel{}{D_2^{}}\stackrel{}{e}\widehat{\stackrel{}{k}}\stackrel{}{D_1}\times \stackrel{}{D_2^{}}\widehat{\stackrel{}{k}}\right)f_{1t}(k^2)`$ (21)
$`+\stackrel{}{e}\widehat{\stackrel{}{k}}\stackrel{}{D_1}\times \stackrel{}{D_2^{}}\widehat{\stackrel{}{k}}f_{1l}(k^2)`$
$`+\left(\stackrel{}{e}\times \widehat{\stackrel{}{k}}\stackrel{}{D_1}\widehat{\stackrel{}{k}}\stackrel{}{D_2^{}}+\stackrel{}{e}\times \widehat{\stackrel{}{k}}\stackrel{}{D_2^{}}\widehat{\stackrel{}{k}}\stackrel{}{D_1}\right)f_2(k^2),`$
where $`\stackrel{}{e}`$ is the polarization of the virtual $`\gamma `$-quantum.
The form factor $`f_{1t}(k^2)\left[f_1\mathrm{}(k^2)\right]`$ describes the absorption of electric dipole virtual photons with transverse $`[`$longitudinal$`]`$ polarization and the form factor $`f_2(k^2)`$, the absorption of a magnetic quadrupole $`\gamma `$-quantum. They have the same fundamental meaning as the elastic electromagnetic form factors of the deuteron.
Generally they are complex functions of $`k^2`$, due to the unitarity condition (Fig. 2) in the variable $`s`$ ( with a $`n+p`$ system in an intermediate state with both nucleons on the mass shell). But their relative phases have to be equal to $`0`$ or $`\pi `$, as a result of $`T`$-invariance of hadron electrodynamics (theorem of Christ and Lee ). In general, they depend also on the $`s`$ variable, so that $`f_i(k^2)f_i(k^2,s)`$.
In order to have a full reconstruction of the spin structure for $`\gamma ^{}+dd+P^0`$, polarization measurements are necessary. A simple one is the tensor polarization of the scattered deuteron (or the tensor analyzing power using a polarized deuteron target).
After summing over the polarization states of the final deuterons the following expressions can be obtained for the threshold $`SF^{^{}}st_1`$ \- $`t_{10}`$ in terms of the electromagnetic threshold form factors $`f_{1t}(k^2)`$, $`f_1\mathrm{}(k^2)`$ and $`f_2(k^2)`$ :
$$3t_1(k^2)=2\left(\left|f_{1t}(k^2)\right|^2+\left|f_2(k^2)\right|^2\right),$$
$$3t_2(k^2)=2\left|f_1\mathrm{}(k^2)\right|^2\text{ ,}$$
$$t_3(k^2)=\frac{1}{2}𝒻_\mathcal{1}\mathrm{}(𝓀^\mathcal{2})\left(𝒻_{\mathcal{1}𝓉}(𝓀^\mathcal{2})+𝒻_\mathcal{2}(𝓀^\mathcal{2})\right)^{},$$
$$t_4(k^2)=\frac{1}{2}\left|f_{1t}(k^2)f_2\left(k^2\right)\right|^2+\frac{1}{2}𝒻_\mathcal{1}\mathrm{}\left(𝒻_{\mathcal{1}𝓉}(𝓀^\mathcal{2})+𝒻_\mathcal{2}(𝓀^\mathcal{2})\right)^{},$$
$$t_5(k^2)=\frac{1}{2}𝓂𝒻_\mathcal{1}\mathrm{}(𝓀^\mathcal{2})\left(𝒻_{\mathcal{1}𝓉}(𝓀^\mathcal{2})+𝒻_\mathcal{2}(𝓀^\mathcal{2})\right)^{},$$
$$3t_6(k^2)=4𝒻_{\mathcal{1}𝓉}(𝓀^\mathcal{2})𝒻_\mathcal{2}^{}\left(𝓀^\mathcal{2}\right),$$
$$3t_7(k^2)=\left|f_{1t}(k^2)f_2(k^2)\right|^22𝒻_\mathcal{1}\mathrm{}(𝓀^\mathcal{2})\left(𝒻_{\mathcal{1}𝓉}(𝓀^\mathcal{2})+𝒻_\mathcal{2}(𝓀^\mathcal{2})\right)^{},$$
$$3t_8(k^2)=\left|f_{1t}(k^2)f_2(k^2)\right|^2,$$
$$3t_9(k^2)=\left|f_{1t}(k^2)f_2(k^2)\right|^2+𝒻_\mathcal{1}\mathrm{}(𝓀^\mathcal{2})\left(𝒻_{\mathcal{1}𝓉}(𝓀^\mathcal{2})+𝒻_\mathcal{2}(𝓀^\mathcal{2})\right)^{}$$
$$3t_{10}(k^2)=𝓂𝒻_\mathcal{1}\mathrm{}(𝓀^\mathcal{2})\left(𝒻_{\mathcal{1}𝓉}(𝓀^\mathcal{2})+𝒻_\mathcal{2}(𝓀^\mathcal{2})\right)^{}\text{.}$$
(22)
This a strong simplification compared to the 41 real SF’s, depending on 13 complex amplitudes, which are necessary in the general case.
The $`SF`$ $`t_5(k^2)`$ is related to the asymmetry of unpolarized electrons scattered by a vector polarized deuteron target (with polarization orthogonal to the electron scattering plane), while the SF $`h_{10}(k^2)`$ is related to the asymmetry of longitudinally polarized electrons scattered by a deuteron target with tensor polarization. These two $`SF^{^{}}s`$ are determined by the interference of the longitudinal $`\left(f_1\mathrm{}(k^2)\right)`$ and both transverse ($`f_{1t}`$ and $`f_2`$) form factors of the threshold transition $`\gamma ^{}+dd+P^0`$. They define the T-odd polarization observables and must vanish if the relative phase of the longitudinal and transverse form factors is equal to $`0`$ or $`\pi `$ . A dedicated experiment at SLAC for the search of T-odd asymmetry of unpolarized electrons (and positrons) by a polarized proton target - with negative result - remains the best test of T-invariance in hadron electrodynamics (at moderate energies). No similar experiments have been done with a polarized deuteron target but an attempt to detect a nonzero vector deuteron polarization in elastic $`ed`$-scattering has been tried, with a negative result too.
From the expressions, obtained for the $`SF^{^{}}s`$ in terms of the corresponding threshold form factors, one can find an optimal strategy for performing a full experiment on $`P^0`$-meson electroproduction on deuteron near threshold. One must first perform a Rosenbluth separation for the differential cross section of unpolarized electron scattering by an unpolarized target, which allows to find the structure functions $`t_1(k^2)`$ and $`t_2(k^2)`$. These $`SF^{^{}}s`$ determine the total cross sections for the absorption of virtual photons with transverse and longitudinal polarizations. It is straightforward then to deduce, from the longitudinal structure function $`t_2(k^2)`$, the $`k^2`$-dependence of the form factor $`f_1\mathrm{}(k^2)`$ \- for absorption of electric dipole longitudinal virtual photons.
The transverse structure function $`t_1(k^2)`$ contains the contributions of both transverse electromagnetic form factors, namely $`|f_{1t}|^2`$ and $`|f_2(k^2)|^2`$. If we interchange the transverse and longitudinal structure functions, we have a situation similar to elastic $`ed`$-scattering : for elastic $`ed`$-scattering the transverse structure function contains only the contribution of the magnetic form factor, so its $`k^2`$-dependence can be found directly (after a Rosenbluth fit), but the longitudinal structure function contains the contributions of the charge and quadrupole electromagnetic form factors of deuteron. To separate these contributions it is necessary to measure the tensor polarization of scattered deuterons or the tensor analyzing power . From this we can conclude that the measurement of the tensor polarization of the final deuteron in $`e+de+d+P^0`$ near threshold, will allow to separate the contributions due to $`f_{1t}(k^2)`$ and $`f_2(k^2)`$.
This procedure, however, does not give the sign of the threshold form factors. For elastic $`ed`$-scattering, using the well known values of the static electromagnetic characteristics of the deuteron : its electric charge, magnetic and quadrupole moments, it is possible to extrapolate the sign step by step for any values of the momentum transfer square $`k^2`$. We can use the same method for $`\gamma ^{}+dd+P^0`$, using at $`k^2=0`$ the signs of the amplitude for $`\gamma +dd+P^0`$ which can be deduced, in principle, from the signs of the threshold amplitudes for the elementary processes $`\gamma +NN+P^0`$.
We can also find the sign of the $`f_{1t}(k^2)`$, $`f_1\mathrm{}(k^2)`$ and $`f_2(k^2)`$ form factors at any value $`k^2`$ by using their relation with the form factors of the $`\gamma ^{}+NN+P^0`$ reactions at threshold. The matrix element for S-wave $`P^0`$-meson production on a nucleon can be parametrized in terms of two form factors, namely :
$$(\gamma ^{}NNP^0)=\chi _2^+[(\stackrel{}{\sigma }\stackrel{}{e}\stackrel{}{e}\widehat{\stackrel{}{k}}\stackrel{}{\sigma }\widehat{\stackrel{}{k}})f_t(k^2)+f_{\mathrm{}}(k^2)\stackrel{}{e}\widehat{\stackrel{}{k}}\stackrel{}{\sigma }\widehat{\stackrel{}{k}}]\chi _1,$$
(23)
where $`\chi _1`$ and $`\chi _2`$ are the two component spinors of the initial and final nucleons ; $`f_t(k^2)`$ and $`f_{\mathrm{}}(k^2)`$ are the threshold electromagnetic form factors, corresponding to the absorption of electric dipole virtual photons with transverse and longitudinal polarizations. At $`k^2=0`$, $`f_{\mathrm{}}(0)=0`$ and $`f_t(0)=E_{0+}`$ is the threshold electric dipole amplitude for $`\gamma +NN+\pi `$ (with real photons).
In the framework of the IA (Fig. 3) the form factors $`f_{1t}(k^2)`$, $`f_1\mathrm{}(k^2)`$ and $`f_2(k^2)`$ for $`\gamma ^{}+dd+P^0`$ can be directly related to the form factors $`f_{\mathrm{}}(k^2)`$ and $`f_t(k^2)`$ for $`\gamma ^{}+NN+P^0`$.
## IV Impulse Approximation
The most conventional starting point of possible mechanisms for pion electroproduction on the deuteron is the IA . This is, for example, the main mechanism in the region of $`\mathrm{\Delta }`$-excitation, where the rescattering effects for $`\gamma +dd+\pi ^0`$ are negligible . A special attention has to be devoted to the threshold region, for $`\gamma (\gamma ^{})+dd+\pi ^0`$, in particular for pion electroproduction in S-state where the rescattering effects may play an important role. Nevertheless, it is possible to show, in a model independent way, using only the Pauli principle, that the main rescattering contribution due to the following two step process: $`\gamma +dp+p+\pi ^{}(\text{and}n+n+\pi ^+)d+\pi ^0`$ vanishes, when the two nucleons in the $`NN\pi `$-intermediate state are on mass shell. We plan to discuss this problem in a separate paper.
### A Isospin structure of the $`\gamma ^{}+dd+P^0`$ and $`\gamma ^{}+dp+n+P^0`$ reactions
As it is well known, isospin is not conserved in electromagnetic interactions of hadrons, but the hadron electromagnetic current has definite transformation properties relative to isospin symmetry. In general, this current contains an isoscalar and an isovector components. The isotopic spin of deuteron is equal to zero, therefore the amplitude of the $`\gamma ^{}+dd+\pi ^0`$ process is defined by the isovector part of the electromagnetic current only. On the other hand the amplitude of the $`\gamma ^{}+dd+\eta `$ reaction is defined by its isoscalar part.
If the amplitude of the $`\gamma +NN+\eta `$ reaction (with real photons) is driven in the near threshold region by the $`S_{11}(1535)`$ contribution, which is dominated by the isovector part, then the amplitude of $`F(\gamma dd\eta )`$ must be small. However the first $`\gamma dd\eta `$ experiment found a very large cross section . During seventeen years any attempt to resolve this contradiction in the framework of quark models and multipole analyses of the $`\gamma ^{}N\pi N`$ reaction, taking into account effects like rescattering, were unsuccessful. A dedicated experiment with a tagged photon beam, found that the main contribution to the $`d(\gamma ,\eta )X`$ reaction is due, in fact, to the inelastic deuteron break-up $`\gamma +d\eta +n+p`$. In this process, the isovector nature of the transition $`\gamma +NS_{11}(1535)`$ results in the production of a $`(np)`$-system with isotopic spin $`I=1`$. Therefore, near the threshold of the $`\gamma +d\eta +n+p`$ reaction, it must be produced in a singlet state with $`𝒥=0`$. This simplifies drastically the spin structure of the amplitude of the $`\gamma +dd^{}+\eta `$ , $`d^{}=(n+p)_{𝒥=0}`$ process since its coherent part must be determined essentially by the isovector ($`i.e.`$ large) part of the elementary $`\gamma +NN+\eta `$ process (in the framework of IA, (Fig. 4)).
In general, the amplitude for $`\gamma +dn+p+P^0`$ (Fig. 5) contains an isoscalar and an isovector part :
$$F(\gamma dnpP^0)=F_d(t)F(\gamma ppP^0)F_d(u)F(\gamma nnP^0),$$
(24)
where $`F_d`$ is a generalized deuteron form factor, the variables $`t`$ and $`u`$ are the virtual $`p`$ and $`n`$ four-momentum squared. The minus sign in Eq. (24) is the consequence of the specific isotopic structure of the $`dp^{}+n`$ and $`dn^{}+p`$ vertices (with one virtual nucleon $`N^{}`$).
At threshold the $`u`$ and $`t`$ variables are equal: $`u_0=t_0=m^2m_p^2{\displaystyle \frac{m}{2m+m_p}}`$ (in the limit: $`M=2m`$, $`m`$ is the nucleon mass ). Above threshold $`u`$ and $`t`$ are no longer equal.
Rewriting Eq. (24):
$`F(\gamma dnpP)=`$ $`{\displaystyle \frac{1}{2}}\left[F_d(t)+F_d(u)\right]\left[F(\gamma ppP^0)F(\gamma nnP^0)\right]`$ (26)
$`+{\displaystyle \frac{1}{2}}\left[F_d(t)F_d(u)\right]\left[F(\gamma ppP^0)+F(\gamma nnP^0)\right],`$
it is possible, by changing the variables $`u`$ and $`t`$ to control the relative role of the isoscalar and isovector contributions.
As mentioned above, the isotopic structure of the threshold amplitudes for $`\gamma +NN+\pi ^0`$ is a very actual problem. Both coherent processes, $`\gamma +dd+\pi ^0`$ and $`\gamma +dd^{}+\pi ^0`$, are sensitive to this structure but the $`F(\gamma pp\pi ^0)`$ and $`F(\gamma nn\pi ^0)`$ amplitudes contribute differently to these processes. Therefore, the ratio of their cross sections near threshold will be essentially sensitive to the (small) $`E_{0+}`$ electric dipole absorption amplitude in $`\gamma +nn+\pi ^0`$. This ratio can be calculated using the existing experimental value for $`\gamma +pp+\pi ^0`$ : $`E_{0+}(\gamma pp\pi ^0)=(1,32\pm 0.05\pm 0,06){\displaystyle \frac{e}{m_\pi }}10^3`$ and the theoretical predictions for $`\gamma +nn+\pi ^0`$. For example, using the ChPT value as calculated in : $`E_{0+}(\gamma nn\pi ^0)=2.13{\displaystyle \frac{e}{m_\pi }}10^3`$ one would get:
$$R=\frac{\sigma (\gamma dd^{}\pi ^0)}{\sigma (\gamma dd\pi ^0)}=\frac{\left|S\right|^2}{\left|V\right|^2}=\frac{\left|1.32+2.13\right|^2}{\left|1.322.13\right|^2}18.$$
(27)
If instead of ChPT predictions for $`\gamma +nn+\pi ^0`$ we had taken dispersion relations calculations , we would get $`R373`$. We should notice that the dispersion relation calculation for neutron seems to be less stable than the one for proton. In any case these very large variations emphasize the large sensitivity of $`R`$ to the isotopic structure of the $`\gamma +NN+\pi ^0`$ amplitudes.
Besides the real photon point, the $`k^2`$-dependence of $`E_{0+}`$ for both the $`\gamma ^{}+pp+\pi ^0`$ and $`\gamma ^{}+nn+\pi ^0`$ reactions is also very interesting .
### B Relationship between the $`\gamma ^{}+dd+P^0`$ and $`\gamma ^{}+NN+P^0`$ amplitudes
In the framework of IA (Fig. 3), the matrix element $`(\gamma ^{}ddP^0)`$ for the $`\gamma ^{}+dd+P^0`$ process can be written:
$``$ $`=`$ $`2{\displaystyle d^3\stackrel{}{p}𝒯r\phi ^+\left(\left|\stackrel{}{p}+\frac{1}{4}\stackrel{}{Q}\right|\right)\widehat{F}\phi \left(\left|\stackrel{}{p}\frac{1}{4}\stackrel{}{Q}\right|\right)},`$ (28)
$`\stackrel{}{Q}`$ $`=`$ $`\stackrel{}{k}\stackrel{}{q},2\stackrel{}{p}=\stackrel{}{p_1}\stackrel{}{p_2}+{\displaystyle \frac{1}{2}}\stackrel{}{Q},`$ (29)
where $`\stackrel{}{P_1}=\stackrel{}{p_1}+\stackrel{}{p_2}`$, $`\stackrel{}{P_2}=\stackrel{}{p_1}^{^{}}+\stackrel{}{p_2}`$ , and $`\stackrel{}{k}+\stackrel{}{p_1}=\stackrel{}{q}+\stackrel{}{p_1}^{^{}}`$ (the notation is explained in Fig. 3),
$$F(\gamma NNP^0)=\chi _2^+\widehat{F}\chi _1,$$
$$\widehat{F}=\left(\stackrel{}{\sigma }\stackrel{}{K}+L\right)/2$$
(30)
and $`\stackrel{}{K}`$, $`L`$ are the spin-dependent and spin-independent contributions to the matrix $`\widehat{F}`$.
For the deuteron wave function we shall use the following representation, which takes into account the $`S`$\- and $`D`$-waves in the $`np`$-system :
$`\phi (\stackrel{}{p})`$ $`=`$ $`{\displaystyle \frac{1}{(2\pi )^{\frac{3}{2}}}}{\displaystyle d^3\stackrel{}{r}e^{i\stackrel{}{p}\stackrel{}{r}}\phi (r)},`$ (31)
$`\phi (r)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{4\pi }}}\left[\stackrel{}{\sigma }\stackrel{}{D}{\displaystyle \frac{u(r)}{r}}+{\displaystyle \frac{w(r)}{\sqrt{2}r}}\left(3{\displaystyle \frac{\stackrel{}{\sigma }\stackrel{}{r}\stackrel{}{D}\stackrel{}{r}}{r^2}}\stackrel{}{\sigma }\stackrel{}{D}\right)\right]{\displaystyle \frac{i\sigma _2}{\sqrt{2}}},`$ (32)
where $`u(r)`$ and $`w(r)`$ are the standard wave functions of the $`S`$\- and $`D`$-states in deuteron. Expression (18) is particularly convenient to establish the spin structure of the amplitude of the $`\gamma ^{}+dd+P^0`$ process. Since in general the amplitudes $`\stackrel{}{K}`$ and $`L`$ (for the processes $`\gamma ^{}+NN+P^0`$) depend on the integration momentum $`\stackrel{}{p}`$ in (18), the wave functions $`u`$ and $`w`$ of the initial and final deuteron will not depend on the same variable. Indeed, due to the nonlocality of $`\gamma ^{}NNP^0`$ vertex, the coordinates $`\stackrel{}{r}`$ and $`\stackrel{}{r}^{^{}}`$ of the initial and final deuterons do not coincide. However, choosing the $`\stackrel{}{K}`$ and $`L`$ amplitudes at a particular value of the internal momentum $`\stackrel{}{p_1}`$, $`\widehat{F}`$ can be taken outside the integration symbol. This allows to express the quantity $``$ in terms of a definite combination of deuteron form factors, multiplied by the isovector (isoscalar) amplitudes for the $`\gamma ^{}+NN+\pi ^0`$ $`(\gamma ^{}+NN+\eta )`$ reaction (factorization hypothesis).
This procedure is usually justified by a rapid fall-off of $`\phi (|\stackrel{}{p})|`$ when $`\left|\stackrel{}{p}\right|`$ increases and by a (relatively) weak dependence of the $`\stackrel{}{K}`$ and $`L`$ amplitudes on $`\left|\stackrel{}{p}\right|`$.
After some transformations, Eq (18) becomes:
$`(\gamma ^{}ddP^0)`$ $`=`$ $`\stackrel{}{D_1}\stackrel{}{D_2^{}}\widehat{L}F_1(\stackrel{}{Q^2})+2\left(3\stackrel{}{D_1}\widehat{\stackrel{}{Q}}\stackrel{}{D_2^{}}\stackrel{}{Q}\stackrel{}{D_1}\stackrel{}{D_2}\right)\widehat{L}F_2(\stackrel{}{Q^2})`$ (34)
$`+i\widehat{\stackrel{}{K}}\stackrel{}{D_1}\times \stackrel{}{D_2^{}}\left(F_3(\stackrel{}{Q^2})+F_4(\stackrel{}{Q^2})\right)3i\widehat{\stackrel{}{K}}\widehat{\stackrel{}{Q}}\widehat{\stackrel{}{Q}}\stackrel{}{D_1}\times \stackrel{}{D_2^{}}F_4(\stackrel{}{Q^2}),`$
with $`\widehat{\stackrel{}{Q}}=\left(\stackrel{}{k}\stackrel{}{q}\right)/|\stackrel{}{k}\stackrel{}{q}|`$, where $`\widehat{\stackrel{}{K}}`$ and $`\widehat{L}`$ are the values of $`\stackrel{}{K}`$ and $`L`$ for a definite value of $`\stackrel{}{p_1}`$ (see below).
The generalized deuteron form factors $`F_i(\stackrel{}{Q^2})`$ are defined by :
$$F_1(\stackrel{}{Q^2})=_0^{\mathrm{}}𝑑rj_0\left(\frac{Qr}{2}\right)\left[u^2(r)+w^2(r)\right],$$
$$F_2(\stackrel{}{Q^2})=_0^{\mathrm{}}𝑑rj_2\left(\frac{Qr}{2}\right)\left[u(r)\frac{w(r)}{\sqrt{8}}\right]w(r),$$
$$F_3(\stackrel{}{Q^2})=_0^{\mathrm{}}𝑑rj_0\left(\frac{Qr}{2}\right)\left[u^2(r)\frac{1}{2}w^2(r)\right],$$
(35)
$$F_4(\stackrel{}{Q^2})=_0^{\mathrm{}}𝑑rj_2\left(\frac{Qr}{2}\right)\left[u(r)+\frac{1}{\sqrt{2}}w(r)\right]w(r),$$
$$j_0(x)=\frac{\mathrm{sin}x}{x},j_2(x)=\mathrm{sin}x\left(\frac{3}{x^3}\frac{1}{x}\right)3\frac{\mathrm{cos}x}{x^2}.$$
The combinations of the deuteron wave functions $`u(r)`$ and $`w(r)`$ in $`F_i(\stackrel{}{Q^2})`$ define the charge, the magnetic and quadrupole form factors of the deuteron. The fourth form factor $`F_4`$ in Eq. (22) is associated with a nonconservation of the current of the transition $`dd+\pi ^0`$, due to the specific structure of the triangle diagram contribution.
The calculated form factors, $`F_i(\stackrel{}{Q^2})`$, using Bonn and Paris deuteron wave functions, are shown in Fig. 6.
The quantity $`\stackrel{}{Q^2}`$ characterizes the value of the four-momentum transfer squared $`t`$ in the $`\gamma ^{}+dd+P^0`$ reaction,
$`t=(kq)^2=2M\left(M\sqrt{M^2+\stackrel{}{Q^2}}\right),`$
so that $`t\stackrel{}{Q^2}`$, when $`\left|\stackrel{}{Q}\right|M`$.
Obviously, the structure of the $`\gamma ^{}+dd+P^0`$ amplitude, Eq. (21), is not the most general one, even in the case of arbitrary values of $`\stackrel{}{K}`$ and $`L`$ and deuteron form factors $`F_i(\stackrel{}{Q^2})`$. Let us consider first the general spin structure of the amplitude for $`\gamma ^{}+NN+P^0`$ process :
$$(\gamma ^{}NN\pi )=\chi _2^{}\chi _1,$$
$$=i\stackrel{}{e}\widehat{\stackrel{}{k}}\times \widehat{\stackrel{}{q}}f_1+\stackrel{}{\sigma }\stackrel{}{e}f_2+\stackrel{}{\sigma }\widehat{\stackrel{}{k}}\stackrel{}{e}\widehat{\stackrel{}{q}}f_3+\stackrel{}{\sigma }\widehat{\stackrel{}{q}}\stackrel{}{e}\widehat{\stackrel{}{q}}f_4$$
(36)
$$+\stackrel{}{e}\widehat{\stackrel{}{k}}(\stackrel{}{\sigma }\widehat{\stackrel{}{k}}f_5+\stackrel{}{\sigma }\widehat{\stackrel{}{q}}f_6),$$
where $`f_i=f_i(s_1,t,k^2)`$ are the scalar amplitudes for $`\gamma ^{}+NN+P^0`$, so that
$$L=if_1\stackrel{}{e}\widehat{\stackrel{}{k}}\times \widehat{\stackrel{}{q}},$$
$$\stackrel{}{K}=\stackrel{}{e}f_2+\widehat{\stackrel{}{k}}\left(\stackrel{}{e}\widehat{\stackrel{}{q}}f_3+\stackrel{}{e}\widehat{\stackrel{}{k}}f_5\right)+\widehat{\stackrel{}{q}}\left(\stackrel{}{e}\widehat{\stackrel{}{q}}f_4+\stackrel{}{e}\widehat{\stackrel{}{k}}f_6\right).$$
(37)
Comparing the expression (21) for the amplitude $`(\gamma ^{}ddP^0)`$ in IA with the general spin structure of the amplitude one can establish a definite connection between both sets of scalar amplitudes, namely $`g_i`$, $`i=113`$, for $`\gamma ^{}+dd+P^0`$ (on one side) and $`f_k`$, $`k=16`$, for $`\gamma ^{}+NN+P^0`$ (on another side). Their exact relations are given below:
$`g_1`$ $`=`$ $`g_3=\mathrm{sin}\theta (f_3+\mathrm{cos}\theta f_4)\left(F_3\left(\stackrel{}{Q}^2\right)+F_4\left(\stackrel{}{Q}^2\right)\right)`$ (39)
$`3Q_k\left(Q_mf_2+Q_k\mathrm{sin}\theta f_3+Q_q\mathrm{sin}\theta f_4\right)F_4\left(\stackrel{}{Q}^2\right),`$
$`g_2`$ $`=`$ $`g_4=(f_2+\mathrm{sin}^2\theta f_4)\left(F_3\left(\stackrel{}{Q}^2\right)+F_4\left(\stackrel{}{Q}^2\right)\right)`$ (41)
$`+3Q_m\left(Q_mf_2+Q_k\mathrm{sin}\theta +Q_q\mathrm{sin}\theta f_4\right)F_4\left(\stackrel{}{Q}^2\right),`$
$`g_5`$ $`=`$ $`\mathrm{sin}\theta f_1\left[F_1\left(\stackrel{}{Q}^2\right)+2F_2\left(\stackrel{}{Q}^2\right)\left(3Q_m^21\right)\right],`$ (42)
$`g_6`$ $`=`$ $`\mathrm{sin}\theta f_1\left[F_1\left(\stackrel{}{Q}^2\right)2F_2\left(\stackrel{}{Q}^2\right)\right],`$ (43)
$`g_7`$ $`=`$ $`\mathrm{sin}\theta f_1\left[F_1\left(\stackrel{}{Q}^2\right)+2F_2\left(\stackrel{}{Q}^2\right)\left(3Q_k^21\right)\right],`$ (44)
$`g_8`$ $`=`$ $`6\mathrm{sin}\theta Q_mQ_kf_1F_2\left(\stackrel{}{Q}^2\right)f_2\left(F_3\left(\stackrel{}{Q}^2\right)+F_4\left(\stackrel{}{Q}^2\right)\right),`$ (45)
$`g_9`$ $`=`$ $`6\mathrm{sin}\theta Q_mQ_kf_1F_2\left(\stackrel{}{Q}^2\right)+f_2\left(F_3\left(\stackrel{}{Q}^2\right)+F_4\left(\stackrel{}{Q}^2\right)\right),`$ (46)
$`g_{10}`$ $`=`$ $`g_{12}=(f_2+f_3\mathrm{cos}\theta +f_4\mathrm{cos}^2\theta +f_5+f_6\mathrm{cos}\theta (F_3\left(\stackrel{}{Q}^2\right)+F_4\left(\stackrel{}{Q}^2\right))`$ (48)
$`3Q_k[Q_k(f_2+f_3\mathrm{cos}\theta +f_5)+Q_q(\mathrm{cos}\theta f_4+f_6)]F_4\left(\stackrel{}{Q}^2\right),`$
$`g_{11}`$ $`=`$ $`g_{13}=\mathrm{sin}\theta (cos\theta f_4+f_6)[F_3\left(\stackrel{}{Q}^2\right)+F_4(\stackrel{}{Q}^2)]+3Q_m[Q_k(f_2+\mathrm{cos}\theta f_3+f_5)+`$ (50)
$`Q_q(\mathrm{cos}\theta f_4+f_6)F_4\left(\stackrel{}{Q}^2\right),`$
where $`Q_m=\widehat{\stackrel{}{Q}}\widehat{\stackrel{}{m}}`$, $`Q_k=\widehat{\stackrel{}{Q}}\stackrel{}{k}`$, $`Q_m^2+Q_k^2=1`$, $`Q_m^2=\mathrm{sin}^2\theta {\displaystyle \frac{\stackrel{}{q}^2}{|\stackrel{}{k}\stackrel{}{q}|^2}}`$, and $`\theta `$ is the $`P^0`$-meson production angle in $`CMS`$ of $`\gamma ^{}+dd+P^0`$ process.
Note that the relations $`g_1+g_3=g_2+g_4=g_{10}+g_{12}=g_{11}+g_{13}=0,`$ which are correct for any amplitude $`f_k`$, result from the factorization hypothesis.
Neglecting the $`D`$-wave contribution, we can predict that the following ratios:
$$\left(H_{xx}^{(0)}H_{yy}^{(0)}\right)/\left(H_{xx}^{(0)}+H_{yy}^{(0)}\right)_0=\left(|f_3|^2+|f_4|^2|f_1|^2|f_2|^2\right)/\left(|f_1|^2+|f_2|^2+|f_3|^2+|f_4|^2\right),$$
$$\left(H_{zz}^{(0)}\right)/\left(H_{xx}^{(0)}+H_{yy}^{(0)}\right)_0=\left(|f_5|^2+|f_6|^2\right)/\left(|f_1|^2+|f_2|^2+|f_3|^2+|f_4|^2\right),$$
(51)
which do not depend on deuteron form factors and therefore on the deuteron structure.
### C Threshold $`\pi ^0`$ electroproduction within the impulse approximation
At threshold the $`L`$ and $`\stackrel{}{K}`$ amplitudes for $`\gamma ^{}+NN+P^0`$ reduce to:
$$\stackrel{}{K}=f_t\left(\stackrel{}{e}\widehat{\stackrel{}{k}}\stackrel{}{e}\widehat{\stackrel{}{k}}\right)+f_{\mathrm{}}\widehat{\stackrel{}{k}}\stackrel{}{e}\widehat{\stackrel{}{k}},L=0,$$
where $`f_{\mathrm{}}(k^2)`$ and $`f_t(k^2)`$ are the threshold form factors for $`\gamma ^{}+NN+\pi ^0`$, corresponding to absorption of electric dipole virtual photons with longitudinal and transverse polarizations.
Taking into account the fact that, at threshold, $`\widehat{\stackrel{}{Q}}=\widehat{\stackrel{}{k}}`$, one obtains the following form factors for $`\gamma ^{}+dd+\pi ^0`$, which are correct in the framework of the IA :
$$f_{1t}(k^2)=f_t(k^2)\left(F_3+F_4\right),f_1\mathrm{}(k^2)=f_{\mathrm{}}(k^2)\left(F_32F_4\right),f_2(k^2)=0,$$
$`i.e.`$ the magnetic quadrupole form factor $`f_2(k^2)`$ is equal to zero in this approximation, independently of the deuteron structure.
## V Model for $`\gamma ^{}+NN+\pi ^0`$
In order to calculate the scalar amplitudes $`g_i,i=113`$, for the $`\gamma ^{}+dd+\pi ^0`$ process in framework of IA , it is necessary to know the $`\stackrel{}{Q}^2`$-dependence of the deuteron form factors $`F_j(Q^2),j=14`$, from one side, and the elementary amplitudes $`f_k,k=16`$, for the process $`\gamma ^{}+NN+\pi ^0`$, from another side. In order to calculate the isovector part of the amplitudes $`f_k`$ for $`\gamma ^{}+NN+\pi ^0`$, we shall use the effective Lagrangian approach-with a standard set of contributions (Fig. 7). Such model has successfully reproduced the experimental data , for the process $`e+pe+p+\pi ^0`$ in the following kinematical conditions: $`1.1W1.4`$ GeV and $`k^2=2.8`$ and 4.0 (GeV/c)<sup>2</sup>, in the whole domain of cos $`\theta _\pi `$ and azymuthal angle $`\varphi `$. The main ingredients of this calculation were the s- and u-channel contributions of N and $`\mathrm{\Delta }`$, with particular attention to the ’off-shell’ properties of the $`\mathrm{\Delta }`$isobar. The comparison with the experimental data allowed to determine the following values of the electromagnetic form factors for the $`\gamma ^{}+p\mathrm{\Delta }^+`$ transition:
$$G_M^{}/3G_d,G_d=(1k^2/0.71\text{GeV}^2)^2,R_{EM}=E_{1+}/M_{1+},R_{SM}=S_{1+}/M_{1+},$$
where $`M_{1+},E_{1+}`$ and $`S_{1+}`$ denote the magnitude of the magnetic dipole, electric (transversal) quadrupole and Coulomb (longitudinal) quadrupole amplitudes or transition form factors for the $`\gamma ^{}+N\mathrm{\Delta }`$ excitation. In our analysis we will use the two following results of the experiment :
* The magnetic dipole form factor $`G_M^{}(k^2)`$ dominates, i.e. the ratios $`R_{EM}`$ and $`R_{SM}`$ are small (in absolute value);
* The magnetic dipole form factor $`G_M^{}`$ decreases with $`k^2`$ faster than the dipole formula.
We parametrize the inelastic magnetic form factor $`G_M^{}(k^2)G(k^2)`$ for the $`N\mathrm{\Delta }`$ electromagnetic transition with the help of the following formula:
$$G(k^2)=\frac{G(0)G_d(k^2)}{(1k^2/m_x^2)}.$$
Using the last experimental data about the ratio $`G(k^2)/3G_d`$ we find $`m_x^2=5.75`$ (GeV/c)<sup>2</sup>, in agreement with previous estimates .
Note in this connection, that the new JLab data about the electric proton form factor $`G_{Ep}(k^2)`$ show also a deviation from the dipole formula, with a similar value of the parameter $`m_x`$.
In order to calculate the amplitudes $`f_i`$, $`i=1,\mathrm{..6}`$, for the elementary processes $`e^{}+Ne^{}+N+\pi ^0`$, $`N=p`$ or $`n`$, we will use a model similar to but with the following modifications:
* we introduce a term describing the exchange of $`\omega `$-meson in the $`t`$channel;
* the $`s`$-channel contribution of the $`\mathrm{\Delta }`$ isobar is parametrized in such a form to avoid any off-mass shell effects (such as the admixture of $`1/2^\pm `$ or $`3/2^{}`$ states).
* the $`u`$channel of the $`\mathrm{\Delta }`$isobar is neglected.
In order to justify the last option, let us note the essential difference between the $`u`$channel contributions of $`N`$ and $`\mathrm{\Delta }`$. The necessity to introduce the $`u`$-channel contribution from the proton exchange in the process $`\gamma ^{}+pp+\pi ^0`$ is dictated by the gauge invariance of the electromagnetic interaction. As a byproduct, it derives the crossing symmetry for the resulting $`s+u`$ proton exchange. In case of $`\mathrm{\Delta }`$-exchange, there is a different situation with respect to the above mentioned symmetry properties: the gauge invariance and the crossing symmetry. Due to the non-diagonality of the electromagnetic transition $`\gamma ^{}+N\mathrm{\Delta }`$, it is possible to parametrize this vertex in a gauge invariant form independently from the virtuality of the $`\mathrm{\Delta }`$. Therefore, the $`\mathrm{\Delta }`$-contribution only in the $`s`$-channel, is gauge invariant, independently from the $`u`$-channel $`\mathrm{\Delta }`$-contribution. This means that for the $`\mathrm{\Delta }`$-contribution there is no direct connection between the gauge invariance and the crossing symmetry, as for the proton exchange. Moreover, even the $`\mathrm{\Delta }`$-contribution in $`s`$ and $`u`$\- channels simultaneously will not induce crossing symmetry. Namely due to the presence of the $`\mathrm{\Delta }`$-pole in the physical region of s-channel, it is necessary to introduce the $`\mathrm{\Delta }`$width in the corresponding propagator- with resulting complex amplitudes, whereas the $`u`$channel $`\mathrm{\Delta }`$-contribution is characterized by real amplitudes. It turns out that we do not have exact crossing symmetry for the $`\mathrm{\Delta }`$-contributions, even for the sum of $`u`$-and $`s`$ diagrams with $`\mathrm{\Delta }`$-exchange.
We will consider only the $`s`$ channel $`\mathrm{\Delta }`$-contribution. In order to avoid problems with off-mass shell effects, we write the matrix element for the $`\mathrm{\Delta }`$-contribution following the two-component formalism for the description of the spin structure of both vertices, $`\mathrm{\Delta }N+\pi `$ and $`\gamma ^{}+N\mathrm{\Delta }`$. Therefore we can write:
$$\gamma +N\mathrm{\Delta }:\stackrel{}{e}\times \stackrel{}{k}\stackrel{}{\chi }^{}\chi _1,\text{M1 transition, only!}$$
$$\mathrm{\Delta }N+\pi :\chi _2^{}\stackrel{}{\chi }\stackrel{}{q},$$
where $``$ is the identity matrix. Each component of the vector $`\stackrel{}{\chi }`$ is a 2-component spinor, satisfying the condition $`\stackrel{}{\sigma }\stackrel{}{\chi }=0`$, in order to avoid any spin 1/2 contribution. Using for the $`\mathrm{\Delta }`$ density matrix the following expression:
$$\rho _{ab}=\frac{2}{3}(\delta _{ab}\frac{i}{2}ϵ_{abc}\sigma _c),$$
we can write the matrix element for the $`\mathrm{\Delta }`$-contribution in the CMS of $`\gamma ^{}+NN+\pi ^0`$ as follows:
$$_\mathrm{\Delta }=\frac{eG(k^2)|\stackrel{}{q}|}{M_\mathrm{\Delta }^2si\mathrm{\Gamma }_\mathrm{\Delta }M_\mathrm{\Delta }}\chi _2^{}(2i\stackrel{}{e}\widehat{\stackrel{}{k}}\times \widehat{\stackrel{}{q}}+cos\theta _\pi \stackrel{}{\sigma }\stackrel{}{e}\stackrel{}{\sigma }\widehat{\stackrel{}{k}}\stackrel{}{e}\widehat{\stackrel{}{q}})\chi _1\sqrt{(E_1+m)(E_2+m)},$$
(52)
where $`M_\mathrm{\Delta }(\mathrm{\Gamma }_\mathrm{\Delta })`$ is the mass (width) of $`\mathrm{\Delta }`$.
The following $`\mathrm{\Delta }`$ contributions to the scalar amplitudes, $`f_{i\mathrm{\Delta }},i=16`$, can be derived:
$$f_{1\mathrm{\Delta }}=2\mathrm{\Pi }(s,k^2),$$
$$f_{2\mathrm{\Delta }}=\mathrm{cos}\theta _\pi \mathrm{\Pi }(s,k^2),$$
$$f_{3\mathrm{\Delta }}=\mathrm{\Pi }(s,k^2),$$
(53)
$$f_{4\mathrm{\Delta }}=f_{5\mathrm{\Delta }}=f_{6\mathrm{\Delta }}=0,$$
where we use the notation:
$$\mathrm{\Pi }(s,k^2)=\frac{G(k^2)|\stackrel{}{q}|}{M_\mathrm{\Delta }^2si\mathrm{\Gamma }_\mathrm{\Delta }M_\mathrm{\Delta }}.$$
The normalization constant $`G(0)`$ can be deduced from the value of the total cross section for the reaction $`\gamma +pp+\pi ^0`$ (with real photons) at $`s=M_\mathrm{\Delta }^2`$:
$$\sigma _T(\gamma pp\pi ^0)=\frac{\alpha }{2}G^2(0)\frac{|\stackrel{}{q}|^3}{|\stackrel{}{k}|}\frac{(E_{1\mathrm{\Delta }}+m)(E_{2\mathrm{\Delta }}+m)}{M_\mathrm{\Delta }^4\mathrm{\Gamma }_\mathrm{\Delta }^2},$$
where
$$E_{1\mathrm{\Delta }}=\frac{M_\mathrm{\Delta }^2+m^2}{2M_\mathrm{\Delta }},E_{2\mathrm{\Delta }}=\frac{M_\mathrm{\Delta }^2+m^2m_\pi ^2}{2M_\mathrm{\Delta }},$$
$$|\stackrel{}{k}|=\frac{M_\mathrm{\Delta }^2m^2}{2M_\mathrm{\Delta }},|\stackrel{}{q}|=\sqrt{E_{2\mathrm{\Delta }}^2m^2}.$$
Using the spin structure of the resonance amplitude (27), we obtain the following structure for the resonance contribution to the matrix element of the process $`\gamma ^{}+dd+\pi ^0`$:
$$_\mathrm{\Delta }(\gamma ^{}dd\pi )=\frac{1}{2}\mathrm{\Pi }(s,k^2)\{2\mathrm{sin}\theta \stackrel{}{e}\widehat{\stackrel{}{n}}[F_1\left(\stackrel{}{Q^2}\right)\stackrel{}{D_1}\stackrel{}{D_2}^{}+F_2\left(\stackrel{}{Q^2}\right)(3\stackrel{}{D_1}\widehat{\stackrel{}{Q}}\stackrel{}{D_2}^{}\widehat{\stackrel{}{Q}}\stackrel{}{D_1}\stackrel{}{D_2}^{})]$$
$$+\left[\left(\stackrel{}{e}\widehat{\stackrel{}{m}}\widehat{\stackrel{}{m}}\stackrel{}{D_1}\times \stackrel{}{D_2}^{}+\stackrel{}{e}\widehat{\stackrel{}{n}}\widehat{\stackrel{}{n}}\stackrel{}{D_1}\times \stackrel{}{D_2}^{}\right)\mathrm{cos}\theta \stackrel{}{e}\widehat{\stackrel{}{q}}\widehat{\stackrel{}{k}}\stackrel{}{D_1}\times \stackrel{}{D_2}^{}\right]\left(F_3\left(\stackrel{}{Q^2}\right)+F_4\left(\stackrel{}{Q^2}\right)\right)$$
$$3\mathrm{cos}\theta F_4\left(\stackrel{}{Q^2}\right)\stackrel{}{e}\widehat{\stackrel{}{m}}\widehat{\stackrel{}{Q}}\stackrel{}{D_1}\times \stackrel{}{D_2}^{}Q_m+3F_4\left(\stackrel{}{Q^2}\right)\stackrel{}{e}\widehat{\stackrel{}{q}}\stackrel{}{Q}\stackrel{}{D_1}\times \stackrel{}{D_2}^{}Q_k\}.$$
Taking into account only the S-wave component of the deuteron wave function it is possible to predict the $`\theta `$ dependence for the simplest polarization observables for $`\gamma ^{}+dd+\pi ^0`$:
$$\left(H_{xx}^{(0)}H_{yy}^{(0)}\right)/\left(H_{xx}^{(0)}+H_{yy}^{(0)}\right)=3\frac{\mathrm{sin}^2\theta }{32\mathrm{cos}^2\theta },H_{xz}^{(0)}=H_{zz}^{(0)}=0.$$
and in the case of tensor polarized deuterons :
$`\left(H_{xx}^{(2)}+H_{yy}^{(2)}\right)/\left(H_{xx}^{(0)}+H_{yy}^{(0)}\right)_0`$ $`=`$ $`Q_{zz}{\displaystyle \frac{\mathrm{cos}^2\theta }{4\left(32\mathrm{cos}^2\theta \right)}},`$
$`\left(H_{xx}^{(2)}H_{yy}^{(2)}\right)/\left(H_{xx}^{(0)}+H_{yy}^{(0)}\right)_0`$ $`=`$ $`\left(Q_{xx}Q_{yy}\right){\displaystyle \frac{\mathrm{cos}^2\theta }{4\left(32\mathrm{cos}^2\theta \right)}}\text{.}`$
For comparison, note that in the case of the process $`e+pe+p+\pi ^0`$ we have (for an unpolarized proton target):
$`\left(H_{xx}^{(0)}H_{yy}^{(0)}\right)/\left(H_{xx}^{(0)}+H_{yy}^{(0)}\right)={\displaystyle \frac{5\mathrm{sin}^2\theta }{53\mathrm{cos}^2\theta }}.`$
The matrix element $`_\omega `$ for the $`\omega `$-exchange in $`\gamma ^{}+NN+\pi ^0`$ can be written in the following form:
$$_\omega =\frac{g_\omega G_\omega (k^2)}{m_\omega (tm_\omega ^2)}ϵ_{\mu \nu \rho \sigma }e_\mu k_\nu q_\sigma u(p_2)\left[\gamma _\rho \frac{\kappa _\omega }{2m}\sigma _{\rho \beta }(kq)_\beta \right]u(p_1)$$
The constants $`\kappa _\omega `$ and $`g_\omega `$ are fixed by the Bonn potential : $`\kappa _\omega =0,`$ $`g_\omega ^2/4\pi =20.`$ The VDM suggests the following parametrization for the form factor $`G_\omega (k^2)`$:
$$G_\omega (k^2)=\frac{G_\omega (0)}{1k^2/m_\rho ^2}.$$
The value $`G_\omega (0)`$ can be fixed by the width of the radiative decay $`\omega \pi \gamma `$, through the following formula:
$$\mathrm{\Gamma }(\omega \pi \gamma )=\frac{\alpha }{24}G_\omega ^2(0)\left(1\frac{m_\pi ^2}{m_\omega ^2}\right)^3m_\omega ,$$
where $`BR(\omega \pi ^0\gamma )=\mathrm{\Gamma }(\omega \pi ^0\gamma )/\mathrm{\Gamma }_\omega =(8.5\pm 1.5)\%`$, $`\mathrm{\Gamma }_\omega =(8.81\pm 0.09)`$ MeV and $`m_\omega =782`$ MeV.
Concerning vector meson exchange in $`e^{}+Ne^{}+N+\pi `$, it is known , that the vector meson exchange is important for the processes $`\gamma +NN+\pi `$, in the considered region of $`W`$. Due to the isovector nature of the electromagnetic current in $`\gamma ^{}+dd+\pi ^0`$, the $`\rho ^0`$-contribution to $`\gamma ^{}+NN+\pi ^0`$ is exactly cancelled. The VDM parametrization of the electromagnetic form factors suggested above for the $`\gamma ^{}\pi \omega `$-vertex as to be considered as a simplified possibility for the space-like region of momentum transfer, where there is no experimental information. However, in the region of time-like momentum transfer, different pieces of information exist. Let us mention three of them. The decay $`\omega \pi +\mathrm{}^++\mathrm{}^{}`$ allows to measure this form factor in the following region $`4m_{\mathrm{}}k^2(m_\omega m_\pi )^2,`$ where $`m_{\mathrm{}}`$ is the lepton mass. The process $`e^++e^{}\pi ^0+\omega `$ is driven by the considered form factor in another time-like region, namely for $`k^2(m_\omega +m_\pi )^2`$. For completeness we mention the $`\tau ^{}\nu _\tau +\pi ^{}+\omega `$ decay . The presence of the same factor $`G_\omega (k^2)`$ in processes so different as $`e^++e^{}\pi ^0+\omega `$ and $`\tau \nu _\tau +\pi ^{}+\omega `$ results from the well known CVC hypothesis (Conservation of Vector Current for the weak semileptonic processes).
Note also that we have taken a ’hard’ expression for the VDM form factor, $`G_\omega (k^2)`$, which is assumed to reproduce at best the structure function $`A(k^2)`$ of elastic $`ed`$ scattering, through the calculations of the meson exchange current due to $`\pi \rho `$ exchange. However this conclusion is correlated to the properties of the nucleon form factor, especially with the behavior of the isoscalar electric form factor, $`G_{Es}=(G_{Ep}+G_{En})/2.`$ New $`G_{Ep}`$ data (with large deviation from the previously assumed dipole behavior) will also favor a hard form factor $`G_\omega (k^2)`$ for the good description of the $`k^2`$ dependence of $`A(k^2)`$ at large momentum transfer. However a satisfactory description will depend also on the large $`k^2`$-dependence of the neutron electric form factors, which will be measured in the next future up to $`\left|k^2\right|=2`$ (GeV/c)<sup>2</sup> . It is then expected that the different observables in the processes $`e+Ne+N+\pi ^0`$ and $`e+de+d+\pi ^0`$ at relatively large momentum transfer are sensitive to the parametrizations of the form factor $`G_\omega (k^2)`$. For example, the VDM parametrization for $`G_\omega (k^2)`$ shows that this form factor is ’harder’ in comparison with nucleon and $`N\mathrm{\Delta }`$ form factors. Therefore, in this case, the relative role of $`\omega `$-exchange will be essentially increased at large momentum transfer.
## VI Results and discussion
In order to test the model for $`\pi ^0`$ electroproduction on deuterons, we compared our calculation to experimental data on $`\pi ^0`$ and $`\pi ^+`$ photoproduction on proton in the $`\mathrm{\Delta }`$-resonance region. The angular distributions at different energies of the real photon reproduce quite well the existing data, a sample of which is shown in Fig. 8. This agreement justifies the generalization of the model in case of $`\pi ^0`$-electroproduction on nucleons, $`e^{}+Ne^{}+N+\pi ^0`$, by introducing the corresponding electromagnetic form factors in the different photon-hadron vertices (see Fig. 7). Note also, that the resulting electromagnetic current for the process $`\gamma ^{}+NN+\pi ^0`$ (with virtual photon) still satisfies the gauge invariance, for any parametrization of the electromagnetic form factors, and for any values of the kinematical variables $`k^2`$, $`W`$ and $`\mathrm{cos}\theta _\pi `$. However this model does not satisfy the T-invariance of the electromagnetic interaction, but here we will consider only T-even observables, such as the different contributions to the $`d(e,e\pi ^0)d`$ differential cross section ( with unpolarized particles in the initial and final states). This problem, which is common to all modern approaches of pion photo- and electro-production on nucleons, is generally not discussed in the existing literature.
In the framework of IA , as it was shown before, the deuteron structure is described by by four inelastic form factors $`F_i(\stackrel{}{Q}^2),i=14`$, where the argument $`\stackrel{}{Q^2}`$ depends on all the three kinematical variables, $`k^2`$, $`W`$ and $`\mathrm{cos}\theta _\pi `$, which characterize the process $`\gamma ^{}+NN+\pi `$:
$$\stackrel{}{Q}^2=(\stackrel{}{k}\stackrel{}{q})^2=\stackrel{}{k}^2+\stackrel{}{q^2}2|\stackrel{}{k}||\stackrel{}{q}|\mathrm{cos}\theta _\pi ,$$
with
$$\stackrel{}{k}^2=k_0^2k^2,k_0=\frac{W^2+k^2m^2}{2W},$$
$$\stackrel{}{q}^2=E_\pi ^2m_\pi ^2,E_\pi =\frac{W^2+m_\pi ^2m^2}{2W}.$$
Fig. 9 illustrates the dependence of the variable $`\stackrel{}{Q}^2`$ on $`\mathrm{cos}\theta _\pi `$ at fixed values of $`k^2`$ and $`W`$, at W=1.2 GeV and W=1.137 GeV (which corresponds to $`E_\gamma =220`$ MeV, see Fig. 8). This dependence is similar for all values of $`k^2`$, in the interval $`\left|k^2\right|=0.5÷2.0`$ (GeV/c)<sup>2</sup>. Note that $`\stackrel{}{Q}_{max}^23`$ (GeV/c)<sup>2</sup> at $`k^2=2`$ (GeV/c)<sup>2</sup> , so, at the same value of four momentum transfer, the process $`\gamma ^{}+dd+\pi ^0`$ is driven by the deuteron form factors at higher momentum transfer in comparison with elastic $`ed`$-scattering.
Comparing Fig. 9 and Fig. 6 (which shows the $`\stackrel{}{Q}^2`$-dependence of the deuteron form factors, in the interval $`0\stackrel{}{Q^2}3`$ (GeV/c)<sup>2</sup>), one can see that in the range $`k^2=0.5÷2.0`$ (GeV/c)<sup>2</sup>, the deuteron form factors are very sensitive to the behavior of the deuteron wave function calculated in different NN-potentials.
The $`\theta _\pi `$-dependence of all four contributions to the inclusive $`d(e,e\pi ^0)d`$ cross section, namely $`H_{xx}\pm H_{yy}`$, $`H_{zz}`$ and $`H_{xz}+H_{zx}`$, for different values of $`k^2`$ and $`W`$ is shown in Figs. 10 and 11 Note that in our normalization, Eq. (2), all components $`H_{ab}`$ are dimensionless numbers.. In order to show the relative role of the different mechanisms for the elementary processes $`\gamma ^{}+NN+\pi `$ ( in the considered kinematical region for the variables $`k^2`$ and $`W`$), each picture shows four curves: $`\mathrm{\Delta }`$ contribution only, $`\mathrm{\Delta }+s+u`$ (nucleon diagrams) and $`\mathrm{\Delta }+s+u\pm \omega `$. The calculations are shown for both relative signs of the vector meson contribution in order to stress the importance of the $`\omega `$ contribution. The positive sign has been choosen from the comparison with experimental data on $`\gamma +pp+\pi ^0`$ (real photons).
The $`\omega `$ contribution is important for all the four considered observables, in particular for the $`H_{xx}\pm H_{yy}`$ terms at $`\theta _\pi 80^o`$; in the case of $`H_{zz}`$ the largest sensitivity appears for backward $`\pi ^0`$ electroproduction.
The relative role of the absorption of virtual photon with longitudinal and transversal polarizations depends essentially on the variables $`k^2`$ and $`W`$, with an increase of the ratio $`H_{zz}/(H_{xx}+H_{yy})`$ with $`k^2`$. At W=1.2 GeV, where the $`\mathrm{\Delta }`$-contribution (with absorption of transversal virtual photons) dominates, the relative role of $`H_{zz}`$ is weaker in comparison with $`H_{xx}+H_{yy}`$. However for $`k^21`$ (GeV/c)<sup>2</sup> $`H_{zz}`$ exceeds $`H_{xx}+H_{yy}`$, even in the resonance region.
The ratio $`(H_{xx}H_{yy})/(H_{xx}+H_{yy})`$ is negative (due to the dominance of the transversal $`\mathrm{\Delta }`$ and $`\omega `$-contributions) and has a $`sin^2\theta _\pi `$ behavior. The longitudinal-transversal interference contribution, $`H_{xz}+H_{zx}`$, shows a particular sensitivity to the different ingredients of the model, with strong $`\theta _\pi `$-dependence, in the whole considered kinematical domain.
In view of the importance of the $`\omega `$ contribution to all observables for the $`d(e,e\pi ^0)d`$ process, we studied the sensitivity to the choice of the electromagnetic $`\gamma ^{}\omega \pi `$-vertex form factor. For this aim we used two parametrizations, a $`hard`$ monopole form, $`G_\omega ^{(h)}(k^2)`$, predicted by the standard VDM, and a $`soft`$ dipole form $`G_\omega ^{(s)}(k^2)`$:
$$G_\omega ^{(h)}(k^2)=\frac{G_\omega (0)}{1{\displaystyle \frac{k^2}{m_\rho ^2}}},G_\omega ^{(s)}(k^2)=\frac{G_\omega (0)}{\left(1{\displaystyle \frac{k^2}{m_\rho ^2}}\right)^2}.$$
Fig. 12 shows the $`\theta _\pi `$-dependence of the following ratios:
$$r_\pm (\mathrm{cos}\theta _\pi )=\frac{(H_{xx}\pm H_{yy})_{hard}(H_{xx}\pm H_{yy})_{soft}}{(H_{xx}\pm H_{yy})_{hard}+(H_{xx}\pm H_{yy})_{soft}}$$
for two different values of $`k^2`$ ($`k^2=0.5`$ and 2 (GeV/c)<sup>2</sup>) and $`W=1.137`$. For $`W=1.2`$ GeV (Fig. 13) the largest sensitivity to the choice of the form factor $`G_\omega (k^2)`$ appears at forward angles for $`\pi ^0`$-production, whereas at $`W=1.137`$ GeV all angles are equally sensitive to this choice. At the $`\mathrm{\Delta }`$-resonance this sensitivity increases slightly with $`k^2`$.
The absolute measurements of the different contributions to the inclusive cross section for $`d(e,e\pi ^0)d`$ will help in defining the appropriate $`k^2`$-dependence of the form factor $`G_\omega (k^2)`$. However , as we can see on Fig. 14, the absolute values of the $`H_{xx}\pm H_{yy}`$ contributions, the shape and absolute values of $`H_{zz}`$ and $`H_{xz}+H_{zx}`$ are also sensitive to the existing $`NN`$potentials, in particular at large $`k^2`$. In Figs 15, 16, 17 and 18, we illustrate the behavior of the four observables, for different parametrizations of the following ingredients:
* the deuteron wave function: for the Bonn and Paris potentials,
* the electromagnetic form factors for the $`\gamma ^{}\pi \omega `$-vertex: $`hard`$ (VDM) and $`soft`$ (dipole) parametrizations;
* the electromagnetic form factor of the proton: dipole or a ’softer’ parametrization based on recent data on the proton electric form factor.
The differences between the different parametrizations increase at large momentum transfer.
The inclusive cross section for $`d(e,e)\pi ^0d`$ is characterized by two contributions, only. After integration over $`d\mathrm{\Omega }_\pi `$, we have:
$$H_t(k^2,W)=_1^{+1}d\mathrm{cos}\theta _\pi (H_{xx}+H_{yy}),$$
$$H_{\mathrm{}}(k^2,W)=_1^{+1}d\mathrm{cos}\theta _\pi H_{zz}.$$
The three-dimensional plot of Fig. 19 shows the dependence of these inclusive functions, on $`k^2`$ and $`W`$. The calculation is done here, for the $`hard`$ form factor $`G_\omega `$, the dipole form factor $`G_{Ep}`$ and the Bonn deuteron wave function.
## VII Conclusions
We have made a general analysis of coherent pseudoscalar neutral mesons production on deuterons, $`e+de+d+P^0`$, which holds for any kinematics of the discussed processes. Threshold $`P^0`$-meson production (at any value of momentum transfer square $`k^2`$ and for the minimum value of the effective mass of the produced hadronic system) is especially interesting due to the essential simplification of the spin structure of the corresponding amplitudes and to the decreasing number of independent kinematical variables. Another kinematical region, which is interesting for the process $`\gamma ^{}+dd+\pi ^0`$, is the $`\mathrm{\Delta }`$-isobar excitation on the nucleons.
Coherent $`P^0`$-meson production is interesting due to its special sensitivity to the isotopic structure of the threshold amplitude for the elementary processes $`\gamma ^{}+NN+P^0`$.
The $`\pi ^0`$-meson electroproduction on the deuteron allows to measure the threshold amplitude for $`\gamma ^{}+nn+\pi ^0`$, which is important for testing hadron electrodynamics .
The $`\eta `$-meson electroproduction on the deuteron could be important for the study of $`\eta N`$\- and $`\eta d`$-interactions, in particular after the finding of a strong energy dependence of the cross section of $`n+pd+\eta `$ process near threshold.
The IA can be considered as a good starting point for the discussion of corrections such as mesonic exchange currents, isobar configurations in deuteron, quark degrees of freedom, etc., but rescattering effects will also have to be discussed, in particular for $`\eta `$-production near threshold.
Using an adequate model for the elementary processes of $`\pi ^0`$-electroproduction on nucleons, $`e^{}+Ne^{}+N+\pi ^0`$, which satisfactorily reproduces the angular dependence of the differential cross section for the processes $`\gamma +pp+\pi ^0`$ and $`\gamma +pn+\pi ^+`$ (in the $`\mathrm{\Delta }`$-resonance region), we estimated the four standard contributions to the exclusive differential cross section for the reaction $`d(e,e\pi ^0)d`$ as a function of the variables $`k^2,W`$ and $`\theta _\pi `$. These calculations were done at relatively large momentum transfer square, $`k^2=0.5÷2.0`$ (GeV/c)<sup>2</sup>, where recent data exist. All observables show a large sensitivity to the parametrization of electromagnetic form factors, in the considered model. A special attention was devoted to the study of the effects of soft and hard parametrizations of form factor for the $`\pi \omega \gamma ^{}`$-vertex, as well as to possible deviation of the proton electric form factor from the dipole fit. Moreover, as it is well known for elastic $`ed`$-scattering, we find here, too, a large dependence of all the observables to the choice of $`NN`$potential. The large sensitivity of the $`d(e,e\pi ^0)d`$ cross section to the $`\omega `$-exchange contribution can be used, in principle, to study the corresponding electromagnetic form factors in the space-like momentum transfer region.
Acknowledgments
We thank J.-M. Laget for interesting discussions on rescattering effects. One of the authors (M. P. R.) is very indebted to the hospitality of Saturne where part of this work was done.
Appendix
We present here the expressions for the structure functions $`h_1h_{41}`$ in terms of the scalar amplitudes $`g_1g_{13}`$. The SF’s $`h_1h_5`$ corresponding to the interaction with an unpolarized deuteron target can be written as:
$`3h_1`$ $`=`$ $`\left|g_1\right|^2+\left|g_2\right|^2+\left|g_3\right|^2+\left|g_4\right|^2,`$
$`3h_2`$ $`=`$ $`\left|g_5\right|^2+\left|g_6\right|^2+\left|g_7\right|^2+\left|g_8\right|^2+\left|g_9\right|^2,`$
$`3h_3`$ $`=`$ $`\left|g_{10}\right|^2+\left|g_{11}\right|^2+\left|g_{12}\right|^2+\left|g_{13}\right|^2,`$
$`3h_4`$ $`=`$ $`\left(_\mathcal{1}_{\mathcal{10}}^{}+_\mathcal{2}_{\mathcal{11}}^{}+_\mathcal{3}_{\mathcal{12}}^{}+_\mathcal{4}_{\mathcal{13}}^{}\right),`$
$`3h_5`$ $`=`$ $`𝓂\left(_\mathcal{1}_{\mathcal{10}}^{}+_\mathcal{2}_{\mathcal{11}}^{}+_\mathcal{3}_{\mathcal{12}}^{}+_\mathcal{4}_{\mathcal{13}}^{}\right),`$
We derive fhe following expressions for the $`SF^{^{}}sh_6`$ \- $`h_{18}`$, which characterize the effects of the target vector polarization :
$`h_6`$ $`=`$ $`𝓂\left(_\mathcal{2}_\mathcal{6}^{}_\mathcal{3}_\mathcal{9}^{}_\mathcal{4}_\mathcal{7}^{}\right),`$
$`h_7`$ $`=`$ $`𝓂\left(_\mathcal{6}_{\mathcal{11}}^{}+_\mathcal{7}_{\mathcal{13}}^{}+_\mathcal{9}_{\mathcal{12}}^{}\right),`$
$`h_8`$ $`=`$ $`\left(_\mathcal{2}_\mathcal{6}^{}_\mathcal{3}_\mathcal{9}^{}_\mathcal{4}_\mathcal{7}^{}\right),`$
$`h_9`$ $`=`$ $`\left(_\mathcal{6}_{\mathcal{11}}^{}_\mathcal{7}_{\mathcal{13}}^{}_\mathcal{9}_{\mathcal{12}}^{}\right),`$
$`h_{10}`$ $`=`$ $`2𝓂_\mathcal{1}_\mathcal{2}^{},`$
$`h_{11}`$ $`=`$ $`2𝓂\left(_\mathcal{5}_\mathcal{9}^{}_\mathcal{7}_\mathcal{8}^{}\right),`$
$`h_{12}`$ $`=`$ $`22𝓂_{\mathcal{10}}_{\mathcal{11}}^{},`$
$`h_{13}`$ $`=`$ $`𝓂\left(_\mathcal{1}_{\mathcal{11}}^{}_\mathcal{2}_{\mathcal{10}}^{}\right),`$
$`h_{14}`$ $`=`$ $`\left(_\mathcal{1}_{\mathcal{11}}^{}_\mathcal{2}_{\mathcal{10}}^{}\right),`$
$`h_{15}`$ $`=`$ $`𝓂\left(_\mathcal{1}_\mathcal{6}^{}_\mathcal{3}_\mathcal{5}^{}_\mathcal{4}_\mathcal{8}^{}\right),`$
$`h_{16}`$ $`=`$ $`𝓂\left(_\mathcal{5}_{\mathcal{12}}^{}_\mathcal{6}_{\mathcal{10}}^{}_\mathcal{8}_{\mathcal{13}}^{}\right),`$
$`h_{17}`$ $`=`$ $`\left(_\mathcal{1}_\mathcal{6}^{}_\mathcal{3}_\mathcal{5}^{}_\mathcal{4}_\mathcal{8}^{}\right),`$
$`h_{18}`$ $`=`$ $`\left(_\mathcal{5}_{\mathcal{12}}^{}_\mathcal{6}_{\mathcal{10}}^{}+_\mathcal{8}_{\mathcal{13}}^{}\right),`$
Finally for the $`SF^{^{}}sh_{19}`$ \- $`h_{41}`$, which describe the effects on tensor target polarization, one obtains :
$`3h_{19}`$ $`=`$ $`\left|g_1\right|^2+\left|g_2\right|^2+\left|g_7\right|^2\left|g_8\right|^2,`$
$`3h_{20}`$ $`=`$ $`\left|g_5\right|^2+\left|g_9\right|^2,`$
$`3h_{21}`$ $`=`$ $`\left|g_{10}\right|^2+\left|g_{11}\right|^2,`$
$`3h_{22}`$ $`=`$ $`\left(_\mathcal{1}_{\mathcal{10}}^{}_\mathcal{2}_{\mathcal{11}}^{}\right),`$
$`3h_{23}`$ $`=`$ $`𝓂\left(_\mathcal{1}_{\mathcal{10}}^{}_\mathcal{2}_{\mathcal{11}}^{}\right),`$
$`3h_{24}`$ $`=`$ $`\left|g_2\right|^2\left|g_3\right|^2\left|g_4\right|^2,`$
$`3h_{25}`$ $`=`$ $`\left|g_6\right|^2+\left|g_7\right|^2+\left|g_9\right|^2,`$
$`3h_{26}`$ $`=`$ $`\left|g_{11}\right|^2\left|g_{12}\right|^2\left|g_{13}\right|^2,`$
$`3h_{27}`$ $`=`$ $`\left(_\mathcal{2}_{\mathcal{11}}^{}_\mathcal{3}_{\mathcal{12}}^{}_\mathcal{4}_{\mathcal{13}}^{}\right),`$
$`3h_{28}`$ $`=`$ $`𝓂\left(_\mathcal{2}_{\mathcal{11}}^{}_\mathcal{3}_{\mathcal{12}}^{}_\mathcal{4}_{\mathcal{13}}^{}\right),`$
$`3h_{29}`$ $`=`$ $`2_\mathcal{1}_\mathcal{2}^{},`$
$`3h_{30}`$ $`=`$ $`2\left(_\mathcal{5}_\mathcal{9}^{}+_\mathcal{7}_\mathcal{8}^{}\right),`$
$`3h_{31}`$ $`=`$ $`2_{\mathcal{10}}_{\mathcal{11}}^{},`$
$`3h_{32}`$ $`=`$ $`\left(_\mathcal{1}_{\mathcal{11}}^{}+_\mathcal{2}_{\mathcal{10}}^{}\right),`$
$`3h_{33}`$ $`=`$ $`𝓂\left(_\mathcal{1}_{\mathcal{11}}^{}+_\mathcal{2}_{\mathcal{10}}^{}\right),`$
$`3h_{34}`$ $`=`$ $`\left(_\mathcal{1}_\mathcal{6}^{}+_\mathcal{3}_\mathcal{5}^{}+_\mathcal{4}_\mathcal{8}^{}\right),`$
$`3h_{35}`$ $`=`$ $`\left(_\mathcal{5}_{\mathcal{12}}^{}+_\mathcal{6}_{\mathcal{10}}^{}+_\mathcal{8}_{\mathcal{13}}^{}\right),`$
$`3h_{36}`$ $`=`$ $`𝓂\left(_\mathcal{1}_\mathcal{6}^{}+_\mathcal{3}_\mathcal{5}^{}+_\mathcal{4}_\mathcal{8}^{}\right),`$
$`3h_{37}`$ $`=`$ $`𝓂\left(_\mathcal{5}_{\mathcal{12}}^{}+_\mathcal{6}_{\mathcal{10}}^{}+_\mathcal{8}_{\mathcal{13}}^{}\right),`$
$`3h_{38}`$ $`=`$ $`\left(_\mathcal{2}_\mathcal{6}^{}+_\mathcal{3}_\mathcal{9}^{}+_\mathcal{4}_\mathcal{7}^{}\right),`$
$`3h_{39}`$ $`=`$ $`\left(_\mathcal{6}_{\mathcal{11}}^{}+_\mathcal{7}_{\mathcal{13}}^{}+_\mathcal{9}_{\mathcal{12}}^{}\right),`$
$`3h_{40}`$ $`=`$ $`𝓂\left(_\mathcal{2}_\mathcal{6}^{}+_\mathcal{3}_\mathcal{9}^{}+_\mathcal{4}_\mathcal{7}^{}\right),`$
$`3h_{41}`$ $`=`$ $`𝓂\left(_\mathcal{6}_{\mathcal{11}}^{}+_\mathcal{9}_{\mathcal{12}}^{}+_\mathcal{7}_{\mathcal{13}}^{}\right).`$
|
warning/0003/nucl-th0003024.html
|
ar5iv
|
text
|
# Kaon-baryon coupling constants in the QCD sum rule approach
## I Introduction
To understand kaon-nuclear physics, it is important to know the hadronic coupling strengths involving the kaons. Among them, $`g_{KN\mathrm{\Lambda }}`$ and $`g_{KN\mathrm{\Sigma }}`$ are the most relevant coupling constants. In contrast to $`g_{\pi NN}`$, however, the determination of these kaon couplings has some difficulties both in the experimental side and in the theoretical side, e.g. see .
Among other theoretical approaches, QCD sum rule method has been used to extract these kaon couplings. However, compared to the large number of works devoted to $`g_{\pi NN}`$, there have been only few QCD sum rule estimates on $`g_{KN\mathrm{\Lambda }}`$ and $`g_{KN\mathrm{\Sigma }}`$, for which there are still ambiguities in among the calculations. Thus the results are quite different from each other. More detailed analyses are needed both experimentally and theoretically to understand this discrepancy, and to understand kaon-nuclear physics.
In Ref. , the OPE was calculated only up to the leading term coming from the quark condensate and to leading order in $`m_s`$ in the sum rule structure proportional to $`\text{/}qi\gamma _5`$. However, the next leading term, dimension 5 $`\overline{q}g_s\sigma Gq`$ may contribute to the OPE side with considerable amount as in nucleon mass sum rule . In addition, operators of dimension 7 may also be important in the OPE side as a further power correction. Thus, in this paper we re-analyze our QCD sum rule calculation including higher dimensional condensates, such as $`\overline{q}g_s\sigma Gq`$ and $`\overline{q}q\frac{\alpha _s}{\pi }G^2`$, and study the contribution of these condensates on the previous results.
In Sec. II we present our sum rules for $`g_{KN\mathrm{\Lambda }}`$ and $`g_{KN\mathrm{\Sigma }}`$, and Sec. III we discuss some uncertainties in our sum rules and summarize our results.
## II QCD sum rules for $`g_{KN\mathrm{\Lambda }}`$ and $`g_{KN\mathrm{\Sigma }}`$
We will closely follow the procedures given in Refs.. Consider the three point function constructed of the two baryon currents $`\eta _B`$, $`\eta _B^{}`$ and the pseudoscalar meson current $`j_5`$.
$`A(p,p^{},q)={\displaystyle 𝑑x𝑑y0|T(\eta _B^{}(x)j_5(y)\overline{\eta }_B(0))|0e^{i(p^{}xqy)}}.`$ (1)
In order to obtain $`g_{KN\mathrm{\Lambda }}`$, we will use the following currents for the nucleon and the $`\mathrm{\Lambda }`$ .
$`\eta _N`$ $`=`$ $`ϵ_{abc}(u_a^TC\gamma _\mu u_b)\gamma _5\gamma ^\mu d_c,`$ (2)
$`\eta _\mathrm{\Lambda }`$ $`=`$ $`\sqrt{{\displaystyle \frac{2}{3}}}ϵ_{abc}\left[(u_a^TC\gamma _\mu s_b)\gamma _5\gamma ^\mu d_c(d_a^TC\gamma _\mu s_b)\gamma _5\gamma ^\mu u_c\right],`$ (3)
where u and d are the up and down quark fields ($`a,b`$ and $`c`$ are color indices), $`T`$ denotes the transpose in Dirac space, and $`C`$ is the charge conjugation matrix. For the $`K^{}`$ we choose the current
$`j_K^{}=\overline{s}i\gamma _5u.`$ (4)
The general expression for $`A(p,p^{},q)`$ has the following form
$`A(p,p^{},q)`$ $`=`$ $`F_1(p^2,p^2,q^2)i\gamma _5+F_2(p^2,p^2,q^2)\text{/}qi\gamma _5`$ (5)
$`+`$ $`F_3(p^2,p^2,q^2)\text{/}Pi\gamma _5+F_4(p^2,p^2,q^2)\sigma ^{\mu \nu }\gamma _5q_\mu p_\nu ^{},`$ (6)
where $`q=p^{}p`$ and $`P=\frac{p+p^{}}{2}`$. Recently, in Ref. it was reported that in the case of $`g_{\pi NN}`$ the $`\sigma ^{\mu \nu }\gamma _5`$ structure is independent of the effective models employed in the phenomenological side and further provides the $`\pi NN`$ coupling with less uncertainties from QCD parameters. Motivated by this result $`g_{KN\mathrm{\Lambda }}`$ and $`g_{KN\mathrm{\Sigma }}`$ were calculated from this structure in Refs. . In this paper, however, we construct the sum rule for only the $`\text{/}qi\gamma _5`$ structure as before, and compare this with our previous one.
On the phenomenological side, keeping the first two terms we have
$`\lambda _N\lambda _\mathrm{\Lambda }{\displaystyle \frac{M_B}{(p^2M_N^2)(p^2M_\mathrm{\Lambda }^2)}}(\text{/}qi\gamma _5)g_{KN\mathrm{\Lambda }}{\displaystyle \frac{1}{q^2m_K^2}}{\displaystyle \frac{f_Km_K^2}{2m_q}}`$ (7)
$`+`$ $`\lambda _N\lambda _\mathrm{\Lambda }^{}{\displaystyle \frac{M_B^{}}{(p^2M_N^2)(p^2M_\mathrm{\Lambda }^{}^2)}}(\text{/}qi\gamma _5)g_{KN\mathrm{\Lambda }^{}}{\displaystyle \frac{1}{q^2m_K^2}}{\displaystyle \frac{f_Km_K^2}{2m_q}}`$ (8)
$`+`$ $`\mathrm{higher}\mathrm{resonances},`$ (9)
where $`M_B=\frac{1}{2}(M_N+M_\mathrm{\Lambda })`$, and $`M_B^{}=\frac{1}{2}(M_NM_\mathrm{\Lambda }^{})`$. Here $`\mathrm{\Lambda }^{}`$ means the $`\mathrm{\Lambda }`$(1405), and we introduce (–) sign for the $`\mathrm{\Lambda }`$ (1405) mass because it is a negative parity state. However, this is not relevant in the following calculation. $`\lambda _N`$, $`\lambda _\mathrm{\Lambda }`$ and $`\lambda _\mathrm{\Lambda }^{}`$ are the coupling strengths of the baryons to their currents. $`m_q`$ is the average of the quark masses, $`f_K`$ the kaon decay constant and $`m_K`$ the kaon mass. We take $`f_K`$ = 0.160 GeV and $`m_s`$ = 0.150 GeV.
As for the OPE side, the new contribution from the quark-gluon condensates is given by
$`\sqrt{{\displaystyle \frac{2}{3}}}{\displaystyle \frac{7}{2^43\pi ^2}}\mathrm{ln}(p^2)(\overline{q}g_s\sigma Gq+\overline{s}g_s\sigma Gs),`$ (10)
and from dimension 7 ops.
$`+\sqrt{{\displaystyle \frac{2}{3}}}{\displaystyle \frac{5}{2^33^2}}{\displaystyle \frac{1}{p^2}}(\overline{q}q+\overline{s}s){\displaystyle \frac{\alpha _s}{\pi }}G^2,`$ (11)
where we take the limit $`p^2p^2`$ and let $`\overline{u}u=\overline{d}d\overline{q}q`$, $`\overline{u}g_s\sigma Gu=\overline{d}g_s\sigma Gd\overline{q}g_s\sigma Gq`$. Here we collect only the terms which contribute to the $`\text{/}q/q^2`$ structure such as Figs. 1 and 2. Using the standard values for $`\overline{s}g_s\sigma Gs=0.8\overline{q}g_s\sigma Gq`$ and $`\overline{q}g_s\sigma Gq=m_0^2\overline{q}q=0.8\overline{q}q`$ the sum rule after Borel transformation to $`p^2=p^2`$ becomes
$`\lambda _N\lambda _\mathrm{\Lambda }{\displaystyle \frac{M_B}{M_\mathrm{\Lambda }^2M_N^2}}\left(e^{M_N^2/M^2}e^{M_\mathrm{\Lambda }^2/M^2}\right)g_{KN\mathrm{\Lambda }}{\displaystyle \frac{f_Km_K^2}{2m_q}}+A\left(e^{M_N^2/M^2}e^{M_\mathrm{\Lambda }^{}^2/M^2}\right)=`$ (12)
$`\sqrt{{\displaystyle \frac{2}{3}}}\left({\displaystyle \frac{33}{40\pi ^2}}E_1M^4+({\displaystyle \frac{11m_s^2}{60\pi ^2}}{\displaystyle \frac{21}{100\pi ^2}})E_0M^2+({\displaystyle \frac{m_s}{3}}\overline{s}s+{\displaystyle \frac{1}{8}}{\displaystyle \frac{\alpha _s}{\pi }}G^2)\right)\overline{q}q.`$ (13)
Here, A is the unknown constant coming from $`\lambda _\mathrm{\Lambda }^{}g_{KN\mathrm{\Lambda }^{}}`$, and
$`E_i=1{\displaystyle \underset{k=0}{\overset{i}{}}}{\displaystyle \frac{s_0^k}{k!(M^2)^k}}e^{\frac{s_0}{M^2}},`$ (14)
where $`s_0`$ is a continuum threshold. One should be cautious, however, that there may be non-accounted terms, which can not be neglected by using this simple Borel transformation.
For $`\lambda _N`$ and $`\lambda _\mathrm{\Lambda }`$, we use the values obtained from the following baryon sum rules for the $`N`$ and $`\mathrm{\Lambda }`$ :
$`E_2^NM^6+bE_0^NM^2+{\displaystyle \frac{4}{3}}a^2=2(2\pi )^4\lambda _N^2e^{M_N^2/M^2},`$ (15)
$`E_2^\mathrm{\Lambda }M^6+{\displaystyle \frac{2}{3}}am_s(13\gamma )E_0^\mathrm{\Lambda }M^2+bE_0^\mathrm{\Lambda }M^2+{\displaystyle \frac{4}{9}}a^2(3+4\gamma )=2(2\pi )^4\lambda _\mathrm{\Lambda }^2e^{M_\mathrm{\Lambda }^2/M^2},`$ (16)
where $`a(2\pi )^2\overline{q}q`$, $`b\pi ^2\frac{\alpha _s}{\pi }G^2`$, and $`\gamma \overline{s}s/\overline{q}q10.2`$. We use different thresholds for $`\lambda _N`$ and $`\lambda _\mathrm{\Lambda }`$ in Eqs. (15) and (16). We take $`s_N`$ = (1.440 GeV)<sup>2</sup> for the nucleon sum rule and $`s_\mathrm{\Lambda }`$ = (1.405 GeV)<sup>2</sup> for the $`\mathrm{\Lambda }`$ sum rule considering the next excited nucleon and $`\mathrm{\Lambda }`$ state, respectively.
$`g_{KN\mathrm{\Lambda }}`$ , however, does not display a plateau as a function of the Borel mass. This is because there is no usual power correction term like ($`\frac{a}{M^2},\frac{b}{M^4}`$, and so on) in the r.h.s. of Eq. (13) even including up to dimension 7 operators. We need more higher dimensional operators to get those terms. Thus, in this case we prefer to use a best-fit method. Eq. (13) has the following form :
$`g_{KN\mathrm{\Lambda }}f_1(M^2)+Af_2(M^2)=f_3(M^2).`$ (17)
Then, we get $`g_{KN\mathrm{\Lambda }}`$ and the unknown constant A by minimizing $`(g_{KN\mathrm{\Lambda }}f_1+Af_2f_3)^2`$ with a fixed $`s_0`$ and an appropriate Borel interval:
$`{\displaystyle _{M_{min}^2}^{M_{max}^2}}(g_{KN\mathrm{\Lambda }}f_1+Af_2f_3)^2𝑑M^2=\mathrm{minimum}.`$ (18)
We fix the continuum threshold $`s_0`$ = 2.074 GeV<sup>2</sup> taking into account the next term from the N(1440), i.e., $`N(1440)\mathrm{\Lambda }`$, in the phenomenological side.
The Borel interval $`M^2`$ is restricted by the following conditions : OPE convergence and pole dominance. The lower limit of $`M^2`$, $`M_{min}^2`$ is determined as the value at which the contribution of the highest dimensional operators is less than 10$`\%`$ of total OPE. The upper limit $`M_{max}^2`$ is determined by restricting the continuum contribution to be less than 50$`\%`$. Then, we get
$`|g_{KN\mathrm{\Lambda }}|`$ $`=`$ $`2.49,`$ (19)
$`|A|`$ $`=`$ $`0.00174\mathrm{GeV}^7,`$ (20)
and the Borel interval (0.478, 1.068) GeV<sup>2</sup> for basic inputs (i.e., $`\overline{q}q`$ = – (0.230 GeV)<sup>3</sup>, $`\frac{\alpha _s}{\pi }G^2`$ = 0.012 GeV<sup>4</sup>, $`m_s`$ = 0.150 GeV, and $`m_0^2`$ = 0.8 GeV<sup>2</sup>). Here we denote the absolute value because we can not determine signs of the coupling strengths ($`\lambda _N`$, $`\lambda _\mathrm{\Lambda }`$ and $`\lambda _\mathrm{\Lambda }^{}`$) in the baryon sum rules. We also calculate the average deviation $`\overline{\delta }`$ $``$ $`_i^N\left|1RHS(M_i^2)/LHS(M_i^2)\right|/N`$ = 8.8 $`\times `$ 10<sup>-2</sup> to test the reliability of our fitting, and it shows that the deviation is less than 10 %.
Table I shows variations of $`g_{KN\mathrm{\Lambda }}`$ for other input parameters, which are coming from the uncertainty of the basic inputs. For example, the first line in Table I shows that $`|g_{KN\mathrm{\Lambda }}|`$ = 3.11 (or 1.93) if we change the quark condensate to $`\overline{q}q`$ = –(0.210 GeV)<sup>3</sup> (or –(0.250 GeV)<sup>3</sup>) while other basic inputs are fixed. In the last line we take $`m_0^2`$ = 0.6 GeV<sup>2</sup> from the lowest value of the standard QCD sum rule estimate , and 1.4 GeV<sup>2</sup> which was evaluated in the instanton vacuum in Ref. . Total variation is about $`\pm `$ 1.25 on the above $`g_{KN\mathrm{\Lambda }}`$ value. On the other hand, the unknown constant $`|A|`$ varies from 0.00120 to 0.00203 GeV<sup>7</sup>.
Next, consider $`g_{KN\mathrm{\Sigma }}`$. The current of $`\mathrm{\Sigma }^{}`$ is obtained by making an SU(3) rotation from the nucleon current
$`\eta _\mathrm{\Sigma }=\sqrt{2}ϵ_{abc}\left[(u_a^TC\gamma _\mu s_b)\gamma _5\gamma ^\mu d_c+(d_a^TC\gamma _\mu s_b)\gamma _5\gamma ^\mu u_c\right].`$ (21)
In this case the contribution of the quark-gluon condensate is given by
$`\sqrt{2}{\displaystyle \frac{1}{2^43\pi ^2}}\mathrm{ln}(p^2)(\overline{q}g_s\sigma Gq+\overline{s}g_s\sigma Gs),`$ (22)
and from dimension 7 ops.
$`+\sqrt{2}{\displaystyle \frac{1}{2^33^2}}{\displaystyle \frac{1}{p^2}}(\overline{q}q+\overline{s}s){\displaystyle \frac{\alpha _s}{\pi }}G^2.`$ (23)
Then, within the same approximation as before we get the following sum rule.
$`\lambda _N\lambda _\mathrm{\Sigma }{\displaystyle \frac{M_B}{M_\mathrm{\Sigma }^2M_N^2}}\left(e^{M_N^2/M^2}e^{M_\mathrm{\Sigma }^2/M^2}\right)g_{KN\mathrm{\Sigma }}{\displaystyle \frac{f_Km_K^2}{2m_q}}+B\left(e^{M_N^{}^2/M^2}e^{M_\mathrm{\Sigma }^2/M^2}\right)=`$ (24)
$`+\sqrt{2}\left({\displaystyle \frac{3}{40\pi ^2}}E_1M^4+({\displaystyle \frac{m_s^2}{60\pi ^2}}+{\displaystyle \frac{3}{100\pi ^2}})E_0M^2{\displaystyle \frac{1}{40}}{\displaystyle \frac{\alpha _s}{\pi }}G^2\right)\overline{q}q,`$ (25)
where $`M_B=\frac{1}{2}(M_N+M_\mathrm{\Sigma })`$ and $`N^{}`$ is N(1440). B is the unknown constant coming from $`\lambda _N^{}g_{KN^{}\mathrm{\Sigma }}`$. Again for $`\lambda _\mathrm{\Sigma }`$, we take the value from the following sum rule for the $`\mathrm{\Sigma }`$:
$`E_2^\mathrm{\Sigma }M^62am_s(1+\gamma )E_0^\mathrm{\Sigma }M^2+bE_0^\mathrm{\Sigma }M^2+{\displaystyle \frac{4}{3}}a^2=2(2\pi )^4\lambda _\mathrm{\Sigma }^2e^{M_\mathrm{\Sigma }^2/M^2}.`$ (26)
We fix the continuum threshold $`s_\mathrm{\Sigma }`$ = (1.660 GeV)<sup>2</sup> considering the next $`\mathrm{\Sigma }`$ state, $`\mathrm{\Sigma }`$ (1660).
Using the continuum threshold $`s_0`$ = 2.356 GeV<sup>2</sup> taking into account the next term from the N(1535), i.e., $`N(1535)\mathrm{\Sigma }`$, in the phenomenological side we get
$`|g_{KN\mathrm{\Sigma }}|`$ $`=`$ $`0.395,`$ (27)
$`|B|`$ $`=`$ $`0.00148\mathrm{GeV}^7`$ (28)
for the same basic inputs. The Borel interval is (0.488, 1.584) GeV<sup>2</sup> and the average deviation of the fit $`\overline{\delta }`$ is 9.7 % in this case. We present the variation of $`g_{KN\mathrm{\Sigma }}`$ on other parameters in Table II. The total variation is about $`\pm `$ 0.377. On the other hand, $`|B|`$ varies from 0.00117 to 0.00184 GeV<sup>7</sup>.
## III Discussion
SU(3) symmetry, using de Swart’s convention, predicts
$`g_{KN\mathrm{\Lambda }}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{3}}}(32\alpha _D)g_{\pi NN},`$ (29)
$`g_{KN\mathrm{\Sigma }}`$ $`=`$ $`+(2\alpha _D1)g_{\pi NN},`$ (30)
where $`\alpha _D`$ is the fraction of the D type coupling, $`\alpha _D=\frac{D}{D+F}`$. In Table III we compare our results with previous QCD sum rule estimates and an SU(3) symmetry prediction, where we denote the error-bar allowing for SU(3) symmetry breaking at the 20 % level. Here we take $`\alpha _D`$ from a recent analysis of hyperon semi-leptonic decay data by Ratcliffe, $`\alpha _D`$=0.64, and $`g_{\pi NN}`$ from an analysis of the $`np`$ data by Ericson et al., $`g_{\pi NN}`$=13.43. A comparison to fitting analyses of experimental data is also provided. SU(3) symmetry predicts $`|g_{KN\mathrm{\Lambda }}/g_{KN\mathrm{\Sigma }}|`$ = 3.55 taking $`\alpha _D`$ = 0.64, while our results show that this ratio is 6.30 using the basic inputs, and the order of SU(3) symmetry breaking is rather huge.
Let us remark on $`g_{\pi NN}`$ which was calculated in Ref. using the three-point function method. After including dimension 5 and 7 condensates as in the previous section the sum rule becomes
$`\lambda _N^2{\displaystyle \frac{e^{M_N^2/M^2}}{M^2}}M_Ng_{\pi NN}{\displaystyle \frac{f_\pi m_\pi ^2}{\sqrt{2}m_q}}+C\left(e^{M_N^{}^2/M^2}e^{M_N^2/M^2}\right)=`$ (31)
$`\left({\displaystyle \frac{1}{\pi ^2}}E_1M^4{\displaystyle \frac{1}{5\pi ^2}}E_0M^2+{\displaystyle \frac{1}{9}}{\displaystyle \frac{\alpha _s}{\pi }}G^2\right)\overline{q}q,`$ (32)
where C is the unknown constant from $`\lambda _N^{}g_{\pi NN^{}}`$ and $`f_\pi `$ = 0.133 GeV. The contribution of the quark-gluon condensates in the OPE side is important as in the $`g_{KN\mathrm{\Lambda }}`$ and $`g_{KN\mathrm{\Sigma }}`$ sum rules. In this case we use the PCAC relation $`f_\pi ^2m_\pi ^2=4m_q\overline{q}q`$ first, then the quark condensate becomes an overall factor on both sides. However, the coupling strength $`\lambda _N`$ is still related to the quark condensate as shown in Eq. (15).
Following the same method in the previous section, and using $`\overline{q}q`$ = –(0.230 GeV)<sup>3</sup>, $`\frac{\alpha _s}{\pi }G^2`$ = 0.012 GeV<sup>4</sup>, and $`s_0`$ = 2.074 GeV<sup>2</sup> as a pure continuum threshold we get
$`|g_{\pi NN}|`$ $`=`$ $`3.65\pm 2.31,`$ (33)
$`|C|`$ $`=`$ $`0.00261\pm 0.00091\mathrm{GeV}^7,`$ (34)
the Borel interval (0.460, 1.110) GeV<sup>2</sup>, and the average deviation of the fit $`\overline{\delta }`$ 9.3 % at the central value. Here the uncertainty comes from using different input parameters, i.e. $`\overline{q}q`$ = –(0.210 GeV)<sup>3</sup> (or –(0.250 GeV)<sup>3</sup>) , $`\frac{\alpha _s}{\pi }G^2`$ = 0.015 GeV<sup>4</sup>, and $`m_0^2`$ = 0.6 (or 1.4 GeV<sup>2</sup>) as before. In this case the most error bar comes from uncertainties of the quark condensate, i.e. from the coupling strength $`\lambda _N`$.
Now, let us discuss some uncertainties in our sum rules. In Eqs. (13), (25) and (32) the contribution of the quark-gluon condensate is about 25 %, 40 %, and 20 %, respectively, of the leading term at $`M^2`$ = 1 GeV<sup>2</sup>. Thus the accurate value of this condensate is one of important factors in our sum rules, and a more precise estimate may be needed (e.g., see Ref.).
As we mentioned before, we need more higher dimensional operators to get some power correction terms in our sum rules. Their contribution will be much smaller than that of dimension 7 operators at the relevant Borel region around $`M^2`$ $``$ 1 GeV<sup>2</sup>. However, those operators may contribute because the lower limit of the Borel interval for each coupling constant is much less than 1 GeV<sup>2</sup> in our sum rules.
We find that the coupling constants become 2 or 3 times larger than the previous ones if we take the coupling strengths ($`\lambda _N`$, $`\lambda _\mathrm{\Lambda }`$ and $`\lambda _\mathrm{\Sigma }`$) from the chiral-odd baryon sum rules . For example, we get 7.04, 0.890, and 14.49 for $`|g_{KN\mathrm{\Lambda }}|`$, $`|g_{KN\mathrm{\Sigma }}|`$, and $`|g_{\pi NN}|`$, respectively, for the basic inputs. Because the coupling strengths from each baryon sum rule (the chiral-even and chiral-odd) are not the same in the whole Borel region and the discrepancy between the coupling strengths is larger in the low Borel region, we get quite different coupling constants. Of course, it should be judged by the stability of the sum rule whether one chooses the coupling strengths from the chiral-even sum rules or those from the chiral-odd sum rules.
As a final remark, in the case of $`g_{\pi NN}`$ it was shown that there is a higher pseudoscalar resonance contamination from the $`\pi `$ (1300) and $`\pi `$ (1800) in the three-point function method . Maybe there is a similar contamination from the K(1460) and K(1830) on the kaon-baryon couplings. Although the masses of the K(1460) and K(1830) are quite uncertain and these states need further experimental confirmation, we can briefly estimate the contribution of the K(1460) as done in Ref. . Using the parameters from recent works , we get
$`\left[{\displaystyle \frac{f_Mm_M^2}{Q^2+m_M^2}}\right]_{Q^2=1\mathrm{GeV}^2}=21.3\mathrm{and}2.2\mathrm{MeV}`$ (35)
for the kaon and K(1460), respectively. Here $`f_M`$ is the decay constant and $`m_M`$ is the meson mass. We take $`f_K`$ = 108 MeV, $`f_{K(1460)}`$ = 3.3 MeV and $`m_K`$ = 496 MeV, $`m_{K(1460)}`$ = 1.45 GeV in Ref. . Comparing the values in Eq. (35) to those for the pion and $`\pi `$ (1300), i.e. 1.7 and 0.4 MeV, the contamination from the excited kaon state on the kaon-baryon couplings seems smaller than that from the excited pion state on $`g_{\pi NN}`$.
In summary, including higher dimensional condensates we re-analyze our previous QCD sum rule estimate on $`g_{KN\mathrm{\Lambda }}`$ and $`g_{KN\mathrm{\Sigma }}`$ in the $`\text{/}qi\gamma _5`$ structure. The contribution of dimension 5 quark-gluon condensates is comparable to that of the leading term, and the present result is much different from the previous one.
###### Acknowledgements.
The author thanks Prof. Su H. Lee for valuable discussions and comments. This work is supported by Research Fellowship of the Japan Society for the Promotion of Science (JSPS).
|
warning/0003/physics0003095.html
|
ar5iv
|
text
|
# Variational–Wavelet Approach to RMS Envelope Equations
## 1 Introduction
In this paper we consider the applications of a new numerical-analytical technique which is based on the methods of local nonlinear Fourier analysis or Wavelet analysis to the nonlinear beam/accelerator physics problems related to root-mean-square (rms) envelope dynamics . Such approach may be useful in all models in which it is possible and reasonable to reduce all complicated problems related with statistical distributions to the problems described by systems of nonlinear ordinary/partial differential equations. In this paper we consider approach based on the second moments of the distribution functions for the calculation of evolution of rms envelope of a beam.
The rms envelope equations are the most useful for analysis of the beam self–forces (space–charge) effects and also allow to consider both transverse and longitudinal dynamics of space-charge-dominated relativistic high–bright ness axisymmetric/asymmetric beams, which under short laser pulse–driven radio-frequency photoinjectors have fast transition from nonrelativistic to relativistic regime -.
From the formal point of view we may consider rms envelope equations after straightforward transformations to standard Cauchy form as a system of nonlinear differential equations which are not more than rational (in dynamical variables). Such rational type of nonlinearities allow us to consider some extension of results from -, which are based on application of wavelet analysis technique to variational formulation of initial nonlinear problem.
Wavelet analysis is a relatively novel set of mathematical methods, which gives us a possibility to work with well-localized bases in functional spaces and give for the general type of operators (differential, integral, pseudodifferential) in such bases the maximum sparse forms.
An example of such type of basis is demonstrated on Fig. 1.
Our approach in this paper is based on the generalization of variational-wavelet approach from -, which allows us to consider not only polynomial but rational type of nonlinearities.
So, our variational-multiresolution approach gives us possibility to construct explicit numerical-analytical solution for the following systems of nonlinear differential equations
$$\dot{z}=R(z,t)\text{or}Q(z,t)\dot{z}=P(z,t),$$
(1)
where $`z(t)=(z_1(t),\mathrm{},z_n(t))`$ is the vector of dynamical variables $`z_i(t)`$,
$`R(z,t)`$ is not more than rational function of z,
$`P(z,t),Q(z,t)`$ are not more than polynomial functions of z and P,Q,R have arbitrary dependence of time.
The solution has the following form
$$z(t)=z_N^{slow}(t)+\underset{jN}{}z_j(\omega _jt),\omega _j2^j$$
(2)
which corresponds to the full multiresolution expansion in all time scales. Formula (2) gives us expansion into a slow part $`z_N^{slow}`$ and fast oscillating parts for arbitrary N. So, we may move from coarse scales of resolution to the finest one for obtaining more detailed information about our dynamical process. The first term in the RHS of equation (2) corresponds on the global level of function space decomposition to resolution space and the second one to detail space. In this way we give contribution to our full solution from each scale of resolution or each time scale (detailed description we give in part 3.2 and numerical illustration in part 7 below). The same is correct for the contribution to power spectral density (energy spectrum): we can take into account contributions from each level/scale of resolution.
In part 2 we describe the different forms of rms equations. Starting in part 3.1 from variational formulation of initial dynamical problem we construct via multiresolution analysis (3.2) explicit representation for all dynamical variables in the base of compactly supported (Daubechies) wavelets. Our solutions (3.3) are parametrized by solutions of a number of reduced algebraical problems one from which is nonlinear with the same degree of nonlinearity and the rest are the linear problems which correspond to particular method of calculation of scalar products of functions from wavelet bases and their derivatives. Then we consider further extension of our previous results. In part 4 we consider modification of our construction to the periodic case, in part 5 we consider generalization of our approach to variational formulation in the biorthogonal bases of compactly supported wavelets and in part 6 to the case of variable coefficients. In part 7 we consider results of numerical calculations.
## 2 RMS Equations
Below we consider a number of different forms of RMS envelope equations, which are from the formal point of view not more than nonlinear differential equations with rational nonlinearities and variable coefficients. Let $`f(x_1,x_2)`$ be the distribution function which gives full information about noninteracting ensemble of beam particles regarding to trace space or transverse phase coordinates $`(x_1,x_2)`$. Then (n,m) moments are:
$$x_1^nx_2^mf(x_1,x_2)dx_1dx_2$$
(3)
The (0,0) moment gives normalization condition on the distribution. The (1,0) and (0,1) moments vanish when a beam is aligned to its axis. Then we may extract the first nontrivial bit of ‘dynamical information’ from the second moments
$`\sigma _{x_1}^2`$ $`=`$ $`<x_1^2>={\displaystyle x_1^2f(x_1,x_2)dx_1dx_2}`$
$`\sigma _{x_2}^2`$ $`=`$ $`<x_2^2>={\displaystyle x_2^2f(x_1,x_2)dx_1dx_2}`$ (4)
$`\sigma _{x_1x_2}^2`$ $`=`$ $`<x_1x_2>={\displaystyle x_1x_2f(x_1,x_2)dx_1dx_2}`$
RMS emittance ellipse is given by
$$\epsilon _{x,rms}^2=<x_1^2><x_2^2><x_1x_2>^2$$
(5)
Expressions for twiss parameters are also based on the second moments.
We will consider the following particular cases of rms envelope equations, which described evolution of the moments (4) (- for full designation):
for asymmetric beams we have the system of two envelope equations of the second order for $`\sigma _{x_1}`$ and $`\sigma _{x_2}`$:
$`\sigma _{x_1}^{^{\prime \prime }}+\sigma _{x_1}^{^{}}{\displaystyle \frac{\gamma ^{}}{\gamma }}+\mathrm{\Omega }_{x_1}^2\left({\displaystyle \frac{\gamma ^{}}{\gamma }}\right)^2\sigma _{x_1}`$ $`=`$ $`{\displaystyle \frac{I}{I_0(\sigma _{x_1}+\sigma _{x_2})\gamma ^3}}+{\displaystyle \frac{\epsilon _{nx_1}^2}{\sigma _{x_1}^3\gamma ^2}},`$
$`\sigma _{x_2}^{^{\prime \prime }}+\sigma _{x_2}^{^{}}{\displaystyle \frac{\gamma ^{}}{\gamma }}+\mathrm{\Omega }_{x_2}^2\left({\displaystyle \frac{\gamma ^{}}{\gamma }}\right)^2\sigma _{x_2}`$ $`=`$ $`{\displaystyle \frac{I}{I_0(\sigma _{x_1}+\sigma _{x_2})\gamma ^3}}+{\displaystyle \frac{\epsilon _{nx_2}^2}{\sigma _{x_2}^3\gamma ^2}}`$ (6)
the envelope equation for an axisymmetric beam is
$$\sigma ^{^{\prime \prime }}+\sigma ^{^{}}\frac{\gamma ^{}}{\gamma }+\mathrm{\Omega }^2\left(\frac{\gamma ^{}}{\gamma }\right)^2\sigma =\frac{k_s}{\sigma \gamma ^3}+\frac{\epsilon _{n,th}^2}{\sigma ^3\gamma ^2}$$
(7)
Also we have related Lawson’s equation for evolution of the rms envelope in the paraxial limit, which governs evolution of cylindrical symmetric envelope under external linear focusing channel of strenghts $`K_r`$:
$$\sigma ^{^{\prime \prime }}+\sigma ^{^{}}\left(\frac{\gamma ^{}}{\beta ^2\gamma }\right)+K_r\sigma =\frac{k_s}{\sigma \beta ^3\gamma ^3}+\frac{\epsilon _n^2}{\sigma ^3\beta ^2\gamma ^2},$$
where
$$K_rF_r/r\beta ^2\gamma mc^2,\beta \nu _b/c=\sqrt{1\gamma ^2}$$
After transformations to Cauchy form we can see that all this equations from the formal point of view are not more than ordinary differential equations with rational nonlinearities and variable coefficients and correspond to the form (1) (also,we may consider regimes in which $`\gamma `$, $`\gamma ^{}`$ are not fixed functions/constants but satisfy some additional differential constraint/equation, but this case does not change our general approach).
## 3 Rational Dynamics
The first main part of our consideration is some variational approach to this problem, which reduces initial problem to the problem of solution of functional equations at the first stage and some algebraical problems at the second stage. We have the solution in a compactly supported wavelet basis. Multiresolution expansion is the second main part of our construction. The solution is parameterized by solutions of two reduced algebraical problems, one is nonlinear and the second are some linear problems, which are obtained from one of the next wavelet constructions: the method of Connection Coefficients (CC), Stationary Subdivision Schemes (SSS).
### 3.1 Variational Method
Our problems may be formulated as the systems of ordinary differential equations
$`Q_i(x){\displaystyle \frac{\mathrm{d}x_i}{\mathrm{d}t}}=P_i(x,t),x=(x_1,\mathrm{},x_n),`$ (8)
$`i=1,\mathrm{},n,\underset{i}{\mathrm{max}}degP_i=p,\underset{i}{\mathrm{max}}degQ_i=q`$
with fixed initial conditions $`x_i(0)`$, where $`P_i,Q_i`$ are not more than polynomial functions of dynamical variables $`x_j`$ and have arbitrary dependence of time. Because of time dilation we can consider only next time interval: $`0t1`$. Let us consider a set of functions
$`\mathrm{\Phi }_i(t)=x_i{\displaystyle \frac{\mathrm{d}}{\mathrm{d}t}}(Q_iy_i)+P_iy_i`$ (9)
and a set of functionals
$`F_i(x)={\displaystyle _0^1}\mathrm{\Phi }_i(t)𝑑tQ_ix_iy_i_0^1,`$ (10)
where $`y_i(t)(y_i(0)=0)`$ are dual (variational) variables. It is obvious that the initial system and the system
$$F_i(x)=0$$
(11)
are equivalent. Of course, we consider such $`Q_i(x)`$ which do not lead to the singular problem with $`Q_i(x)`$, when $`t=0`$ or $`t=1`$, i.e. $`Q_i(x(0)),Q_i(x(1))\mathrm{}`$.
In part 5 we consider more general approach, which is based on possibility taking into account underlying symplectic structure and on more useful and flexible analytical approach, related to bilinear structure of initial functional. Now we consider formal expansions for $`x_i,y_i`$:
$`x_i(t)=x_i(0)+{\displaystyle \underset{k}{}}\lambda _i^k\phi _k(t)y_j(t)={\displaystyle \underset{r}{}}\eta _j^r\phi _r(t),`$ (12)
where $`\phi _k(t)`$ are useful basis functions of some functional space ($`L^2,L^p`$, Sobolev, etc) corresponding to concrete problem and because of initial conditions we need only $`\phi _k(0)=0`$.
$$\lambda =\{\lambda _i\}=\{\lambda _i^r\}=(\lambda _i^1,\lambda _i^2,\mathrm{},\lambda _i^N),r=1,\mathrm{},N,i=1,\mathrm{},n,$$
(13)
where the lower index i corresponds to expansion of dynamical variable with index i, i.e. $`x_i`$ and the upper index $`r`$ corresponds to the numbers of terms in the expansion of dynamical variables in the formal series. Then we put (12) into the functional equations (11) and as result we have the following reduced algebraical system of equations on the set of unknown coefficients $`\lambda _i^k`$ of expansions (12):
$`L(Q_{ij},\lambda ,\alpha _I)=M(P_{ij},\lambda ,\beta _J),`$ (14)
where operators L and M are algebraization of RHS and LHS of initial problem (8), where $`\lambda `$ (13) are unknowns of reduced system of algebraical equations (RSAE)(14).
$`Q_{ij}`$ are coefficients (with possible time dependence) of LHS of initial system of differential equations (8) and as consequence are coefficients of RSAE.
$`P_{ij}`$ are coefficients (with possible time dependence) of RHS of initial system of differential equations (8) and as consequence are coefficients of RSAE.
$`I=(i_1,\mathrm{},i_{q+2}),J=(j_1,\mathrm{},j_{p+1})`$ are multiindexes, by which are labelled $`\alpha _I`$ and $`\beta _I`$ — other coefficients of RSAE (14):
$$\beta _J=\{\beta _{j_1\mathrm{}j_{p+1}}\}=\underset{1j_kp+1}{}\phi _{j_k},$$
(15)
where p is the degree of polinomial operator P (8)
$$\alpha _I=\{\alpha _{i_1}\mathrm{}\alpha _{i_{q+2}}\}=\underset{i_1,\mathrm{},i_{q+2}}{}\phi _{i_1}\mathrm{}\dot{\phi _{i_s}}\mathrm{}\phi _{i_{q+2}},$$
(16)
where q is the degree of polynomial operator Q (8), $`i_{\mathrm{}}=(1,\mathrm{},q+2)`$, $`\dot{\phi _{i_s}}=\mathrm{d}\phi _{i_s}/\mathrm{d}t`$.
Now, when we solve RSAE (14) and determine unknown coefficients from formal expansion (12) we therefore obtain the solution of our initial problem. It should be noted if we consider only truncated expansion (12) with N terms then we have from (14) the system of $`N\times n`$ algebraical equations with degree $`\mathrm{}=max\{p,q\}`$ and the degree of this algebraical system coincides with degree of initial differential system. So, we have the solution of the initial nonlinear (rational) problem in the form
$`x_i(t)=x_i(0)+{\displaystyle \underset{k=1}{\overset{N}{}}}\lambda _i^kX_k(t),`$ (17)
where coefficients $`\lambda _i^k`$ are roots of the corresponding reduced algebraical (polynomial) problem RSAE (14). Consequently, we have a parametrization of solution of initial problem by solution of reduced algebraical problem (14). The first main problem is a problem of computations of coefficients $`\alpha _I`$ (16), $`\beta _J`$ (15) of reduced algebraical system. As we will see, these problems may be explicitly solved in wavelet approach.
Next we consider the construction of explicit time solution for our problem. The obtained solutions are given in the form (17), where $`X_k(t)`$ are basis functions and $`\lambda _k^i`$ are roots of reduced system of equations. In our first wavelet case $`X_k(t)`$ are obtained via multiresolution expansions and represented by compactly supported wavelets and $`\lambda _k^i`$ are the roots of corresponding general polynomial system (14) with coefficients, which are given by CC or SSS constructions. According to the variational method to give the reduction from differential to algebraical system of equations we need compute the objects $`\alpha _I`$ and $`\beta _J`$.
### 3.2 Wavelet Framework
Our constructions are based on multiresolution approach. Because affine group of translation and dilations is inside the approach, this method resembles the action of a microscope. We have contribution to final result from each scale of resolution from the whole infinite scale of spaces. More exactly, the closed subspace $`V_j(j𝐙)`$ corresponds to level j of resolution, or to scale j. We consider a r-regular multiresolution analysis of $`L^2(𝐑^n)`$ (of course, we may consider any different functional space) which is a sequence of increasing closed subspaces $`V_j`$:
$$\mathrm{}V_2V_1V_0V_1V_2\mathrm{}$$
(18)
satisfying the following properties:
$`{\displaystyle \underset{j𝐙}{}}V_j=0,\overline{{\displaystyle \underset{j𝐙}{}}}V_j=L^2(𝐑^n),`$
$`f(x)V_j<=>f(2x)V_{j+1},`$
$`f(x)V_0<=>f(xk)V_0,,k𝐙^n.`$ (19)
There exists a function $`\phi V_0`$ such that {$`\phi _{0,k}(x)=\phi (xk),k𝐙^n`$} forms a Riesz basis for $`V_0`$.
The function $`\phi `$ is regular and localized: $`\phi `$ is $`C^{r1},\phi ^{(r1)}`$ is almost everywhere differentiable and for almost every $`x𝐑^n`$, for every integer $`\alpha r`$ and for all integer p there exists constant $`C_p`$ such that
$$^\alpha \phi (x)C_p(1+|x|)^p$$
(20)
Let $`\phi (x)`$ be a scaling function, $`\psi (x)`$ is a wavelet function and $`\phi _i(x)=\phi (xi)`$. Scaling relations that define $`\phi ,\psi `$ are
$`\phi (x)`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{N1}{}}}a_k\phi (2xk)={\displaystyle \underset{k=0}{\overset{N1}{}}}a_k\phi _k(2x),`$ (21)
$`\psi (x)`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{N2}{}}}(1)^ka_{k+1}\phi (2x+k).`$ (22)
Let indices $`\mathrm{},j`$ represent translation and scaling, respectively and
$$\phi _{jl}(x)=2^{j/2}\phi (2^jx\mathrm{})$$
(23)
then the set $`\{\phi _{j,k}\},k𝐙^n`$ forms a Riesz basis for $`V_j`$. The wavelet function $`\psi `$ is used to encode the details between two successive levels of approximation. Let $`W_j`$ be the orthonormal complement of $`V_j`$ with respect to $`V_{j+1}`$:
$$V_{j+1}=V_jW_j.$$
(24)
Then just as $`V_j`$ is spanned by dilation and translations of the scaling function, so are $`W_j`$ spanned by translations and dilation of the mother wavelet $`\psi _{jk}(x)`$, where
$$\psi _{jk}(x)=2^{j/2}\psi (2^jxk).$$
(25)
All expansions which we used are based on the following properties:
$`\{\psi _{jk}\},j,k𝐙\text{is a Hilbertian basis of }L^2(𝐑)`$
$`\{\phi _{jk}\}_{j0,k𝐙}\text{is an orthonormal basis for}L^2(𝐑),`$
$`L^2(𝐑)=\overline{V_0{\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}W_j},`$ (26)
or
$`\{\phi _{0,k},\psi _{j,k}\}_{j0,k𝐙}\text{is an orthonormal basis for}L^2(𝐑).`$
### 3.3 Wavelet Computations
Now we give construction for computations of objects (17),(18) in the wavelet case. We use compactly supported wavelet basis: orthonormal basis for functions in $`L^2(𝐑)`$.
Let be $`f:𝐑𝐂`$ and the wavelet expansion is
$`f(x)={\displaystyle \underset{\mathrm{}𝐙}{}}c_{\mathrm{}}\phi _{\mathrm{}}(x)+{\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{k𝐙}{}}c_{jk}\psi _{jk}(x)`$ (27)
If in formulae (29) $`c_{jk}=0`$ for $`jJ`$, then $`f(x)`$ has an alternative expansion in terms of dilated scaling functions only $`f(x)=\underset{\mathrm{}𝐙}{}c_J\mathrm{}\phi _J\mathrm{}(x)`$. This is a finite wavelet expansion, it can be written solely in terms of translated scaling functions. Also we have the shortest possible support: scaling function $`DN`$ (where $`N`$ is even integer) will have support $`[0,N1]`$ and $`N/2`$ vanishing moments. There exists $`\lambda >0`$ such that $`DN`$ has $`\lambda N`$ continuous derivatives; for small $`N,\lambda 0.55`$. To solve our second associated linear problem we need to evaluate derivatives of $`f(x)`$ in terms of $`\phi (x)`$. Let be $`\phi _{\mathrm{}}^n=\mathrm{d}^n\phi _{\mathrm{}}(x)/\mathrm{d}x^n`$. We consider computation of the wavelet - Galerkin integrals. Let $`f^d(x)`$ be d-derivative of function $`f(x)`$, then we have $`f^d(x)=_{\mathrm{}}c_l\phi _{\mathrm{}}^d(x)`$, and values $`\phi _{\mathrm{}}^d(x)`$ can be expanded in terms of $`\phi (x)`$
$`\varphi _{\mathrm{}}^d(x)`$ $`=`$ $`{\displaystyle \underset{m}{}}\lambda _m\phi _m(x),`$ (28)
$`\lambda _m`$ $`=`$ $`{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}\phi _{\mathrm{}}^d(x)\phi _m(x)dx,`$
where $`\lambda _m`$ are wavelet-Galerkin integrals. The coefficients $`\lambda _m`$ are 2-term connection coefficients. In general we need to find $`(d_i0)`$
$`\mathrm{\Lambda }_{\mathrm{}_1\mathrm{}_2\mathrm{}\mathrm{}_n}^{d_1d_2\mathrm{}d_n}={\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \phi _\mathrm{}_i^{d_i}(x)dx}`$ (29)
For Riccati case we need to evaluate two and three connection coefficients
$`\mathrm{\Lambda }_{\mathrm{}}^{d_1d_2}={\displaystyle _{\mathrm{}}^{\mathrm{}}}\phi ^{d_1}(x)\phi _{\mathrm{}}^{d_2}(x)𝑑x,\mathrm{\Lambda }^{d_1d_2d_3}={\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}\phi ^{d_1}(x)\phi _{\mathrm{}}^{d_2}(x)\phi _m^{d_3}(x)𝑑x`$ (30)
According to CC method we use the next construction. When $`N`$ in scaling equation is a finite even positive integer the function $`\phi (x)`$ has compact support contained in $`[0,N1]`$. For a fixed triple $`(d_1,d_2,d_3)`$ only some $`\mathrm{\Lambda }_\mathrm{}m^{d_1d_2d_3}`$ are nonzero: $`2N\mathrm{}N2,2NmN2,|\mathrm{}m|N2`$. There are $`M=3N^29N+7`$ such pairs $`(\mathrm{},m)`$. Let $`\mathrm{\Lambda }^{d_1d_2d_3}`$ be an M-vector, whose components are numbers $`\mathrm{\Lambda }_\mathrm{}m^{d_1d_2d_3}`$. Then we have the first reduced algebraical system : $`\mathrm{\Lambda }`$ satisfy the system of equations $`(d=d_1+d_2+d_3)`$
$`A\mathrm{\Lambda }^{d_1d_2d_3}=2^{1d}\mathrm{\Lambda }^{d_1d_2d_3},A_{\mathrm{},m;q,r}={\displaystyle \underset{p}{}}a_pa_{q2\mathrm{}+p}a_{r2m+p}`$ (31)
By moment equations we have created a system of $`M+d+1`$ equations in $`M`$ unknowns. It has rank $`M`$ and we can obtain unique solution by combination of LU decomposition and QR algorithm. The second reduced algebraical system gives us the 2-term connection coefficients:
$`A\mathrm{\Lambda }^{d_1d_2}=2^{1d}\mathrm{\Lambda }^{d_1d_2},d=d_1+d_2,A_{\mathrm{},q}={\displaystyle \underset{p}{}}a_pa_{q2\mathrm{}+p}`$ (32)
For nonquadratic case we have analogously additional linear problems for objects (31). Solving these linear problems we obtain the coefficients of nonlinear algebraical system (16) and after that we obtain the coefficients of wavelet expansion (19). As a result we obtained the explicit time solution of our problem in the base of compactly supported wavelets. We use for modelling D6, D8, D10 functions and programs RADAU and DOPRI for testing.
In the following we consider extension of this approach to the case of periodic boundary conditions, the case of presence of arbitrary variable coefficients and more flexible biorthogonal wavelet approach.
## 4 Variational Wavelet Approach for Periodic Trajectories
We start with extension of our approach to the case of periodic trajectories. The equations of motion corresponding to our problems may be formulated as a particular case of the general system of ordinary differential equations $`dx_i/dt=f_i(x_j,t)`$, $`(i,j=1,\mathrm{},n)`$, $`0t1`$, where $`f_i`$ are not more than rational functions of dynamical variables $`x_j`$ and have arbitrary dependence of time but with periodic boundary conditions. According to our variational approach we have the solution in the following form
$`x_i(t)=x_i(0)+{\displaystyle \underset{k}{}}\lambda _i^k\phi _k(t),x_i(0)=x_i(1),`$ (33)
where $`\lambda _i^k`$ are again the roots of reduced algebraical systems of equations with the same degree of nonlinearity and $`\phi _k(t)`$ corresponds to useful type of wavelet bases (frames). It should be noted that coefficients of reduced algebraical system are the solutions of additional linear problem and also depend on particular type of wavelet construction and type of bases.
This linear problem is our second reduced algebraical problem. We need to find in general situation objects
$`\mathrm{\Lambda }_{\mathrm{}_1\mathrm{}_2\mathrm{}\mathrm{}_n}^{d_1d_2\mathrm{}d_n}={\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \phi _\mathrm{}_i^{d_i}(x)\mathrm{d}x},`$ (34)
but now in the case of periodic boundary conditions. Now we consider the procedure of their calculations in case of periodic boundary conditions in the base of periodic wavelet functions on the interval and corresponding expansion (35) inside our variational approach. Periodization procedure gives us
$`\widehat{\phi }_{j,k}(x)`$ $``$ $`{\displaystyle \underset{\mathrm{}Z}{}}\phi _{j,k}(x\mathrm{})`$ (35)
$`\widehat{\psi }_{j,k}(x)`$ $`=`$ $`{\displaystyle \underset{\mathrm{}Z}{}}\psi _{j,k}(x\mathrm{})`$
So, $`\widehat{\phi },\widehat{\psi }`$ are periodic functions on the interval . Because $`\phi _{j,k}=\phi _{j,k^{}}`$ if $`k=k^{}\mathrm{mod}(2^j)`$, we may consider only $`0k2^j`$ and as consequence our multiresolution has the form $`{\displaystyle \underset{j0}{}}\widehat{V}_j=L^2[0,1]`$ with $`\widehat{V}_j=\mathrm{span}\{\widehat{\phi }_{j,k}\}_{k=0}^{2j1}`$ . Integration by parts and periodicity gives useful relations between objects (36) in particular quadratic case $`(d=d_1+d_2)`$:
$`\mathrm{\Lambda }_{k_1,k_2}^{d_1,d_2}=(1)^{d_1}\mathrm{\Lambda }_{k_1,k_2}^{0,d_2+d_1},\mathrm{\Lambda }_{k_1,k_2}^{0,d}=\mathrm{\Lambda }_{0,k_2k_1}^{0,d}\mathrm{\Lambda }_{k_2k_1}^d`$ (36)
So, any 2-tuple can be represent by $`\mathrm{\Lambda }_k^d`$. Then our second additional linear problem is reduced to the eigenvalue problem for $`\{\mathrm{\Lambda }_k^d\}_{0k2^j}`$ by creating a system of $`2^j`$ homogeneous relations in $`\mathrm{\Lambda }_k^d`$ and inhomogeneous equations. So, if we have dilation equation in the form $`\phi (x)=\sqrt{2}_{kZ}h_k\phi (2xk)`$, then we have the following homogeneous relations
$$\mathrm{\Lambda }_k^d=2^d\underset{m=0}{\overset{N1}{}}\underset{\mathrm{}=0}{\overset{N1}{}}h_mh_{\mathrm{}}\mathrm{\Lambda }_{\mathrm{}+2km}^d,$$
(37)
or in such form $`A\lambda ^d=2^d\lambda ^d`$, where $`\lambda ^d=\{\mathrm{\Lambda }_k^d\}_{0k2^j}`$. Inhomogeneous equations are:
$$\underset{\mathrm{}}{}M_{\mathrm{}}^d\mathrm{\Lambda }_{\mathrm{}}^d=d!2^{j/2},$$
(38)
where objects $`M_{\mathrm{}}^d(|\mathrm{}|N2)`$ can be computed by recursive procedure
$$M_{\mathrm{}}^d=2^{j(2d+1)/2}\stackrel{~}{M_{\mathrm{}}^d},\stackrel{~}{M_{\mathrm{}}^k}=<x^k,\phi _{0,\mathrm{}}>=\underset{j=0}{\overset{k}{}}\left(\genfrac{}{}{0pt}{}{k}{j}\right)n^{kj}M_0^j,\stackrel{~}{M_0^{\mathrm{}}}=1.$$
(39)
So, we reduced our last problem to standard linear algebraical problem. Then we use the same methods as in part 3.3. As a result we obtained for closed trajectories of orbital dynamics the explicit time solution (35) in the base of periodized wavelets (37).
## 5 Variational Approach in Biorthogonal <br>Wavelet Bases
Now we consider further generalization of our variational wavelet approach. Because integrand of variational functionals is represented by bilinear form (scalar product) it seems more reasonable to consider wavelet constructions which take into account all advantages of this structure.
The action functional for loops in the phase space is
$$F(\gamma )=_\gamma p𝑑q_0^1H(t,\gamma (t))𝑑t$$
(40)
The critical points of $`F`$ are those loops $`\gamma `$, which solve the Hamiltonian equations associated with the Hamiltonian $`H`$ and hence are periodic orbits.
Let us consider the loop space $`\mathrm{\Omega }=C^{\mathrm{}}(S^1,R^{2n})`$, where $`S^1=R/𝐙`$, of smooth loops in $`R^{2n}`$. Let us define a function $`\mathrm{\Phi }:\mathrm{\Omega }R`$ by setting
$$\mathrm{\Phi }(x)=_0^1\frac{1}{2}<J\dot{x},x>dt_0^1H(x(t))𝑑t,x\mathrm{\Omega }$$
(41)
The critical points of $`\mathrm{\Phi }`$ are the periodic solutions of $`\dot{x}=X_H(x)`$. Computing the derivative at $`x\mathrm{\Omega }`$ in the direction of $`y\mathrm{\Omega }`$, we find
$`\mathrm{\Phi }^{}(x)(y)={\displaystyle \frac{d}{dϵ}}\mathrm{\Phi }(x+ϵy)|_{ϵ=0}={\displaystyle _0^1}<J\dot{x}H(x),y>dt`$ (42)
Consequently, $`\mathrm{\Phi }^{}(x)(y)=0`$ for all $`y\mathrm{\Omega }`$ iff the loop $`x`$ satisfies the equation
$$J\dot{x}(t)H(x(t))=0,$$
(43)
i.e. $`x(t)`$ is a solution of the Hamiltonian equations, which also satisfies $`x(0)=x(1)`$, i.e. periodic of period 1.
But now we need to take into account underlying bilinear structure via wavelets.
We started with two hierarchical sequences of approximations spaces :
$`\mathrm{}V_2V_1V_0V_1V_2\mathrm{},\mathrm{}\stackrel{~}{V}_2\stackrel{~}{V}_1\stackrel{~}{V}_0\stackrel{~}{V}_1\stackrel{~}{V}_2\mathrm{},`$
and as usually, $`W_0`$ is complement to $`V_0`$ in $`V_1`$, but now not necessarily orthogonal complement. New orthogonality conditions have now the following form:
$$\stackrel{~}{W}_0V_0,W_0\stackrel{~}{V}_0,V_j\stackrel{~}{W}_j,\stackrel{~}{V}_jW_j$$
(44)
translates of $`\psi `$ $`\mathrm{span}`$ $`W_0`$, translates of $`\stackrel{~}{\psi }\mathrm{span}\stackrel{~}{W}_0`$. Biorthogonality conditions are
$$<\psi _{jk},\stackrel{~}{\psi }_{j^{}k^{}}>=_{\mathrm{}}^{\mathrm{}}\psi _{jk}(x)\stackrel{~}{\psi }_{j^{}k^{}}(x)dx=\delta _{kk^{}}\delta _{jj^{}},$$
(45)
where $`\psi _{jk}(x)=2^{j/2}\psi (2^jxk)`$. Functions $`\phi (x),\stackrel{~}{\phi }(xk)`$ form dual pair:
$$<\phi (xk),\stackrel{~}{\phi }(x\mathrm{})>=\delta _{kl},<\phi (xk),\stackrel{~}{\psi }(x\mathrm{})>=0\text{for}k,\mathrm{}.$$
(46)
Functions $`\phi ,\stackrel{~}{\phi }`$ generate a multiresolution analysis. $`\phi (xk)`$, $`\psi (xk)`$ are synthesis functions, $`\stackrel{~}{\phi }(x\mathrm{})`$, $`\stackrel{~}{\psi }(x\mathrm{})`$ are analysis functions. Synthesis functions are biorthogonal to analysis functions. Scaling spaces are orthogonal to dual wavelet spaces. Two multiresolutions are intertwining $`V_j+W_j=V_{j+1},\stackrel{~}{V}_j+\stackrel{~}{W}_j=\stackrel{~}{V}_{j+1}`$. These are direct sums but not orthogonal sums.
So, our representation for solution has now the form
$$f(t)=\underset{j,k}{}\stackrel{~}{b}_{jk}\psi _{jk}(t),$$
(47)
where synthesis wavelets are used to synthesize the function. But $`\stackrel{~}{b}_{jk}`$ come from inner products with analysis wavelets. Biorthogonality yields
$$\stackrel{~}{b}_\mathrm{}m=f(t)\stackrel{~}{\psi }_\mathrm{}m(t)dt.$$
(48)
So, now we can introduce this more complicated construction into our variational approach. We have modification only on the level of computing coefficients of reduced nonlinear algebraical system. This new construction is more flexible. Biorthogonal point of view is more stable under the action of large class of operators while orthogonal (one scale for multiresolution) is fragile, all computations are much more simpler and we accelerate the rate of convergence. In all types of (Hamiltonian) calculation, which are based on some bilinear structures (symplectic or Poissonian structures, bilinear form of integrand in variational integral) this framework leads to greater success.
## 6 Variable Coefficients
In the case when we have situation when our problem is described by a system of nonlinear (rational) differential equations, we need to consider extension of our previous approach which can take into account any type of variable coefficients (periodic, regular or singular). We can produce such approach if we add in our construction additional refinement equation, which encoded all information about variable coefficients . According to our variational approach we need to compute integrals of the form
$$_Db_{ij}(t)(\phi _1)^{d_1}(2^mtk_1)(\phi _2)^{d_2}(2^mtk_2)dx,$$
(49)
where now $`b_{ij}(t)`$ are arbitrary functions of time, where trial functions $`\phi _1,\phi _2`$ satisfy a refinement equations:
$$\phi _i(t)=\underset{k𝐙}{}a_{ik}\phi _i(2tk)$$
(50)
If we consider all computations in the class of compactly supported wavelets then only a finite number of coefficients do not vanish. To approximate the non-constant coefficients, we need choose a different refinable function $`\phi _3`$ along with some local approximation scheme
$$(B_{\mathrm{}}f)(x):=\underset{\alpha 𝐙}{}F_{\mathrm{},k}(f)\phi _3(2^{\mathrm{}}tk),$$
(51)
where $`F_{\mathrm{},k}`$ are suitable functionals supported in a small neighborhood of $`2^{\mathrm{}}k`$ and then replace $`b_{ij}`$ in (49) by $`B_{\mathrm{}}b_{ij}(t)`$. In particular case one can take a characteristic function and can thus approximate non-smooth coefficients locally. To guarantee sufficient accuracy of the resulting approximation to (49) it is important to have the flexibility of choosing $`\phi _3`$ different from $`\phi _1,\phi _2`$. In the case when D is some domain, we can write
$$b_{ij}(t)_D=\underset{0k2^{\mathrm{}}}{}b_{ij}(t)\chi _D(2^{\mathrm{}}tk),$$
(52)
where $`\chi _D`$ is characteristic function of D. So, if we take $`\phi _4=\chi _D`$, which is again a refinable function, then the problem of computation of (49) is reduced to the problem of calculation of integral
$`H(k_1,k_2,k_3,k_4)=H(k)=`$ (53)
$`{\displaystyle _{𝐑^s}}\phi _4(2^jtk_1)\phi _3(2^{\mathrm{}}tk_2)\phi _1^{d_1}(2^rtk_3)\phi _2^{d_2}(2^stk_4)dx`$
The key point is that these integrals also satisfy some sort of refinement equation :
$$2^{|\mu |}H(k)=\underset{\mathrm{}𝐙}{}b_{2k\mathrm{}}H(\mathrm{}),\mu =d_1+d_2.$$
(54)
This equation can be interpreted as the problem of computing an eigenvector. Thus, we reduced the problem of extension of our method to the case of variable coefficients to the same standard algebraical problem as in the preceding sections. So, the general scheme is the same one and we have only one more additional linear algebraic problem by which we in the same way can parameterize the solutions of corresponding problem.
## 7 Numerical Calculations
In this part we consider numerical illustrations of previous analytical approach. Our numerical calculations are based on compactly supported Daubechies wavelets and related wavelet families.
On Fig. 2 we present according to formulae (2) contributions to approximation of our dynamical evolution (top row on the Fig. 3) starting from the coarse approximation, corresponding to scale $`2^0`$ (bottom row) to the finest one corresponding to the scales from $`2^1`$ to $`2^5`$ or from slow to fast components (5 frequencies) as details for approximation. Then on Fig. 3, from bottom to top, we demonstrate the summation of contributions from corresponding levels of resolution given on Fig. 2 and as result we restore via 5 scales (frequencies) approximation our dynamical process(top row on Fig. 3 ). In this particular model case we considered for approximation simple two frequencies harmonic process. But the same situation we have on the Fig. 5 and Fig. 6 in case when we added to previous 2-frequencies harmonic process the noise as perturbation. Again, our dynamical process under investigation (top row of Fig. 6) is recovered via 5 scales contributions (Fig. 5) to approximations (Fig. 6). The same decomposition/approximation we produce also on the level of power spectral density in the process without noise (Fig. 4) and with noise (Fig. 7). On Fig. 8 we demonstrate the family of localized contributions to beam motion, which we also may consider for such type of approximation.
It should be noted that complexity of such algorithms are minimal regarding other possible. Of course, we may use different multiresolution analysis schemes, which are based on different families of generating wavelets and apply such schemes of numerical–analytical calculations to any dynamical process which may be described by systems of ordinary/partial differential equations with rational nonlinearities .
## Acknowledgments
We would like to thank Professor James B. Rosenzweig and Mrs. Melinda Laraneta for nice hospitality, help, support and discussions before and during Workshop and all participants for interesting discussions.
|
warning/0003/physics0003006.html
|
ar5iv
|
text
|
# On mixing times for stratified walks on the 𝑑-cube
## 1 Introduction.
The stratified random walk (SRW) on the $`d`$-cube $`Q_d`$ is the Markov chain whose state space is the set of vertices of the $`d`$-cube and whose transition probabilities are defined thus:
Given a set of non-zero probabilities $`p=(p_0,p_1,\mathrm{},p_{d1})`$, from any vertex with $`k`$ 1’s, the process moves either to any neighboring vertex with $`k+1`$ 1’s with probability $`\frac{p_k}{d}`$; or to any neighboring vertex with $`k1`$ 1’s with probability $`\frac{p_{k1}}{d}`$; or to itself with the remaining probability. The simple random walk on the $`d`$-cube corresponds to the choice $`p_k=1`$ for all $`k`$.
Vaguely speaking, the mixing time of a Markov chain is the time it takes the chain to have its distribution close to the stationary distribution under some measure of closeness. Chung and Graham studied the SRW on the $`d`$-cube in , mainly with algebraic methods, and found bounds for the mixing times under total variation and relative pointwise distances. Here we use non-algebraic methods, the electric and coupling approaches, in order to study the same SRW and get exact results for maximal commute times and bounds for cover times and mixing times under total variation distance. We take advantage of the fact that there seems to be some inequality or another linking hitting times, commute times, cover times and any definition of mixing time with any other under any measure of closeness (see Aldous and Fill and Lovász and Winkler ).
## 2 The electric approach
On a connected undirected graph $`G=(V,E)`$ such that the edge between vertices $`i`$ and $`j`$ is given a resistance $`r_{ij}`$ (or equivalently, a conductance $`C_{ij}=1/r_{ij}`$), we can define the random walk on $`G`$ as the Markov chain $`𝐗=\{𝐗(n)\}_{n0}`$ that from its current vertex $`v`$ jumps to the neighboring vertex $`w`$ with probability $`p_{vw}=C_{vw}/C(v)`$, where $`C(v)=_{w:wv}C_{vw}`$, and $`wv`$ means that $`w`$ is a neighbor of $`v`$. There may be a conductance $`C_{zz}`$ from a vertex $`z`$ to itself, giving rise to a transition probability from $`z`$ to itself. Some notation: $`𝔼_aT_b`$ and $`𝔼_aC`$ denote the expected value, starting from the vertex $`a`$, of respectively, the hitting time $`T_b`$ of the vertex $`b`$ and the cover time $`C`$, i. e., the number of jumps needed to visit all the states in $`V`$; $`R_{ab}`$ is the effective resistance, as computed by means of Ohm’s law, between vertices $`a`$ and $`b`$.
A Markov chain is reversible if $`\pi _i(i,j)=\pi _j(j,i)`$ for all $`i,j`$, where $`\{\pi _.\}`$ is the stationary distribution and $`(,)`$ are the transition probabilities. Such a reversible Markov chain can be described as a random walk on a graph if we define conductances thus:
$$C_{ij}=\pi _i(i,j).$$
(2.1)
We will be interested in finding a closed form expression for the commute time $`𝔼_0T_d+𝔼_dT_0`$ between the origin, denoted by 0, and its opposite vertex, denoted by $`d`$.
Notice first that the transition matrix for $`𝐗=\{𝐗(n),n0\}`$, the SRW on the $`d`$-cube, is doubly stochastic and therefore its stationary distribution is uniform. If we now collapse all vertices in the cube with the same number of 1’s into a single vertex, and we look at the SRW on this collapsed graph, we obtain a new reversible Markov chain $`𝐒=\{𝐒(n),n0\}`$, a birth-and-death chain in fact, on the state space $`\{0,1,\mathrm{}d\}`$, with transition probabilities
$`(k,k+1)`$ $`=`$ $`{\displaystyle \frac{dk}{d}}p_k,`$ (2.2)
$`(k,k1)`$ $`=`$ $`{\displaystyle \frac{k}{d}}p_{k1},`$ (2.3)
$`(k,k)`$ $`=`$ $`1(k,k+1)(k,k1).`$ (2.4)
It is plain to see that the stationary distribution of this new chain is the Binomial with parameters $`d`$ and $`\frac{1}{2}`$. It is also clear that the commute time between vertices 0 and $`d`$ is the same for both $`𝐗`$ and $`𝐒`$. For the latter we use the electric machinery described above, namely, we think of a linear electric circuit from $`0`$ to $`d`$ with conductances given by (2.1) for $`0id`$, $`j=i1,i,i+1`$, and where $`\pi _i=\left({\displaystyle \genfrac{}{}{0pt}{}{d}{i}}\right){\displaystyle \frac{1}{2^d}}`$.
It is well known (at least since Chandra et al. proved it in ) that
$$𝔼_aT_b+𝔼_aT_b=R_{ab}\underset{z}{}C(z),$$
(2.5)
where $`R_{ab}`$ is the effective resistance between vertices $`a`$ and $`b`$.
If this formula is applied to a reversible chain whose conductances are given as in (2.1), then it is clear that
$$C(z)=\pi _z$$
and therefore the summation in (2.5) equals 1. We get then this compact formula for the commute time:
$$𝔼_aT_b+𝔼_aT_b=R_{ab},$$
(2.6)
where the effective resistance is computed with the individual resistors having resistances
$$r_{ij}=\frac{1}{C_{ij}}=\frac{1}{\pi _i(i,j)}.$$
In our particular case of the collapsed chain, because it is a linear circuit, the effective resistance $`R_{0d}`$ equals the sum of all the individual resistances $`r_{i,i+1}`$, so that (2.6) yields
$$𝔼_0T_d+𝔼_dT_0=R_{0d}=2^d\underset{k=0}{\overset{d1}{}}\frac{1}{p_k\left(\genfrac{}{}{0pt}{}{d1}{k}\right)}.$$
(2.7)
Because of the particular nature of the chain under consideration, it is clear that $`𝔼_0T_d+𝔼_dT_0`$ equals the maximal commute time ($`\tau ^{}`$ in the terminology of Aldous ) between any two vertices.
(i) For simple random walk, formula (2.7) is simplified by taking all $`p_k=1`$. This particular formula was obtained in with a more direct argument, and it was argued there that
$$\underset{k=0}{\overset{d1}{}}\frac{1}{\left(\genfrac{}{}{0pt}{}{d1}{k}\right)}=2+o(1).$$
An application of Matthews’ result (see ), linking maximum and minimum expected hitting times with expected cover times, yields immediately that the expected cover time is $`𝔼_vC=\mathrm{\Theta }(|V|\mathrm{log}|V|),`$ which is the asymptotic value of the lower bound for cover times of walks on a graph $`G=(V,E)`$ (see ). Thus we could say this SRW is a “rapidly covered” walk.
(ii) The so-called Aldous cube (see ) corresponds to the choice $`p_k=\frac{k}{d1}`$. This walk takes place in the “punctured cube” that excludes the origin. Formula (2.7) thus, must exclude $`k=0`$ in this case, for which we still get a closed-form expression for the commute time between vertex $`d`$, all of whose coordinates are 1, and vertex $`s`$ which consists of the collapse of all vertices with a single 1:
$$𝔼_sT_d+𝔼_dT_s=2^d\underset{k=1}{\overset{d1}{}}\frac{1}{\left(\genfrac{}{}{0pt}{}{d2}{k1}\right)}.$$
(2.8)
The same argument used in (i) tells us that the summation in (2.8) equals $`2+o(1)`$ and, once again, Matthews’ result tells us that the walk on the Aldous cube has a cover time of order $`|V|\mathrm{log}|V|`$.
(iii) The choice $`p_k=\frac{1}{\left(\genfrac{}{}{0pt}{}{d1}{k}\right)}`$ would be in the terminology of a “slow walk”: the commute time is seen to be exactly equal to $`|V|\mathrm{log}_2|V|`$ and thus the expected cover time is $`O(|V|\mathrm{log}^2|V|)`$. In general, the SRW will be rapidly covered if and only if
$$\underset{k=0}{\overset{d1}{}}\frac{1}{p_k\left(\genfrac{}{}{0pt}{}{d1}{k}\right)}=c+o(1),$$
for some constant $`c`$.
Remark. A formula as compact as (2.7) could be easily obtained through the commute time formula (2.6). It does not seem that it could be obtained that easily, by just adding the individual hitting times $`𝔼_iT_{i+1}`$. (A procedure that is done, for instance, in , , , and in the next section).
## 3 The coupling approach
In order to asess the rate of convergence of the SRW on the cube $`Q_d`$ to the uniform stationary distribution $`\pi `$, we will bound the mixing time $`\tau `$ defined as
$$\tau =\mathrm{min}\{t:d(t^{})\frac{1}{2e},\text{ for all }t^{}>t\},$$
where
$$d(t)=\underset{𝐱Q_d}{\mathrm{max}}P_𝐱(𝐗(t)=)\pi (),$$
and $`\theta _1\theta _2`$ is the variation distance between probability distributions $`\theta _1`$ and $`\theta _2`$, one of whose alternative definitions is (see Aldous and Fill ), chapter 2):
$$\theta _1\theta _2=\mathrm{min}(V_1V_2),$$
where the minimum is taken over random pairs $`(V_1,V_2)`$ such that $`V_m`$ has distribution $`\theta _m,m=1,2.`$
The bound for the mixing time is achieved using a coupling argument that goes as follows: let $`\{𝐗(t),t0\}`$ and $`\{𝐘(t),t0\}`$ be two versions of the SRW on $`Q_d`$ such that $`𝐗(0)=𝐱`$ and $`𝐘(0)\pi `$. Then
$$_𝐱(𝐗(t)=)\pi ()(𝐗(t)𝐘(t)).$$
(3.1)
A coupling between the processes $`𝐗`$ and $`𝐘`$ is a bivariate process such that its marginals have the distributions of the original processes and such that once the bivariate process enters the diagonal, it stays there forever. If we denote by
$$T_𝐱=inf\{t;𝐗(t)=𝐘(t)\}$$
the coupling time, i. e., the hitting time of the diagonal, then (3.1) translates as
$$_𝐱(𝐗(t)=)\pi ()(T_𝐱>t),$$
(3.2)
and therefore,
$$d(t)\underset{𝐱Q_d}{\mathrm{max}}(T_𝐱>t).$$
(3.3)
If we can find a coupling such that $`𝔼T_𝐱=O(f(d))`$, for all $`xQ_d`$ and for a certain function $`f`$ of the dimension $`d`$, then we will also have $`\tau =O(f(d))`$. Indeed, if we take $`t=2ef(d)`$, then (3.3) and Markov’s inequality imply that $`d(t)1/2e`$ and the definition of $`\tau `$ implies $`\tau =O(f(d))`$.
We will split $`T_𝐱`$ as $`T_𝐱=T_𝐱^1+T_𝐱^2`$, where $`T_𝐱^1`$ is a coupling time for the birth-and-death process $`𝐒`$, and $`T_𝐱^2`$ is another coupling time for the whole process X, once the bivariate $`𝐒`$ process enters the diagonal, and we will bound the values of $`𝔼T_𝐱^1`$ and $`𝔼T_𝐱^2`$.
More formally, for any $`𝐱,𝐲Q_d`$ define
$`s(𝐱)={\displaystyle \underset{i=1}{\overset{d}{}}}x_i`$ (3.4)
$`d(𝐱,𝐲)={\displaystyle \underset{i=1}{\overset{d}{}}}|x_iy_i|.`$ (3.5)
Define also for the birth-and-death process $`𝐒(t)=s(𝐗(t))`$ its own mixing time:
$$\tau ^{(S)}=inf\{t;d_S(t)\frac{1}{2e}\},$$
where $`d_S(t)=\mathrm{max}_i_i(𝐒(t)=)\pi _S()`$, and $`\pi _S`$ is the stationary distribution of $`𝐒`$. Notice that $`s(𝐘(0))\pi _S`$ since $`𝐘(0)\pi `$.
Now we will prove that $`\tau ^{(S)}=O(f_S(d)),`$ for a certain function $`f_S`$ of the expected hitting times of the “central states”, and that this bound implies an analogous bound for $`𝔼T_𝐱^1`$. Indeed, as shown by Aldous , we can bound $`\tau ^{(S)}`$ by a more convenient stopping time
$$\tau ^{(S)}K_2\tau _1^{(3)}$$
(3.6)
where $`\tau _1^{(3)}=\mathrm{min}_\mu \mathrm{max}_i\mathrm{min}_{U_i}𝔼_iU_i`$ and the innermost minimum is taken over stopping times $`U_i`$ such that $`_i(S(U_i))=\mu ()`$. In particular,
$`\tau _1^{(3)}`$ $``$ $`\underset{b}{\mathrm{min}}\underset{i}{\mathrm{max}}\underset{U_i^b}{\mathrm{min}}𝔼_iU_i^b`$ (3.7)
$``$ $`\underset{b}{\mathrm{min}}\underset{i}{\mathrm{max}}𝔼_iT_b`$ (3.8)
$`=`$ $`\mathrm{max}(𝔼_0T_{d/2},𝔼_dT_{d/2})`$ (3.9)
where the innermost minimum in (3.7) is taken over stopping times $`U_i^b`$ such that $`_i(S(U_i)=b)=1`$. Expression (3.9) follows from (3.8) since we are dealing with birth and death chains. Therefore, combining (3.6) and (3.9) we have
$$\tau ^{(S)}K_2\mathrm{max}(𝔼_0T_{d/2},𝔼_dT_{d/2}):=f_S(d).$$
(3.10)
In general, a coupling implies an inequality like (3.2). However, the inequality becomes an equality for a certain maximal (non-Markovian) coupling, described by Griffeath . Let $`T_𝐱^1`$ be the coupling time for the maximal coupling between $`s(𝐗(t))`$ and $`s(𝐘(t))`$ such that
$$P_𝐱(S(𝐗(t))=)\pi _S()=(T_𝐱^1>t).$$
Then
$$d_S(t)=\underset{𝐱Q_d}{\mathrm{max}}P_𝐱(S(𝐗(t))=)\pi _S()=\underset{𝐱Q_d}{\mathrm{max}}(T_𝐱^1>t).$$
By the definition of $`\tau ^{(S)}`$ it is clear that $`P(T_𝐱^1>\tau ^{(S)})\frac{1}{2e}`$. Moreover, by the “submultiplicativity” property (see Aldous and Fill , chapter 2)
$$d(s+t)2d(s)d(t),s,t0,$$
(3.11)
we have that
$$P(T_𝐱^1>k\tau ^{(S)})\frac{1}{2e^k},k1.$$
(3.12)
Thus
$`𝔼T_𝐱^1`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}𝔼\left(T_𝐱^1\mathrm{𝟏}((k1)\tau ^{(S)}<T_𝐱^1k\tau ^{(S)})\right)`$
$``$ $`\tau ^{(S)}{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}kP((k1)\tau ^{(S)}<T_𝐱^1k\tau ^{(S)})`$
$``$ $`\tau ^{(S)}{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}kP((k1)\tau ^{(S)}<T_𝐱^1)`$
$``$ $`\tau ^{(S)}\left(1+{\displaystyle \underset{k=2}{\overset{\mathrm{}}{}}}k{\displaystyle \frac{1}{2e^{k1}}}\right).`$
Since the series in the right hand side converges we have
$$𝔼T_𝐱^1=O(f_S(d)).$$
Once the bivariate $`𝐒`$ process hits the diagonal
$$𝐃=\{(𝐱,𝐲)Q_d\times Q_d;\underset{i=1}{\overset{d}{}}x_i=\underset{i=1}{\overset{d}{}}y_i\},$$
(3.13)
we devise one obvious coupling that forces the bivariate $`𝐗`$ process to stay in $`𝐃`$ and such that the distance defined in (3.5) between the marginal processes does not decrease. In words: we select one coordinate at random; if the marginal processes coincide in that coordinate, we allow them to evolve together; otherwise we select another coordinate in order to force two new coincidences. Formally, for each ($`𝐗(t),𝐘(t))𝐃`$, let $`I_1,I_2`$ and $`I_3`$ be the partition of $`\{0,1,\mathrm{},d\}`$ such that
$`I_1`$ $`=`$ $`\{i;X_i(t)=Y_i(t)\}`$
$`I_2`$ $`=`$ $`\{i;X_i(t)=0,Y_i(t)=1\}`$
$`I_3`$ $`=`$ $`\{i;X_i(t)=1,Y_i(t)=0\}`$
Given $`(𝐗(t),𝐘(t))𝐃`$, choose $`i`$ u.a.r. from $`\{0,1,\mathrm{},d\}`$.
* If $`iI_1`$;
1. If $`X_i(t)=1`$ then make $`X_i(t+1)=Y_i(t+1)=0`$ with probability $`p_{s(𝐗(t))1}`$; otherwise $`X_i(t+1)=Y_i(t+1)=1`$.
2. If $`X_i(t)=0`$ then make $`X_i(t+1)=Y_i(t+1)=1`$ with probability $`p_{s(𝐗(t))}`$; otherwise $`X_i(t+1)=Y_i(t+1)=0`$.
* If $`iI_2`$;
1. Select $`jI_3`$ u.a.r.;
2. Make $`X_i(t+1)=Y_j(t+1)=1`$ with probability $`p_{s(𝐗(t))}`$; otherwise $`X_i(t+1)=Y_j(t+1)=0`$.
* If $`iI_3`$;
1. Select $`jI_2`$ u.a.r.;
2. Make $`X_i(t+1)=Y_j(t+1)=0`$ with probability $`p_{s(𝐗(t))1}`$; otherwise $`X_i(t+1)=Y_j(t+1)=1`$.
Then, it is easy to check that $`(𝐗(t+1),𝐘(t+1))𝐃`$ and $`d(𝐗(t+1),𝐘(t+1))d(𝐗(t),𝐘(t))`$. Moreover, noticing that $`|I_2|=|I_3|=d(𝐗(t),𝐘(t))/2`$, we have
$$\left(d(𝐗(t+1),𝐘(t+1))=d(𝐗(t),𝐘(t))2|𝐗(t),𝐘(t)\right)=\frac{d(𝐗(t),𝐘(t))}{2d}(p_{s(𝐗(t))}+p_{s(𝐗(t))1})$$
(3.14)
$$\left(d(𝐗(t+1),𝐘(t+1))=d(𝐗(t),𝐘(t))|𝐗(t),𝐘(t)\right)=1\frac{d(𝐗(t),𝐘(t))}{2d}(p_{s(𝐗(t))}+p_{s(𝐗(t))1}).$$
(3.15)
In this case, it is straightforward to compute
$`m(i,s(𝐗(t)))`$ $`=`$ $`i𝔼\left[d(𝐗(t+1),𝐘(t+1))|d(𝐗(t),𝐘(t))=i,𝐗(t),𝐘(t)\right]`$ (3.16)
$`=`$ $`{\displaystyle \frac{i}{d}}(p_{s(𝐗(t))}+p_{s(𝐗(t))1}).`$
Let $`T_𝐱^2`$ be the coupling time for the second coupling just described. That is, let $`T_𝐱^2=inf\{t>T_𝐱^1:d(𝐗(t),𝐘(t))=0\}`$. Then, as a consequence of the optional sampling theorem for martingales we have the following comparison lemma (cf. Aldous and Fill , Chapter 2).
###### Lemma 3.17
$$𝔼\left[T_𝐱^2|d(𝐗(T_𝐱^1),𝐘(T_𝐱^1))=L,(𝐗(T_𝐱^1),𝐘(T_𝐱^1))𝐃,s(𝐗(T_𝐱^1))=s\right]\underset{i=1}{\overset{L}{}}\frac{d}{i(p_s+p_{s1})}$$
(3.18)
for all $`s=0,1,\mathrm{},d`$.
Proof. Define $`(𝐗^{}(t),𝐘^{}(t))=(𝐗(t+T_𝐱^1),𝐘(t+T_𝐱^1)`$ for all $`t0`$. Define $`Z(t)=d(𝐗^{}(t),𝐘^{}(t))`$ and $`_t=\sigma (𝐗^{}(t),𝐘^{}(t))`$. Then, it follows from (3.16) that
$$m(i,s)=iE\left[Z(1)|Z(0)=i,s(𝐗^{}(0))=s\right].$$
(3.19)
Also, for all $`s\{1,\mathrm{},d1\}`$, $`0<m(1,s)m(2,s)\mathrm{}m(d,s)`$. Fix $`s\{1,\mathrm{},d1\}`$ and write
$$h(i)=\underset{j=1}{\overset{i}{}}\frac{1}{m(i,s)}$$
(3.20)
and extend $`h`$ by linear interpolation for all real $`0xd`$. Then $`h`$ is concave and for all $`i1`$
$`𝔼\left[h(Z(1))|Z(0)=i,s(𝐗^{}(0))=s\right]`$ $``$ $`h(im(i,s))`$
$``$ $`h(i)m(i,s)h^{}(i)`$
$`=`$ $`h(i)1,`$
where the first inequality follows from the concavity of $`h`$ and $`h^{}`$ is the first derivative of $`h`$. Now, defining $`\stackrel{~}{h}`$ such that
$$h(i)=1+\underset{j}{}\left[h(Z(1))|Z(0)=i,s(𝐗^{}(0))=s\right]h(j)+\stackrel{~}{h}(i)$$
(3.21)
and
$$M(t)=t+h(Z(t))+\underset{u=0}{\overset{t1}{}}\stackrel{~}{h}(Z(u))$$
(3.22)
we have that $`M`$ is an $`_t`$-martingale and applying the optional sampling theorem to the stopping time $`T_0=inf\{t;Z(t)=0\}`$ we have
$$𝔼\left[M(T_0)|Z(0)=i,s(𝐗^{}(0))=s\right]=𝔼\left[M(0)|Z(0)=i,s(𝐗^{}(0))=s\right]=h(i).$$
(3.23)
Noticing that $`M(T_0)T_0`$ and $`T_0=T_𝐱^2`$, we obtain the desired result $``$
Since $`s(𝐗(t))`$ is distributed as $`\pi _S`$, we can write:
$$𝔼\left[T_𝐱^2|d(𝐗(T_𝐱^1),𝐘(T_𝐱^1))=L,(𝐗(T_𝐱^1),𝐘(T_𝐱^1))𝐃\right]\underset{s=0}{\overset{d}{}}\pi _S(s)\underset{i=1}{\overset{L}{}}\frac{d}{i(p_s+p_{s1})}:=g(d).$$
(3.24)
Putting the pieces together, we have found a coupling time $`T_𝐱`$ for the whole process such that
$$𝔼T_𝐱f_S(d)+g(d).$$
The task now is to find explicit bounds for $`f_S(d)`$ and $`g(d)`$ for particular workable cases.
To avoid unnecessary complications, we will assume $`d=2m`$, and compute only the hitting times for the $`𝐒`$ process of the type $`E_0T_m`$. Hitting times in birth-and-death processes assume the following closed-form (see for an electrical derivation):
$$𝔼_kT_{k+1}=\frac{1}{\pi _kP(k,k+1)}\underset{i=0}{\overset{k}{}}\pi _i,0kd1,$$
and in our case this expression turns into
$$𝔼_kT_{k+1}=\frac{\underset{i=0}{\overset{k}{}}\left(\genfrac{}{}{0pt}{}{d}{i}\right)}{\left(\genfrac{}{}{0pt}{}{d1}{k}\right)p_k}.$$
Therefore
$$𝔼_0T_m=\underset{k=0}{\overset{m1}{}}\frac{\underset{i=0}{\overset{k}{}}\left(\genfrac{}{}{0pt}{}{2m}{i}\right)}{\left(\genfrac{}{}{0pt}{}{2m1}{k}\right)p_k}.$$
(3.25)
(i) In case all $`p_k=1`$, we have the simple random walk on the cube, and it turns out there is an even more compact expression of (3.25), namely:
$$𝔼_0T_m=\underset{k=0}{\overset{m1}{}}\frac{\underset{i=0}{\overset{k}{}}\left(\genfrac{}{}{0pt}{}{2m}{i}\right)}{\left(\genfrac{}{}{0pt}{}{2m1}{k}\right)}=m\underset{k=0}{\overset{m1}{}}\frac{1}{2k+1},$$
(3.26)
as was proved in , and the right hand side of (3.26) equals $`m[H(2m){\displaystyle \frac{1}{2}}H(m)]`$, where $`H(n)=1+\frac{1}{2}+\mathrm{}+\frac{1}{n}`$, allowing us to conclude immediately that in this case $`𝔼_0T_m=𝔼_0T_{d/2}\frac{d}{4}\mathrm{log}d+\frac{d}{4}\mathrm{log}2`$. Also, we have
$$𝔼[T_𝐱^2|d(𝐗(T_𝐱^1),𝐘(T_𝐱^1))=L,(𝐗(T_𝐱^1),𝐘(T_𝐱^1))𝐃]\frac{d}{2p}\underset{i=1}{\overset{L}{}}\frac{1}{i}\frac{d}{2p}O(\mathrm{log}L).$$
(3.27)
Thus in this case both $`f_S(d)`$ and $`g(d)`$, and a fortiori $`𝔼T_𝐱`$ and $`\tau `$, are $`O(d\mathrm{log}d)`$.
(ii) For the Aldous cube, $`p_k=\frac{k}{d1}`$, and (3.25) becomes (recall this cube excludes the origin):
$$𝔼_1T_m=\underset{k=1}{\overset{m1}{}}\frac{\underset{i=0}{\overset{k}{}}\left(\genfrac{}{}{0pt}{}{2m}{i}\right)}{\left(\genfrac{}{}{0pt}{}{2m2}{k1}\right)}=\underset{k=0}{\overset{m2}{}}\frac{\underset{i=0}{\overset{k}{}}\left(\genfrac{}{}{0pt}{}{2m}{i}\right)}{\left(\genfrac{}{}{0pt}{}{2m2}{k}\right)}+\underset{k=0}{\overset{m2}{}}\frac{\left(\genfrac{}{}{0pt}{}{2m}{k+1}\right)}{\left(\genfrac{}{}{0pt}{}{2m2}{k}\right)}.$$
(3.28)
After some algebra, it can be shown that the second summand in (3.28) equals $`(2m1)[H(2m1){\displaystyle \frac{1}{m}}]`$, and the first summand can be bounded by twice the expression in (3.26), on account of the fact that $`\left({\displaystyle \genfrac{}{}{0pt}{}{2m1}{k}}\right)2\left({\displaystyle \genfrac{}{}{0pt}{}{2m2}{k}}\right),`$ for $`0km1`$. Therefore, we can write
$$𝔼_1T_{d/2}\frac{3}{2}d\mathrm{log}d+\text{ smaller terms},$$
thus improving by a factor of $`\frac{1}{2}`$ the computation of the same hitting time in .
Also, we have
$$𝔼[T_𝐱^2|d(𝐗(T_𝐱^1),𝐘(T_𝐱^1))=L,(𝐗(T_𝐱^1),𝐘(T_𝐱^1))𝐃,s(𝐗(T_𝐱^1))=s]\underset{i=1}{\overset{L}{}}\frac{d(d1)}{i(2s1)}$$
(3.29)
Thus, in this case
$`𝔼[T_𝐱^2|d(𝐗(T_𝐱^1),𝐘(T_𝐱^1))=d,(𝐗(T_𝐱^1),𝐘(T_𝐱^1))𝐃]`$ (3.30)
$``$ $`{\displaystyle \underset{s=1}{\overset{d}{}}}\pi _S(s){\displaystyle \underset{i=1}{\overset{d}{}}}{\displaystyle \frac{d(d1)}{i(2s1)}}`$
$``$ $`\mathrm{\Phi }(\sqrt{d}/3){\displaystyle \underset{i=1}{\overset{d}{}}}{\displaystyle \frac{d(d1)}{i}}+(1\mathrm{\Phi }(\sqrt{d}/3)){\displaystyle \underset{i=1}{\overset{d}{}}}{\displaystyle \frac{d(d1)}{i(2d/31)}}`$
$``$ $`e^{d/9}d(d1)\mathrm{log}d+(1e^{d/9})d\mathrm{log}d.`$
And so $`\tau =O(d\mathrm{log}d)`$ also in this case.
(iii) Slower walks. Consider the case when the probability $`p_k`$ grows exponentially in $`k`$, more specifically
$$p_k=\left(\frac{k+1}{n+1}\right)^\alpha $$
(3.31)
with $`\alpha >1`$. In this case, it seems that (3.25) is useless to get a closed expression for $`𝔼_0T_{d/2}`$. However, Graham and Chung provide the following bound
$$𝔼_iT_{d/2}c_0(\alpha )d^\alpha ,\text{ for all }dd_0(\alpha ),0id$$
(3.32)
where $`c_0(\alpha )`$ and $`d_0(\alpha )`$ are constants depending only on $`\alpha `$. Moreover, (3.24) becomes
$`g(d)`$ $`=`$ $`{\displaystyle \underset{s=0}{\overset{d}{}}}\pi _S(s){\displaystyle \underset{i=1}{\overset{d}{}}}{\displaystyle \frac{d(d+1)^\alpha }{i((s+1)^\alpha +s^\alpha )}}`$ (3.33)
$`=`$ $`d(d+1)^\alpha {\displaystyle \underset{s=0}{\overset{d}{}}}{\displaystyle \frac{\pi _S(s)}{(s+1)^\alpha +s^\alpha }}{\displaystyle \underset{i=1}{\overset{d}{}}}{\displaystyle \frac{1}{i}}`$
$``$ $`d(d+1)^\alpha \mathrm{log}d{\displaystyle \underset{s=0}{\overset{d}{}}}{\displaystyle \frac{\pi _S(s)}{(s+1)^\alpha +s^\alpha }}`$
$``$ $`d(d+1)^\alpha \mathrm{log}d{\displaystyle \underset{s=0}{\overset{d}{}}}{\displaystyle \frac{\pi _S(s)}{(s+1)^\alpha }}`$
$`=`$ $`d(d+1)^\alpha \mathrm{log}d𝔼\left[{\displaystyle \frac{1}{(1+X)^\alpha }}\right],`$
where $`X`$ is a Binomial ($`d,\frac{1}{2}`$) random variable. Jensen’s inequality and the same argument that lead to (3.30) show that $`𝔼[(1+X)^\alpha ]O(d^\alpha )`$ and (3.33) can be bound by $`O(d\mathrm{log}d)`$. This fact together with (3.32) allows us to conclude that $`\tau =O(d^\alpha )`$ in this case, thus improving on the rate of the mixing time provided in by a factor of $`\mathrm{log}d`$.
Acknowledgments. This paper was initiated when both authors were visiting scholars at the Statistics Department of the University of California at Berkeley. The authors wish to thank their hosts, especially Dr. David Aldous with whom they shared many fruitful discussions. This work was partially supported by FAPESP 99/00260-3.
|
warning/0003/hep-ph0003010.html
|
ar5iv
|
text
|
# 1 Introduction
## 1 Introduction
Various observational facts, notably the dynamics of spiral galaxies, clusters and superclusters, and their theoretical explanations in modern cosmology and astrophysics can be reconciled only by postulating the existence of the non-luminous matter at amounts exceeding the total matter density provided by baryons. Thus, if gravity does not undergo some drastic changes at distances larger than few kpc (a highly improbable option from various points of view), the explanation of the rotational curves of spiral galaxies requires the existence of halos, composed of dark matter objects.
One of the most intriguing options from the point of view of particle physics is the possibility that dark matter is composed of massive particles of yet unknown identity. On the one hand these (presumably) weakly interacting massive particles (WIMPs) are the subject of intense experimental searches at numerous underground detectors . On the other hand, some of the extensions of the Standard Model naturally incorporate such stable particles. For example, in supersymmetric extensions of SM the lightest supersymmetric particle can be stable and neutral, thus providing a viable dark matter candidate. Extensive theoretical studies of stable SUSY particles, in particular neutralinos, have been performed over the last two decades for wide ranging values of the supersymmetric masses and couplings. We believe, however, that unless supersymmetry is experimentally verified and the stability of the lightest superpartner is established in collider experiments, other possibilities for WIMPs should not be discarded.
In this letter we adopt a completely different, yet fully justified, phenomenological approach to the physics of WIMPs. We do not make any particular assumptions about the origin and possible identity of these particles, apart from the assumption that their masses are larger than a few GeV. We then assess what can be inferred about the properties of WIMPs directly from experimental limits implied by various searches. An additional incentive to address WIMP physics in a model independent way stems from recently discovered problems with subgalactic structure formation within the non-interacting cold dark matter scenario . A “generalized” WIMP with sufficient self-interaction may avoid these problems .
Among important low-energy properties of WIMPs are their electromagnetic form factors. It has been known for a long time that the possibility of charged WIMPs is strongly disfavored and stringent limits exist in the case of a fractional charge . Here we assume that WIMPs are neutral and discuss other electro-magnetic form factors such as magnetic and electric dipole moments, anapole moment and polarizability. We put strong constraints on the possible electromagnetic form factors of WIMPs using the result of direct experimental searches aimed at the detection of the WIMP-nucleus recoil signal. These constraints depend on the number density and average velocities of WIMPs in the halo. The limits on fluxes of energetic neutrinos, originating from the annihilation of WIMPs in the center of the Sun and the Earth, can also be used to obtain indirect constraints on the form factors. These constraints are less certain, as they require additional assumptions about annihilation rates and neutrino branching ratios.
This paper is organized as follows. The calculations of the WIMP-nuclei cross sections are performed in the next section. In section 3 we present direct limits on the form factors, using the results of various underground experiments. In section 4 we obtain indirect limits on the form factors, using limits on neutrino fluxes. Our conclusions are reserved for section 5.
## 2 Electromagnetic form factors of WIMPs
The interaction of any compact object and a slowly varying electro-magnetic field can be parametrized via a set of electro-magnetic form factors. In our case, the object is a WIMP with spin $`S`$, and its interactions, linear in external fields $`𝐁`$ and $`𝐄`$, can be written as
$`H=\mu 𝐁{\displaystyle \frac{𝐒}{S}}d𝐄{\displaystyle \frac{𝐒}{S}}a𝐣{\displaystyle \frac{𝐒}{S}}`$ (1)
$`{\displaystyle \frac{1}{4S(2S1)}}[S_iS_j+S_jS_i{\displaystyle \frac{2}{3}}\delta _{ij}S(S+1)]\left(Q{\displaystyle \frac{}{x_i}}E_j+M{\displaystyle \frac{}{x_i}}B_j\right)+\mathrm{}.`$
Here $`\mu `$ and $`d`$ are magnetic and electric dipole moments, and $`M`$ and $`Q`$ are magnetic and electric quadrupole moments. The anapole moment $`a`$ is the form factor which describes the contact interaction with the external current density $`𝐣`$ .
If the spin of the WIMP is zero, these moments do not exist and the interaction with electromagnetic field is given by the charge radius $`r_D`$ of the WIMP and the polarizabilities
$$H=\frac{1}{6}er_D^2\frac{}{x_i}E_i\frac{1}{2}\chi _EE^2\frac{1}{2}\chi _BB^2\chi _{EB}𝐄𝐁+\mathrm{}$$
(2)
If time reversal is a good symmetry, then $`d`$, $`M`$ and $`\chi _{EB}`$ necessarily vanish. Conserved parity also forbids the existence of $`a`$. Here we do not assume that these symmetries are preserved a priori, and we allow all of the form factors to exist.
In what follows we calculate the cross section of the elastic scattering of a WIMP off a nucleus due to the existence of the electromagnetic form factors. We start with the scattering due to the magnetic moment which is arguably the most interesting case.
The calculation can be done in the non-relativistic limit, as $`v/c10^3`$ on average for WIMPs in the halo. We assume that the possible size of the form factors is of the order of, or smaller than, the respective characteristic values for nuclei, $`\mu <\mu _n`$, $`Q<e\mathrm{fm}^3`$, and so on. It can then easily be checked that this range of velocities also ensures the applicability of the Born approximation. In our units, $`\mathrm{}=1`$, $`c=1`$, $`e^2=\alpha =1/137`$ and $`\mu _n=e/2m_p`$.
The cross-section for the scattering of a Wimp through the interaction of its magnetic moment with the magnetic moment of the nucleus is
$$\sigma =\frac{32\pi }{3}M_R^2(N,D)\mu _D^2\mu _N^2\frac{\left(S+1\right)\left(I+1\right)}{3SI},$$
(3)
where $`M_R(N,D)`$ is the reduced mass of the WIMP-nucleus system, $`\mu _D`$ and $`\mu _N`$ are the magnetic moments of the WIMP and the nucleus, respectively, and $`I`$ is the spin of the nucleus. The cross-section for the interaction between the electric dipole moment of the WIMP and the charge of the nucleus is
$$\sigma =8\pi Z^2\left(\frac{d}{e}\right)^2\left(\frac{\alpha }{v}\right)^2\frac{S+1}{3S}\mathrm{ln}\frac{q_{max}}{q_{min}}.$$
(4)
It is logarithmically divergent at large distances. In practice, $`q_{max}`$ can be chosen to be of the order of $`M_R(N,D)v`$ and $`q_{min}`$ is determined either by the inverse of the impact parameters at which the nuclear charge is screened by the electrons or by the lowest momentum transfer that can be detected in the experiment, whichever is larger. The cross-section for the interaction between the nuclear electric current and the anapole of the WIMP is given by
$$\sigma =\frac{1}{2\pi }Z^2\alpha ^2\frac{M_R^2(N,D)}{4m_p^2}\left(\frac{a}{\mu _n}\right)^2v^2\frac{S+1}{3S},$$
(5)
where $`\mu _n`$ is the nuclear magneton, and $`m_p`$ is the proton mass. The contribution of the nuclear spin current is negligibly small. The cross-section for scattering through the interaction between the electric quadrupole moment of the WIMP and the electric charge of the nucleus takes the form
$$\sigma =\frac{8\pi }{9}M_R^2(N,D)Z^2\alpha ^2\left(\frac{Q}{e}\right)^2\frac{\left(S+1\right)\left(2S+3\right)}{10S\left(2S1\right)}.$$
(6)
In Eqs. (4)-(6) the normalization of the WIMP spin dependence is such that it is equal to 1 for the lowest spin for which these moments may exist, $`(S+1)/3S=1`$ for $`S=1/2`$ and $`(S+1)(2S+3)/10S(2S1)=1`$ for $`S=1`$. The cross-section for scattering of a WIMP because of its charge radius in the electric field of the nucleus is given by
$$\sigma =\frac{4\pi }{9}M_R^2(N,D)Z^2\alpha ^2r_D^4.$$
(7)
Finally, the cross-section for the scattering of a WIMP through its polarizability in the electric field of the nucleus, is
$$\sigma \frac{144\pi }{25}M_R^2(N,D)Z^4\alpha ^2\frac{\chi _E^2}{r_0^2}.$$
(8)
This cross-section diverges for a point-like nucleus. We therefore took the charge distribution of the nucleus to be homogeneous in a sphere with radius $`r_01.2\mathrm{fm}\sqrt[3]{A}`$.
Scattering due to the magnetic quadrupole moments and magnetic and mixed polarizabilities, $`\chi _{EB}`$, was also considered. These effects turned out to be additionally suppressed and do not lead to interesting limits. We therefore do not quote them here.
## 3 Direct constraints
We illustrate the bounds on the electro-magnetic form factors of WIMPs which are implied by the limit on the recoil signal with results of the NaI detector of the DAMA experiment , and Ge detector of the CDMS collaboration . Other experiments aimed at the direct detection of dark matter can be used as well to get limits which would typically be a few times weaker (see also Ref. for more complete listing of experiments). Traditionally, the results of the direct searches are presented in the form of a normalized bound on the “WIMP–proton”cross-section in the spin dependent case, and a normalized “WIMP–nucleon”cross-section in the spin independent case. Such normalized bounds are useful to compare results for different types of nuclei, to compare results from different experiments, and to compare theoretical predictions with experimental data, but for our purpose of determining bounds on the WIMP electro-magnetic form factors, it is necessary to reverse this normalization. The bounds, quoted in are, in some sense, specifically formulated for neutralino-nucleon scattering, which necessarily implies a nuclear theory input when switching from nucleus-neutralino to nucleon-neutralino description. It is clear that in the case of the scattering due to the WIMP electro-magnetic form factors nuclear theory is not involved, and much of the nuclear physics uncertainty is not there. Thus, in order to obtain bounds on the form factors, we recover the nucleus-WIMP cross section from the nucleon-WIMP by the use of the equation
$$\sigma =\frac{M_R^2(N,D)}{M_R^2(p,D)}\frac{I(N)}{I(p,n)}\sigma _{p,n}.$$
(9)
Here $`p`$ and $`n`$ stand for proton and neutron, respectively, and the ratio of the reduced masses takes into account a kinematic factor. In the spin dependent case, $`I(N)/I(p)`$ is a spin factor, which takes the values $`I(N)/I(p)=0.009`$ for Iodine, and $`I(N)/I(p)=0.055`$ for Sodium . In the spin independent case, $`I(N)/I(n)=(AZ)^2`$ for Ref. and $`I(N)/I(n)=A^2`$ for Ref. . For simplicity, we set all nuclear form factors equal to one, which is a better approximation for Sodium than for Iodine.<sup>1</sup><sup>1</sup>1For a more rigorous treatment electric and magnetic form factors of nuclei should be included, for many of which the experimental results are available. A more accurate analysis could be done by the experimental collaborations using their original data. The ensuing bounds on the WIMP electro-magnetic form factors are listed in Table(1).
The bounds presented in this table are obtained with the standard assumptions that the mass density of WIMPs in our halo is 0.3 Gev/cm<sup>3</sup> and that their velocities are of the order of $`200`$ km/sec. The data obtained with germanium detectors are sensitive to all electro-magnetic form factors discussed here except the magnetic dipole moment, as the most abundant germanium isotope has no spin. Let us assume that the recently reported annual modulation of the DAMA signal, consistent with $`m_D50`$ GeV , and the absence of the signal, resulting in strong limits on spin-independent cross section, quoted by CDMS , are both correct. Then the presence of a WIMP magnetic moment at the level of few$`\times \mathrm{\hspace{0.17em}10}^5`$$`1\times 10^4`$ compared to the size of the nuclear magneton can reconcile the two experimental results.
It is instructive to consider how these bounds are relaxed in the scaling regime, when the WIMP mass is larger than nuclear masses and their abundance is kept as a free parameter. To this end we introduce the parameter $`\delta _D`$, characterizing the fraction of WIMPs of species $`D`$ in the halo dark matter density ($`\delta _D=1`$ in Table 1). The counting rate in a detector is proportional to the number density of the WIMP particles and therefore should scale with $`M_D`$ and $`\delta _D`$ according to the following formula:
$$R_{N,D}(\delta _D,M_D)=R_N^0\delta _D\left(\frac{M_D^0}{M_D}\right),$$
(10)
where $`M_D^0`$ is the lowest mass of WIMPs at which the scaling sets in and $`R_N^0`$ is the counting rate, corresponding to $`M_D=M_D^0`$ and $`\delta _D=1`$. It is easy to see from the figures in Refs. that $`M_D=100`$ GeV is close to the mass where the scaling behaviour begins , and therefore we can simply rescale the bounds given in Table 1 for an arbitrary $`\delta _D`$ and arbitrary and heavy $`M_D`$, by multiplying them by $`\sqrt{\delta _DM_D/100\mathrm{G}\mathrm{e}\mathrm{V}}`$.
## 4 Indirect constraints
The elastic scattering of WIMPs in the solar (and earth) media leads to gravitational trapping. Due to the scattering, WIMPs dissipate energy and eventually they accumulate in the center of the sun (and earth). Their density grows until the accumulation rate is counter-balanced by the pair annihilation rate, and an equilibrium is reached. In the process of annihilation, among other SM products, energetic neutrinos can be produced. The interaction of these neutrinos in the earth in turn leads to the production of energetic muons, which can be detected with neutrino telescopes.
Thus limits on the flux of energetic neutrinos can be used to set bounds on the electro-magnetic form factors of WIMPs. For WIMPS interacting through their electro-magnetic form factors, the trapping rate in the center of the sun and earth is determined from the elastic WIMP-nucleus cross sections that we calculated in Eqs.(3)-(8).
We use the results of Ref. , which follows the original treatment in Refs. , and estimate the neutrino flux at the surface of the earth due to WIMP annihilation in the sun as
$$\varphi _\nu (560cm^2s^1)N_{eff}\delta _D\sigma _{p,36}\frac{\text{GeV}^2}{M_D^2}.$$
(11)
In this formula $`N_{eff}`$ is the average number of neutrinos per annihilation event and $`\sigma _{p,36}`$ is the WIMP-proton elastic cross section in units of $`10^{36}`$ cm<sup>2</sup>. It is assumed that the equilibrium number of WIMPs in the sun has been reached, so that the annihilation and capture rates are equal.
The experimental limit on the flux of energetic upward muons (for example obtained by the Kamioka, MACRO and Baksan experiments, ), is $`\varphi _\mu <1.4\times 10^{14}`$ cm<sup>-2</sup>s<sup>-1</sup>, and the probability that a neutrino directed towards the detector produces a muon at the detector is $`P(100\mathrm{G}\mathrm{e}\mathrm{V})10^7`$ . Using Eq.(11), we deduce the following indirect limits on the electro-magnetic form factors of WIMPs with $`m_D=100`$ GeV,
$`\left|\mu _D/\mu _n\right|\sqrt{\left(S+1\right)/3S}`$ $`<`$ $`510^6`$
$`\left|d/e\right|\sqrt{\left(S+1\right)/3S}`$ $`<`$ $`210^{22}\mathrm{cm}`$
$`\left|Q/e\right|\sqrt{\left(S+1\right)\left(2S+3\right)/10S\left(2S1\right)}`$ $`<`$ $`10^8\mathrm{fm}^2`$ (12)
$`\left|a/\mu _n\right|\sqrt{\left(S+1\right)/3S}`$ $`<`$ $`10^2\mathrm{fm}`$
$`r_D^2`$ $`<`$ $`1.410^8\mathrm{fm}^2`$
$`\left|\chi \right|`$ $`<`$ $`5.610^7\mathrm{fm}^3,`$
where $`N_{eff}`$ is taken to be 0.1 and $`\delta _D=1`$. Note that for higher WIMP masses the absorption of neutrinos in the core of the sun becomes important .
If the WIMP-nucleus elastic scattering cross-sections in the earth are comparable to the WIMP-proton cross-section, than the neutrino signal from the center of the earth is approximately four orders of magnitude weaker than the signal from the sun . However, the limit on $`\chi _E`$ from the neutrino flux from the center of the earth is almost one order of magnitude better than the limit derived from the neutrino flux from the sun, as the corresponding WIMP-nucleus elastic scattering cross-sections are enhanced by a large factor, $`Z^4510^5`$, for heavy (Fe) nuclei (see Eq.(8)) comprising the core of the earth.
Although indirect bounds are definitely stronger, they do require a number of additional assumptions. In particular, these bounds require that annihilation actually occurs, which in practice may not be true. For example, if a WIMP caries a conserved quantum number, and there are no “anti-WIMPs” - such as heavy Dirac neutrinos in the absence of antineutrinos - the annihilation is forbidden.
## 5 Discussion
We have shown that limits on recoil signals from underground detectors can be translated into bounds on the electro-magnetic form factors of WIMPs. Similar bounds can be inferred from the limits on energetic neutrino fluxes from the center of the earth and the sun, generated by annihilating WIMPs. Our analysis shows that all electromagnetic form factors have to be very small when expressed in characteristic nuclear units. In particular, the magnetic moment of a WIMP has to be five orders of magnitude smaller than the magnetic moment of a proton.
Unfortunately, all quoted bounds are trivially satisfied for the case of neutralino dark matter. Indeed, because it is a spin-1/2 particle, a neutralino cannot possess a quadrupole moment. In addition, because it is a Majorana fermion, a neutralino cannot have a magnetic or electric dipole moment either ($`\overline{\chi }\sigma _{\mu \nu }\chi 0`$). Finally, an anapole moment and a polarizability can be generated only radiatively. This results only in a small correction to the dominant neutralino-nucleus elastic cross section induced by squark exchange .
A magnetic moment can arise naturally if WIMPs are Dirac fermions. If we assume for a moment that the WIMP is a heavy Dirac neutrino, a large magnetic moment can be induced due to virtual $`W`$-charged lepton exchange. The value of the induced magnetic moment grows quadratically with the size of the Yukawa coupling until the point where the Yukawa interaction becomes essentially strong. In latter case a natural estimate for the size of the magnetic moment is $`e/m_\nu `$, with $`m_\nu `$ presumably of the order of 1 TeV or heavier. In this case, the interaction through the magnetic moment can be very important and, for light nuclei, can even dominate $`Z`$-boson exchange. It is interesting to note that the electric dipole moment of a WIMP is also very constrained. The limit of $`\mathrm{5\hspace{0.17em}10}^{21}`$ e cm is stronger than any other EDM constraint, apart from EDMs of electrons, neutrons and heavy atoms. The method we used to obtain this limit is in the spirit of the original paper by Purcell and Ramsey , where the first bound on the neutron EDM was inferred from neutron-nucleus scattering.
A model independent approach to the physics of WIMPs is partly motivated by recently discovered problems within the cold dark matter model. While large-scale structure can be described reasonably well, the numerical simulations of galactic substructures appear to be in conflict with observational data . A recently proposed cure, invoking self-interacting dark matter , is incompatible with the very restrictive neutralino model. This fact gives further impetus to our approach to study WIMPs in a manner as model independent as possible. The limits on the electro-magnetic form factors obtained in this work show that the amount of self-interaction, which could be induced by these form factors, is not sufficient to generate a cross section of WIMP-WIMP scattering at the level of $`10^{25}`$ cm<sup>2</sup> as required by the self-interacting WIMP scenario.
The authors would like to thank Keith Olive, Arkady Vainshtein and Viktor Zacek for usefull discussions. This work was supported in part by the Department of Energy under Grant No. DE-FG-02-94-ER-40823.
|
warning/0003/hep-ph0003291.html
|
ar5iv
|
text
|
# Pair production of neutralinos via gluon-gluon collisions 11footnote 1The project supported by National Natural Science Foundation of China
Abstract
The production of a neutralino pair via gluon-gluon fusion is studied in the minimal supersymmetric model(MSSM) at proton-proton colliders. The numerical analysis of their production rates are carried out in the mSUGRA scenario. The results show that this cross section may reach about 80 femto barn for $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_2^0`$ pair production and 23 femto barn for $`\stackrel{~}{\chi }_2^0\stackrel{~}{\chi }_2^0`$ pair production with suitable input parameters at the future LHC collider. It shows that this loop mediated process can be competitive with the quark-antiquark annihilation process at the LHC.
PACS number(s): 14.80.Ly, 12.15.Ji, 12.60.Jv
1. Introduction
The Standard Model(SM) is a successful theory of strong and electroweak interactions up to the present accessible energies. The hierarchy problem suggests that in principle the SM is the low energy effective theory of a more fundamental one. Supersymmetry (SUSY) is presently the most popular attempt to solve the hierarchy problem of the SM, where the cancellation of quadratic divergences is guaranteed and hence any mass scale is stable under radiative corrections. The most favorable candidate for a realistic extension of the SM is the minimal supersymmetric standard model(MSSM). In the MSSM, a discrete symmetry called R-parity is kept in order to assure baryon and lepton number conservations, because the gauge-coupling unification supports conservation of R-parity. It implies that there exists an absolutely stable lightest supersymmetric particle(LSP). Proper electroweak symmetry breaking induces the right properties of the LSP being a natural weakly interacting cold dark matter candidate, which can explain many astrophysical observations. In most cases the LSP in the MSSM may be the lightest Majorana fermionic neutralino. The neutralinos are mass eigenstates which are model dependent linear combinations of neutral gauginos and higgsinos. They are determined by diagonalizing the corresponding mass matrix. In the MSSM the mass matrix depends on four unknown parameters, namely $`\mu `$, $`M_2`$, $`M_1`$, and $`\mathrm{tan}\beta =v_2/v_1`$, the ratio of the vacuum expectation values of the two Higgs fields. $`\mu `$ is the supersymmetric Higgs-boson-mass parameter and $`M_2`$ and $`M_1`$ are the gaugino mass parameters associated with the $`SU(2)`$ and $`U(1)`$ subgroups, respectively.
Direct search of supersymmetric particles in experiment is one of the promising tasks for present and future colliders. The multi-TeV Large Hadron Collider(LHC) at CERN and the possible future Next Linear Collider(NLC) are elaborately designed in order to study the symmetry-breaking mechanism and new physics beyond the SM. When the LEP2 running will be terminated, the hadron colliders Tevatron and LHC will be the machines left in searching for supersymmetric particles. Therefore it is necessary to give a proper and full understanding of the production mechanisms of supersymmetric particles at hadron colliders. If supersymmetry really exists at TeV scale, SUSY particles should be discovered and it will be possible to make accurate measurements to determine their masses and other parameters of the Lagrangian at the LHC, and then we will have a better understanding of the supersymmetric model. We know that there are several mechanisms inducing the production of a chargino/neutralino pair at hadron colliders. One is through the quark-antiquark annihilation called Drell-Yan process, and another is via gluon-gluon fusion. Although the neutralino pair production via gluon-gluon fusion is a process induced by one-loop Feynman diagrams, the production rate can be still significant due to the large gluon luminosity in hadron colliders. The direct production channels of chargino/neutralinos and sleptons at the hadron colliders Tevatron and LHC $`p\overline{p}/pp\stackrel{~}{\chi }_i\stackrel{~}{\chi }_j+X`$ and $`\stackrel{~}{l}\overline{\stackrel{~}{l^{}}}+X`$ via quark-antiquark annihilation are investigated by Beenakker et al. A recent work showed that the one-loop process of the lightest chargino pair production via gluon-gluon fusion at the LHC can be considered as part of the NLO QCD correction to the quark-antiquark annihilation process. The neutralino pair production via quark-antiquark annihilation at the LHC was also considered by Han et al.. The measurement of CP violating supersymmetric phases in the chargino and neutralino pair production at the NLC was investigated by Barger et al. The lower mass limit of $`29.1GeV`$ on the lightest neutralino is obtained experimentally by analyzing the DELPHI results with the assumption that $`M_1/M_2{}_{}{}^{>}0.5`$,$`tan\beta =1`$, $`\mu =62.3GeV`$ and $`M_2=46.0GeV`$.
In this paper we concentrate on the direct neutralino pair production via gluon-gluon collisions at the LHC in the framework of the MSSM with complete one-loop Feynman diagrams. The numerical calculation will be illustrated in the CP conserving mSUGRA scenario with five input parameters, namely $`m_{1/2}`$, $`m_0`$, $`A_0`$, $`\mu `$ and $`\mathrm{tan}\beta `$, where $`m_{1/2}`$, $`m_0`$ and $`A_0`$ are the universal gaugino mass, scalar mass at GUT scale and the trilinear soft breaking parameter in the superpotential respectively. From these five parameters, all the masses and couplings of the model are determined by the evolution from the GUT scale down to the low electroweak scale. The paper is organized as follows: In section 2, we introduce the relevant features of the model and the analysis of the cross section in this work. In section 3, we discuss the numerical results of the cross sections and finally, a short summary is presented.
2. The Calculation of $`ppgg+X\stackrel{~}{\chi }_i^0\stackrel{~}{\chi }_j^0+X`$
In the MSSM the physical neutralino mass eigenstates $`\stackrel{~}{\chi }_i^0(i=1,2,3,4)`$ are the combinations of the neutral gauginos, $`\stackrel{~}{B}`$ and $`\stackrel{~}{W}^3`$, and the neutral higgsino, $`\stackrel{~}{H}_1^0`$ and $`\stackrel{~}{H}_2^0`$. In the two-component fermion fields $`\psi _j^0=(i\lambda ^{^{}},i\lambda ^3,\psi _{H_1^0},\psi _{H_2^0})`$, where $`\lambda ^{^{}}`$ is the bino and $`\lambda ^3`$ is the neutral wino, the neutralino mass term in the Lagrangian is given by
$$_M=\frac{1}{2}\left(\psi ^0\right)^TY\psi ^0+h.c.,$$
$`\left(2.1\right)`$
where the matrix $`Y`$ reads
$$Y=\left(\begin{array}{cccc}M_1& 0& m_Z\mathrm{sin}\theta _W\mathrm{cos}\beta & m_Z\mathrm{sin}\theta _W\mathrm{sin}\beta \\ 0& M_2& m_Z\mathrm{cos}\theta _W\mathrm{cos}\beta & m_Z\mathrm{cos}\theta _W\mathrm{sin}\beta \\ m_Z\mathrm{sin}\theta _W\mathrm{cos}\beta & m_Z\mathrm{cos}\theta _W\mathrm{cos}\beta & 0& \mu \\ m_Z\mathrm{sin}\theta _W\mathrm{sin}\beta & m_Z\mathrm{cos}\theta _W\mathrm{sin}\beta & \mu & 0\end{array}\right)$$
$`\left(2.2\right)`$
$`M_1`$, $`M_2`$ and $`\mu `$ can be complex and introduce CP-violating phases. By reparametrization of the fields, $`M_2`$ can be real and positive without loss of generality. The matrix $`Y`$ is symmetric and can be diagonalized by one unitary matrix $`N`$ such that $`N_D=N^{}YN^+=diag(m_{\stackrel{~}{\chi }_1^0},m_{\stackrel{~}{\chi }_2^0},m_{\stackrel{~}{\chi }_3^0},m_{\stackrel{~}{\chi }_4^0})`$ with the order of $`m_{\stackrel{~}{\chi }_1^0}m_{\stackrel{~}{\chi }_2^0}m_{\stackrel{~}{\chi }_3^0}m_{\stackrel{~}{\chi }_4^0}`$. Then the two-component mass eigenstates can be
$$\chi _i^0=N_{ij}\psi _j^0,i,j=1,\mathrm{},4.$$
$`\left(2.3\right)`$
The proper four-component mass eigenstates are the neutralinos which are defined in terms of two-component fields as
$$\stackrel{~}{\chi }_i^0=\left(\begin{array}{c}\chi _i^0\\ \overline{\chi }_i^0\end{array}\right)\left(i=1,\mathrm{},4\right),$$
$`\left(2.4\right)`$
and the mass term becomes
$$_m=\frac{1}{2}\underset{i}{}\stackrel{~}{M}_i\overline{\stackrel{~}{\chi }_i^0}\stackrel{~}{\chi }_i^0,$$
$`\left(2.5\right)`$
where $`\stackrel{~}{M}_i`$ are the diagonal elements of $`N_D`$.
The neutralino pair via gluon-gluon collisions can only be produced through one-loop diagrams in the lowest order. The calculation for this process can be simply carried out by summing all unrenormalized reducible and irreducible diagrams and the result will be finite and gauge invariant. In this work, we perform the calculation in the ’t Hooft-Feynman gauge. The generic Feynman diagrams contributing to the subprocess $`gg\stackrel{~}{\chi }_i^0\stackrel{~}{\chi }_j^0`$ in the MSSM are depicted in Fig.1, where the exchange of incoming gluons in Fig.1(a.1$``$3) and Fig.1(c.1$``$2) are not shown. All the one-loop diagrams can be divided into three groups: (1) box diagrams shown in Fig.1(a.1$``$3), (2) quartic interaction diagrams in Fig.1(b.1$``$2), (3) triangle diagrams shown in Fig.1(c.1$``$2). The $`Z^0`$ boson intermediated s-channel diagrams with quark or squark loops, as shown in Fig.1(b.2), Fig.1(c.1) and Fig.1(c.2), cannot contribute to the cross section. That can be explained in two fields. Firstly, the vector component of the $`Z^0`$ boson wave function does not couple to the initial gg state as the result of the Laudau-Yang theorem. Secondly, the CP-odd scalar component of the $`Z^0`$ boson does not couple to the invariant CP-even $`\stackrel{~}{\chi }_i^0\stackrel{~}{\chi }_j^0`$ state. We should mention that there are also some diagrams not contributing to the process, which we do not draw them in Fig.1(for example the s-channel diagrams with trilinear gluon interactions). Since the vertices of $`A^0(G^0)\stackrel{~}{q}\stackrel{~}{q}`$ vanish, there is no diagrams with a triangle squark loop coupling with an $`A^0`$ or $`G^0`$ Higgs boson.
We denote the reaction of neutralino pair production via gluon-gluon collision as:
$$g(p_3,\mu )g(p_4,\nu )\stackrel{~}{\chi }_i^0\left(p_1\right)\stackrel{~}{\chi }_j^0\left(p_2\right).$$
$`\left(2.6\right)`$
The corresponding matrix element for each of the diagrams can be written as
$$=^{\widehat{s}}+^{\widehat{t}}+^{\widehat{u}}$$
$`\left(2.7\right)`$
where $`^{\widehat{s}}`$, $`^{\widehat{t}}`$ and $`^{\widehat{u}}`$ are the matrix elements from s-channel, t-channel and u-channel diagrams, respectively. The amplitude parts from the u-channel box and triangle vertex interaction diagrams can be obtained from the t-channel’s by the following:
$$^{\widehat{u}}=(1)^{\delta _{ij}}^{\widehat{t}}(\widehat{t}\widehat{u},p_3p_4,\mu \nu ),$$
$`\left(2.8\right)`$
where the indices i and j are for the final neutralino. When $`i=j`$, there is a relative minus sign due to the Fermi statistics, which requires the amplitude to be antisymmetric under interchange of the two final identical fermions. The cross section for this subprocess at one-loop order in unpolarized gluon collisions can be obtained by
$$\widehat{\sigma }\left(\widehat{s},gg\stackrel{~}{\chi }_i^0\stackrel{~}{\chi }_j^0\right)=\frac{1}{16\pi \widehat{s}^2}\left(\frac{1}{2}\right)^{\delta _{ij}}_{\widehat{t}^{}}^{\widehat{t}^+}𝑑\widehat{t}\stackrel{}{}\left|\right|^2,$$
$`\left(2.9\right)`$
where $`\widehat{t}^\pm =1/2\left[(m_{\stackrel{~}{\chi }_i}^2+m_{\stackrel{~}{\chi }_j}^2\widehat{s})\pm \sqrt{(m_{\stackrel{~}{\chi }_i}^2+m_{\stackrel{~}{\chi }_j}^2\widehat{s})^24m_{\stackrel{~}{\chi }_i}^2m_{\stackrel{~}{\chi }_j}^2}\right]`$. The factor $`(\frac{1}{2})^{\delta _{ij}}`$ is due to the two identical particles in the final states. The bar over the sum means average over initial spins and colors.
With the results from Eq.(2.9), we can easily obtain the total cross section at $`pp`$ collider by folding the cross section of the subprocess $`\widehat{\sigma }(gg\stackrel{~}{\chi }_i^0\stackrel{~}{\chi }_j^0)`$ with the gluon luminosity,
$$\sigma \left(ppgg+X\stackrel{~}{\chi }_i^0\stackrel{~}{\chi }_j^0+X\right)=_{\left(m_{\stackrel{~}{\chi }_i^0}+m_{\stackrel{~}{\chi }_j^0}\right)^2/s}^1𝑑\tau \frac{d_{gg}}{d\tau }\widehat{\sigma }\left(gg\stackrel{~}{\chi }_i^0\stackrel{~}{\chi }_j^0,at\widehat{s}=\tau s\right).$$
$`\left(2.10\right)`$
where $`\sqrt{s}`$ and $`\sqrt{\widehat{s}}`$ denote the proton-proton and gluon-gluon c.m.s. energies respectively and $`\frac{d_{gg}}{d\tau }`$ is the gluon luminosity, which is defined as
$$\frac{d_{gg}}{d\tau }=_\tau ^1\frac{dx_1}{x_1}\left[f_g(x_1,Q^2)f_g(\frac{\tau }{x_1},Q^2)\right].$$
$`\left(2.11\right)`$
Here we used $`\tau =x_1x_2`$, the definitions of $`x_1`$ and $`x_2`$ are from Ref.. In our numerical calculation we adopt the MRST(mode 2) parton distribution function $`f_g(x_i,Q^2)`$ , and ignore the supersymmetric QCD corrections to the parton distribution functions for simplicity. The factorization scale $`Q`$ was chosen as the average of the final particle masses $`\frac{1}{2}(m_{\stackrel{~}{\chi }_i}+m_{\stackrel{~}{\chi }_j})`$. The numerical calculation is carried out for the LHC at the energy $`14TeV`$.
3. Numerical results and discussions
In this section, we present some numerical results of the total cross section from the full one-loop diagrams involving virtual (s)quarks for the process $`ppgg+X\stackrel{~}{\chi }_i^0\stackrel{~}{\chi }_j^0+X`$, respectively. The input parameters are chosen as: $`m_t=173.8GeV`$, $`m_Z=91.187GeV`$, $`m_b=4.5GeV`$, $`\mathrm{sin}^2\theta _W=0.2315`$, and $`\alpha =1/128`$. We adopt a simple one-loop formula for the running strong coupling constant $`\alpha _s`$
$$\alpha _s\left(\mu \right)=\frac{\alpha _s\left(m_Z\right)}{1+\frac{332n_f}{6\pi }\alpha _s\left(m_Z\right)\mathrm{ln}\left(\frac{\mu }{m_Z}\right)}.$$
$`\left(3.1\right)`$
where $`\alpha _s(m_Z)=0.117`$ and $`n_f`$ is the number of active flavors at energy scale $`\mu `$. In our numerical calculation to get the low energy scenario from the mSUGRA, the renormalization group equations(RGE’s) are run from the weak scale $`M_Z`$ up to the GUT scale, taking all threshold into account. We use two loop RGE’s only for the gauge couplings and the one-loop RGE’s for the other supersymmetric parameters. The GUT scale boundary conditions are imposed and the RGE’s are run back to $`M_Z`$, again taking threshold into account. The decay widths of the intermediate Higgs bosons are considered at the tree level and these formula can be found in ref..
The neutralino pair production cross sections dependence on the factorization scale Q is illustrated in Fig.2, where we choose the mSUGRA parameters as $`m_0=100GeV`$, $`m_{1/2}=150GeV`$, $`A=300GeV`$, $`\mathrm{tan}\beta =4`$ and $`\mu >0`$. The masses of $`\stackrel{~}{\chi }_1^0`$ $`\stackrel{~}{\chi }_2^0`$ are $`51.4GeV`$ and $`98GeV`$. For the $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_2^0`$ pair production, the cross section changes from $`28fb`$ to $`30fb`$ when the scale Q changes from $`0.2m_{\stackrel{~}{x}_2^0}`$ to $`m_{\stackrel{~}{x}_2^0}`$. For the $`\stackrel{~}{\chi }_2^0\stackrel{~}{\chi }_2^0`$ pair production, the cross section changes from $`14fb`$ to $`15.5fb`$ when the Q changes from $`0.2m_{\stackrel{~}{x}_2^0}`$ to $`m_{\stackrel{~}{x}_2^0}`$. When Q is larger than $`m_{\stackrel{~}{x}_2^0}`$, the cross sections of these two processes are nearly independent of the factorization scale.
In Fig.3 we show the cross sections of gaugino-like neutralino pair production and the higgsino-like neutralino pair production. For the gaugino-like neutralino pair production, we choose the same values of mSUGRA parameters adopted in Fig.2. But in the higgsino-like neutralino case, we choose $`M_1=150GeV`$, $`M_2=210GeV`$ and $`\mu =90GeV`$ in order to keep the masses of $`\stackrel{~}{\chi }_1^0`$ and $`\stackrel{~}{\chi }_2^0`$ being the same as those in the previous gaugino-like neutralino case, the other parameters of the higgsino-like neutralino case at the weak scale are chosen being the same as those in the gaugino-like neutralino case. Then the difference of the cross sections between gaugino-like neutralino pair production and higgsino-like neutralino pair production only comes from the change of the coupling of $`\stackrel{~}{\chi }_i^0\stackrel{~}{q}q`$ and $`\stackrel{~}{\chi }_i^0\stackrel{~}{\chi }_j^0H^0(A^0,h^0,G^0)`$ where the matrix elements of N are changed. For the $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_2^0`$ pair production, the cross section in the higgsino-like case is smaller than that in the gaugino-like case, but for the $`\stackrel{~}{\chi }_2^0\stackrel{~}{\chi }_2^0`$ pair production, the cross section in the higgsino-like case is larger than that in the gaugino-like case.
The cross sections for $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_2^0`$ and $`\stackrel{~}{\chi }_2^0\stackrel{~}{\chi }_2^0`$ productions at the LHC versus the mass of $`\stackrel{~}{\chi }_2^0`$ is shown in Fig.4. We assume that $`\stackrel{~}{\chi }_1^0`$ is the LSP and escapes detection. The input parameters are chosen as $`m_0=100GeV`$, $`A_0=300GeV`$, $`\mathrm{tan}\beta =4`$, $`\mu >0`$ and $`m_{1/2}`$ varying from $`130GeV`$ to $`325GeV`$. In the framework of the mSUGRA, the masses of $`\stackrel{~}{\chi }_1^0`$ and $`\stackrel{~}{\chi }_2^0`$ increase from $`41GeV`$ to $`129GeV`$ and from $`81GeV`$ to $`251GeV`$ respectively, and the masses of Higgs boson $`H^0`$ and $`A^0`$ increase from $`230GeV`$ to $`509GeV`$ and from $`226GeV`$ to $`507GeV`$ respectively with our input parameters. In our case, the mass of the Higgs boson $`h^0`$ is always smaller than $`m_{\stackrel{~}{\chi }_1^0}+m_{\stackrel{~}{\chi }_2^0}`$ and the contribution to the cross section of $`h^0`$ is very small due to the s-channel suppression. Since the masses of $`H^0`$ and $`A^0`$ are larger than $`m_{\stackrel{~}{\chi }_1^0}+m_{\stackrel{~}{\chi }_2^0}`$, the cross section will be strongly enhanced by the s-channel resonance effects when the total energy $`\sqrt{\widehat{s}}`$ of the subprocess approaches the mass of $`H^0`$ or $`A^0`$. Since $`m_{\stackrel{~}{\chi }_2^0}>m_{\stackrel{~}{\chi }_1^0}`$, the pair production of $`\stackrel{~}{\chi }_2^0\stackrel{~}{\chi }_2^0`$ compared to $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_2^0`$ is kinematically suppressed. For the $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_2^0`$ pair production, the cross section can reach 43 femto barn when $`m_{1/2}=130GeV`$, $`m_{\stackrel{~}{\chi }_1^0}=41GeV`$ and $`m_{\stackrel{~}{\chi }_2^0}=81GeV`$. The cross section of $`\stackrel{~}{\chi }_2^0\stackrel{~}{\chi }_2^0`$ pair production can reach 23 femto barn when $`m_{1/2}=130GeV`$. From the results of , we know that the cross sections of $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_2^0`$ and $`\stackrel{~}{\chi }_2^0\stackrel{~}{\chi }_2^0`$ pair production via quark-antiquark annihilation at the LHC are about 70 femto barn and 300 femto barn respectively, when the input parameters are chosen as $`m_0=150GeV`$, $`A_0=300GeV`$, $`\mathrm{tan}\beta =4`$, $`\mu >0`$ and $`m_{\stackrel{~}{\chi }_2^0}=81GeV`$. The pair production rate of $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_2^0`$ via gluon-gluon fusion is about few ten percent of that via quark-antiquark annihilation at the LHC with the same input parameters. The cross section of $`\stackrel{~}{\chi }_2^0\stackrel{~}{\chi }_2^0`$ via gluon-gluon fusion is about few percent of that via quark-antiquark annihilation at the LHC. The neutralino pair production via gluon-gluon fusion at the LHC is of the same order of the next leading order(NLO) QCD correction to the quark-antiquark production process. The neutralino pair production rate via gluon-gluon fusion should be considered in studying the pair pair productions of neutralinos at proton-proton colliders.
In Fig.5 we present the cross sections of neutralino pair productions versus $`\mathrm{tan}\beta `$ where the other four input parameters are chosen as $`m_0=100GeV`$, $`A_0=300GeV`$, $`\mu >0`$ and $`m_{1/2}=150GeV`$. When $`\mathrm{tan}\beta `$ increases from 4 to 32, the mass of $`\stackrel{~}{\chi }_1^0`$ increases from $`51GeV`$ to $`57GeV`$ and the mass of $`\stackrel{~}{\chi }_2^0`$ increases from $`98GeV`$ to $`102.5GeV`$. So the masses of neutralinos are nearly independent of $`\mathrm{tan}\beta `$. The masses of Higgs boson $`H^0`$ and $`A^0`$ depend on $`\mathrm{tan}\beta `$ and decrease from $`258GeV`$ to $`165GeV`$ and from $`254GeV`$ to $`165GeV`$ respectively. The mass of $`h^0`$ is a function of $`\mathrm{tan}\beta `$ too, but keeps $`m_{h^0}<m_{\stackrel{~}{\chi }_1^0}+m_{\stackrel{~}{\chi }_2^0}`$. Since the couplings of Higgs bosons to quarks and squarks pair are related to the ratio of the vacuum expectation values, $`\mathrm{tan}\beta `$ should effect the cross sections substantially. The cross sections of the neutralino pair productions decrease when $`\mathrm{tan}\beta `$ goes from 4 to 6, and increase when $`20>\mathrm{tan}\beta >6`$ due to the effect of the coupling strength of Higgs bosons to quarks and squarks pair. For the $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_2^0`$ pair production, when $`tan\beta <27`$, the masses of $`H^0`$ and $`A^0`$ are larger than $`m_{\stackrel{~}{\chi }_1^0}+m_{\stackrel{~}{\chi }_2^0}`$, the s-channel resonance effects of $`H^0`$ and $`A^0`$ on the cross section can be enhanced, while when $`tan\beta >27`$, the masses of $`H^0`$ and $`A^0`$ become smaller than $`m_{\stackrel{~}{\chi }_1^0}+m_{\stackrel{~}{\chi }_2^0}`$, then there will be no s-channel resonance effects and the cross section goes down rapidly. For the similar reason, the cross section of the $`\stackrel{~}{\chi }_2^0\stackrel{~}{\chi }_2^0`$ pair production increase in the range of $`20>\mathrm{tan}\beta >6`$ due to the effect of the Yukawa couplings of Higgs bosons to quarks and squarks pair and the s-channel resonance effects of the intermediate Higgs bosons, but it goes down again when $`\mathrm{tan}\beta >20`$ where the masses of $`H^0`$ and $`A^0`$ become smaller than $`m_{\stackrel{~}{\chi }_2^0}+m_{\stackrel{~}{\chi }_2^0}`$ and the s-channel resonance effects of the Higgs bosons disappear. From the figure we can see that if we choose suitable input parameters, the cross section can be largely enhanced. Especially the cross section of the $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_2^0`$ pair production can even reach about 80 femto barn when $`\mathrm{tan}\beta 27`$.
In our calculation we find that the discrepancy between our results from the newer sets of gluon densities in the MRST(mode 2) scheme and the older sets in the MRS(set G) scheme is very small and cannot be distinguished in figures.
4. Summary
In this paper, we studied the pair production process of the neutralino via gluon-gluon fusion at the LHC collider. The numerical analysis of their production rates is carried out in the mSUGRA scenario with some typical parameter sets. The results show that the cross section of the neutralino pair-production via gluon-gluon fusion can reach about 80 femto barn for $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_2^0`$ pair production and 23 femto barn for $`\stackrel{~}{\chi }_2^0\stackrel{~}{\chi }_2^0`$ pair production at the LHC collider. In some c.m.s energy regions of incoming gluons, where the s-channel resonance conditions are satisfied in the parameter space, we can see observable enhancement effects of the cross sections. We conclude that neutralino pair production via gluon-gluon fusion can be competitive with the quark-antiquark annihilation production process at the LHC and should be considered as a part of the NLO QCD correction to the quark-antiquark annihilation production process.
Acknowledgement: This work was supported in part by the National Natural Science Foundation of China(project number: 19875049), the Youth Science Foundation of the University of Science and Technology of China, a grant from the Education Ministry of China and the exchange program between China and Austria(Project number V.A.12). One of the authors Y. Jiang would like to thank Prof. F.F. Schöberl for valuable discussion in this context.
Figure Captions
Fig.1 The Feynman diagrams of the subprocess $`gg\stackrel{~}{\chi }_i^0\stackrel{~}{\chi }_j^0`$.
Fig.2 Dependence of the cross sections for the productions of $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_2^0`$ and $`\stackrel{~}{\chi }_2^0\stackrel{~}{\chi }_2^0`$ pairs on the factorization scale Q with the mSUGRA parameters $`m_0=100GeV`$, $`m_{1/2}=150GeV`$, $`A=300GeV`$, $`\mathrm{tan}\beta =4`$ and $`\mu >0`$. The masses of $`\stackrel{~}{\chi }_1^0`$ $`\stackrel{~}{\chi }_2^0`$ are $`51.4GeV`$ and $`98GeV`$.
Fig.3 The cross sections for the neutralino pair production at the gaugino-like neutralino case and higgsino-like neutralino case at the LHC. For the gaugino-like neutralino case , the parameters are chosen in mSUGRA scenario with $`m_0=100GeV`$, $`m_{1/2}=150GeV`$, $`A=300GeV`$, $`\mathrm{tan}\beta =4`$ and $`\mu >0`$. For the higgsino-like neutralino case, we choose $`M_1=150GeV`$, $`M_2=210GeV`$, $`\mu =90GeV`$ and the other parameters the same as those in the gaugino-like neutralino case.
Fig.4 Total cross sections of the process $`ppgg+X\stackrel{~}{\chi }_i^0\stackrel{~}{\chi }_j^0+X`$ as function of $`m_{\stackrel{~}{\chi }_2^0}`$ at the LHC. The solid curve is for $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_2^0`$ pair production and the dashed curve is for $`\stackrel{~}{\chi }_2^0\stackrel{~}{\chi }_2^0`$ pair production.
Fig.5 Total cross sections of the process $`ppgg+X\stackrel{~}{\chi }_i^0\stackrel{~}{\chi }_j^0+X`$ as function of $`\mathrm{tan}\beta `$ at the LHC. The solid curve is for $`\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_2^0`$ pair production and the dashed curve is for $`\stackrel{~}{\chi }_2^0\stackrel{~}{\chi }_2^0`$ pair production.
|
warning/0003/cond-mat0003010.html
|
ar5iv
|
text
|
# Dependence of the physical properties of 𝑁𝑑_0.5𝐶𝑎_0.5𝑀𝑛𝑂_{3+𝛿} on the oxydation state of 𝑀𝑛.
## Abstract
We present a comparative study of the magnetic, transport and structural properties of $`Nd_{0.5}Ca_{0.5}MnO_{3+\delta }`$ \[$`\delta =`$ 0.02(1) and 0.04(1)\]. We have found significant differences between the low temperature magnetic orders and the magnetization curves below $`T_{CO}250K`$ of the two samples. In particular one component of the magnetic moment presents a ferromagnetic coupling between the $`(\mathrm{0\hspace{0.33em}0\hspace{0.33em}1})`$ planes ($`Pbnm`$ setting) deviating the angle between neighboring $`Mn`$ ions from $`180^o`$ (perfect CE order) to $`150^o`$ \[$`\delta =0.02(1)`$\] and $`130^o`$ \[$`\delta =0.04(1)`$\]. These results imply a remarkable $`\delta `$ dependence which is discussed in the light of a non-random spatial distribution of defects in the perfect charge order scheme.
A considerable research effort concerning the charge order (CO) phenomenon in doped manganites ($`Ln_{1x}A_xMnO_3`$ with $`Ln=`$ rare earth and $`A=`$ alkaline earth) has been made during the last years. Ingredients such as Coulomb repulsion, the effect of Jahn-Teller distortion on the $`e_g`$ energy levels and lattice distortions play an important role on the stability of the CO state. This state is usually accompanied by a real space ordering of the orbitals occupied by the $`e_g`$ electrons. When this orbital order (OO) occurs, the $`Mn`$-$`O`$-$`Mn`$ superexchange interactions are FM if a half-filled $`d_{3r_i^2r^2}`$ ($`r_i=x,y`$ or $`z`$) and an empty $`e_g`$ orbitals are involved and antiferromagnetic (AFM) if the involved ones are two $`t_{2g}`$ orbitals. This makes the magnetic behavior to be strongly dependent on (i) the mean valence state of the $`Mn`$ ions (or, equivalently, on the density of $`e_g`$ electrons) and (ii) the real space OO. When the formal ratio $`\frac{Mn^{+4}}{Mn^{+3}}`$ ions is 1, and for a rich variety of rare earths the low temperature magnetic structure found is the CE. This structure is usually explained as resulting from a well established CO and OO. Some previous works have been devoted to study how the off-stoichiometry of the $`\frac{Mn^{+4}}{Mn^{+3}}=1`$ ratio affects the CO. One interesting effect is the variation in the CO modulation wave vector \[$`𝑸=\frac{2\pi }{a}(\mathrm{0\hspace{0.33em}1}/2\alpha \mathrm{\hspace{0.33em}0})`$ in the $`Pbnm`$ setting\] that different annealing conditions produce, for example, in $`Pr_{0.5}Ca_{0.5}MnO_3`$ and in $`Sm_{0.5}Ca_{0.5}MnO_3`$. Incommensurability of the CO in $`La_{0.5}Ca_{0.5}MnO_3`$ have also been reported in the temperature range of the structural transformation, but it disappears at lower temperature. Barnabé et al. (Ref. ) argue that in an under-doped sample (presenting CO) the extra $`e_g`$ electrons are randomly placed in the $`Mn`$ positions leaving the CO commensurated ($`\alpha =0`$), but the extra $`Mn^{+4}`$ ions of an over-doped sample introduce incommensurability in the CO structure ($`\alpha >0`$) even when the excess of $`Mn^{+4}`$ is very small. Using electron diffraction data, Barnabé et al. conclude that the extra $`Mn^{+4}`$ ions form ordered $`(\mathrm{0\hspace{0.33em}1\hspace{0.33em}0})`$ planes. Such an incommensurability of an over-doped sample has been also reported by Chen et al. in Ref. .
In this brief report we present the effects of changing the oxydation state of $`Mn`$ in $`Nd_{0.5}Ca_{0.5}MnO_{3+\delta }`$. The notation we use here is convenient from the point of view of the chemical formulation, but the actual structural defects are not interstitial oxygens, they should correspond to cation vacancies. This study is based on the comparison of the structural, magnetic and transport properties. The main changes correspond to the magnetic structure. Its evolution with $`\delta `$ indicates that the extra holes due to $`\delta `$ are not randomly distributed.
Two different polycrystalline samples have been sintered by standard solid state reaction by mixing high purity powders of $`CaCO_3`$, $`Mn_2O_3`$ and $`Nd_2O_3`$ in appropiate ratios for the nominal composition $`Nd_{1/2}Ca_{1/2}MnO_3`$. After some intermediate treatments, the two samples were pressed into pellets, fired and ground again for several times. For the first sample (air sample) the firings were done at $`1450^oC`$ in air followed by a rapid quench to RT ($`500^oC/hour`$). For the second sample ($`O_2`$ sample) the firings were done at $`1400^oC`$ in an atmosphere of flowing $`O_2`$ and the final one was followed by a slow cooling to RT ($`50^oC/hour`$). X-ray powder diffraction confirmed that both samples are well crystallized in a single phase. Agreeing with the synthesis conditions, thermogravimetric analysis (TGA), used to determine the oxygen content of the samples ($`Nd_{0.5}Ca_{0.5}MnO_{3+\delta }`$), evidenced a larger content of oxygen in the $`O_2`$ ($`\delta =0.04(1)`$, $`Mn^{+4}=58(2)\%`$) than in air sample ($`\delta =0.02(1)`$, $`Mn^{+4}=54(2)\%`$). Neutron diffraction (ND) patterns of the air sample were collected at the Institute Laue Langevin (Grenoble) using D2B ($`\lambda =1.594\mathrm{\AA }`$) and D1B ($`\lambda =2.52\mathrm{\AA }`$) diffractometers. ND patterns of the $`O_2`$ sample were collected at the Laboratoire Leon Brillouin (Paris) using 3T2 ($`\lambda =1.227\mathrm{\AA }`$) and G4.2 ($`\lambda =2.426\mathrm{\AA }`$) diffractometers. For both samples, ND patterns were collected for several temperatures in the range $`1.5K`$ to room temperature (RT). They were analyzed by the Rietveld method using the program FULLPROF. Resistivity was measured by the four-probe method using a commercial PPMS (Quantum Design). Magnetization measurements have been carried out using a commercial SQUID (Quantum Design).
Figure 1 shows the high-resolution ND patterns, collected at RT, refined using an average $`Pbnm`$ structure. The refined lattice parameters, $`Mn`$-$`O`$ bond lengths and $`Mn`$-$`O`$-$`Mn`$ bond angles are listed in Tab. I for comparison. There is a good agreement between the structural parameters obtained for $`Nd_{0.5}Ca_{0.5}MnO_{3.02(1)}`$ with those for $`Nd_{0.5}Ca_{0.5}MnO_{3.04(1)}`$ and with those previously reported for $`Nd_{0.5}Ca_{0.5}MnO_3`$. The lattice parameters show a negligible compression of the $`c`$-axis. The $`MnO_6`$ octahedra are very regular and without any appreciable apical compression. In order to check the mean oxidation state of the $`Mn`$ ions using ND data, we refined the occupation factors of $`Nd`$, $`Ca`$ and $`Mn`$ ions assuming that no interstitial oxygens are present (that is, fixing the oxygen stoichiometry to 3). It is known that any excess of oxygen in perovskites always corresponds to the presence of cation vacancies since there is no space left in the perovskite structure to allocate interstitial oxygens. We observed no improvement of the refinement factors and no significant differences between the nominal composition and the refined values: $`Nd_{0.502(8)}Ca_{0.492(8)}Mn_{1.001(9)}O_3`$ for the air sample $`Nd_{0.496(8)}Ca_{0.492(8)}Mn_{0.992(8)}O_3`$ for the $`O_2`$ sample.
ND patterns of both samples present superlattice peaks below $`T_{CO}250K`$, attributed to the development of the CO/OO state. A small incommensurability of the charge order can be appreciated through the position of the $`(\mathrm{2\hspace{0.33em}1}/2\alpha \mathrm{\hspace{0.33em}2})`$ peak for the $`Nd_{0.5}Ca_{0.5}MnO_{3.04(1)}`$ sample \[$`\alpha =0.03(1)`$\] but not, within the resolution of our ND data, for the $`Nd_{0.5}Ca_{0.5}MnO_{3.02(1)}`$ sample ($`\alpha 0`$). The inset of Fig. 2(a) shows the growth of the $`(\mathrm{2\hspace{0.33em}1}/2\alpha \mathrm{\hspace{0.33em}2})`$ peak found for $`Nd_{0.5}Ca_{0.5}MnO_{3.02(1)}`$. Its integrated intensity relative to that of the $`(\mathrm{0\hspace{0.33em}2\hspace{0.33em}2}),(\mathrm{2\hspace{0.33em}0\hspace{0.33em}2})`$ positions is, at low temperature, $`I_{(\mathrm{2\hspace{0.33em}1}/2\alpha \mathrm{\hspace{0.33em}2})}/I_{(\mathrm{0\hspace{0.33em}2\hspace{0.33em}2}),(\mathrm{2\hspace{0.33em}0\hspace{0.33em}2})}=\mathrm{7.4\hspace{0.17em}10}^3`$ for $`Nd_{0.5}Ca_{0.5}MnO_{3.02(1)}`$ and $`\mathrm{7.1\hspace{0.17em}10}^3`$ for $`Nd_{0.5}Ca_{0.5}MnO_{3.04(1)}`$, indicating a very similar degree of long range ordering in both samples. The effect on the lattice parameters of the CO/OO is shown in the insets of Fig. 1. The strong compression of the $`c`$ lattice parameter is due to the localization of the $`e_g`$ electrons in $`d_{3x^2r^2}`$ or $`d_{3y^2r^2}`$ orbitals contained in the $`(\mathrm{0\hspace{0.33em}0\hspace{0.33em}1})`$ planes. The low temperature ND patterns have been refined using the $`Pbnm`$ description of the single cell. The obtained lattice parameters, $`Mn`$-$`O`$ bond distances and $`Mn`$-$`O`$-$`Mn`$ bond angles are listed in Tab. I. The aforementioned in-plane localization of the $`e_g`$ electrons causes an apical compression of the $`MnO_6`$ octahedra that can be quantified through the parameter $`ϵ_d=\left|1\frac{d_{MnO1}}{d_{MnO2}}\right|\times 10^4`$. $`ϵ_d`$ is, at low temperature, about four times larger than at RT for both samples. $`La_{1/2}Ca_{1/2}MnO_3`$ ($`ϵ_d=205`$) and $`Nd_{1/2}Sr_{1/2}MnO_3`$ ($`ϵ_d=215`$), both presenting CO and a CE-type magnetic structure, present an apical compression comparable to $`Nd_{0.5}Ca_{0.5}MnO_{3.02(1)}`$ ($`ϵ_d=217`$). A little lower value is found for $`Nd_{0.5}Ca_{0.5}MnO_{3.04(1)}`$ ($`ϵ_d=192`$), reflecting the partial lack of $`e_g`$ electrons.
Figure 2(a) shows the temperature evolution of the low angle region of the ND patterns of $`Nd_{0.5}Ca_{0.5}MnO_{3.02(1)}`$. The apparition of magnetic reflections is observed to occur at $`T_{N1}=165K`$ \[$`Nd_{0.5}Ca_{0.5}MnO_{3.02(1)}`$\] and $`T_{N2}=150K`$ \[$`Nd_{0.5}Ca_{0.5}MnO_{3.04(1)}`$\]. The low angle regions of the refined ND patterns at low temperature are plotted in Fig. 3 for both samples. The magnetic peaks have been indexed on the basis of the $`Pbnm`$ setting. In both samples, there are two families of magnetic peaks displaying $`(hkl)`$ Miller indices with $`h`$ half integer and k both integer or half integer. The first family of peaks has $`l=`$odd \[$`(1/\mathrm{2\hspace{0.33em}0\hspace{0.33em}1})`$, $`(1/\mathrm{2\hspace{0.33em}1}/\mathrm{2\hspace{0.33em}1})`$,…\] and is associated to the CE-type AFM structure. The second family of peaks has $`l=`$even \[$`(1/\mathrm{2\hspace{0.33em}0\hspace{0.33em}0})`$, $`(1/\mathrm{2\hspace{0.33em}1}/\mathrm{2\hspace{0.33em}0})`$,…\] and is usually associated to the pseudo CE-type AFM structure. These two families of peaks are also present in $`Nd_{0.45}Ca_{0.55}MnO_3`$. We have obtained good fits to the ND patterns assuming that the components of the magnetic moment contained in the $`(\mathrm{0\hspace{0.33em}0\hspace{0.33em}1})`$ planes ($`m_x`$ and $`m_y`$) present a CE-type order structure, while the component of the magnetic moment in the $`[\mathrm{0\hspace{0.33em}0\hspace{0.33em}1}]`$ direction ($`m_z`$) presents a pseudo CE-type AFM order. The obtained magnetic structures are described in Tab. II and schematically plotted in Fig. 2(b). In both cases the ordered moments are well below the expected values for perfectly ordered moments. The effect of the considered out-of-plane component of the magnetic moment can be interpreted as a deviation from the perfect value ($`180^o`$) of the angles formed by the magnetic moments in neighboring $`(\mathrm{0\hspace{0.33em}0\hspace{0.33em}1})`$ planes \[e.g. between the magnetic moments of the $`Mn^{+3}`$ ions in positions $`(\frac{1}{4},\frac{1}{2},0)`$ and $`(\frac{1}{4},\frac{1}{2},\frac{1}{2})`$ of Tab. II as is schematically depicted in Fig. 2(b)\] that is about $`30^o`$ in $`Nd_{0.5}Ca_{0.5}MnO_{3.02(1)}`$ and $`50^o`$ in $`Nd_{0.5}Ca_{0.5}MnO_{3.04(1)}`$.
The magnetization \[M(T)\] and resistivity \[$`\rho (T)`$\] curves are shown in Fig. 4 and its inset respectively. The CO transition can be well appreciated as a sudden increase in the $`\rho (T)`$ curves, caused by the localization of the charges, and as a local maximum in the $`M(T)`$ curves. This local maximum is due to the transition from the FM $`Mn`$-$`Mn`$ correlations above $`T_{CO}`$ (double-exchange) to the AFM $`Mn`$-$`Mn`$ correlations below $`T_{CO}`$ ($`Mn`$-$`O`$-$`Mn`$ superexchange). There is scarcely any variance in $`\rho (T)`$ from one sample to the other. In this sense, the magnetization curves clearly show two regimes: above $`T_{CO}`$ both curves are almost identical, but below this point the curves systematically start to come away.
The results presented in the previous paragraphs show that the two samples are indeed very similar: the different oxydation state of the $`Mn`$ does not significantly change the structure, the transport and the magnetic properties of the compound above $`T_{CO}`$. Consequently, the differences observed at lower temperatures cannot be attributed to intrinsic differences in the structure due to the different concentration of cation vacancies.
In contrast with the behavior above $`T_{CO}`$, substantial discrepancies in the $`M(T)`$ curves start at this temperature, just when the $`Mn`$-$`O`$-$`Mn`$ superexchange interactions are dominated by the real space location of the $`e_g`$ electrons. This indicates that the difference in the magnetic behavior of the samples is driven by the defects introduced by the lack of $`e_g`$ electrons in the OO state. The different density of carriers causes differences in the real space distribution of $`Mn`$-$`O`$-$`Mn`$ superexchange interactions giving rise to a dissimilar $`M(T)`$ evolution. The presence of defects in the OO structure gives rise to frustrated magnetic coupling and, hence, differences between the magnetic structure and the perfect CE-type. To be emphasized is that such great changes in the obtained magnetic order and in the $`M(T)`$ are hard to explain in the scenario of a random location of the holes among $`Mn`$ positions. The last will simply lead to a CE structure with disorder. Consequently, a certain spatial grouping of the defects may enhance their effects and also cause the small incommensurability of the CO.
Of special interest is to determine if the two families of magnetic peaks are sustained by the same structural lattice or they correspond to two slightly different lattices. In $`Nd_{0.5}Sr_{0.5}MnO_3`$, the coexistence of an A-type and a CE-type AFM phases below the CO transition, supported by different crystallographic cells, has been attributed to inhomogeneities in the cation distribution. In the present case (for both samples) the two families of magnetic peaks can be very well reproduced using the same lattice parameters and appear at the same temperature \[see Fig. 2(a)\]. Within our resolution, a single set of cell parameters allows a proper indexation of all the magnetic intensities (see Fig. 3).
In summary, off-stoichiometric samples of $`Nd_{0.5}Ca_{0.5}MnO_{3.02(1)}`$ \[$`\%Mn^{+4}=54(2)`$\] and $`Nd_{0.5}Ca_{0.5}MnO_{3.04(1)}`$ \[$`\%Mn^{+4}=58(2)`$\] samples have been investigated in comparison with the stoichiometric $`Nd_{0.5}Ca_{0.5}MnO_3`$. The same structural features and macroscopic behavior are observed above $`T_{CO}`$. A detailed ND study reveals pronounced influence of the different, but very similar, oxidation state of $`Mn`$ upon the magnetic long range ordering. The intensities of the magnetic $`(h/2k/\mathrm{2\hspace{0.33em}0})`$ peaks (absent in the CE-type order) increase significantly with the value of $`\delta `$. As a result the angle formed by the spins of successive $`MnO_2`$ $`(\mathrm{0\hspace{0.33em}0\hspace{0.33em}1})`$ planes change from $`\theta =180^o`$ for $`\delta =0`$ to $`\theta =130^o`$ for $`\delta =0.04`$. These remarkable changes in the magnetic structure are unusually important compared to the effects of a small off-stoichiometry in most of the transition metal oxides. They indicate, in connection with the incommensurability of the CO detected for the $`\delta =0.04(1)`$ sample, that the extra holes are not randomly placed in the $`Mn`$ positions but they form spatial sub-structures.
This work has been done with financial support from the CICyT (MAT97-0669) MEC (PB97-1175) and Generalitat de Catalunya (GRQ95-8029). A.L. acknowledges financial support form the Oxide Spin Electronics Network (EU TMR) program. The ILL and LLB are acknowledged for making available the beam time.
|
warning/0003/hep-lat0003014.html
|
ar5iv
|
text
|
# Percolation and Magnetization in the Continuous Spin Ising Model
## 1 Introduction
It has recently been proposed that the deconfinement transition in SU(2) gauge theory can be characterized as percolation of Polyakov loop clusters \[forsatz\]. This idea is based on the fact that SU(2) gauge theory and the Ising model belong to the same universality class \[Svet\], and that the magnetization transition in the Ising model can be specified as percolation of clusters defined through local bond weights \[Fortuin\]. It thus seems natural that the same holds for the corresponding Polyakov loop clusters in SU(2) gauge theory, provided suitable bond weights can be defined. In \[forsatz\], it was shown that in a lattice formulation effectively corresponding to the strong coupling limit, SU(2) gauge theory indeed leads to the predicted Ising critical exponents.
In the present paper, we want to consider the classical continuous spin model on $`Z^2`$ introduced by Griffiths \[Gr\], an Ising model with spins taking values continuously between $`1`$ and $`+1`$, and prove that in this case magnetization and percolation transitions coincide. We also present a detailed numerical study of the model with simulations on the 2d square lattice.
For SU(2) lattice gauge theory it was shown in the strong coupling limit \[Ka\] that the partition function can be written in a form which, apart from a factor which depends on the group measure, is the partition function of the classical spin model, with spins varying continuously in some bounded real interval \[Gr\]. These continuous spins thus are the natural counterparts of the Polyakov loops in SU(2) gauge theory. Following the pioneering work \[Fortuin\], the relationship between the thermodynamical features of spin models and the properties of the corresponding geometrical clusters have in the past years received considerable attention (see \[SW\] \- \[Cha2\]). We want to address here the random percolation content of the continuous spin Ising model \[Gr\] obtained through a Fortuin-Kasteleyn transformation \[Fortuin\] of the partition function. A specific class of geometrical clusters \[Wo\] shows a percolation transition whose critical behaviour matches the conventional Ising counterparts. In particular, it can be shown that up to finite constants, the spin magnetization equals the probability of long range cluster connectivity; moreover, the spin susceptibility becomes equal (up to the same constants) to some mean cluster size \[BCG\].
The numerical method to be used here is the Wolff algorithm \[Wo\]. It is based on the formulation of a single cluster update algorithm for spin systems with unprecedented performances regarding the problem of critical slowing down behaviour near phase transition \[Wo\]. It turns out that the Wolff clusters exhibit critical percolation behaviour at the point of the Ising phase transition. Very recently, this specific property of Wolff clusters has been put into a rigorous framework \[Cha1\], leading to a concept of Wolff measures which provides the theoretical basis for the method to be used below. We shall extend this work to the 2-d continuous spin model and corroborate the results through numerical simulations.
## 2 The Model.
Let $`\mathrm{\Lambda }Z^2`$ be a finite lattice with edge set $`E`$ and vertex set $`V=\{1,..,N\}`$ and $`𝐒`$ a random configuration of spin variables $`𝐒=\{S_i\}\mathrm{\Sigma }^V`$, where the $`S_{i,iV}`$ take values independently in $`\mathrm{\Sigma }=[1,+1]`$. Consider the Hamiltonian
$$H(𝐒)=\underset{<i,j>}{\mathbf{}}J_{i,j}S_iS_j,$$
(2.1)
where the sum is over all nearest neighbour pairs ($`<i,j>`$) of spins and $`J_{i,j}>0`$. Rewriting the spin variables $`S_i=l_i\sigma _i`$, where $`\sigma _i=sign(S_i)=\pm 1`$ and $`l_i=|S_i|[0,1]`$, one has
$$H(𝐒)=H(𝝈,𝒍)=\underset{i,j}{\mathbf{}}J_{i,j}l_il_j\sigma _i\sigma _j$$
(2.2)
### 2.1 Joint Distribution over Spin Amplitudes and Bond Variables
Following Edwards and Sokal \[ES\], the random cluster representation ( RCR ) \[Fortuin, G\] of this model is derived in the following way. Let $`𝒏\{0,1\}^E`$ be a configuration of bond variables on $`\mathrm{\Lambda }`$ where,
1. $`n_{ij}=1`$ means that there is a bond between sites $`i`$ and $`j`$ in $`V`$,
2. $`n_{ij}=0`$ is the opposite event.
Consider the following joint probability distribution of the random variables $`(\sigma ,l,n)`$,
$`\mathrm{I}\mathrm{P}(𝐒,𝒏)=\mathrm{I}\mathrm{P}(𝒍,𝝈,𝒏)`$ $`=`$ $`Z_\mathrm{\Lambda }^1{\displaystyle \underset{<i,j>}{\mathbf{}}}e^{\beta J_{i,j}l_il_j}\left[p_{i,j}\delta _{n_{ij,1}}\delta _{\sigma _i\sigma _j}+(1p_{ij})\delta _{n_{ij,0}}\right]`$ (2.3)
$`=`$ $`Z_\mathrm{\Lambda }^1{\displaystyle \underset{<i,j>}{\mathbf{}}}w_{\beta ,J}(𝒍,𝝈,𝒏),`$
where $`p_{i,j}=1e^{2\beta J_{i,j}l_il_j}`$ is a weighting factor to be interpreted as the probability for a link $`n_{ij}`$ to exist ($`n_{ij}=1`$) between two nearest neighbour spins $`S_i`$ and $`S_j`$ that share the same sign ($`\delta _{\sigma _i\sigma _j}=1`$), and $`(1p_{i,j})`$ the probability of having no link ($`n_{ij}=0`$) between them. $`Z_\mathrm{\Lambda }`$ is the partition function of the model given by
$$Z_\mathrm{\Lambda }=\underset{\{\sigma \}}{\mathbf{}}\mathbf{}\underset{i}{\mathbf{}}df_i(l_i)\mathrm{exp}[\beta \underset{<i,j>}{\mathbf{}}J_{i,j}l_il_j\sigma _i\sigma _j],$$
(2.4)
where $`df_i`$ is the probability distribution of the spin amplitudes $`l_i`$ and $`\beta `$ is the inverse temperature.
Starting with the joint distribution (2.3), it is straightforward to express the weights $`w_{\beta ,J}(𝒍,𝝈,𝒏)`$ in the following way
$$w_{\beta ,J}(𝒍,𝝈,𝒏)=\underset{<i,j>}{\mathbf{}}e^{\beta J_{i,j}l_il_j}\underset{<i,j>}{\mathbf{}}_{n_{ij}=1}p_{i,j}\delta _{\sigma _i\sigma _j}\underset{<i,j>}{\mathbf{}}_{n_{ij}=0}(1p_{ij}).$$
Now, we write $`\stackrel{~}{𝒏}`$ for a bond configuration fulfilling the compatibility condition that a bond $`\stackrel{~}{𝐧}_{i,j}`$ exists between sites $`i`$ and $`j`$ ($`\stackrel{~}{𝐧}_{i,j}=1`$) iff $`\sigma _i=\sigma _j`$ and call $`c(\stackrel{~}{𝒏})`$ the number of clusters of bonds in the configuration $`\stackrel{~}{𝒏}`$. Then, the weights $`w_{\beta ,J}`$ take the form (FK-representation)
$$w_{\beta ,J}(𝒍,\stackrel{~}{𝒏})=\underset{<i,j>}{\mathbf{}}e^{\beta J_{i,j}l_il_j}\left(2^{c(\stackrel{~}{𝒏})}\underset{<i,j>\stackrel{~}{𝒏}}{\mathbf{}}p_{i,j}\underset{<i,j>\stackrel{~}{𝒏}}{\mathbf{}}(1p_{ij})\right).$$
(2.5)
These weights are called the Wolff weights of the random cluster distribution. The joint distribution on the configurations ($`𝒍,\stackrel{~}{𝒏}`$) provide the right framework to relate the critical behavior in the original spin system and its associated percolation representation.
###### Remark 2.1
The Wolff weights $`w_{\beta ,J}(𝐥,\stackrel{~}{𝐧})`$ differ from the (FK) weights of the Ising case (see \[Fortuin, CK\]) by the term $`{\displaystyle \underset{<i,j>}{\mathbf{}}}e^{\beta J_{i,j}l_il_j}`$ that reflect the measure on the spin amplitudes in the continuous case.
## 3 Wolff Cluster Algorithm and Distributions
Non-local cluster Monte Carlo (MC ) algorithms have brought significant improvements in the simulation of Ising models near criticality. Starting from the ground-breaking work of Fortuin and Kasteleyn \[Fortuin\], which relates the partition function of spin systems with that of a correlated percolation model, Swendsen and Wang \[SW\] derived a non-local cluster algorithm which drastically reduces the critical slowing down phenomena near the transition point, with a dynamical critical exponent near $`0.25`$, where $`z`$ is defined by
$$\tau =\xi ^z.$$
(3.1)
Here $`\tau `$ is the correlation time in MC simulations (measured in MC steps per site) and $`\xi `$ is the correlation length. For a system of size $`L`$, near the critical point, $`\xi `$ scales like $`L^2`$ in two dimensions. In the case of a local update algorithm like Metropolis or Heatbath, $`z`$ is found to be close to $`2`$, so the required time to reach stable configurations at criticality in this case is of order $`L^4`$, reducing noticeably the possible size of samples to study.
Building on the Swendsen-Wang idea, Wolff \[Wo\] improved the method (see below) and derived a non-local update MC scheme with a dynamical critical exponent $`z`$ smaller than or equal to Swendsen-Wang’s in any dimension and furthermore easier to implement.
### 3.1 The Wolff Algorithm
First let us recall briefly the main features of the Wolff method. Consider the $`2`$-$`D`$ ferromagnetic Ising Model on $`Z^2`$ with coupling constant $`J`$ and inverse temperature $`\beta `$. Starting from a randomly chosen spin $`\sigma _0`$, visit all nearest neighbours and
1. with probability $`p_b=1exp(2\beta J)`$,
2. and only if they have the same orientation as $`\sigma _0`$, (III.1)
include them in the same cluster as $`\sigma _0`$; spins not satisfying both conditions are excluded. Repeat iteratively this procedure with newly added spins in the cluster until no more neighbours fulfill the above compatibility condition ($`III.1`$). Now flip all the spins in that cluster with probability $`1`$. After that, erase all the bonds and start this procedure again. It turns out that this dynamics verifies the detailed balance condition, i.e. it samples the Gibbs distribution of the Ising model (see \[Wo\]). The distinguishing feature of the Wolff method compared to Swendsen-Wang’s is that, in the latter (following Fortuin and Kasteleyn \[Fortuin\]) one needs to build, with the same growing probability $`p_b`$ as before, all possible clusters of like spins and then, with probability 1/2, flip all the spins in those clusters. Then all the bonds are erased and one starts again from the newly created spin configuration. It can also be shown that this method verifies the detailed balance condition \[SW\]. In summary, the following remarks can be made:
1. The building procedure of an individual cluster in both methods is clearly the same, thus the Wolff cluster belongs to the set of Swendsen-Wang clusters.
2. When the Wolff cluster is built, the randomly chosen spin has obviously higher probability to fall in a large Swendsen-Wang cluster than in a smaller one.
It results that the distribution of Wolff clusters is given by the distribution of Swendsen-Wang clusters (see below) modified by an additional weight that takes into account the size of the clusters.
|
warning/0003/astro-ph0003291.html
|
ar5iv
|
text
|
# Detection of H𝛼 Emission in a Methane (T-type) Brown Dwarf
## 1 Introduction
Activity is an important parameter in the study of stellar populations. Numerous investigations of late-type (F-M) stars have shown correlations between emission (e.g., Ca II H & K, Mg II h & k, Balmer series, etc.) and fundamental parameters such as age, rotation, and metallicity (Hawley, Reid, & Gizis, 2000). It is generally believed that the majority of this optical emission occurs in the chromosphere via collisional heating by ions and electrons along magnetic field lines. Indeed, this hypothesis is supported by the observed correlation of activity and rotation in late-type stars (Kraft, 1967; Noyes et al., 1984; Basri, 1987), which is expected if magnetic fields are generated by an internal dynamo (e.g. $`\alpha `$-$`\mathrm{\Omega }`$ dynamo; Parker 1955). Decrease in activity as stars age can be attributed to spin-down due to angular momentum loss in stellar winds (Stauffer & Hartmann, 1986).
As we examine cooler M and L dwarfs, however, these relations begin to break down. As stars become fully convective ($``$ 0.3 M), the $`\alpha `$-$`\mathrm{\Omega }`$ dynamo mechanism becomes ineffective, as it requires a low-buoyancy, radiative/convective boundary to anchor flux lines (Spiegel & Weiss, 1980). However, the observed activity level remains roughly constant around this transition point (Hawley, Gizis, & Reid, 1996), suggesting a turbulent dynamo as an alternate magnetic field source (Durney, De Young, & Roxburgh, 1993). Indeed, flaring activity, which is magnetically driven, is seen in objects as late as M9.5 (Liebert et al., 1999; Reid et al., 1999), supporting the existence of substantial magnetic fields beyond the convective cut-off. Alternately, acoustic heating could sufficiently heat the chromosphere (Schrijver, 1987; Mathioudakis & Doyle, 1992) to produce a “basal flux” of H$`\alpha `$ emission. In either case, Gizis et al. (2000) have shown that the fraction of objects with measurable emission rises to 100% at spectral type M7, then rapidly declines, so that no emission is seen in types L5 or later (Kirkpatrick et al., 2000). The decrease in (steady) activity even encompasses objects with rapid rotation (Basri & Marcy, 1995; Tinney & Reid, 1998), at odds with trends in hotter stars. Whether this drop in emission is due to ineffective chromospheric heating, decreased magnetic activity, or some other mechanism is unclear, but the end of the main sequence appears to mark a change in activity.
Based on the results of Gizis et al. (2000), we would not expect significant activity in T dwarfs, brown dwarfs that show CH<sub>4</sub> absorption bands at 1.6 and 2.2 $`\mathrm{\mu m}`$ (Kirkpatrick et al., 1999). Nonetheless, we have observed H$`\alpha `$ in emission in the T dwarf 2MASSW J1237392+652615 (Burgasser et al., 1999, herefter 2MASS J1237+65), identified from the Two Micron All Sky Survey (Skrutskie et al., 1997). We describe the optical observations of this and two other T dwarfs in $`\mathrm{\S }`$2; in $`\mathrm{\S }`$3 we discuss the H$`\alpha `$ detection in 2MASS J1237+65 and possible emission mechanisms; in $`\mathrm{\S }`$4 we compare 2MASS J1237+65 to the unusual M9.5Ve PC 0025+0447; we summarize our results in $`\mathrm{\S }`$5.
## 2 Optical Spectroscopy
### 2.1 Observations
2MASS J1237+65 was observed on three consecutive nights, 1999 July 16, 17, and 18 (UT) using the Low Resolution Imaging Spectrograph (Oke et al., 1995, hereafter LRIS) on the Keck II 10-meter telescope. On each occasion, conditions were transparent with seeing of 0$`\stackrel{}{\mathrm{.}}`$8 to 1$`\stackrel{}{\mathrm{.}}`$0, and we employed a 1$`\mathrm{}`$ wide slit. The target was acquired via blind offset from field stars, as it was invisible in the guiding imager. The July 16 and 17 observations were made using the 400 lines mm<sup>-1</sup> grating blazed at 8500 Å, covering the wavelength range 6300 to 10100 Å at 9 Å resolution. Total integrations of 4800 (1800 + 1800 + 1200) and 1800 s were obtained, respectively. The July 17 observation was plagued by poor target centering, reducing the observed flux by a factor of $``$ 2; the data for this night are thus omitted, although H$`\alpha `$ was detected. The July 18 observation was an 1800 s exposure using the 300 lines mm<sup>-1</sup> grating blazed at 5000 Å, covering 3800 to 8600 Å at 12 Å resolution. Additional 3600 and 2700 s observations (total integration time) of T dwarfs SDSSp J162414.37+002915.6 (Strauss et al., 1999, hereafter SDSS 1624+00) and SDSSp J134646.45-003150.4 (Tsvetanov et al., 2000, hereafter SDSS 1346-00) were made on July 16 and 17, using the 400 lines mm<sup>-1</sup> grating blazed at 8500 Å.
The data for all three objects were reduced and calibrated using standard IRAF<sup>1</sup><sup>1</sup>1IRAF is distributed by the National Optical Astronomy Observatories, which are operated by the Association of Universities for Research in Astronomy, Inc., under cooperative agreement with the National Science Foundation. routines. A one-second dark exposure was used to remove the bias, and quartz-lamp flat-field exposures were used to normalize the response of the detector. The individual spectra were extracted using the APEXTRACT routine, allowing for slight curvature of the point-source dispersion line viewed through the LRIS optics. Due to the low flux levels of the T dwarf targets, we used an extraction template derived from observations of standard stars. Wavelength calibration was achieved using Hg-Ne-Ar arc lamp exposures taken after each object exposure. Finally, the spectra were flux calibrated using observations of the DC white dwarf standard LTT 9491 (Hamuy et al., 1994). Data have not been corrected for telluric absorption, so the atmospheric H<sub>2</sub>O bands at 8161-8282, 8950-9300, and 9300-9650 Å are still present in the spectra (telluric H<sub>2</sub>O and O<sub>2</sub> bands are not visible shortward of 8000 Å due to low flux levels).
### 2.2 Optical Spectra
Spectra from 6300-10100 Å for all three objects are shown in Figure 1. Data blueward of 6300 Å are not shown, as our July 18 observations failed to detect any flux, continuum or emission, in this regime. Features identified by Oppenheimer et al. (1998) in Gl 229B are indicated. The optical spectrum of SDSS 1624+00 is discussed in detail in Liebert et al. (2000). Cs I lines at 8521 and 8944 Å are seen strongly in SDSS 1624+00, while both SDSS 1346-00 and 2MASS J1237+65 have progressively weaker lines. Broadened H<sub>2</sub>O at 9250 Å and a weak CH<sub>4</sub> feature centered near 8950 Å are noted, both of which are stronger in 2MASS J1237+65 and SDSS 1346-00 than in SDSS 1624+00. The pressure-broadened resonant K I doublet at 7665 and 7699 Å, which is prominent in the spectra of L dwarfs (Kirkpatrick et al., 1999), is masking most of the flux between 7300 and 8100 Å in all three objects (Burrows, Marley, & Sharp, 2000; Liebert et al., 2000). We note a feature in SDSS 1624+00 centered around 9900 Å which we identify as the 9896 Å 0-0 A<sup>4</sup>$`\mathrm{\Delta }`$-X<sup>4</sup>$`\mathrm{\Delta }`$ band of FeH, a feature that weakens from L6V and later (Kirkpatrick et al., 1999). This detection supports the claims of Nakajima et al. (2000) and Liebert et al. (2000) that SDSS 1624+00 is probably warmer than Gl 229B, which shows no FeH band. Hydride bands in SDSS 1346-00 and 2MASS J1237+65 are either weak or absent, suggesting that they in turn are cooler than SDSS 1624+00. The remaining feature of interest, H$`\alpha `$ emission at 6563 Å in 2MASS J1237+65, is the subject of the remainder of this article.
## 3 H$`\alpha `$ in 2MASS J1237+65
The H$`\alpha `$ emission seen in 2MASS J1237+65 (Figure 2) is quite unexpected given the cool nature of this object, and great care was taken to verify its presence over three nights of observation. The emission spike was clearly seen in the raw data and is not spatially extended, thus ruling out the possibility of a background HI region or telluric emission. H$`\alpha `$ emission was detected in five observations, spanning three nights and two different instrument settings, with the emission full-width at half-maximum equivalent to the instrumental resolution ($``$ 5 pixels). Hence, we can rule out the possibility of a chance cosmic ray. Finally, since no emission at 6563 Å is seen for the other T dwarfs, which were reduced using the same bias and flat-field exposures, we can rule out detector features. The H$`\alpha `$ emission line is thus quite real.
Measurements of the H$`\alpha `$ line for all three T dwarfs are given in Table 1. The integrated line luminosity at 6563 Å averaged over the three July 16 observations is $`f_{H\alpha }`$ = (6.6$`\pm `$0.6) $`\times `$ 10<sup>-17</sup> erg cm<sup>-2</sup> s<sup>-1</sup>. An equivalent width measurement is not possible as no continuum flux was detected in this spectral region. Note that there appears to be an increase in H$`\alpha `$ flux on July 18 to $`f_{H\alpha }`$ = (8.5$`\pm `$0.3) $`\times `$ 10<sup>-17</sup> erg cm<sup>-2</sup> s<sup>-1</sup>. We estimate the bolometric flux of 2MASS J1237+65 using the 2MASS J-band magnitude of 15.90$`\pm `$0.06 and assuming a Gl 229B bolometric correction of BC<sub>J</sub> = 2.3 (Matthews et al., 1996). This yields m<sub>bol</sub> = 18.2 and f<sub>bol</sub> $``$ 1.3 $`\times `$ 10<sup>-12</sup> erg cm<sup>-2</sup> s<sup>-1</sup>; thus, log($`f_{H\alpha }`$/$`f_{bol}`$) = log(L/L<sub>bol</sub>) $``$ $``$4.3. A limit on H$`\beta `$ emission at 4861 Å is estimated at $`f_{H\beta }`$ $``$ 2.0 $`\times `$ 10<sup>-17</sup> erg cm<sup>-2</sup> s<sup>-1</sup>, resulting in a Balmer emission ratio on 1999 July 18 (UT) of $`f_{H\beta }`$/$`f_{H\alpha }`$ $``$ 0.24.
### 3.1 Unusually Strong Magnetic Activity
The detection of H$`\alpha `$ is clearly inconsistent with the emission trend seen in objects later than L5V and indeed in brown dwarfs in general (Gizis et al., 2000). The level of activity for 2MASS J1237+65 can be placed in context with other late-type emission-line objects by comparing relative flux emitted through the H$`\alpha `$ line. Hawley, Gizis, & Reid (1996) find a mean log(L/L<sub>bol</sub>) $``$ $``$3.8 for field dMe stars with M<sub>bol</sub> $`<`$ 12 (spectral type M5V or earlier), while Gizis et al. (2000) show that this falls to below $``$5 for L dwarfs. Hence, even the activity level of 2MASS J1237+65 is inconsistent with these trends, while SDSS 1346-00 and SDSS 1624+00 have upper limits of log(L/L<sub>bol</sub>) $``$ $``$5.3 and $``$5.7, respectively. Does 2MASS J1237+65 have an unusually strong magnetic field, perhaps powered by a dynamo mechanism different than that of M and L dwarfs? Further examination of activity statistics in this cool regime are clearly needed.
### 3.2 Flaring
While sustained activity may appear to diminish toward later spectral types, flaring does not. A good example of this is BRI 0021-0214, a rapid M9.5V rotator (Basri & Marcy, 1995) which in quiescence shows little or no H$`\alpha `$ emission (Basri & Marcy, 1995; Tinney, Delfosse, & Forveille, 1997; Tinney & Reid, 1998), yet was seen to flare by Reid et al. (1999) with log(L/L<sub>bol</sub>) = $``$4.2, similar to 2MASS J1237+65. The BRI 0021-0214 flare also showed no continuum emission, again similar to the activity in 2MASS J1237+65. While flaring has not yet been observed in L dwarfs, the occurrence of flares in the latest M dwarfs (Reid et al., 1999; Liebert et al., 1999; Fleming, Giampapa, & Schmitt, 2000) suggests that they may simply have not yet been observed. However, most strong M-dwarf flares persist only a few hours (Hawley & Petterson, 1991), while smaller flares (which would show weaker continuum emission) have timescales on the order of minutes (Nelson et al., 1986). Thus, prolonged emission of 2MASS J1237+65 requires either a continuous flaring mechanism or fortuitous timing on our part. Further observations of the emission line are warranted to investigate more fully the temporal behavior of this emission.
### 3.3 Acoustic Flux
Schrijver (1987) first pointed out that a minimum “basal flux” seen in late-type stars could be attributed to chromospheric heating from acoustic waves. We can estimate the amount of acoustic flux from 2MASS J1237+65 by extrapolating the models of Ulmschneider et al. (1996), which improved on pivotal work done by Bohn (1984). Estimating T<sub>eff</sub> $``$ 1000 K, log g(cm s<sup>-2</sup>) $``$ 5, and using the scaling F<sub>acoustic</sub> $``$ T<sub>eff</sub><sup>20</sup> from the coolest points in their Table 1, we estimate F<sub>acoustic</sub> $``$ 4 $`\times `$ 10<sup>-8</sup> erg cm<sup>-2</sup> s<sup>-1</sup>, or log(L<sub>acoustic</sub>/L<sub>bol</sub>) $``$ $``$15, well below the observed activity level. As such, acoustic energy is not likely the source of activity in 2MASS J1237+65.
### 3.4 An Interacting Binary
Finally, it is possible that 2MASS J1237+65 is active as a result of a binary interaction with an equal magnitude or fainter companion. The binarity of brown dwarfs is well established, (Martín et al., 1998; Basri & Martín, 1999; Koerner et al., 1999), and the identification of three equal magnitude binaries in an L dwarf sample of ten by Koerner et al. (1999) is consistent with the binary fraction seen in late main sequence stars ($``$ 35%; Fischer & Marcy 1992; Henry & McCarthy 1993). Hence, it is reasonable to consider that 2MASS J1237+65 could itself be double. In this case, the companion would have to be equal-mass or smaller, as the optical continuum shows no evidence of a warmer component nor any significant variation over three days.
An unusual feature of brown dwarf binaries is the possibility of sustained Roche lobe overflow. Because of its degenerate interior, if a brown dwarf loses mass on a dynamical timescale ($`\tau `$ $``$ ($`\frac{R^3}{GM})^{\frac{1}{2}}`$ $``$ 1 hr), its radius increases approximately as dlnR/dlnM $``$ -$`\frac{1}{3}`$ when its mass lies in the range 5 - 70 M<sub>Jup</sub><sup>2</sup><sup>2</sup>2M<sub>Jup</sub> = 1.9 $`\times `$ 10<sup>30</sup> g = 0.00095 M, R<sub>Jup</sub> = 7.1 $`\times `$ 10<sup>9</sup> cm = 0.10 R (Allen, 1973). (Burrows & Liebert, 1993). Sustained conservative mass transfer then occurs if the change in the Roche lobe of the secondary satisfies
$$\frac{dlnR_L}{dlnM_2}=\frac{8}{3}q\frac{4}{3}(q+1)\frac{0.4q^{\frac{2}{3}}+\frac{1}{3}q^{\frac{1}{3}}(1+q^{\frac{1}{3}})^1}{0.6q^{\frac{2}{3}}+ln(1+q^{\frac{1}{3}})}<\frac{1}{3}.$$
(1)
We have taken $`q`$ = $`\frac{M_2}{M_1}`$ as the mass fraction, $`M_1`$ and $`M_2`$ the primary and secondary masses, and $`R_L`$ the secondary Roche lobe radius, as approximated by Eggleton (1983)
$$R_L=a\frac{0.49q^{\frac{2}{3}}}{0.6q^{\frac{2}{3}}+ln(1+q^{\frac{1}{3}})}.$$
(2)
Here $`a`$ is the binary separation. Equation (1) is satisfied when $`q`$ $`<`$ 0.63. We have calculated possible brown dwarf binary separations, periods, and edge-on primary radial velocities as a function of $`q`$ (Figure 3) for primary masses $`M_1`$ = 30 and 70 M<sub>Jup</sub> and (Zapolsky & Salpeter, 1969)
$$R_L=R_2=\frac{\pi M_{(Jup)}^{\frac{1}{3}}}{(1+1.8M_{(Jup)}^{\frac{1}{2}})^{\frac{4}{3}}}R_{(Jup)}.$$
(3)
M<sub>(Jup)</sub> and R<sub>(Jup)</sub> are the mass and radius of the secondary in units normalized to Jupiter. This rough model shows that sustained outflow requires a separation $`a`$ $``$ 4 - 20 R<sub>Jup</sub> (Figure 3a). This close separation may seem problematic; however, the brown dwarf spectroscopic binary PPl 15 has been shown to have an orbital separation of about 60 R<sub>Jup</sub> (Basri & Martín, 1999). Note that this scenario results in at least partial eclipsing for orbital inclination $`i`$ $``$ 6 - 30$`\mathrm{°}`$, suggesting that photometric monitoring could detect a transit event in a period of $``$ 1 - 10 hours. Alternatively, radial velocity monitoring of the H$`\alpha `$ line can produce a measurable variation (v<sub>1</sub> $``$ 10 - 90 cos($`i`$) km s<sup>-1</sup>) over this period.
The emission mechanism from overflow accretion is uncertain, as a hot accretion disk is ruled out by the lack of thermal continuum. Ionization along a shock front or magnetic streaming onto the primary’s pole are possibilities. Nonetheless, we find this an intriguing and, more importantly, observationally constrainable hypothesis.
## 4 PC 0025+0447: An Analogue?
A possible analogue to 2MASS J1237+65 is the much warmer M9.5Ve PC 0025+0447, which was identified by Schneider et al. (1991) in a search for high redshift quasars due to its unusually strong Balmer line emission. The contribution of H$`\alpha `$ luminosity to total luminosity (log(L/L<sub>bol</sub>) = $``$3.4; Schneider et al. 1991) is a full order of magnitude stronger in this object than 2MASS J1237+65, and appears to be steady over a timescale of at least 7 - 8 years with no indication of flaring activity (Schneider et al., 1991; Mould et al., 1994; Martín, Basri, & Zapatero Osorio, 1999). Both Schneider et al. (1991) and Martín, Basri, & Zapatero Osorio (1999) argue that PC 0025+0447 is a young (600 Myr) brown dwarf; if this were the case, then it is possible that 2MASS J1237+65 is an evolved version of this object that has slowly declined in activity with age. Martín, Basri, & Zapatero Osorio (1999) propose various mechanisms for the activity in PC 0025+0447, including sustained flaring and emission from a highly active, low-density corona. Mould et al. (1994), however, argue against PC 0025+0447 being a brown dwarf based on significant Li depletion, and propose that it is a Hyades-age ($``$ 700 Myr) main sequence star with normal photospheric thermal emission.
We propose a scenario in which PC 0025+0447 is itself an interacting binary, consisting of a $``$ 0.1 M M dwarf and a brown dwarf companion losing mass to the primary. This scenario provides a natural explanation for the observed phenomena: Roche lobe overflow from the companion provides a steady H$`\alpha `$ emission source, accretion may lead to the observed variable veiling of the M dwarf spectrum (Martín, Basri, & Zapatero Osorio, 1999), and the accreted material from the brown dwarf companion may retain lithium at primordial abundance, leading to the intermittent appearance of Li 6707 Å absorption seen by Martín, Basri, & Zapatero Osorio (1999), without requiring modification of the measured trigonometric parallax. Extrapolation from Figure 3 leads to similar maximal separations (3 - 25 R<sub>Jup</sub>) and periods (1 - 14 hours) as for 2MASS J1237+65. Unfortunately, despite many photometric and spectroscopic observations of PC 0025+0447, no time-resolved data is available to constrain this hypothesis. More detailed modeling is required to determine the feasibility of these mechanisms.
## 5 Summary
We have reported the detection of H$`\alpha `$ emission in the T dwarf 2MASS J1237+65, at a level of log(L/L<sub>bol</sub>) = $``$4.3. This emission is intriguing, as it is a salient exception to the cool dwarf temperature-activity relations identified thus far. We have proposed various activity mechanisms, including a strong magnetic field, continuous flaring, and an interacting brown dwarf binary, but these are speculative guesses at best. Comparison can be drawn with the M9.5Ve PC 0025+0447, which, if it is a brown dwarf, could be a warm analogue to 2MASS J1237+65. Both objects could also be close binary systems with lower-mass brown dwarf companions, which are steadily transferring mass to their primaries by Roche lobe overflow. Nonetheless, the mechanism for both objects remains unclear. Further investigation of the temporal stability of the H$`\alpha `$ line in 2MASS J1237+65 and searches for emission in other T dwarfs is clearly warranted.
###### Acknowledgements.
A. J. B. acknowledges the assistance of Terry Stickel and Wayne Wack at the telescope, useful discussions with P. Goldreich and M. Marley, helpful comments from our referee S. Hawley, and the foresight of S. Phinney for assigning the brown dwarf binary problem in class (equation 1). A. J. B., J. D. K., I. N. R., J. E. G., and J. L. acknowledge funding through a NASA/JPL grant to 2MASS Core Project science. A. J. B., J. D. K., and J. E. G. acknowledge the support of the Jet Propulsion Laboratory, California Institute of Technology, which is operated under contract with the National Aeronautics and Space Administration. Data presented herein were obtained at the W. M. Keck Observatory which is operated as a scientific partnership among the California Institute of Technology, the University of California, and the National Aeronautics and Space Administration. The Observatory was made possible by generous financial support of the W. M. Keck Foundation.
|
warning/0003/nlin0003033.html
|
ar5iv
|
text
|
# 1 Introduction and the Statement of Main Results
## 1 Introduction and the Statement of Main Results
In the present work we consider the billiards on two-dimensional surfaces $`M`$ of constant Gaussian curvature $`r`$ in the presence of a homogeneous magnetic field of magnitude $`\beta `$, which is perpendicular to $`M`$. Inside the billiard domain $`Q`$ the pointlike particle of unit charge and mass moves at unit velocity along curves of constant geodesic curvature $`\beta `$ and reflects elastically at the boundary $`Q`$. In the following, we will call these dynamical systems as magnetic billiards and $`M`$ as magnetic surface if $`\beta 0`$.
Magnetic billiards on the plane have been considered in many works \[BR\], \[K\], \[BK\], \[Ta1\] and on the hyperbolic plane in \[Ta2\]. The study of such billiards is strongly motivated by mesoscopic physics, where such billiard models are used as simplified version of the mesoscopic devices in the presence of magnetic fields. In the present paper we treat the magnetic billiards simultaneously on all surfaces of constant curvature (sphere, plane and hyperbolic plane). For all values of $`r`$ and $`\beta `$ we establish a common criterion for hyperbolicity of the billiard dynamics, whose geometric realization depends only on the type of linearized dynamics (geometric optics) on $`M`$. This extends our recent results \[GSG\] for non-magnetic surfaces of constant curvature.
The dynamics on $`M`$ depend crucially both on the curvature of the surface and on the strength of magnetic field. Firstly, let us consider the case $`\beta =0`$. On the plane the neighboring trajectories separate only linearly with time, so that the motion of the point mass between collisions with the boundary is neutral. Exponential separation of billiard trajectories can only occur if the successive reflections from the boundaries introduce sufficient instability. On the hyperbolic plane the negative curvature induces exponential divergence of geodesics. Thus, the boundary of hyperbolic billiard can be neutral (i. e., with zero curvature), and the chaotic dynamics will be provided by the metric. On the sphere, in contrast, any two geodesics intersect twice at focal points. Thus, the boundary reflections have to compensate for the focusing effect of the sphere, in order to produce chaotic dynamics. We will call the mentioned above types of linearized dynamics arising on the plane, hyperbolic plane and sphere as parabolic , hyperbolic and elliptic respectively.
In the presence of a constant magnetic field, an additional focusing effect appears. To simplify matters, let us discuss first the planar case when $`\beta 0`$. Consider an infinitesimal beam of trajectories which emerges from the same point at the time $`t=0`$ (fig. 1a). Then, by elementary calculations (see e.g., \[K\]) the curvature of the infinitesimal beam at the time $`t=s`$ is given by
$$\chi (s)=\beta \mathrm{cot}(\beta s).$$
(1)
Consider now the reflection of infinitesimal beams at the boundary. Let $`m`$ be the bouncing point, let $`C`$ be the osculating circle at $`m`$ and let $`m^{}`$ be the second point in which the particle trajectory intersects $`C`$. Denote by $`\chi _{}`$, $`\chi _+`$ the curvatures of the infinitesimal beam immediately before and after reflection. Then, the change of the curvature under the reflection \[Ta1\], \[K\] at the bouncing point $`m`$ (see fig. 1b) is given by
$$\chi _+\chi _{}=2\beta \mathrm{cot}(\beta d),$$
(2)
where $`2d=\overline{mm^{}}`$ is the signed length of particle trajectory between $`m`$ and $`m^{}`$ (when the curvature of the boundary at $`m`$ is positive, $`2d`$ is simply the time which the particle spends after (before) reflection in the osculating circle $`C`$). It is a simple observation, that eqs. (1), (2) are actually the same ones, which have appeared in \[GSG\] for the non-magnetic billiards on the sphere of radius $`\beta ^1`$, where the parameters $`s`$ and $`d`$ have the same geometric meaning. Thus, the geometric optics (or linearized dynamics) on the plane in the presence of a magnetic field $`\beta `$ is equivalent to the geometric optics on the non-magnetic sphere of radius $`\beta ^1`$. More generally, we demonstrate in the body of the paper the equivalence of geometric optics on the surfaces of constant Gaussian curvature $`r`$ in the presence of magnetic field $`\beta `$ to the geometric optics on the non-magnetic surfaces of constant Gaussian curvature $`r_{eff}=r+\beta ^2`$. As a consequence, the type of linearized dynamics (elliptic, parabolic or hyperbolic) depends only on the sign of $`r_{eff}`$. We will call the parameter $`r_{eff}`$ as effective curvature of the surface and will refer to the cases $`r_{eff}=0`$, $`r_{eff}<0`$, $`r_{eff}>0`$ as Parabolic (P), Hyperbolic (H) and Elliptic (E) respectively.
Up to now, the study of hyperbolic billiards has been mainly restricted to the Euclidean plane (see \[Tab\] for review). See, however, \[Ta2\] for some results on chaotic billiards on the hyperbolic plane, and \[Vet1\], \[Vet2\], \[KSS\] for some results on hyperbolic billiards on a general Riemannian surface. In our recent work \[GSG\] we have generalized Wojtkowski’s criterion of hyperbolicity for planar billiards \[Wo2\] to billiards on arbitrary surfaces of constant curvature. On the basis of the equivalence of geometric optics on magnetic and non-magnetic surfaces of constant curvature, we extend in the present paper the criterion of \[GSG\] to the case of magnetic billiards.
The hyperbolicity criterion in \[GSG\] can be formulated in terms of a special class of trajectories, which generalize two-periodic orbits. Let $`Q`$ be a billiard table on a surface of constant curvature. The billiard map $`\varphi :VV`$ acts on the phase space $`V`$, which consists of pairs $`v=(m,\theta )`$. Here $`m`$ is the position of the ball on the boundary $`Q`$ of $`Q`$, and $`\theta `$ is the angle between the outgoing velocity and the tangent to $`Q`$ at $`m`$. The billiard map preserves a natural probability measure $`\mu `$ on $`V`$. We denote the images of $`v`$ after $`n`$ iterations by $`(m_{n+1},\theta _{n+1})=\varphi ^n(v)`$. The trajectory $`\varphi ^n(v)`$ is a generalized two-periodic trajectory (g.t.p.t.) if the following conditions are satisfied:
1. The incidence angle and the curvature of the boundary $`\kappa _n`$ at the bouncing points have period 2: $`\theta _{2n}=\theta _2`$, $`\theta _{2n+1}=\theta _1`$, $`\kappa _{2n}=\kappa _2`$, $`\kappa _{2n+1}=\kappa _1`$;
2. The length of trajectory between consecutive bouncing points is constant: $`s=\overline{m_nm}_{n+1}`$ (see fig. 2).
If $`\theta _i=\pi /2`$, the g.t.p.t. is an usual two-periodic orbit.
Along a g.t.p.t. the linearized map $`D_v\varphi `$ is two-periodic, and the stability of a g.t.p.t. is determined by $`D_v\varphi ^2`$. For each surface of constant curvature, the stability type of a g.t.p.t. is completely determined by the triplet of parameters $`(d_1,d_2,s)`$, where $`2d_1`$ (resp. $`2d_2`$) is the signed length of the chord generated by the intersection of the trajectory $`m_1m_2`$ with the osculating circle at $`m_1`$ (resp. $`m_2`$). We shall use the symbol $`T(d_1,d_2,s)`$ for the g.t.p.t. with parameters $`(d_1,d_2,s)`$.
Let us consider g.t.p.t.s for planar, non-magnetic billiards in some detail. Here $`s`$ is the euclidean distance between consecutive bouncing points, and $`d_i=r_i\mathrm{sin}\theta _i`$, $`i=1,2`$, where $`r_i`$ are the radii of curvature of the boundary $`Q`$ at the respective points. If the curvature of the boundary at the bouncing point is zero we take $`r_i=\mathrm{}`$ as the radius of curvature and $`d_i=\mathrm{}`$ respectively. By an elementary computation, $`T(d_1,d_2,s)`$ is unstable if and only if
$$s\{\begin{array}{cc}[d_1,d_2][d_1+d_2,\mathrm{})\hfill & \text{if }d_1,d_20\hfill \\ [0,\mathrm{})\hfill & \text{if }d_1,d_20\hfill \\ [0,d_1+d_2][d_1,\mathrm{})\hfill & \text{if }d_10,d_20\text{.}\hfill \end{array}$$
(3)
Moreover, the trajectory is hyperbolic (i. e., strictly unstable) if $`s`$ is in the interior of the corresponding interval, and the trajectory is parabolic if $`s`$ is a boundary point (in the limiting case $`d_1=d_2=\mathrm{}`$ the trajectory is parabolic for any value of $`s`$).
We consider two classes of unstable g.t.p.t.s. The g.t.p.t. $`T(d_1,d_2,s)`$ is B-unstable if in eq. (3) $`s`$ belongs to a “big interval”:
$$s\{\begin{array}{cc}[d_1+d_2,\mathrm{})\hfill & \text{if }d_1,d_20\hfill \\ [0,\mathrm{})\hfill & \text{if }d_1,d_20\hfill \\ [d_1,\mathrm{})\hfill & \text{if }d_10,d_20\text{.}\hfill \end{array}$$
(4)
On the contrary, if $`s`$ belongs to a “small interval”, then $`T(d_1,d_2,s)`$ is S-unstable:
$$s\{\begin{array}{cc}[d_1,d_2]\hfill & \text{if }d_1,d_20\hfill \\ [0,d_1+d_2]\hfill & \text{if }d_10,d_20\text{.}\hfill \end{array}$$
(5)
Note that a small interval shrinks to a point when $`|d_1|=|d_2|`$.
It has been demonstrated in \[GSG\], that the notions of B-unstable and S-unstable g.t.p.t.s are generalized to arbitrary surfaces of constant curvature, where the analogs of (3,4,5) exist. The concept of g.t.p.t.s and the associated structures make sense for billiard on any surface immersed in a magnetic field. Since the stability properties of the trajectories depend only on the linearized dynamics (geometrical optics) of the system, one has essentially the same stability intervals (in terms of the parameters $`(s,d_1,d_2)`$) for g.t.p.t.s on the surface of constant Gaussian curvature $`r`$ immersed in the magnetic field $`\beta `$ and for g.t.p.t.s on the non-magnetic surface of constant Gaussian curvature $`r_{eff}=r+\beta ^2`$. As a consequence, one can extend the notions of B-unstable and S-unstable g.t.p.t.s to magnetic surfaces of constant curvature as well.
Equipped with the mentioned above definitions, we are ready now to formulate the main results of the present paper. Let $`Q`$ be a billiard table on a magnetic surface of constant curvature, and let $`\lambda (v)0`$ be the Lyapunov exponent of the billiard. With any point $`v=(m_1,\theta _1)V`$ of the phase space we associate a formal g.t.p.t. $`T(v)`$. Let $`\varphi (v)=(m_2,\theta _2)`$. We set $`d_1=d(v)`$, $`d_2=d(\varphi (v))`$ and $`s=\overline{m_1m_2}`$ for the length of particle trajectory between $`m_1`$ and $`m_2`$. Then $`T(v)`$ is determined by the triple $`(d_1,d_2,s)`$. The formal g.t.p.t. $`T(v)`$ can be realized as an actual g.t.p.t. $`T(d_1,d_2,s)`$ in an auxiliary billiard table $`Q_v`$, constructed from the boundary $`Q`$ around $`m_i`$ (see \[GSG\]). Let $`\varphi _v`$ be the map in $`Q_v`$ corresponding to the g.t.p.t. $`T(v)`$, and let $`\overline{\lambda }(v)=lim_{n\pm \mathrm{}}\frac{1}{n}\mathrm{log}D\varphi _v^n0`$ be the Lyapunov exponent of $`T(v)`$. Then the sufficient condition for hyperbolic dynamics in $`Q`$ can be formulated as follows (see also Theorem 1 for the alternative formulation in the body of the paper).
Main Theorem. If for $`\mu `$ almost every point $`vV`$, $`T(v)`$ is B-unstable, and for $`\mu `$ almost every point $`vV`$, $`T(\varphi ^n(v))`$ is strictly B-unstable for some $`n`$, then the billiard in $`Q`$ is hyperbolic ($`\lambda (v)`$ is positive $`\mu `$ almost everywhere).
After deriving the conditions, which insure that g.t.p.t.s are B-unstable (analogs of (4)), this theorem turns to be a geometric criterion for hyperbolicity of the billiard dynamics. In particular, for planar non-magnetic billiards, the Main Theorem yields Wojtkowski’s criterion for hyperbolic dynamics \[Wo2\]. Let $`Q`$ be a billiard table satisfying the assumptions of the Main Theorem. Following the approach of Wojtkowski’s \[Wo2\], the metric entropy $`h(Q)=_V\lambda (v)𝑑\mu `$ of the billiard can be actually estimated from below (Theorem 2 in the body of the paper):
$$h(Q)_V\overline{\lambda }(v)𝑑\mu .$$
(6)
The paper is organized as follows. In Section 2 we provide the necessary preliminaries and establish the relationship between the geometric optics (i. e., the rules of propagation and reflection of infinitesimal light beams) on magnetic and non-magnetic surfaces of constant curvature. In Section 3 we apply these results to obtain explicit analogs of (3-5) for all magnetic (non-magnetic) surfaces of constant Gaussian curvature. We obtain a linear instability conditions for g.t.p.t.s and show that they distinguish between B-unstable and S-unstable trajectories in a natural way. In Section 4, we reformulate the Main Theorem in a slightly different way and prove it using the invariant cone fields method. We define our cone fields for magnetic billiards on all surfaces of constant curvature exactly as in \[GSG\]. Using the geometric optics language, the proof of the preservation of the cone field under the assumptions of the Main Theorem is reduced to the corresponding non-magnetic problem. The Lyapunov exponent estimation (6) follows from the results of \[GSG\] by the same arguments. In Section 5, for each type of linearized dynamics we derive the criterion of hyperbolicity for elementary billiard tables (the boundary of these billiards consists of arcs of constant geodesic curvature). We apply it to construct several classes of magnetic billiard tables with hyperbolic dynamics on surfaces of constant curvature. Finally, for each type of linearized dynamics we formulate some general principles for the design of the magnetic billiard tables satisfying the conditions of the Main Theorem. Several examples of billiards satisfying these principles are also given here. The derivation of the differential condition on the boundary of billiards satisfying the Main Theorem is given in the Appendix.
## 2 Geometric Optics
Let $`M`$ be a surface of constant Gaussian curvature $`r`$. We will distinguish three cases: $`M=𝐑^\mathrm{𝟐}`$ -plane ($`r=0`$), $`M=𝐒^\mathrm{𝟐}`$ -sphere ($`r>0`$) and $`M=𝐇^\mathrm{𝟐}`$ -hyperbolic plane ($`r<0`$). Let us consider the dynamical system arising on $`M`$ from the motion of particle of unit charge, mass and velocity in the presence of homogeneous magnetic field perpendicular to the surface. We will denote by $`M(r,\beta )`$ the corresponding surface $`M`$ in the presence of magnetic field of strength $`\beta `$. We can assume, without loss of generality, that $`\beta 0`$. The particle trajectories on $`M(r,\beta )`$ are curves of constant geodesic curvature $`\beta `$ (circles, paracycles or hypercycles). We call $`g^s`$ the corresponding flow on $`M(r,\beta )`$.
Let $`Q`$ be a connected domain in $`M(r,\beta )`$, with a piecewise smooth boundary $`Q`$. The billiard in $`Q`$ is the dynamical system arising from the motion of a point mass inside $`Q`$ under the action of magnetic field $`\beta `$, with specular reflections at the boundary. The phase space $`V`$ of the billiard consists of unit tangent vectors, with origin points on $`Q`$, pointing inside $`Q`$. The first return associated with $`V`$ is the billiard map, $`\varphi :VV`$. We denote the probabilistic invariant measure on $`V`$ as $`\mu `$ (for its realization in the planar case see e.g., \[BK\]). We will use the standard coordinates $`(l,\theta )`$ on $`V`$, where $`l`$ is the arclength parameter on $`Q`$ and $`\theta ,0\theta \pi `$, is the angle between the vector and $`Q`$.
Let $`D_v\varphi :T_vVT_{\varphi (v)}V`$ denote the derivative of $`\varphi `$. In what follows, we are interested in the action of $`D\varphi `$ on the projectivization $`B`$ of the tangent space $`TV`$, see \[Wo1\], \[Wo2\]. The space $`B=B_v`$, which is the set of straight lines in the tangent planes $`T_vV`$, $`vV`$, can be conveniently represented using the language of geometric optics. An oriented curve $`\gamma M`$, of class $`C^2`$, defines a “light beam”, i.e., the family of particle trajectories orthogonal to $`\gamma `$. The trajectories intersecting $`\gamma `$ infinitesimally close to a point, $`m\gamma `$, form an “infinitesimal beam”, which is completely determined by the normal unit vector $`\stackrel{}{n}T_mM`$ to $`\gamma `$, and by the geodesic curvature $`\chi `$ of $`\gamma `$ at $`m`$. Let $`\stackrel{}{n}=vV`$ (the point $`m`$ belongs to the boundary of $`Q`$). We denote by $`(v,\chi )`$ the infinitesimal beam determined by the pair $`(v,\chi )`$, $`vV`$, $`\chi 𝐑\mathrm{}`$.
On the other hand, each pair $`(v,\chi )`$ uniquely defines the line $`b(v,\chi )B`$, see e.g., \[Wo2\]. Here the vector $`v`$ defines the corresponding plane $`T_vV`$, and the curvature $`\chi `$ defines the direction of the line. Thus, infinitesimal beams yield a geometric representation of $`B=B_v`$, where $`\chi `$ can be used as a projective coordinate on the space $`B_v`$, see \[Wo2\], \[GSG\]. As a result, one can study the action of $`D\varphi `$ on $`B`$ in terms of the action of $`\varphi `$ on the curvature of infinitesimal beams $``$.
Let $`\rho _m:T_mMT_mM`$, $`mQ`$ be the linear reflection about the tangent line to $`Q`$. We will use the same letters, $`\varphi `$, $`\rho `$, and $`g`$, for the differentials of these mappings. Since the billiard map is the composition:
$$\varphi =\rho g,$$
(7)
it remains to compute the action of $`g`$ and $`\rho `$ on the curvature of infinitesimal beams. In other words, we need to know how the geodesic curvature of an infinitesimal beam changes: a) along free-flight trajectory on $`M(r,\beta )`$; b) under reflection.
a) Bouncless propagation. Let us consider the change of the geodesic curvature $`\chi (s)`$ of an infinitesimal beam along a particle trajectory on $`M(r,\beta )`$. As it has been shown in \[Ta2\], \[Ta3\] (see also \[K\] for the planar case) the free-flight evolution of $`\chi `$ satisfies the Ricatti equation:
$$\dot{\chi }=r_{eff}+\chi ^2,$$
(8)
where $`r_{eff}=r+\beta ^2`$ is the effective curvature of $`M(r,\beta )`$. For convenience we introduce also the parameter $`\xi =|r_{eff}|^{\frac{1}{2}}`$ (then $`\xi ^1`$ is the radius of the non-magnetic surface, which has the same linearized dynamics as $`M(r,\beta )`$). Let $`\chi =\chi (0)`$ be the geodesic curvature of an infinitesimal beam at the initial point, set $`\chi ^{}=\chi (s)=g^s\chi `$ be the geodesic curvature after the time $`s`$ of free-flight evolution. Using eq. (8) one can immediately obtain the relation between $`\chi `$ and $`\chi ^{}`$. We will distinguish three cases:
(E) -Elliptic $`(r_{eff}>0)`$. $`M`$ is either $`𝐒^\mathrm{𝟐}`$, or $`𝐑^\mathrm{𝟐}`$, or $`𝐇^\mathrm{𝟐}`$ with $`\beta ^2>|r|`$ (strong field).
$$\chi ^{}/\xi =\mathrm{cot}(\xi s)+\frac{\mathrm{sin}^2(\xi s)}{\mathrm{cot}(\xi s)\chi /\xi };$$
(9)
(P) -Parabolic $`(r_{eff}=0)`$. $`M`$ is $`𝐇^\mathrm{𝟐}`$ and $`\beta ^2=|r|`$, or $`M`$ is $`𝐑^\mathrm{𝟐}`$ and $`\beta =0`$.
$$\chi ^{}=s^1+\frac{s^2}{s^1\chi };$$
(10)
(H) -Hyperbolic $`(r_{eff}<0)`$. $`M`$ is $`𝐇^\mathrm{𝟐}`$ and $`\beta ^2<|r|`$ (weak field).
$$\chi ^{}/\xi =\mathrm{coth}(\xi s)+\frac{\mathrm{sinh}^2(\xi s)}{\mathrm{coth}(\xi s)\chi /\xi }.$$
(11)
Note, for $`M=𝐇^\mathrm{𝟐}`$ in the case (E) the particle trajectories are circles (trajectories of finite length), in the case (P) they are paracycles (trajectories of infinite length, which touch the boundary of Poincar$`\stackrel{´}{\mathrm{e}}`$ disc) and in the case (H) they are hypercycles (trajectories of infinite length, which have two points on the boundary of Poincar$`\stackrel{´}{\mathrm{e}}`$ disc).
b) Reflection. Let $`\chi _{}`$, $`v_{}T_mM`$ be the geodesic curvature and the direction of the infinitesimal beam just before the reflection at the point $`mM`$. We denote $`\chi _+=\rho \chi _{}`$, $`vv_+=\rho v_{}V`$ to be the geodesic curvature and the direction of the infinitesimal beam immediately after reflection. Let $`\kappa `$ be the curvature of $`Q`$ at $`m`$, and let $`\theta `$ be the angle between $`v`$ and the positive tangent vector to $`Q`$ at $`m`$. Then, the extension of the well known formula for planar billiards (see e.g., \[Si\]) to magnetic billiards on arbitrary surfaces (see \[Ta2\], \[Ta3\]) gives
$$\chi _+=\chi _{}+2K(v),\text{ where }K(v)=\frac{\kappa +\beta \mathrm{cos}\theta }{\mathrm{sin}\theta }.$$
(12)
Using classical formulas for surfaces of constant curvature (\[Vi\], see also \[Ta1\], \[Ta2\]), it is possible to give a geometric interpretation of the function $`K()`$.
Let $`vV`$, and let $`lQ`$ be the origin point of $`v`$. Let $`C(l)M`$ be the osculating circle (resp. paracycle or hypercycle if $`M=𝐇^\mathrm{𝟐}`$ and $`|\kappa (l)|\xi `$) of $`Q`$. The free-flight particle trajectory, $`G(v)`$, corresponding to $`v`$ intersects $`C(l)`$ at $`l`$ and at another point $`l^{}`$. In the cases (E) and (P) ($`r_{eff}0`$) let $`d(v)`$ be one half of the signed distance between $`l`$ and $`l^{}`$, along $`G(v)`$. To eliminate the ambiguity (when $`K(v)=0)`$, we choose the following intervals for $`d(v)`$: $`d(v)[\mathrm{},\mathrm{})`$ in the case (P) and $`\xi d(v)[\pi /2,\pi /2)`$ in the case (E). In the case (H) ($`r_{eff}<0`$) we will use the following classification of points of the phase space $`V`$, see fig. 3. We say that $`vV`$ is of type $`A`$ (resp. $`B`$) if $`|K(v)|\xi `$ (resp. $`|K(v)|<\xi `$). Let $`V^A,V^B`$ be the corresponding subsets of $`V`$. Then $`V=V^AV^B`$ is a partition. We denote by $`d^A(v)`$ ($`d^B(v)`$) one half of the signed distance between $`l`$ and $`l^{}`$ along $`G(v)`$ if $`vV^A`$ (resp. $`vV^B`$). To unite both cases we will use the notation:
$`d(v)=\{\begin{array}{cc}d^A(v)\hfill & \text{if }vV^A\hfill \\ d^B(v)+i\frac{\pi }{2}\hfill & \text{if }vV^B\text{.}\hfill \end{array}`$ (13)
Then we have
$`K^1(v)=\{\begin{array}{cc}d(v)\hfill & \text{in the case (P)}\hfill \\ \xi ^1\mathrm{tan}(\xi d(v))\hfill & \text{in the case (E)}\hfill \\ \xi ^1\mathrm{tanh}(\xi d(v))\hfill & \text{in the case (H).}\hfill \end{array}`$ (14)
As one can see from eqs. (9-14), the geometric optics (described in terms of parameters $`d`$, $`s`$) depend only on the value $`r_{eff}`$. This fact allows to study the linearized dynamics problems on $`M(r,\beta )`$ using the corresponding results for the non-magnetic surfaces of the constant Gaussian curvature $`r_{eff}=r+\beta ^2`$. The corresponding transition $`M(r,\beta )M(r_{eff},0)`$ is schematically illustrated by fig. 4.
## 3 Stability of Generalized Two-Periodic Trajectories
Let $`Q`$ be a billiard table on $`M(r,\beta )`$. For each $`vV`$ let $`t(v)`$ be the corresponding past semitrajectory in $`Q`$. Consider the curvature evolution of an infinitesimal beam along $`t(v)`$. Starting with $`(\varphi ^kv,\chi )`$ for arbitrary $`\chi `$, we obtain after $`k`$ steps forward the infinitesimal beam $`(v,\chi ^{(k)})`$, $`\chi ^{(k)}=\varphi ^k\chi `$. Eqs. (9-11) and (12) describe the action of the billiard map on the curvature of infinitesimal beams. Assuming $`k`$ to be infinity, we obtain a formal continued fraction corresponding to the semitrajectory $`t(v)`$:
$$c(v)=\chi ^{(\mathrm{})}=a_0+\frac{b_0}{a_1+{\displaystyle \frac{b_1}{a_2\mathrm{}}}}.$$
(15)
The coefficients of the continued fraction are determined by $`d_i=d(\varphi ^iv)`$, and by the lengths $`s_i`$ of consecutive billiard segments as follows:
$$\begin{array}{ccc}(P)\hfill & a_i=2s_i^1+2d_i^1,\hfill & b_i=s_i^2;\hfill \\ (E)\hfill & a_i=2\mathrm{cot}(\xi s_i)+2\mathrm{cot}(\xi d_i),\hfill & b_i=\mathrm{sin}^2(\xi s_i);\hfill \\ (H)\hfill & a_i=2\mathrm{coth}(\xi s_i)+2\mathrm{coth}(\xi d_i),\hfill & b_i=\mathrm{sinh}^2(\xi s_i).\hfill \end{array}$$
The continued fractions (15) determines the stability type of the trajectory: $`t(v)`$ is unstable if $`c(v)`$ is convergent (see e.g., \[Si\]). Since for a given sequence of $`d_i`$ and $`s_i`$, $`c(v)`$ is completely determined by $`r_{eff}`$, one can reduce the problem of stability trajectories on $`M(r,\beta )`$ to the corresponding “non-magnetic” problem on $`M(r_{eff},0)`$.
As it has been mentioned in the introduction, we are interested in the stability properties of generalized two periodic trajectories. A trajectory is a generalized two periodic trajectory (g.t.p.t.) if its parameters $`d_i`$ are periodic: $`d_{2i+1}=d_1`$, $`d_{2i}=d_2`$ and $`s_i=s`$ are the same along the trajectory (see fig. 2). Obviously, a g.t.p.t. yields a periodic continued fraction. We denote by $`T(d_1,d_2,s)`$ the g.t.p.t. with parameters $`(d_1,d_2,s)`$ and by $`c(d_1,d_2,s)`$ the associated continued fraction.
The stability of $`T(d_1,d_2,s)`$, or equivalently, the convergence of the two periodic continued fraction $`c(d_1,d_2,s)`$ has been studied in \[GSG\] for non-magnetic surfaces of constant curvature. On the basis of the equivalence between the magnetic and non-magnetic problems we can immediately generalize the results of \[GSG\] to the case $`\beta 0`$.
Proposition 1. The continued fraction $`c(d_1,d_2,s)`$ converges if and only if the following inequalities are satisfied.
$$\begin{array}{cc}(P)\hfill & (sd_1)(sd_2)(sd_1d_2)s0;\\ (E)\hfill & \mathrm{sin}(\xi (sd_1))\mathrm{sin}(\xi (sd_2))\mathrm{sin}(\xi (sd_1d_2))\mathrm{sin}(\xi s)0;\\ (H)\hfill & \mathrm{sinh}(\xi (sd_1))\mathrm{sinh}(\xi (sd_2))\mathrm{sinh}(\xi (sd_1d_2))\mathrm{sinh}(\xi s)0.\end{array}$$
Below we reformulate Proposition 1 explicitly as conditions for the instability of the corresponding g.t.p.t.
(P) $`T(d_1,d_2,s)`$ is unstable if and only if
$$s\{\begin{array}{cc}[d_1,d_2][d_1+d_2,\mathrm{})\hfill & \text{if }d_1,d_20\hfill \\ [0,\mathrm{})\hfill & \text{if }d_1,d_20\hfill \\ [0,d_1+d_2][d_1,\mathrm{})\hfill & \text{if }d_10,d_20.\hfill \end{array}$$
(16)
(E) In this case $`0\xi s2\pi `$, and we set $`\overline{\pi }=\pi \xi ^1`$,
$$s\mathrm{mod}\overline{\pi }=\{\begin{array}{cc}\mathrm{s}\hfill & \text{if }s\overline{\pi }\hfill \\ \mathrm{s}\overline{\pi }\hfill & \text{if }s>\overline{\pi }.\hfill \end{array}$$
Then $`T(d_1,d_2,s)`$ is unstable if and only if
$$s\mathrm{mod}\overline{\pi }\{\begin{array}{cc}[\mathrm{d}_1+\mathrm{d}_2,\overline{\pi }][\mathrm{d}_1,\mathrm{d}_2]\hfill & \text{if }d_1,d_20\hfill \\ [0,\mathrm{d}_1+\mathrm{d}_2+\overline{\pi }][\overline{\pi }\mathrm{d}_1,\overline{\pi }\mathrm{d}_2]\hfill & \text{if }d_1,d_20\hfill \\ [\mathrm{d}_2,\overline{\pi }+\mathrm{d}_1][0,\mathrm{d}_1+\mathrm{d}_2]\hfill & \text{if }d_10,d_20,|d_2||d_1|\hfill \\ [\mathrm{d}_2,\overline{\pi }+\mathrm{d}_1][\overline{\pi }+\mathrm{d}_2+\mathrm{d}_1,\overline{\pi }]\hfill & \text{if }d_10,d_20,|d_2||d_1|.\hfill \end{array}$$
(17)
(H) It matters whether $`v_iV^A`$ or $`v_iV^B`$ for $`i=1,2`$. We say that $`T(d_1,d_2,s)`$ is of type $`(AA)`$ if $`v_1V^A`$ and $`v_2V^A`$. The other types: $`(AB)`$, $`(BA)`$, and $`(BB)`$ are defined analogously. We formulate the explicit criteria of instability for $`T(d_1,d_2,s)`$ type-by-type.
Type $`(AA)`$:
$$s\{\begin{array}{cc}[d_1^A,d_2^A][d_1^A+d_2^A,\mathrm{})\hfill & \text{if }d_1^A,d_2^A0\hfill \\ [0,\mathrm{})\hfill & \text{if }d_1^A,d_2^A0\hfill \\ [0,d_1^A+d_2^A][d_1^A,\mathrm{})\hfill & \text{if }d_1^A0,d_2^A0\text{,}\hfill \end{array}$$
(18)
Type $`(BB)`$:
$$s\{\begin{array}{cc}[d_1^B+d_2^B,\mathrm{})\hfill & \text{if }d_1^B+d_2^B0\hfill \\ [0,\mathrm{})\hfill & \text{if }d_1^B+d_2^B0\text{,}\hfill \end{array}$$
(19)
Types $`(AB)`$ or $`(BA)`$:
$$s\{\begin{array}{cc}[d_1^A,\mathrm{})\hfill & \text{if }d_1^A0\hfill \\ [0,\mathrm{})\hfill & \text{if }d_1^A0\text{,}\hfill \end{array}$$
(20)
It is worth mentioning that in Proposition 1 (resp. eqs. (16-20)) the hyperbolicity of $`T(d_1,d_2,s)`$ corresponds to strict inequalities (resp. inclusions in the interior). The equality case (resp. boundary case) corresponds to the parabolicity of $`T(d_1,d_2,s)`$. There are also two special cases when $`T(d_1,d_2,s)`$ is parabolic independently of the value of $`s`$: (P), $`d_1=d_2=\mathrm{}`$ and (H), $`|d_1^A|=|d_2^A|=\mathrm{}`$.
We call the right hand side of eqs. (16-20) the instability set of $`T(d_1,d_2,s)`$. In general, it is a union of two intervals, where one of them degenerates when $`|d_1|=|d_2|`$, while the other is always nontrivial. Following the terminology of our previous work \[GSG\], we will say that the interval which persists is a “big interval”, while the other one is a “small interval”. We will say that $`T(d_1,d_2,s)`$ is (strictly) B-unstable if $`s`$ belongs to the (interior of the) big interval of instability. The proposition below makes this terminology explicit.
Proposition 2. The g.t.p.t. $`T(d_1,d_2,s)`$ is B-unstable if (and only if) the triple $`(d_1,d_2,s)`$ satisfies the following conditions:
(P)
$$s\{\begin{array}{cc}[d_1+d_2,\mathrm{})\hfill & \text{if }d_1,d_20\hfill \\ [0,\mathrm{})\hfill & \text{if }d_1,d_20\hfill \\ [d_1,\mathrm{})\hfill & \text{if }d_10,d_20\hfill \end{array}$$
(21)
(E)
$$s\mathrm{mod}\overline{\pi }\{\begin{array}{cc}[\mathrm{d}_1+\mathrm{d}_2,\overline{\pi }]\hfill & \text{if }d_1,d_20\hfill \\ [0,\mathrm{d}_1+\mathrm{d}_2+\overline{\pi }]\hfill & \text{if }d_1,d_20\hfill \\ [\mathrm{d}_2,\overline{\pi }+\mathrm{d}_1]\hfill & \text{if }d_10,d_20\hfill \end{array}$$
(22)
(H) The case $`(AA)`$
$$s\{\begin{array}{cc}[d_1^A+d_2^A,\mathrm{})\hfill & \text{if }d_1^A,d_2^A0\hfill \\ [0,\mathrm{})\hfill & \text{if }d_1^A,d_2^A0\hfill \\ [d_1^A,\mathrm{})\hfill & \text{if }d_1^A0,d_2^A0,\hfill \end{array}$$
(23)
or $`|d_1^A|=|d_2^A|=\mathrm{}`$ and arbitrary $`s`$.
(H) The case $`(BB)`$
$$s\{\begin{array}{cc}[d_1^B+d_2^B,\mathrm{})\hfill & \text{if }d_1^B+d_2^B0\hfill \\ [0,\mathrm{})\hfill & \text{if }d_1^B+d_2^B0\hfill \end{array}$$
(24)
(H) The cases $`(AB)`$ or $`(BA)`$
$$s\{\begin{array}{cc}[d_1^A,\mathrm{})\hfill & \text{if }d_1^A0\hfill \\ [0,\mathrm{})\hfill & \text{if }d_1^A0\text{.}\hfill \end{array}$$
(25)
Obviously, the conditions (21-25) for B-unstable g.t.p.t.s are the same as those which appeared in \[GSG\] for the corresponding non-magnetic cases.
## 4 The Main Theorem
Let $`Q`$ be a billiard table and $`vV`$ be an arbitrary point in the phase space of the billiard map. Set $`v_1=v,v_2=\varphi (v),d_i=d(v_i),i=1,2`$, and let $`s=s(v)`$ be the length of the particle trajectory between the origin points of $`v_1`$ and $`v_2`$ respectively. We will associate with $`v`$ a formal g.t.p.t. $`T(v)=T(d_1,d_2,s)`$, which parameters are defined by the triplet $`(d_1,d_2,s)`$. We denote by $`\lambda (v)`$ the Lyapunov exponent of the billiard $`Q`$ and by $`\overline{\lambda }(v)`$ the Lyapunov exponent of $`T(v)`$ (see Sect. 1), which are defined for $`\mu `$-almost all $`vV`$.
Using Proposition 2 we introduce the following special class of points of the phase space of the billiard map.
Definition 1. A point $`vV`$ of the billiard phase space is
a) B-hyperbolic (or strictly B-unstable) if the corresponding g.t.p.t. $`T(v)`$ is strictly B-unstable;
b) B-parabolic if the corresponding g.t.p.t. $`T(v)`$ is B-unstable and parabolic (i.e., $`s`$ belongs to the boundary of the appropriate interval (21-25) );
c) B-unstable if the corresponding g.t.p.t. $`T(v)`$ is B-unstable (i.e., B-parabolic or B-hyperbolic);
d) eventually strictly B-unstable if there is some integer $`n`$ such that $`T(\varphi ^i(v))`$ is B-unstable for $`0i<n`$ and $`T(\varphi ^n(v))`$ is strictly B-unstable.
Below we formulate the main theorem of the present work.
Theorem 1. Let $`Q`$ be a billiard table on $`M(r,\beta )`$. If $`\mu `$-almost every point of the billiard phase space is eventually strictly B-unstable, then the Lyapunov exponent $`\lambda `$ is positive $`\mu `$-almost everywhere.
Proof. The proof of the theorem is based on the cone field method which has been initially applied to the planar billiards in \[Wo1\], \[Wo2\].
A cone in $`T_vV`$ corresponds to an interval in the projectivization $`B_v`$. Therefore, a cone field, $`𝒲`$, is determined by a function, $`W()`$, on $`V`$, where each $`W(v)`$ is an interval in the projective coordinate $`\chi `$. We define the function $`W(v)`$ as in \[GSG\]. For completeness, we repeat this definition below.
$$\begin{array}{cc}\text{ (P) and (E) }& W(v)=\{\begin{array}{cc}[K(v),+\mathrm{}]\hfill & \text{if }K(v)0\hfill \\ [\mathrm{},K(v)]\hfill & \text{if }K(v)0\hfill \end{array}\hfill \\ & \\ \text{ (H) }& W(v)=\{\begin{array}{cc}[K(v),+\mathrm{}]\hfill & \text{if }K(v)\xi \hfill \\ [\mathrm{},K(v)]\hfill & \text{if }K(v)\xi \hfill \end{array}\hfill \end{array}$$
As it follows from Lemma 2 in \[GSG\], this cone field is eventually strictly preserved by the billiard map if the conditions of Theorem 1 are satisfied. By this fact the proof of the theorem follows immediately from Wojtkowski’s theorem (Theorem 1 in \[Wo2\]). $`\mathrm{}`$
Applying the method developed in \[Wo2\], one can actually estimate from below the Lyapunov exponent using the cone field defined above. The result is given by the next theorem.
Theorem 2. Let $`Q`$ be a billiard table satisfying the assumptions of Theorem 1, then
$$h(\varphi )=_V\lambda (v)𝑑\mu _V\overline{\lambda }(v)𝑑\mu .$$
Proof. The proof follows immediately by the repetition of calculations given in the proof of the analogous theorem for the non-magnetic case (see Theorem 2 in \[GSG\]). $`\mathrm{}`$
## 5 Applications and Examples
Theorem 1 together with Proposition 2 lead to a simple geometric criterion for billiard tables with hyperbolic dynamics. In this section we apply this criterion to construct various classes of hyperbolic billiards on $`M(r,\beta )`$.
### 5.1 Elementary billiard tables
There is a class of billiard tables, where the application of Theorem 1 gives an especially simple criterion for hyperbolicity. This class consists of billiard tables $`Q`$, whose boundary is a finite union of arcs, $`\mathrm{\Gamma }_i`$, of constant geodesic curvature, $`\kappa (\mathrm{\Gamma }_i)=\kappa _i`$. We call these tables elementary. We will use the notation $`\mathrm{\Gamma }_i^+`$ (resp. $`\mathrm{\Gamma }_i^{}`$) if $`\kappa (\mathrm{\Gamma }_i)>0`$ (resp. $`\kappa (\mathrm{\Gamma }_i)0`$). Let $`C_i`$ be the curve of constant geodesic curvature such that $`\mathrm{\Gamma }_iC_i`$ and $`D_iM`$ be the corresponding disk ($`C_i=D_i`$). Since the representation $`Q=_{i=1}^N\mathrm{\Gamma }_i`$ is unique, we call $`\mathrm{\Gamma }_i`$ the components of $`Q`$. In the following, we consider elementary billiard tables for which $`|\kappa _i|\beta `$. One may easily see that the fulfillment of this inequality is necessary for billiards satisfying the conditions of Theorem 1 (see discussion in the Section 5.2 for billiards with boundaries of general type).
(E) Elliptic case ($`r_{eff}>0`$). Let $`DM`$ be a disc such that $`D`$ is the circle whose geodesic curvature $`\kappa `$ satisfies $`\kappa \beta `$. We define the component $`DM`$ as set of the points $`m^{}M`$ satisfying the condition $`\overline{m^{}m}=\overline{\pi }`$ for some point $`mD`$, where $`\overline{m^{}m}`$ is the length of the particle trajectory between the points $`m`$, $`m^{}`$. We will refer to $`D`$ as the dual component of $`D+D`$. Straightforward analysis shows that $`D`$ is the ring whose width equals to the diameter of $`D`$ and its radius is defined by $`\xi `$ (for $`M=𝐑^\mathrm{𝟐}`$ its radius is $`\beta ^1`$), see figs. 5a,b,c. When $`M=𝐒^\mathrm{𝟐}`$ and $`\beta =0`$, $`D`$ is the disk obtained from $`D`$ by reflection about the center of $`𝐒^\mathrm{𝟐}`$, as it has been defined in \[GSG\].
Let us also introduce the terminology: If $`RSM`$ are regions with piecewise $`C^1`$ boundaries, we call an inclusion $`RS`$ proper if $`RintS\mathrm{}`$. The application of Theorem 1 to the elementary billiard tables in the case $`r_{eff}>0`$ leads to the following criterion for hyperbolicity.
Corollary E. Let $`QM`$ be an elementary billiard table whose boundary consists of $`N>1`$ components of type plus or minus. Suppose $`Q`$ satisfies the following conditions:
Condition E1. For every component $`\mathrm{\Gamma }_i^+`$ of $`Q`$ we have $`D_iQ`$. Besides, either $`D_iQ`$, or $`D_iMQ`$, where the inclusions are proper;
Condition E2. For every component $`\mathrm{\Gamma }_j^{}`$ we have $`D_jMQ`$, and the inclusions $`D_jMQ`$, or $`D_jQ`$ are proper.
Then the billiard in $`Q`$ is hyperbolic.
Outline of proof: The assumptions of Corollary E imply those of Theorem 1.
Remark. Suppose $`Q^{}=MQ`$ is connected. If $`Q`$ satisfies Conditions E1 and E2, then $`Q^{}`$ also does, and hence the billiard in $`Q^{}`$ is hyperbolic.
Examples. “Lorenz gas” billiards. Such billiards are obtained by removing from $`M`$ a number of disjoint discs $`D_i`$, so that $`Q=MD_i`$. If all the intersections $`D_i\pm D_j`$ $`ij`$, are empty, then the billiard in $`Q`$ is hyperbolic by Corollary E. The simplest example of such hyperbolic billiard is obtained by removing two disks from the magnetic plane, see fig. 6a.
The intersections $`D_i\pm D_j`$ $`ij`$, are always empty, if all the discs are contained inside of a free-flight particle trajectory (i.e., if all the discs lie inside a circle of geodesic curvature $`\beta `$). Such billiards are the “magnetic” analogs of the non-magnetic hyperbolic billiard tables on the sphere, obtained by removing a finite number of disjoint disks from one hemisphere \[GSG\]. The examples of hyperbolic billiards of this type on $`𝐒^\mathrm{𝟐}`$, $`𝐑^\mathrm{𝟐}`$ and $`𝐇^\mathrm{𝟐}`$ are shown in fig. 6b,c,d.
One can consider also unbounded billiard tables $`Q`$ obtained by removing an infinite number of disjoint disks from $`𝐑^\mathrm{𝟐}`$, $`𝐇^\mathrm{𝟐}`$. The simplest example of this type is obtained by removing a chain of equal disks from $`M=𝐑^\mathrm{𝟐}`$, as shown in fig. 7a (this billiard can be also seen as cylinder with one hole). Because of the translation symmetry, one needs to check the non-intersection condition only for one disk. The non-intersection condition is also necessary for hyperbolicity of such billiards. If it is not satisfied, then $`Q`$ has at least two stable g.t.p.t.s (see fig. 7a).
Another type of unbounded hyperbolic billiard tables can be obtained by removing a lattice of the disks from $`M=𝐑^\mathrm{𝟐},𝐇^\mathrm{𝟐}`$. The example of such billiard shown in fig. 7b, is equivalent to the torus with one hole. Here, again, because of the translation symmetry, one has to check the non-intersection condition only for a single disk.
“Flowers” like billiards. Consider a simply connected billiard table $`Q`$, whose boundary consists of several circular arcs of positive and negative curvature satisfying the condition $`|\kappa _i|\beta `$. Such billiards were originally introduced by Bunimovich \[Bu1\] \[Bu2\] as examples of planar (non-magnetic) hyperbolic billiards with convex boundary. It has been demonstrated that for $`r=0`$, $`\beta =0`$ such billiards are hyperbolic if the conditions $`D_iQ`$ are satisfied for each convex component of the boundary. For $`r_{eff}>0`$ we have by Corollary E the additional requirement: $`QD_i=\mathrm{}`$ for each component of the boundary (compare with the analogous conditions in \[GSG\] for the case $`\beta =0`$, $`M=𝐒^\mathrm{𝟐}`$). The examples of hyperbolic billiards $`Q`$ on $`𝐒^\mathrm{𝟐}`$, $`𝐑^\mathrm{𝟐}`$, $`𝐇^\mathrm{𝟐}`$ of “flower” type satisfying the conditions of Corollary E are shown in fig. 8a,b,c. It follows from the remark above that billiards in the domain $`Q^{}=MQ`$ are also hyperbolic.
(H+P) Hyperbolic and parabolic cases ($`r_{eff}0`$). The criterion for hyperbolicity in this case is given by the following corollary.
Corollary H. Let $`QM`$ be an elementary billiard table, and let $`Q`$ consist of $`N>1`$ components. If $`Q`$ satisfies conditions:
Condition H1. For every convex component $`\mathrm{\Gamma }_i^+`$ of $`Q`$, we have $`D_iQ`$;
Condition H2. For every concave component of $`Q`$, we have $`\kappa (\mathrm{\Gamma }_i^{})\beta `$ and for every convex component $`\kappa (\mathrm{\Gamma }_i^+)\xi `$.
Then the billiard in $`Q`$ is hyperbolic.
Outline of proof: The assumptions of Corollary H imply those of the Theorem 1.
Remark. When $`\beta 0`$ and $`r0`$, the condition H2 is automatically fulfilled and Corollary H turns to be the classical criterion of Bunimovich \[Bu2\] for hyperbolicity of planar, non-magnetic billiard tables.
Examples. Analogs of Sinai billiards. The boundary of these billiards consists of concave arcs $`\mathrm{\Gamma }_i^{}`$ of constant curvature (see fig. 9). If the condition $`\kappa (\mathrm{\Gamma }_i^{})\beta `$ is satisfied for each component of the boundary, then the billiard is hyperbolic by Corollary H.
Analogs of Bunimovich billiards. The example of hyperbolic billiard table with convex components satisfying the conditions H1, H2 is shown in fig. 10.
Remark. The assumptions in Corollaries E and H that $`N>1`$ and that the inclusions be proper are needed only to exclude certain degenerate situations, where each $`vV`$ is B-parabolic. This is the case, for instance, if $`Q`$ is a disc, or the annulus between concentric circles.
### 5.2 Hyperbolic billiard tables with boundary of general type
Let us consider billiard tables on $`M(r,\beta )`$ with piecewise smooth boundary, $`Q=_i\gamma _i`$ of general type. The components $`\gamma _i`$ are $`C^2`$ smooth curves parameterized by the arclength $`l`$, whose curvature $`\kappa _i(l)`$ has the same sign along each $`\gamma _i`$. We will refer to $`\gamma _i`$ as convex component if $`\kappa _i(l)>0`$, or as concave component if $`\kappa _i(l)0`$. Let us denote $`\kappa (\gamma _i)=\mathrm{max}\{\kappa _i(l),l\gamma _i\}`$ for the convex components, and $`\kappa (\gamma _i)=\mathrm{min}\{\kappa _i(l),l\gamma _i\}`$ for the concave components.
Following the terminology in \[Wo2\], we introduce the class of convex scattering curves on $`M(r,\beta )`$.
Definition 2. A smooth convex curve $`\gamma M`$ is (strictly) convex scattering if for any $`vV`$, such that the origin points of $`v`$ and $`\varphi (v)`$ belong to $`\gamma `$, the corresponding g.t.p.t. $`T(v)`$ is (strictly) B-unstable.
A curve $`\gamma `$ is convex scattering if one of the relevant conditions (21-25) is satisfied for each pair of points on $`\gamma `$. Regarding the planar non-magnetic case, this leads to the definition of Wojtkowski \[Wo2\] for convex scattering curve. Let us introduce the parameter $`R(l)=(\kappa (l)\beta )^1`$. Considering the infinitesimally close points on $`\gamma `$ we show in Appendix that the condition $`R^{\prime \prime }(l)0`$ is necessary for $`\gamma `$ to be convex scattering. It should be noted, that this condition is also sufficient in the planar, non-magnetic case (see \[Wo2\]), but not for generic parameters $`r`$, $`\beta `$ (see \[GSG\] for $`\beta =0`$ case).
In what follows, we formulate the principles for design of hyperbolic billiards satisfying the conditions of Theorem 1. Let $`Q`$ be a billiard table satisfying the conditions of Theorem 1. Then each convex component of $`Q`$ has to be convex scattering and consequently, the condition $`R^{\prime \prime }0`$ holds along each convex component of the boundary. There is an additional restriction on the curves $`\gamma _i`$ which compose the boundary of $`Q`$. It follows from Proposition 2 that for billiards satisfying the conditions of Theorem 1 the sign of $`K(v)`$ ($`d(v)`$) depends only on the origin point of $`v`$ (there is no dependence on $`\theta `$) for any $`vV`$, i.e., $`K(v)`$ ($`d(v)`$) has the same sign along $`\gamma _i`$ as $`\kappa (\gamma _i)`$. This happens if for each component $`\gamma _i`$, $`|\kappa (\gamma _i)|\beta `$ (the magnetic field is sufficiently weak). Thus, in what follows we particularly exclude from our consideration the magnetic billiards with flat boundaries. Such billiards do not satisfy the conditions of Theorem 1.
Design of hyperbolic billiard tables in the (E) case.
By Definition 2 a curve $`\gamma `$ is convex scattering if it is convex and the condition
$$d_1+d_2s\overline{\pi },$$
(26)
holds for any pair of points on $`\gamma `$. For simplicity of exposition, we will restrict our attention for $`M=𝐑^\mathrm{𝟐},𝐇^\mathrm{𝟐}`$ to the bounded billiard tables and for $`M=𝐒^\mathrm{𝟐}`$ to the billiard tables which can be placed in a hemisphere. Theorem 1 yields the following principles for the design of piecewise billiard tables with hyperbolic dynamics in (E) case:
P1: $`|\kappa (\gamma _i)|\beta `$ for all components.
P2: All convex components of $`Q`$ are convex scattering.
P3: Any convex component of $`Q`$ has to be “sufficiently far”, but not “too far”, from any other component. Any concave component has to be not “too far”, from any other concave component.
The precise meaning of P3 is that the parameters of any two consecutive bouncing points, which belong to different components of the boundary, satisfy the condition (22). In particular it implies the set of restrictions on the angles between consecutive components of the boundary. It can be formulated as an additional principle.
P4: Let $`\gamma _i,\gamma _{i+1}Q`$ be two adjacent components, meeting at a vertex. If both $`\gamma _i`$ and $`\gamma _{i+1}`$ are convex, then the interior angle at the vertex is greater than $`\pi `$. If $`\gamma _i`$ and $`\gamma _{i+1}`$ have different sign of curvature, then the angle in question is greater or equal to $`\pi `$.
Another restriction which arises from P3 is that the length (equivalently the time) of free-flight between any two consequent bouncing points on the boundary of the billiard has to be not greater than $`\overline{\pi }`$. In other words, the billiard table has to be “smaller” than circle drawn by a free-flight particle on $`M(r,\beta )`$.
Examples. The examples of the hyperbolic billiards on $`𝐑^\mathrm{𝟐}`$ satisfying the above principles are shown in fig. 11a,b. A bounded Sinai-like billiard, whose boundary consists of (strictly) concave components (fig. 11a) always satisfies the principles P1-P4 for sufficiently weak magnetic field.
The example of a convex billiard is shown in fig. 11b. It is a cardioid, whose boundary is strictly convex scattering curve for $`\beta =0`$ (see \[Wo2\]). For $`\beta =0`$ this billiard is hyperbolic, as it follows from Theorem 1. Since strictly convex scattering curve remains to be such under small perturbations of $`\beta `$, the billiard in fig. 11b is hyperbolic for sufficiently weak magnetic field.
Design of hyperbolic billiard tables in the (P+H) case.
Definition 2 leads to the following geometric conditions on the convex scattering curve in the (P+H) case. A convex curve $`\gamma `$ is convex scattering if $`\kappa (\gamma )\xi `$ and for each pair of points on $`\gamma `$
$$d_1+d_2s.$$
(27)
Theorem 1 yields the following principles for the design of billiard tables with hyperbolic dynamics in the (P+H) case:
P1: $`\kappa (\gamma _i)\xi `$ for any convex component of $`Q`$ and $`\kappa (\gamma _i)\beta `$ for any concave component of $`Q`$.
P2: All convex components of $`Q`$ are convex scattering.
P3: Any convex component of $`Q`$ is “sufficiently far” from any other component.
More precisely, condition P3 means that any two consecutive bouncing points of the billiard ball, which belong to different components, satisfy eqs. (23-25). In particular, this yields, the same inequalities (P4) as in (E) case, for the interior angles between consecutive components of $`Q`$.
Examples. In (P+H) case, any concave billiard is hyperbolic if the condition $`\kappa (\gamma _i)\beta `$ is fulfilled for each component of the boundary. As in the case (E), the examples of the convex hyperbolic billiards can be obtained from their non-magnetic counterparts satisfying the conditions of Theorem 1.
Finally, it should be noted, that formulated above principles for design of hyperbolic billiards on $`M(r,\beta )`$ are robust under small perturbations of $`\beta `$, $`r`$ and the billiard wall. Generally, one can construct hyperbolic billiards on magnetic surfaces of constant curvature on the basis of the corresponding non-magnetic planar billiards satisfying Wojtkowski’s criterion.
## 6 Conclusions
In the present paper we have formulated the criterion for hyperbolic dynamics in billiards on surfaces of constant Gaussian curvature $`r`$ in the presence of a homogeneous magnetic field $`\beta `$ perpendicular to the surface. The criterion is valid for all values of $`r`$, $`\beta `$ and its geometric realization depends only on the type of linearized dynamics (elliptic, parabolic or hyperbolic). In this way we extend our recent results in \[GSG\] to the case of magnetic surfaces of constant curvature. The basic property, which allows unification of the hyperbolicity criteria for the magnetic and non-magnetic billiards on surfaces of constant curvature, is the equivalence between the geometric optics in both cases. In fact, in terms of special parameters $`d_i`$, $`s_i`$ the geometric optics depend only on the effective curvature $`r_{eff}=r+\beta ^2`$ of the surface. It is important to stress, that the dynamics in magnetic and non-magnetic billiards are very different (e.g., the magnetic field breaks time reversal symmetry). It is the only linearized dynamics, which are the same for the considered systems. Applying the hyperbolicity criterion, we were able to construct the different classes of hyperbolic billiards for each type of the linearized dynamics (equivalently for each of the signs of $`r_{eff}`$).
There are two types of necessary conditions which arise for hyperbolic billiards satisfying our criterion. The first one is a requirement for the convex components of the boundary to be convex scattering. As a consequence, the inequality $`R^{\prime \prime }(l)0`$ has to be satisfied along each convex component. This inequality is generalization of well-known Wojtkowski’s condition \[Wo2\] for convex component of planar (non-magnetic) hyperbolic billiard. It has been demonstrated for planar non-magnetic billiards in \[Bu3\], \[Bu4\], \[Do\] that Wojtkowski’s criterion can be considerably strengthened. This suggests, in particular, that condition $`R^{\prime \prime }(l)0`$ can be relaxed for general parameters $`r`$, $`\beta `$ by employing invariant cone fields, different from the one used in the present paper (see discussion in \[GSG\]). The second type of conditions is specific for magnetic billiards. This is a requirement of “weakness” for the magnetic field compared to the curvature of the billiard boundary. For generic systems, such condition is expected, in order to prevent stable skipping orbits close to the boundary. It has been shown in \[BR\], (see also \[BK\]) that billiard with sufficiently smooth boundary possesses invariant tori corresponding to skipping trajectories. It seems that in the strong field regime a part of stable periodic orbits has to survive even if the smoothness of the boundary is broken. It remains, however an open question, whether the condition $`|\kappa _i|\beta `$ can be relaxed for generic billiard.
The positive Lyapunov exponent for a billiard implies strong mixing properties: countable number of ergodic components, positive entropy, Bernoulli property etc. It should be pointed out, however, that ergodicity does not automatically follow from the positivety of Lyapunov exponent. Nevertheless, one can expect that billiards satisfying the conditions of Theorem 1 will be typically ergodic. It seems that the methods developed for the proof of ergodicity of planar hyperbolic billiards can be extended to the class of billiards considered in the present paper.
## Acknowledgment
The author is indebted to Professor U. Smilansky for proposing this investigation and for critically reading the manuscript. The author would like to thank Andrey Shapiro de Brosh for interesting and inspiring discussions, and various valuable remarks.
This work was supported by the Minerva Center for Nonlinear Physics of Complex Systems.
## 7 Appendix
We will investigate the conditions under which a convex arc on the surface of constant curvature $`M`$ in the presence of magnetic field $`\beta `$ is convex scattering.
For simplicity of exposition, we consider the case, when $`M`$ is magnetic plane. Let $`\gamma (l)M`$ be any smooth curve parameterized by arclength $`l`$, and let $`\kappa (l)`$ be the geodesic curvature of $`\gamma `$. Let now $`\gamma (l_0)`$ and $`\gamma (l_1)`$ be two points on $`\gamma `$, such that the arc of $`\gamma `$ between $`\gamma (l_0)`$ and $`\gamma (l_1)`$ lies entirely on one side of straight line passing through $`\gamma (l_0)`$ and $`\gamma (l_1)`$. We choose cartezian coordinate system $`(x,y)`$ in such a way that $`y(l_0)=y(l_1)=0`$, $`x(l_0)=x(l_1)`$ and the arc of $`\gamma `$ between $`\gamma (l_0)`$ and $`\gamma (l_1)`$ lies above $`x`$-axis, see fig. 12. Let $`\alpha (l)`$ be the angle, which $`\frac{d\gamma }{dl}`$ makes with $`x`$-axis , then
$$\frac{dx}{dl}=\mathrm{cos}\alpha ;\frac{dy}{dl}=\mathrm{sin}\alpha ;\frac{d\alpha }{dl}=\kappa .$$
(28)
We introduce also an auxiliary variable $`\delta `$, such that $`\beta x=\mathrm{sin}\delta `$.
For $`\beta >0`$ there are two different particle trajectories connecting the points $`\gamma (l_0)`$ and $`\gamma (l_1)`$ (resp. two different g.t.p.t.s corresponding to these points), see fig. 12. Below, we consider the trajectory which lies in the lower halfplane. Then, the results for trajectory in the upper halfplane are obtained by the change of the sign of $`\beta `$ to the opposite. Let $`\theta =\alpha +\delta `$. Then, at the points $`l_{0,1}`$, $`\theta (l_{0,1})`$ are the angles between $`\gamma `$ and the particle trajectory connecting $`\gamma (l_0)`$ and $`\gamma (l_1)`$. Set $`\mathrm{\Delta }=sd_1d_2`$. By eq. 14 we get
$`\mathrm{\Delta }=\beta ^1{\displaystyle \left[d\left(\mathrm{arctan}\left(\frac{\beta \mathrm{sin}\theta }{\kappa \beta \mathrm{cos}\theta }\right)\right)+d\delta \right]}`$
$$=𝑑l\left(\frac{\kappa ^{}\mathrm{sin}\theta +\kappa \left(\kappa +\beta \frac{\mathrm{sin}\alpha }{\mathrm{sin}\delta }\right)\left(\frac{\mathrm{cos}\alpha }{\mathrm{cos}\delta }\mathrm{cos}\theta \right)}{\kappa ^2+\beta ^22\beta \kappa \mathrm{cos}\theta }\right)$$
(29)
We separate the last integral into the sum of two parts. The first one is
$`I={\displaystyle 𝑑l\left(\frac{\kappa ^{}\mathrm{sin}\theta }{\kappa ^2+\beta ^22\beta \kappa \mathrm{cos}\theta }\right)}={\displaystyle 𝑑y\left(\frac{R^{}}{1+4R^2\kappa \beta \mathrm{sin}^2\theta /2}\right)},`$
where $`R^1(l,\beta )=\kappa (l)\beta `$. Since $`y(l_0)=y(l_1)=0`$, we obtain
$$I=𝑑l\left(\frac{yR^{\prime \prime }}{1+4R^2\kappa \beta \mathrm{sin}^2\frac{\theta }{2}}\frac{yR^{}(4R^2\kappa \beta \mathrm{sin}^2\frac{\theta }{2})^{}}{(1+4R^2\kappa \beta \mathrm{sin}^2\frac{\theta }{2})^2}\right)=\frac{R^{\prime \prime }L^3\kappa }{12}+O(L^4),$$
(30)
where $`L=l_1l_0`$ is the length of the curve between the points $`\gamma (l_0)`$, $`\gamma (l_1)`$. Analogously, for second part we have
$`II={\displaystyle 𝑑l\left(\frac{\kappa \left(\kappa +\beta \frac{\mathrm{sin}\alpha }{\mathrm{sin}\delta }\right)\left(\frac{\mathrm{cos}\alpha }{\mathrm{cos}\delta }\mathrm{cos}\theta \right)}{\kappa ^2+\beta ^22\beta \kappa \mathrm{cos}\theta }\right)}`$
$`={\displaystyle 𝑑y\left(\frac{\kappa \frac{\mathrm{sin}\theta }{\mathrm{cos}\delta }\left(\kappa \frac{\mathrm{sin}\delta }{\mathrm{sin}\alpha }+\beta \right)}{\kappa ^2+\beta ^22\beta \kappa \mathrm{cos}\theta }\right)}=O(L^4).`$
Adding both parts we obtain finally
$$\mathrm{\Delta }=I+II=\frac{R^{\prime \prime }L^3\kappa }{12}+O(L^4).$$
(31)
Thus, if the curve $`\gamma `$ is convex scattering, then the condition $`R^{\prime \prime }(l,\beta )0`$ holds everywhere on $`\gamma `$. Considering trajectories of the second type (i.e., trajectories which lay in the upper halfplane), we obtain the condition $`R^{\prime \prime }(l,\beta )0`$ for convex scattering curves. However, it easy to see, that $`R^{\prime \prime }(l,\beta )0`$ actually implies $`R^{\prime \prime }(l,\beta )0`$.
Repeating the same analysis for general $`M(r,\beta )`$ we have found (see also \[GSG\] for $`\beta =0`$ case) that eq. 31 holds for all surfaces of constant curvature. As a consequence, $`R^{\prime \prime }0`$ is a necessary condition for convex scattering on $`M(r,\beta )`$. On the contrary, if the strict inequality $`R^{\prime \prime }<0`$ holds along $`\gamma `$, then by eq. 31, any sufficiently small piece of $`\gamma `$ is convex scattering.
|
warning/0003/astro-ph0003449.html
|
ar5iv
|
text
|
# Limits on magnetic field strength of the extended nebula of PSR B1706-44 from optical, X-ray and TeV observations
## 1 Introduction
PSR B1706-44 is a young Vela-like pulsar (spin-down age 17.5 kyr), with a period of 102 ms and a large spin-down power (Ė = 3.4 x 10<sup>36</sup> ergs/s). It was discovered during a 20 cm radio pulsar survey of the southern Galactic plane (Johnston et al. 1992), and later detected in soft X-rays during the ROSAT mission (Becker et al. 1992, 1995) and identified as a pulsed GeV source by EGRET (Thompson et al. 1992). Very high energy (VHE) $`\gamma `$-ray observations above 1 TeV from CANGAROO (Kifune et al. 1995; Kifune 1997) and above 0.3 TeV from Durham (Chadwick et al. 1997) detected and confirmed the existence of unpulsed radiation of statistical significance from this source. The CANGAROO detection is consistent with a point source, with angular extend not exceeding $`0.^{}12`$ – the pixel size of the imaging camera. A dispersion based distance measure of Taylor & Cordes (1993) places PSR B1706-44 at $``$1.8 kpc. A possible SNR G 343.1-2.3 association with the pulsar was proposed by McAdam et al. (1993), but this was later found to be unlikely by Frail et al. (1994) and Nicastro et al. (1996).
Chakrabarty & Kaspi (1998) gave a red-band upper limit of $`R>\mathrm{\hspace{0.33em}18}`$ to the pulsar. A 3$`\sigma `$ upper limit to the pulsar magnitude of $`V=24.5`$ was given by Lundqvist et al. (1999), and this is consistent with the theoretical prediction of $`V=24.12`$ (Urama & Okeke 1998, and references therein). Mignani et al. (1999), using the same data as Lundqvist et al., obtained an upper limit of $`V>\mathrm{\hspace{0.33em}27.5}`$.
The photon flux above 1 TeV from this unpulsed source is only two times smaller compared to the flux of the Crab Nebula at the same energy. This is remarkable in the context of a synchrotron-inverse Compton (on external photon fields) interpretation for the TeV $`\gamma `$-rays from PSR B1706-44, since the synchrotron nebula of this source is very weak compared to the Crab Nebula’s synchrotron intensity. Whereas the Crab Nebula’s intense synchrotron emission is the result of a large $`B`$ ($`10^4`$ G) field, we must have a much weaker field for this plerion to avoid a bright synchrotron nebula.
De Jager (1995) speculated that a Vela-like compact synchrotron nebula ($`<1^{}`$, Harnden et al. 1985) may be present, which would account for most of the unpulsed X-ray emission from this source, since the pulsar wind magnetic field (which scales as $`B1/R`$) would predict a compact synchrotron nebula, if the conditions are similar to Vela. The corresponding particle density required for the synchrotron nebula would however be too low to produce a detectable inverse Compton compact nebula, since the target photon density from the CMBR and Galactic disk would be too small.
A compact nebula was indeed discovered (Finley et al. 1998), which confirms the abovementioned interpretation. The scenario of Harding & de Jager (1997), with a detailed analysis by Aharonian et al. (1997) may apply in this case: electrons are streaming away from the compact X-ray nebula into a low-B extended plerion. If the diffusion coefficient for electrons in the extended plerion is small enough, it may be possible to trap enough electrons to account for the TeV emission by IC scattering.
Frail et al. (1994) detected a radio synchrotron nebula associated with the pulsar, which has a radius of $`2^{}`$ at 20 cm wavelength. This may represent the low energy counterpart of the TeV nebula, since the associated $`B`$ in the radio nebula is smaller than the pulsar driven field in the compact X-ray nebula, and the size is just smaller than the angular resolution of the CANGAROO telescope.
If the TeV $`\gamma `$-rays originate from the extended (radio?) nebula, we will expect a bright extended nebula in optical and/or X-rays, if the field strength is large enough. No extended X-ray nebula was however seen, resulting in an upper limit of four times the flux of the point source/compact X-ray nebula at 1 keV, for $`R=2^{}`$ (K. Brazier 1998, personal communication). This is consistent with the expectation that the pulsar field strength (which drops as $`1/R`$ outside the light cylinder) should start to drop below the ambient field strength of a few $`\mu `$G at a distance of a few arcminutes (de Jager & Harding 1998). The consequence of this is that both the synchrotron brightness and synchrotron characteristic frequencies should drop with increasing $`R`$. Scaling from the Durham detection above 0.3 TeV, and the expected $`B3\mu `$G in the extended nebula (de Jager & Harding 1998), we may use equation (5) of Aharonian et al. (1997) to calculate the synchrotron frequency which corresponds to the Durham detected $`\gamma `$-rays near 0.3 TeV, assuming that the CMBR is mostly responsible for the inverse Compton scattering. This gives a frequency of
$$\nu 1.5\times 10^{15}\left(\frac{E_\gamma }{0.3TeV}\right)\left(\frac{B}{3\mu G}\right)Hz,$$
which is within a factor of 3 from the V-band. The non-detection of the unpulsed component at 20 GeV (Thompson et al. 1996) hints at a marginal turnover in the $`\gamma `$-ray spectrum (de Jager & Harding 1998), which corresponds to a turnover frequency as low as $`10^{14}`$ Hz, which is well below the V-band.
At a distance of 1.8 kpc to the source, we expect significant interstellar absorption in the blue and UV bands. We therefore take wide field V-band images as a compromise between interstellar absorption and avoiding a spectral turnover at too low frequencies. By taking overlapping images, we cover a total field-of-view of about $`10^{}`$, which should include the expected TeV source. Even if the full extend of the TeV source is $`10^{}`$, we still expect a typically centrally brightened image, resulting in a radial gradient from the source. Failure to detect the optical counterpart will allow us to set upper limits on the magnetic field strength in the extended nebula, since any stronger field should have resulted in a brighter optical nebula.
## 2 Observations and data analyses
V band CCD observations of PSR B1706-44 were carried out on May 24, 1999 at the South African Astronomical Observatory (SAAO) 1.0 m telescope in Sutherland. Five frames, each $`5.^{^{}}3\times 5.^{^{}}3`$ (pixel size = 0.31 arcsec), were taken. The first one was centered on the radio pulsar position (Johnston et al. 1992, 1995), whereas the other four frames were offset in such a way that they have an overlap of $``$10% with each other around the pulsar position. The overall region observed around the pulsar is $`9.^{^{}}5\times 9.^{^{}}5`$. The observation time for each frame was 12 minutes. Table 1 gives the position (RA, Dec), the starting time (MJD), $`K`$ values and the duration time (seconds) log for the five frames.
Observations were affected by the moonlight (about 60% bright). The average counts for each frame were decreasing with time, probably due to the moon’s effect across the sky. The data were cleaned and flat-fielded using the IRAF image-processing software. Fig. 1 shows a V band CCD image of the central frame with the pulsar and the standard star positions labeled + and 1 (Mignani et al. 1999, their Figs. 1 & 2), respectively. Only the inner part (25% of the total area) of Frame 1 is shown on Fig. 1.
For each frame, we made a histogram of the number of pixels versus intensity per pixel using a logarithmic binning procedure prescribed by
$$K=N_o\mathrm{log}(N/N_{min}),$$
where $`N_o`$ ($`=99/\mathrm{log}(N_{max}/N_{min})`$) is a normalizing constant, as shown in Fig. 2 for the central frame. This binning procedure ensures that all intensities are confined to the interval $`K=1`$ and 100.
Whereas the diffuse emission is expected to be confined to a rather narrow low-$`K`$ peak in the histogram, the pixels contaminated by point sources should show a broad tail. Fig. 2 indeed shows such behaviour. It can be shown that the weak plerionic emission is expected to contribute to the low $`K`$-values ($`K<\mathrm{\hspace{0.33em}25}`$), so that we can reject the star–like objects by cutting at $`K30`$ (for frame 1) as shown in Fig. 2. This rejects only 10% of the pixels, so that we preserve 90% of the plerionic component. The threshold value of K for each frame is given in Table 1, and probably depends on the phase of the moon and the background bright stars (in each frame), which contribute to the background for each frame. Note that the final results are not sensitive to the threshold $`K`$.
By renormalising the count rate of each frame to the count rate of the central frame, we were able to construct a composite image of a $`10^{^{}}\times 10^{^{}}`$ region around the pulsar. Since the plerionic structure is expected to be circular around the pulsar (the radio nebula is at least circular), we constructed annuli around the pulsar and made a plot of the mean count rate per pixel versus radius of the annuli. These count rates are constant with radius within statistics for $`R=0`$ to $`5^{^{}}`$, which means that there is no evidence for a centrally brightened source, and hence no evidence for a plerionic structure centered on the pulsar.
## 3 Upper limits
The 3$`\sigma `$ upper limits for plerionic emission confined in circles, as well as annuli, are shown in Fig. 3. These upper limits were converted to absolute intensities, $`S_V`$, using the observed counts from Star 1 (see Fig. 1) of Chakrabarty & Kaspi (1998) with $`V=17.4`$ mag.
The 3$`\sigma `$ $`S_V`$ upper limits increase with $`R`$ as expected, and at $`R=0.^{^{}}5`$ (compact X-ray nebula), our 3$`\sigma `$ upper limit is $`V15`$, while for an extended nebula of about 4 arcminutes, the upper limit is $`V>\mathrm{\hspace{0.33em}12}`$.
Assuming that the synchrotron spectral index, $`\alpha `$, of the nebula is the same in both the optical and X-ray regions, the monochromatic X-ray flux, $`S_{1keV}`$, at 1 keV is given by $`S_{1keV}=S_V(ϵ_x/ϵ_o)^\alpha ,`$ where $`ϵ_x`$ (= 1 keV) and $`ϵ_o`$ (= 2.3 eV) are the X-ray and optical energies, respectively. The observed optical flux can be expressed as $`S(V)=S_V\mathrm{exp}(\tau _\nu )`$, where $`\tau _\nu `$ is the optical depth due to interstellar absorption, with $`\tau _\nu 1.8`$ for $`d=1.8`$ kpc. Using the above relations and the de Jager et al. (1995) expression for the TeV $`\gamma `$-ray IC energy spectrum (their Equation 2), we calculate the upper limits to the magnetic fields for different values of $`\alpha `$ between 1.0 and 2.5, inclusively. Fig. 4 shows the curves for $`\alpha `$ = 1.0, 1.5, 2.0 and 2.5, as well as the magnetic field limits from the X-ray upper limits for $`R=2^{}`$ for similar values of $`\alpha `$. The magnetohydrodynamic curve of Kennel & Coroniti (1984) model, with $`\sigma =1.0`$ (equipartition between particles and fields), is superimposed on Fig. 4, and it is considered as the lower limit to the extended nebula at large radii ($`R>\mathrm{\hspace{0.33em}2}.^{^{}}0`$). Also shown in Fig. 4 are the Galactic field strengths at the source according to Heiles (1996) and Broadbent et al. (1990).
## 4 Discussions
Our search for an extended nebula synchrotron counterpart to the TeV nebula did not reveal any nebular structure around PSR B1706-44. The TeV observations do hint at the existence of such a plerionic structure, but our results constrain the nebular magnetic field strength if we assume that the TeV $`\gamma `$-rays originate from the IC scattering of energetic pulsar wind electrons on the CMBR. At $`R=2^{}`$, we find that the optical and X-ray limits are complementary: the spectral index is constrained to $`\alpha =1.5`$ to 2.3, whereas $`B3`$ $`\mu `$G. This limit is barely above the pulsar wind solution, but below the average Galactic field values shown in Figure 4. Furthermore, our constraints on $`\alpha `$ are marginally consistent with $`\alpha =1.3\pm 0.3`$ measured for the compact X-ray nebula (Finley et al. 1998).
Our optical observations and the information gathered from other wave bands (mainly from X-ray and $`\gamma `$-ray energies) show that better quality observations (dark moon) should improve these limits, or, result in a detection. Polarimetry of the “extended nebula” of PSR B1706-44 may also help to clarify the status of synchrotron emission.
|
warning/0003/gr-qc0003082.html
|
ar5iv
|
text
|
# Closed Universes With Black Holes But No Event Horizons As a Solution to the Black Hole Information ProblemThis work was supported in part by the National Science Foundation under grant number DMS-97-32058.
## 1 Introduction
One of the outstanding questions of black hole physics is determining what happens to the information that falls into a black hole. Hawking has shown that a black hole radiates away its mass, and he pointed out that if a black hole were to completely evaporate — which it inevitably will in a universe that exists forever in cosmological proper time (a Planck size remnant would probably be inconsistent with the Bekenstein Bound (, , ).) — then any information exclusively inside the black hole would disappear from the universe, violating unitarity. Many solutions have been proposed to resolve this paradox. Hawking himself believes that unitarity is indeed violated, but it has been argued that such a resolution would be inconsistent with locality and/or conservation of energy . Susskind (, ) and ’t Hooft propose that all information inside a black hole is also completely encoded on its surface, so there is no net information inside the black hole. But at the semi-classical level, this “holographic principle” would not resolve the paradox, because the generators of an event horizon — the black hole surface — cannot end in spacetime, but at a singularity which itself would annihilate any information on the horizon. To avoid unitarity violation, the information must get outside the black hole event horizon.
Bekenstein has pointed out that the Hawking radiation is not quite a blackbody spectrum, and thus it carries some information about the initial state of the black hole. He has shown that the permitted information outflow rate can be as large as the rate of black hole’s entropy decrease, and hence it is possible for information to gradually leak out of a black hole during evaporation. However, Bekenstein emphasized that he had not demonstrated that all the information got out, just that it was possible that it did, and if only one bit of information fails to escape, unitarity will be violated. Bekenstein also did not address the semi-classical event horizon problem.
We shall show in this paper that a purely classical gravity solution to the black hole information problem is possible and consistent with all observations: the universe may have no event horizons at all. In such a universe, there would be no black hole event horizons to prevent the exchange of information between one part of the cosmos and another. A spacetime with no event horizons has a future c-boundary (, pp. 217–221) which is topologically a single point, and hence has been called an Omega Point Spacetime. It can be shown that if a spacetime’s future c-boundary is a single point, then the spacetime necessarily admits compact Cauchy surfaces, and the global spacetime topology is $`S\times R^1`$, where $`S`$ is the topology of any Cauchy surface. Even in a universe with compact Cauchy surfaces, we would expect black holes to evaporate to completion if the universe were to expand forever. Hence a spacetime which avoids the black hole information paradox because of the absence of event horizons would have to end in a final singularity before any black hole would have time to evaporate. Since the expected black hole lifetime is $`10^{64}(M/M_{})^3`$ years, , our universe would have to expand to a maximum size and recontract before $`10^{64}`$ years have passed. It can be shown (, ) that the only two simple topologies possible for universe with a maximal Cauchy hypersurface and satisfying the weak energy condition are $`S^3`$ and $`S^2\times S^1`$.
We shall construct in this paper a spherically symmetric $`S^3`$ universe with a black hole but with no event horizons. The spacetime will be shown to satisfy all the standard energy conditions. Indeed, the stress energy tensor for the spacetime with be just pressure-free dust in its expanding phase. We shall discuss various definitions of “black hole” in closed universes, and show that the spacetime we construct has a black hole by any of these definitions. The parameters of the constructed spacetime can be chosen so that the black hole is identical to a black hole in any dust spherically symmetric (Tolman-Bondi) $`S^3`$ universe — and hence it would be in appearance a black hole with event horizons according to any observations that could be carried out in the expanding phase of a closed universe. The null generators of what is apparently the event horizons stay close to the trapped surfaces during the expanding phase of the closed universe, and only expand out into the universe at large very close to the final singularity. This means the standard astrophysical analysis of black holes and their collisions (e.g., , , ) can be trusted, since they are in the short run the same as in asymptotically flat spacetime. Thus our proposal is quite different from many proposals which eliminate event horizons by eliminating black hole type trapped surfaces. In our proposal, black hole trapped surfaces exist as usual, but they do not give rise to event horizons.
Our paper will be organized as follows. In Section 2 we shall construct a Friedmann-Robertson-Walker (FRW) universe which is a standard dust FRW closed universe until arbitrarily near the final singularity when we join it to a metric which satisfies all the energy conditions, which has a final singularity, but which has no event horizons. In Section 3 we show that the no event horizon metric constructed in Section 2 satisfies the Einstein equations for a scalar field with a suitably chosen exponential potential. In Section 4, we shall generalize the FRW $`w=1/3`$ universe to the spherically symmetric case, obtaining an inhomogeneous (but spherically symmetric) spacetime which satisfies all the standard energy conditions, yet has a future c-boundary which is a single point. In Section 5, we join this modified version of the FRW event horizonless metric to a general Tolman-Bondi closed universe. In Section 6, we discuss various definitions for a black hole in a closed universe, and show that the Tolman-Bondi universe parameters of the metric in section 5 can be chosen so that by any of these definitions, the expanding phase of the universe has a black hole. In Section 7, we show that the recent supernova observations which strongly suggest that the universe is currently accelerating are consistent with a universe which recollapses to a final singularity before any black hole has time to completely evaporate, provided the acceleration is due to quintessence with certain specified properties. Finally in our concluding Section 8, we shall point out how our “no event horizon” solution to the black hole information paradox naturally complements the “holographic principle” resolution, which assumes that all information in a black hole interior is coded also on its surface.
## 2 A 3-sphere FRW Universe with Final Singularity But No Event Horizons
The Friedmann equation for an $`S^3`$ closed universe is
$$\left(\frac{1}{R}\frac{dR}{dt}\right)^2=\frac{8\pi GM}{3R^{3(1+w)}}\frac{1}{R^2}$$
(2.1)
where the pressure $`p=w\rho `$, with $`w=\gamma 1`$, where $`\gamma `$ is the adiabatic index and $`\rho `$ the mass density. If $`w=1/3`$, then
$$R(t)=\sqrt{(8\pi GM/3)1}(t_ft)$$
(2.2)
is a solution to (2.1) for $`t<t_f`$ with a final singularity at $`t=t_f`$, provided $`(8\pi GM/3)>1`$. The second order equation for the Friedmann universe,
$$\frac{1}{R}\frac{d^2R}{dt^2}=\frac{4\pi G}{3}\left(\rho +3p\right)$$
(2.3)
is automatically satisfied for $`p=(1/3)\rho `$ and $`d^2R/dt^2=0`$.
The closed FRW universe with the scale factor (2.2), namely
$$ds^2=dt^2+R_0^2(t_ft)^2[d\chi ^2+\mathrm{sin}^2\chi (d\theta ^2+\mathrm{sin}^2\theta d\phi ^2)],$$
(2.4)
has no event horizons; that is, its future c-boundary consists of a single point — the Omega Point. Indeed, the equation for future directed null geodesics, $`ds^2=0`$ can be integrated for radial null geodesics to give
$$\mathrm{\Delta }\chi =^{t_f}\frac{dt}{R(t)}=+\mathrm{}$$
(2.5)
which shows that radial null geodesics circumnavigate the universe an infinite number of times as the future c-boundary at $`t=t_f`$ is approached. By homogeneity and isotropy, we can transpose the coordinate system so that any spatial location $`(\chi ,\theta ,\varphi )`$ can reach any other location $`(\chi ^{},\theta ^{},\varphi ^{})`$ via a radial null geodesic segment, and (2.5) shows such radial geodesics can be exchanged an infinite number of times. Hence all future endless timelike curves define the same c-boundary point: the future c-boundary is a single point.
A perfect fluid with $`w=1/3`$ satisfies the weak, the strong, and the dominant energy conditions, since Hawking and Ellis have shown (, pp. 89-95) that for diagonalizable stress energy tensors (Type I matter), the weak energy condition will hold if $`\rho 0`$ and $`\rho +p0`$; the strong energy condition will hold if $`\rho +p0`$ and $`\rho +3p0`$; and the dominant energy condition will hold if $`\rho 0`$ and $`\rho p\rho `$.
It is possible to join the metric with scale factor (2.2) to any closed FRW universe at any time in the collapsing phase. Consider for example the dust ($`w=0`$) scale factor:
$`R(\tau )={\displaystyle \frac{R_{max}}{2}}(1\mathrm{cos}\tau )`$ (2.6)
$`t(\tau )={\displaystyle \frac{R_{max}}{2}}(\tau \mathrm{sin}\tau )`$ (2.7)
where $`0<\tau <2\pi `$ is the conformal time, and $`R_{max}`$ is the radius of the universe at maximum expansion, which occurs when $`\tau =\pi `$. We will make the join at conformal time $`\pi t_{join}<2\pi `$, which by (2.7) gives a proper time of
$$t_{join}t_{join}(\tau _{join})=\frac{R_{max}}{2}(\tau _{join}\mathrm{sin}\tau _{join})$$
The standard junction conditions require continuity of the metric and its first derivatives at the join:
$`R(\tau _{join})={\displaystyle \frac{R_{max}}{2}}(1\mathrm{cos}\tau _{join})=R(t_{join})=Qt_{join}+A`$ (2.8)
$`dR/dt|_{\tau _{join}}=\left({\displaystyle \frac{dR/d\tau }{dt/d\tau }}\right)|_{\tau _{join}}={\displaystyle \frac{\mathrm{sin}\tau _{join}}{1\mathrm{cos}\tau _{join}}}=dR/dt|_{t_{join}}=Q`$ (2.9)
where $`Q\sqrt{8\pi GM/31}`$ and $`At_fQ`$. Solving for $`t_f`$ and $`M`$ yield
$`{\displaystyle \frac{8\pi GM}{3}}={\displaystyle \frac{2}{1\mathrm{cos}\tau _{join}}}`$ (2.10)
$`t_f={\displaystyle \frac{R_{max}}{2}}\left(\tau _{join}+{\displaystyle \frac{2[\mathrm{cos}\tau _{join}1]}{\mathrm{sin}\tau _{join}}}\right)`$ (2.11)
Notice that as $`\tau _{join}\pi ^+`$, we have $`t_f+\mathrm{}`$, which means joining at the time of maximum expansion would yield the Einstein static universe thereafter. As $`\tau _{join}2\pi `$, we have $`t_f\pi R_{max}`$ and $`M+\mathrm{}`$, which means that an arbitrarily large total mass of the $`w=1/3`$ matter is required to eliminate event horizons if the join is made arbitrarily close to the usual proper time end of a dust FRW universe, $`t_f=\pi R_{max}`$.
The standard junction conditions yield a global metric which is $`C^{\mathrm{}}`$, except at the join, where it is $`C^1`$. One can smooth this metric to one which is $`C^{\mathrm{}}`$ everywhere and which satisfies the energy conditions everywhere by allowing $`w`$ to vary smoothly from $`0`$ to $`1/3`$ in a neighborhood $`[t_{join},t_{join}\mathrm{\Delta }t)`$. By the constraint FRW equation $`(R^1dR/dt)^2=R^2+8\pi G\rho /3`$ and the dynamical FRW equation $`2R^1d^2R/dt^2=(R^1dR/dt)^2R^28\pi Gp`$, continuity in $`p`$ and $`\rho `$ would insure that $`R(t)`$ is $`C^2`$, and repeatedly differentiating the dynamical equation would yield that $`R(t)`$ is $`C^{\mathrm{}}`$ (the constraint and dynamical FRW equations imply the conservation equation $`T_{;\nu }^{\mu \nu }=0`$).
Since $`R_\mu ^\mu =8\pi GT_\mu ^\mu =16\pi G\rho =16\pi GM/R^2`$ for the $`w=1/3`$ equation of state, where $`R_\mu ^\mu `$ is the Ricci scalar, the Omega Point singularity at $`t=t_f`$ (at $`R=0`$) is a p.p. curvature singularity (, p. 260). The spacetime is a counter-example to a conjecture by R.K. Sachs that the only Omega Point spacetimes are formed by suitably identifying Minkowski space, which would have locally extendible singularities.
## 3 A $`w=1/3`$ Perfect Fluid Can Be Generated By a Quintessence Scalar Field With an Exponential Potential
We shall now show that a scalar field with exponential potential will generate, at least in a FRW universe, a $`w=1/3`$ perfect fluid behavior near the final singularity. That is, the $`w=1/3`$ perfect fluid behaviour will be seen if the potential for the scalar field $`\varphi `$ is of the form $`V(\varphi )=V_0e^{B\varphi }`$, where $`V_0`$ and $`B`$ are constants. Such a potential is often discussed as a particularly plausible potential for the inflaton field which is thought to be responsible for inflation in the early universe, and as a model of the quintessence field which is responsible for the cosmological acceleration in the present epoch. This will show that a $`w=1/3`$ equation of state is physically plausible near the final singularity of a closed universe, and thus that the absence of event horizons is physically possible.
The stress energy tensor for a scalar field $`\varphi `$ with potential $`V(\varphi )`$ is
$$T_{\alpha \beta }=\left[\varphi _{;\alpha }\varphi _{;\beta }\frac{1}{2}g_{\alpha \beta }\left(\varphi _{;\mu }\varphi _{;\nu }g^{\mu \nu }+2V(\varphi )\right)\right]$$
(3.1)
In the FRW universe, we have $`\varphi =\varphi (t)`$, so $`\varphi _{;i}=0`$ and $`\varphi _{;0}=\varphi _{,0}`$, where the $`i`$ denotes a spatial coordinate, $`0`$ the time coordinate $`t`$, and the semicolon and comma denote the covariant and partial derivatives respectively. In a local orthonormal frame we obtain
$$T_{\widehat{0}\widehat{0}}=\frac{1}{2}(\varphi _{,\widehat{0}})^2+V(\varphi )$$
(3.2)
and
$$T_{\widehat{0}\widehat{0}}+3T_{\widehat{i}\widehat{i}}=2\left[(\varphi _{,\widehat{0}})^2V(\varphi )\right]$$
(3.3)
If $`w=1/3`$, $`T_{\widehat{0}\widehat{0}}+3T_{\widehat{i}\widehat{i}}=0`$, which means
$$V(\varphi )=(\varphi _{,\widehat{0}})^2=(\varphi _{,0})^2$$
(3.4)
where we have used $`\varphi _{,\widehat{0}}=\varphi _{,0}=d\varphi /dt`$. Thus
$$8\pi GT_{\widehat{0}\widehat{0}}=12\pi G(\varphi _{,0})^2=G_{\widehat{0}\widehat{0}}=\frac{3((R_{,0})^2+1)}{R^2}=\frac{3(R_0^2+1)}{R_0^2(t_ft)^2}$$
(3.5)
Taking the square root gives
$$\frac{d\varphi }{dt}=\frac{\sqrt{(R_0^2+1)/4\pi G}}{R_0(t_ft)}$$
(3.6)
which can be immediately integrated to yield
$$\varphi _0\varphi =\sqrt{(1/4\pi G)(1+1/R_0^2)}\mathrm{ln}(t_ft)$$
(3.7)
where $`\varphi _0`$ is a constant. Equation (3.7) can be written
$$(t_ft)^1=\mathrm{exp}\left[(\varphi \varphi _0)/\sqrt{(1/4\pi G)(1+1/R_0^2)}\right]$$
(3.8)
We thus obtain for the potential
$$V(\varphi )=(\varphi _{,0})^2=\frac{(R_0^2+1)}{4\pi GR_0^2}\left[\frac{1}{(t_ft)^2}\right]=V_0e^{B\varphi }$$
(3.9)
where
$$B=\sqrt{\frac{16\pi GR_0^2}{R_0^2+1}}$$
(3.10)
and
$$V_0=\frac{(R_0^2+1)}{4\pi GR_0^2}e^{B\varphi _0}$$
(3.11)
It was pointed out in the previous Section that the join between the dust (or radiation) dominated FRW part of the universe and the $`w=1/3`$ portion can be made at any time and at any radius. If the constants $`V_0`$ and $`B`$ are fixed by the laws of physics, then as the above relation between these constants and the constants $`R_0`$ and $`\varphi _0`$ indicate, the physical laws would also restrict the radius of the join, and the value of the scalar field at the join.
It is interesting to confirm that the potential (3.9) satisfies the second order equation of motion for a scalar field in the FRW universe with $`R(t)=R_0(t_ft)`$. The equation of motion with arbitrary scalar potential $`V(\varphi )`$ is (, p. 466; , p. 431):
$$\varphi _{;\alpha }^{;\alpha }=\frac{V(\varphi )}{\varphi }$$
(3.12)
In the FRW universe we have $`\varphi _{;\alpha }^{;\alpha }=(\varphi _{,0})^{;0}=(\varphi ^{,0})_{;0}=(\varphi _{,0})_{;0}`$, and in a coordinate basis, the identity $`A_{;\alpha }^\alpha =(1/\sqrt{g})(\sqrt{g}A^\alpha )_{,\alpha }`$, for any vector field $`A^\alpha `$, applies. Thus in a FRW coordinate basis, the scalar field equation of motion can be written
$$\frac{1}{R^3}\left(R^3(\varphi _{,0})\right)_{,0}=\frac{V(\varphi )}{\varphi }$$
(3.13)
which can be reduced to the standard expression $`\ddot{\varphi }+3H\dot{\varphi }+V^{}(\varphi )=0`$ as follows. In both a coordinate basis, and in an orthonormal basis, we have $`A_{,0}=dA/dt`$, for any function $`A`$. Thus, for an exponential potential, we have $`V(\varphi )/\varphi =BV(\varphi )=B(\varphi _{,0})^2`$. Using $`R=R_0(t_ft)`$ and the expression (3.6) for $`\varphi _{,0}`$, it is confirmed that (3.13) is indeed an identity.
An alternative derivation of the fact that the $`w=1/3`$ equation of state can be generated by a scalar field with exponential potential would be to make use of Barrow’s work (, ) on scalar fields with exponential potentials in flat space (FRW k = 0). Barrow in fact noted that in the far future, an exponential potential could give rise to the $`w=1/3`$ equation of state in the k = +1 case, but he did not attempt to derive the constants ($`B`$ and $`V_0`$ above) that would allow the $`w=1/3`$ equation of state to be joined to a dust equation of state for earlier times, which is why we did the calculation above. In addition, Vilenkin (see also ) has pointed out that a $`w=1/3`$ equation of state can be generated by a tangled network of very light cosmic strings.
In joining two metrics with different equations of state, one effectively assumes that one form of matter disappears and is replaced by the other. More realistically, if a scalar field were to be present near a final singularity, we would expect it to be in addition to dust or radiation already present. In such a situation, a pure exponential potential uncoupled to the other forms of matter would not give rise to a single c-boundary point, if its stress-energy tensor increased as $`R^2`$, since dust and radiation would increase as $`R^3`$ and $`R^4`$ respectively; such a universe would inevitably become radiation dominated sufficiently near the singularity. However, in a FRW universe, we can always find, for any assumed mixture of dust and radiation, a suitable potential $`V(\varphi )`$ which would have the effect of cancelling out the gravitational force of the dust and matter fields, leaving an effective pure exponential scalar field (this is in effect what happens after the join between the $`w=1/3`$ equation of state and the $`w=0`$ equation of state fluids).
But the actual universe is not expected to be FRW near the final singularity. Even if the universe were FRW in the beginning, we would expect it to become curvature dominated near the final singularity, since the “effective energy density” curvature perturbations around FRW grow as $`R^6`$, much faster than the densities of dust or radiation (, p. 807). So in the actual universe, the elimination of event horizons would have to be carried out by the global collective interactions (of known forces) which give rise to the Misner mixmaster horizon elimination mechanism, as described in .
On the other hand, a pure scalar inflaton (quintessence field) with exponential potential might be expected to be the entire matter content in the very early universe, and the initial singularity might be expected to be FRW. In such a case the effect of such an inflaton field would be to eliminate the particle horizons. In other words, with an exponential inflaton (quintessence) field, the horizon problem of cosmology would be automatically resolved.
## 4 Generalizing the FRW $`w=1/3`$ Omega Point Spacetime to the Spherically Symmetric Case
The approach used in Section 2 for creating spacetimes with no event horizons can be generalized to yield a wider class of such spacetimes. Instead of using the metric (2.4), we introduce functions $`N(\chi )`$ and $`Z(\chi )`$, where $`N`$ is positive on $`[0,\pi ]`$ and $`Z`$ is positive on $`(0,\pi )`$, vanishing at $`0`$ and $`\pi `$. The metric we then use is
$$ds^2=dt^2+(t_ft)^2[N^2d\chi ^2+Z^2(d\theta ^2+\mathrm{sin}^2\theta d\phi ^2)]$$
(4.1)
###### Proposition 1
A Tolman-Bondi spacetime with metric (4.1) has a c-boundary which is a single point.
Proof. To check that this spacetime actually has no event horizons, we mimic the calculation of the same proposition for the $`w=1/3`$ FRW universe in section 2. Let $`N_{max}`$ be the maximum value of $`N`$ on $`[0,\pi ]`$. Then
$$\mathrm{\Delta }\chi =^{t_f}\frac{dt}{N(\chi )(t_ft)}^{t_f}\frac{dt}{N_{max}(t_ft)}=+\mathrm{}.$$
(4.2)
Thus in this class of spacetimes, radial null geodesics are capable of hitting every value of $`\chi `$ an infinite number of times. In order to conclude that every point in space can communicate with every other point, however, we must refine the argument given in Section 2 a bit, for we no longer have the symmetry of the 3-sphere to exploit. We do, however, still have (2-)spherical symmetry. Therefore we can say that a null geodesic may be sent from the origin to any $`(\chi ,\theta ,\phi )`$, and vice versa. Hence, given points $`P_1=(\chi _1,\theta _1,\phi _1)`$ and $`P_2=(\chi _2,\theta _2,\phi _2)`$ which desire to communicate with one another, there exists a piecewise $`C^{\mathrm{}}`$ null curve from $`P_1`$ to $`P_2`$, consisting of a null curve from $`P_1`$ to the origin and then a null curve from the origin to $`P_2`$. Applying an elementary result of Penrose (, Lemma 2.16), we conclude that there exists a timelike or null curve from $`P_1`$ to $`P_2`$, which is precisely what we wanted. QED.
We would like the spacetime (4.1) to satisfy the weak, dominant, and strong energy conditions . Let $`G`$ be the Einstein tensor of this spacetime. Using the equations for the nonzero components of the Einstein tensor in , we can compute $`G`$ in the orthonormal basis $`\omega ^{\widehat{ı}}`$, where:
$$\omega ^{\widehat{0}}=dt,\omega ^{\widehat{1}}=N(t_ft)d\chi ,\omega ^{\widehat{2}}=Z(t_ft)d\theta ,\omega ^{\widehat{3}}=Z(t_ft)\mathrm{sin}\theta d\phi .$$
In this basis, all off-diagonal terms of $`G`$ are zero. Thus all matter is Type I , and the energy conditions will hold if the following six conditions are satisfied:
$`G^{\widehat{0}\widehat{0}}0`$ (4.3)
$`G^{\widehat{0}\widehat{0}}+G^{\widehat{1}\widehat{1}}+G^{\widehat{2}\widehat{2}}+G^{\widehat{3}\widehat{3}}0`$ (4.4)
$`G^{\widehat{0}\widehat{0}}+G^{\widehat{1}\widehat{1}}0`$ (4.5)
$`G^{\widehat{0}\widehat{0}}G^{\widehat{1}\widehat{1}}0`$ (4.6)
$`G^{\widehat{0}\widehat{0}}+G^{\widehat{2}\widehat{2}}=G^{\widehat{0}\widehat{0}}+G^{\widehat{3}\widehat{3}}0`$ (4.7)
$`G^{\widehat{0}\widehat{0}}G^{\widehat{2}\widehat{2}}=G^{\widehat{0}\widehat{0}}G^{\widehat{3}\widehat{3}}0`$ (4.8)
The strong energy condition is (4.3) and (4.4), the weak is (4.3), (4.5), and (4.7), and the dominant is (4.3) and (4.5)-(4.8). Computing these expressions with our metric (4.1), we see that $`G^{00}+G^{11}+G^{22}+G^{33}=0`$ identically, and (4.3), (4.5), (4.6), (4.7), and (4.8) are equivalent to, respectively:
$`3+{\displaystyle \frac{1}{Z^2}}{\displaystyle \frac{1}{N^2}}\left({\displaystyle \frac{2Z^{\prime \prime }}{Z}}{\displaystyle \frac{2Z^{}N^{}}{ZN}}+{\displaystyle \frac{(Z^{})^2}{Z^2}}\right)0`$ (4.9)
$`2{\displaystyle \frac{2}{N^2}}\left({\displaystyle \frac{Z^{\prime \prime }}{Z}}{\displaystyle \frac{Z^{}N^{}}{ZN}}\right)0`$ (4.10)
$`4+{\displaystyle \frac{2}{Z^2}}{\displaystyle \frac{2}{N^2}}\left({\displaystyle \frac{Z^{\prime \prime }}{Z}}{\displaystyle \frac{Z^{}N^{}}{ZN}}+{\displaystyle \frac{(Z^{})^2}{Z^2}}\right)0`$ (4.11)
$`2+{\displaystyle \frac{1}{Z^2}}{\displaystyle \frac{1}{N^2}}\left({\displaystyle \frac{Z^{\prime \prime }}{Z}}{\displaystyle \frac{Z^{}N^{}}{ZN}}+{\displaystyle \frac{(Z^{})^2}{Z^2}}\right)0`$ (4.12)
$`4+{\displaystyle \frac{1}{Z^2}}{\displaystyle \frac{1}{N^2}}\left({\displaystyle \frac{3Z^{\prime \prime }}{Z}}{\displaystyle \frac{3Z^{}N^{}}{ZN}}+{\displaystyle \frac{(Z^{})^2}{Z^2}}\right)0`$ (4.13)
where prime () denotes differentiation with respect to $`\chi `$ (everything is a function of $`\chi `$). In the FRW case, $`N=R_0`$, $`Z=R_0\mathrm{sin}\chi `$, and the energy condition equations are all are equivalent to:
$$1+\frac{1}{R_0^2}0.$$
In other words, in the FRW case, the energy conditions are always satisfied.
We shall now show that if the metric (4.1) defines a universe that is “sufficiently large,” it will automatically satisfy the energy conditions. Suppose we are given functions $`N_0(\chi )`$ and $`Z_0(\chi )`$ such that there exist constants $`R_1,R_2,ϵ_1,ϵ_2>0`$ with $`N_0(\chi )=R_1`$ and $`Z_0(\chi )=R_1\mathrm{sin}\chi `$ for $`0\chi <ϵ_1`$ and $`N_0(\chi )=R_2`$ and $`Z_0(\chi )=R_2\mathrm{sin}\chi `$ for $`\pi ϵ_2<\chi \pi `$. In other words, $`N_0`$ and $`Z_0`$ look like the $`N`$ and $`Z`$ from FRW universes near $`\chi =0`$ and $`\chi =\pi `$. Then we know that near $`\chi =0`$ and $`\chi =\pi `$, the energy conditions are satisfied for $`N=RN_0`$ and $`Z=RZ_0`$, where $`R`$ is an arbitrary positive constant. Then since the expressions multiplied by $`\frac{1}{N^2}`$ in the energy conditions are bounded for $`\chi [ϵ_1,\pi ϵ_2]`$, we may find a constant multiplier $`R`$ such that the metric (4.1) with $`N=RN_0`$ and $`Z=RZ_0`$ satisfies all the energy conditions everywhere. The current observational evidence indicates that the universe is very close to being spatially flat, so the actual universe satisfies the “sufficiently large” criterion.
## 5 A 3-Sphere Universe Containing a Black Hole But Having No Event Horizons
We will produce a large class of $`S^3\times R^1`$ spacetimes which are in their expanding phase, special cases of the general spherically symmetric dust solution and which are eventually joined to a spacetime of the type described in section 4, so that they end with a c-boundary of a point (and hence have no event horizons), and satisfy the energy conditions everywhere.
### 5.1 General Dust Solution
The general spherically symmetric pressureless dust solution is:
$$ds^2=dt^2+(1f^2)^1\left(\frac{Y}{\chi }\right)_t^2d\chi ^2+Y^2(d\theta ^2+\mathrm{sin}^2\theta d\phi ^2)$$
(5.1)
where the notation $`(\frac{Y}{\chi })_t`$ denotes differentiation of $`Y`$ with respect to $`\chi `$ where the independent variables are $`t`$, $`\chi `$, $`\theta `$, and $`\phi `$ (subscripts to differentials in general will specify independent variables, with the assumption that $`\theta `$ and $`\phi `$ are always independent), $`f`$ is an arbitrary function of $`\chi `$ alone taking values in $`[0,1]`$ and $`Y`$ and $`t`$ are given by:
$`t=t_0(\chi )+(\eta \mathrm{sin}\eta ){\displaystyle \frac{m(\chi )}{f(\chi )^3}}`$ (5.2)
$`Y=(1\mathrm{cos}\eta ){\displaystyle \frac{m(\chi )}{f(\chi )^2}}.`$ (5.3)
In the above expressions, $`t_0`$ is an arbitrary function of $`\chi `$ alone, $`m`$ is another arbitrary function of $`\chi `$ positive on $`(0,\pi )`$, and $`\eta `$ is defined by (5.2). The only restrictions on these free functions are that to maintain the nondegeneracy of the metric in a closed universe, $`f`$ should equal 1 at one $`\chi `$-value in the interior of $`[0,\pi ]`$, at which point $`m^{}`$, $`f^{}`$ and $`t_0^{}`$ should all be zero. The general dust metric becomes degenerate whenever $`Y^{}=0`$ and $`f1`$ or $`Y^{}0`$ and $`f=1`$. Such 2-spheres of degeneracy correspond to shell-crossing singularities, and if these degeneracy spheres occur before the final singularity at $`\eta =2\pi `$, they will give rise to a breakdown in global hyperbolicity, as is well-known. We shall assume that the free functions $`m,f,t_0`$ are so chosen that this does not occur.
The dust case of the FRW metric (the case where $`w=0`$) is a special case of this general metric. Letting
$$f=\mathrm{sin}\chi ,m=\frac{R_{max}}{2}\mathrm{sin}^3\chi ,\mathrm{and}t_0=0,$$
one obtains
$$t=\frac{R_{max}}{2}(\eta \mathrm{sin}\eta ),Y=\frac{R_{max}}{2}(1\mathrm{cos}\eta )\mathrm{sin}\chi ,$$
$$\mathrm{and}\left(\frac{Y}{\chi }\right)_t=\frac{R_{max}}{2}(1\mathrm{cos}\eta ).$$
The resulting metric is
$$ds^2=dt^2+\left[\frac{R_{max}}{2}(1\mathrm{cos}\eta )\right]^2[d\chi ^2+\mathrm{sin}^2\chi (d\theta ^2+\mathrm{sin}^2\theta d\phi ^2],$$
precisely the Friedmann collapsing dust $`S^3`$ solution.
### 5.2 The Join
We have shown above that (2.4) can be joined in a $`C^1`$ manner to any collapsing dust FRW $`S^3`$ universe at any time in the collapsing phase by a suitable choice of the constants $`R_0`$ and $`t_f`$. We shall now generalize this construction substantially, joining a certain class of Tolman-Bondi pressureless dust solutions (including the FRW $`S^3`$ collapsing dust solution) to universes of the sort (4.1), so that we produce a large class of universes which start with pressureless dust and in the end have no event horizons.
In making this join, we will allow the hypersurface $`𝒥`$ along which the two metrics are joined to vary as a free function. For convenience, we will take $`𝒥`$ to be spherically symmetric, parametrized as $`(t_𝒥,\chi ,\theta ,\phi )`$ in the dust universe, where $`\eta _𝒥=\eta _𝒥(\chi )`$ is a free function of $`\chi `$, and
$$t_𝒥t(\eta _𝒥(\chi ),\chi )=t_0(\chi )+(\eta _𝒥\mathrm{sin}\eta _𝒥)m(\chi )/f^3(\chi )$$
(5.4)
We will also assume that the $`t`$, $`\chi `$, $`\theta `$, and $`\phi `$ coordinates agree across $`𝒥`$. Therefore, we will have six degrees of freedom altogether: $`t_0`$, $`m`$, $`f`$, $`\eta _𝒥`$, $`N`$, and $`Z`$, all free functions of $`\chi `$.
To make the join $`C^1`$, we first must make it continuous across $`𝒥`$. This in particular means that the metric coefficients will agree along $`𝒥`$ itself. Therefore the metric coefficients will agree in all derivatives along vectors tangent to $`𝒥`$, and thus in order to check that the join is $`C^1`$, one must only check that the first derivatives of the metric coefficients agree in a direction independent to the tangent spaces of $`𝒥`$. The direction we choose is $`(/t)_\chi `$. This direction is linearly independent to the tangent spaces of $`𝒥`$ because $`(/\eta )_\chi `$ is never tangent to $`𝒥`$ (since $`𝒥`$ is parametrized with $`\eta `$ a function of $`\chi `$), and
$$\left(\frac{}{t}\right)_\chi =\frac{f^3}{m(1\mathrm{cos}\eta )}\left(\frac{}{\eta }\right)_\chi .$$
Since we are assuming that the coordinates are the same on the dust universe as on the universe (4.1), the off-diagonal coefficients already agree (they are 0 on both sides of $`𝒥`$), and $`g_{tt}`$ is -1 on both sides of $`𝒥`$. Furthermore, since we have spherical symmetry on both sides of $`𝒥`$, we need check only one of $`g_{\theta \theta }`$ and $`g_{\phi \phi }`$. We are left with four junction conditions:
$`{\displaystyle \frac{(Y^{})^2}{1f^2}}=N^2(t_ft_𝒥)^2`$ $`(g_{\chi \chi }\mathrm{continuous})`$ (5.5)
$`Y^2=Z^2(t_ft_𝒥)^2`$ $`(g_{\theta \theta }\mathrm{continuous})`$ (5.6)
$`{\displaystyle \frac{2Y^{}\dot{Y}^{}}{1f^2}}=2N^2(t_ft_𝒥)`$ $`(g_{\chi \chi }C^1)`$ (5.7)
$`2Y\dot{Y}=2Z^2(t_ft_𝒥)`$ $`(g_{\theta \theta }C^1).`$ (5.8)
Here the dot ($`\dot{}`$) denotes application of $`(/t)_\chi `$ and the prime () denotes application of $`(/\chi )_t`$. These tangent vectors arise from the coordinate system $`(t,\chi ,\theta ,\phi )`$, and thus they commute with one another, so that (5.7) makes sense.
Therefore we have six free functions — $`t_0(\chi ),m(\chi ),f(\chi )`$ on the Tolman-Bondi side of the join, and $`\eta _𝒥(\chi ),N(\chi ),Z(\chi )`$ on the final singularity side of the join — and four differential equations relating them. One would expect that these four equations would determine four of the functions in terms of the other two, and this is exactly what happens. We find it convenient to choose $`m(\chi )`$ and $`f(\chi )`$ as the arbitrary functions, and expressing all other functions in terms of these two. After some manipulation of the above equations we obtain
$`{\displaystyle \frac{(1+\mathrm{cos}\eta _𝒥)^2}{|\mathrm{sin}^3\eta _𝒥|}}=C{\displaystyle \frac{m}{f^3}}`$ (5.9)
$`t_0=t_f+\left({\displaystyle \frac{2(1\mathrm{cos}\eta _𝒥)}{\mathrm{sin}\eta _𝒥}}\eta _𝒥\right){\displaystyle \frac{m}{f^3}}`$ (5.10)
$`N={\displaystyle \frac{|\dot{Y}^{}(\eta _𝒥,\chi )|}{\sqrt{1f^2}}}`$ (5.11)
$`Z=|\dot{Y}(\eta _𝒥,\chi )|.`$ (5.12)
where $`C`$ is a constant of integration. Equation (5.9) can be inverted to give $`\eta _𝒥(\chi )`$. Then $`\eta _𝒥`$ is inserted into equations (5.10), (5.11), and (5.12)to yield $`t_0(\chi )`$, $`N(\chi )`$, and $`Z(\chi )`$ respectively in terms of the arbitrary functions$`m(\chi )`$ and $`f(\chi )`$, and the constants $`t_f`$ and $`C`$ . Notice that the allowed Tolman-Bondi dust metrics are no longer completely general, since the function $`t_0`$ is now fixed rather than being completely arbitrary. The constants $`t_f`$ and $`C`$ allow the join to be made as far in the future in proper time (the constant $`t_f`$) and in $`\eta `$ time (the constant $`C`$) as one wishes. To see the latter, note that equation (5.9) is of the form $`F(\eta _𝒥)=Cm/f^3`$, where $`F`$ increases monotonically from 0 to $`+\mathrm{}`$ as $`\eta _𝒥`$ ranges from $`\pi `$ to $`2\pi `$. Thus if we make $`C`$ arbitrarily large, $`\eta _𝒥`$ can be made arbitrarily close to $`2\pi `$, i.e. as close as we wish to the final singularity. Furthermore, one observes that the boundary requirement of $`t_0^{}`$ (that it is 0 whenever $`Y^{}`$ and hence $`m^{}`$ and $`f^{}`$ are 0) is automatically satisfied, since by the junction conditions
$$t_0^{}=\frac{m}{f^3}\left(\frac{m^{}}{m}3\frac{f^{}}{f}\right)\left(\frac{(1\mathrm{cos}\eta _𝒥)\mathrm{sin}\eta _𝒥}{2+\mathrm{cos}\eta _𝒥}\eta _𝒥+\mathrm{sin}\eta _𝒥\right).$$
The join in Section 2 is in fact a special case of this construction. Consider the FRW choices for $`m`$ and $`f`$:
$$m(\chi )=\frac{R_{max}}{2}\mathrm{sin}^3\chi ,\mathrm{and}f(\chi )=\mathrm{sin}\chi .$$
Then $`m/f^3=1`$ identically, so that by (5.9) $`\eta _𝒥`$ will be a constant, and then by (5.10) $`t_0`$ will be constant. Choosing $`t_f`$ appropriately, we may make $`t_0`$ identically zero, so that our Tolman-Bondi universe is in fact the FRW collapsing dust universe. As noted above, we may make the join at any time in the collapsing phase, just as in the FRW construction of Section 2, and a simple calculation reveals that the $`N`$ and $`Z`$ forced by the join are precisely those which give the $`w=1/3`$ universe.
### 5.3 Join With a Possible “Weak” Shell-Crossing Singularity
In order to make the metric coefficients differentiable across $`𝒥`$, we had to impose four conditions on six functions. We would, however, like to join a completely arbitrary Tolman-Bondi metric to the metric (4.1), and this will require eliminating one of the equations. The junction condition which on physical grounds is the least important is (5.7), the requirement that $`g_{\chi \chi }`$ be $`C^1`$. If $`g_{\chi \chi }`$ is not $`C^1`$ at $`𝒥`$, then the curvature will be a $`\delta `$ function on $`𝒥`$, but this $`\delta `$ function will correspond to a shell-crossing singularity, a singularity that is generally agreed to be unphysical. (Notice also that requiring $`g_{\chi \chi }`$ be $`C^1`$ across $`𝒥`$ actually requires that the radii $`Y`$ of the constant $`\chi `$ spheres have one of its second derivatives, $`\dot{Y}^{}`$, be continuous across $`𝒥`$.) So we drop the junction condition (5.7). A little manipulation yields
$`\eta _𝒥2\mathrm{tan}{\displaystyle \frac{\eta _𝒥}{2}}={\displaystyle \frac{f^3(t_ft_0)}{m}}`$ (5.13)
$`t_𝒥=t_0(\chi )+(\eta _𝒥\mathrm{sin}\eta _𝒥){\displaystyle \frac{m(\chi )}{f^3(\chi )}}`$ (5.14)
$`N(\chi )={\displaystyle \frac{|Y^{}(\eta _𝒥,\chi )|}{(t_ft_𝒥)\sqrt{1f^2}}}`$ (5.15)
$`Z(\chi )=|\dot{Y}(\eta _𝒥,\chi )|`$ (5.16)
We proceed as in the previous section, solving first for $`\eta _𝒥`$ and then substituting this into the other equations. Note that in order to invert equation (5.13), solving for $`\eta _𝒥`$, the fact that $`t_0`$ must be less than $`t_f`$ (the final $`t`$-value of the universe) implies that the LHS should be positive for all values of $`\chi `$. But then since the positive values of $`u2\mathrm{tan}(u/2)`$ for $`u[0,\mathrm{\hspace{0.17em}2}\pi ]`$ are all at least $`2\pi `$, we must therefore have the RHS being at least $`2\pi `$ for all values of $`\chi `$. Since $`m`$, $`f`$, and $`t_0`$ are all defined on the same compact interval, we may (and must) choose the constant $`t_f`$ so large that the RHS is always at least $`2\pi `$. Notice that the differentiability of $`N`$ in equation 5.15 will be guaranteed when $`Y^{}=0`$ and $`f=1`$ because the value of $`(Y^{})/(1f^2)`$ as $`\chi `$ approaches such a point is well defined and positive, and thus $`(t_ft_0)`$ will be just the square root of this positive value. This means that the more standard type of shell-crossing singularity ($`Y^{}=0`$ but $`f1`$, so that $`g_{\chi \chi }=0`$) is assumed not to occur on $`𝒥`$. For this reason, we called the allowed singularity a “weak” shell-crossing singularity.
Therefore, if we make the appropriate choice of $`t_f`$ as described above, we may join an arbitrary Tolman-Bondi closed dust metric to the metric (4.1) provided we allow for a possible shell-crossing singularity. The only additional restriction we must impose on the Tolman-Bondi functions $`m`$, $`f`$, and $`t_0`$ is that they must be chosen to make the universe “sufficiently large” as discussed in Section 4.
## 6 Black Holes
One interesting consequence of the above constructions is that they provide examples of spacetimes satisfying the energy conditions which can contain black holes, but do not contain event horizons. In order for this statement to make sense, however, we need a good definition of a black hole in a closed universe, for in a closed universe the black hole singularity is actually just a component of the final singularity (cf. ). We will discuss in detail three such definitions, the first due to Hayward , the second due to Tipler , and the third due to Wheeler .
In the standard definition of a black hole (cf. Wald (, p. 300, or , p. 924), the black hole $`B`$ is the spacetime region $`BMJ^{}(^+)`$, where $`M`$ is the spacetime manifold, $`^+`$ is “scri plus” — future null infinity, and $`J^{}(S)`$ is the causal past of a set $`S`$, which is to say that $`J^{}(S)`$ is the set of all spacetime points $`p`$ which can be reached by a past-directed timelike or null curve from $`S`$ to $`p`$. (Discussions of global general relativity and the definitions of concepts used in this discipline can be found in Wald (, chapter 8; Misner, Thorne and Wheeler , Chapter 34; or () and ().) This definition cannot be applied in a closed universe, because $`^+`$ does not exist in a closed universe with a final singularity. However, this standard definition of a black hole is never used in practice. When astrophysicists search for black holes, they look for gravitational fields implying the presence of trapped or marginally trapped surfaces. In asymptotically flat spacetimes, (1) all trapped surfaces can be proven to be inside of a black hole (in the standard definition), and (2) black holes are expected to evolve rapidly to a Schwarzschild or Kerr black hole, in which there are trapped surfaces arbitrarily close to the boundary of the black hole $`J^{}(^+)`$ — the event horizon. Now trapped surfaces can be in closed universes.
### 6.1 Trapped and Marginally Trapped Surfaces
Thus, the fundamental concept in Hayward’s and Tipler’s definitions of a black hole is that of a trapped surface (cf. Hayward ). Let $`𝒮`$ be a compact spacelike 2-surface embedded in our spacetime manifold. Let $`P`$ be a point of $`𝒮`$. There are precisely two null directions normal to $`𝒮`$ at $`P`$. Suppose furthermore that these null directions can be expressed as two vector fields $`N_+`$ and $`N_{}`$ defined on all of $`𝒮`$. We can choose both of $`N_+`$ and $`N_{}`$ to be future-directed. Now, allowing $`𝒮`$ to evolve along $`N_+`$ and $`N_{}`$, we can measure its area at every instant, and logarithmically differentiate the resulting function with respect to the evolution parameter. Call these quantities $`\theta _+`$ and $`\theta _{}`$, respectively.
###### Definition 1
$`𝒮`$ is called a (future) trapped surface if both $`\theta _+<0`$ and $`\theta _{}<0`$. If one of these quantities is zero and the other is negative, then $`𝒮`$ is called a marginally trapped surface.
The intuition here is that light rays emitted from a trapped surface will converge, no matter whether they are sent “outward” or “inward.” This is certainly a necessary property of a black hole, and it would be sufficient if it weren’t for the fact that the cosmological singularity produces a wealth of trapped surfaces as the universe collapses. In a FRW closed universe, for example, there is a trapped surface passing through every spatial point in the collapsing phase. In order to distinguish these cosmological trapped surfaces from non-cosmological ones — black hole type trapped surfaces — we need additional criteria.
### 6.2 Hayward’s Black Hole Definition
Marginally trapped surfaces are important because in some sense they are where the horizon of a prospective black hole should be. To distinguish trapped surfaces arising from black holes from those arising from the collapse of the universe, therefore, Hayward considers what should be happening in a neighborhood of a marginally trapped surface which arises because of a black hole. Without loss of generality let $`\theta _+=0`$, $`\theta _{}<0`$ along the marginally trapped surface $`𝒮`$. Then $`N_+`$ can sensibly be called “outward,” $`N_{}`$ “inward.” Outward-directed light rays run instantaneously parallel to the surface, and inward-directed light rays converge. In the case of black hole-based marginally trapped surfaces, however, we would like to say that outward light rays just outside $`𝒮`$ diverge, while outward light rays just inside $`𝒮`$ converge. This can be accomplished mathematically by extending the embedding of $`𝒮`$ to a “double-null foliation” (cf. Hayward ) in the direction of $`N_+`$ and $`N_{}`$, extending $`\theta _+`$ appropriately, and computing the sign of $`_{}\theta _+`$, where $`_{}`$ denotes the Lie derivative in the direction of $`N_{}`$.
###### Definition 2
A marginally trapped surface with $`\theta _+=0`$ (resp. $`\theta _{}=0`$) is called inner if $`_{}\theta _+`$ (resp. $`_+\theta _{}`$) is positive, outer if $`_{}\theta _+`$ (resp. $`_+\theta _{}`$) is negative, and degenerate otherwise.
As makes sense, inner marginally trapped surfaces correspond to cosmological collapse, and outer marginally trapped surfaces correspond to non-cosmological collapse, i.e., to the marginally trapped surfaces we would expect to find inside black holes.
As mentioned above, in asymptotically flat spacetimes, all the trapped surfaces are inside the black hole, and furthermore, all future directed causal (timelike or null) curves from any trapped surface $`𝒯_i`$ can also be shown to be inside the black hole. Thus if $`B`$ is the black hole region, we must have $`J^+(_iT_i)B`$ in order to capture the astrophysically defining feature of a black hole in a black hole definition applicable to a closed universe. Also, in asymptotically flat spacetimes, any spacetime point $`p`$ whose causal future eventually enters the causal future of a trapped surface can be proven to be inside a black hole. This means that we should also include in the black hole $`B`$ all points $`p`$ such that $`J^+(J^+(p)J^{}(𝒯_i))J^{}(𝒯_i)`$.
This gives
###### Definition 3
(Hayward) a black hole is the set of all spacetime points $`p`$ such that $`J^+(J^+(p)J^{}(𝒯_i))J^{}(𝒯_i)`$. where $`𝒯_i`$ is the union of all outer marginally trapped surfaces
### 6.3 Tipler’s Black Hole Definition
Tipler’s criterion () is related to Hayward’s, but perhaps is a bit simpler (see Hayward for a short discussion of how these criteria relate).
Instead of using a double-null foliation to test whether a given marginally trapped surface corresponds to the cosmological collapse or to a local black hole, Tipler instead supposes that the marginally trapped surface in question is contained in the boundary of a spacelike hypersurface-with-boundary $`𝒯`$ whose interior $`𝒯𝒯`$ is foliated by trapped surfaces. He then (assuming without loss of generality that $`\theta _+=0`$, $`\theta _{}<0`$) makes the following
###### Definition 4
If the family of null vectors $`N_{}`$ (which are all on $`𝒯`$) point in the direction of $`𝒯`$, then all trapped surfaces which can be obtained from trapped surfaces in $`𝒯`$ by an acausal homotopy foliated by trapped surfaces will be called non-cosmological.
In particular, any trapped surface in $`𝒯`$ is non-cosmological in this case. Thus black hole type trapped surfaces would be non-cosmological trapped surfaces, and we have
###### Definition 5
(Tipler) a black hole is the set of all spacetime points $`p`$ such that $`J^+(J^+(p)J^{}(𝒯_i))J^{}(𝒯_i)`$, where $`𝒯_i`$ is the union of all non-cosmological trapped surfaces
### 6.4 Hayward-Tipler Black Holes in Tolman-Bondi Closed Universes
Specializing to the Tolman Dust case, we already have a convenient foliation by 2-spheres, and we will use this foliation to assist in evaluating the two definitions outlined in the previous section.
First of all, note that the normal bundle to this foliation is spanned by the vector fields $`/t`$ and $`/\chi `$. Therefore we may set
$$N_\pm =\frac{}{t}\pm \frac{\sqrt{1f^2}}{Y^{}}\frac{}{\chi },$$
since $`/t`$ is future-directed and $`g(/t,N_\pm )=1`$. Since the area of the 2-sphere at $`(t,\chi )`$ is $`4\pi Y^2`$, we can then easily compute
$$\theta _\pm =N_\pm (\mathrm{log}(4\pi Y^2))=\frac{2f^3}{(1\mathrm{cos}\eta )m}\left(\mathrm{cot}\frac{\eta }{2}\pm \frac{\sqrt{1f^2}}{f}\right).$$
Note that in order for the 2-sphere at $`(t,\chi )`$ to be marginally trapped, we are forced to have $`\theta _+=0`$ and $`\theta _{}<0`$ (since $`\theta _+>\theta _{}`$).
Let us consider first Tipler’s definition. Choose $`𝒯`$ to be a constant-$`t`$ hypersurface. An example of a vector pointing in the direction of $`𝒯`$ is $`\nu =(\theta _+)^{}(/\chi )`$ (here primes are as in the join conditions above), since $`\theta _+`$ will decrease to become negative on $`𝒯`$, where the 2-spheres will be trapped surfaces. Thus a sufficient condition for cosmological trapped surfaces in a neighborhood of a marginally trapped surface at $`(t,\chi )`$ is
$$0<g(\nu ,N_{})=(\theta _+)^{}\frac{Y^{}}{\sqrt{1f^2}}.$$
Expanding this a bit, letting $`Q=m/f^3`$, ignoring positive multipliers and using the fact that $`\theta _+=0`$, we obtain the condition
$$\frac{Y^{}}{\sqrt{1f^2}}\left[\frac{t_0^{}+(\eta \mathrm{sin}\eta )Q^{}}{4f^4Q}\frac{f^{}}{f^2\sqrt{1f^2}}\right]>0.$$
We now consider how Hayward’s definition applies. We are taking the derivative in the $`N_{}`$ direction of $`\theta _+`$, and determining its sign. To that end, we will have a Hayward black-hole-type marginally trapped surface if
$$0>_{}\theta _+=\frac{1}{Y^{}}\left\{\frac{1}{2f^3Q}\left(\frac{f^{}}{f}\frac{Q^{}}{Q}\right)\frac{\sqrt{1f^2}}{2f^6Q^2}[t_0^{}+(\eta \mathrm{sin}\eta )Q^{}]\right\}.$$
### 6.5 Black Holes in a Joined Universe
Now let us suppose that we are looking for a black hole in a pressureless dust universe which can be joined to the $`N`$-$`Z`$ universe defined above. We first consider the case of the differentiable join. Leaving $`m`$ and $`f`$ free as in our derivation of the join conditions, we compute Tipler’s criterion to be
$$\frac{Y^{}}{\sqrt{1f^2}}\left[\frac{(T+\eta \mathrm{sin}\eta )Q^{}}{4f^4Q}\frac{f^{}}{f^2\sqrt{1f^2}}\right]>0,$$
and Hayward’s to be
$$\frac{1}{2Y^{}f^3Q}\left\{\frac{f^{}}{f}\frac{Q^{}}{Q}\frac{Q^{}(T+\eta \mathrm{sin}\eta )\sqrt{1f^2}}{Qf^3}\right\}<0,$$
where $`T(\chi )`$ is defined
$$T=\frac{3\mathrm{sin}\eta _𝒥}{2+\mathrm{cos}\eta _𝒥}\eta _𝒥.$$
It is clear that we can choose $`m,f,t_f`$ in the $`C^1`$ join in such a way that the resulting joined universe contains black holes in either the Hayward or Tipler sense, and satisfies the energy conditions.
### 6.6 Wheeler’s Black Hole Definition
Following Wheeler and Qadir , we will consider black holes in the spherically symmetric dust universe to be regions in which the universe is collapsing much faster than elsewhere. An intuitive measure of the rate of collapse of the universe in any region is a measurement of the elapsed time between the initial and final singularities.
We can use physical models to approximate the elapsed time between singularities in black hole regions. For simplicity, we will consider a one solar mass black hole. Misner, Thorne, and Wheeler give the time for a particle to fall from radius $`R_i`$ to the singularity in a standard Schwarzschild black hole as $`\pi (R_i^3/8GM)^{1/2}`$. Taking $`R_i`$ to be the Schwarzschild event horizon $`R_i=2GM/c^2`$, we have that the elapsed time from the penetration of the event horizon to the final singularity is approximately $`t_{horizon}=5\times 10^6\mathrm{sec}(M/M_{})`$. A dust cloud generating such a black hole would have a total lifetime of twice this, giving us the elapsed time from initial to final singularity as
$$t_{total}=10^5\mathrm{sec}\left(\frac{M}{M_{}}\right)$$
(6.1)
in the vicinity of a black hole.
Misner, Thorne and Wheeler cite the illustrative time that might be expected to elapse from the beginning to the end of a typical closed FRW universe (without black hole regions)as
$$t_{total}=60\times 10^9\mathrm{years}=1.8\times 10^{17}\mathrm{sec}.$$
Combining this result with the above calculation, we conclude that the elapsed time between initial and final singularities in the vicinity of a 1 $`M_{}`$ black hole will be on the order of $`10^{22}`$ times smaller than that in non-black hole regions of the universe. Since an upper bound to the mass of a black hole in the current epoch of universe history is believed to be $`10^{10}M_{}`$, we would expect that the elapsed time between the initial and final singularities inside the largest black hole in existence today would be on the order of $`10^{12}`$ times smaller than that in non-black hole regions of the universe, since by (6.1), a black hole lifetime scales linearly with its mass.
The elapsed time from the beginning to the end of the universe along any timelike curves of constant $`\chi `$ can be obtained from (5.2) by computing $`t(2\pi ,\chi )t(0,\chi )`$. Thus the total time elapsed from initial to final singularity along a timelike curve of constant $`\chi `$ is simply $`2\pi m(\chi )/f(\chi )^3`$. Therefore we can assert that if $`\chi _a`$ corresponds to a black hole region of some hypersurface and $`\chi _b`$ corresponds to the cosmological region, we should obtain that:
$$\frac{t_{total}(\chi _b)}{t_{total}(\chi _a)}=\frac{m(\chi _b)/f(\chi _b)^3}{m(\chi _a)/f(\chi _a)^3}>10^{13}.$$
(6.2)
To apply this notion of black holes in the dust universe to our joined metric, we consider any 2-sphere with coordinate radius $`\chi _a`$ to be inside a black hole if (6.2) is satisfied when $`\chi _b`$ is the coordinate radius of a 2-sphere whose size evolves like the 2-spheres of spherical symmetry in a FRW universe. It will be sufficient that there exist a pair of 2-spheres with coordinate radii $`\chi _a`$ and $`\chi _b`$ such that (6.2) is satisfied. Alternatively, we could simply restrict attention to $`1M_{}`$ black holes and fix $`\chi _b`$. An elementary calculation shows that a solar mass black hole is a 2-sphere with radius corresponding to the radial value $`\chi =10^{23}`$. A third way of picking the pair of 2-spheres is to require that the $`\chi _a`$ 2-sphere is the “largest” 2-sphere in the universe “today”. This would mean that the $`\chi _a`$ 2-sphere is an extremal — maximal — 2-sphere embedded in the 3-sphere corresponding to “today”. We have to require the 2-sphere to be extremal since one can construct a non-extremal 2-sphere of arbitrary size embedded in a 3-sphere.
Wheeler points out that the natural meaning of “today” — the choice of a spacelike hypersurface though the Earth — is the constant mean curvature hypersurface through the Earth today. Tipler (, p. 440) has shown that if the strong energy condition holds, and if the universe began close to homogeneity and isotropy, an Omega Point spacetime can be uniquely foliated by constant mean curvature hypersurfaces, so Wheeler’s proposal does indeed define a unique “today” over the entire universe (Tipler also shows that a constant mean curvature hypersurface probably coincides with the rest frame of the CBR at any event, so Wheeler’s “today across the entire universe” is even easy to locate experimentally.) Putting all of these criteria together yields
###### Definition 6
(Wheeler) a black hole is the set of all spacetime points $`p`$ such that $`J^+(J^+(p)J^{}(𝒮))J^{}(𝒮)`$, where $`𝒮`$ is any 2-sphere with coordinate radius $`\chi _a`$ for which (6.2) holds when $`\chi _b`$ is the coordinate radius of a maximal 2-sphere in the constant mean curvature hypersurface which includes the 2-sphere with coordinate radius $`\chi _a`$.
It is clear that since (6.2) does not depend on the function $`t_0`$, just on the functions $`m,f`$, it is possible to construct even in the case of the $`C^1`$ join a universe which is essentially a closed dust FRW everywhere outside a small black hole by Wheeler’s definition, a black hole which is centered at the origin of coordinates $`\chi =0`$.
## 7 Quintessence and Recollapse
The “no event horizon solution” to the black hole information problem requires that the universe recollapse to a final singularity before black holes have time to evaporate. However, the best observations , independently confirmed by a number of groups, indicate that the universe is currently accelerating. Furthermore, the observed structure is best explained (given a Hubble constant of $`65\pm 5`$ km/sec-Mpc and spatial flatness) via a $`\mathrm{\Lambda }`$CDM model . If this acceleration were to continue — as it would if it were due to a positive cosmological constant — then the universe would expand forever, and our proposed solution to the BH unitarity problem would be incorrect: unitarity would be violated.
However, Barrow (, see also ) was the first to point out that an accelerating universe today need not preclude a recollapse in the far future of a closed universe. Since the acceleration of the scale factor $`R`$ is given by equation (2.3), namely $`\ddot{R}=(4\pi G/3)R(\rho +3p)=(4\pi G/3)R(1+3w)\rho `$, acceleration today implies that $`w<1/3`$ today (the data give $`w<5/9`$ today at the 95% confidence level ), but if eventually $`w>1/3`$ for all time greater than some far future value $`t_{future}`$, then the recollapse theorems of Barrow, Galloway and Tipler will apply, and recollapse will occur.
Thus unitarity implies that the observed acceleration “today” (meaning most of past proper time) cannot be due to a positive cosmological constant, but must instead be due to quintessence. This is of course the general expectation of cosmologists, since the only plausible non-zero values of the cosmological constant are near the Planck density of $`(10^{19}\mathrm{GeV})^4`$, or near the density of the SM Higgs field at its minimum $`(200\mathrm{GeV})^4`$, whereas the observed density of the material causing the acceleration is of the order of the closure density, $`(10^3\mathrm{eV})^4`$ (, ).
The “standard model” of quintessence (, ) is a scalar field $`\varphi `$ with a very shallow potential $`V(\varphi )`$ in the present epoch, resulting in scalar field excitations of very small mass, $`m_\varphi \sqrt{V^{\prime \prime }(\varphi )/2}H_010^{33}\mathrm{eV}`$. Since we know very little more than this about $`V(\varphi )`$, the potential could have a minimum around which the field will oscillate in the far future. In such a case, in the far future the leading term in the expansion of the potential about the minimum would be $`\frac{1}{2}m_\varphi ^2\varphi ^2`$, yielding an oscillation frequency $`\omega =m_\varphi `$ and $`w0`$ in the far future where $`m_\varphi >>H(t)`$. With such a potential for the quintessence, recollapse would occur, since the curvature term in the Friedmann equation decreases as $`R^2`$ whereas the quintessence term would eventually decrease like the matter, $`R^3`$.
Or the potential could be an exponential with no minimum. These are the most popular current models of quintessence, since such potentials are suggested by supersymmetry. There are many exponential potentials which allow recollapse, as established by Barrow . For example, if the potential dies off sufficiently fast with $`\varphi `$, then in the far future, the density will drop off as $`R^6`$ as does a massless scalar field, the universe will become matter and then curvature dominated in the far future, and recollapse will result.
In summary, there are many quintessence models consistent with all current observations which allow recollapse in the far future. Thus the scenario of horizon elimination proposed in this paper is consistent with all current astronomical observations.
## 8 Conclusion
The “holographic principle” (, , ) claims that all physics on a manifold —especially quantum gravity — can be completely described by a theory defined only on the boundary of that manifold. This is a completely reasonable principle in the case that the boundary of the manifold is a Cauchy surface for the manifold, because in this situation the data on the boundary uniquely determines the manifold and the properties of all physical fields defined on the manifold. For a classical black hole which forms by collapse in an asymptotically flat spacetime and then settles down to a Schwarzschild exterior in the far future, the black hole event horizon is indeed a Cauchy surface for the interior. More generally, if the spacetime is globally hyperbolic, we would expect the event horizon to still be a Cauchy surface for the black hole interior. If we include the c-boundary points in anti deSitter space, then the Cauchy horizons surrounding a region plus the points on the c-boundary where the horizon generators terminate form a Cauchy surface for the interior spacetime region enclosed by the Cauchy horizons; once again we would expect the holographic principle to be valid.
But there are problems with the holographic principle in the case of black holes which evaporate to completion. Since the entire spacetime is no longer globally hyperbolic, it is not clear that the event horizon is a Cauchy surface for the interior. There are problems with the c-boundary completion: looked at from inside a black hole, the c-boundary inside a spherically symmetric black hole is a 2-sphere (the TIPs define a 2-sphere), whereas looked at from the future after the evaporation is complete, the c-boundary is a single point (the TIFs define a point). That is, the causal completion does not define the boundary of the interior manifold uniquely. Even if the event horizon were actually a Cauchy surface for the black hole interior, the information can never leave the horizon to the exterior spacetime, since the event horizon generators must terminate at the singularity which ends the black hole evaporation.
This problem is obviated in an Omega Point spacetime. The null generators of the black hole apparent horizon will actually be a Cauchy horizon for the entire spacetime, for it can be shown that $`I^+(p)`$ is a Cauchy surface for the entire spacetime for any point $`p`$ in the spacetime (Lemma 1 in, p. 436). Thus the holographic principle is true for all manifolds which are future sets (sets for which $`I^+(S)S`$). In particular, for all points $`p_i`$ we wish to include in “black holes” (by any of the defintions given above), the boundary $`I^+(p_i)`$ will be a Cauchy surface for the spacetime, and so the holographic principle will hold for the surfaces of black holes in Omega Point spacetimes.
Another area of general relativity that is naturally complemented by the no-event-horizons resolution of the black hole information problem is the computation of gravitational radiation from colliding black holes. Matzner et al have noted that the computer simulation of a black hole collision is much simpler if characteristic evolution is used in the black hole exterior, because in asymptotically flat spacetimes, the characteristic formulation can be compactified. In an Omega Point spacetime, the characteristic formulation is automatically compactified: the null boundary $`I^+(p)`$ of any point $`p`$ in an Omega Point spacetime has been shown by Tipler to be a compact Cauchy surface (, p. 436), as we pointed out above. We conjecture that the calculation would be even easier done in an Omega Point background space , such as the spherically symmetric Omega Point spacetimes exhibited in Sections 2 and 4. In an Omega Point spacetime, it is not necessary to add the c-boundary points to compactify characteristic null surfaces like $`I^+(p)`$. This is important, because as York has recently emphasized, in general coordinate systems, the initial value problem cannot be well-posed in general relativity. However, Tipler has shown (, p. 440) that Omega Point spacetimes which satisfy the strong energy condition and begin in a “crushing” singularity” (all FRW singularities are of this type, as are all “stable” singularities) possess a unique foliation by constant mean curvature hypersurfaces, and York has shown that the initial data problemis well posed on such a hypersurface. (If the universe is currently accelerating, the strong energy condition will not hold everywhere, but nevertheless a constant mean curvature foliation will still exist, (, p. 439). However, the foliation may only be unique in the very early universe and in the very late universe where the strong energy condition will hold.) As we have emphasized repeatedly, we define black holes operationally in terms of trapped surfaces, just as is done by the groups trying to compute the amount of radiation emitted from colliding black holes. Locally, their calculations of the black hole surfaces would be the same in asymptotically flat space and in an Omega Point spacetime. No quantum effects would effect the location or the size or the existence of trapped surfaces evolved in the black hole collision calculations.
Finally, we point out that many of the well-known difficulties associated with doing quantum field theory in curved spacetimes disappear in Omega Point spacetimes. As we mentioned in Section 1, Omega Point spacetimes necessarily are foliated by compact Cauchy surfaces, and in spacetimes with compact Cauchy surfaces — i.e., in closed universes — there exists a natural unitary equivalence class of quantum field theory constructions, specifically, those constructed from all the Hadamard vacuum states (, p.96). (Roughly speaking, a “Hadamard state” is one in which the short distance singularity structure of the two point function in curved space is the same as it is in Minkowski space , pp. 92–95). In spacetimes without compact Cauchy surface, there are no unitarily equivalent representations of the quantum field algebra, and it was this fact which led many relativists to give up the postulate of unitarity. In an Omega Point spacetime, it is not necessary to give up unitarity.
It is not even necessary to give up the notion of “particle” or “vacuum state” in a curved Omega Point spacetime, as many relativists have previously believed (e.g. , p. 59 and p. 96). The method of Hamiltonian diagonalization (, p. 65) will define a unique vacuum and Fock space with respect to any given Cauchy surface, and we pointed out above that a unique foliation of the spacetime by constant mean curvature exists in a physically realistic Omega Point spacetime (unique except possibly in the periods where the universe is accelerating). These constant mean curvature hypersurfaces are the natural “rest frame” of the universe, and are the natural corresponding frames to the global Lorentz frames in Minkowski space. In FRW universes, the constant mean curvature hypersurfaces are the “rest frames” of the CBR — observers on worldlines normal to these hypersurfaces would measure isotropic CBR temperature. With respect to such a unique global foliation, the notion of “particle” and ”vacuum state” is defined and is unique in curved spacetime.
And such notions must be defined if quantum field theory in curved spacetime is to be a legitimate low energy limit of the (still unknown) quantum theory of gravity. Weinberg (, p. 2) has pointed out that quantum field theories are regarded today as “mere effective field theories,” just low energy approximations to a more fundamental theory. Quantum field theories are not themselves fundamental, but we use them only because any relativistic quantum theory will closely approximate a quantum field theory when applied to particles at a low enough energy. If this is true, then quantum particles are more fundamental than quantum fields, and thus a semi-classical theory like quantum field theory in a classical curved space background must contain a natural definition of the more fundamental entity, the particle.
In short, assuming the universe to end in a c-boundary which is a single point — assuming the universe to be an Omega Point spacetime — solves the black hole information problem, allows the standard concepts of relativistic quantum mechanics to be carried over into curved spacetimes, simplifies the characteristic initial value problem, and is consistent with all astronomical observations. The actual universe may indeed be an Omega Point spacetime.
|
warning/0003/cs0003016.html
|
ar5iv
|
text
|
# Abductive and Consistency-Based Diagnosis Revisited: a Modeling Perspective
## Introduction
Several characterizations have been given for Model-Based Diagnosis (?). All approaches assume that a model of the system to be diagnosed is available: either a model of the correct behavior of the system, or a model of its abnormal behavior, or both.
Diagnostic reasoning has been characterized as a form of nonmonotonic reasoning: either as consistency-based reasoning, or abductive reasoning. In the first case a set of assumptions of correct behavior must be rejected in order to restore consistency with (abnormal) observations; in the second case, a set of assumptions of abnormal behavior must be introduced to entail the abnormal observations.
In (?) the two definitions are shown to be two extremes of a spectrum whose intermediate points may also be relevant, depending on the assumptions about the completeness of the model. Poole also pointed out (?) the importance of the representation problem for logic-based diagnosis, i.e. what has to be represented about the modeled system in order to use the different conceptualizations.
In spite of such previous work, several confusions remain in the field, for example, the confusion of declarative issues with computational issues, such as backward vs forward chaining along the model, and the lack of acknowledgement that in some significant cases the approaches are equivalent.
Moreover, in Model-Based Diagnosis a common view is that modeling is the problem; in particular, any model is an abstraction and the problem is in finding the right way of abstracting the behavior of the system to be modeled. This is a particularly significant issue since one of the claimed advantages of model-based systems is that they can rely on the same model of the system for different reasoning tasks, e.g. planning, diagnosis, configuration, reconfiguration after failure; but, unfortunately, which is the right abstraction, and then the right model, may depend on the task.
This paper, based on a general notion of prediction, illustrates, summarizing and complementing several views in the literature, how the appropriate notion of diagnosis and explanation depends on the predictiveness of the model. In particular, for deterministic models abduction and consistency-based explanation are equivalent, while for nondeterministic models, even if at first sight consistency seems to be the one providing the correct diagnoses, abduction can usually be adapted to provide the same correct diagnoses with a better (at least for someone) notion of explanation.
## Basic definitions
In this paper we assume to rely on a component-based model of the system to be diagnosed; such a model describes the normal and/or abnormal behavior of the system in terms of the normal and/or abnormal behavior of its components. In the classical Reiter’s approach (?) there was no distinction between different abnormal behaviors, and no model or constraints for the abnormal behavior of a component; in later papers (??) the concept of behavioral mode was introduced, and we similarly assume that:
* the system is composed of a set $`COMPS`$ of components;
* each component has different, mutually exclusive, behavioral modes, typically one normal mode and several abnormal modes.
In the classical example of combinatorial circuits, a component could be an AND gate and one of its (abnormal) behavioral modes could be “stuck-at-0”.
In a logical representation of the model, the fact that AND gate $`a`$ is in mode $`stuckat0`$ would be represented with the atomic formula $`stuckat0(a)`$, which would occur in the formulae defining or constraining such behavioral mode, in this case $`andgate(a)stuckat0(a)output(a,0)`$
Here we do not assume that the model is represented in logic, it could also be, e.g., a set of qualitative equations. In this case the equation for the ok mode of the AND gate would be $`out(a)=and(inp1(a),inp2(a))`$, with an appropriate definition of the function $`and`$, while the equation corresponding to the $`stuckat0`$ mode would be $`out(a)=0`$.
In any case we are interested in the notion of parameter in the model. In a model including AND gates, a parameter could be, e.g., the output of AND gate $`a_1`$, which in logic would be the lambda expression<sup>1</sup><sup>1</sup>1This is the expression that, for example, applied to 0 gives the formula $`output(a_1,0)`$ $`\lambda x.output(a_1,x)`$ while in an equational model would be the variable $`out(a_1)`$. A subset of the parameters is the set of observable parameters. Each parameter has a domain of possible values; the parameter in the gate example has domain $`\{0,1\}`$ and in general, in the qualitative models commonly used in Model-Based Diagnosis, the domain would be finite. The granularity of such a domain is a major modeling choice, as, in general, is the choice of the appropriate qualitative abstraction; a first step in providing support for this is given in (?).
As mentioned above, we do not make unnecessary restrictions on the way the model is described and, therefore, on which is the basic inference mechanism. E.g. the model could be a set of logical formulae, with entailment as the inference mechanism, or a set of equations or (more generally) constraints on finite domains, in which case constraint propagation would be the inference mechanism. What the model is required to provide is a notion of prediction relating component behavior to observations as described in the following.
A diagnostic problem is characterized by a set of observations, i.e. an assignment of values to some or all the observable parameters. As in (?) we distinguish between a set $`CXT`$ of “contextual” observations (i.e. “inputs” to the system, e.g. to a circuit) and a set $`OBS`$ of observations to be explained by a diagnosis.
A mode assignment is an assignment of one behavior mode to each component in $`COMPS`$.
We assume that the way the system is modeled provides a notion of prediction, i.e. states whether a mode assignment $`F`$ predicts the set $`S`$ of values for parameter $`p`$ in context $`CXT`$. This means that, in context $`CXT`$ and given the assumptions $`F`$ on the behavior of components, the model of the system implies that $`p`$ takes one of the values in $`S`$.
For example, in a logical framework, where the system is modeled in a set $`MODEL`$ of logical formulae, for a finite set $`S=\{v_1,\mathrm{},v_n\}`$ of values, $`F`$ predicts values $`S`$ for $`p`$ in context $`CXT`$, iff $`MODELFCXTp(v_1)\mathrm{}p(v_n)`$ and the same condition does not hold for any $`S^{}S`$, since we are interested in the most specific prediction.
In an equational model, given the equations $`MODEL`$ representing the system, and the equations $`F`$ and $`CXT`$ corresponding to the mode assignment and the context, the prediction for $`p`$ will be the set $`S`$ of values parameter $`p`$ takes in the solutions of the system of equations $`MODELFCXT`$.
A particularly significant case is of course the one where $`S`$ is a singleton, i.e. the model is able to predict an exact value for the parameter. We will refer to this case as a deterministic prediction. For example, any mode assignment that gives the “stuck-at-0” mode to the andgate $`A_1`$ predicts the value $`0`$ for the output of $`A_1`$, while of course for the “ok” mode the prediction will depend on the context and on the mode of other components.
At the other extreme is the case of a fault making no prediction on a parameter $`p`$, which we intend to coincide with the case where the prediction is the whole domain of the parameter.
An observation on a parameter $`p`$ will, in general, be a set of values $`O`$; in a precise observation (at least as precise as the domain granularity of $`p`$) such a set will be a singleton.
###### Definition 1
Given the prediction $`S`$ on parameter $`p`$ of a mode assignment $`F`$ in context $`CXT`$ and the observation $`O`$ for $`p`$, we say that:
* the prediction is consistent with the observation if $`SO\mathrm{}`$
* the prediction implies the observation if $`SO`$
Obviously, if prediction $`S`$ implies $`O`$ then it is consistent with it<sup>2</sup><sup>2</sup>2Since, as noticed above, we consider no prediction as predicting the whole domain of $`p`$, $`S`$ cannot be empty..
Moreover, in the particular case of a precise observation $`O=\{v\}`$, a consistent prediction must include $`O`$, while to imply $`O`$ a prediction must coincide with $`O`$, i.e. predicting value $`v`$ fr $`p`$.
In (?) a spectrum of definitions of diagnoses is introduced, i.e. a definition with a parameter $`OBS^+OBS`$ representing the subset of the observations to be explained abductively, while for all (other) observation consistency is required. In the terminology introduced above, such a definition can be reformulated as follows.
###### Definition 2
Given diagnostic problem characterized by observations $`CXT`$ and $`OBS`$, and a subset $`OBS^+`$ of $`OBS`$, a diagnosis is a mode assignment $`F`$ such that
1. the predictions of $`F`$ are consistent with $`OBS`$
2. the predictions of $`F`$ imply observations $`OBS^+`$
If $`OBS^+=\mathrm{}`$ this definition is consistency-based diagnosis, if $`OBS^+=OBS`$ it is abductive diagnosis (thus imposing, in general, a stronger requirement than consistency-based diagnosis, reducing therefore the set of diagnoses), and all intermediate choices are possible; in (?) some guidelines are given for this choice, based on the completeness of the model, i.e. on the fact that all the possible explanations for the observations have been provided on the model; this corresponds to the idea of “anticipating explanations” in (?), and will be referred here as backward completeness for the reasons that will be clear in the following.
## Fully predictive models
In the literature there are results on conditions that make abductive and consistency based diagnoses coincide (??), but such results are only formulated for logical representations where predictions are truth values (i.e. a prediction on $`p`$ is either entailing $`p`$ or entailing $`\neg p`$).
In the context of the previous section, we introduce the notion of a model whose prediction is deterministic on all observable parameters.
###### Definition 3
A model is fully predictive if for any context $`CXT`$, for any mode assignment $`F`$, for any observable parameter, $`F`$ makes a deterministic prediction on $`p`$ in context $`CXT`$.
Of course this condition can hardly be met if the domain of a parameter is the set of real numbers, but it is more sensible for e.g. the binary domain of combinatorial circuits and for qualitative abstractions of domains (moreover, in some cases a model can be sensibly expressed in a form that makes it fully predictive, as we will discuss in the next sections).
This condition corresponds to another notion of completeness of the model: a model where the consequences of assumptions can be given non-ambiguously, and, moreover, all the consequences of assumptions have been written down: therefore we can refer to it as forward completeness, as opposed to backward completeness mentioned above.
A trivial result is the following.
###### Property 1
For a fully predictive model, the choice of $`OBS^+`$ in definition 2 is irrelevant, i.e. abductive and consistency-based diagnosis coincide.
Proof. For a fully predictive model, a prediction on $`p`$ is a singleton $`\{v\}`$. If it is consistent with the observation, such observation on $`p`$ must be a set $`S`$ including $`v`$ (or the set $`\{v\}`$ itself). Therefore the prediction implies the observation in $`p`$.
## Nondeterministic models
Even with the underlying assumption that the system to be diagnosed behaves, at a macroscopic level<sup>3</sup><sup>3</sup>3I.e. the one of classical physics, rather than quantum physics., deterministically<sup>4</sup><sup>4</sup>4A faulty system can be intermittently faulty, or, more generally, its behavior could be time-varying. Here we do not consider the temporal dimension — see (?) for the different ways of taking it into account — and we mean that the system behaves deterministically in a time interval during which it is in the same behavioral mode., a model cannot be assumed to be able to predict observations with infinite precision. This is the reason why we mentioned above that fully predictiveness can be too restrictive if the domain of observations are the real numbers.
Models are, usually, a convenient abstraction of a system, and qualitative abstractions of real values have emerged in AI as a meaningful and (sometimes) useful abstraction. Such an abstraction may have pragmatical value, allowing to represent a whole class of possible worlds at a time, and also some cognitive value, since humans (including, sometimes, scientists and engineers) are able to perform qualitative inferences or even tend to reason qualitatively, at least in some stage in the analysis of a system.
Our view is therefore that the system behaves deterministically; it is abstraction which makes the model nondeterministic. Sometimes this abstraction is due to lack of knowledge, sometimes it is just a modeling convenience. For example, we introduce the “flat” fault mode for a battery to include a whole range of values for its voltage, which would avoid predicting exactly the voltage of a flat battery (we will return to this example below).
The problem of abductive hypotheses being insufficiently predictive to entail observations has been discussed early in (?)<sup>5</sup><sup>5</sup>5Dating back to 1987 as a Technical Report. in the framework of plan recognition, proposing as a simple example the fact that the plan of getting food does not explain why the agent, whose plan we are trying to recognize, is going to a specific supermarket: the plan will presumably only entail going to some supermarket. For this reason Kautz rejects the explicit use of abduction and relies on a form of closure to achieve the intended explanations. Such a closure is similar to the ones in (???) (and also explanatory closure in (?)) which in fact provide results showing the equivalence of abduction to deduction (and consistency-based reasoning is a form of deduction) under an appropriate closure.
A different solution can be given based on the following idea in (?). Given an entity, e.g. $`lube`$ $`oil`$, more specific than another, e.g. $`fluid`$, (so that $`lube\mathrm{\_}oil(x)fluid(x)`$), if we want to explain $`lube\mathrm{\_}oil(a)`$, but our assumption only entails $`fluid(a)`$, we rely on transforming the implication $`lube\mathrm{\_}oil(x)fluid(x)`$ into the equivalence $`fluid(x)etc(x)lube\mathrm{\_}oil(x)`$, where $`etc`$ is assumable — corresponding to the assumption that the fluid happens to be lube oil. Assuming $`etc(a)`$ allows entailing $`fluid(a)`$.
Similar ideas are part of the representation methodology in (?), in particular the idea of “anticipating explanations” which is also combined with the idea of “ parametrizing assumptions”. The latter is used for example to represent the model of a flat battery with
$$battery(B)flat(B,V)voltage(B,T)$$
$$flat(B,V)0V1.2$$
where $`flat`$ is assumable, so that e.g. $`voltage(b1,0.8)`$, given $`battery(b1)`$, is explained by $`flat(b1,0.8)`$ which is also consistent with the constraint. This achieves the result of being able to describe concisely a class of faults as well as reasoning, when necessary, on a specific instance of the class.
The parametrized assumption methodology has been adopted e.g. in (?) for the same reason, and in (?) for temporal constraints between the temporal extent of “causes” and “effects”, or, in general, explanans and explanandum. In that case, in fact, the same problem arises unless the temporal extent of the cause uniquely determines the one of the effect.
In general, this problem arises with any case where there are non-functional constraints between parameters of the explanans and those of the explanandum, i.e. the constraints do not allow to entail exact values for the latter even given exact values for the former.
In (?), however, the parametrizing assumptions methodology is rejected due to the potential proliferation of hypotheses, and a notion of explanation is introduced where “explaining” is not entailing, in particular, having a prediction which is an abstraction of the observation is sufficient. Further, a notion of conditional explanation is introduced where prediction and observation must have something that is more specific than both (e.g. $`A`$ for $`AB`$ and $`AC`$). In terms of predicting and observing values for parameters (as in the definitions above) this coincides with prediction and observation having an intersection, i.e. with consistency-based explanation.
## Qualitative deviations
In this section we discuss abductive and consistency-based explanations in the context of a representational abstraction we have used in recent years for modeling systems in the Vehicle Model-Based Diagnosis (VMBD) project: qualitative deviations (???).
The system is modeled in terms of differential equations that include appropriate parameters for components, whose values correspond to different (correct or faulty) behavior modes of the component. From these equations<sup>6</sup><sup>6</sup>6Due to the qualitative abstraction, the exact form of the quantitative equation is not necessary to build the qualitative model; it could however be necessary for fault detection., corresponding equations for qualitative deviations are derived:
* for each variable x, its deviation $`\mathrm{\Delta }x(t)`$ is defined as $`\mathrm{\Delta }x(t)=x(t)x_{ref}(t)`$, where $`x_{ref}(t)`$ is a reference behavior (one choice is to consider the correct behavior of all the system as the reference behavior);
* from any equation $`A=B`$, the corresponding equation $`\mathrm{\Delta }A`$ $`=`$ $`\mathrm{\Delta }B`$ is derived;
* finally, the corresponding sign equation \[$`\mathrm{\Delta }A`$\] = \[$`\mathrm{\Delta }B`$\] is derived; it equates the signs of the two deviations. There are rules for expressing this equation in terms of signs of deviations of individual variables rather than expressions.
These models are useful to express concisely a number of dependencies. For example, for a tank containing a liquid, with an input flow $`in`$ and an output flow $`out`$, the equation
$$\mathrm{\Delta }level=[\mathrm{\Delta }in][\mathrm{\Delta }out]$$
where $`\mathrm{\Delta }level`$ is the sign of the derivative of the deviation of $`level`$ and $``$ is subtraction in the sign algebra, expresses how the level of liquid deviates from the expected value in a wide (and exhaustive) range of cases. In particular, the reference behavior of the system need not be a stable state: even in such a behavior, flows and level may be continuously varying. What the model states is, for example, that:
* starting from the reference behavior, the level does not deviate from it, if the flows do not;
* if the outflow has no deviation and the inflow has a negative deviation, the level will start deviating negatively — this includes the cases where it decreases instead of being constant, decreases more than expected, increases less than expected, becomes constant instead of increasing, and decreases instead of increasing.
But what if the outflow also has a negative deviation? Subtracting a negative number from another can give a positive, negative or zero result. Again, the qualitative abstraction makes the model nondeterministic, at least in some case.
We do not discuss here an appropriate notion of diagnosis for dynamic systems (see (?)), nor a complete model for a system including the tank, but suppose that the inflow is due to a pump (so that a pump fault makes the inflow negatively deviated) and the outflow is regulated by a control system.
If the pump fault occurs, there is a negative deviation of the inflow, so a negative deviation of the level’s derivative, and then of the level; then (we omit the model) the control system reacts with a negative deviation of the outflow, and, given the qualitative model, any result is possible for the level deviation’s derivative. Therefore the fault will not entail any observation on the further trend of the level; this means that it would not be an abductive diagnosis for any observation, while it would be a consistency-based diagnosis for any observation. But would it be considered a good explanation?
We can do something better without abandoning the convenience of qualitative abstractions and using abduction. Suppose e.g. we observe that the level deviation’s derivative changes sign: the pump fault is consistent, but does not predict this; what is the least presumptive assumption that predicts this? The assumption that $`\mathrm{\Delta }out<\mathrm{\Delta }in`$ (which for absolute values means $`|\mathrm{\Delta }out|>|\mathrm{\Delta }in|`$), so that $`[\mathrm{\Delta }in][\mathrm{\Delta }out]`$ gives ’+’. A natural language expression of this explanation is that the control system is compensating for the fault. Notice that the assumption $`\mathrm{\Delta }out<\mathrm{\Delta }in`$ is still a qualitative assumption and could correspond to a natural way of interpreting what is going on in the system.
Thus, in general, for all cases where the result of sign operators is ambiguous, we can introduce assumptions that make the result unambiguous, e.g. for the sum $`[a][b]`$ where $`a`$ and $`b`$ have opposite signs, we can use the assumptions $`|a|<|b|`$, $`|a|=|b|`$, $`|a|>|b|`$ to predict values ’-’,0,’+’ respectively.
## Conclusions
In this paper we have reviewed several points of view in the literature on the relation of the notion of diagnosis with properties of the model of the system to be diagnosed, in particular, properties related to the abstractions in the model with respect to the real system.
We have pointed out that ideas that have appeared in the literature on this subject can be successfully integrated with modeling abstractions that have become widely used more recently.
The problem of identifying the right modeling abstraction for a task is far from being solved, but some work on this is being done in the model-based reasoning community (see e.g. (?) for the temporal dimension in diagnosis) and we expect in the near future more work in this direction, starting e.g. from the results in (?).
|
warning/0003/cond-mat0003460.html
|
ar5iv
|
text
|
# Topological and energetic factors: what determines the structural details of the transition state ensemble and “on–route” intermediates for protein folding? An investigation for small globular proteins.
## I Introduction
Our understanding of the protein folding problem has been thoroughly changed by the new view that has emerged in the last decade. This new view, based on the energy landscape theory and funnel concept , describes folding as the progressive evolution of an ensemble of partially folded structures through which the protein moves on its way to the native structure. The existence of a deep energy funnel in natural proteins and the relatively simple connectivity between most conformational states which are structurally close makes this description possible even when only a few simple reaction coordinates that measure similarity to the native structure are used. The folding mechanism is controlled by both the shape of this free energy landscape and the roughness on it, which arises from the conflicts among interactions that stabilize the folded state and therefore can create non–native conformational traps .
The energetic roughness, however, is not the only limiting factor in determining a sequence’s foldability. Even if the energetic roughness could be completely removed, the folding landscape would not be completely smooth. Theoretical and experimental advances indicate that the final structure of the protein also plays a major role in determining a protein’s foldability. Some particular folding motifs may be intrinsically more designable than others. To address this difference in foldability which is not dependent on energetic frustration, we have introduced the concept of “topological frustration” .
Let us imagine an ideal situation for which the order of native contact formation during folding is not biased. In this “ideal” situation, there are an enormously large number of equivalent folding pathways, and an analysis of the transition state ensemble would show that for this ensemble nearly all parts of the protein have a similar probability of participation. The structure in the transition ensemble has been estimated by analogy with minimalist lattice models made to reproduce the global landscape features of small, fast folding proteins: similar Levinthal entropies, stabilities and energetic roughness as gauged by the glass transition temperatures. These models show a transition state ensemble about half way through the unfolded and folded states . In this ideal case, all the contacts in this transition ensemble would exist with the same probability.
Although the average amount of native formation in the transition ensemble is about 50%, the lattice simulations show that, even when the sequence is designed to have substantially reduced energetic frustration, there are variations in the amount of nativeness of specific contacts in the transition state ensemble . Real proteins display similar heterogeneity in contact formation. In systems with no energetic frustration and equal native interactions, these variations in the transition state ensemble are created solely by the folding motif and polymeric constraints that make certain contacts more geometrically accessible and stable than others. This variation in frequency that some contacts are made in the transition state ensemble generally reduces the entropy of the transition state and, when determined by the native motif, is a gauge of the amount of “topological frustration” in the system. Although this type of frustration can be modified by some design tricks , it cannot be completely eliminated: it reflects an intrinsic difficulty in folding to a particularly chosen shape. Minimalist models have shown how this heterogeneity leads to a transition ensemble that is a collection of diffuse nuclei which have various levels of native contact participation . The minimalist models calibrated to real proteins show similar overall levels of contact heterogeneity as in real proteins . This picture of a transition state composed of several diffuse nuclei has been confirmed by other lattice and off–lattice studies . In addition to selecting sequences which have low levels of energetic frustration, evolution appears to have selected for a particular set of folding motifs which have reduced levels of “topological frustration”, discarding other structures to which it is too difficult to fold .
Guided by theoretical folding studies on lattice, off–lattice, and all–atom simulations (see for instance , , , , , ) as well as recent experimental evidence , we suggest that real proteins, and especially small, fast folding (sub–millisecond), two–state like proteins, have sequences with a sufficiently reduced level of energetic frustration that the experimentally observed “structural polarization” of the transition state ensemble (viz. the variation in the amount of local native structure) is primarily determined by the topological constraints. That is, in well designed sequences, the variations are more determined by the type of native fold than by differences in sequence which leave the native fold relatively unchanged and the energetic frustration small.
The amount of native local structure in the transition state can be experimentally measured by using single and double point mutants as probes in the $`\mathrm{\Phi }`$ value technique . If the topology is a dominant source of heterogeneity in transition state structure, then the majority of evolved sequences which fold to the same motif would exhibit similar local structure in the transition state ensemble. We provide evidence in this paper that not only is this the case, that much of the transition state ensemble is determined by the final folded form, but, also for larger proteins that are not two–state folders, some “on–route” intermediates are determined by topological effects as well. Thus it appears that the dominance of topology in folding extends even into some larger, slower folding proteins with intermediates. This fact is consistent with some recent observations by Plaxco and collaborators that reveal a substantial correlation between the average sequence separation between contacting residues in the native structure and the folding rates for single domain proteins .
To ascertain the extent of topological control of the folding behavior, we create several simplified energetic models of small, globular proteins using potentials created to minimize energetic frustration. We show that these energetically unfrustrated models reproduce nearly all the known global features of the transition states of the real proteins on whose native structures they are based, including the structure of folding intermediates. We directly compare the structure of the transition state ensemble experimentally determined by $`\mathrm{\Phi }`$ value measurements with the numerically determined one. The simulated transition state ensemble is inferred from structures sampled in equilibrium around the free energy barrier between the folded and unfolded states. This free energy is computed as a function of a single reaction coordinate that measures the fraction of formed native contacts. The validity of this method has been demonstrated in references .
The organization of the paper is as follows: in section II we present in some detail the physical concepts underlying this work in the light of recent experimental results. In section III we present results for a sample of five small, globular proteins, and compared these results against the available experimental data. The off–lattice model used in our study is presented in the Appendix. In order to investigate the relevance of the topology, we chose a model which reproduces the topological features of a given real protein and eliminates most of the energetic frustration and variations in the strength of native residue–residue contacts. The predicted transition state for these proteins are in good agreement with experimental evidences, supporting our hypothesis of the major role played by topology.
## II Checking the Folding Mechanism by Analyzing the Transition State Ensemble
How do we know what the folding transition state ensemble looks like? Experimental analysis of folding transition state ensembles has been largely performed using the $`\mathrm{\Phi }`$–value analysis technique introduced by Fersht and co–workers . $`\mathrm{\Phi }`$ values measure the effects that a mutation at a given position along the chain has on the folding rate and stability:
$$\mathrm{\Phi }\frac{RT\mathrm{ln}(k_{mut}/k_{wt})}{\mathrm{\Delta }\mathrm{\Delta }G^0}$$
(1)
where $`k_{mut}`$ and $`k_{wt}`$ are the mutant and wild–type folding rate respectively, $`R`$ is the ideal gas constant, $`T`$ is the absolute temperature, and $`\mathrm{\Delta }\mathrm{\Delta }G^0`$ is the difference in the total stability between the mutant and wild–type proteins in $`kcal/mol`$.
Because the folding event of small fast folding proteins is well described as a diffusive process over a barrier determined by the free energy profile, the folding rate can be written as a Kramer’s–like equation
$$k=k_0\mathrm{exp}[\mathrm{\Delta }G^{}/RT]$$
(2)
where $`k_0`$ is a factor depending on the barrier shape and the configurational diffusion coefficient of the system. If $`k_0`$ is insensitive to small sequence changes, what appears to be true for reasonably unfrustrated sequences the $`\mathrm{\Phi }`$ value is then seen to be a ratio of free energy changes of the folding barrier to stability:
$$\mathrm{\Phi }=\frac{\mathrm{\Delta }\mathrm{\Delta }G^{}}{\mathrm{\Delta }\mathrm{\Delta }G^0}$$
(3)
where $`\mathrm{\Delta }\mathrm{\Delta }G^{}`$ is given by
$$\mathrm{\Delta }\mathrm{\Delta }G^{}=\mathrm{\Delta }G_{mut}^{}\mathrm{\Delta }G_{wt}^{}=RT\mathrm{ln}\frac{k_{mut}}{k_{wt}}.$$
(4)
When this relationship is valid and the mutation can be considered a small perturbation, the $`\mathrm{\Phi }`$ value is a convenient measure of the fraction of native structure which is formed in the transition state ensemble around the site of the mutation. A $`\mathrm{\Phi }`$ value close to 1 means that the free energy change between the mutant and the wild type is almost the same in the transition state and native state, indicating that native contacts involving the mutated residue are already formed at the transition state. Inversely, a $`\mathrm{\Phi }`$ close to $`0`$ means that the free energy change is the same in the transition state and unfolded states, so the local environment of the residue is probably unfolded–like. A detailed analysis of the mutation is needed to determine exactly what contacts are disrupted under mutation. Ideally, mutations are made which eliminate small hydrophobic side–groups. Studies using $`\mathrm{\Phi }`$ values with multiple same–site mutations generally support the accuracy of $`\mathrm{\Phi }`$ value as a structural measurement of the transition ensemble , although sizable changes in the transition state structure have been induced in at least one protein through a single point mutation . In interpreting $`\mathrm{\Phi }`$ values, it is also important to remember that they only measure the relative change in structure, not the absolute amount of structure. This leads to the possibility that some mutants with low $`\mathrm{\Phi }`$ values may have nearly native local environments in the transition state, a possibility seen clearly in the experimental studies of Procarboxypeptidase A2 .
The validity of $`\mathrm{\Phi }`$ values as structural measurements clearly supports the Kramer’s–like description of the folding rate and the fact that the $`\mathrm{\Phi }`$ can be properly understood as a ratio of the free energy change of the transition ensemble over the change of the native ensemble (equation 3). This latter equation is very convenient as a starting point for computing $`\mathrm{\Phi }`$ values. In several recent simulation papers for lattice and off–lattice protein models, we have investigated this issue at length . All these studies concluded that as long as the systems present a weak or moderate level of energetic frustration (such as the Gō–like models in this work), $`\mathrm{\Phi }`$ values determined from changes in the free energy barrier,determined using a single simple reaction coordinate, yield quantitatively correct $`\mathrm{\Phi }`$ values. Therefore, all the calculations performed in this work were done utilizing eq. 3 — no actual kinetics was performed but only the appropriate sampling of the protein configurational space (see Appendix and refs. ,, ,, for example, for details). Technically, as long as the folding barriers are of a few $`k_BT`$ or more and the displacement of the barrier position along this reaction coordinate under mutation is sufficiently small, the $`\mathrm{\Phi }`$ values can be computed using free energy perturbation:
$$\mathrm{\Phi }=\frac{\mathrm{\Delta }\mathrm{\Delta }G^{TS}\mathrm{\Delta }\mathrm{\Delta }G^U}{\mathrm{\Delta }\mathrm{\Delta }G^F\mathrm{\Delta }\mathrm{\Delta }G^U}=\frac{\mathrm{ln}e^{\mathrm{\Delta }E/RT}_{TS}\mathrm{ln}e^{\mathrm{\Delta }E/RT}_U}{\mathrm{ln}e^{\mathrm{\Delta }E/RT}_F\mathrm{ln}e^{\mathrm{\Delta }E/RT}_U}.$$
(5)
We use equation 5 to compute $`\mathrm{\Phi }`$ values for our protein models using fixed transition, unfolded, and folded regions identified by the free energy profile viewed using a single order parameter: $`Q`$, the fraction of native contacts formed in a given conformation.
What experimental evidence exists as to the role of topology in determining the average structure in the folding transition state ensemble? The clearest evidence to date of the role of topology comes from comparisons of the transition state structure of two homologues of the SH3 domain (src SH3 and $`\alpha `$–spectrin SH3). These two homologues have only weak identity ($`30`$% identity with gaps), but $`\mathrm{\Phi }`$ values at corresponding sequence positions are highly correlated , supporting the degeneracy in the folding behavior for these two sequences. Furthermore, one of these sequences has a strained $`\mathrm{\Phi }`$$`\mathrm{\Psi }`$ conformation in the high $`\mathrm{\Phi }`$ region of the distal turn. The fact that this strain does not detectably lower the $`\mathrm{\Phi }`$ values in the local neighborhood , suggests that the sequence details and local stability are less important for determining how structured a region is in the transition state ensemble than its location in the final folded conformation. Other evidence indicates that these results may be more generally applicable than simply for SH3 or $`\beta `$–sheet proteins. Sequence conservation has been shown not to correlate with $`\mathrm{\Phi }`$ values , indicating that in general sequence changes at a given position in a protein weakly affect the $`\mathrm{\Phi }`$ value at that position.
Results for some small fast folding proteins (such as CI2 and the $`\lambda `$repressor) suggest that the transition state is an expanded version of the native state, with a certain degree of additional inhomogeneity over the structure (similar to the theoretical predications for small $`\alpha `$–helical proteins ), while results for other proteins (as the $`\beta `$–sheet SH3 domain) show apparently larger structural heterogeneity in the transition state . This difference in the degree of “structural polarization” that is emerging between small $`\alpha `$–helix and $`\beta `$–sheet proteins suggests that the folding mechanism of a given protein is fundamentally tied to the type of secondary structural elements and their native arrangement. Current studies using $`\mathrm{\Phi }`$ value technique have been made of src SH3 , $`\alpha `$spectrin SH3 , CI2 , Barnase , Barstar , $`\lambda `$repressor , CheY , protein L , Procarboxypeptidase A2 , RNase H and the tetrameric protein domain from tumor suppressor p53 .
In this paper, we analyze five proteins (SH3, CI2, Barnase, RNase H and CheY) that have been extensively studied experimentally and for which, therefore, details of their transition state ensemble are quite well known. We generate sequences (and potentials) for simulating these different globular proteins. These sequences have the native backbone folds of real experimentally studied globular proteins but sequence and potential interactions designed to drastically reduce the energetic frustration and heterogeneity in residue–residue interactions. By comparing the transition state structures of these unfrustrated models with the experimental studies of their real protein cousins, we quantify the effects of the native topology. If topology completely determines how folding occurs, then the model and real proteins should have identical folding behavior and $`\mathrm{\Phi }`$ values. If energetic frustration and heterogeneity are critical for determining the folding mechanism, then the variations in $`\mathrm{\Phi }`$ values with position should bear far reduced similarity to those in the real proteins on which the computer homologues are based.
Two of the five studied proteins are simple two–state like fast folding proteins (SH3 and CI2), while the other three (Barnase, RNase H and CheY) are known to fold through the formation of an intermediate state. We show not only that our simple models can reproduce most of the $`\mathrm{\Phi }`$ value structure, but also that models for Barnase, RNase H and CheY correctly reproduce the folding intermediates of these proteins, suggesting that many of the “on–route” intermediates are also largely determined by the type of native fold.
We represent the five globular proteins using a simplified C<sub>α</sub> model with a Gō–like Hamiltonian as detailed in the Appendix. This potential is in its details unlike that of real proteins, which have residue–residue interactions with many components (Coulomb interactions, hydrogen bonding, solvent mediated interactions, etc., etc.). The crucial features of this potential are its low level of energetic frustration, that characterizes good folders and a native conformation equal to the real protein. The ability of this model to reproduce features of the real transition state ensemble and real folding intermediates is a strong indication that the retention of the topology is enough to determine the global features of their folding mechanism. Using these models, we simulate the dynamics of a protein starting from its native structure, for several temperatures. To monitor the thermodynamics of the system, we group the configurations obtained during a simulation as a function of the reaction coordinate, $`Q`$, defined as the fraction of the native contacts formed in a conformation ($`Q=0`$ at the fully unfolded state and $`Q=1`$ at the folded state). The choice of $`Q`$ as order parameter for the folding is motivated by the fact that in a funnel–like energy landscape, a well designed sequence has the energy of its conformations reasonably correlated to degree of nativeness, and the parameter $`Q`$ is a good measure of the degree of similarity with the native structure. Our Gō–like potential is minimally frustrated for the chosen native structure, and the prediction of transition state ensemble structures and folding rates for these Gō–like systems has been shown to be quite accurate . From the free energy profile as a function of $`Q`$, it is easy to locate the unfolded, folded and transition state ensembles, as it is shown in next section. Since these models consider totally unfrustrated sequences, they may not reproduce the precise energetics of the real proteins, such as the value of the barrier heights and the stability of the intermediates, nonetheless they are able to determine the general structure of these ensembles.
In order to compare the folding process simulated using our model to the actual process for a given protein (as obtained from experimental $`\mathrm{\Phi }`$–values analysis), we need to choose a “mutation” protocol to compute $`\mathrm{\Phi }`$ values. Experimentally, the ideal mutation is typically one that removes a small hydrophobic side–group such as a methyl group that makes well–defined and identifiable residue–residue contacts in the native state. The $`\mathrm{\Phi }`$ value is then sensitive to this known contact. Our computational mutation is the removal of a single native bond, so our computer $`\mathrm{\Phi }`$ values are sensitive to the fractional formation of this bond $`Q_{ij}`$ between residues $`i`$ and $`j`$. We make these mutations because, as in most real mutations, they are sensitive to the formation of specific contacts, rather than being averages over interactions with many parts of the native structure. They mostly resemble the interaction $`\mathrm{\Phi }_{int}`$ value made by making double cycle mutants . $`\mathrm{\Phi }`$ values are computed from equation 5. In an ideal, perfectly smooth funnel–like energy landscape, all the $`\mathrm{\Phi }`$ values should be equal; in an energetically unfrustrated situation, $`\mathrm{\Phi }`$ values variations are due to the structure of the native conformation.
## III Determining the Transition State Ensemble of Small Globular Proteins
We have discussed the idea of “topological frustration” and its role in determining the structural heterogeneity of the transition state ensemble. We explore its role directly by creating protein models which drastically reduce the energetic frustration and energetic heterogeneity among residue–residue native interactions leaving the topology as the primary source of the residual frustration. Results obtained with these models, constructed using a C<sub>α</sub> level of resolution with a Gō–like potential designed to fold to the native trace of chosen proteins, are then compared against the experimental data of those proteins. Five proteins with different folding motifs and different amounts of transition state heterogeneity (variation in $`\mathrm{\Phi }`$ values) and/or intermediates have been investigated.
We first analyze Chymotrypsin Inhibitor II (CI2), a mixed $`\alpha `$$`\beta `$ protein with a broad distribution of $`\mathrm{\Phi }`$ values (nearly uniform from 0 to 1). Then we present an analysis for the src SH3 domain, a largely $`\beta `$–sheet protein with a more polarized transition state structure (a substantial number of large $`\mathrm{\Phi }`$ values). We then apply the same technique to Barnase, Ribonuclease H (RNase H) and CheY, three other mixed $`\alpha `$$`\beta `$ proteins which fold via a folding intermediate. Although these proteins are not two–state folding proteins, we demonstrate that topology is also the dominant determinant of their folding behavior. We show that the topology plays a major role not only in the transition state ensemble, but it is also largely responsible for the existence and general structure of the folding intermediate. This result may be quite common for “on–route” folding intermediates and could provide a computational method for distinguishing between “on–pathway” and “off–pathway” structures which are inferred from experiments. To check the applicability of this method, the same approach presented in this paper has been extended elsewhere to a pair of larger proteins (Dihydrofolate Reductase and Interleukin–1$`\beta `$). Even for these very large proteins we found that the overall structure of the transition state and intermediate ensembles experimentally observed can be obtained utilizing similar simplified models.
### A Analysis of two–state folders: CI2 and SH3
#### 1 CI2
The Chymotrypsin Inhibitor 2 (CI2) protein is a 64 residue protein, consisting of six $`\beta `$–sheets packed against an $`\alpha `$–helix to form a hydrophobic core. Experimental studies have established that CI2 folding and unfolding can be modeled by simple two–state kinetics. The structure of the transition state for this protein has been extensively characterized by protein engineering , by free energy functional approaches , by a geometrical variational principle , and by all–atom molecular dynamics simulations . These studies have shown the transition state has roughly half of the native interactions formed in the transition state ensemble and a broad distribution of $`\mathrm{\Phi }`$ values in agreement with the general predictions of the energy landscape theory used with a law of corresponding states for small proteins . The broad distribution of $`\mathrm{\Phi }`$ values suggests that most hydrophobic contacts are represented at a level of about 50% in the transition state ensemble.
We constructed a Gō–like C<sub>α</sub> model of CI2 as described in the Appendix. Several fixed temperature simulations were made and combined using the WHAM algorithm to generate a specific heat versus temperature profile and a plot of the potential of mean force as a function of the folding order parameter $`Q`$ (see figure 1). From the free energy profile, we identified the dominant barrier, and used the thermal ensemble of states at its location to generate $`\mathrm{\Phi }`$ values from equation 5. The ranges of values of $`Q`$ used to determine each of these ensembles are shaded in figure 1. The mutations have been implemented by the removal of single attractive interactions (they are replaced with the same short ranged repulsive interactions used between residues without native interactions). The values computed via this method are shown in figure 2. Also shown in this figure is the fractional formation of individual native contacts in the transition state. The small difference between these two figures is primarily due to the fact that in the $`\mathrm{\Phi }`$ calculations the native contact formation at the folded and unfolded states are also taken into account. Because of the higher concentration of contacts between residues near–by in sequence and the local conformational preferences, the unfolded state shows a high level of local structure. The inaccurate representation of local contacts in the unfolded state makes the short range $`\mathrm{\Phi }`$ values less reliable as transition structure estimates than long range $`\mathrm{\Phi }`$ values.
From the calculations, we detect three significant regions of large $`\mathrm{\Phi }`$ values: the $`\alpha `$–helix, the mini–core defined by strands 3 and 4 and their connecting loop, and between the C–terminus of strand 4 and the N–terminus of strand 5. These regions generally have $`\mathrm{\Phi }`$ values in excess of $`0.6`$. Slightly smaller values of about $`0.5`$ exist for the short range contacts between the N–terminal of strand 3 and the C–terminal of the $`\alpha `$–helix and for contacts between strand 3 and strand 4. All other regions lack a consistent set of large $`\mathrm{\Phi }`$ values. Despite the large number of native contacts between strands 1 and 2 and the $`\alpha `$–helix and between strands 5 and 6 and the $`\alpha `$–helix, only low $`\mathrm{\Phi }`$ values are observed in this region (nearly all below $`0.2`$ in value). A comparison between these data and the exhaustive analysis of Fersht and colleagues shows excellent overall agreement. They have found that “$`\beta `$–strands 1, 5 and 6 … are not structured in the transition state….”. Strand 2 also shows a highly reduced amount of structure. Furthermore, “the central residues of $`\beta `$–strands 3 and 4 interact with the $`\alpha `$–helix to form the major hydrophobic core of CI2.” The hydrophobic mini–core in this region (defined as the cluster formed by side–chains of residues 32, 38, and 50) is detected by single mutant and double mutant $`\mathrm{\Phi }`$–values to be at least 30% formed in transition ensemble. Similarly, they found the $`\alpha `$–helix, particularly the N–capping region, to be highly ordered.
In summary, we see a quite good overall agreement except for a discrepancy in the short range interactions in the loop region between strands 4 and 5. This protein shows generally higher $`\mathrm{\Phi }`$ values between interactions which are more local in sequence and lower $`\mathrm{\Phi }`$ values between interactions which are distant in sequence. The results are thus consistent with the picture of the transition state as a collection of non–specific and somewhat diffuse nuclei . This overall low level of frustration suggests a low level of “topological frustration” in this model as well and a particularly designable motif.
#### 2 src SH3 domain
Src SH3 is the 57 residue fragment of Tyrosine–Protein Kinase that stretches from T84 to S140. It has five $`\beta `$ strands (and a short 3–10 helix) in an anti–parallel arrangement, forming a partial $`\beta `$ sandwich. Experimental measurements have shown that the SH3 domain folds using a rapid, apparently two–state mechanism. A $`\mathrm{\Phi }`$ value analysis reveals that the distal loop hairpin and diverging turn regions are both highly structured and docked together at the transition state; the hydrophobic interactions between the base of the hairpin and the strand following the diverging turn are partially formed, while other regions of src SH3 appear only weakly ordered in the transition state ensemble. The overall representation of the transition state structure of src SH3 —having the distal loop and diverging turn largely formed and other regions weakly formed— agrees with studies of $`\alpha `$spectrin SH3 which has a similar backbone structure but a dissimilar sequence (approx 30% identity with gaps). This observed similarity along with evidence of a strained backbone conformation in the distal loop of the $`\alpha `$spectrin SH3 supports the concept of “topological” dominance in folding .
Fig. 3 shows the folding behavior as obtained from our dynamics simulations of the Gō–like analogous of the src SH3. The free energy barrier defining the transition state location is evident in the figure. As before, we have computed $`\mathrm{\Phi }`$ values from equation 5 by mutating (removing) every native residue–residue attractive contact. The results of this calculation are shown in figure 4. In addition to $`\mathrm{\Phi }`$ values, the contact formation probability at the transition state ensemble have been calculated. Our previous caveats concerning $`\mathrm{\Phi }`$ values for local interactions still apply. We observe the highest collection of off–diagonal (long range) $`\mathrm{\Phi }`$ values is in the diverging turn —distal loop interaction exactly as seen from the experimental $`\mathrm{\Phi }`$ value measurements. We see very low values in the RT loop region, in accord with the two mutants in this loop. We also see medium to high values between the two $`\beta `$ strands which are connected by the distal loop. The transition state structure of the SH3 presents a substantially larger degree of structural polarization than CI2, where the $`\mathrm{\Phi }`$ values are much more uniform. This suggests that SH3 has a backbone conformation which is intrinsically more difficult to fold, i.e., there is a greater level of “topological frustration” in this structure. Nevertheless the transition state composition is well reproduced for both the two proteins.
### B Analysis of three proteins which fold throughout the formation of an intermediate state: Barnase, RNase H and CheY
Barnase, RNase H and CheY are three small $`\alpha `$$`\beta `$ proteins (although larger than the previous two proteins): Barnase is a 110 residue protein, composed by three $`\alpha `$–helices (located in the first 45 residues) followed by five $`\beta `$–strands; RNase H consists of 155 residues which arrange themselves in five $`\alpha `$–helices and five $`\beta `$–strands; CheY is a 129 residues, classic $`\alpha \beta `$ parallel fold in which five $`\beta `$–strands are surrounded by five $`\alpha `$–helices. Experimental results show that these three proteins do not fold by following a simple two–state kinetics directly from the unfolded state to the native structure, but fold through the formation of a metastable intermediate which interconverts into the native state. This brings up an interesting question: is topology alone able to determine the presence of an intermediate in the folding process? In Figs. 5, 7 and 9 we show evidence for the first time that such intermediates can be created solely from a Gō–like minimalist model which preserves the native topology. The presence of this intermediate during these protein’s folding events is a requirement of the native protein motifs. The free energy changes upon mutations of a wild–type three–state protein are experimentally measured both for the intermediate and the transition state, to define two different sets of $`\mathrm{\Phi }`$–values for the protein:
$$\begin{array}{c}\mathrm{\Phi }_I=\frac{\mathrm{\Delta }\mathrm{\Delta }G_I\mathrm{\Delta }\mathrm{\Delta }G_U}{\mathrm{\Delta }\mathrm{\Delta }G_F\mathrm{\Delta }\mathrm{\Delta }G_U}\\ \mathrm{\Phi }_{TS}=\frac{\mathrm{\Delta }\mathrm{\Delta }G_{TS}\mathrm{\Delta }\mathrm{\Delta }G_U}{\mathrm{\Delta }\mathrm{\Delta }G_F\mathrm{\Delta }\mathrm{\Delta }G_U}\end{array}$$
(6)
where $`\mathrm{\Phi }_I`$ provides information about the structural composition of the intermediate state (I), and $`\mathrm{\Phi }_{TS}`$ of the transition state (TS). Following we discuss in some details the results for the three proteins. Since, as for the first two proteins, the $`\mathrm{\Phi }`$–values and the native contact probabilities provide somewhat similar information, for simplicity, we show only the results obtained for the native contact probabilities (for safety we have checked the $`\mathrm{\Phi }`$–values and determined that similar information is recovered).
#### 1 Barnase
The analysis of experimentally obtained $`\mathrm{\Phi }`$ values for the Barnase shows that some relevant regions of the structure are fully unfolded in the intermediate while other regions are fully folded.
Fig. 6 shows the intermediate and the transition state structure obtained from the Gō–like model. The intermediate shows substantial structural heterogeneity: there are very high probability values for interactions within the $`\beta `$–sheet region and its included loops, and very low values for interactions within the $`\alpha `$–helices and their loops and between the $`\alpha `$–helical and $`\beta `$–sheet regions. Some local short range helical interactions are formed. The transition state ensemble structure shows the same structure as the intermediate with the addition of strong interactions within helices 2 and 3; between helix 2, helix 3, the first $`\beta `$–strand, and the intervening loops; and between the second $`\beta `$–strand and the second helix.
Comparing these simulation results with extensive mutagenesis studies of reference , we observe a good qualitative agreement. The $`\beta `$–sheet region is highly structured in the intermediate as it is the core region 3 (consisting of the packing of loop 3, that joins strands 1 and 2, and of loop 5, that joins strands 4 and 5, with the other side of the $`\beta `$–sheet). In agreement with experiments, the earliest formed part of the protein appears to be the $`\beta `$–sheet region. Also the core region 2 (formed by the hydrophobic residues from helix 2, helix 3, the first strand, and the first two loops) is found to be only weakly formed in the intermediate and the transition state.
There are two minor discrepancies between the Barnase model and the experimental data. First, we slightly overestimate the formation of core region 2 in the transition state ensemble. Second, we underestimate the amount of structure in core region 1 (formed by the packing of the first helix against a side of the $`\beta `$–sheet) in both intermediate and transition ensemble. In particular, we under–represent the interaction between helix 1 and the $`\beta `$–sheet region. The experimentally observed early packing of helix 1 against the rest of the structure is not reproduced by our model. Clearly there are some important energetic factors which have been neglected by the simple model. These may be inferred from the Barnase crystal structure. For example, one can see that helix 1 is largely solvent exposed, with interactions between it and the remainder of the protein formed by only five of the eleven helix residues. 83 % of the interactions reside on the hydrophobic residues PHE7, ALA11, LEU14 and GLN15, and the 17 % of the interactions are formed by the charged residues ASP8 and ASP12, while the solvent exposed part of the helix is composed of polar residues. Large stabilizing interactions other than tertiary (most hydrophobic) interactions are neglected in the model, being probably responsible for the failure in predicting the formation of the structural parts involving helix 1. In this structural detail, it appears that the topological factors are not the leading determinant of the folding behavior.
#### 2 Ribonuclease H
Kinetic studies of the wild–type RNase H have shown that an intermediate state is populated in the folding process, and the structure of this intermediate has been extensively investigated by circular dichroism, fluorescence and hydrogen exchange methods and by protein engineering . Fig. 7 shows that, consistently with the experimental evidences, we find an intermediate state in the folding process of the RNase H model. Experimental results indicate that the most stable region of the protein intermediate involves the $`\alpha `$–helix 1, the strand 4, the $`\alpha `$–helix 4 and the $`\alpha `$–helix 2. Hydrogen exchange experiments have shown that the $`\alpha `$–helix 1 is the region of the protein most protected from exchange, suggesting that most of the interactions involving the $`\alpha `$–helix 1 are already significantly formed at the intermediate state of the folding process. The helix 4 and the $`\beta `$–strand 4 are the next most protected regions, while the $`\alpha `$–helix 5 has low to moderate level of protection. After the completion of the this intermediate structure, the rate–limiting transition state involves the ordering of the $`\beta `$–sheet and the $`\alpha `$–helix 5. The packing of helix 5 across the sheet is found to be the latest folding event.
The results of the model for RNase H show a good agreement with the experimental evidences. As shown in Fig. 8, we find that the formation of contacts involving the helix 1 is the earliest event in the folding process. Contacts arising from the $`\alpha `$–helix 4 and the $`\beta `$–strand 4 are then formed at the intermediate state and consolidated at the transition state. In agreement with the experimental results, we find that, at the transition state, interactions between the $`\alpha `$–helix 1, the strand 4 and the rest of the protein are mostly formed; the $`\alpha `$–helix 4 is also well structured and interactions between the helix 4 and the other parts of the protein are partly formed. Interactions among the strands are almost all formed, but the sheet is not yet docked to the helix 5.
#### 3 CheY
Utilizing protein engineering , the transition state of CheY has been characterized and it can be described as a combination of two subdomains: the first half of the protein (subdomain 1), comprising the $`\alpha `$–helices 1 and 2 and the $`\beta `$–strands 1–3, is substantially folded whereas the second half (subdomain 2) is completely disorganized. The helix 1 seems to play the role of a nucleation site around which subdomain 1 begins to form. Moreover, an intermediate has been detected at the early stage of the folding process where all the five $`\alpha `$–helices are rather structured. The last two helices, however, are very unstructured in the later occurring transition state. From this result it has been suggested that a misfolded species is visited at the beginning of the folding process.
Our simple model detects two possible intermediates for this protein, one of them is an “on–route” intermediate that is short–living and occurs just before the transition state ensemble ($`Q`$ around 0.6 in Fig. 9). Surprisingly, the unfrustrated model is also able to detected a “misfolded” trap in the folding of CheY. Since non–native interactions are not allowed in the model, this trap is a long–living partially folded state created by the topological constrains. There is no direct connection between this trap state and the fully folded state. The structure of this trap is shown in Fig. 10 and it agrees with the experimental observation of all helices well structured. Differently from the previously discussed proteins, the model of CheY seems to have a tendency to first form a “wrong” part of the protein and, when this happens, a partial unfolding must occur before the folding can be completed.
Finally, analyzing the transition state structure, we find a good agreement with the experimental data. As shown in Fig. 10, the first part of the protein (subdomain 1) is almost fully folded at the transition state ensemble, while subdomain 2 is completely unfolded.
## IV Conclusions
Recent theoretical studies and experimental results suggest that the folding mechanism for small fast folding proteins is strongly determined by the native state topology. The amount of energetic frustration, arising from the residual conflict among the amino–acid interactions, appears largely reduced for these proteins so that topological constraints are important factors in governing the folding process. Towards exploring this topological influence in real proteins, we analyzed the folding process of the Gō–like analogous of five real proteins. Since we have used Gō–like potentials, the energetic frustration is effectively removed from the system, while the native fold topology is taken into account. It is important to highlight that the results from such studies exhibit the overall topological features of the folding mechanism, although we do not expect the precise energetic values for barrier heights and intermediate state stabilities. For example, real proteins are not necessarily totally unfrustrated and they have only to minimize energetic frustration to a sufficiently reduced level in order to be good folders. Also, as long as energetic frustration is small enough, creating some heterogeneity at the native interactions may help to reduce topological frustration , and that will energetically favor some contacts over others.
The effective use of a small number of global order parameters as reaction coordinates, in interpreting real data or studying more detailed protein folding model, depends critically on the degree of frustration present in real proteins . Since our results show that general structural features of the transition state ensemble in real proteins, at least for this class of fast folding proteins, is reproducible by using a substantially unfrustrated potential, several different global order parameters should work to explain the folding mechanism. For this reason, it should not be a surprise the fact that, utilizing energy landscape ideas and the funnel concept, some very simple models with approximate order parameters determined by single or few sequence approximation have been successful in predicting qualitative features of the transition state ensemble.
Again, we have compared in details the structure of the transition state ensemble of the five proteins resulting from our simulations with experimental data. The agreement between our results and the experimental data supports the idea that energetic frustration is indeed sufficiently reduced and the protein folding mechanism, at least for small globular proteins, is strongly dependent on topological effects. The structure of the transition state ensemble of the CI2 presents a broad distribution of $`\mathrm{\Phi }`$ values —i.e. a reduced degree of structural polarization— in agreement with predictions based on the energy landscape theory (see , ). On the other hand, the structure of the SH3 transition state ensemble shows a higher degree of polarization. Nevertheless, by using our simplified Gō–like model, we have reproduced the transition state composition for both proteins, demonstrating that topology is largely responsible for the observed experimental differences. The last three proteins we have analyzed, (Barnase, RNase H and CheY) are known to fold through a three–state kinetics, involving the formation of an intermediate structure. Our Gō–like model of these proteins also fold with a three–state kinetics with intermediates that are analogous to the ones detected experimentally. This fact suggests that topology is also a dominant factor in determining the “on–route” intermediates.
## V acknowledgments
We thank Viara Grantcharova and David Baker for informations about the SH3 structure. We also thank Vladimir Sobolev for the CSU software. We are indebted to Angel García, Peter Wolynes, Steve Plotkin, Jorge Chahine, Joan Shea, Margaret Cheung, Charlie Brooks, Amos Maritan and Jayanth Banavar for helpful discussions. One of us (C.C.) expresses her gratitude to Giovanni Fossati for his suggestions and for carefully reading the manuscript, and to the Center for Astrophysics & Space Sciences of UCSD for the usage of graphics facilities and computer time. This work has been supported by the NSF (Grant #96–03839), by the La Jolla Interfaces in Science program (sponsored by the Burroughs Wellcome Fund) and by the Molecular Biophysics training grant program (NIH T32 GN08326).
## Model and Method
In order to investigate how the native state topology affects the folding of a given protein we follow the dynamics of the protein by using a Gō–like Hamiltonian to describe the energy of the protein in a given configuration. A Gō–like Hamiltonian takes into account only native interactions, and each of these interactions enters in the energy balance with the same weight. It means that the system gains energy as much as any amino acid pair involved in a native contact is close to its native configuration, no matter how strong the actual interaction is in the real protein. Residues in a given protein are represented as single beads centered in their C–$`\alpha `$ positions. Adjacent beads are strung together into a polymer chain by mean of bond and angle interactions, while the geometry of the native state is encoded in the dihedral angle potential and a non–local potential. The energy of a configuration $`\mathrm{\Gamma }`$ of a protein having the configuration $`\mathrm{\Gamma }_0`$ as its native state is thus given by the expression:
$$\begin{array}{c}E(\mathrm{\Gamma },\mathrm{\Gamma }_0)=_{bonds}K_r\left(rr_0\right)^2+_{angles}K_\theta \left(\theta \theta _0\right)^2+\\ _{dihedral}K_\varphi ^{(n)}\left[1+\mathrm{cos}\left(n\times (\varphi \varphi _0)\right)\right]+\\ _{i<j3}\{ϵ(i,j)[5\left(\frac{\sigma _{ij}}{r_{ij}}\right)^{12}6\left(\frac{\sigma _{ij}}{r_{ij}}\right)^{10}]+ϵ_2(i,j)\left(\frac{\sigma _{ij}}{r_{ij}}\right)^{12}\}.\end{array}$$
(7)
In the previous expression $`r`$ and $`r_0`$ represent the distances between two subsequent residues at, respectively, the configuration $`\mathrm{\Gamma }`$ and the native state $`\mathrm{\Gamma }_0`$. Analogously, $`\theta `$ ($`\theta _0`$) and $`\varphi `$ ($`\varphi _0`$) represent the angles formed by three subsequent residues and the dihedral angle defined by four subsequent residues along the chain at the configuration $`\mathrm{\Gamma }`$ ($`\mathrm{\Gamma }_0`$). The dihedral potential consists of a sum of two terms for every four adjacent $`C_\alpha `$ atoms, one with period $`n=1`$ and one with $`n=3`$. The last term in Eq. (7) contains the non–local native interactions and a short range repulsive term for non–native pairs (i.e. $`ϵ(i,j)=constant>0`$ and $`ϵ_2(i,j)=0`$ if $`i`$$`j`$ is a native pair while $`ϵ(i,j)=0`$ and $`ϵ_2(i,j)=constant>0`$ if $`i`$$`j`$ is a non–native pair). The parameter $`\sigma _{ij}`$ is taken equal to $`i`$$`j`$ distance at the native state for native interactions, while $`\sigma _{ij}=4`$ Å for non–native (i.e. repulsive) interactions. Parameters $`K_r`$, $`K_\theta `$, $`K_\varphi `$, $`ϵ`$ weight the relative strength of each kind of interaction entering in the energy and they are taken to be $`K_r=100ϵ`$, $`K_\theta =20ϵ`$, $`K_\varphi ^{(1)}=ϵ`$ and $`K_\varphi ^{(3)}=0.5ϵ`$. With this choice of the parameters we found that the stabilizing energy residing in the tertiary contacts is approximately twice the stabilizing energy residing in the torsional degrees of freedom. This balance among the energy terms is optimal to study the folding of our Gō–like protein models. The native contact map of a protein is derived with the CSU software based upon the approach developed in ref. . Native contacts between pairs of residues $`(i,j)`$ with $`ji+3`$ are discarded from the native map as any three and four subsequent residues are already interacting in the angle and dihedral terms. A contact between two residues $`(i,j)`$ is considered formed if the distance between the $`C_\alpha `$’s is shorter than $`\gamma `$ times their native distance $`\sigma _{ij}`$. It has been shown that the results are not strongly dependent on the choice made for the cut–off distance $`\gamma `$. In this work we used $`\gamma =1.2`$. We have used Molecular Dynamics (entailing the numerical integration of Newton’s laws of motion) for simulating the kinetics of the protein models. We employed the simulation package AMBER (Version 4.1) at constant temperature, i.e. using Berendsen algorithm for coupling the system to an external bath . Both temperature and energy are measured in units of the folding temperature $`T_f`$ in the simulations.
For each protein model, several constant temperature simulations were made and combined using the WHAM algorithm to generate a specific heat profile versus temperature and a free energy $`F(Q)`$ as a function of the folding reaction coordinate Q. This algorithm is based on the fact that the logarithm of probability distribution $`P(Q)`$ of the values taken by a certain variable Q (e.g. the order parameter) at fixed temperature T may serve as an estimate for the the free energy profile $`F(Q)`$ at that temperature. In fact, the probability to have a certain value $`Q_1`$ for the variable Q, at temperature $`T=1/\beta `$, in the canonical ensemble is given by:
$$P_\beta (Q_1)=\frac{W(Q_1)e^{\beta E(Q_1)}}{Z_\beta }$$
(8)
where $`W(Q)`$ is the density of configurations at a point $`Q`$ in the configurational space, $`Z_\beta `$ is the canonical partition function at temperature $`T=1/\beta `$ and $`E(Q)`$ is the energy of the system at the value Q of the reaction coordinate<sup>*</sup><sup>*</sup>*Since our model is almost energetically unfrustrated, the energy fluctuations for a set of configurations with fixed Q are strongly reduced such that the energy in a given configuration could be considered as a function of $`Q`$.. Since the free energy $`F`$ is
$$F(Q)=E(Q)TS(Q)$$
(9)
and the entropy S(Q) is related to the configurational density $`W(Q)`$
$$W(Q)e^{S(Q)/k}$$
(10)
where $`k`$ is the Boltzmann constant, it follows that
$$\frac{P_\beta (Q_1)}{P_\beta (Q_2)}=\frac{e^{\beta F(Q_1)}}{e^{\beta F(Q_2)}}$$
(11)
and free energy differences can be computed by
$$\beta (F(Q_1)F(Q_2))=\mathrm{log}\frac{P_\beta (Q_1)}{P_\beta (Q_2)}.$$
(12)
By using the procedure of refs. , data from a finite set of simulations can be used to obtain complete thermodynamic information over a large parameter region.
Probability distributions are obtained by sampling the configurational space during Molecular Dynamics simulations.
For the smaller proteins (CI2 and SH3) we have determined the errors on the estimates of the transition temperature and contact probabilities (or $`\mathrm{\Phi }`$ values). This has been accomplished by computing these quantities from several (more than 10) uncorrelated sets of simulations. We found that the standard deviation for each single contact probability is 0.06 for CI2 and 0.05 for SH3, while the transition temperature is determined in both cases with an uncertainty smaller than 0.5%. These errors are obtained using about 200 uncorrelated conformations in the transition state ensemble. Since Barnase, RNase H and CheY have twice to three times the number of tertiary contacts of SH3 and CI2, in order to have appropriate statistics, we have sampled about 500 uncorrelated conformations (thermally weighted) for every transition state ensemble or intermediate.
### Captions to the figures
Fig. 1. (a) Free energy $`F(Q)`$ as a function of the reaction coordinate $`Q`$ around the folding temperature for the model of CI2. Free energies are measured in units of $`k_BT_f`$, where $`T_f`$ is the folding temperature. The unfolded, folded and transition state regions are shown in the light blue shaded areas. (b) A typical sample simulation at a temperature around the folding temperature. The reaction coordinate Q as a function of time (measured in arbitrary unit of molecular dynamics steps) is shown. The two–state behaviour is apparent from the data. The unfolded and folded states are equally populated at the folding temperature. (c) Heat capacity as a function of the temperature (units of folding temperature).
Fig. 2. The results for the transition state structure from the simulations for CI2. The probability of native contact formation at the transition state (left panel), and bond $`\mathrm{\Phi }`$–values (right panel) are shown. Different colors indicate different values from 0 to 1, as quantified by the color scale. The $`\alpha `$–helix, the interactions between the strands 4 and 5, and the minicore (i.e. interactions between residues 32,38 and 50) are the parts of the structure formed with the highest probability, although they are not fully formed. Overall, the transition state ensemble appears as an expanded version of the native state where most contacts have a similar probability of participation, but some interactions are less like to occur. These results agree with the transition state structure experimentally obtained.
Fig. 3. (a) Free energy $`F(Q)`$ as a function of the reaction coordinate $`Q`$ for a set of temperatures around the folding temperature. Free energies are measured in units of $`k_BT_f`$. The choices for the unfolded, folded and transition state regions are marked as shaded regions. (b) The reaction coordinate Q as a function of time (unit of molecular dynamics steps), from a typical sample simulation around the folding temperature. As in Fig. 1, the two–state behaviour is apparent. At the transition temperature the model protein has equal probability to be found in the unfolded or in the folded state. (c) Heat capacity as a function of the temperature, in units of folding temperature.
Fig. 4. The transition state structure as obtained from the simulations for SH3. Panel in the left represents the probability for a native contact to be formed at the transition state, while the panel in the right shows the results for bond $`\mathrm{\Phi }`$–values. Different colors indicate different values from 0 to 1, as quantified by the color scale. Diverging turn and distal loop are marked on the contact map. The interactions within and between these two parts of the protein chain appear to be formed with high probability. The interactions between the two strands joined by the distal loop are partially formed, while the contacts involving the first 20 residues do not contribute to the transition state structure. This description of the transition state is in agreement with experimental results.
Fig. 5. (a) Free energy $`F(Q)`$ of Barnase protein as a function of the reaction coordinate $`Q`$ around the folding temperature. Free energies are measured in units of $`k_BT_f`$. The unfolded, folded and intermediate state regions are marked in green, while the top of the two barriers are marked in light blue. The local minimum in the free energy profile between the unfolded and folded minima locates the folding intermediate state. The presence of a folding intermediate state is also evident from panel (b), where the order parameter $`Q`$ is plotted as a function of time for a typical molecular dynamics simulation around the folding temperature. In the interval $`Q(0.40.5)`$, the same state (i.e. with the same average structure) is visited both from the unfolded and folded structures.
Fig. 6. The probability of native contact formation for the intermediate (left panel) and transition state (right panel) structures as obtained from our simulations of Barnase. Different colors indicate different values from 0 to 1, as quantified by the color scale. The earliest formed part of the protein appears to be the $`\beta `$–sheet region, in agreement with experimental results. The core 3 (formed by loops 3 and 5 to the $`\beta `$–sheet) is formed at the intermediate and transition state, while core 1 (the packing of the helix 1 against the $`\beta `$–sheet) and the core 2 (the interactions between the hydrophobic residues from the helices 2 and 3, the strand 1, and the first two loops) start to form only after the transition state. The formation of the $`\alpha `$–helix 1 occurs as a late event of the folding from our simulations, while from experimental results it seems to be already formed at the intermediate and transition state. The early formation of the $`\alpha `$–helix is most probably due to energetic factors rather then from topology requirements (and then beyond the prediction possibility of this model), as detailed in the text.
Fig. 7. (a) Free energy $`F(Q)`$ of the model of RNase H as a function of the reaction coordinate $`Q`$ around the folding temperature. Free energies are measured in units of $`k_BT_f`$. The regions corresponding to the unfolded, folded and intermediate state are marked in green, while the top of the two barriers are marked in light blue. A folding intermediate is detected as a local minimum in the free energy between the unfolded and folded minima. In panel (b) the fraction of native contacts formed, $`Q`$, is plotted versus the simulation time for a sample of our simulations (at a temperature $`T=0.99T_f`$) where the transition from unfolded to folded state is observed. The local minimum of panel (a) corresponds to a transiently populated intermediate (located at $`Q`$ around 0.4) that later evolves to the fully folded state.
Fig. 8. The probability of native contact formation at the intermediate (left panel) and transition state (right panel) structure, as observed for the RNase H model. Different colors indicate different values from 0 to 1, as quantified by the color scale. In agreement with experimental results, we found that interactions involving the $`\alpha `$–helix 1 are the first formed in the folding process. Contacts between the $`\alpha `$–helix 1 and the strand 4 are highly probably formed at the intermediate. Also the $`\alpha `$–helix 4 is well structured and the $`\beta `$–sheet is partly formed. These interactions strengthen at the transition state where also the $`\beta `$–sheet is almost completely formed, while the packing of helix 5 across the sheet is not yet accomplished.
Fig. 9. (a) Free energy $`F(Q)`$ profile for the model of CheY plotted as a function of the reaction coordinate $`Q`$ for a set of temperatures around the folding temperature. Free energies are measured in units of $`k_BT_f`$. Differently from the corresponding figures of Barnase (Fig. 5) and RNase H (Fig. 7), two different structures are populated between the folded and unfolded states. In addition to the “on–route” intermediate state (marked in green as the regions corresponding to the folded and unfolded states), a “misfolded” intermediate structure (marked in brown at Q around 0.4) is transiently visited from the unfolded state. The top of the two barriers are marked in light blue. In agreement with experimental results, we found that in this “misfolded” structure, all the five $`\alpha `$–helices are rather structured while, in the later occurring “on–route” intermediate and transition state ensemble, the helices 4–5 are completely unstructured (see fig. 10). Panel (b) shows a typical sample of the simulation around the folding temperature, in a region where the folding occurs. The first transiently populated intermediate state corresponds to a structure where all the helices are formed. Before to proceed to the folded state, a partial unfolding occurs.
Fig. 10. The probability of the native CheY contacts to be formed in the “misfolded” intermediate (left panel) and transition state (right panel) for the model protein. Different colors indicate different values from 0 to 1, as quantified by the color scale. In agreement with experimental data, all the helices are mostly formed in the transiently populated “misfolded” structure, while helices 4 and 5 are rather unstructured at the transition state. The two subdomains experimentally detected in the CheY transition state are evident in the figure: the first part of the protein (all interactions arising from the $`\alpha `$–helices 1–2 and the $`\beta `$–strands 1–3) is folded, while the second part (interactions among the $`\alpha `$–helices 4–5 and the $`\beta `$–strands 4–5) is completely unfolded. The helix 3 is structured but the interactions between the helix 3 and the rest of the protein are not completely formed.
|
warning/0003/nucl-ex0003009.html
|
ar5iv
|
text
|
# Measurements of Light Nuclei Production in 11.5 A GeV/c Au + Pb Heavy-Ion Collisions
## I Introduction
Relativistic heavy ion collisions are believed to reach energy densities an order of magnitude greater than that of normal nuclear matter. These collisions allow the examination of the strong interaction in a novel environment as well as providing a possible doorway to new states of matter. In order to understand the dynamics of the collision system, one must use the only available tools-the species and momenta of the particles which exit the collision region. The use of emitted hadrons to probe the collision system is complicated because these hadrons rescatter many times as they traverse the collision region, and consequently lose some of their direct information about the earlier stages of the time evolution. However, the final space-time extent of the system at freeze-out (the time when strong interactions cease) and position-momentum correlations of the emitted particles contain much information about the entire time evolution. In principle, these carry information about the equation-of-state of the early collision region.
In order to extract information about both the momentum and position distributions of the source at freeze-out, it is necessary to measure multi-particle correlations. One widely used technique is Hanbury-Brown-Twiss interferometry for which the correlations between particles are due to quantum statistics. Another method is through the measured yields of light nuclei, which are formed by the coalescence of individual nucleons.
Because of the violence of heavy ion collisions, it is highly improbable for a nuclear cluster near center-of-mass rapidity $`y_{cm}`$=1.6 in a collision at these energies to be a fragment of the beam or target nucleus . This would involve a cluster suffering a momentum loss of several GeV/c per nucleon that does not destroy the cluster, which is typically bound by only a few MeV per nucleon. These nuclei then are formed by coalescence and so represent correlations of several nucleons. As the mass of measured nuclei increases, of course, so does the number of particles involved in the correlation, and so does the sensitivity to features of the freeze-out distribution.
In part due to the fragility of these states, the observed light nuclei are believed to be formed only near freeze-out of the collision system, at which time the mean free path of a bound cluster is long enough for it to escape without further collision. It is this notion that gives rise to a class of models of light nuclei production, the coalescence models (for example, ). In general, these models assume a phase space distribution of nucleons at freeze-out and impose some coalescence conditions on the freeze-out positions and momenta of the nucleons in order to calculate the yields of nuclei. These models differ both in their assumptions about the phase space profiles and coalescence conditions. Their differences are often characterized by their predictions of the invariant coalescence, or $`B_A`$, parameters which are defined as
$$B_A\frac{(E\frac{d^3N_A}{dP^3})}{(E\frac{d^3N_{neutron}}{dp^3})^N(E\frac{d^3N_{proton}}{dp^3})^Z}$$
(1)
where a nucleus with baryon number $`A`$ and momentum $`P=pA`$ is formed out of $`Z`$ protons and $`N`$ neutrons.
In the simplest momentum space coalescence models, coalescence is assumed to take place between any nucleons with a small enough momentum difference. Early experimental results at the Bevalac with beam energies of $`500A`$MeV and high energy proton induced reactions revealed $`B_A`$ parameters which were approximately constant for these different collision systems ( contain useful summaries). The assumption of such simple models is that if the collision spatial volume is similar to size of the cluster, all nucleons whose momentum difference is less than a fixed value will fuse to form a nuclear cluster. Thus the experimentally observed constant $`B_A`$ values may indicate collision volumes in these systems that are not substantially larger than the RMS size of a deuteron.
In more advanced models, assuming a quantum mechanical sudden approximation and using a density matrix formalism, accounting for both the positions and momenta of the nucleons , the $`B_A`$ parameters take on a relationship with the source volume $`V`$ of $`B_A(1/V)^{(A1)}`$. Heavy-ion experimental results at higher energies at the AGS ($`10AGeV`$) and CERN ($`160AGeV`$) revealed $`B_A`$ values that decreased with beam energy. This observation was understood to be a sign of significant expansion in the collision volume before freeze-out. This larger source volume creates a situation where some nucleons with small relative momentum will have too large a spatial separation to coalesce, thus reducing $`B_A`$.
The density matrix formalism assumes that although the collision volume can expand significantly, there is no correlation between the momentum and the position of a given particle. This assumption then leads to a prediction of no kinematic dependence of the $`B_A`$ parameter. However, there is a great deal of evidence that collective motion is present leading to expansion of the collision volume and significant position-momentum correlations . Although the overall expansion of the system tends to decrease $`B_A`$ values by spatially isolating nucleons from each other, collective motion makes it more likely that nucleons that are spatially close together also have similar momenta, which to some extent works in the opposite direction by increasing $`B_A`$. Other coalescence models have made an effort to include the effects of both larger source volumes and collective motion, by including coalescence as a ’afterburner’ in collision cascade models such as RQMD and by analytical calculations .
Light nuclei production can also be calculated from thermal models which assume at least local thermal equilibrium of the system and thus that particle production (including, in this case, composite particles) is governed by a single temperature and chemical potentials. This gives rise to expressions for the $`B_A`$ with the same $`(1/V)^{(A1)}`$ dependence found in some coalescence models. Collective expansion can be included in these models ; it affects only the amount of energy available in local rest frames for particle production.
There has been much experimental effort in the measurement of light nuclei in heavy ion reactions. Previous measurements at AGS energies, including clusters of $`A4`$, show values for $`B_A`$ that are considerably lower than at Bevalac energies , indicating a much greater expansion of the system. This is also consistent with AGS results showing that the $`B_A`$ become smaller with more central collisions and larger target nuclei . At CERN-SPS energies, production of secondary particles (chiefly pions) is several times as large as at the AGS, leading to a larger expansion before hadronic freeze-out. Coalescence of high mass clusters is thus much less probable, as indicated by yet lower values for the $`B_A`$ as measured for example by experiment NA52 .
In this paper we describe and report the results of measurements by E864 of the yields of light nuclei in collisions of $`{}_{}{}^{197}Au`$ with beam momentum of 11.5 A GeV/c on targets of $`{}_{}{}^{208}Pb`$ and $`{}_{}{}^{197}Pt`$. The yields are measured for nuclei from baryon number A=1-7. In Section II we briefly describe the experimental apparatus and the analyses used to produce our final invariant multiplicities. In Section III we report the results and compare them with measurements of other experiments where such measurements overlap. Finally in Section IV we examine trends in the data and discuss interpretations in the context of several different models of light nuclei production.
## II Experiment 864
### A Apparatus
Brookhaven AGS Experiment 864 is an open geometry, high data rate spectrometer which was chiefly designed to search for rarely produced objects in Au+Pb collisions. Figure 1 shows a schematic view of the experimental apparatus, a thorough description of which is given in . Event centrality (impact parameter) is characterized by the charged particle multiplicity measured in an annular scintillator counter located approximately 10 cm downstream of the target, which subtends an angular region from 16.6<sup>o</sup> to 45<sup>o</sup> in azimuth when viewed from the target. The products of an interaction travel downstream through two dipole spectrometer magnets, M1 and M2. Charged particle identification is performed using information from the scintillator hodoscope walls (H1, H2, and H3) and the straw tube stations (S2 and S3). The hodoscopes walls each consist of 206 1cm thick scintillator slats placed vertically. They provide information about the charge, time-of-flight, and position of each charged particle hit, and this information is used to identify candidate charged particle tracks. These tracks are then rejected or confirmed and further refined by spatial hit information provided by the straw tube stations. Under the assumption that the track originates in the target, a rigidity is assigned to the track by a look-up table generated from a GEANT simulation of the experimental apparatus (using a technique described in as applied for the PHENIX experiment at RHIC). With information on rigidity, time-of-flight, and charge, a mass can then be assigned to the track, providing particle identification. A typical charge one mass distribution with a field of 0.45 Tesla in our spectrometer magnets is shown in Figure 2. Mass resolutions of 3 to 7% RMS are typical for particles with velocity $`\beta .985`$.
At the downstream end of the apparatus is our hadronic calorimeter ; an array of 754 towers, each measuring 10cm x 10cm on the front face, each of which provides energy and time information. This is the essential piece of the apparatus in our analyses of yields of neutral particles, for which the tracking detectors serve only to provide a veto. In addition, the energy and timing information from each tower is used to provide the input for a level II high mass trigger (the Late-Energy Trigger or LET ), which provides an enhancement of approximately a factor of 50 in our searches for rare high mass states.
Measurements reported in this paper are from a variety of experimental conditions, including different trigger conditions and different magnetic field settings in M1 and M2. The different data sets and their experimental conditions are listed in Table I. Because of the large acceptance open geometry design of the experiment, the different data sets often have significant regions of overlap with one another, allowing a consistency check on the measurements-see for example the measurements of alpha particles introduced in Section III.
### B Data Analysis
In order to measure the yield of a given species, a mass plot analogous to Figure 2 is made for each kinematic bin. The number of tracks which lie within its mass peak are determined. Background under each peak is then estimated, generally with fits to signal plus background, and subtracted away from this count. The invariant multiplicity in a given kinematic bin is then determined by correcting this number of raw counts for the geometric acceptance, trigger efficiency, charge cut efficiency, track quality cut efficiency, and detector efficiency.
Geometric acceptances are generally 25% or lower as shown in Figure 3 for protons. Charge cut and track quality cut efficiencies are typically 90 to 95%. The three redundant charge measurements for each track allow easy calculation of the charge cut efficiency in each hodoscope simply by examining charge measurements made in each hodoscope against results from the other two; charge misidentification in each of the three hodoscopes is less than one percent, so fewer than one track in a million is assigned an incorrect charge. In general, the track quality cuts are determined by comparison with monte-carlo simulations. In cases where the efficiencies are particularly high, they are determined directly from the data. The total track detector efficiency is 85 to 90%; this is determined by excluding each detector in turn from the track-finding process.
The trigger efficiency is significant only for those measurements which were made using the LET. This mass and momentum dependent efficiency ranges from approximately 40% to 90% for the measurements reported here. The LET efficiency is determined by two methods. The first, used mainly for slow, high mass states for which the efficiency is very high, is done simply using a monte-carlo simulation of the shower generated by the object and knowledge of the LET look-up table in each tower. The second method is from the equation $`ϵ_{LET}=N_{LET}/(N_{LET}+R\times N_{nonLET})`$ where R is the rejection factor provided by the LET (i.e. number of events which fire the LET divided by the total number of events). $`N_{LET}`$ and $`N_{nonLET}`$ are the numbers of particles of interest which do and do not produce LET triggers in LET triggered events, respectively.
Sources of systematic error that we have quantified include possible error in the determination of the efficiencies listed above as well as error in background subtraction (particularly relevant for the deuterons and alpha particles at their highest rapidities). Also examined were the effects of changing the assumed input distribution for each particle species in the determination of geometrical acceptances and efficiencies and effects of possible differences in the magnetic field with the field maps that were used in reconstruction of tracks.
Overall errors are generally dominated by systematics, particularly for the lighter states. Statistical errors can be significant for the heavier states, particularly in the determination of LET efficiencies in which the number of particles of interest which do not fire the LET generally has the largest statistical error.
## III results
Measurements of invariant multiplicities for protons, neutrons, deuterons, $`{}_{}{}^{3}He`$, and $`{}_{}{}^{4}He`$ are shown in Figures 4 through 9, and the values of the data shown in these figures are listed in tables in the Appendix. Figure 4 displays proton invariant multiplicities for three different bins of collision centrality. Figure 5 shows only proton yields from 10% most central Au+Pb collisions along with neutron multiplicities measured by E864 (see Reference ) for comparison. Figures 6 and 7 display deuteron and $`{}_{}{}^{3}He`$ invariant multiplicities for the same centrality bins used for the proton measurements. Tritons, $`{}_{}{}^{4}He`$, $`{}_{}{}^{6}He`$, $`{}_{}{}^{6}Li`$, $`{}_{}{}^{7}Li`$, and $`{}_{}{}^{7}Be`$ are measured by E864 only in 10% most central collisions; yields for tritons and alpha particles are shown in Figure 8 and Figure 9, respectively, while yields for the heavier nuclei are listed in tables in the Appendix. Added detail concerning most of these measurements may be found in the Ph.D. theses listed in Reference .
### A Contributions from Hyperon Decays
For comparison to other experimental results and calculations of coalescence parameters, it is important to quantify the contribution to proton yields that is made by protons which come from decays of hyperon states, which to E864 are indistinguishable from primordial protons. There are three dominant hyperon decays which produce protons: $`\mathrm{\Lambda }p+\pi ^{}`$, $`\mathrm{\Sigma }^0\mathrm{\Lambda }+\gamma p+\pi ^{}+\gamma `$, and $`\mathrm{\Sigma }^+p+\pi ^0`$. The contributions from these decays were evaluated using a GEANT simulation of the experiment with an input distribution taken from measurements of E891 for the Lambda and an input distribution from RQMDv2.3 for the Sigma. From this simulation it was determined that protons from hyperon decays account for approximately 12% of the measured yields of protons with only a slight kinematic dependence. The proton and neutron yields from E864 shown in Figures 4 and 5 have not been corrected for hyperon feed-down (nor have the values listed in the appendix tables).
### B Comparisons with Other Experimental Results
Figure 10 shows a comparison of the light nuclei measurements from AGS experiments E864, E877 , and E878 . Because of the different beam momenta of the experiments (10.8 A GeV/c in E878), the yields are plotted versus beam normalized rapidity. The yields shown for E864 and E877 are average yields at approximately $`p_T/A`$ = 150 MeV/c, while the E878 yields are measurements at $`p_T0`$. Other caveats to the comparisons of the yields shown in Figure 10 are noted in the figure caption.
Proton yields measured by E864 are clearly higher than measurements by E878. Some of this difference can be attributed to protons which are feed-down from hyperon decay. The acceptance of E878 for these feed-down protons is only about 10% of what it is for primordial protons, while in E864 the two acceptances are nearly the same. When this difference is taken into account (see Section III A) and the E864 yields of primordial protons are lowered by approximately 12%, the results of the two experiments are different by approximately 25% at midrapidity; this lies within the range of systematic errors for the two experiments. At higher rapidities, the agreement is better.
Comparison of the three experiments’ measurements of deuterons and tritons shows close agreement and the measurements of E864 and E878 of alpha particle yields also agree within errors.
## IV Discussion
### A General Trends of the Spectra
#### 1 Transverse Dependence of Yields
E864 has sufficient coverage in transverse momentum for us to extract measurements of inverse slope parameters of the different light nuclei yields in the rapidity range $`2.2y2.4`$. Shown in Figure 11 are yields of protons, deuterons, $`{}_{}{}^{3}He`$, and $`{}_{}{}^{4}He`$ as a function of transverse mass $`m_Tm_0=\sqrt{p_T^2+m_0^2}m_0`$ in this rapidity slice. Overlayed on the measurements are fits of each species to a Boltzmann distribution in transverse mass, from each of which we extract an inverse slope parameter, $`T`$, as noted in Figure 11. The fits from which we extract $`T`$ are linear fits to the log of the invariant multiplicities divided by $`m_T`$; this is the same fitting method used for determining the slope parameters for neutrons in . The $`\chi ^2`$ values are less than one per degree of freedom for all these fits.
These slope parameters, as well as those for neutrons in this same rapidity range, are displayed in Figure 12 as a function of mass number. Polleri et. al. have demonstrated the sensitivity of these trends in the inverse slopes to the density and velocity profiles of the nucleons at the time when coalescence occurs. To make this point, they have performed calculations of the behaviour of these trends for different assumptions about the source distributions. Two of these assumptions give rise to the two curves shown overlayed on the data in Figure 12. The first has a ’box’ spatial profile and a linear velocity profile, and the second has a Gaussian spatial distribution and a velocity profile $`v(r_T)r_T^{(1/2)}`$. In determining the curves shown in Figure 12 we have fit the data to these two different functional forms with the numerical constraint of T=100 MeV for zero mass. Neither of these sets of model parameters provides an adequate description of the data. These curves are meant only to illustrate the sensitivity of these measurements: clearly, the sensitivity to the differences in these assumptions increases with increasing mass of the measured nuclei.
Shown in Figure 13 are these same trends in light nuclei slope parameters including centralities other than the 10% most central collisions. For the more peripheral events, the trends are consistent with a linear dependence of slope parameter on mass (matching the calculation in Reference including box density profiles). Only for the most central events is there a clear rollover in the slope as a function of mass.
#### 2 Longitudinal Dependence of Yields
In order to examine the rapidity dependence of the yields of light nuclei, we observe the trends in multiplicities in the $`p_T/A`$ range from 100 to 200 MeV/c. (Our transverse coverage is not sufficient to integrate in transverse momentum for a full measurement of dN/dy over this entire rapidity range.) In Figure 14 we show the invariant multiplicities of deuterons for this low $`p_T`$ range for three different centralities as a function of rapidity. The yields are concave as a function of rapidity (i.e. they are lowest at center-of-mass and increase toward beam or target rapidity) and become more concave for the more peripheral events. We can parameterize this concavity by fitting each set of yields to a quadratic $`a+b(yy_{cm})^2`$. These fits are shown overlayed on the data in Figure 14. The ratio of coefficients, $`b/a`$, from this parameterization serves as a measure of the relative concavity of a species’ spectrum as a function of rapidity at low $`p_T`$ .
Values of this ratio $`b/a`$ are plotted in Figure 15 as a function of mass number $`A`$ for protons, deuterons and $`{}_{}{}^{3}He`$ in three different collision centralities. We observe that the concavity of a spectrum increases with the mass of a species and with collision centrality. This can be understood as some form of collective motion in the longitudinal direction, either expansion or incomplete stopping, which increases in peripheral events. It also is consistent with a rapidity dependent transverse expansion which pushes the nuclei out of the transverse momentum range measured here.
### B $`B_A`$ parameters
The trends noted in Section IV A will also be observable in a study of the behaviour of the $`B_A`$ parameters. For simplicity, we can from our neutron measurements characterize the neutron spectrum as a factor of 1.19$`\pm `$.08 greater than the proton spectrum (feed-down from hyperons is taken into account in determining this ratio), and so we evaluate $`B_A`$ as
$$B_A\frac{(E\frac{d^3N_A}{dP^3})}{(1.19)^N(E\frac{d^3N_{proton}}{dp^3})^A}.$$
(2)
for a nucleus of baryon number $`A`$ with $`N`$ neutrons and $`Z`$ protons. Again, all invariant yields are evaluated at a common velocity.
In Figures 16 and 17 we show measurements of $`B_2`$ and $`B_3`$ as a function of rapidity and transverse momentum. Because nucleons which are feed-down from weak decays are not available as nucleons for coalescence, we have subtracted the contribution to proton and neutron invariant yields from hyperon feed down for the calculation of $`B_A`$.
The values of the $`B_A`$ in Figures 16 and 17 are clearly not momentum independent as is assumed in many early coalescence models ( and , for example). This is expected given the large amount of evidence for collective flow in heavy-ion collisions . The values of the $`B_A`$ parameters seem to increase slightly with increasing transverse momentum and there is a clear increase away from center-of-mass rapidity. Both of these increases in $`B_A`$ away from center-of-mass momentum are consistent with expectations from an expanding source, although the longitudinal dependence can also be interpreted as a sign of incomplete stopping. Indeed the fact that the invariant yields of antiprotons near $`p_T`$=0 are strongly peaked near center-of-mass rapidity as measured by E864 while the proton yields are essentially flat may be taken as further evidence of incomplete stopping.
### C Source Size Calculations
The $`B_A`$ parameters can be related through coalescence models to source sizes. Following the model of Sato and Yazaki which uses a density matrix representation of the source distribution and projects it onto a representation of the deuteron wave function, we can relate the $`B_A`$ parameters to the RMS source radius through
$$B_A=(\frac{2J_A+1}{2^A})\frac{A^{5/2}}{m^{(A1)}}(4\pi \frac{\nu _A\nu }{\nu _A+\nu })^{(3/2)(A1)}$$
(3)
with $`\nu _A`$ the size parameter for a cluster with baryon number $`A`$ and spin $`J_A`$ ($`m`$ represents the nucleon mass). This model assumes the absence of collective motion of the nucleons, and therefore a radius (given by $`R_{RMS}=\sqrt{3/2\nu }`$) which is independent of momentum. We evaluate the source size using $`B_2`$, $`B_3`$, $`B_4`$, $`B_6`$, and $`B_7`$ values from our measurements nearest the collision center-of-mass (again assuming a neutron to proton ratio of 1.19) and list the results in Table II. We have used values for the $`\nu _A`$ from for $`A`$ = 2,3 and 4 and following have done a polynomial extrapolation to determine $`\nu _6`$ and $`\nu _7`$. This extrapolation for $`\nu _A`$ may be suspect particularly for the halo nucleus $`{}_{}{}^{6}He`$, but the final value for $`R_{RMS}`$ is quite insensitive to the value for $`\nu _6`$. The extracted radius for deuterons $`A`$=2 is considerably larger than the initial size of the colliding nuclei, and the radius parameters decrease for clusters of increasing mass.
A model by Scheibl and Heinz which includes the effect of collective flow in a density matrix prescription for coalescence leads to an expression for source dimensions:
$$B_2=\frac{3\pi ^{3/2}<C_d>}{2m_TR_{}^2(m_T)R_{}(m_T)}e^{2(m_Tm)(\frac{1}{T_p^{}}\frac{1}{T_d^{}})}$$
(4)
where $`R_{}`$ and $`R_{}`$ are the transverse and longitudinal dimensions of the fraction of the source which contributes to deuteron emission (comparable to the radius parameters extracted in the YKP parameterization of HBT interferometry), $`<C_d>`$ is a quantum mechanical correction factor for the finite size of the deuteron which is evaluated by the authors under various assumptions about the source, and $`T_p^{}`$ and $`T_d^{}`$ are the inverse slope parameters for protons and deuterons. Equation 4 as written assumes a box density profile for the source; a Gaussian profile would result in the absence of the final exponential factor. Plugging in our results for $`B_2`$, we can extract values for $`(R_{}^2(m_T)R_{}(m_T))^{1/3}`$ which are shown in Figure 18 as the solid circles. For the calculations shown here we have used 0.75 for $`<C_d>`$ and the values for $`T_d^{}`$ and $`T_p^{}`$ as measured at rapidity 2.3.
Shown also in Figure 18 as hollow circles are the results of source size calculations with the similar fragment coalescence model of Llope et. al. via the equation
$$R_c^3=\pi ^{3/2}\frac{(2J_c+1)}{(2J_a+1)(2J_b+1)}\frac{m_c}{m_am_b}\frac{E_a\frac{d^3N_a}{dp_a^3}E_b\frac{d^3N_b}{dp_b^3}}{E_c\frac{d^3N_c}{dp_c^3}}$$
(5)
which relates the effective source radius $`R`$ in the frame of a cluster $`c`$ which may be formed through the coalescence of smaller clusters $`a`$ and $`b`$. (Note that the radius parameters $`R`$ shown in Figure 18 are meant to describe sources of the form $`\rho (r)exp(r^2/2R^2)`$ and so correspond to RMS radii of $`R_{RMS}=\sqrt{3}R`$.)
We observe in Figure 18 that there is a decrease in source size with increasing distance away from the center-of-mass (again, as expected for a source with radial expansion), but our measurements do not give a clear picture concerning scaling of the source size with transverse mass such as has been noted in results for sizes from HBT two-particle correlations and in measurements of $`B_2`$ at the CERN SPS .
### D Comparison with RQMD
We can also compare our results with predictions from the cascade model RQMD version 2.3. The complex many-body processes by which light nuclei are formed are not included in RQMD, rather an afterburner is used to calculate the coalescence of these states based upon an input of the positions and momenta of the nucleons at freeze-out. This model then explicitly includes the position-momentum correlations due to expansion, etc. that are present in RQMD.
Figure 19 displays the E864 measurements along with predictions from RQMD with the coalescence afterburner for comparison over the transverse momentum range $`0.1p_T/A0.2`$ GeV/c as a function of rapidity. For comparison, RQMDv2.3 was run under two different conditions, one including the effect of repulsive mean-field potentials (potentials mode) and one not (cascade mode). We note that there is an increasing disagreement with increasing mass as noted previously in Reference . Calculations of $`{}_{}{}^{4}He`$ production in cascade mode are at a level of approximately 100 lower than our measurements; with potentials on, the discrepancy is still larger by about a factor of two. The level of disagreement varies somewhat in with transverse momentum (the slopes are in fact generally better predicted with potentials on than off) but particularly for the heavier states the change is slight compared to the overall level of disagreement.
### E Scaling of Yields Versus Mass
Shown in Figure 20 are the invariant yields in a small kinematic region at or near $`y=1.9,p_T300MeV`$. Over ten orders of magnitude, the yields in this kinematic bin fit very closely to an exponential dependence with a penalty factor of approximately 48 for each nucleon added (see and references therein for a discussion of such exponential behavior of cluster yields at lower collision energies). Of course, as we have seen above, the yields for different species have different kinematic dependences due to collective motion, and so this should not be taken as a penalty factor which governs the integrated yields of the various species, which from the kinematic dependences of the $`B_A`$ parameters discussed in Section IV B can clearly be different.
We can make a crude estimate of the difference in these two penalty factors by parameterizing the rapidity and mass dependences of the inverse slope parameters as follows: from our neutron measurements which extend up to beam rapidity, we can roughly parameterize the distribution of T in rapidity as a Gaussian with a width of 1.1 units. From Figure 13, we also can make a rough parameterization of the mass dependence of the inverse slope as $`T(A+1.0)`$ With these two parameterizations and our penalty factor of 48 at $`y=1.9,p_T300MeV`$, we calculate a penalty factor of 25 in the integrated yields for the addition of an extra baryon to a coalesced state. This should be considered as a lower limit on such a penalty factor since as previously noted the inverse slope parameter rolls over as a function of mass rather than following the linear dependence we have used in this estimate .
We can also estimate this penalty factor in overall yields by following section VI of Reference . Equation 6.10 in this reference allows us effectively to make an estimate of the difference between the penalty factor near $`p_T`$ = 0 and the overall penalty factor in two limiting cases: a static, homogeneous fireball (which would translate our penalty factor of 48 at $`p_T`$ =0 to an overall penalty factor of 72) and a rapidly expanding system (which would result in an overall penalty factor of 39).
### F Spin and Isospin Dependences of Yields
In Figure 21, we display three ratios as a function of transverse momentum in the rapidity ranges y=1.8-2.0 and y=2.0-2.2. The three ratios are the ratios of invariant yield of neutrons over invariant yield of protons, the ratio of yield of tritons over yield of $`{}_{}{}^{3}He`$ and the ratio of $`{}_{}{}^{6}He`$ to $`{}_{}{}^{6}Li`$. The $`n/p`$ and $`t/^3He`$ ratios are consistent with a value of approximately 1.2 (in Reference , we extract values of $`1.19\pm .08`$ and $`1.23\pm .04`$, respectively, for the two ratios in the range $`1.6y2.4`$.). In contrast, the $`{}_{}{}^{6}He/^6Li`$ is much nearer to a value of 0.3. $`p`$, $`n`$, $`{}_{}{}^{3}He`$ and $`t`$ are all spin J=1/2 states while $`{}_{}{}^{6}He`$ is spin 0 and $`{}_{}{}^{6}Li`$ is spin 1. We take this as evidence that the yields scale as the degeneracy factor 2J+1 which is commonly predicted in thermal and coalescence models.
With the dependences upon mass number, isospin and spin divided away one can examine the yields for other dependences. This topic, including a possible dependence on binding energy per nucleon with an inverse slope parameter dependence of a few MeV, is discussed in Reference .
## V Summary
We have shown results of measurements of light nuclei from $`A`$=1 to $`A`$=7. The increase with mass of light nuclei inverse slope parameters appears to roll over at approximately $`A`$=3 in central events but not necessarily in more peripheral events. Also we have shown that the yields near $`p_T`$=0 are concave as a function of rapidity and that this relative concavity increases in more peripheral events and in higher mass nuclei, consistent with both radial expansion and incomplete stopping. These trends are also evident from our observations of the kinematic dependences of the $`B_A`$ parameters. From these parameters we have extracted source dimensions from various models. Efforts to extract more quantitative information about the source from these measurements using the cascade model RQMD with a coalescence afterburner were unsuccessful as predictions of the model differ from our results by an amount that increases with mass and reaches a level of 100 or more by $`A`$=4.
We have also examined the overall scaling of the yields up to $`A`$=7, extracting a penalty factor of about 48 to add a nucleon to a coalesced state near midrapidity at low transverse momentum. This likely translates into a somewhat smaller penalty factor in overall yields for the addition of a nucleon, but we have argued that this is unlikely to differ by as much as a factor of two from our measured penalty of 48.
## VI Acknowledgements
We gratefully acknowledge the efforts of the AGS staff in providing the beam. This work was supported in part by grants from the Department of Energy (DOE) High Energy Physics Division, the DOE Nuclear Physics Division, and the National Science Foundation.
## VII APPENDIX
Shown in Table III through Table XVI are results of measurements by E864 of invariant multiplicities of light nuclei from $`A`$=1 through $`A`$=7 in 10% most central Au+Pb collisions. Results for less central data are also listed for protons, deuterons and $`{}_{}{}^{3}He`$.
|
warning/0003/astro-ph0003071.html
|
ar5iv
|
text
|
# Stellar parameters from very low resolution spectra and medium band filters
## 1 Background and Objectives
An understanding of the origin, properties and evolution of our Galaxy requires a careful census of its constituents, in particular its stellar members. Of special importance are the intrinsic physical properties of these stars. The fundamental properties are mass, age and abundances, as these determine a star’s history and future development. However, ages are not observable, and masses can only be directly obtained from some multiple systems. Thus we must indirectly gain this information via the stellar spectrum, and a number of atmospheric parameters have been defined for this purpose. The main ones are the effective temperature, $`\mathrm{T}_{\mathrm{eff}}`$, the surface gravity, $`\mathrm{log}g`$, and the metallicity, \[M/H\]. To these can also be added the alpha abundances, {$`\alpha _i`$} (which measure the devations away from the ‘standard’ abundance ratios), the photospheric microturbulence velocity, $`V_{\mathrm{micro}}`$, and the extinction by the interstellar medium, A($`\lambda )`$(although not intrinsic to the star, it is necessary for determining its luminosity). Masses and ages can then be determined from stellar structure and evolution models and with calibration via binary systems. It is important to realise that this modelling is complex, and a number of assumptions have to be made. There is, therefore, a limit to the precision with which we can determine physical properties.
Historically, spectroscopic parameters have been measured indirectly through the MK classification system (Morgan et al. morgan\_43a (1943)) or via colour-magnitude and colour-colour diagrams. In the MK system, the two parameters spectral type and luminosity class act as proxies for $`\mathrm{T}_{\mathrm{eff}}`$ and $`\mathrm{log}g`$. Originally a qualitative system relying on a visual match between observed spectra and a system of standards, much progress has been made in quantifying it with automated techniques (e.g. Weaver & Torres-Dodgen weaver\_97a (1997); Bailer-Jones et al. bailerjones\_98a (1998)). The most commonly used classification techniques have been neural networks and $`\chi ^2`$ matching to templates (or more generally, minimum distance methods). A summary of recent progress in this area is given by von Hippel & Bailer-Jones (vonhippel\_00a (2000)).
Despite this focus on the MK system, it is not well suited to classifying data from the deep surveys which will be central to the future development of Galactic astrophysics. This is for a number of reasons, but in particular because it lacks a measure of metallicity. Although MK does make allowance for various ‘peculiar’ stars, these are defined as exceptions, and the notation is not suited to a statistical, quantifiable analysis. This is problematic given the significance of metal poor halo stars in a deep survey. There is also now no good reason why we should not determine physical parameters directly from the observational data.
Some attempts have been made to determine the physical parameters of real spectra directly by training neural networks on synthetic spectra. Gulati et al. (gulati\_97a (1997)) used this approach to determine the effective temperatures of ten solar metallicity G and K dwarfs. Taking the “true” effective temperature of these stars as those given by Gray & Corbally gray\_94a (1994), they found a mean “error” in the network-assigned temperatures of 125 K. Bailer-Jones et al. (bailerjones\_97a (1997)) determined $`\mathrm{T}_{\mathrm{eff}}`$ for over 5000 dwarfs and giants in the range B5–K5, and also showed evidence of sensitivity of the parametrization models to metallicity.
The accuracy with which physical parameters can be determined from a stellar spectrum depends upon, amongst other things, the wavelength coverage, spectral resolution and signal-to-noise ratio (SNR). From the point of view of designing a stellar survey project it is essential to know how well the stellar parameters can be determined for a given set of these observational parameters. Moreover, given that there is always a limit to the collecting area and integration time available, there is always a trade-off between spectral resolution, sensitivity and sky coverage.
The goal of this paper is to determine the accuracy with which physical stellar parameters can be determined from spectroscopic data at a range of SNRs and resolutions which could realistically be achieved in a deep survey mission. This specification rules out high resolution spectra. The parametrization work has been carried out using neural networks (Sect. 3) because they have been shown to be one of the best approaches for this kind of work. This is not to presuppose, however, that some other approach may not ultimately be better. The simulations have been made using a large database of synthetic spectra generated from Kurucz atmospheric models (Sect. 4). While these spectra do not show the full range of variation in real stellar spectra, they are adequate for a realistic demonstration of what is possible as a function of SNR and resolution. The results are presented in Sect. 5 and summarised and discussed in Sect. 6. Finally, the requirements for a complete survey-oriented classification system are given in Sect. 7.
## 2 The GAIA Galactic Survey Mission
The simulations in this paper were partially inspired by the need to produce an optimal photometric/spectroscopic system for the GAIA Galactic survey mission. GAIA is a candidate for the ESA cornerstone 5 mission for launch in 2009 (ESA, in preparation). It is primarily an astrometric mission with a precision of a few microarcseconds, and will survey the entire sky down to V=20, thus observing c. 10<sup>9</sup> stars in our Galaxy. Radial velocities will be obtained on board down to V=17.5, thus providing a 6D phase space survey (three spatial and three velocity co-ordinates) for stars brighter than this limit. A survey of this size will have a profound impact on Galactic astrophysics, but to achieve this it is essential that the physical characteristics of the target objects are measured and correlated with their spatial and kinematic properties. As GAIA is a continuously scanning satellite, a fixed total amount of integration time is available for each object, so there is a trade-off between resolution, signal-to-noise ratio and wavelength coverage. For various reasons, the current GAIA design does not include a spectrograph (other than a 1.5 Å resolution region between 8470 and 8700 Å intended for radial velocity measurements), but instead will image all objects in several medium and broad band filters (Table 1). Three filter systems are shown: the system nominally selected for the mission plus two alternatives. The profiles of the two alternatives are represented as Gaussians in this paper. The ability of these filter system to determine stellar parameters will be compared with that for spectra of various resolutions.
## 3 The network model
A neural network is an algorithm which performs a non-linear parametrized mapping between an input vector, $`𝐱`$, and an output vector, $`𝐲`$. (The term ‘neural’ is misleading: although originally developed to be very simplified models of brain function, many neural networks have nothing to do with brain research and are better described in purely mathematical terms.) The network used in this paper is a feedforward multilayer perceptron with two ‘hidden layers’. These hidden layers form non-linear combinations of their inputs. The output from the first hidden layer is the vector $`𝐩`$, the elements of which are given by
$$p_j=\mathrm{tanh}\left(\underset{i}{}w_{i,j}x_i\right)$$
These values are then passed through a second hidden layer which performs a similar mapping, the output from that layer being the vector $`𝐪`$
$$q_k=\mathrm{tanh}\left(\underset{j}{}w_{j,k}p_j\right)$$
The output from the network, $`𝐲`$, is then the weighted sum of these
$$y_l=\underset{k}{}w_{k,l}q_k$$
The $`\mathrm{tanh}`$ function provides the non-linear capability of the network, and the weights, $`𝐰`$, are its free parameters. The model is supervised, which means that in order for it to give the required input–output mapping it must be trained on a set of representative data patterns. These are inputs (stellar spectra) for which the true target outputs (stellar parameters) are known. The training is a numerical least-squares minimisation: Starting with random values for the weights, a set of spectra are fed through the network and the error in the actual outputs with respect to the desired (target) outputs calculated. The gradient of this error with respect to each of the N weights is then used to iteratively perturb the weights towards a minimum of the error function. Thus the training is a minimisation problem in an N-dimensional space, and the resulting input–output mapping can be regarded as a non-linear interpolation of the training data. Once the network has been trained the weights are fixed and the network used to obtain physical stellar parameters for new spectra.
The results in this paper use a network code written by the author consisting of five and ten hidden nodes in the first and second hidden layers respectively. The complexity of the network is determined by the number of hidden nodes and layers. While networks with a single hidden layer can provide non-linear mappings, experience has shown that a second hidden layer can lead to considerable improvement in performance (Bailer-Jones et al. bailerjones\_98a (1998)). This has been confirmed with the data in this paper. Significant further improvement is not expected through the addition of more hidden nodes/layers. The network has three outputs, one for each of the parameters $`\mathrm{T}_{\mathrm{eff}}`$, $`\mathrm{log}g`$ and \[M/H\]. The error which is minimised is the commonly-used sum-of-squares error (the sum being over all training patterns and outputs), except that the error contribution from each output is weighted by a factor related to the precision with which that parameter can be determined.
I stress that a neural network is not fundamentally different from many other parameter fitting algorithms. Its strengths are that it has a fast and straight-forward training algorithm, can map arbitrarily complex functions (given sufficient data to determine the function), and can be parallelised in software or hardware to achieve considerable increases in speed. One of the common criticisms of neural networks is that it is difficult to interpret their weights and get an idea of exactly how they achieve their results. While this is essentially true, part of this difficulty stems from the fact that the models are problem-independent: they are purely mathematical models that do not explicitly take into account the physics of the problem. Moreover, in order to fully understand the model it would be necessary to simplify it, and this in turn would reduce its performance. This “interpretability–complexity” trade-off is inherent to almost any type of heuristic model.
## 4 Synthetic spectra
A large grid of synthetic spectra have been generated using Kurucz atmospheric models (Kurucz kurucz\_92a (1992)) and the synthetic spectral generation program of Gray (Gray & Corbally gray\_94a (1994)). The parameter grid consists of 36 $`\mathrm{T}_{\mathrm{eff}}`$ values between 4000 K and 30 000 K (step sizes between 250 K and 5000 K), 7 values of $`\mathrm{log}g`$ between 2.0 and 5.0 dex (in 0.5 steps) and 15 values of \[M/H\] between $``$3.0 and $`+`$1.0 dex (step sizes between 0.1 and 0.5). The microturbulence velocity was fixed at 2.0 kms<sup>-1</sup>. This yielded an (almost complete) grid of 3537 atmospheric models. Contiguous spectra were calculated between 3000 and 10 000 Å in 0.05 Å steps with a line list of over 900 000 atomic and molecular lines. The resolution, $`r`$, of these spectra was then degraded to 25, 50, 100, 200 and 400 Å FWHM by Gaussian convolution. (Each resolution element is sampled by two pixels, so these resolutions correspond to 560, 280, 140, 70 and 35 inputs to the network respectively.) These resolutions are considerably lower then the 1–5 Å generally used for MK classification. The spectra were also combined with the transmission curves of the filters (Table 1) to produce three sets of filter fluxes. Poisson noise was added to all data sets to simulate signal-to-noise ratios of 5, 10, 20, 50 and 1000 per resolution element. The result is 40 sets of 3537 absolute spectral energy distributions at each combination of resolution and SNR. The absolute flux information is retained.
It is noted that Kurucz models do not produce highly accurate spectra for all types of stars. This is particularly true at low $`\mathrm{T}_{\mathrm{eff}}`$ as they exclude water opacity (and there are no H<sub>2</sub>O lines in the line lists). For this reason spectra have not been calculated below 4000 K. Furthermore, the models lack chromospheres and so do not reproduce features such as emission in the cores of the CaII H & K absorption lines. For the present investigation, however, it is not necessary to have highly accurate individual spectra, but spectra which reflect differences of the appropriate scale and complexity.
## 5 Spectral parametrization results
As the neural network is a parameter fitting algorithm, it is essential that its performance is evaluated on an independent set of data from that on which it is trained. For this purpose, each of the 40 data sets was randomly split into two halves and one used for training (1760 spectra) and the other for testing (1759 spectra). $`\mathrm{log}_{10}`$$`\mathrm{T}_{\mathrm{eff}}`$ (rather than $`\mathrm{T}_{\mathrm{eff}}`$) is used as a target in the networks to reduce the dynamic range of this parameter and give a better representation of the uncertainties. The input and output parameters are scaled to have zero mean and unit standard deviation to prevent ‘saturation’ of the network during training.
For each data set a committee of three identical networks was trained from different initial random weights. The resultant parameter for any star is then the average from the three networks. This helps to reduce the effects of imperfect training convergence. Each network was trained with a conjugate gradient algorithm for 10 000 iterations and used weight decay regularisation to avoid overtraining. More training did not reduce the error further. The longest training time (for the largest input vector) was about one day on a Sun SPARC Enterprise (no parallelisation of the code). The time to parametrize was of order $`10^3`$ seconds per spectrum.
The precision to which physical parameters can be determined from a stellar spectrum depends not only on the SNR and resolution, but also on the type of star. For example, it is more difficult to determine the metallicity of hot stars on account of the almost complete absence of metal lines. Therefore, I summarise the performance of each data set for three different temperature ranges (for all $`\mathrm{log}g`$ and \[M/H\]):
1. $`\mathrm{T}_{\mathrm{eff}}`$ $`<`$ 5800 K (stars later than solar – 408 spectra in the test subset)
2. 5800 $`<`$ $`\mathrm{T}_{\mathrm{eff}}`$ $`<`$ 10 000 (A and F stars – 888 spectra in the test subset)
3. $`\mathrm{T}_{\mathrm{eff}}`$ $`>`$ 10 000 K (O and B stars – 463 spectra in the test subset)
The error measure I used is the average absolute error, $`ϵ`$, of each parameter, i.e. the absolute difference between the network output and the target value averaged over all stars in the test subset for that temperature range. This error is more robust than the often-used RMS error because it is less distorted by outliers and more characteristic of the majority of the error distribution. For a Gaussian distribution $`1\sigma =1.25ϵ`$, although some of the error distributions deviate significantly from Gaussian.
The results of the parametrization process are shown in Figs 13 and tabulated in Tables 24. Before interpreting these results we should consider the limits which the data themselves place on the performance. First, the network will be unable to produce errors smaller than the smallest variations in the data set. If, to take a hypothetical example, the spectra did not change as the metallicity changed by 1.0 dex, we could not expect the network to determine \[M/H\] to much better than 0.5 dex. Second, the grid of atmospheric models represents the physical parameters at a finite sampling, e.g. a constant step size of 0.5 dex for $`\mathrm{log}g`$. This sampling does not in itself limit the precision achievable; it is perfectly possible for the network to legitimately give an error much smaller than the sampling because the network is minimising a continuous error function and not just obtaining the best match between a spectrum and a set of templates. Nonetheless, the network input–output mapping is an interpolation of the training data, and the more coarsely sampled the parameter grid the harder it is for the network to get a reliable interpolation. Consequently, while the network may be able to achieve sub-sampling accuracy, we should not be surprised if it cannot. Thus to avoid over-interpreting these results we should not compare two errors which are both smaller than half the sampling level. The average ‘half-sampling’ values for \[M/H\] and $`\mathrm{log}g`$ are 0.2 and 0.25 respectively, and for log $`\mathrm{T}_{\mathrm{eff}}`$ in the three temperature ranges (cool, intermediate and hot) are 0.01, 0.01 and 0.03 respectively. The implication is that, if the network produces errors smaller than these half-sampling values (as it does), we cannot know whether the performance is limited by the network model or by the data themselves. A distinction will only be possible with a more sensitive and finely sampled grid of atmospheric models.
With the above caveat taken into account, I draw attention to some interesting features in Figs 13.
1. Good $`\mathrm{T}_{\mathrm{eff}}`$ determination is possible with all resolutions/filter systems and SNRs. The larger error in $`\mathrm{T}_{\mathrm{eff}}`$ above 10 000 K may be an artifact of the larger half-sampling value in this region ($``$1000 K).
2. Only at high resolution can $`\mathrm{log}g`$ be determined for the coolest stars and even then the determination is poor relative to the hotter stars. This is understandable, at least in part, because the $`\mathrm{log}g`$ spectral signature is primarily in the line widths which are only apparent at high resolution.
3. Although the three filter systems differ somewhat, they give essentially the same performance as each other.
4. The filter systems (each with 10–15 input parameters) have similar $`\mathrm{log}g`$ and $`\mathrm{T}_{\mathrm{eff}}`$ as the $`r`$=400 Å spectra (35 inputs).
5. At low SNR, the $`r`$=400 Å spectra and the filters give poor \[M/H\] and very poor $`\mathrm{log}g`$ determination for all three temperature ranges.
6. At high SNR (1000) all resolutions/filter systems appear to be equally good at determining any of the parameters. Differences will probably become apparent with a more sensitive training grid.
7. At higher temperatures the accuracy is more sensitive to SNR than at lower temperatures.
8. Metallicity determination is more difficult at higher temperatures, especially for the filters and low resolution spectra. This is understandable as at high temperature there are fewer and weaker metal lines which are only significant at high SNR and/or resolution.
9. In most cases there is little difference between the performances of the $`r`$=25, 50 and 100 Å spectra, at least for this data grid.
## 6 Summary and Discussion
The results demonstrate that a fully automated neural network can accurately determine the three principal physical parameters from spectroscopic or photometric stellar data, something which has not previously been demonstrated. Moreover, this work has used spectra of considerably lower resolution than have been used before in automated classifiers. Even at low resolution (50–100 Å FWHM) and SNR (5–10 per resolution element), neural networks can yield good determinations of $`\mathrm{T}_{\mathrm{eff}}`$ and \[M/H\], and even for $`\mathrm{log}g`$ for stars earlier than solar. Still lower resolutions permit good results provided the SNR is high enough ($`50`$). That good $`\mathrm{T}_{\mathrm{eff}}`$ can be achieved even at low resolution and SNR is perhaps not surprising when we consider that the spectra have absolute fluxes, which will be the case with high precision parallax missions such as GAIA. However, the more distant objects will have lower precision parallaxes and hence errors in the mean flux level. But even if we completely ignore distance information (and flux normalise the spectra), the shape of the spectrum is still a strong indicator of $`\mathrm{T}_{\mathrm{eff}}`$: For example, Bailer-Jones et al. (bailerjones\_98a (1998)) obtained an MK spectral type precision of 0.8 subtypes ($`\mathrm{\Delta }\mathrm{log}`$$`\mathrm{T}_{\mathrm{eff}}`$=0.010–0.015) across a wide range of spectral types (B2–M7) using flux normalised spectra. This is similar to what can be achieved from broad band photometry, implying that $`\mathrm{T}_{\mathrm{eff}}`$ determination only requires very low resolution.
The good performance of ‘high’ resolution spectroscopy (25 Å) at very low SNR ($`\sqrt{5}`$ per pixel) was not expected. It seems to imply that for a given amount of integration time it may be better to sacrifice SNR for resolution. It is noteworthy that while the filters provide good $`\mathrm{T}_{\mathrm{eff}}`$, their ability to determine \[M/H\] and especially $`\mathrm{log}g`$ is very limited at low SNR.
How do these results compare with classical parametrization methods? Gray (gray\_92a (1992)) compiles results showing that with photometric errors below 0.01 magnitudes, the B$``$V colour calibrates $`\mathrm{T}_{\mathrm{eff}}`$ to 2–3% (4% for hotter stars) in the absence of reddening. Slightly better precision can be obtained from the slope of the Paschen continuum and size of the Balmer discontinuity. The latter may also be used to measure $`\mathrm{log}g`$ to $`\pm 0.2`$ dex. With spectra at a few Å resolution over a similar wavelength range to that used here, Cacciari et al. (cacciari\_87a (1987)) obtained uncertainties in log $`\mathrm{T}_{\mathrm{eff}}`$ and $`\mathrm{log}g`$ of 0.01 and 0.04 respectively. Sinnerstad (sinnerstad\_80a (1980)) made uvby,$`\beta `$ photometric measurements of B stars, and for uncertainties of 0.005 in $`\beta `$ and of 0.01 in u$``$b (i.e. SNR $``$ 200), infers errors in log $`\mathrm{T}_{\mathrm{eff}}`$ and $`\mathrm{log}g`$ of 0.004 and 0.08 respectively. These are similar to or slightly better than the results for similar stars in Tables 24 ($`ϵ_3`$) at the highest resolutions. High resolution ($`r0.1`$ Å) spectra have generally been used to determine metallicity, and in a review, Cayrel de Strobel (cayreldestrobel\_85a (1985)) notes that metallicity can be determined to $`\pm 0.07`$ dex at SNR=250 (but only $`\pm 0.2`$ dex at SNR=50) provided the effective temperature and gravity are approximately known. At lower SNR (10–20), Jones et al. (jones\_96a (1996)) could determine \[Fe/H\] to $`\pm 0.2`$ dex for G stars using a set of spectroscopic indices measured at 1 Å resolution in the range 4000–5000 Å, again using a known effective temperature.
More recently, Katz et al. (katz\_98a (1998)) have used a minimum distance method to parametrize spectra by finding the closest matching template spectrum. The template grid consisted of 211 flux calibrated spectra (3900–6800 Å, $`r0.1`$ Å) with 4000 K $``$$`\mathrm{T}_{\mathrm{eff}}`$$``$ 6300 K, $`0.29`$ \[Fe/H\] $`+0.35`$, and $`\mathrm{log}g`$ for dwarfs and giants. The internal accuracy of the method for log $`\mathrm{T}_{\mathrm{eff}}`$, $`\mathrm{log}g`$ and \[M/H\] was 0.008, 0.28 dex, and 0.16 dex respectively at SNR=100, and 0.009, 0.29 dex and 0.17 dex at SNR=10. As expected, their results for $`\mathrm{log}g`$ are much better than those in this paper at the similar temperature range ($`ϵ_1`$ in Table 3), presumably due to their much higher resolution. In contrast, their performance for \[M/H\] is similar and for $`\mathrm{T}_{\mathrm{eff}}`$ somewhat worse than that in this paper at 500 times lower resolution. Their results also confirm that at high resolution a lower SNR leads to very little loss in performance. Snider et al. (snider\_00a (2000)) trained and tested neural networks on a set of 182 real F,G and K spectra over the range 3630–4890 Å at intermediate resolution ($``$1Å), and achieved 1$`\sigma `$ errors in log $`\mathrm{T}_{\mathrm{eff}}`$, $`\mathrm{log}g`$ and \[M/H\] of 0.015, 0.41 dex and 0.22 dex respectively, based on training and testing a network with a set of 182 real F,G and K spectra.
When judging the relative values of the different resolution/SNR combinations in this paper, we must also take account of their implementation ‘costs’, specifically the relative integration times required. Usually for a survey, a fixed total amount of integration time is available for all filters/spectra. In the case of GAIA – which is continuously rotating – a star moves across a focal plane covered with a mosaic of CCDs which are clocked at the rotation rate. The different filters are fixed to different CCDs, so that as a star moves across the mosaic it is recorded in different wavelength ranges. Thus less numerous and/or broader filters would achieve a higher SNR than more or narrower filters. Some filters could be replaced with a slitless spectrograph (e.g. a prism or grism). This disperses every point on the sky and thus gives the full integration time for all wavelengths, but at the expense of increased sky noise and object confusion. These could be reduced by using one or more dichroics to redirect the light to two or more focal planes. (Confusion would be reduced further with GAIA by the fact that each area of sky is observed at many different position angles over the mission life.) An alternative approach is a set of many medium band filters ($`100`$ for $`r`$=100 Å over the complete wavelength range, although omission of some filters could be achieved). While this avoids the two principal disadvantages of the slitless spectrograph, the integration time per wavelength interval is dramatically reduced.
## 7 Development of a survey parametrization system
The development of a complete parametrization system will require further research, much of which needs to be directed at taking better account of the true nature of the observational data. Directions and suggestions for the course of this work are now given.
### 7.1 Object selection
Essentially all of the work in the literature on automated classification deals with preselected objects. In contrast, an unpointed survey will pick up a whole range of objects, necessitating a filtering system to select the stars. Such a system could make use of both object morphology and spectral features, and systems based on neural networks (e.g. Odewahn et al. odewahn\_93a (1993); Miller & Coe miller\_96a (1996); Serra-Ricart et al. serraricart\_96a (1996)) and Principal Components Analysis (Bailer-Jones et al. bailerjones\_98a (1998)) have been proposed. Such a system must be relatively robust and always allow for ‘unknown’ objects which can be dealt with manually.
### 7.2 Model training
It will be necessary to have a stellar database for training which takes better account of the larger range of variation present in the Galactic stellar population. Ideally, a large set of real spectra across a wide range of physical parameters should be obtained for this purpose. Good atmospheric models and synthetic spectra are nonetheless still required for determining their physical parameters and thus for training the network. There are two possible approaches to training. The first is to train on synthetic spectra suitably preprocessed to be in the same form as the observed spectra (e.g. Bailer-Jones et al. bailerjones\_97a (1997)). The alternative is to obtain a representative sample of real spectra with the survey system, calibrate them, and then use them to train a network. In theory the latter method gives a better sampling of the true cosmic variance in the spectra, but of course requires that a representative sample is selected from the survey data. This sample could be improved as the survey progressed. Neural networks are fast to train and apply, so it is realistic to expect that even for a database of $`10^9`$ objects the network could be retrained and applied to the whole database in less than a day.
### 7.3 Improved stellar models
More advanced model atmospheres are required for a number of reasons:
1. $`\mathrm{T}_{\mathrm{eff}}`$, \[M/H\] and $`\mathrm{log}g`$ do not uniquely describe a true spectrum. Models sensitive to different abundance ratios and which include chromospheres (for example) are necessary.
2. Kurucz models assume LTE which is known to break down in a number of regimes (e.g. for very hot stars).
3. Both the atmospheric models and the line lists lack water opacity, known to be important for cool stars, thus setting the current lower $`\mathrm{T}_{\mathrm{eff}}`$ limit of about 4000 K.
4. Yet more advanced models (which include dust) are required for very cool stars (L and T dwarfs) and brown dwarfs, of which many will be found by GAIA.
### 7.4 Reddening
Of particular importance is insterstellar extinction (reddening), especially in deep surveys. The extinction can, in theory, be determined by the network by training it on artificially reddened synthetic spectra and providing the network with a “reddening” output parameter (or parameters). This has been demonstrated on limited data sets by Weaver & Torres-Dodgen weaver\_95a (1995) and Gulati, Gupta & Singh gulati\_97a (1997), who determined E(B$``$V) to within 0.05 and 0.08 magnitudes respectively. The latter made use of 6 Å resolution UV spectra (4850–5380 Å). The former used red spectra (5800–8900 Å) at 15 Å resolution and found that the spectral type and luminosity class classifications did not degrade much as reddening was added. It is therefore to be expected that the parametrizations in this paper will be robust to reddening, particularly as the spectra have a much larger wavelength coverage. The filter systems proposed for GAIA were of course designed with interstellar extinction in mind, and a study of its impact has been carried out (ESA, in preparation). This work shows that suitable Q parameters (non-linear combinations of the filter fluxes) used to determine the physical parameters are largely insensitive to reddening. It also claims that narrow band filters are not necessary for overcoming reddening. In some parts of the parameter space, reddening is more problematic (e.g. for K stars), largely due to a degeneracy between it and $`\mathrm{T}_{\mathrm{eff}}`$ and $`\mathrm{log}g`$. However, at intermediate and high Galactic latitudes it is expected that E(B$``$V) can be determined to within 0.002 magnitudes. Munari (munari\_99a (1999)) similarly shows that reddening-free indices exist for the Asiago filter system. As a neural network also forms non-linear combinations of the filter fluxes, it is reasonable to suppose that it too will be robust to redenning, although this will be the subject of future work.
### 7.5 Binary systems
The parametrization model used in a real survey must confront the fact that most stars are in spatially unresolved multiple systems. Independent measurement of the physical properties of each component is desirable and in principle achievable – when the brightness ratio is large enough – by training the network with composite spectra. In this case the network model would need to have multiple sets of outputs to deal with each component. An alternative approach is to use ‘probabilistic outputs’ in which the single output for, say, $`\mathrm{T}_{\mathrm{eff}}`$, is replaced with a series of outputs in which each value of $`\mathrm{T}_{\mathrm{eff}}`$ (6000, 6250, 6500 etc.) is represented separately. The network then evaluates the probability that each temperature is present in the input spectrum. This method is not recommended, however, as it eliminates the intrinsically continuous nature of the physical parameters. It would also greatly increase the number of outputs and hence the number of free parameters (weights) in the network.
### 7.6 Incomplete data
Object confusion should not result in any overlapped spectrum being rejected entirely. Rather, it would be better to have a parametrization model which is robust to missing data. This is a major challenge for the feedforward network models used in this and most other papers on automated classification, and will presumably require some transformation of the input spectrum. An analysis of the effect of wavelength coverage on the parameter determination accuracy is important because a smaller spectral coverage (or coverages – it need not be contiguous) would also reduce this confusion.
Finally, the model should make use of all available data. In the case of GAIA, this means including the data from the high resolution spectrograph (8470–8700 Å at 0.75 Å$`/\mathrm{pix}^1`$) used to measure radial velocities. As the inputs to the network need not be homogenous, there should be no problem incorporating different types of data.
## Acknowledgements
I would like to thank Fabio Favata, Gerry Gilmore and Michael Perryman for useful discussions on this work, in particular within the context of the GAIA mission. I am also grateful to Robert Kurucz for use of his model atmospheres and compiled line lists, and to Richard Gray for the use of his synthetic spectra generation program.
|
warning/0003/astro-ph0003414.html
|
ar5iv
|
text
|
# Cosmological Surveys at Submillimetre Wavelengths
## 1. Evidence for Massive Star-Formation at High-Redshift
In addressing the question ‘what is the main epoch of metal production in the universe?’, or equivalently, ‘when did the cosmic star-formation rate reach its peak value?’, a number of separate lines of evidence suggest that a high-rate of star-formation ($`100M_{}\mathrm{yr}^1`$) must have occurred in massive systems at $`z3`$. This evidence includes (i) the demonstration by Renzini (1998), using clusters of galaxies as probes of the past star-formation and metal production history, that 30–50% of the present-day baryons are currently locked up in massive structures which appear to have formed at $`z>3`$; (ii) the peak in the co-moving number density of AGN (radio galaxies and quasars) at $`z2`$, AGN whose counterparts at low-redshift are hosted in luminous, massive elliptical galaxies ($`>2L^{}`$ \- Taylor et al. 1996, McClure et al. 1999). At $`z2`$ the universe is only 3–4 Gyrs old, which implies that a sustained star-formation rate (SFR) $`>200M_{}\mathrm{yr}^1`$ is required to build a massive elliptical galaxy by $`z2`$ (assuming that galaxies hosting high-z AGN have already converted the major fraction of their mass into stars); (iii) the recent discovery of elliptical galaxies at $`z1.5`$ which contain stellar populations with ages of 3–4 Gyrs (Dunlop et al. 1996, Peacock et al. 1998). Regardless of the cosmological model this requires an extreme formation redshift ($`z>5`$) for the initial starburst in these galaxies; (iv) the dramatic increase in the number of star-forming galaxies at high-redshift identified in ground-based and HST faint galaxy samples. Using a Lyman-break colour selection technique (Steidel et al. 1996), more than 3000 galaxies now have photometric redshifts with $``$ 700 galaxies already spectroscopically confirmed at $`z2`$ (Adelberger priv.comm.), with SFRs $`15h^2M_{}\mathrm{yr}^1`$. However the attenuating effects of dust, inevitably associated with star-formation, means that SFRs estimated from these rest-frame UV luminosities must be treated as strict lower-limits. Near-IR observations of rest-frame Balmer-line emission suggest an upward correction factor to the SFRs of $`215\times `$ (Pettini et al. 1998), whilst more robust measurements of SFRs, derived from rest-frame FIR luminosities, imply SFRs 600 times greater than that estimated from the UV luminosity (Hughes et al. 1998, Cimatti et al. 1998); (v) similar evolution seen in both the radio-source population and the local starburst population, implying that radio source evolution is a good tracer of the star-formation history of the Universe, suggests that the SFR density derived from Lyman-limit galaxies at $`z34`$ is under-estimated by a factor of $`5`$ (Dunlop 1998), and therefore that once again the star-formation activity in the Universe peaked at $`z>2`$; (vi) the small, but increasing number of submm continuum and CO detections of high-$`z`$ quasars and radio galaxies, indicate that the host galaxies of these powerful AGN contain large quantities of metal-enriched molecular gas ($`110\times 10^{10}M_{}`$, after correcting for gravitational amplification) which can fuel massive bursts of star-formation (Omont et al. 1996, Hughes et al. 1997, Combes et al. 1999).
Taken together, the observational evidence suggests that much of the on-going star-formation in the young Universe may be hidden by dust from optical surveys and possibly also from IR surveys. Hence the transparent view of the Universe provided by submm observations, which now have the instrumental sensitivity to detect high-$`z`$ dust-enshrouded galaxies forming stars at a rate $`>100M_{}\mathrm{yr}^1`$, and the preliminary evidence that galaxies (particularly massive spheroidal systems) exhibit strong luminosity evolution at submm wavelengths, demonstrate that comprehensive submm surveys will provide an important alternative measurement of the star-formation history of high-$`z`$ galaxies unhindered by the effects of dust.
## 2. Submillimetre Cosmological Surveys
The star-formation history of the high-$`z`$ starburst galaxy population can be determined from an accurate measure of the integral submm source-counts, the luminosities and redshift distribution of the submm-selected galaxies. The contribution of the submm sources to the total FIR–mm background measured by COBE (Hauser et al. 1998, Fixen et al. 1998) places an additional strong constraint on the possible evolution. By designing a series of cosmological submm surveys, covering a sufficiently wide range of complementary depths and areas, it is possible to discriminate between competing models of galaxy evolution and the epochs of formation of massive galaxies.
The possibility of conducting cosmological surveys at submm wavelengths has been realised in the last few years with rapid technological advances in semiconductor materials, wafer fabrication, filter design and the temperature stability and performance of cryogenic systems operating at $`100400`$ mK. This has led to the development and successful commissioning of sensitive bolometer arrays (e.g. SCUBA, SHARC, MAMBO, BOLOCAM), all of which will be upgraded within the next few years. These cameras operate on largest telescopes (10-m CSO, 15-m JCMT, 30-m IRAM) and exploit the best ground-based mm and submm observing conditions.
Despite these recent advances in instrumental sensitivity the primary reason that submm observations of galaxies at cosmological redshifts are at all possible is illustrated in Fig 1. When attempting to observe galaxies out to extreme redshifts, observational cosmologists usually suffer the combined effect of cosmological dimming due to the vast distances, a steeply declining luminosity function and positive K-corrections. However the steep spectral index of the Rayleigh-Jeans emission from dust ($`\mathrm{S}_\nu \nu ^{2+\beta },\beta 1.5`$) heated by young massive stars or AGN, which radiates at 30–70 K and dominates the FIR-submm luminosity, produces a negative K-correction at submm and mm wavelengths of sufficient strength to completely compensate for cosmological dimming at redshifts $`z1`$ with the result that, in an Einstein-de Sitter universe, a dust-enshrouded starburst galaxy of a given luminosity should be as easy to detect at an extreme redshift $`z8`$ as at $`z1`$. Note that at millimetre wavelengths sources actually become brighter with increasing redshift. The situation is inevitably less favourable (by a factor of 2-3) for low values of $`\mathrm{\Omega }`$, but nevertheless this relative “ease” of access to the very high-redshift universe remains unique to submm cosmology (Blain & Longair 1993, Blain & Longair 1996, Hughes & Dunlop 1998).
By early 2001, the initial extensive programme of extragalactic SCUBA (850 $`\mu `$m) surveys conducted on the 15-m JCMT will be completed, covering areas of 0.002–0.12 deg<sup>2</sup> with respective $`3\sigma `$ depths in the range $`1.5\mathrm{mJy}<\mathrm{S}_{850\mu \mathrm{m}}<8\mathrm{mJy}`$. To ensure these submm data are fully exploited, the current SCUBA surveys have been restricted to fields extensively studied at other wavelengths and which contain deep X-ray, optical, IR and radio imaging data (e.g. Hubble Deep Field, Hawaii Deep Fields, low-$`z`$ lensing clusters, CFRS, Lockman Hole, ELAIS). The results from these first submm surveys (outlined below) confirm that an era of massive dust-enshrouded star-formation exists at early epochs ($`z>1`$), with an intensity previously underestimated at optical wavelengths (Smail et al. 1997, Hughes et al. 1998, Eales et al. 1999, Blain et al. 1999, Barger et al. 1999b). However despite the enhancement of observed sub-mm fluxes of high-$`z`$ starburst galaxies by factors of 3–10 at 850$`\mu `$m, the deepest SCUBA surveys todate are still only sensitive to high-$`z`$ galaxies with SFRs comparable to the most luminous local ULIRGs ($`100M_{}\mathrm{yr}^1`$).
The following preliminary results from the on-going SCUBA surveys have already made a significant impact on several cosmological questions.
* The faint submm source-counts at 850 $`\mu `$m are reasonably well determined between 1–10 mJy (e.g. Barger et al. 1999b) and significantly exceed a no-evolution model, requiring roughly $`(1+z)^3`$ luminosity evolution out to $`z12`$, however a variety of models are consisitent with the data. The submm background measured by COBE requires that the SCUBA source-counts must converge at $`S_{850\mu m}0.5`$ mJy.
* Submm sources generally appear to be associated with $`z>1`$ galaxies, although it not yet clear whether they necessarily have optical, IR and radio counterparts. There is still much debate about the fraction of submm sources at $`z2`$, and the fraction of submm-selected galaxies that contain an AGN. Currently the most accurate identifications and spectroscopic redshifts are found in the SCUBA survey of lensing clusters (Frayer et al. 1998, 1999, Ivison et al. 1998, Barger et al. 1999a).
* Approximately 30–50% of the submm background has been resolved into individual high-$`z`$ galaxies at flux densities $`S_{850\mu \mathrm{m}}>2`$ mJy, and therefore existing unlensed submm surveys, which are confusion-limited at about this flux level, are only within a factor $`4`$ in sensitivity of resolving the entire submm/FIR background. In the SCUBA surveys of highly-lensed clusters it is possible to measure the source-counts down to a reduced confusion limit of $`S_{850\mu \mathrm{m}}0.5\mathrm{mJy}`$ (Blain et al. 1998). It now appears that the majority of submm background is due to a population of high-$`z`$ ultraluminous ($`L_{\mathrm{FIR}}>10^{12}L_{}`$) dusty galaxies forming stars at rates $`>100M_{}\mathrm{yr}^1`$, i.e. similar to the local ULIRG galaxies although the surface density of the high-$`z`$ submm population is significantly higher.
* At high-redshift ($`2<z<4`$) the submm surveys find $`5`$ times the star formation rate observed in the initial optical surveys (Madau 1997). However new optical surveys now agree more closely with the earlier submm result, finding no significant evidence for a decline in the star formation density between $`z=2`$ and $`z=4`$, when corrections are applied to the rest-frame UV luminosities to account for obscuration by dust (Pettini et al. 1998, Steidel et al. 1998).
## 3. Limitations on an Understanding of High-z Galaxy Evolution
Despite the success of the first SCUBA surveys, a number of deficiencies can be identified in the submm data which prevent a more accurate understanding of the star-formation history of high-$`z`$ galaxies. This review describes these deficiencies and outlines the future observations which will alleviate the following problems.
### 3.1. Constraining the evolutionary models
To improve the constraints on the competing evolutionary models provided by the current submm source-counts, it is necessary to (1) extend the restricted wavelength range of the surveys, (2) increase the dynamic range of the flux densities over which accurate source-counts are measured, and (3) increase the number of sources detected at a given flux level by surveying greater areas.
All these goals can be achieved by conducting future surveys with a combination of more sensitive, larger format bolometer arrays operating at 200$`\mu `$m – 3 mm on larger ground-based and airborne telescopes. Ground-based surveys at mm wavelengths can take advantage of a more stable and transparent atmosphere which will provide increased available integration time (to gain deeper survey sensitivity or greater survey area) and increased flux calibration accuracy. Future surveys with more sensitive and larger format arrays (e.g. BOLOCAM) will allow significantly greater areas to be covered (hence more sources detected) and will increase the range of the flux densities over which sources are detected. Furthermore conducting surveys with larger diameter telescopes (e.g. 50-m Large Millimetre Telescope (LMT/GTM), 100-m Green Bank Telescope) will reduce the beam-size, hence decrease the depth of the confusion limit (allowing deeper surveys) and improve the positional accuracy.
### 3.2. Future millimetre cosmological surveys
The SCUBA 850$`\mu `$m surveys have indicated that the extragalactic submillimetre background is dominated by a high-$`z`$ population of sources with $`S_{850\mu m}1`$ mJy. At $`z18`$ the same population of sources will have 1.1mm flux densities in the range $`0.30.6`$ mJy respectively. The extragalactic confusion limit at 850$`\mu `$m, estimated from the deepest SCUBA surveys, occurs at a depth of $`3\sigma <2\mathrm{m}\mathrm{J}\mathrm{y}`$ and corresponds to a source density N$`(S)4000\pm 1500\mathrm{deg}^2`$ at a resolution of 15 arcsecs. The ratio of 850$`\mu `$m (JCMT) and 1.1mm (LMT) beam-areas ($`6^{\prime \prime }`$) implies that confusion at 1.1 mm on the LMT will begin to become significant at a source density of $`24000\mathrm{deg}^2`$. Extrapolating models that adequately describe the measured 850$`\mu `$m source-counts to longer wavelengths suggests that at 1.1 mm confusion occurs at $`0.050.1`$ mJy. Consequently a deep 1.1 mm BOLOCAM/LMT survey has sufficient sensitivity and resolution to detect the entire submm-mm extragalactic background at a level above the confusion limit.
The current submm source-counts are measured with varying degrees of precision at 850$`\mu `$m between 1–10 mJy (Fig. 2,), whilst the future combination of BOLOCAM on the 10-m CSO (2000) and later ($`2002`$) on the 50-m LMT will provide accurate source counts at 1.1 mm between 0.1 and 100 mJy. Table 1 illustrates the predicted number of galaxies detected in a series of future 50-hour LMT surveys at 1.1 mm with BOLOCAM covering areas of 0.001–70 sq. degrees during the commissioning phase of the telescope and later, with an improved surface ($`\eta 70\mu `$m) during routine operation. For example in a similar duration to the original SCUBA HDF survey, a 50-hour LMT survey at 1.1 mm (with a $`3\sigma `$ detection limit of 2 mJy) will detect $`2000`$ galaxies equivalent to HDF850.1, the brightest submm source in the HDF ($`S_{850\mu \mathrm{m}}7`$ mJy).
### 3.3. Ambiguity in the counterparts & redshifts of submm galaxies
The current SCUBA surveys (with $`15^{\prime \prime }`$ resolution at 850$`\mu `$m) are struggling to unambiguously identify the submm sources with their optical/IR/radio counterparts. Hence the redshift distribution and luminosities of the submm sources are still uncertain. This results directly from the submm positional errors of $`23^{\prime \prime }`$ that are typical for even the highest S/N submm detections, and from the lack of submm data measuring the redshifted FIR spectral peak at 200–450 $`\mu `$m.
The positions of the brightest SCUBA sources ($`S_{850\mu m}>8`$ mJy) can be improved with mm-interferometric observations. However an IRAM Plateau de Bure follow-up of the brightest source in the Hubble Deep Field has demonstrated that even with $`2^{\prime \prime }`$ resolution and sub-arcsec positional errors, an ambiguous optical identification, and hence ambiguous redshift remains (Downes et al. 1999). It should be no surprise that submm selected galaxies, including those with mm-interferometric detections, do not always have optical counterparts, since high-$`z`$ galaxies observed in the earliest stages of formation may be heavily obscured by dust. Indeed this is the most compelling reason for conducting the submm surveys in the first instance and therefore searches for the counterparts may be more successful at near-infrared wavelengths. This was recently demonstrated by Smail et al. (1999) who took deep near-IR ($`2\mu `$m) images of two lensed clusters, previously observed by SCUBA (Smail et al. 1997). The original counterparts were identified as two bright low-redshift ($`z0.4`$) galaxies 5–10 arcsecs distant from the submm sources. However the new IR images revealed two high-$`z`$ ($`z>2`$) IR galaxies, with no optical counterparts, within 2-3 arcsecs of the SCUBA sources. The consequence of these mis-identifications is an inaccurate determination of star-formation history of high-$`z`$ galaxies.
The uncertainty in the redshift distribution of the submm-selected galaxies can be significantly reduced by measuring the mid-IR to radio SEDs of the individual sources. The power of using mid-IR to radio flux ratios (e.g. 15/850$`\mu `$m, 450/850$`\mu `$m, 850/1300$`\mu `$m, 850$`\mu `$m/1.4 GHz) as a crude measure of the redshift of submm-selected galaxies was demonstrated by Hughes et al. (1998) during the SCUBA survey of the Hubble Deep Field and has since been described elsewhere (e.g. Carilli & Yun 1999, Blain 1999). The overall similarity of the IR-radio SEDs of starburst galaxies, ULIRGS and radio-quiet AGN in the low-$`z`$ universe provides a useful template with which to compare the colours of high-$`z`$ submm population (particularly in the absence of information regarding the relative starburst/AGN contributions). Hence given sufficient instrumental sensitivity, the FIR–submm–radio colours of a submm source can discriminate between optical/IR counterparts which are equally probable on positional grounds alone, but which have significantly different redshifts, $`\delta z1.5`$ (Fig. 3).
This important technique, and the necessity for sensitive short submm data (200–500$`\mu `$m) measuring the rest-frame FIR SEDs of the individual high-z submm galaxies, without which it remains impossible to constrain their bolometric luminosities and SFRs, provide the major scientific justifications behind BLAST, a future ($`2003`$) NASA long-duration balloon-borne large-aperture submm telescope (P.I. M. Devlin, UPenn) operating at an altitude of 140,000 ft. Table 2 describes a series of possible BLAST surveys which demonstrate that even a single 50-hour survey will be able to follow-up all the wide-area shallow SCUBA surveys observed todate. For example if there are 5$`\sigma `$ SCUBA sources ($`S_{850\mu \mathrm{m}}>13`$ mJy) with no BLAST 5$`\sigma `$ counterparts at 300$`\mu `$m, i.e. $`S_{300\mu \mathrm{m}}<50`$ mJy, then the 300/850$`\mu `$m flux ratio must be $`4`$. This implies that the SCUBA source is most likely a galaxy at $`z3`$ for all typical starburst SEDs (Fig. 4). BLAST will also be an ideal complement to the future BOLOCAM bright mm-surveys ($`S_{1.1mm}>4\mathrm{m}\mathrm{J}\mathrm{y}`$) on the CSO and LMT. A measurement of the confusion noise due to extragalactic sources at 200–500$`\mu `$m with a $`2`$-m class telescope is an important secondary goal since the result will influence future FIRST survey strategies beyond 2007.
### 3.4. Millimetre CO-line spectroscopic redshifts
An accurate determination of the redshift distribution of submm-selected galaxies can ultimately be achieved through the measurement of mm-wavelength CO spectral-line redshifts, without recourse to having first identified the correct optical or IR counterparts. In the high-$`z`$ Universe the frequency separation of adjacent mm-wavelength CO transitions is $`\delta \nu _{\mathrm{J},\mathrm{J}1}115/(1+\mathrm{z})`$ GHz. Hence at redshifts $`>2`$, any adjacent pair of CO transitions are separated by $`<40`$ GHz, similar to the width of the 3 mm (75–110 GHz) atmospheric window. At these frequencies one can expect to detect the most luminous redshifted CO transitions from starbursts ($`J`$=6–5 $``$ $`J`$=3–2). Therefore, provided one can first pre-select from submm surveys those galaxies with sufficiently high (but still unknown) redshifts, using their FIR–mm colours, the availability of a “CO redshift machine” with a large instantaneous bandwidth ($`\mathrm{\Delta }\nu 35`$ GHz), operating on large single-dish mm-wavelength telescopes (e.g. 50-m LMT, 100-m GBT), will offer an incredibly powerful and more efficient alternative method to determine the accurate redshift distribution of the submm population.
Whilst future submm surveys will undoubtably detect increasing numbers of high-$`z`$ galaxies, an accurate description of their evolutionary history will not be possible without accurate redshifts and constraints on their bolometric luminosities (and star-formation rates). The overall strategy requires the follow-up of submm surveys with sensitive FIR–submm airborne (SOFIA), balloon-borne (BLAST) and satellite observations (SIRTF, FIRST), together with wide-band mm-wavelength spectroscopic measurements.
## References
Barger, A. et al. 1998, Nature, 394, 248
Barger, A. etal. 1999a, astro-ph/9903142, AJ, in press
Barger, A. etal. 1999b, ApJ, 518, L5
Blain, A. & Longair, M. 1993, MNRAS, 264, 509
Blain, A. & Longair, M. 1996, MNRAS, 279, 847
Blain, A., 1999, astro-ph/9906438, MNRAS, in press
Blain, A. et al. 1998, astro-ph/9812412
Carilli, C., Yun, M.S., 1999, astro-ph/9812251, ApJ, in press
Cimatti, A. et al. 1998, Nature, 392, 895
Combes, F. et al. 1999, astro-ph/9902286
Downes, D. et al. 1999, A&A, 347, 809
Dunlop, J.S. et al. 1996, Nature, 381, 581
Dunlop, J.S. 1998, in ‘Observational Cosmology with New Radio Surveys’, p.157
Eales, S.A. et al. 1999, ApJ, 515, 518
Fixsen, D.J. et al. 1998, ApJ, 508, 123
Frayer, D.T. et al. 1998, ApJ, 506, L7
Frayer, D.T. et al. 1999, astro-ph/9901311
Hauser, M.G. et al. 1998, ApJ, 508, 25
Hughes, D.H., Dunlop, J.S., Rawlings, S. 1997, MNRAS, 289, 766
Hughes, D.H. et al. 1998, Nature, 394, 241
Hughes, D.H. & Dunlop, J.S. 1998, astro-ph/9802260
Ivison, R.J. et al. 1998, MNRAS, 298, 583
Lilly, S. et al. 1999, ApJ, 518, 441
Madau, P. 1997, astro-ph/9709147
McLure, R. 1999, MNRAS, 308, 377
Omont, A. et al. 1996, A&A, 315, 10
Peacock, J.A. et al. 1998, MNRAS, 298, 1089
Pettini, M. et al. 1998, ApJ, 508, 539
Renzini A. 1998, astro-ph/9801209
Smail, I. et al. 1997, MNRAS, 490, L5
Smail, I. et al. 1999, astro-ph/9905246, MNRAS, in press
Steidel, C. et al. 1996, ApJ, 462, 17
Steidel, C. et al. 1998, astro-ph/9811399
Taylor, G. et al. 1996, MNRAS, 283, 930
|
warning/0003/math0003002.html
|
ar5iv
|
text
|
# Hyperelliptic jacobians and modular representations
## 1. Introduction
In the author proved that in characteristic $`0`$ the jacobian $`J(C)=J(C_f)`$ of a hyperelliptic curve
$$C=C_f:y^2=f(x)$$
has only trivial endomorphisms over an algebraic closure $`K_a`$ of the ground field $`K`$ if the Galois group $`\mathrm{Gal}(f)`$ of the irreducible polynomial $`fK[x]`$ is “very big”. Namely, if $`n=\mathrm{deg}(f)5`$ and $`\mathrm{Gal}(f)`$ is either the symmetric group $`𝐒_n`$ or the alternating group $`𝐀_n`$ then the ring $`\mathrm{End}(J(C_f))`$ of $`K_a`$-endomorphisms of $`J(C_f)`$ coincides with $`𝐙`$. The proof was based on an explicit description of the Galois module $`J(C_f)_2`$ of points of order $`2`$ on $`J(C_f)`$. Namely, the action of the Galois group $`\mathrm{Gal}(K)`$ factors through $`\mathrm{Gal}(f)`$ and the $`\mathrm{Gal}(f)`$-module $`J(C_f)_2`$ could be easily described in terms of the (transitive) action of $`\mathrm{Gal}(f)`$ on the set $`_f`$ of roots of $`f`$.
It turns out that if $`\mathrm{Gal}(f)`$ contains $`𝐀_n`$ then the Galois module $`J(C_f)_2`$ enjoys the following property ():
(\*):each subalgebra in $`\mathrm{End}_{𝐅_2}(J(C_f)_2)`$ which contains the identity operator and is stable under the conjugation by Galois automorphisms either consists of scalars or coincides with $`\mathrm{End}_{𝐅_2}(J(C_f)_2)`$.
Applying (\*) to the subalgebra $`\mathrm{End}(J(C_f))𝐙/2𝐙`$, one concludes that it consists of scalars, i.e., $`\mathrm{End}(J(C_f)`$ is a free abelian group of rank $`1`$ and therefore coincides with $`𝐙`$. (The case of $`\mathrm{End}(J(C))𝐙/2𝐙=\mathrm{End}_{𝐅_2}(J(C_f)_2)`$ could not occur in characteristic zero.)
The proof of (\*) was based on the well-known explicit description of $`J(C_f)_2`$ , and elementary properties of $`𝐀_n`$ and its simplest nontrivial representation in characteristic $`2`$ of dimension $`n1`$ or $`n2`$ (depending on whether $`n`$ is odd or even).
In this paper we study property (\*) itself from the point of view of representation theory over $`𝐅_2`$. Our results allow, in principle, to check the validity of (\*) even if $`\mathrm{Gal}(f)`$ does not contain $`𝐀_n`$. We prove that $`\mathrm{End}(J(C_f))=𝐙`$ for an infinite series of $`\mathrm{Gal}(f)=𝐋_2(2^r):=\mathrm{PSL}_2(𝐅_{2^r})`$ and $`n=2^r+1`$ (with $`r3`$ and $`\mathrm{dim}(J(C_f))=2^{r1}`$) or when $`\mathrm{Gal}(f)`$ is the Suzuki group $`\mathrm{𝐒𝐳}(2^{2r+1})`$ and $`n=2^{2(2r+1)}+1`$ (with $`\mathrm{dim}(J(C_f))=2^{4r+1}`$).
We refer the reader to , , , , , for a discussion of known results about, and examples of, hyperelliptic jacobians without complex multiplication.
The paper is organized as follows. In §2 we state the main results and begin the discussion of linear representations for which an analogue of the property (\*) holds true; we call such representations very simple. In §3 we prove that the very simplicity of the Galois module $`X_{\mathrm{}}`$ of points of prime order $`\mathrm{}`$ on an abelian variety $`X`$ implies in characteristic zero that $`X`$ does not have nontrivial endomorphisms. In §4 we remind basic facts about permutation groups and corresponding ordinary representations and modular representations over $`𝐅_2`$. We use them in §5 in order to restate the main results as assertions about the very simplicity of certain permutation modules using an explicit description of points of order $`\mathrm{}=2`$ on hyperelliptic jacobians. It turns out that all these permutations modules are Steinberg representations. In §6 we prove that the Steinberg representations are the only absolutely irreducible nontrivial representations (up to an isomorphism) over $`𝐅_2`$ for groups $`𝐋_2(2^r)`$ and $`\mathrm{𝐒𝐳}(2^{2r+1})`$. In §7 we study very simple linear representations; in particular, we prove that all the (modular) Steinberg representations discussed in Section 6 are very simple. This ends the proof of main results.
## 2. Main results
Throughout this paper we assume that $`K`$ is a field. We fix its algebraic closure $`K_a`$ and write $`\mathrm{Gal}(K)`$ for the absolute Galois group $`\mathrm{Aut}(K_a/K)`$. If $`X`$ is an abelian variety of dimension $`g`$ defined over $`K`$ then for each prime $`\mathrm{}\mathrm{char}(K)`$ we write $`X_{\mathrm{}}`$ for the kernel of multiplication by $`\mathrm{}`$ in $`X(K_a)`$. It is well-known that $`X_{\mathrm{}}`$ is a $`2g`$-dimensional $`𝐅_{\mathrm{}}`$-vector space provided with a natural structure of $`\mathrm{Gal}(K)`$-module. We write $`\mathrm{End}(X)`$ for the ring of $`K_a`$-endomorphisms of $`X`$ and $`\mathrm{End}^0(X)`$ for the corresponding finite-dimensional $`𝐐`$-algebra $`\mathrm{End}(X)𝐐`$.
The following notion plays a crucial role in this paper and will be discussed in detail in §7.
###### Definition 2.1.
Let $`V`$ be a vector space over a field $`𝐅`$, let $`G`$ be a group and $`\rho :G\mathrm{Aut}_𝐅(V)`$ a linear representation of $`G`$ in $`V`$. We say that the $`G`$-module $`V`$ is very simple if it enjoys the following property:
If $`R\mathrm{End}_𝐅(V)`$ be an $`𝐅`$-subalgebra containing the identity operator $`\mathrm{Id}`$ such that
$$\rho (\sigma )R\rho (\sigma )^1R\sigma G$$
then either $`R=𝐅\mathrm{Id}`$ or $`R=\mathrm{End}_𝐅(V)`$.
###### Remarks 2.2.
1. Clearly, the $`G`$-module $`V`$ is very simple if and only if the corresponding $`\rho (G)`$-module $`V`$ is very simple.
2. Clearly, if $`V`$ is very simple then the corresponding algebra homomorphism
$$𝐅[G]\mathrm{End}_𝐅(V)$$
is surjective. Here $`𝐅[G]`$ stands for the group algebra of $`G`$. In particular, a very simple module is absolutely simple.
3. If $`G^{}`$ is a subgroup of $`G`$ and the $`G^{}`$-module $`V`$ is very simple then the $`G`$-module $`V`$ is also very simple.
4. Let $`G^{}`$ be a normal subgroup of $`G`$. If $`V`$ is a faithful very simple $`G`$-module then either $`G^{}\mathrm{Aut}_𝐅(V)`$ consists of scalars (i.e., lies in $`𝐅\mathrm{Id}`$) or the $`G^{}`$-module $`V`$ is also very simple.
###### Lemma 2.3.
Let $`X`$ be an abelian variety of positive dimension $`g`$ over $`K`$. Let $`\mathrm{}`$ be a prime different from $`\mathrm{char}(K)`$. Assume that the $`\mathrm{Gal}(K)`$-module $`X_{\mathrm{}}`$ is very simple. Then either $`\mathrm{End}(X)=𝐙`$ or $`\mathrm{char}(K)>0`$ and $`X`$ is a supersingular abelian variety.
We prove Lemma 2.3 in §3.
###### Theorem 2.4.
Let $`K`$ be a field with $`\mathrm{char}(K)2`$, $`f(x)K[x]`$ an irreducible separable polynomial of degree $`n5`$. Let $`=_fK_a`$ be the set of roots of $`f`$, let $`K(_f)=K()`$ be the splitting field of $`f`$ and $`\mathrm{Gal}(f):=\mathrm{Gal}(K()/K)`$ the Galois group of $`f`$, viewed as a subgroup of $`\mathrm{Perm}()`$. Let $`C_f`$ be the hyperelliptic curve $`y^2=f(x)`$. Let $`J(C_f)`$ be its jacobian, $`\mathrm{End}(J(C_f))`$ the ring of $`K_a`$-endomorphisms of $`J(C_f)`$. Assume that $`n`$ and $`\mathrm{Gal}(f)`$ enjoy one of the following properties:
1. $`n=2^m+19`$ and the Galois group $`\mathrm{Gal}(f)`$ of $`f`$ contains a subgroup isomorphic to $`𝐋_2(2^m)`$;
2. For some positive integer $`k`$ we have $`n=2^{2(2k+1)}+1`$ and the Galois group $`\mathrm{Gal}(f)`$ of $`f`$ is isomorphic to $`\mathrm{𝐒𝐳}(2^{2k+1})`$;
Then:
1. The $`\mathrm{Gal}(K)`$-module $`J(C)_2`$ is very simple;
2. Either $`\mathrm{End}(J(C_f))=𝐙`$ or $`\mathrm{char}(K)>0`$ and $`J(C_f)`$ is a supersingular abelian variety.
###### Remark 2.5.
It follows from Lemma 2.3 that in order to prove Theorem 2.4, it suffices to check only the assertion a).
## 3. Proof of Lemma 2.3
Recall that $`\mathrm{dim}_𝐅_{\mathrm{}}(X_{\mathrm{}})=2g`$. Since $`X`$ is defined over $`K`$, one may associate with every $`u\mathrm{End}(X)`$ and $`\sigma \mathrm{Gal}(K)`$ an endomorphism $`{}_{}{}^{\sigma }u\mathrm{End}(X)`$ such that
$${}_{}{}^{\sigma }u(x)=\sigma u(\sigma ^1x)xX(F_a).$$
Let us put
$$R:=\mathrm{End}(X)𝐙/\mathrm{}𝐙\mathrm{End}_𝐅_{\mathrm{}}(X_{\mathrm{}}).$$
Clearly, $`R`$ satisfies all the conditions of Lemma 2.3. This implies that either $`R=𝐅_{\mathrm{}}\mathrm{Id}`$ or $`R=\mathrm{End}_𝐅_{\mathrm{}}(X_{\mathrm{}})`$. If $`\mathrm{End}(X)𝐙/\mathrm{}𝐙=R=𝐅_{\mathrm{}}\mathrm{Id}`$ then the free abelian group $`\mathrm{End}(X)`$ has rank $`1`$ and therefore coincides with $`𝐙`$. If $`\mathrm{End}(X)𝐙/\mathrm{}𝐙=R=\mathrm{End}_𝐅_{\mathrm{}}(X_{\mathrm{}})`$ then the free abelian group $`\mathrm{End}(X)`$ has rank $`(2\mathrm{d}\mathrm{i}\mathrm{m}(X))^2=(2g)^2`$ and therefore the $`𝐐`$-algebra $`\mathrm{End}^0(X)`$ has dimension $`(2g)^2`$.
Now Lemma 2.3 becomes an immediate corollary of the following assertion proven in (see Lemma 3.1).
###### Lemma 3.1.
Let $`Y`$ be an abelian variety of dimension $`g`$ over an algebraically closed field $`K_a`$. Assume that the semisimple $`𝐐`$-algebra $`\mathrm{End}^0(Y)=\mathrm{End}(Y)𝐐`$ has dimension $`(2g)^2`$. Then $`\mathrm{char}(K_a)>0`$ and $`Y`$ is supersingular.
## 4. Permutation groups and permutation modules
Let $`B`$ be a finite set consisting of $`n5`$ elements. We write $`\mathrm{Perm}(B)`$ for the group of permutations of $`B`$. A choice of ordering on $`B`$ gives rise to an isomorphism
$$\mathrm{Perm}(B)𝐒_n.$$
Let $`G`$ be a subgroup of $`\mathrm{Perm}(B)`$. For each $`bB`$ we write $`G_b`$ for the stabilizer of $`b`$ in $`G`$; it is a subgroup of $`G`$.
###### Remark 4.1.
Assume that the action of $`G`$ on $`B`$ is transitive. It is well-known that each $`G_b`$ is a subgroup of index $`n`$ in $`G`$ and all the $`G_b`$’s are conjugate one to another in $`G`$. Each conjugate of $`G_b`$ in $`G`$ is the stabilizer of a point in $`B`$. In addition, one may identify the $`G`$-set $`B`$ with the set of cosets $`G/G_b`$ with the standard action by $`G`$.
Let $`𝐅`$ be a field. We write $`𝐅^B`$ for the $`n`$-dimensional $`𝐅`$-vector space of maps $`h:B𝐅`$. The space $`𝐅^B`$ is provided with a natural action of $`\mathrm{Perm}(B)`$ defined as follows. Each $`s\mathrm{Perm}(B)`$ sends a map $`h:B𝐅`$ into $`sh:bh(s^1(b))`$. The permutation module $`𝐅^B`$ contains the $`\mathrm{Perm}(B)`$-stable hyperplane
$$(𝐅^B)^0=\{h:B𝐅\underset{bB}{}h(b)=0\}$$
and the $`\mathrm{Perm}(B)`$-invariant line $`𝐅1_B`$ where $`1_B`$ is the constant function $`1`$. The quotient $`𝐅^B/(𝐅^B)^0`$ is a trivial $`1`$-dimensional $`\mathrm{Perm}(B)`$-module.
Clearly, $`(𝐅^B)^0`$ contains $`𝐅1_B`$ if and only if $`\mathrm{char}(𝐅)`$ divides $`n`$. If this is not the case then there is a $`\mathrm{Perm}(B)`$-invariant splitting
$$𝐅^B=(𝐅^B)^0𝐅1_B.$$
Clearly, $`𝐅^B`$ and $`(𝐅^B)^0`$ carry natural structures of $`G`$-modules. Their characters depend only on characteristic of $`𝐅`$.
Let us consider the case of $`𝐅=𝐐`$. Then the character of $`𝐐^B`$ sends each $`gG`$ into the number of fixed points of $`g`$ (, ex. 2.2, p. 12); it is called the permutation character. Let us denote by $`\chi =\chi _B:G𝐐`$ the character of $`(𝐐^B)^0`$. It is known that the $`𝐐[G]`$-module $`(𝐐^B)^0`$ is absolutely simple if and only if $`G`$ acts doubly transitively on $`B`$ (, ex. 2.6, p. 17). Clearly, $`1+\chi `$ is the permutation character.
Now, let us consider the case of $`𝐅=𝐅_2`$. It is well-known that one may view $`𝐅_2^B`$ as the $`𝐅_2`$-vector space of all subsets of $`B`$ with symmetric difference as a sum. Namely, a subset $`T`$ corresponds to its characteristic function $`\chi _T:B\{0,1\}=𝐅_2`$ and a function $`h:B𝐅_2`$ corresponds to its support $`\mathrm{supp}(h)=\{xBh(x)=1\}`$. Under this identification each $`sG\mathrm{Perm}(B)`$ sends $`T`$ into $`s(T)=\{s(b)bT\}`$.
Under this identification the hyperplane $`(𝐅_2^B)^0`$ corresponds to the $`𝐅_2`$-vector space of all subsets of $`B`$ of even cardinality with symmetric difference as a sum.
If $`n`$ is even then let us define the $`\mathrm{Perm}(B)`$-module
$$Q_B:=(𝐅_2^B)^0/(𝐅_21_B).$$
If $`n`$ is odd then let us put
$$Q_B:=(𝐅_2^B)^0.$$
When $`n`$ is even, the quotient $`Q_B`$ corresponds to the $`n2`$-dimensional $`𝐅_2`$-vector space of all subsets of $`B`$ of even cardinality with symmetric difference as a sum where each subset $`TB`$ of even cardinality is identified with its complement $`BT`$.
###### Remark 4.2.
Clearly, $`\mathrm{dim}_{𝐅_2}(Q_B)=n1`$ if $`n`$ is odd and $`\mathrm{dim}_{𝐅_2}(Q_B)=n2`$ if $`n`$ is even. In both cases $`Q_B`$ is a faithful $`G`$-module.
Let $`G^{(2)}`$ be the set of $`2`$-regular elements of $`G`$. Clearly, the Brauer character of the $`G`$-module $`𝐅_2^B`$ coincides with the restriction of $`1+\chi _B`$ to $`G^{(2)}`$. This implies easily that the Brauer character of the $`G`$-module $`(𝐅_2^B)^0`$ coincides with the restriction of $`\chi _B`$ to $`G^{(2)}`$.
###### Remark 4.3.
Let us denote by $`\varphi _B=\varphi `$ the Brauer character of the $`G`$-module $`Q_B`$. One may easily check that $`\varphi _B`$ coincides with the restriction of $`\chi _B`$ to $`G^{(2)}`$ if $`n`$ is odd and with the restriction of $`\chi _B1`$ to $`G^{(2)}`$ if $`n`$ is even.
###### Remark 4.4.
Assume that $`n=\mathrm{\#}(B)`$ is even. Let us choose $`bB`$ and let $`G^{}:=G_b`$ and $`B^{}=B\{b\}`$. Then $`n^{}=\mathrm{\#}(B^{})=n1`$ is odd and there is a canonical isomorphism of $`G^{}`$-modules $`Q_B^{}Q_B`$ defined as follows. First, there is a natural $`G^{}`$-equivariant embedding $`𝐅_2^B^{}𝐅_2^B`$ which could be obtained by extending each $`h:B^{}𝐅_2`$ to $`B`$ by letting $`h(b)=0`$. Second, this embedding identifies $`(𝐅_2^B^{})^0`$ with a hyperplane of $`(𝐅_2^B)^0`$ which does not contain $`1_B`$. Now the desired isomorphism is given by the composition
$$Q_B^{}=(𝐅_2^B^{})^0(𝐅_2^B)^0(𝐅_2^B)^0/(𝐅_21_B)=Q_B.$$
This implies that if the $`G^{}`$-module $`Q_B^{}`$ is very simple then the $`G`$-module $`Q_B`$ is also very simple.
###### Remark 4.5.
Assume that $`G`$ acts on $`B`$ doubly transitively, $`\mathrm{\#}(B)`$ is odd and $`\mathrm{\#}(B)1=\mathrm{dim}_𝐐((𝐐_B)^0)`$ coincides with the largest power of $`2`$ dividing $`\mathrm{\#}(G)`$. Then it follows from a theorem of Brauer-Nesbitt (, Sect. 16.4, pp. 136–137 ; , p. 249) that $`Q_B`$ is an absolutely simple $`𝐅_2[G]`$-module. In particular, $`Q_B`$ is (the reduction of) the Steinberg representation .
## 5. Points of order $`2`$ on hyperelliptic jacobians
We keep all notations of Section 2. In addition, we assume that $`K`$ is a field of characteristic different from $`2`$. Let $`C`$ be a hyperelliptic curve over $`K`$ defined by an equation $`y^2=f(x)`$ where $`f(x)K[x]`$ is a polynomial of degree $`n5`$ without multiple roots. The rational function $`xK(C)`$ defines a canonical double cover $`\pi :C𝐏^1`$. Let $`B^{}C(K_a)`$ be the set of ramification points of $`\pi `$ (Weierstraß points). Clearly, the restriction of $`\pi `$ to $`B^{}`$ is an injective map $`\pi :B^{}𝐏^1(K_a)`$, whose image is either the set $`=_f`$ of roots of $`f`$ if $`n`$ is even or the disjoint union of $`\mathrm{}`$ and $``$ if $`n`$ is odd. By abuse of notation, we also denote by $`\mathrm{}`$ the ramification point lying above $`\mathrm{}`$ if $`n`$ is odd and by $`\mathrm{}_1`$ and $`\mathrm{}_2`$ two unramified points lying above $`\mathrm{}`$ if $`n`$ is even. Clearly, if $`n`$ is odd then $`\mathrm{}C(K)`$. If $`n`$ is even then the $`2`$-element set $`\{\mathrm{}_1,\mathrm{}_2\}`$ is stable under the action of $`\mathrm{Gal}(K)`$.
Let us put
$$B=\{(\alpha ,0)f(\alpha )=0\}C(K_a).$$
Then $`\pi `$ defines a bijection between $`B`$ and $``$ which commutes with the action of $`\mathrm{Gal}(K)`$. If $`n`$ is even then $`B`$ coincides with $`B^{}`$. In the case of odd $`n`$ the set $`B^{}`$ is the disjoint union of $`B`$ and $`\mathrm{}`$.
###### Theorem 5.1.
Suppose $`n`$ is an integer which is greater than or equal to $`5`$. Suppose $`f(x)K[x]`$ is a separable polynomial of degree $`n`$, $`K_a`$ the set of roots of $`f`$, let $`K()`$ be the splitting field of $`f`$ and $`\mathrm{Gal}(f):=\mathrm{Gal}(K()/K)`$ the Galois group of $`f`$.
Suppose $`C`$ is the hyperelliptic curve $`y^2=f(x)`$ of genus $`g=[\frac{n1}{2}]`$ over $`K`$. Suppose $`J(C)`$ is the jacobian of $`C`$ and $`J(C)_2`$ is the group of its points of order $`2`$, viewed as a $`2g`$-dimensional $`𝐅_2`$-vector space provided with the natural action of $`\mathrm{Gal}(K)`$. Then the homomorphism $`\mathrm{Gal}(K)\mathrm{Aut}_{𝐅_2}(J(C)_2)`$ factors through the canonical surjection $`\mathrm{Gal}(K)\mathrm{Gal}(K()/K)=\mathrm{Gal}(f)`$ and the $`\mathrm{Gal}(f)`$-modules $`J(C)_2`$ and $`Q_{}`$ are isomorphic. In particular, the $`G(K)`$-module $`J(C)_2`$ is very simple if and only if the $`\mathrm{Gal}(f)`$-module $`Q_B`$ is very simple.
###### Remark 5.2.
Clearly, $`\mathrm{Gal}(K)`$ acts on $`B`$ through the canonical surjective homomorphism $`\mathrm{Gal}(K)\mathrm{Gal}(f)`$, because all points of $`B`$ are defined over $`K()`$ and the natural homomorphism $`\mathrm{Gal}(f)\mathrm{Perm}(B)`$ is injective. Clearly, $`\pi :B`$ is a bijection of $`\mathrm{Gal}(f)`$-sets. This implies easily that the $`\mathrm{Gal}(f)`$-modules $`Q_B`$ and $`Q_{}`$ are isomorphic. So, in order to prove Theorem 5.1 it suffices to check that the $`\mathrm{Gal}(K)`$-modules $`Q_B`$ and $`J(C)_2`$ are isomorphic.
###### Proof of Theorem 5.1.
Here is a well-known explicit description of the group $`J(C)_2`$ of points of order $`2`$ on $`J(C)`$. Let us denote by $`L`$ the $`K`$-divisor $`2(\mathrm{})`$ on $`C`$ if $`n`$ is odd and the $`K`$-divisor $`(\mathrm{}_1)+(\mathrm{}_2)`$ if $`n`$ is even. In both cases $`L`$ is an effective divisor of degree $`2`$. Namely, let $`TB^{}`$ be a subset of even cardinality. Then (, Ch. IIIa, Sect. 2, Lemma 2.4; , pp. 190–191; see also ) the divisor $`e_T=_{PT}(P)\frac{\mathrm{\#}(T)}{2}L`$ on $`C`$ has degree $`0`$ and $`2e_T`$ is principal. If $`T_1,T_2`$ are two subsets of even cardinality in $`B^{}`$ then the divisors $`e_{T_1}`$ and $`e_{T_2}`$ are linearly equivalent if and only if either $`T_1=T_2`$ or $`T_2=B^{}T_1`$. Also, if $`T=T_1\mathrm{}T_2`$ then the divisor $`e_T`$ is linearly equivalent to $`e_{T_1}+e_{T_2}`$. Hereafter we use the symbol $`\mathrm{}`$ for the symmetric difference of two sets. Counting arguments imply easily that each point of $`J(C)_2`$ is the class of $`e_T`$ for some $`T`$. We know that such a choice is not unique. However, in the case of odd $`n`$ if we demand that $`T`$ does not contain $`\mathrm{}`$ then such a choice always exists and unique. This observation leads to a canonical group isomorphism
$$Q_B=(𝐅_2^B)^0J(C)_2,T\mathrm{cl}(e_T)$$
in the case of odd $`n`$. Here $`\mathrm{cl}`$ stands for the linear equivalence class of a divisor. In the case of even $`n`$ we are still able to define a canonical surjective group homomorphism
$$(𝐅_2^B)^0J(C)_2,T\mathrm{cl}(e_T)$$
and one may easily check that the kernel of this map is the line generated by the set $`B`$, i.e., the line generated by the constant function $`1_B`$. This gives rise to the injective homomorphism
$$Q_B=(𝐅_2^B)^0/(𝐅_21_B)J(C)_2,$$
which is an isomorphism, by counting arguments. So, in both (odd and even) cases we get a canonical isomorphism $`Q_BJ(C)_2`$, which obviously commutes with the actions of $`\mathrm{Gal}(K)`$. In other words, we constructed an isomorphism of $`\mathrm{Gal}(K)`$-modiles $`Q_B`$ and $`J(C)_2`$. In light of Remark 5.2, this ends the proof of Theorem 5.1. ∎
Combining Theorem 5.1 and Lemma 2.3 (for $`\mathrm{}=2`$), we obtain the following corollary.
###### Corollary 5.3.
Let $`K`$ be a field with $`\mathrm{char}(K)2`$, $`K_a`$ its algebraic closure, $`f(x)K[x]`$ an irreducible separable polynomial of degree $`n5`$. Let $`=_fK_a`$ be the set of roots of $`f`$, let $`K(_f)=K()`$ be the splitting field of $`f`$ and $`\mathrm{Gal}(f):=\mathrm{Gal}(K()/K)`$ the Galois group of $`f`$, viewed as a subgroup of $`\mathrm{Perm}()`$. Let $`C_f`$ be the hyperelliptic curve $`y^2=f(x)`$. Let $`J(C_f)`$ be its jacobian, $`\mathrm{End}(J(C_f))`$ the ring of $`K_a`$-endomorphisms of $`J(C_f)`$. Assume that the $`\mathrm{Gal}(f)`$-module $`Q_{}`$ is very simple. Then either $`\mathrm{End}(J(C_f))=𝐙`$ or $`\mathrm{char}(K)>0`$ and $`J(C_f)`$ is a supersingular abelian variety.
Notice that in order to prove Theorem 2.4, it suffices to check the following statement.
###### Theorem 5.4.
Let $`n`$ be a positive integer, $`B`$ a $`n`$-element set, $`H\mathrm{Perm}(B)`$ a permutation group. Assume that $`(n,H)`$ enjoy one of the following properties:
1. $`n=2^m+19`$ and $`H`$ contains a subgroup isomorphic to $`𝐋_2(2^m)`$;
2. For some positive integer $`k`$ we have $`n=2^{2(2k+1)}+1`$ and $`H`$ contains a subgroup isomorphic to $`\mathrm{𝐒𝐳}(2^{2k+1})`$;
Then the $`H`$-module $`Q_B`$ is very simple.
###### Proof of Theorem 2.4 modulo Theorem 5.4.
Let us put
$$n=\mathrm{deg}(f),B=,H=\mathrm{Gal}(f).$$
It follows from Theorem 5.4 that the $`\mathrm{Gal}(f)`$-module $`𝐐_{}`$ is very simple. Now the result follows readily from Corollary 5.3. ∎
We prove Theorem 5.4 at the end of §7.
## 6. Steinberg representation
In this section we prove that the Steinberg representation is the only nontrivial absolutely irreducible representation over $`𝐅_2`$ (up to an isomorphism) of groups $`𝐋_2(2^m))`$ and $`\mathrm{𝐒𝐳}(2^{2k+1})`$. We refer to for basic properties of Steinberg representations.
Let us fix an algebraic closure of $`𝐅_2`$ and denote it by $``$. We write $`\varphi :`$ for the Frobenius automorphism $`xx^2`$. Let $`q=2^m`$ be a positive integral power of two. Then the subfield of invariants of $`\varphi ^m:`$ is a finite field $`𝐅_q`$ consisting of $`q`$ elements. Let $`q^{}`$ be an integral positive power of $`q`$. If $`d`$ is a positive integer and $`i`$ is a non-negative integer then for each matrix $`u\mathrm{GL}_d()`$ we write $`u^{(i)}`$ for the matrix obtained by raising each entry of $`u`$ to the $`2^i`$th power.
Recall that an element $`\alpha 𝐅_q`$ is called primitive if $`\alpha 0`$ and has multiplicative order $`q1`$ in the cyclic multiplicative group $`𝐅_q^{}`$.
###### Lemma 6.1.
Let $`q>2`$, let $`d`$ be a positive integer and let $`G`$ be a subgroup of $`\mathrm{GL}_d(𝐅_q^{})`$. Assume that there exists an element $`uG\mathrm{GL}_d(𝐅_q^{})`$, whose trace $`\alpha `$ lies in $`𝐅_q^{}`$ and has multiplicative order $`q1`$. Let $`V_0=^d`$ and let $`\rho _0:G\mathrm{GL}_d(𝐅_q^{})\mathrm{GL}_d()=\mathrm{Aut}_{}(V_0)`$ be the natural $`d`$-dimensional representation of $`G`$ over $``$. For each positive integer $`i<m`$ we define a $`d`$-dimensional $``$-representation
$$\rho _i:G\mathrm{Aut}(V_i)$$
as the composition of
$$G\mathrm{GL}_d(𝐅_q^{}),xx^{(i)}$$
and the inclusion map
$$\mathrm{GL}_d(𝐅_q^{})\mathrm{GL}_d()\mathrm{Aut}_{}(V_i).$$
Let $`S`$ be a subset of $`\{0,1,\mathrm{}m1\}`$. Let us define a $`d^{\mathrm{\#}(S)}`$-dimensional $``$-representation $`\rho _S`$ of $`G`$ as the tensor product of representations $`\rho _i`$ for all $`iS`$. If $`S`$ is a proper subset of $`\{0,1,\mathrm{}m1\}`$ then there exists an element $`uG`$ such that the trace of $`\rho _S(u)`$ does not belong to $`𝐅_2`$. In particular, $`\rho _S`$ could not be obtained by extension of scalars to $``$ from a representation of $`G`$ over $`𝐅_2`$.
###### Proof.
Clearly,
$$\mathrm{tr}(\rho _i(u))=(\mathrm{tr}(\rho _0(u))^{2^i}uG.$$
This implies easily that
$$\mathrm{tr}(\rho _S(u))=\underset{iS}{}\mathrm{tr}(\rho _i(u))=(\mathrm{tr}(\rho _0(u))^M$$
where $`M=_{iS}2^i`$. Since $`S`$ is a proper subset of $`\{0,1,\mathrm{}m1\}`$, we have
$$0<M<\underset{i=0}{\overset{m1}{}}2^i=2^m1=\mathrm{\#}(𝐅_q^{}).$$
Recall that there exists $`uG`$ such that $`\alpha =\mathrm{tr}(\rho _0(u))`$ lies in $`𝐅_q^{}`$ and the exact multiplicative order of $`\alpha `$ is $`q1=2^m1`$.
This implies that $`0\alpha ^M1`$. Since $`𝐅_2=\{0,1\}`$, we conclude that $`\alpha ^M𝐅_2`$. Therefore $`\mathrm{tr}(\rho _S(u))=(\mathrm{tr}(\rho _0(u))^M=\alpha ^M𝐅_2`$. ∎
###### Theorem 6.2.
Let $`q8`$ be a power of $`2`$ and $`G=𝐋_2(q)=\mathrm{PSL}_2(𝐅_q)=\mathrm{SL}_2(𝐅_q)`$. Let $`\rho :G\mathrm{Aut}(V)`$ be an absolutely irreducible faithful representation of $`G`$ over $``$. If the trace map $`\mathrm{tr}_\rho :G`$ takes on values in $`𝐅_2`$ then $`\mathrm{dim}_{}(V)=q`$. In particular, $`\rho `$ is the Steinberg representation of $`G`$.
###### Proof.
Let us put $`q^{}=q`$. We have
$$G=\mathrm{SL}_2(𝐅_q)\mathrm{GL}_2(𝐅_q).$$
Clearly, for each $`\alpha 𝐅_q`$ one may find a $`2\times 2`$ matrix with determinant $`1`$ and trace $`\alpha `$. This implies that $`G`$ satisfies the conditions of Lemma 6.1.
The construction described in Lemma 6.1 allows us to construct a $`d^{\mathrm{\#}(S)}`$-dimensional $``$-representation $`\rho _S`$ of $`G`$ for each subset $`S`$ of of $`\{0,1,\mathrm{}m1\}`$. It is well-known (, pp. 588-589) that $`\rho _S`$’s exhaust the list of all absolutely irreducible $``$-representations of $`G=\mathrm{SL}_2(𝐅_q)`$ and therefore $`\rho `$ is isomorphic to $`\rho _S`$ for some $`S`$. It follows from Lemma 6.1 that either $`S`$ is empty or $`S=\{0,1,\mathrm{}m1\}`$. The case of empty $`S`$ corresponds to the trivial $`1`$-dimensional representation. Therefore $`S=\{0,1,\mathrm{}m1\}`$ and $`\rho `$ is $`2^m=q`$-dimensional. ∎
Suppose $`m=2k+13`$ is an odd integer. Let $`q=2^m=2^{2k+1}`$ and $`d=4`$. Recall (. pp. 182–194) that the Suzuki group $`\mathrm{𝐒𝐳}(q)`$ is the subgroup of $`\mathrm{GL}_4(𝐅_q)`$ generated by the matrices $`S(a,b),M(\lambda ),T`$ defined as follows. For each $`a,b𝐅_q`$ the matrix $`S(a,b)`$ is defined by
$$S(a,b)=\left(\begin{array}{cccc}1& 0& 0& 0\\ a& 1& 0& 0\\ b& a^{2^{k+1}}& 1& 0\\ a^{2^{k+1}+2}+ab+b^{2^{k+1}}& a^{2^{k+1}+1}+b& a& 1\end{array}\right)$$
and for each $`\lambda 𝐅_q^{}`$ the matrix $`M(\lambda )`$ is defined by
$$M(\lambda )=\left(\begin{array}{cccc}\lambda ^{1+2^k}& 0& 0& 0\\ 0& \lambda ^{2^k}& 0& 0\\ 0& 0& \lambda ^{2^k}& 0\\ 0& 0& 0& \lambda ^{1+2^k}\end{array}\right).$$
The matrix $`T`$ is defined by
$$T=\left(\begin{array}{cccc}0& 0& 0& 1\\ 0& 0& 1& 0\\ 0& 1& 0& 0\\ 1& 0& 0& 0\end{array}\right).$$
Notice that the trace of $`S(0,b)T`$ is $`b^{2^{k+1}}`$. This implies easily that for each $`\alpha 𝐅_q`$ one may find an element of $`\mathrm{𝐒𝐳}(q)\mathrm{GL}_4(𝐅_q)`$ with trace $`\alpha `$. This implies that
$$G=\mathrm{𝐒𝐳}(q)\mathrm{GL}_4(𝐅_q)$$
satisfies the conditions of Lemma 6.1. Notice also that $`\mathrm{\#}(\mathrm{𝐒𝐳}(q))=(q^2+1)q^2(q1)`$ (, p. 187).
###### Theorem 6.3.
Let $`\rho :\mathrm{𝐒𝐳}(q)\mathrm{Aut}(V)`$ be an absolutely irreducible faithful representation of $`\mathrm{𝐒𝐳}(q)`$ over $``$. If the trace map $`\mathrm{tr}_\rho :\mathrm{𝐒𝐳}(q)`$ takes on values in $`𝐅_2`$ then $`\mathrm{dim}_{}(V)=q^2`$. In particular, $`\rho `$ is the Steinberg representation of $`G`$.
###### Proof.
Let us put $`q^{}=q`$. We know that $`G=\mathrm{𝐒𝐳}(q)\mathrm{GL}_4(𝐅_q)`$ satisfies the conditions of Lemma 6.1.
The construction described in Lemma 6.1 allows us to construct a $`4^{\mathrm{\#}(S)}`$-dimensional $``$-representation $`\rho _S`$ of $`G`$ for each subset $`S`$ of of $`\{0,1,\mathrm{}m1\}`$. It is known (, pp. 56–57) that $`\rho _S`$’s exhaust the list of all absolutely irreducible $``$-representations of $`G=\mathrm{SL}_2(𝐅_q)`$ and therefore $`\rho `$ is isomorphic to $`\rho _S`$ for some $`S`$. It follows from Lemma 6.1 that either $`S`$ is empty or $`S=\{0,1,\mathrm{}m1\}`$. The case of empty $`S`$ corresponds to trivial $`1`$-dimensional representation. Therefore $`S=\{0,1,\mathrm{}m1\}`$ and $`\rho `$ is $`4^m=q^2`$-dimensional. ∎
###### Remark 6.4.
Assume that in the case 5.4(i) (resp. 5.4(ii)) that $`H=𝐋_2(2^m))`$ (resp. $`\mathrm{𝐒𝐳}(2^{2k+1})`$). It follows from Remark 4.5 that $`Q_B`$ is the Steinberg representation of $`H`$.
## 7. Very simple representations
###### Examples 7.1.
1. If $`\mathrm{dim}(V)=1`$ then $`V`$ is always very simple.
2. Assume that there exist $`G`$-modules $`V_1`$ and $`V_2`$ such that $`\mathrm{dim}(V_1)>1,\mathrm{dim}(V_2)>1`$ and the $`G`$-module $`V`$ is isomorphic to $`V_1_𝐅V_2`$. Then $`V`$ is not very simple. Indeed, the subalgebra
$$R=\mathrm{End}_𝐅(V_1)𝐅\mathrm{Id}_{V_2}\mathrm{End}_𝐅(V_1)_𝐅\mathrm{End}_𝐅(V_2)=\mathrm{End}_𝐅(V)$$
is stable under the conjugation by elements of $`G`$ but coincides neither with $`𝐅\mathrm{Id}`$ nor with $`\mathrm{End}_𝐅(V)`$. (Here $`\mathrm{Id}_{V_2}`$ stands for the identity operator in $`V_2`$.)
3. Let $`XG`$ be a central extension of $`G`$. Assume that there exist $`X`$-modules $`V_1`$ and $`V_2`$ such that $`\mathrm{dim}(V_1)>1,\mathrm{dim}(V_2)>1`$ and $`V`$, viewed as $`X`$-module, is isomorphic to $`V_1_𝐅V_2`$. Then $`V`$ is not very simple as an $`X`$-module. Since $`X`$ and $`G`$ have the same images in $`\mathrm{Aut}_𝐅(V)`$, the $`G`$-module $`V`$ is also not very simple.
4. Assume that there exists a subgroup $`G^{}`$ in $`G`$ of finite index $`m>1`$ and a $`G^{}`$-module $`V^{}`$ such that the $`𝐅[G]`$-module $`V`$ is induced by the $`𝐅[G^{}]`$-module $`V^{}`$. (In particular, $`m`$ must divide $`\mathrm{dim}(V)`$.) Then $`V`$ is not very simple. Indeed, one may view $`W`$ as a $`G^{}`$-submodule of $`V`$ such that $`V`$ coincides with the direct sum $`_{\sigma G/G^{}}\sigma W`$. Let $`R=_{\sigma G/G^{}}\mathrm{End}_𝐅(\sigma W)`$ be the algebra of all operators sending each $`\sigma W`$ into itself. Then $`R`$ is stable under the conjugation by elements of $`G`$ but coincides neither with $`𝐅\mathrm{Id}`$ nor with $`\mathrm{End}_𝐅(V)`$.
###### Example 7.2.
Let $`n5`$ be an integer, $`B`$ a $`n`$-element set. Suppose $`G`$ is either $`\mathrm{Perm}(B)𝐒_n`$ or the only subgroup in $`\mathrm{Perm}(B)`$ of index $`2`$ (isomorphic to $`𝐀_n`$). Then the $`G`$-module $`Q_B`$ is very simple. If $`n`$ is odd then this assertion is proven in , Th. 4.1. If $`n`$ is even then $`n6`$, $`n^{}=n15`$ is odd and the result follows from the odd case combined with Remarks 2.2 and 4.4.
###### Remarks 7.3.
Assume that there exist $`G`$-modules $`V_1`$ and $`V_2`$ such that $`\mathrm{dim}(V_1)>1,\mathrm{dim}(V_2)>1`$ and the $`G`$-module $`V`$ is isomorphic to $`V_1_𝐅V_2`$.
1. If $`V`$ is simple then both $`V_1`$ and $`V_2`$ are also simple. Indeed, if say, $`V^{}`$ is a proper $`G`$-stable subspace in $`V_1`$ then $`V^{}_𝐅V_2`$ is a proper $`G`$-stable subspace in $`V_1_𝐅V_2=V`$.
2. If $`V`$ is absolutely simple then both $`V_1`$ and $`V_2`$ are also absolutely simple. Indeed, assume that say, $`R_1:=\mathrm{End}_G(V_1)`$ has $`𝐅`$-dimension greater than $`1`$. Then $`\mathrm{End}_G(V)=\mathrm{End}_G(V_1_𝐅V_2)`$ contains $`R_1\mathrm{Id}_{V_2}R_1`$ and therefore also has dimension greater than $`1`$.
###### Lemma 7.4.
Let $`H`$ be a group, $`𝐅`$ a field and $`V`$ a simple $`𝐅[H]`$-module of finite $`𝐅`$-dimension $`N`$. Let $`R\mathrm{End}_𝐅(V)`$ be an $`𝐅`$-subalgebra containing the identity operator $`\mathrm{Id}`$ and such that
$$uRu^1RuH.$$
Then:
1. The faithful $`R`$-module $`V`$ is semisimple.
2. Either the $`R`$-module $`V`$ is isotypic or there exists a subgroup $`H^{}H`$ of index $`r`$ dividing $`N`$ and a $`H^{}`$-module $`V^{}`$ of finite $`𝐅`$-dimension $`N/r`$ such that $`r>1`$ and the $`H`$-module $`V`$ is induced by $`V^{}`$. In addition, if $`𝐅=𝐅_2`$ then $`r<N`$.
###### Proof.
We may assume that $`N>1`$. Clearly, $`V`$ is a faithful $`R`$-module and
$$uRu^1=RuH.$$
Step 1. $`V`$ is a semisimple $`R`$-module. Indeed, let $`UV`$ be a simple $`R`$-submodule. Then $`U^{}=_{sH}sU`$ is a non-zero $`H`$-stable subspace in $`V`$ and therefore must coincide with $`V`$. On the other hand, each $`sU`$ is also a $`R`$-submodule in $`V`$, because $`s^1Rs=R`$. In addition, if $`WsU`$ is an $`R`$-submodule then $`s^1W`$ is an $`R`$-submodule in $`U`$, because
$$Rs^1W=s^1sRs^1W=s^1RW=s^1W.$$
Since $`U`$ is simple, $`s^1W=\{0\}`$ or $`U`$. This implies that $`sU`$ is also simple. Hence $`V=U^{}`$ is a sum of simple $`R`$-modules and therefore is a semisimple $`R`$-module.
Step 2. The $`R`$-module $`V`$ is either isotypic or induced. Indeed, let us split the semisimple $`R`$-module $`V`$ into the direct sum
$$V=V_1\mathrm{}V_r$$
of its isotypic components. Dimension arguments imply that $`r\mathrm{dim}(V)=N`$. It follows easily from the arguments of the previous step that for each isotypic component $`V_i`$ its image $`sV_i`$ is an isotypic $`R`$-submodule for each $`sH`$ and therefore is contained in some $`V_j`$. Similarly, $`s^1V_j`$ is an isotypic submodule obviously containing $`V_i`$. Since $`V_i`$ is the isotypic component, $`s^1V_j=V_i`$ and therefore $`sV_i=V_j`$. This means that $`s`$ permutes the $`V_i`$; since $`V`$ is $`H`$-simple, $`H`$ permutes them transitively.
This implies that all $`V_i`$ have the same dimension $`N/r`$ and therefore $`r`$ divides $`\mathrm{dim}(V)=N`$. Let $`H^{}=H_i`$ be the stabilizer of $`V_i`$ in $`H`$, i.e.
$$H_i=\{sHsV_i=V_i\}.$$
The transitivity of the action of $`H`$ on $`V_j`$s implies that $`[H:H_i]=r`$.
If $`r=1`$ then $`H=H^{}=H_i`$. This means that $`sV_i=V_i`$ for all $`sH`$ and $`V=V_i`$ is isotypic.
Assume that $`r>1`$ and consider the $`H^{}`$-module $`W=V_i`$. Clearly, $`[H:H^{}]=[H:H_i]=r`$ divides $`N`$ and the $`H`$-module $`V`$ is iduced by $`W`$.
Step 3. Assume that $`r=N`$ and $`𝐅=𝐅_2`$. Then each $`V_i`$ is one-dimensional and contains exactly one non-zero vector say, $`v_i`$. Then the sum $`_{i=1}^Nv_i`$ is a non-zero $`H`$-invariant vector which contradicts the simplicity of the $`H`$-module $`V`$. ∎
###### Theorem 7.5.
Suppose $`H`$ is a group and
$$\rho :H\mathrm{Aut}_{𝐅_2}(V)$$
is an absolutely simple $`𝐅_2[H]`$-module of finite dimension $`N`$. Suppose there exists an $`𝐅_2`$-subalgebra $`R\mathrm{End}_{𝐅_2}(V)`$ containing the identity operator $`\mathrm{Id}`$ and such that
$$uRu^1RuH.$$
Assume, in addition, that $`H`$ does not have nontrivial cyclic quotients of order dividing $`N`$. If the $`R`$-module $`V`$ is isotypic then there exist $`𝐅_2[H]`$-modules $`V_1`$ and $`V_2`$ such that $`V`$, viewed as $`H`$-module, is isomorphic to $`V_1_{𝐅_2}V_2`$ and the image of $`R\mathrm{End}_{𝐅_2}(V)`$ under the induced isomorphism
$$\mathrm{End}_{𝐅_2}(V)=\mathrm{End}_{𝐅_2}(V_1_{𝐅_2}V_2)=\mathrm{End}_{𝐅_2}(V_1)_{𝐅_2}\mathrm{End}_{𝐅_2}(V_2)$$
coincides with $`\mathrm{End}_{𝐅_2}(V_1)\mathrm{Id}_{V_2}`$. In particular, if both $`V_1`$ and $`V_2`$ have dimension greater than $`1`$ then the $`H`$-module $`V`$ is not very simple.
###### Proof.
Since $`V`$ is isotypic, there exist a simple $`R`$-module $`W`$, a positive integer $`d`$ and an isomorphism
$$\psi :VW^d$$
of $`R`$-modules. Let us put
$$V_1=W,V_2=𝐅_2^d.$$
The isomorphism $`\psi `$ gives rise to the isomorphism of $`𝐅_2`$-vector spaces
$$V=W^d=W_{𝐅_2}𝐅_2^d=V_1_{𝐅_2}V_2.$$
We have
$$d\mathrm{dim}(W)=\mathrm{dim}(V)=N.$$
Clearly, $`\mathrm{End}_R(V)`$ is isomorphic to the matrix algebra $`\mathrm{Mat}_d(\mathrm{End}_R(W))`$ of size $`d`$ over $`\mathrm{End}_R(W)`$.
Let us put
$$k=\mathrm{End}_R(W).$$
Since $`W`$ is simple, $`k`$ is a finite-dimensional division algebra over $`𝐅_2`$. Therefore $`k`$ must be a finite field.
We have
$$\mathrm{End}_R(V)\mathrm{Mat}_d(k).$$
Clearly, $`[k:𝐅_2]`$ divides $`\mathrm{dim}_{𝐅_2}(W)`$ and therefore divides $`\mathrm{dim}_{𝐅_2}(V)=N`$. Clearly, $`\mathrm{Aut}(k/𝐅_2)`$ is always a cyclic group of order $`[k:𝐅_2]`$ and therefore has order dividing $`N`$.
Clearly, $`\mathrm{End}_R(V)\mathrm{End}_{𝐅_2}(V)`$ is stable under the adjoint action of $`H`$. This induces a homomorphism
$$\alpha :H\mathrm{Aut}_{𝐅_2}(\mathrm{End}_R(V))=\mathrm{Aut}_{𝐅_2}(\mathrm{Mat}_d(k)).$$
Since $`k`$ is the center of $`\mathrm{Mat}_d(k)`$, it is stable under the action of $`H`$, i.e., we get a homomorphism $`H\mathrm{Aut}(k/𝐅_2)`$, which must be trivial, since $`H`$ is perfect and $`\mathrm{Aut}(k/𝐅_2)`$ is a cyclic group of order dividing $`N`$ and therefore the kernel of the homomorphism must coincide with $`H`$. This implies that the center $`k`$ of $`\mathrm{End}_R(V)`$ commutes with $`H`$. Since $`\mathrm{End}_H(V)=𝐅_2`$, we have $`k=𝐅_2`$. This implies that $`\mathrm{End}_R(V)\mathrm{Mat}_d(𝐅_2)`$ and one may rewrite $`\alpha `$ as
$$\alpha :H\mathrm{Aut}_{𝐅_2}(\mathrm{Mat}_d(𝐅_2))=\mathrm{Aut}(\mathrm{End}_{𝐅_2}(V_2))=\mathrm{Aut}_{𝐅_2}(V_2)/𝐅_2^{}=\mathrm{Aut}_{𝐅_2}(V_2).$$
It follows from the Jacobson density theorem that $`R=\mathrm{End}_{𝐅_2}(W)\mathrm{Mat}_m(𝐅_2)`$ with $`dm=N`$.
The adjoint action of $`H`$ on $`R`$ gives rise to a homomorphism
$$\beta :H\mathrm{Aut}_{𝐅_2}(\mathrm{End}_{𝐅_2}(W))=\mathrm{Aut}_{𝐅_2}(W)/𝐅_2^{}=\mathrm{Aut}_{𝐅_2}(W).$$
Clearly, $`\alpha `$ and $`\beta `$ provide $`V_2`$ and $`V_1`$ respectively with the structure of $`H`$-modules. Notice that
$$R=\mathrm{End}_{𝐅_2}(V_1)=\mathrm{End}_{𝐅_2}(V_1)\mathrm{Id}_{V_2}\mathrm{End}_{𝐅_2}(V_1)_{𝐅_2}\mathrm{End}_{𝐅_2}(V_2)=\mathrm{End}_{𝐅_2}(V).$$
Now our task boils down to comparison of the structures of $`H`$-module on $`V=V_1_{𝐅_2}V_2`$ defined by $`\rho `$ and $`\beta \alpha `$ respectively. I claim that
$$\rho (g)=\beta (g)\alpha (g)gH.$$
Indeed, notice that the conjugation by $`\rho (g)`$ in $`\mathrm{End}_{𝐅_2}(V)=\mathrm{End}_{𝐅_2}(V_1_{𝐅_2}V_2)`$ leaves stable $`R=\mathrm{End}_{𝐅_2}(V_1)_{𝐅_2}\mathrm{Id}_{V_2}`$ and coincides on $`R`$ with the conjugation by $`\alpha (g)\mathrm{Id}_{V_2}`$. Since the centralizer of $`\mathrm{End}_{𝐅_2}(V_1)\mathrm{Id}_{V_2}`$ in
$$\mathrm{End}_{𝐅_2}(V)=\mathrm{End}_{𝐅_2}(V_1)_{𝐅_2}\mathrm{End}_{𝐅_2}(V_2)$$
coincides with $`\mathrm{Id}_{V_1}\mathrm{End}_{𝐅_2}(V_2)`$, there exists $`u\mathrm{Aut}_{𝐅_2}(V_2)`$ such that
$$\rho (g)=\beta (g)u.$$
Since the conjugation by $`\rho (g)`$ leaves stable the centralizer of $`R`$, i.e. $`\mathrm{Id}_{V_1}\mathrm{End}_{𝐅_2}(V_2)`$ and coincides on it with the conjugation by $`\mathrm{Id}_{V_1}\alpha (g)`$, there exists a non-zero constant $`\gamma 𝐅_2^{}`$ such that $`u=\gamma \beta (g)`$. This implies that
$$\rho (g)=\beta (g)u=\gamma \beta (g)\alpha (g).$$
Now one has only to recall that $`𝐅_2^{}=\{1\}`$ and therefore $`\gamma =1`$. ∎
###### Remark 7.6.
In the notations of Th. 7.5 the $`H`$-modules $`V_1`$ and $`V_2`$ must be absolutely simple. It follows easily from Remarks 7.3.
Lemma 7.4 and Theorem 7.5 together with Remark 7.6 imply easily the following criterion of very simplicity over $`𝐅_2`$.
###### Theorem 7.7.
Let $`H`$ be a group and $`V`$ be a $`𝐅_2[H]`$-module of finite dimension $`N`$ over $`𝐅_2`$. Assume, in addition, that $`H`$ does not have nontrivial cyclic quotients of order dividing $`N`$ (e.g., $`H`$ is perfect).
Then $`V`$ is very simple if and only if the following conditions hold:
1. The $`H`$-module $`V`$ is absolutely simple;
2. There do not exist a subgroup $`H^{}H`$ of $`H`$ and a $`𝐅_2[H^{}]`$-module $`V^{}`$ such that $`V`$ is induced by $`V^{}`$;
3. There do not exist absolutely simple $`𝐅_2[H]`$-modules $`V_1`$ and $`V_2`$, both of dimension greater than $`1`$ and such that the $`H`$-module $`V`$ is isomorphic to $`V_1_{𝐅_2}V_2`$.
Combining Theorem 7.7 with Lemma 7.4 and Theorem 7.5 we get easily the following corollary.
###### Corollary 7.8.
Let $`H`$ be a group and $`V`$ be a $`𝐅_2[H]`$-module of finite dimension $`N`$ over $`𝐅_2`$. Then $`V`$ is very simple if the following conditions hold:
1. The $`H`$-module $`V`$ is absolutely simple;
2. $`H`$ does not contain a subgroup of finite index $`r`$ with $`rN`$ and $`1<r<N`$. In addition, $`H`$ does not have cyclic quotients of order $`N`$, i.e., $`H`$ does not have a normal subgroup $`H^{}`$ of index $`N`$ with cyclic quotient $`H/H^{}`$;
3. There do not exist absolutely simple $`𝐅_2[H]`$-modules $`V_1`$ and $`V_2`$, both of dimension greater than $`1`$ and such that the $`H`$-module $`V`$ is isomorphic to $`V_1_{𝐅_2}V_2`$.
The following assertion follows easily from Lemma 7.4 and Theorem 7.5.
###### Corollary 7.9.
Suppose a positive integer $`N>1`$ and a group $`H`$ enjoy the following properties:
* $`H`$ does not contain a subgroup of index dividing $`N`$ except $`H`$ itself.
* Let $`N=ab`$ be a factorization of $`N`$ into a product of two positive integers $`a>1`$ and $`b>1`$. Then either there does not exist an absolutely simple $`𝐅_2[H]`$-module of $`𝐅_2`$-dimension $`a`$ or there does not exist an absolutely simple $`𝐅_2[H]`$-module of $`𝐅_2`$-dimension $`b`$.
Then each absolutely simple $`𝐅_2[H]`$-module of $`𝐅_2`$-dimension $`N`$ is very simple. In other words, in dimension $`N`$ the properties of absolute simplicity and very simplicity over $`𝐅_2`$ are equivalent.
The next two theorems provide examples of very simple Steinberg representations.
###### Theorem 7.10.
Let $`q=2^m8`$ be an integral power of $`2`$, let $`B`$ be a $`(q+1)`$-element set. Let $`G^{}`$ be a group acting faithfully on $`B`$. Assume that $`G^{}`$ contains a subgroup $`G`$ isomorphic to $`𝐋_2(q)`$. Then the $`G^{}`$-module $`Q_B`$ is very simple.
###### Proof.
We have $`𝐋_2(q)=GG^{}\mathrm{Perm}(B)`$. Clearly, it suffices to check that the $`𝐋_2(q)`$-module $`Q_B`$ is very simple.
First, notice that $`𝐋_2(q)`$ acts doubly transitively on $`B`$. Indeed, each subgroup of $`𝐋_2(q)`$ (except $`𝐋_2(q)`$ itself) has index $`q+1=\mathrm{\#}(B)`$ (, (6.27), p. 415). This implies that $`𝐋_2(q)`$ acts transitively on $`B`$. If the stabilizer $`G_b`$ of a point $`bB`$ has index $`q+1`$ then it follows easily from Th. 6.25 on p. 412 of ) that $`G_b`$ in conjugate to the (Borel) subgroup of upper-triangular matrices and therefore the $`𝐋_2(q)`$-set $`B`$ is isomorphic to the projective line $`𝐏^1(𝐅_q)`$ with the standard action of $`𝐋_2(q)`$ which is well-known to be doubly (and even triply) transitive. By Remark 4.5, this implies that the $`𝐅_2[𝐋_2(q)]`$-module $`Q_B`$ is absolutely simple. Recall that
$$\mathrm{dim}_{𝐅_2}(Q_B)=\mathrm{\#}(B)1=q=2^m.$$
By Theorem 6.2, there no absolutely simple nontrivial $`𝐅_2[𝐋_2(q)]`$-modules of dimension $`<2^m`$. This implies that $`Q_B`$ is not isomorphic to a tensor product of absolutely simple $`𝐅_2[𝐋_2(q)]`$-modules of dimension $`>1`$. Recall that all subgroups in $`𝐋_2(q)`$ different from $`𝐋_2(q)`$ have index $`q+1>q=\mathrm{dim}_{𝐅_2}(Q_B)`$. It follows from Corollary 7.8 that the $`G`$-module $`Q_B`$ is very simple. Since $`GG^{}`$, the $`G^{}`$-module $`Q_B`$ is also very simple. ∎
###### Theorem 7.11.
Let $`k`$ be a positive integer and $`q=2^{2k+1}`$, let $`B`$ be a $`(q^2+1)`$-element set. Let $`G^{}`$ be a group acting faithfully on $`B`$. Assume that $`G^{}`$ contains a subgroup $`G`$ isomorphic to $`\mathrm{𝐒𝐳}(q)`$. Then the $`G^{}`$-module $`Q_B`$ is very simple.
###### Proof.
We have $`\mathrm{𝐒𝐳}(q)=GG^{}\mathrm{Perm}(B)`$. First, notice that $`\mathrm{𝐒𝐳}(q)`$ acts doubly transitively on $`B`$. Indeed, the classification of subgroups of Suzuki groups (, Remark 3.12(e), p. 194) implies that each subgroup of $`\mathrm{𝐒𝐳}(q)`$ (except $`\mathrm{𝐒𝐳}(q)`$ itself) has index $`q^2+1=\mathrm{\#}(B)`$. This implies that $`\mathrm{𝐒𝐳}(q)`$ acts transitively on $`B`$. If the stabilizer $`G_b`$ of a point $`bB`$ has index $`q^2+1`$ then it follows easily from the same classification that $`G_b`$ is conjugate to the subgroup $`𝔉`$ generated by all $`S(a,b)`$ and $`M(\lambda )`$ and therefore the $`\mathrm{𝐒𝐳}(q)`$-set $`B`$ is isomorphic to an ovoid $`𝒪=\mathrm{𝐒𝐳}(q)/𝔉`$ where the action of $`\mathrm{𝐒𝐳}(q)`$ is known to be doubly transitive (, Th. 3.3 on pp. 184–185 and steps g) and i) of its proof on p. 187). By Remark 4.5, this implies that the $`𝐅_2[\mathrm{𝐒𝐳}(q)]`$-module $`Q_B`$ is absolutely simple. Recall that $`\mathrm{dim}_{𝐅_2}(Q_B)=\mathrm{\#}(B)1=q^2=2^{2(2k+1)}`$. By Theorem 6.3, there no absolutely simple nontrivial $`𝐅_2[\mathrm{𝐒𝐳}(q)]`$-modules of dimension $`<2^{2(2k+1)}`$. This implies that $`Q_B`$ is not isomorphic to a tensor product of absolutely simple $`𝐅_2[\mathrm{𝐒𝐳}(q)]`$-modules of dimension $`>1`$. Recall that all subgroups in $`G=\mathrm{𝐒𝐳}(q)`$ (except $`\mathrm{𝐒𝐳}(q)`$ itself) have index $`q^2+1>q^2=\mathrm{dim}_{𝐅_2}(Q_B)`$. It follows from Corollary 7.8 that the $`G`$-module $`Q_B`$ is very simple. Since $`GG^{}`$, the $`G^{}`$-module $`Q_B`$ is also very simple. ∎
###### Proof of Theorem 5.4.
The cases (i) and (ii) of Theorem 5.4 follow from Theorems 7.10 and 7.11 respectively applied to $`G^{}=H`$. ∎
In light of Corollary 5.3 it would be interesting to classify all permutation subgroups $`G\mathrm{Perm}(B)`$ with very simple $`G`$-modules $`Q_B`$. We finish the paper by examples of very simple $`Q_B`$ attached to Mathieu groups $`\mathrm{M}_{11}`$ and $`\mathrm{M}_{12}`$ and to related group $`𝐋_2(11)=\mathrm{PSL}_2(11)`$.
###### Theorem 7.12.
Let $`n`$ be a positive integer, $`B`$ a $`n`$-element set, $`G\mathrm{Perm}(B)`$ a permutation group. Assume that $`(n,G)`$ enjoy one of the following properties:
1. $`n=11`$ and $`G`$ is isomorphic either to $`𝐋_2(11)`$ or $`\mathrm{M}_{11}`$;
2. $`n=12`$ and either $`G\mathrm{M}_{12}`$ or $`G\mathrm{M}_{11}`$ and $`G`$ acts transitively on $`B`$.
Then the $`G`$-module $`Q_B`$ is very simple.
###### Proof.
Assume that $`n=11`$. Since $`M_{11}`$ contains a subgroup isomorphic to $`𝐋_2(11)`$ (, p. 18), it suffices to check the case of $`G=𝐋_2(11)`$, in light of Remark 2.2(iii).
The group $`G=𝐋_2(11)`$ has two conjugacy classes of maximal subgroups of index $`11`$ and all other subgroups in $`G`$ have index greater than $`11`$ (, p. 7). Therefore all subgroups in $`G`$ (except $`G`$ itself) have index greater than $`10`$ and the action of $`G`$ on the $`11`$-element set $`B`$ is transitive. The permutation character (in both cases) is (in notations of , p. 7) $`1+\chi _5`$, i.e., $`\chi =\chi _5`$. The restriction of $`\chi _5`$ to the set of $`2`$-regular elements coincides with absolutely irreducible Brauer character $`\phi _4`$ (in notations of , p. 7). In particular, the corresponding $`G`$-module $`Q_B`$ is absolutely simple and has dimension $`10`$. Since $`10=25`$ and $`5`$ is a prime, the very simplicity of the $`G`$-module $`Q_B`$ follows from Th. 5.4 of . This proves the case (i).
Assume that $`n=12`$.
Suppose $`G=\mathrm{M}_{11}`$ and the action of $`G`$ on the $`12`$-element set $`B`$ is transitive. All subgroups $`G_b`$ in $`G`$ of index $`12`$ are isomorphic to $`𝐋_2(11)`$ (, p. 18). It follows from the already proven case (i) for $`𝐋_2(11)`$ and Remark 4.4 that the $`G`$-module $`Q_B`$ is very simple.
Suppose $`G=\mathrm{M}_{12}`$. The action of $`G`$ on the $`12`$-element $`B`$ is transitive, since all subgroups in $`G`$ (except $`G`$ itself) have index $`12`$. All subgroups $`G_b`$ in $`G`$ of index $`12`$ are isomorphic to $`\mathrm{M}_{11}`$ (, p. 33). It follows from the already proven case (i) for $`\mathrm{M}_{11}`$ and Remark 4.4 that the $`G`$-module $`Q_B`$ is very simple. ∎
Combining Corollary 5.3 and Theorem 7.12 (with $`B=,G=\mathrm{Gal}(f)`$ and taking into account that the irreducibility of $`f`$ means that $`\mathrm{Gal}(f)`$ acts transitively on $``$), we obtain the following statement.
###### Theorem 7.13.
Let $`K`$ be a field with $`\mathrm{char}(K)2`$, $`K_a`$ its algebraic closure, $`f(x)K[x]`$ an irreducible separable polynomial of degree $`n5`$. Let $`=_fK_a`$ be the set of roots of $`f`$, let $`K(_f)=K()`$ be the splitting field of $`f`$ and $`\mathrm{Gal}(f):=\mathrm{Gal}(K()/K)`$ the Galois group of $`f`$, viewed as a subgroup of $`\mathrm{Perm}()`$. Let $`C_f`$ be the hyperelliptic curve $`y^2=f(x)`$. Let $`J(C_f)`$ be its jacobian, $`\mathrm{End}(J(C_f))`$ the ring of $`K_a`$-endomorphisms of $`J(C_f)`$.
Assume that $`n`$ and $`\mathrm{Gal}(f)`$ enjoy one of the following properties:
1. $`n=11`$ and $`\mathrm{Gal}(f)`$ is isomorphic either to $`𝐋_2(11)`$ or to $`\mathrm{M}_{11}`$;
2. $`n=12`$ and $`\mathrm{Gal}(f)`$ is isomorphic either to $`\mathrm{M}_{11}`$ or to $`\mathrm{M}_{12}`$;
Then either $`\mathrm{End}(J(C_f))=𝐙`$ or $`\mathrm{char}(K)>0`$ and $`J(C_f)`$ is a supersingular abelian variety.
|
warning/0003/math0003159.html
|
ar5iv
|
text
|
# A remark on quiver varieties and Weyl groups
## 1. Notations and definitions
In this section we give the definition of quiver varieties. Except some minor change all definition are due to Nakajima \[Na1, Na2\].
### 1.1. The graph
Let $`(I,H)`$ be a finite oriented graph: $`I`$ is the set of vertices that we suppose of cardinality $`n`$, $`H`$ the set of arrows and the orientation is given by the two maps
$$hh_0\text{ and }hh_1$$
from $`H`$ to $`I`$. We suppose also that:
1. $`hHh_0h_1`$,
2. an involution $`h\overline{h}`$ of $`H`$ without fixed points and satisfying $`\overline{h}_0=h_1`$ is fixed,
3. a map $`\epsilon :H\{1,1\}`$ is given such that $`\epsilon (\overline{h})=\epsilon (h)`$. We define $`\mathrm{\Omega }=\{hH:\epsilon (h)=1\}`$ and $`\overline{\mathrm{\Omega }}=\{hH:\epsilon (h)=1\}`$.
Observe that given a symmetric graph without loops is always possible to define $`\epsilon `$ and an involution $`\overline{}`$ as above.
### 1.2. The Cartan matrix and the Weyl group
Let $`A`$ be the matrix whose entries are the numbers
$$a_{ij}=card\{hH:h_0=i\text{ and }h_1=j\}.$$
We define a generalized symmetric Cartan matrix by $`C=2IA`$. Following \[Lusztig:QG\] an $`X,Y`$-regular root datum $`(I,X,X^{},<,>)`$ with Cartan matrix equal to $`C`$ is defined in the following way:
1. $`X^{}`$ and $`X`$ are finetely generated free abelian groups,
2. $`<,>:X\times X^{}`$ is a perfect bilinear pairing,
3. two linearly independent sets $`\mathrm{\Pi }=\{\alpha _i:iI\}X`$ and $`\mathrm{\Pi }^{}=\{\alpha _i^{}:iI\}X^{}`$ are fixed and we set $`Q=\mathrm{\Pi }`$ and $`Q^{}=\mathrm{\Pi }^{}`$,
4. $`<\alpha _i,\alpha _j^{}>=c_{ij}`$,
5. (nonstandard) $`\mathrm{rank}X=\mathrm{rank}X^{}=2n\mathrm{rank}C`$,
6. (nonstandard) a linearly independent set $`\{\omega _i:iI\}`$ of $`X`$ such that $`<\omega _i,\alpha _j^{}>=\delta _{ij}`$ is fixed.
Once $`C`$ is given it is easy to construct a data as above. We call $`𝔥`$ the complexification of $`X^{}`$ and we observe that through the bilinear pairing $`<,>`$ we can identify $`𝔥^{}`$ with the complexification of $`X`$. We observe also that the triple $`(𝔥,\mathrm{\Pi },\mathrm{\Pi }^{})`$ is a realization of the Cartan matrix $`C`$ (\[Kac\] pg.1).
The Weyl group $`W`$ attached to $`C`$ is defined as the subgroup of $`Aut(X)GL(𝔥^{})`$ generated by the reflections
$$s_i:xx<x,\alpha _i^{}>\alpha _i.$$
(1)
Observe that the dual action is given by $`s_i(y)=y<\alpha _i,y>\alpha _i^{}`$ and that the lattices $`Q`$ and $`Q^{}`$ are stable for these actions. So the annihilator $`\stackrel{}{Q^{}}=\{xX:<x,y>=0yQ^{}\}`$ is also stable by $`W`$ and we can consider the action of $`W`$on the lattice $`P=X/\stackrel{}{Q^{}}\text{Hom}_{}(Q,)`$ and we call $`x\overline{x}`$ the projection from $`X`$ to $`P`$. We observe also that this projection is an isomorphism from the lattice $`\stackrel{~}{P}`$, that is not $`W`$-stable, spanned by $`\{\omega _i:iI\}`$ and $`P`$. Finally we observe that
$$\overline{\alpha }_i=\underset{jI}{}c_{ij}\overline{\omega }_j.$$
### 1.3. $`d,v`$ and the space of all matrices
For the exposition it will be usefull to identify the set $`I`$ with the set of integers $`\{1,\mathrm{},n\}`$.
Let $`d=(d_1,\mathrm{},d_n)`$ and $`v=(v_1,\mathrm{},v_n)`$ be two $`n`$-tuples of integers. We also think of $`d,v`$ as elements of $`X`$ in the following way:
$$d=\underset{iI}{}d_i\omega _i\text{ and }v=\underset{iI}{}v_i\alpha _i;$$
(2)
and through this identification we define an action of $`W`$ on $`v`$. We define also $`v^{}=_{iI}v_i\alpha _i^{}Q^{}`$. Once $`d,v`$ are fixed we fix complex vector spaces $`D_i`$ and $`V_i`$ of dimensions $`d_i`$ and $`v_i`$ and we define the following spaces of maps:
$`S_\mathrm{\Omega }(d,v)`$ $`={\displaystyle \underset{iI}{}}\text{Hom}(D_i,V_i){\displaystyle \underset{h\mathrm{\Omega }}{}}\text{Hom}(V_{h_0},V_{h_1}),`$ (3a)
$`S_{\overline{\mathrm{\Omega }}}(d,v)`$ $`={\displaystyle \underset{iI}{}}\text{Hom}(V_i,D_i){\displaystyle \underset{h\overline{\mathrm{\Omega }}}{}}\text{Hom}(V_{h_0},V_{h_1}),`$ (3b)
$`S(d,v)`$ $`=S_\mathrm{\Omega }(d,v)S_{\overline{\mathrm{\Omega }}}(d,v).`$ (3c)
More often, when it will not be ambiguous we will write $`S_\mathrm{\Omega },S_{\overline{\mathrm{\Omega }}}`$ and $`S`$ instead of $`S_\mathrm{\Omega }(d,v),S_{\overline{\mathrm{\Omega }}}(d,v)`$ and $`S(d,v)`$.
For each $`hH`$ (resp. $`iI`$) we define the projection $`B_h`$ (resp. $`\gamma _i`$ and $`\delta _i`$) from $`S`$ to $`\text{Hom}(V_{h_0},V_{h_1})`$ (resp. $`\text{Hom}(D_i,V_i)`$ and $`\text{Hom}(V_i,D_i)`$) with respect to the decomposition described in (3).
When an element $`s`$ of $`S`$ is fixed we will often write $`B_h`$ (resp. $`\gamma _i`$, $`\delta _i`$) instead of $`B_h(s)`$ (resp. $`\gamma _i(s)`$ and $`\delta _i(s)`$). We will also use $`\gamma `$ for $`(\gamma _1,\mathrm{},\gamma _n)`$, $`\delta `$ for $`(\delta _1,\mathrm{},\delta _n)`$ and $`B`$ for $`(B_h)_{hH}`$ and often we will write an element of $`S`$ as a triple $`(B,\gamma ,\delta )`$.
Once $`D_i,V_i`$ and an element $`s`$ of $`S`$ are fixed we define also:
$`T_i`$ $`=D_i{\displaystyle \underset{h:h_1=i}{}}V_{h_0},`$ (4a)
$`a_i`$ $`=a_i(s)=(\delta _i(s),(B_{\overline{h}}(s))_{h:h_1=i}):V_iT_i,`$ (4b)
$`b_i`$ $`=b_i(s)=(\gamma _i(s),(\epsilon (h)B_h(s))_{h:h_1=i}):T_iV_i.`$ (4c)
We will identify the dual of space of the $``$-linear maps $`\text{Hom}(E,F)`$ between two finite dimensional vector spaces with $`\text{Hom}(F,E)`$ through the pairing $`<\phi ,\psi >=\mathrm{Tr}(\phi \psi )`$. So we can describe $`S`$ also as $`S_\mathrm{\Omega }S_\mathrm{\Omega }^{}=T^{}S_\mathrm{\Omega }`$ and we observe that a natural symplectic structure $`\omega `$ is defined over $`S`$ by
$$\omega ((s_\mathrm{\Omega },s_{\overline{\mathrm{\Omega }}}),(t_\mathrm{\Omega },t_{\overline{\mathrm{\Omega }}}))=<s_\mathrm{\Omega },t_{\overline{\mathrm{\Omega }}}><t_\mathrm{\Omega },s_{\overline{\mathrm{\Omega }}}>.$$
### 1.4. Hermitian structure on $`S`$
We suppose now that the spaces $`D_i`$, $`V_i`$ are endowed with hermitian metrics. So we can speak of the adjoint $`\phi ^{}`$ of a linear map between these spaces, and we have a positive definite hermitian structure $`h`$ on $`S`$ with explicit formula:
$`h((B,\gamma ,\delta ),(\stackrel{~}{B},\stackrel{~}{\gamma },\stackrel{~}{\delta }))`$ $`={\displaystyle \underset{hH}{}}\mathrm{Tr}(B_h\stackrel{~}{B}_h^{})+{\displaystyle \underset{iI}{}}\mathrm{Tr}(\gamma _i\stackrel{~}{\gamma }_i^{}+\stackrel{~}{\delta }_i^{}\delta _i)`$ (5)
$`={\displaystyle \underset{iI}{}}\mathrm{Tr}(a_i\stackrel{~}{a}_i^{}+\stackrel{~}{b}_i^{}b_i)`$
and an associated real and closed symplectic form $`\omega _I(s,t)=\mathrm{Re}h(\text{i}s,t)\mathrm{Im}h(s,t)`$.
### 1.5. Group actions and moment maps
We can define an action of the groups $`G=GL(V)=GL(V_i)`$ and $`GL(D)=GL(D_i)`$ on the set $`S`$ in the following way:
$`g(B_h,\gamma _i,\delta _i)`$ $`=(g_{h_1}B_hg_{h_0}^1,g_i\gamma _i,\delta _ig_i^1)`$ $`\text{for }g=(g_i)GL(V),`$ (6)
$`g(B_h,\gamma _i,\delta _i)`$ $`=(B_h,\gamma _ig_i^1,g_i\delta _i)`$ $`\text{for }g=(g_i)GL(D).`$ (7)
Observe that these actions commute and that $`\omega `$ is $`GL(V)`$ invariant. Moreover if $`U=U(V)=U(V_i)`$ is the group of unitary trasformations in $`GL(V)`$ the real simplectic form $`\omega _I`$ is $`U(V)`$ invariant.
Define $`\mu ,\mu _I:S𝔤=gl(V_i)`$ by the following explicit formulas:
$`\mu _i(B,\gamma ,\delta )`$ $`={\displaystyle \underset{hH:h_1=i}{}}\epsilon (h)B_hB_{\overline{h}}+\gamma _i\delta _i=b_ia_i,`$
$`\mu _{I,i}(B,\gamma ,\delta )`$ $`={\displaystyle \frac{\text{i}}{2}}\left({\displaystyle \underset{hH:h_1=i}{}}B_hB_h^{}B_{\overline{h}}^{}B_{\overline{h}}+\gamma _i\gamma _i^{}\delta _i^{}\delta _i\right)={\displaystyle \frac{\text{i}}{2}}(b_ib_i^{}a_i^{}a_i),`$
If we identify $`𝔤^{}=\text{Hom}_{}(𝔤,)`$ (resp. $`𝔲^{}=\text{Hom}_{}(𝔲,)`$) with $`𝔤=gl(V_i)`$ (resp. $`𝔲`$) through the pairing $`<(x_i),(y_i)>=_i\mathrm{Tr}(x_iy_i)`$, we can observe that $`\mu `$ is a moment map for the action of $`G`$ on the symplectic manifold $`(S,\omega )`$ and that $`\mu _I`$ is a moment map for the action of $`U`$ on the symplectic manifold $`(S,\omega _I)`$. It is common to group all these moment maps together and to define an hyperKähler moment map
$$\stackrel{~}{\mu }=(\mu _I,\mu ):S𝔲𝔤=()_{}𝔲.$$
### 1.6. Quiver varieties as hyperKähler quotients
Let $`\zeta _i=(\xi _i,\lambda _i)`$ and $`\zeta =(\zeta _1,\mathrm{},\zeta _N)`$. We define:
$$𝔏_\zeta (d,v)=\{sS:\mu (s)\lambda _i\mathrm{Id}_{V_i}=0\text{ and }\mu _{I,i}(s)\text{i}\xi _i\mathrm{Id}_{V_i}=0\}.$$
We observe that $`𝔏_\zeta (d,v)`$ is stable for the action of $`U(V)`$, so, at least as a topological Hausdorf space we can define the *quiver variety of type* $`\zeta `$ as
$$𝔐_\zeta (d,v)=𝔏_\zeta (d,v)/U(V).$$
It will be convenient to define also $`𝔐_\zeta (d,v)=\mathrm{}`$ if $`d,v^n`$ and there exists $`i`$ such that $`v_i<0`$ or $`d_i<0`$ for some $`i`$. We call $`=^n^n`$ and we observe that we can idenyify it to $`()_{}P`$ through:
$$(\xi _1,\mathrm{},\xi _n,\lambda _1,\mathrm{},\lambda _n)\underset{iI}{}(\xi _i,\lambda _i)\overline{\omega }_i.$$
(8)
In particular we consider an action of the Weyl group $`W`$ on $``$ through this identification.
###### Remark 1.
There is a surjective map from $``$ to $`Z_U(𝔲)Z_G(𝔤)`$:
$$(\xi _1,\mathrm{},\xi _n,\lambda _1,\mathrm{},\lambda _n)\underset{iI}{}(\text{i}\xi _i,\lambda _i)\mathrm{Id}_{V_i}$$
Observe that $`𝔏_\zeta `$ is the fiber of $`\stackrel{~}{\mu }`$ over the image of $`\zeta `$ in $`Z_U(𝔲)Z_G(𝔤)`$.
###### Remark 2.
If $`v,d0`$ define: $`I^{}=\{iI:v_i0\}`$, $`H^{}=\{hH:h_0,h_1I^{}\}`$, $`\epsilon ^{}=\epsilon |_H^{}`$, $`v^{}=(v_i)_{iI^{}}`$, $`d^{}=(d_i)_{iI^{}}`$, and $`\zeta ^{}=(\zeta _i)_{iI^{}}`$ then it is clear that
$$𝔏_\zeta ^{}(d^{},v^{})𝔏_\zeta (d,v)\text{ and }𝔐_\zeta ^{}(d^{},v^{})𝔐_\zeta (d,v).$$
Except for the last equivalence which is trivial in our case the following is a general well known fact (\[GIT\] ch. 8).
###### Lemma 3.
Let $`s𝔏_\zeta `$ then
$$d\stackrel{~}{\mu }_s\text{ is surjective}d\mu _s\text{ is surjective}d\mu _I\text{ is surjective }$$
$$dim\mathrm{Stab}_G\{s\}=0\mathrm{Stab}_G\{s\}=\{1_G\}$$
###### Definition 4.
If $`u^n=Q^{}`$ $`AQ^{}`$ we define
$$_u=\{\zeta =(\xi ,\lambda ):<\xi ,u^{}>=<\lambda ,u^{}>=0\}\text{ and }_A=\underset{uA}{}H_u.$$
Let now $`U_v=\{u^n\{0\}\text{ such that }`$ $`0u_iv_i\}`$ and $`=_{U_v}`$. $``$ is a union of a finite number of real subspace of $``$ of codimension $`3`$.
###### Lemma 5 (Nakajima, \[Na1\]).
If $`\zeta `$ and $`\stackrel{~}{\mu }(s)=\zeta `$ then $`\mathrm{Stab}_G\{s\}=1_G`$.
As a consequence of the above lemma and general results on on hyperKähler manifolds (for example \[Hit\] or \[HKLR\]) we obtain the following corollary.
###### Corollary 6.
If $`\zeta 𝔷`$ then if it is not empty $`𝔐_\zeta (d,v)`$ is a smooth hyperKähler manifold of real dimension $`2<2dv,v^{}>`$.
### 1.7. Geometric invariant theory and moment map
In this section we explain the relation between the moment map and the GIT quotient proved by Kempf, Ness \[KeNe\], Kirwan \[GIT\] and others. To be more precise we need a generalization of their results in the case of an action on an affine variety proved by Migliorini \[Migliorini\].
Let $`X`$ be an affine variety over $``$ and $`G`$ a reductive group acting on $`X`$. We can assume that $`X`$ is a closed subvariety of a vector space $`V`$ where $`G`$ acts linearly. Let $`h`$ be an hermitian form on $`V`$ invariant by the action of a maximal compact group $`U`$ of $`G`$ and define a real $`U`$-invariant symplectic form on $`V`$ by
$$\eta (x,y)=\mathrm{Re}h(\text{i}x,y).$$
Then we can define a moment map $`\nu :V𝔲^{}=\text{Hom}_{}(𝔲,)`$:
$$<\nu (x),u>=\frac{1}{2}\eta (ux,x).$$
We observe that the real symplectic form $`\eta `$ resctricted to a complex submanifold is always non degenerate and that $`\mu `$ restricted to the non singular locus of $`X`$ is a moment map for the action of $`U`$ on $`X`$.
Now let $`\chi `$ be a multiplicative character of $`G`$. We observe that for all $`gU`$ we have $`|\chi (g)|=1`$ so $`\text{i}d\chi :u`$. In particular we can think to $`\text{i}d\chi `$ as an element of $`𝔲^{}`$. Morover we observe that it is invariant by the dual adjoint action, hence it makes sense to consider the quotient:
$$𝔐=\nu ^1(id\chi )/U.$$
As we saw our variety are a particular case of this construction.
On the other side we can consider the GIT quotient. Let us remind the definition. If $`\phi `$ is a character of $`G`$ we consider the line bundle $`L_\phi =V\times `$ on $`V`$ with the following $`G`$-linearization:
$$g(x,z)=(gx,\phi (g)z).$$
An invariant section of $`L_\phi `$ is determined by an algebraic function $`f:V`$ such that $`f(gx)=\phi (g)f(x)`$ for all $`gG`$ and $`xV`$. We use the same symbol $`L_\phi `$ also for the restriction of $`L_\phi `$ to $`X`$.
Given a rational action of $`G`$ on $``$-vector space $`A`$ we define
$$A_{\phi ,n}=\{aA:ga=\phi ^n(g)a\text{ for all }gG\},$$
$$A_\phi =\underset{n=0}{\overset{\mathrm{}}{}}A_{\phi ,n}\text{ as a graded vector space.}$$
Hence we have that $`H^0(X,L_\phi )^G=[X]_{\phi ,1}`$. We observe that if $`I`$ is the ideal of algebraic function on $`V`$ vanishing on $`X`$ then
$$H^0(X,L_\phi )^G=\frac{H^0(V,L_\phi )^G}{I_{\phi ,1}}.$$
This last fact can be proved easily for example averaging a $`\phi `$ equivariant function $`f`$ on $`X`$ in the following way:
$$\stackrel{~}{f}(v)=_U\phi ^1(u)f(uv)𝑑u.$$
###### Definition 7.
A point $`x`$ of $`X`$ is said to be $`\chi `$-*semistable* if there exist $`n>0`$ and $`fH^0(X,L_\chi ^n)^G`$ such that $`f(x)0`$. We observe that by the remark above a point of $`X`$ is $`\chi `$-semistable if and only if is $`\chi `$-semistable as a point of $`V`$. We call $`X_\chi ^{ss}`$ (resp. $`V_\chi ^{ss}`$) the open subset of $`\chi `$-semistable points of $`X`$ (resp. $`V`$).
###### Proposition 8 ( \[GIT, Newstead\]).
There exists a good quotient of $`X_\chi ^{ss}`$ by the action of $`G`$ and we have that
$$X_\chi ^{ss}//G=\mathrm{Proj}[X]_\chi .$$
Moreover $`\mathrm{Proj}[X]_\chi `$ is a finetely generated $``$-algebra and a natural projective map
$$\pi :X_\chi ^{ss}//GX//G=\mathrm{Spec}[X]^G$$
is defined.
In the case of $`\chi 1`$ the following fact is well known:
$$\mathrm{Proj}[X]_\chi =\mathrm{Spec}[X]^G=\nu ^1(0)/U.$$
The following result is less well known, and its proof requires some adjustment of the classical proof for the case $`\chi 1`$ (see for example an appendix of \[Migliorini\] or \[TatoPhD\] ).
###### Proposition 9 (Migliorini, \[Migliorini\]).
Let $`xX`$ then
$$gG:\nu (gx)=\text{i}d\chi Gx\text{ is a closed orbit in }X_\chi ^{ss}.$$
###### Proposition 10 (Migliorini, \[Migliorini\]).
The inclusion $`\nu ^1(\text{i}d\chi )X_\chi ^{ss}`$ induces an homeomorphism
$$\nu ^1(\text{i}d\chi )/UX_\chi ^{ss}//G.$$
### 1.8. Quiver varieties as algebraic varieties
If $`m=(m_1,\mathrm{},m_n)^N`$ we define a character $`\chi _m`$ of $`G_v`$ by $`\chi _m=_{iI}det_{GL(V_I)}^{m_i}`$.
If $`\lambda =(\lambda _1,\mathrm{},\lambda _n)^n`$ and $`m=(m_1,\mathrm{},m_n)^n`$ then we define the varieties:
$`\mathrm{\Lambda }_\lambda (d,v)`$ $`=\{sS:\mu _i(s)\lambda _i\mathrm{Id}_{V_i}=0\text{ for all }i\},`$
$`\mathrm{\Lambda }_{m,\lambda }(d,v)`$ $`=\{s\mathrm{\Lambda }_\lambda (d,v):s\text{ is }\chi _m\text{semistable}\}.`$
and the associeted *quiver varieties*
$`M_{m,\lambda }(d,v)`$ $`=\mathrm{\Lambda }_{m,\lambda }(d,v)//G_v\text{ and }`$
$`M_\lambda ^0(d,v)`$ $`=M_{0,\lambda }(d,v)=\mathrm{\Lambda }_\lambda (d,v)//G_v`$
We call $`p_{m,\lambda }^{d,v}:\mathrm{\Lambda }_{m,\lambda }(d,v)M_{m,\lambda }(d,v)`$ the quotient map. Observe that the inclusion $`\mathrm{\Lambda }_{m,\lambda }(d,v)\mathrm{\Lambda }_\lambda (d,v)`$ induces a projective morphism
$$\pi _{m,\lambda }^{d,v}:M_{m,\lambda }(d,v)M_\lambda ^0(d,v).$$
Finally it will be convenient to define $`M_{m,\lambda }(d,v)=\mathrm{}`$ if $`d,v^n`$ and $`v_i<0`$ or $`d_i<0`$ for some $`i`$.
We identify $`^n`$ with $`P`$ and $`Z=^n`$ with $`_{}P`$ through
$$(m_1,\mathrm{},m_n)m_i\overline{\omega }_i\text{ and }(\lambda _1,\mathrm{},\lambda _n)\lambda _i\overline{\omega }_i.$$
###### Remark 11.
As in 1 we have a surjective map from $`Z`$ to $`Z(𝔤)`$ and $`\mathrm{\Lambda }_\lambda (d,v)`$ is the fiber of $`\mu `$ over the image of $`\lambda `$ in $`Z(𝔤)`$.
###### Remark 12.
Remark 2 holds without changes also in this case.
###### Remark 13.
Observe that $`PZ`$. Observe also that the map $`m\chi _m`$ define a surjective morphism from $`P`$ to $`\text{Hom}(G_v,^{})`$ and that the following diagram commute:
In particular we can apply 10 to the action of $`G_v`$ on $`\mathrm{\Lambda }_\lambda (d,v)`$ and we obtain:
$$𝔐_{(m,\lambda )}(d,v)M_{m,\lambda }(d,v).$$
###### Proposition 14.
Let $`(m,\lambda )`$ and $`s\mathrm{\Lambda }_{m,\lambda }(d,v)`$ then $`\mathrm{Stab}_{G_v}(s)=\{1\}`$.
###### Proof.
As we have already claimed it is enough to prove $`dim\mathrm{Stab}_G(s)=\{1\}`$. We know that there is a good quotient of $`\mathrm{\Lambda }_{m,\lambda }(d,v)`$ so it is enough to prove that any closed orbit has maximal dimension. By Proposition 9 if $`G_vs`$ is closed in $`\mathrm{\Lambda }_{m,\lambda }(d,v)`$ then there exists $`gG_v`$ such that $`\mu _I(gs)=\text{i}d\chi `$. The thesis follows now form $`(m,\lambda )`$ and Lemma 5. ∎
###### Corollary 15.
If $`(m,\lambda )`$ and $`M_{m,\lambda }(d,v)\mathrm{}`$ then it is a smooth algebraic variety of dimension $`<v^{},2dv>`$.
### 1.9. Path algebra and $`b`$-path algebra
To describe functions on quiver varieties we need some notation about the path algebra.
###### Definition 16.
A *path* $`\alpha `$ in our graph is a sequence $`h^{(m)}\mathrm{}h^{(1)}`$ such that $`h^{(i)}H`$ and $`h_1^{(i)}=h_0^{(i+1)}`$ for $`i=1,\mathrm{},m1`$. We define also $`\alpha _0=h_0^{(1)}`$, $`\alpha _1=h_1^{(m)}`$ and we say that the length of $`\alpha `$ is $`m`$. If $`\alpha _0=\alpha _1`$ we say that $`\alpha `$ is a closed path. We consider also the empty paths $`\mathrm{}_i`$ for $`iI`$ and we define $`(\mathrm{}_i)_0=(\mathrm{}_i)_1=i`$. The product of path is defined in the obvious way.
A $`b`$-*path* $`[\beta ]`$ in our graph is a sequence $`[i_{m+1}^{r_{m+1}}\alpha ^{(m)}i_m^{r_m}\mathrm{}\alpha ^{(1)}i_1^{r_1}]`$, that we write between square brackets such that $`i_jI`$, $`\alpha ^{(j)}`$ are $`B`$-path, $`r_j`$ and $`\alpha _0^{(j)}=i_j`$ and $`\alpha _1^{(j)}=i_{j+1}`$ for $`j=1,\mathrm{},m`$. We consider also the “empty” $`b`$-paths indiced by elements of $`I`$: $`[\mathrm{}_i]`$. We define $`[\beta ]_0=i_1`$, $`[\beta ]_1=i_{m+1}`$ and $`[\mathrm{}_i]_0=[\mathrm{}_i]_1=i`$. The length of $`[\beta ]`$ is $`_{j=1}^{m+1}r_j+_{j=1}^mlength(\alpha ^j)`$ and the product of $`b`$-paths is defined in the obvious way:
$$[\beta ][\beta ^{}]=\{\begin{array}{cc}0\hfill & \text{if }[\beta ^{}]_1[\beta ]_0\hfill \\ [\beta \beta ^{}]\hfill & \text{if }[\beta ^{}]_1=[\beta ]_0=i\hfill \end{array}$$
Given a path $`\alpha =h^{(m)}\mathrm{}h^{(1)}`$ and a $`b`$-path $`\beta =[i_{m+1}^{r_{m+1}}\alpha ^{(m)}\mathrm{}`$ $`\mathrm{}i_1^{r_1}]`$ we define an evaluation of $`\alpha `$ and $`\beta `$ on $`S`$ in the following way: if $`s=(B,\gamma ,\delta )S`$ then
$`\mathrm{}_i(s)`$ $`=0\text{Hom}(V_i,V_i)\text{ and }[\mathrm{}_i](s)=0\text{Hom}(V_i,V_i),`$
$`\alpha (s)`$ $`=B_{h^{(m)}}\mathrm{}B_{h^{(1)}}\text{Hom}(V_{\alpha _0},V_{\alpha _1}),`$
$`\beta (s)`$ $`=(\gamma _{i_{m+1}}\delta _{i_{m+1}})^{r_{m+1}}\alpha ^{(m)}(s)(\gamma _{i_m}\delta _{i_m})^{r_m}\mathrm{}`$
$`\mathrm{}\alpha ^{(1)}(s)(\gamma _{i_1}\delta _{i_1})^{r_1}\text{Hom}(D_{\beta _0},D_{\beta _1}).`$
The path algebra $``$ is the vector space spanned by paths with the product induced by the product of paths. If $`i,jI`$ we say that an element in $``$ is of type $`(i,j)`$ if it is in the linear span of the paths with source in $`i`$ and target in $`j`$.
The $`b`$-path algebra $`𝒬`$ is the vector space spanned by $`b`$-paths with the product induced by the product of $`b`$-paths described above. If $`i,jI`$ we say that an element in $``$ is of type $`(i,j)`$ if it is in the linear span of the $`b`$-paths with source in $`i`$ and target in $`j`$.
###### Remark 17.
We observe that the evaluation on $`S`$ is a morphism of algebra from $``$ to the algebra defined by the morphisms of the category of vector spaces. Moreover if $`f`$ is of type $`(i,j)`$ we observe that $`f(s)\text{Hom}(V_i,V_j)`$.
## 2. Generators of the projective ring of a quiver variety
In this section we want to describe a set of generators of the graded ring $`[S]_\chi `$ and by consequence of the projective ring of a quiver variety $`[\mathrm{\Lambda }_\lambda ]_\chi `$. More precisely we will give a set of generators as $`[S]^G`$-module of its $`l`$-homogeneous part: $`[S]_{\chi ,l}`$. This result is a generalization of the one obtained by Lusztig in the case of invariants: $`\chi 1`$. First of all recall his result.
###### Theorem 18 (Lusztig, \[Lu:Q3\] theorem 1.3).
The ring $`[S]^G`$ is generated by the polynomials:
$$s\mathrm{Tr}\left(\alpha (s)\right)\text{ and }s\phi \left(\delta _{\beta _1}(s)\beta (s)\gamma _{\beta _0}(s)\right)$$
for $`\alpha `$ a closed path, $`\beta `$ a path and $`\phi \left(\text{Hom}(D_{\beta _0},D_{\beta _1})\right)^{}`$.
#### 2.0.1. Determinants
To describe our result we do first some general remark. Forget for a moment our quiver, and suppose to have a finite set of finite dimensional vector spaces $`X_1,\mathrm{},X_k`$ of dimensions $`u_1,\mathrm{},u_k`$ and a pair of nonnegative integers $`(m_i^+,m_i^{})`$ for each of them. Finally let $`m^+,m^{}`$ two nonnegative integers such that
$$N=\underset{i=1}{\overset{k}{}}m_i^+u_i+m^+=\underset{i=1}{\overset{k}{}}m_i^{}u_i+m^{},$$
and two vector spaces $`M^+`$ and $`M^{}`$ of dimension $`m^+,m^{}`$. Construct the vector spaces:
$$Y=\underset{i=1}{\overset{k}{}}^{m_i^{}}X_iM^{},Z=\underset{i=1}{\overset{k}{}}^{m_i^+}X_iM^+$$
and observe that $`dimY=dimZ=N`$. Define an action of the general linear group $`GL(X_i)`$ of $`X_i`$ on $`Y`$ by
$$g_i(\underset{j=1}{\overset{k}{}}v_jx_j+m)=\underset{ji}{}v_jx_j+m+v_ig_ix_i,$$
and also a similar action on $`Z`$. Hence the vector space $`\text{Hom}(Y,Z)`$ acquires a natural structure of $`G_X=_{i=1}^kGL(X_i)`$ module. If we choose an isomorphism $`\sigma `$ between $`\text{Hom}(^NY,^NZ)`$ and $``$ we can define a function $`det`$ on $`\text{Hom}(Y,Z)`$ by
$$det(A)=\sigma \left(\stackrel{n}{}A\right).$$
For simplicity we do not emphasize the role of $`\sigma `$ on this definition, so strictly speaking, $`det`$ is a function defined only up to a nontrivial constant factor. We observe also that
$`^nY\left(^{u_i}X_1\right)^{m_1^{}}\mathrm{}\left(^{u_k}X_k\right)^{m_k^{}}^m^{}M^{}`$
(and similarly for $`Z`$) so an isomorphism $`\sigma `$ is determined if we choose orientations, or basis, of $`X_j,M^+,M^{}`$. Finally observe that for any $`g=(g_j)G_X`$ we have
$$det(gA)=\underset{i=1}{\overset{k}{}}(det_{GL(X_i)}(g_i))^{m_i^+m_i^{}}det(A).$$
#### 2.0.2. Description of generators
We go back now to our quiver and we describe a set of covariant polynomials on $`S`$. Any character $`\chi `$ of the group $`G_v=GL(V)`$ is of the form $`\chi =\chi _m=_{iI}det_{GL(V_i)}^{m_i}`$. We fix such a character and we define
$`I^+`$ $`=\{iI:m_i>0\}\text{ and }\stackrel{~}{m}_i=m_i\text{ if }iI^+,`$
$`I^0`$ $`=\{iI:m_i=0\}\text{ and }\stackrel{~}{m}_i=0\text{ if }iI^0,`$
$`I^{}`$ $`=\{iI:m_i<0\}\text{ and }\stackrel{~}{m}_i=m_i\text{ if }iI^{}.`$
We use now the construction explained in 2.0.1 in the case $`X_i=V_i`$ and $`m_i^+m_i^{}=m_i`$. We choose ordered sets $`A=(a_1,\mathrm{},a_m^{})(D_i)^m^{}`$ and $`B=(b_1,\mathrm{},b_{m^+})(D_i^{})^{m^+}`$ and we define a function $`I:A,BI`$ by $`aD_{I(a)}`$, $`bD_{I(b)}^{}`$. In the framework described above it is then possible to set $`M^{}=_{i=1}^m^{}_{a_i}`$ and $`M^+=_{i=1}^{m^+}_{b_i}`$. In particular we have
$$Y=\underset{iI}{}\underset{h=1}{\overset{m_i^{}}{}}V_i^{(h)}\underset{i=1}{\overset{m^{}}{}}_{a_i},Z=\underset{iI}{}\underset{k=1}{\overset{m_i^+}{}}V_i^{[k]}\underset{i=1}{\overset{m^+}{}}_{b_i}$$
where $`V_i^{(l)},V_i^{[l]}`$ are isomorphic copies of $`V_i`$. We choose now elements of the $`b`$-path algebra as follows:
1. for any $`i,jI`$ and for any $`1hm_i^{}`$, $`1km_j^+`$ we choose an element $`\alpha _{j,k}^{i,h}`$ of the $`b`$-path algebra of type $`(i,j)`$,
2. for any $`iI`$, $`1hm_i^{}`$ and for any $`1lm^+`$ we choose an element $`\alpha _l^{i,h}`$ of the $`b`$-path algebra of type $`(i,I(b_l))`$,
3. for any $`1lm^{}`$ and for any $`jI`$, $`1km_j^+`$ we choose an element $`\alpha _{j,k}^l`$ of the $`b`$-path algebra of type $`(I(a_l),j)`$,
4. for any $`1lm^{}`$ and for any $`1l^{}m^{}`$ we choose an element $`\alpha _l^{}^l`$ of the $`b`$-path algebra of type $`(I(a_l),I(b_l^{}))`$.
We call such a data $`\mathrm{\Delta }=(\{(m_i^+,m_i^{})\}_{iI},(m^+,m^{}),A,B,\alpha _{j,k}^{i,h},\alpha _l^{i,h},\alpha _{j,k}^l,\alpha _l^{}^l)`$ a $`\chi `$-data and we attach to it a $`\chi `$-covariant function on $`S`$:
$$f_\mathrm{\Delta }(s)=det\left(\mathrm{\Psi }_\mathrm{\Delta }(s)\right)$$
where $`\mathrm{\Psi }_\mathrm{\Delta }`$ is a linear map from $`Y`$ to $`Z`$ defined by
$`[\mathrm{\Psi }_\mathrm{\Delta }]_{V_j^{[k]}}^{V_i^{(h)}}(s)`$ $`=\alpha _{j,k}^{i,h}(s),`$
$`[\mathrm{\Psi }_\mathrm{\Delta }]_{_{b_l}}^{V_i^{(h)}}(s)`$ $`=b_l\delta _{I(b_l)}\alpha _l^{i,h}(s),`$
$`[\mathrm{\Psi }_\mathrm{\Delta }]_{V_j^{[k]}}^{_{a_l}}(s)`$ $`=\alpha _{j,k}^l(s)\gamma _{I(a_l)}|_{a_l},`$
$`[\mathrm{\Psi }_\mathrm{\Delta }]_{_{b_l^{}}}^{_{a_l}}(s)`$ $`=b_l^{}\delta _{I(b_l^{})}\alpha _l^{}^l(s)\gamma _{I(a_l)}|_{a_l}.`$
The function $`f_\mathrm{\Delta }`$ are a set of generators as $`[S]^G`$-module of $`[S]_{\chi ,1}`$, but we will need to define a smaller set of generators. To define this set we give a notion of good $`\mathrm{\Delta }`$.
###### Definition 19.
A data $`\mathrm{\Delta }`$ as above is said to be $`\chi `$-*good* if it satisfies the following conditions:
1. $`m_i^++m_i^{}=\stackrel{~}{m}_i`$ for all $`iI`$,
2. $`\alpha _l^{}^l=0`$ for all $`l,l^{}`$,
3. $`\alpha _{}^{}`$ is an element of the path algebra (and not just an element of the $`b`$-path algebra which is obviously bigger),
4. $`card\{(j,k):\alpha _{j,k}^{i,h}0\}+card\{l:\alpha _l^{i,h}0\}v_i`$ for all $`i,h`$,
5. $`card\{(i,h):\alpha _{j,k}^{i,h}0\}+card\{l:\alpha _{j,k}^l0\}v_j`$ for all $`j,k`$,
6. for all $`l`$ there exists at most one pair $`(i,h)`$ such that $`\alpha _l^{i,h}0`$,
7. for all $`l`$ there exists at most one pair $`(j,k)`$ such that $`\alpha _{j,k}^l0`$.
For the applications the only important point will be the first one.
###### Proposition 20.
The set of polynomials $`f_\mathrm{\Delta }`$ with $`\mathrm{\Delta }`$ $`\chi `$-good generates $`[S]_{\chi ,1}`$ as a $`[S]^{G_v}`$-module.
###### Remark 21.
Prof. Weyman said me that in the case $`D=0`$ a similar proposition has been proved by him and for arbitrary characteristic.
### 2.1. Some remark on the invariant theory of $`GL(n)`$
If $`V`$ is a finite dimensional representation of a linearly reductive Lie group $`G`$ and $`S`$ is a simple representation of $`S`$ we write $`V[S]`$ for the $`S`$-isotipic component of type $`S`$ of $`V`$.
We now fix $`n`$ and we make some remark on the representations of $`GL(n)`$. To any partition of height less or equal to $`n`$ we associate an irreducible representation of $`GL(n)`$ in the usual way. If we multiply these representations by a power of the inverse of determinant representation we obtain a complete list of irreducible representations of $`GL(n)`$. If $`\lambda `$ is a partition we call $`\lambda `$ the transpose partition as usual and we define $`\lambda ^{op}=(\lambda _1\lambda _n,\lambda _1\lambda _{n1},\mathrm{},\lambda _1\lambda _1)`$. We call $`\delta `$ the determinant representation of $`GL(n)`$ and $`\epsilon =1^n`$ the associated partition. Finally we call $`V`$ the natural representation.
###### Lemma 22.
1. $`L_{\lambda }^{}{}_{}{}^{}=\delta ^{\lambda _1}L_{\lambda ^{op}}`$,
2. $`\text{Hom}_{GL(n)}(\delta ^m,L_\lambda L_\mu )=\{\begin{array}{cc}\hfill & \text{ if }\mathrm{\lambda }=\mu ^{op}+(m\mu _1)\epsilon ,\hfill \\ 0\hfill & \text{ otherwise,}\hfill \end{array}`$
3. $`\text{Hom}_{GL(n)}(\delta ^m,L_\lambda L_\mu ^{})=\{\begin{array}{cc}\hfill & \text{ if }\mathrm{\lambda }=\mu +m\epsilon ,\hfill \\ 0\hfill & \text{ otherwise.}\hfill \end{array}`$
###### Proof.
We prove only 2).
$`Hom_{GL(n)}(\delta ^m,L_\lambda L_\mu )`$ $`=Hom_{GL(n)}(\delta ^mL_{\mu }^{}{}_{}{}^{},L_\lambda )`$
$`=Hom_{GL(n)}(\delta ^{m\mu _1}L_{\mu ^{op}},L_\lambda )`$
If $`m\mu _1`$ the last group is isomorphic to $`Hom_{GL(n)}(L_{\mu ^{op}+(m\mu _1)\epsilon },L_\lambda )`$ and if $`m<\mu _1`$ is isomorphic to $`Hom_{GL(n)}(L_{\mu ^{op}};L_{\lambda +(\mu _1m)\epsilon })`$. In any case the thesis follows. ∎
We want now to describe $`\text{Hom}_{GL(n)}(\delta ^m,V^i\left(V^{}\right)^j)`$. To do it we will use Schur-duality. Remind that the irreducible representations of the groups $`S_m`$ are parametrized by partitions of $`m`$ and call $`S_\lambda `$ the irreducible representation associated with $`\lambda `$. Consider now the action of $`S_m`$ on $`V^m`$ given by permuting the factors. This action commute with the $`GL(n)`$ action. Schur duality asserts that the action of the group $`S_m\times GL(n)`$ on $`V^m`$ decomposes in the following way:
$$V^m=\underset{\begin{array}{c}\lambda m\\ ht(\lambda )n\end{array}}{}S_\lambda L_\lambda .$$
We describe a set of elements of $`\text{Hom}_{GL(n)}(\delta ^m,V^i\left(V^{}\right)^j)`$. Let $`m`$ be a nonnegative integers a choose a permutation $`\sigma `$ of $`\{1,\mathrm{},i+mn\}`$. To $`\sigma `$ we associate maps:
$`\mathrm{\Phi }_\sigma `$ $`:\left(V^i(V^{})^i\right)^{GL(n)}V^{i+mn}(V^{})^i[\delta ^m]`$
$`\mathrm{\Psi }_\sigma `$ $`:\left(V^i(V^{})^i\right)^{GL(n)}V^i(V^{})^{i+mn}[\delta ^m]`$
by
$`\mathrm{\Phi }_\sigma (ts)`$ $`=\sigma (o\mathrm{}ot)s`$
$`\mathrm{\Psi }_\sigma (ts)`$ $`=t\sigma (o^{}\mathrm{}o^{}s)`$
where $`o`$ is a nonzero vector in $`^nV`$ and $`o^{}`$ is a non zero vector in $`^nV^{}`$ and $`tV^i`$, $`s(V^{})^i`$.
###### Lemma 23.
1) If $`ij+mn`$ then
$$\text{Hom}_{GL(N)}(\delta ^m,V^i\left(V^{}\right)^j)=0.$$
2) If $`m>0`$ then
$$V^{i+mn}(V^{})^i[\delta ^m]=\underset{\sigma }{}\mathrm{Im}\mathrm{\Phi }_\sigma .$$
3) If $`m>0`$ then
$$V^i(V^{})^{i+mn}[\delta ^m]=\underset{\sigma }{}\mathrm{Im}\mathrm{\Psi }_\sigma .$$
###### Proof.
1) follows directly from lemma 22.
2) Let $`M=V^{i+mn}(V^{})^i[\delta ^m]`$ and $`N=(V^i(V^{})^i)^G`$. $`N`$ is a $`S_i\times S_i`$ module, $`M`$ is a $`S_{i+mn}\times S_i`$\- module and the maps $`\mathrm{\Phi }_{,\sigma }`$ are equivariant with respect the $`S_i`$ action on $`(V^{})^i`$. In particular it is enough to prove that if $`\lambda `$ is a partition of $`i`$, $`M_\lambda `$ is the $`S_\lambda `$-isotipic component of $`M`$ w.r.t. the $`S_i`$ action and $`N_\lambda `$ the $`S_\lambda `$-isotipic component of $`N`$ w.r.t. the $`S_i`$ action on $`(V^{})^i`$ then
$$M_\lambda =\underset{\sigma }{}\mathrm{\Phi }_\sigma (N_\lambda ).$$
By point 3 of lemma22 we have that
$`M`$ $`={\displaystyle \underset{\begin{array}{c}\lambda i\\ ht(\lambda )n\end{array}}{}}\left(S_{\lambda +m\epsilon }L_{\lambda +m\epsilon }\right)\left(S_\lambda L_\lambda \right)^{}[\delta ^m]`$
$`={\displaystyle \underset{\begin{array}{c}\lambda i\\ ht(\lambda )n\end{array}}{}}S_{\lambda +m\epsilon }S_\lambda \left(\delta ^m\left(L_\lambda L_{\lambda }^{}{}_{}{}^{}\right)^G\right)`$
In particular $`M_\lambda =S_{\lambda +m\epsilon }S_\lambda `$ is an irreducible representation of $`S_{i+mn}\times S_i`$. Observe $`_\sigma \mathrm{\Phi }(N_\lambda )`$ is a $`S_{i+mn}\times S_i`$-submodule of $`M_\lambda `$ and that it is clearly nonzero. So $`M_\lambda =_\sigma \mathrm{\Phi }_\sigma (N_\lambda )`$ as claimed.
The proof of 3) is equal to the previous one. ∎
We want now to give a slightly different formulation of the lemma above. Let $`M=V^i(V^{})^j`$ we want to describe $`M_{\delta ^m}^{}=\{\phi M^{}:g\phi =\delta ^m(g)\phi \}`$. Of course this problem is completely equivalent to the previous one. What we want to do is to reformulate in a more convenient way for our purposes the description of a set of generators of $`M_{\delta ^m}^{}`$. Let $`m0`$ and choose $`=\{I_1,\mathrm{},I_m\}`$ a collection of $`m`$ disjoint subsets of $`\{1,\mathrm{},i+mn\}`$ of cardinality $`n`$. Let $`I_j=\{i_{j1}<\mathrm{}<i_{jn}\}`$ and $`\{1,\mathrm{},i+mn\}=\{j_1<\mathrm{}<j_i\}`$. To $``$ and to a permutation $`\sigma S_i`$ we associate elements
$$\varphi _{,\sigma }\left(V^{i+mn}(V^{})^i\right)_{\delta ^m}^{}\text{ and }\psi _{,\sigma }\left(V^i(V^{})^{i+mn}\right)_{\delta ^m}^{}$$
defined by
$`\varphi _{,\sigma }(v_1\mathrm{}v_{i+mn}\phi _1\mathrm{}\phi _i)`$ $`={\displaystyle \underset{j=1}{\overset{m}{}}}<o^{},v_{j1}\mathrm{}v_{jn}>{\displaystyle \underset{h=1}{\overset{i}{}}}<v_{j_h},\phi _{\sigma _h}>`$
$`\psi _{,\sigma }(v_1\mathrm{}v_i\phi _1\mathrm{}\phi _{i+mn})`$ $`={\displaystyle \underset{j=1}{\overset{m}{}}}<o,\phi _{j1}\mathrm{}\phi _{jn}>{\displaystyle \underset{h=1}{\overset{i}{}}}<v_{\sigma _h},\phi _{j_h}>`$
where $`o`$ is a nonzero vector in $`^nV`$ and $`o^{}`$ is a non zero vector in $`^nV^{}`$.
###### Lemma 24.
1) If $`ij+mn`$ then $`\left(V^i(V^{})^j\right)_{\delta ^m}^{}=0.`$
2) If $`m0`$ then $`\left(V^{i+mn}(V^{})^i\right)_{\delta ^m}^{}`$ is generated by the functions $`\varphi _{,\sigma }.`$
3) If $`m0`$ then $`\left(V^i(V^{})^{i+mn}\right)_{\delta ^m}^{}`$ is generated by the functions $`\psi _{,\sigma }.`$
###### Proof.
The proof is clear by the previous lemma. ∎
### 2.2. A special case
In this section we proove a special case of Proposition 20 in which we are able to give a more precise result. To simplify the exposition of the proof of Proposition 20 we will also prove another lemma.
Here and in the following we will use polarization. If $`V`$ is finite dimensional vector space then we can define a map
$$\mathrm{}:(V^n)^{}S^n(V^{})[V]\text{ through }\mathrm{}(\phi )(v)=\phi (v\mathrm{}v).$$
###### Lemma 25.
$`\mathrm{}`$ is surjective, moreover if $`V`$ is a finite dimensional representation of a reductive group $`\mathrm{\Gamma }`$, and $`\chi `$ is a character of $`\mathrm{\Gamma }`$ then
$$\mathrm{}((V^n)_\chi ^{})=S^n(V^{})_\chi $$
where $`E_\chi `$ is the isotipic component of type $`\chi ^1`$ of a $`G`$ module $`E`$.
###### Lemma 26.
For $`i=1,\mathrm{},n`$ let $`\mathrm{\Gamma }_i`$ be a reductive group, $`\chi _i`$ be a character of $`\mathrm{\Gamma }_i`$ and $`E_i`$ be a f.d.representation of $`\mathrm{\Gamma }_i`$. Let $`\mathrm{\Gamma }=\mathrm{\Gamma }_i`$, then $`E=_iE_i`$ is a representation of $`\mathrm{\Gamma }`$ and $`\chi =\chi _i`$ is a character of $`\mathrm{\Gamma }`$. Then
$$E_\chi ^{}=(E_1)_{\chi _1}^{}\mathrm{}(E_n)_{\chi _n}^{}.$$
Let $`J^+`$, $`J^{}`$ be two sets of indeces and define $`\stackrel{~}{J}^+=\{0\}J^+`$, $`\stackrel{~}{J}^{}=\{0\}J^{}`$ and $`J=J^+\times J^{}`$, $`\stackrel{~}{J}=\stackrel{~}{J}^+\times \stackrel{~}{J}^{}\{(0,0)\}`$. For each $`i\stackrel{~}{J}^+`$ (resp. $`j\stackrel{~}{J}^+`$) choose a vector space $`Y_i`$ (resp. $`X_j`$) and define $`X=_{jJ^{}}X_j`$ and $`Y=_{iJ^+}Y_i`$. Consider the group
$$G_{XY}=\underset{iJ^+}{}GL(Y_i)\times \underset{jJ^{}}{}GL(X_j)$$
and its character $`c=_{iJ^+}det_{GL(Y_i)}\times \left(_{jJ^{}}det_{GL(X_j)}\right)^1`$.
We fix a matrix $`r=(r_{ij})_{i\stackrel{~}{J}^+,j\stackrel{~}{J}^{}}`$ of integers such that $`r_{i0}=1=r_{0j}`$ for all $`i,j`$ and $`r_{00}=1`$ and we consider the vector spaces:
$$H^{XY}=H=\underset{(i,j)\stackrel{~}{J}}{}\text{Hom}(X_j,Y_i)^{r_{ij}}\text{ and }H_0^{XY}=H_0=\underset{(i,j)J}{}\text{Hom}(X_j,Y_i)$$
where we adopt the convention $`E^n=^n=0`$ if $`n<0`$. When the spaces $`X,Y`$ will be clear from the context we will write $`H`$ and $`H_0`$ insted of $`H^{XY}`$ and $`H_0^{XY}`$. We fix a basis $`e_m^{ij}`$ of $`^{r_{ij}}`$ so we have a canonical identification
$$H=\underset{(i,j)\stackrel{~}{J}}{}\text{Hom}(X_j,Y_i)^{r_{ij}}.$$
(9)
We want to study $`c`$-equivariant polynomials on $`H`$. If we choose two finite dimensional vector spaces $`\stackrel{~}{A},\stackrel{~}{B}`$, linear maps $`\alpha :\stackrel{~}{A}X_0`$ , $`\beta :Y_0\stackrel{~}{B}`$, and elements $`\phi _{ij}(^{r_{ij}})^{}`$ for all $`i,j`$ then we can define a map $`\mathrm{\Phi }_{\phi ,\alpha ,\beta }:HH_0\text{Hom}(\stackrel{~}{A},Y)\text{Hom}(X,\stackrel{~}{B})\text{Hom}(X\stackrel{~}{A},Y\stackrel{~}{B})`$ by
$$\mathrm{\Phi }_{\phi ,\alpha ,\beta }(\underset{(i,j)\stackrel{~}{J}}{}A_{ij}v_{ij})=\underset{(i,j)J}{}\phi _{ij}(v_{ij})A^{ij}+\underset{iJ^+}{}A^{i0}\alpha +\underset{jJ^{}}{}\beta A^{0j}$$
(10)
where $`A_{ij}v_{ij}\text{Hom}(X_j,Y_i)^{r_{ij}}`$.
The following is a special version of 20.
###### Lemma 27.
$`[H]_c`$ is generated as a vector space by the following functions:
$$sdet\left(\mathrm{\Phi }_{\phi ,\alpha ,\beta }(s)\right)$$
where $`\mathrm{\Phi }_{\phi ,\alpha ,\beta }:HH_0`$ is as above.
#### 2.2.1. The special case
We will study an even more special case in which we are able to prove a better result that I find nice. In the above setting suppose that $`X_0=Y_0=0`$ and that $`r_{ij}=1`$ for all $`i,j`$.
Define the following set of matrices:
$`𝒮_n`$ $`=\{S=(s_{ij})^{J^+\times J^{}}:{\displaystyle \underset{i,j}{}}s_{ij}=n\}`$
$`𝒮^{XY}`$ $`=𝒮=\{S=(s_{ij})^{J^+\times J^{}}:{\displaystyle \underset{j}{}}s_{ij}=dimY_iiJ^+`$
$`\text{ and }{\displaystyle \underset{i}{}}s_{ij}=dimX_jjJ^{}\}`$
As for $`H`$ we will write $`𝒮`$ when the spaces $`X_j,Y_i`$ will be clear from the context. Observe that $`𝒮=\mathrm{}`$ if $`_jdimX_j_idimY_i`$ and that if $`N=dimX_j=dimY_i`$ then $`𝒮𝒮_N`$. For each $`card(J^+)\times card(J^{})`$ matirx $`S=(s_{ij})`$ we consider $`\phi _{ij}^{}`$ given by $`\phi _{ij}(\lambda )=s_{ij}\lambda `$ and we define
$$\mathrm{\Phi }_S=\mathrm{\Phi }_{\phi ,0,0}\text{ and }f_S=f_S^{XY}=det(\mathrm{\Phi }_S).$$
###### Proposition 28.
$`\{f_S\}_{S𝒮^{XY}}`$ is a basis of $`[H]_c`$.
###### Proof.
We have to compute $`S^n(H^{})_c=(S^n(H))_c^{}`$ for all $`n`$. For all $`S𝒮_n`$ define
$$E_S=\underset{(i,j)J^+\times J^{}}{}S^{s_{ij}}\left(\text{Hom}(X_j,Y_i)\right).$$
Observe that $`S^n(H)=_{S𝒮_n}E_S`$ as a $`G`$-module. So $`S^n(H)_c^{}=_{S𝒮_n}(E_S)_c^{}`$. Observe now that $`E_S`$ is a quotient of
$$\stackrel{~}{E}_S=\underset{(i,j)J^+\times J^{}}{}(X_j^{})^{s_{ij}}Y_i^{s_{ij}}.$$
(11)
By the lemmas in the previous section we have that
$$(\stackrel{~}{E}_S)_c^{}=\{\begin{array}{cc}0\hfill & \text{ if }S𝒮^{XY},\hfill \\ \hfill & \text{ if }S𝒮^{XY}.\hfill \end{array}$$
So in particular $`(E_S)_c^{}=0`$ if $`S𝒮^{XY}`$. Hence $`dimS^n(H)_c^{}card(𝒮^{XY})`$.
The function $`f_S`$ are clearly $`c`$-equivariant so the only thing that we have to prove is that they are linearly independent. To prove it we will prove a generalization of it.
If $`iJ^+`$ and $`jJ^{}`$ let $`E_{ij}`$ be the $`card(J^+)\times card(J^{})`$ matrix with a $`1`$ in the $`(i,j)`$ position and $`0`$ elsewhere.
For each $`iJ^+,jJ^{},m`$ and $`N`$ we consider the following sentence $`P_{i,j,m,N}`$:
> If $`_jX_j=N=_iY_i`$ then $`\{f_{S+mE_{ij}}\}_{S𝒮^{XY}}`$ is lineraly independent.
In the case $`m=0`$ we call this proposition $`P_{0,N}`$ since it does not depend on $`i,j`$ and observe that $`NP_{0,N}`$ is equivalent to our thesis.
For each $`N`$ we consider also the following sentence $`Q_N`$:
> If $`_jX_j=N=_iY_i`$ then $`P_{i,j,m,N}`$ is true for all $`iJ^+`$, $`jJ^{}`$ and $`m`$.
*First remark:* $`N=1`$ is true.
*Second remark:* let $`𝒮_0^{XY}=\{SS^{XY}:s_{ij}=0`$ for all $`iJ^+`$ and $`jJ^{}`$ such that $`dimY_i,dimX_j2`$. Observe that $`\{f_S\}_{S\stackrel{~}{𝒮}^{XY}}`$ is linearly independent.
Now we prove $`Q_N`$ by induction on $`N`$.
*Firts step:* $`Q_{N1}P_{0,N}`$. Suppose that there exists $`c_S`$ such that
$$\underset{S𝒮^{XY}}{}c_Sf_S=0.$$
If $`dimX_{j_0},dimY_{i_0}2`$ choose a nonzero element $`x_{j_0}X_{j_0}`$ (resp. $`y_{i_0}Y_{i_0}`$) and an hyperplane $`X_{j_0}^{}X_{j_0}`$ (resp. $`Y_{i_0}^{}Y_{i_0}`$) such that $`X_{j_0}=x_{j_0}X_{j_0}^{}`$ (resp. $`Y_{j_0}=y_{i_0}Y_{i_0}^{}`$) and define:
$$\stackrel{~}{X}_j=\{\begin{array}{cc}X_j\hfill & \text{ if }jj_0\hfill \\ X_{j_0}^{}\hfill & \text{ if }j=j_0\hfill \end{array}\text{ and }\stackrel{~}{Y}_i=\{\begin{array}{cc}Y_i\hfill & \text{ if }ii_0\hfill \\ Y_{i_0}^{}\hfill & \text{ if }i=i_0\hfill \end{array}$$
(12)
and define $`\mathrm{\Psi }:H^{\stackrel{~}{X}\stackrel{~}{Y}}H^{XY}`$ by
$$\mathrm{\Psi }(T)|_{\stackrel{~}{X}_j}=T\text{ and }\mathrm{\Psi }(T)(x_{j_0})=y_{j_0}.$$
(13)
Then we see that
$$0=\underset{S𝒮^{XY}}{}c_Sf_S^{XY}(\mathrm{\Psi }(T))=\underset{S𝒮^{XY}:s_{i_0j_0}0}{}s_{i_0j_0}c_Sf_S^{\stackrel{~}{X}\stackrel{~}{Y}}(T)$$
$$=\underset{S𝒮^{\stackrel{~}{X}\stackrel{~}{Y}}}{}(s_{i_0j_0}+1)c_{S+E_{i_0j_0}}f_{S+E_{i_0j_0}}^{\stackrel{~}{X}\stackrel{~}{Y}}(T)$$
By induction $`P_{i,j,1,N1}`$ is true for all $`i,j`$ so we see that $`c_S=0`$ for all $`S𝒮`$ such that there exists $`i_0,j_0`$ such that $`s_{i_0,j_0}1`$ and $`dimX_{j_0},dimY_{i_0}2`$. Now we conclude by the second remark.
*Second step:* $`Q_{N1}P_{i_0,j_0,m,N}`$ if $`dimX_{j_0},dimY_{i_0}2`$ and $`m1`$. Suppose that $`_{S𝒮^{XY}}c_Sf_{S+mE_{i_0j_0}}^{XY}=0`$. We can construct $`\stackrel{~}{X}_j,\stackrel{~}{Y}_i,\mathrm{\Psi }`$ as in the first step and we see that
$$0=\underset{S𝒮^{XY}}{}c_Sf_{S+mE_{i_0j_0}}^{XY}(\mathrm{\Psi }(T))=\underset{S𝒮^{XY}}{}(s_{i_0j_0}+m)c_Sf_{S+mE_{i_0j_0}}^{\stackrel{~}{X}\stackrel{~}{Y}}(T)$$
$$=\underset{S𝒮^{\stackrel{~}{X}\stackrel{~}{Y}}}{}(s_{i_0j_0}+m+1)c_{S+E_{i_0j_0}}f_{S+(m+1)E_{i_0j_0}}^{\stackrel{~}{X}\stackrel{~}{Y}}(T)$$
and by $`P_{i_0,j_0,m+1,N1}`$ we deduce $`c_S=0`$ for all $`S`$.
*Third step:* $`Q_{N1}Q_N`$. By the previous two step we have only to prove $`P_{i_0,j_0,m,N}`$ for $`m1`$ and $`dimX_{j_0}=1`$ or $`dimY_{i_0}=1`$. We will suppose $`dimX_{j_0}=1`$, the other case is completely similar. Suppose that $`_{S𝒮^{XY}}c_Sf_{S+mE_{i_0j_0}}=0`$. Set
$$\stackrel{~}{𝒮}_i=\{S𝒮^{XY}:s_{ij_0}=1\}$$
and observe that since $`dimX_{j_0}=1`$ then $`𝒮^{XY}=\stackrel{~}{𝒮}_i`$. Now choose a non zero vector $`x_{j_0}X_{j_0}`$ and for all $`iJ^+`$ choose a non zero vector $`y_iY_i`$ and an hyperplane $`Y_i^{}`$ of $`Y_i`$ such that $`Y_i=y_iY_i^{}`$.
Now fix $`i_1i_0`$ such that $`dimY_{i_1}2`$ and consider $`\stackrel{~}{J}^+=J^+`$ and $`\stackrel{~}{J}^{}=J^{}\{j_0\}`$. For all $`i\stackrel{~}{J}^+`$ and for all $`j\stackrel{~}{J}^{}`$ define:
$$\stackrel{~}{X}_j=X_j,\text{ and }\stackrel{~}{Y}_i=\{\begin{array}{cc}Y_i\hfill & \text{ if }ii_1,\hfill \\ Y_{i_1}^{}\hfill & \text{ if }i=i_1.\hfill \end{array}$$
For any $`S\stackrel{~}{𝒮}_{i_1}`$ we define also $`t(S)𝒮^{\stackrel{~}{X}\stackrel{~}{Y}}`$ by $`t(S)_{ij}=s_{ij}`$ for all $`i\stackrel{~}{J}^+`$, $`j\stackrel{~}{J}^{}`$. $`St(S)`$ is a bijection between $`\stackrel{~}{𝒮}_{i_1}`$ and $`𝒮^{\stackrel{~}{X}\stackrel{~}{Y}}`$: we call $`t^1`$ the inverse map. Finally we define $`\mathrm{\Psi }:H^{\stackrel{~}{X}\stackrel{~}{Y}}H^{XY}`$ as in the previou step and we observe that if $`S𝒮`$ then $`f_{S+mE_{i_0j_0}}\mathrm{\Psi }=0`$ if $`S\stackrel{~}{𝒮}_{i_1}`$. Hence
$$0=\underset{S𝒮^{XY}}{}c_Sf_{S+mE_{i_0j_0}}(\mathrm{\Psi }(T))=\underset{S\stackrel{~}{𝒮}_i}{}c_Sf_{t(S)}(T)$$
$$=\underset{S𝒮^{\stackrel{~}{X}\stackrel{~}{Y}}}{}c_{t^1(S)}f_S(T)$$
and applying $`P_{0,N1}`$ we obtain $`c_S=0`$ for all $`S\stackrel{~}{𝒮}_{i_1}`$ if $`dimY_{i_1}2`$ and $`i_1i_0`$.
In a similar way we prove $`c_S`$ if $`S\stackrel{~}{𝒮}_{i_1}`$ and $`dimY_{i_1}=1`$ and $`i_1i_0`$.
Finally we observe that if $`S\stackrel{~}{𝒮}_{i_0}`$ then $`f_{S+mE_{i_0j_0}}=(m+1)f_S`$, hence $`c_S=0`$ follows now from $`P_{0,N}`$ that we already know to be true. ∎
#### 2.2.2. Proof of Lemma 27
We study first $`(H^n)_c^{}`$ and then we apply polarization. As in the previous section we can decompose $`H^n`$ in summands of the following form:
$$E=\underset{(i,j)\stackrel{~}{J}}{}(\text{Hom}(X_j,Y_i)^{r_{ij}})^{s_{ij}}=\underset{(i,j)\stackrel{~}{J}}{}(X_j^{})^{s_{ij}}Y_i^{s_{ij}}(^{r_{ij}})^{s_{ij}}$$
(14)
where $`s_{ij}`$ are nonnegative integers such that $`_{i,j}s_{ij}=n`$. Observe that the order of the factors is not important for us since we will aply polarization.
We can describe easily $`E_c^{}`$ using the lemma in the previous section. In particular a necessary and sufficient condition for the existence of $`c`$-covariants is $`_{i\stackrel{~}{J}^+}s_{ij}=dimX_j`$ for all $`jJ^{}`$ and $`_{j\stackrel{~}{J}^+}s_{ij}=dimY_i`$ for all $`iJ^+`$. Moreover
$$E_c^{}\underset{(i,j)J}{}\left((^{r_{ij}})^{}\right)^{s_{ij}}\underset{jJ^{}}{}(Y_0^{})^{s_{0j}}\underset{iJ^+}{}X_0^{s_{i0}}$$
To write explicit formulas we choose an order on the factors of $`E`$ for example choosing a lexicographic order in $`i\stackrel{~}{J}^+`$, $`j\stackrel{~}{J}^{}`$ and $`1qs_{ij}`$:
$$E=\underset{q=1}{\underset{}{X_1^{}Y_1^{r_{11}}}}\mathrm{}\underset{q=s_{11}}{\underset{}{X_1^{}Y_1^{r_{11}}}}\underset{q=1}{\underset{}{X_1^{}Y_2^{r_{12}}}}\mathrm{}$$
Once we have chosen such an order we can write an element of $`E`$ as linear combination elements of the form $`_{(i,j,q)\stackrel{~}{K}}x^{i,j,q}y^{i,j,q}v^{i,j,q}`$ with $`x^{i,j,q}X_j^{}`$, $`y^{i,j,q}Y_i`$ $`v^{i,j,q}^{r_{ij}}`$ and we setted $`\stackrel{~}{K}=\{(i,j,q)\stackrel{~}{J}\times :\mathrm{\hspace{0.17em}1}qs_{ij}\}`$. We define also $`K=\{(i,j,q)\stackrel{~}{K}:(i,j)J\}`$. Using this convention if
$$\varphi =\underset{(i,j,q)\stackrel{~}{K}}{}\varphi ^{i,j,q}\underset{(i,j,q)K}{}(^{r_{ij}})^{}\underset{(0,j,q)\stackrel{~}{K}}{}Y_0^{}\underset{(i,0,q)\stackrel{~}{K}}{}X_0$$
(15)
the corresponding $`c`$ equivariant linear function on $`E`$ is defined on an element $`s=_{(i,j,q)\stackrel{~}{K}}x^{i,j,q}y^{i,j,q}v^{i,j,q}`$ by
$`\varphi (s)`$ $`={\displaystyle \underset{iJ^+}{}}<{\displaystyle \underset{\begin{array}{c}\\ (i,j,q)\stackrel{~}{K}\end{array}}{}}y^{i,j,q},o_i^{}>{\displaystyle \underset{jJ^{}}{}}<{\displaystyle \underset{\begin{array}{c}\\ (i,j,q)\stackrel{~}{K}\end{array}}{}}x^{i,j,q},o_j>`$
$`{\displaystyle \underset{(i,j,q)K}{}}\varphi ^{i,j,q}(v^{i,j,q}){\displaystyle \underset{(i,0,q)\stackrel{~}{K}}{}}\varphi ^{i,0,q}(x^{i,0,q}){\displaystyle \underset{(0,j,q)\stackrel{~}{K}}{}}\varphi ^{0,j,q}(y^{0,j,q})`$
Now consider the group
$$𝔖=𝔖_1\times 𝔖_2\times 𝔖_3=\underset{(i,j)J}{}S_{s_{ij}}\times \underset{jJ^{}}{}S_{s_{0j}}\times \underset{iJ^+}{}S_{s_{i0}}.$$
This group acts naturally on $`_{(i,j)J}\left((^{r_{i_1j_1}})^{}\right)^{s_{ij}}_{jJ^{}}(Y_0^{})^{s_{0j}}_{iJ^+}X_0^{s_{i0}}=E_c^{}`$ by permuting the factors and we observe that
$$\mathrm{}\left((\sigma _1,\sigma _2,\sigma _3)\varphi \right)=\epsilon (\sigma _2)\epsilon (\sigma _3)\mathrm{}(\varphi )$$
for all $`\varphi E_c^{}`$ and for all $`(\sigma _1,\sigma _2,\sigma _3)𝔖`$. So we have that
$$\mathrm{}(E_c^{})=\mathrm{}\left(\underset{(i,j)J}{}S^{s_{ij}}\left((^{r_{ij}})^{}\right)\underset{jJ}{}\stackrel{s_{0j}}{}Y_0^{}\underset{iJ^+}{}\stackrel{s_{i0}}{}X_0\right).$$
In particular since $`S^m(V)`$ is spanned by vectors of the form $`v\mathrm{}v`$, $`\mathrm{}(E_c^{})`$ is spanned by the functions $`\mathrm{}(\varphi )`$ with $`\varphi `$ of the following special form:
$$\varphi =\underset{(i,j)J}{}(\varphi ^{i,j})^{s_{ij}}\underset{jJ^{}}{}\varphi ^{0,j,1}\mathrm{}\varphi ^{0,j,s_{0j}}\underset{iJ^+}{}\varphi ^{i,0,1}\mathrm{}\varphi ^{i,0,s_{i0}}.$$
(16)
The lemma now follows from the following claim:
Claim: For each $`\varphi `$ as in (16) $`\mathrm{}(\varphi )`$ is a linear combination of the functions $`det(\varphi _{\phi ,\alpha ,\beta })`$.
We prove the claim as follows: we construct vector spaces $`A_i`$, $`B_j`$ and $`A=_{iJ^+}A_i`$, $`B=_{jJ^{}}B_j`$ and
$$\stackrel{~}{H}_0=H_0\underset{iJ^+}{}\text{Hom}(A_i,Y_i)\underset{jJ^{}}{}\text{Hom}(X_j,B_j)\text{Hom}(XA,YB)=\stackrel{~}{H}.$$
Observe that on $`\stackrel{~}{H},\stackrel{~}{H}_0`$ there is an action of $`\stackrel{~}{G}=G_{XY}\times G_{AB}=G_{XY}\times _{iJ^+}Gl(A_i)\times _{jJ^{}}Gl(B_j)`$ and we call $`\stackrel{~}{c}`$ the character of $`\stackrel{~}{G}`$ given by $`(_jdet_{GL(X_j)}\times _idet_{GL(A_i)})^1\times _idet_{GL(Y_i)}\times _jdet_{GL(B_j)}`$. We have an embedding of $`G_{XY}`$ in $`\stackrel{~}{G}`$ such that $`\sigma ^{}\stackrel{~}{c}=c`$. Observe also that by Proposition 28 we know that the $`\stackrel{~}{c}`$-covariants functions on $`H`$ are generated by the functionts $`det(\mathrm{\Phi }_{\stackrel{~}{S}})`$ with $`S\stackrel{~}{𝒮}`$: we put a tilde to emphasize that we have to consider also the components $`\{A_i\}`$ and $`\{B_j\}`$. Then we construct a $`G_{XY}`$-equivariant map $`\rho :H\stackrel{~}{H}_0`$ such that
1. there exists a $`\stackrel{~}{c}`$-covariant function $`f`$ on $`\stackrel{~}{H}`$ such that $`\mathrm{}(\varphi )=f\rho `$.
2. for all $`\stackrel{~}{S}\stackrel{~}{𝒮}`$ there exists $`\phi ,\alpha ,\beta `$ as in equation (10) such that $`det(\mathrm{\Phi }_{\stackrel{~}{S}})\rho =det(\mathrm{\Phi }_{\phi ,\alpha ,\beta })`$ .
The claim now follows by Proposition 28.
For $`iJ^+`$ and $`jJ^{}`$ define
$$A_i=^{s_{i0}},A=\underset{iJ^+}{}A_i,B_j=^{s_{0j}},B=\underset{jJ^{}}{}B_j.$$
Define also $`\alpha _i:A_iX_0`$ and $`(\beta _j)^t:B_j^{}Y_0^{}`$ by
$`\alpha _i(e_l)`$ $`=\varphi ^{i,0,l},`$ and $`\alpha `$ $`={\displaystyle \underset{iJ^+}{}}\alpha _i:AX_0`$
$`(\beta _i)^t(e^l)`$ $`=\varphi ^{0,j,l}`$ $`\text{and }\beta ^t`$ $`={\displaystyle \underset{jJ^{}}{}}B_j:B^{}Y_0,`$
where $`e_l`$ (resp. $`e^l`$) is the canonical basis of $`^m`$ (resp. $`(^m)^{}`$). We define $`\beta _i`$ (resp. $`\beta `$) as the transpose of $`(\beta _i)^t`$ (resp. $`\beta ^t`$). Now define $`\rho ^{ij}:\text{Hom}(X_j,Y_i)^{r_{ij}}\text{Hom}(X_j,Y_i)`$, $`\rho ^{i0}:\text{Hom}(X_0,Y_i)\text{Hom}(A_i,Y_i)`$, $`\rho ^{0j}:\text{Hom}(X_j,Y_0)\text{Hom}(X_j,B_j)`$ by
$$\rho ^{ij}(Tv)=\varphi ^{i,j}(v)T,\rho ^{i0}(T)=T\alpha _i,\rho ^{0j}(T)=\beta _jT,$$
and finally define $`\rho =_{i,j\stackrel{~}{J}}\rho ^{ij}:H\stackrel{~}{H}_0`$. Observe that $`\rho `$ is $`G_{XY}`$-equivariant.
Observe now that $`\stackrel{~}{H}_0^n=\stackrel{~}{E}_{\stackrel{~}{S}}`$ where $`\stackrel{~}{S}\stackrel{~}{𝒮}`$ and $`\stackrel{~}{E}_{\stackrel{~}{S}}`$ is defined as in (11). In particular we choose the following summund of $`\stackrel{~}{H}_0^n`$:
$$\stackrel{~}{E}=\underset{(i,j)J}{}\text{Hom}(X_j,Y_i)^{s_{ij}}\underset{jJ^{}}{}\text{Hom}(X_j,B_j)^{s_{0j}}\underset{iJ^+}{}\text{Hom}(A_i,Y_i)^{s_{i0}}$$
and we observe that $`(\stackrel{~}{E})_{\stackrel{~}{c}}^{}=`$. Choose a non zero element $`\stackrel{~}{\varphi }(\stackrel{~}{E})_{\stackrel{~}{c}}^{}`$ and observe that up to a scalar we have
$$\mathrm{}_{\stackrel{~}{H}}(\stackrel{~}{\varphi })\rho =\mathrm{}(\varphi ).$$
(17)
To see this choose $`\varphi `$ as in (16), and bases $`y_h^i`$ of $`Y_i`$, $`x_k^j`$ of $`X_j^{}`$ (and its dual basis $`z_k^j`$ of $`X_j`$) . Choose also a bases $`\epsilon _m^{ij}`$ of $`^{r_{ij}}`$ such that $`\varphi ^{i,j}(\epsilon _m^{ij})=\delta _{m,1}`$ and set $`A^{ij}=\rho ^{ij}(s)=_{h,k}a_{hk}^{ij}y_h^ix_k^j`$ for $`sH`$. Then
$`\mathrm{}(\varphi )(t)`$ $`={\displaystyle \underset{h𝒦_Y,k𝒦_X}{}}{\displaystyle \underset{iJ^+}{}}<{\displaystyle \underset{\begin{array}{c}\\ (i,j,q)\stackrel{~}{K}\end{array}}{}}a_{h(i,j,q)k(i,j,q)}^{ij}y_{h(i,j,q)}^i,o_i^{}>{\displaystyle \underset{jJ^{}}{}}<{\displaystyle \underset{\begin{array}{c}\\ (i,j,q)\stackrel{~}{K}\end{array}}{}}x_{k(i,j,q)}^j,o_j>`$
$`{\displaystyle \underset{iJ^+}{}}<{\displaystyle \underset{\begin{array}{c}\\ (i,0,q)\stackrel{~}{K}\end{array}}{}}x_{k(i,0,q)}^0,\varphi ^{i,0,1}\mathrm{}\varphi ^{i,0,s_{i0}}>`$
$`{\displaystyle \underset{jJ^{}}{}}<{\displaystyle \underset{\begin{array}{c}\\ (0,j,q)\stackrel{~}{K}\end{array}}{}}a_{h(0,j,q)k(0,j,q)}^{0j}y_{h(0,j,q)}^0,\varphi ^{0,j,1}\mathrm{}\varphi ^{0,j,s_{i0}}>`$
$`={\displaystyle \underset{k𝒦_X}{}}{\displaystyle \underset{iJ^+}{}}<{\displaystyle \underset{\begin{array}{c}\\ (i,j,q)\stackrel{~}{K}\end{array}}{}}A^{ij}z_{k(i,j,q)}^j,o_i^{}>{\displaystyle \underset{jJ^{}}{}}<{\displaystyle \underset{\begin{array}{c}\\ (i,j,q)\stackrel{~}{K}\end{array}}{}}x_{k(i,j,q)}^j,o_j>`$
$`{\displaystyle \underset{iJ^+}{}}<{\displaystyle \underset{\begin{array}{c}\\ (i,0,q)\stackrel{~}{K}\end{array}}{}}x_{k(i,0,q)}^0,\varphi ^{i,0,1}\mathrm{}\varphi ^{i,0,s_{i0}}>`$
$`{\displaystyle \underset{jJ^{}}{}}<{\displaystyle \underset{\begin{array}{c}\\ (0,j,q)\stackrel{~}{K}\end{array}}{}}A^{0j}z_{k(0,j,q)}^j,\varphi ^{0,j,1}\mathrm{}\varphi ^{0,j,s_{i0}}>`$
where the indeces are as follows:
$`𝒦_X`$ $`=\{k:\stackrel{~}{K}:\mathrm{\hspace{0.17em}1}k(i,j,q)dimX_j\}`$
$`𝒦_Y`$ $`=\{h:\stackrel{~}{K}:\mathrm{\hspace{0.17em}1}h(i,j,q)dimY_i\}.`$
The lefthandside in (17) clearly furnishes the same expression.
Finally if we fix $`\stackrel{~}{S}=(s_{MN})_{N\{X_j\}\{A_i\}\text{ and }M\{Y_i\}\{B_j\}}\stackrel{~}{𝒮}`$ and we choose $`\phi _{ij}=s_{Y_iX_j}\varphi ^{i,j}`$ and $`\alpha =_{iJ^+}s_{Y_iA_i}\alpha _i:AX_0`$ and $`\beta =_{jJ^{}}s_{B_jX_j}\beta _j:Y_0B`$ we have
$$det(\mathrm{\Phi }_{\stackrel{~}{S}})\rho =det(\mathrm{\Phi }_{\phi ,\alpha ,\beta }).$$
###### Remark 29.
The basis of $`[H]_c`$ we have described are different from the polarization of the natural basis of $`E_c^{}`$. The relation between the two basis is given by formulas of the following types
1. If $`A=\left(\begin{array}{cc}a_{11}& a_{12}\\ a_{21}& a_{22}\end{array}\right)`$ and $`B=\left(\begin{array}{cc}b_{11}& b_{12}\\ b_{21}& b_{22}\end{array}\right)`$ then
$$det\left(\begin{array}{cc}a_{11}& b_{12}\\ a_{21}& b_{22}\end{array}\right)+det\left(\begin{array}{cc}b_{11}& a_{12}\\ b_{21}& a_{22}\end{array}\right)=det(A+B)detAdetB.$$
2. If $`A=\left(\begin{array}{cc}a_{11}& a_{12}\\ a_{21}& a_{22}\end{array}\right)`$, $`B=\left(\begin{array}{cc}b_{11}& b_{12}\\ b_{21}& b_{22}\end{array}\right)`$, $`C=\left(\begin{array}{cc}c_{11}& c_{12}\\ c_{21}& c_{22}\end{array}\right)`$ and $`D=\left(\begin{array}{cc}d_{11}& d_{12}\\ d_{21}& d_{22}\end{array}\right)`$ then
$$det\left(\begin{array}{cc}a_{11}& b_{11}\\ a_{21}& b_{21}\end{array}\right)det\left(\begin{array}{cc}c_{12}& d_{12}\\ c_{22}& d_{22}\end{array}\right)det\left(\begin{array}{cc}a_{11}& b_{12}\\ a_{21}& b_{22}\end{array}\right)det\left(\begin{array}{cc}c_{12}& d_{11}\\ c_{22}& d_{21}\end{array}\right)+$$
$$det\left(\begin{array}{cc}a_{12}& b_{11}\\ a_{22}& b_{21}\end{array}\right)det\left(\begin{array}{cc}c_{11}& d_{12}\\ c_{21}& d_{22}\end{array}\right)+det\left(\begin{array}{cc}a_{12}& b_{12}\\ a_{22}& b_{22}\end{array}\right)det\left(\begin{array}{cc}c_{11}& d_{11}\\ c_{21}& d_{21}\end{array}\right)=$$
$$=det\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)+det\left(\begin{array}{cc}A& 0\\ 0& D\end{array}\right)+det\left(\begin{array}{cc}0& B\\ C& 0\end{array}\right)$$
The first type of formula correspond to the reduction of Lemma 27 to the case $`r_{ij}=1`$ and $`X_0=Y_0=1`$. The second type of formula correspond to the case of Proposition 28.
### 2.3. Proof of Proposition 20
Choose basis $`_i`$ (resp. $`_i^{}`$) of the vector spaces $`D_i`$ and $`D_i^{}`$ and we write our vector space $`S(d,v)`$ in the following way:
$$S=\underset{hH}{}V_{h_0}^{}V_{h_1}\underset{\begin{array}{c}iI\\ b_i\end{array}}{}V_{i,b}\underset{iI,b^{}_i^{}}{}V_{i,b^{}}^{}$$
where $`V_{i,b}`$ (resp. $`V_{i,b^{}}^{}`$) is an isomorphic copy of $`V_i`$ (resp. $`V_i^{}`$). We fix also a character $`\chi _m`$ and $`m_i`$, $`\stackrel{~}{m}_i,m_i^+,m_i^{},I^+,I^{},I^0`$ as in 2.0.2 and we describe first the $`\chi _m`$-covariants of $`S^n`$. To do it we observe that we can decompose $`S^n`$ in the following way:
$$S^n=\underset{\mathrm{}}{}E_1^{(\mathrm{})}\mathrm{}E_n^{(\mathrm{})}$$
where each $`E_i^{(\mathrm{})}`$ is a representation of $`G`$ of one of the following types: $`V_{h_0}^{}V_{h_1}`$, $`V_{i,b}`$ or $`V_{i,b^{}}^{}`$. So it is enough to compute the $`\chi `$-covariants of each piece $`E_1^{(\mathrm{})}\mathrm{}E_n^{(\mathrm{})}`$. We fix one of them: $`E=E_1\mathrm{}E_n`$ and we compute $`E_\chi ^{}`$. Let $`I^{}`$ be a copy of $`I`$ and fix an isomorphism $`ii^{}`$ between the two sets. For each $`j=1,\mathrm{},n`$ we define a subset $`𝒮_j`$ of $`II^{}`$ according to the following rule:
$$𝒮_j=\{\begin{array}{cc}\{h_0^{},h_1\}\hfill & \text{ if }E_j=V_{h_0}^{}V_{h_1},\hfill \\ \{i\}\hfill & \text{ if }E_j=V_{i,b}\text{ for some }b_i,\hfill \\ \{i^{}\}\hfill & \text{ if }E_j=V_{i,b^{}}^{}\text{ for some }b^{}_i^{}.\hfill \end{array}$$
Let now be $`𝒮=_{j=1}^n𝒮_j`$. An element of $`𝒮`$ can be thought as a couple $`(i,j)`$ (or $`(i^{},j)`$) where $`i`$ (or $`i^{}`$) is in $`𝒮_j`$. We consider now a special class of partitions of $`𝒮`$: a collection $`𝔉=\{𝒞,_i^{(l)}\text{ for }iI\text{ and }1lm_i\}`$ of disjoint subsets of $`2^𝒮`$ is called $`m`$-*special* if:
1. $`𝔉`$ is a partition of $`𝒮`$,
2. $`C𝒞`$ $`cardC=2`$ and $`iI`$, $`𝒮_{j_1},𝒮_{j_2}`$ such that $`i𝒮_{j_1}`$, $`i^{}𝒮_{j_2}`$ and $`C=\{(i,j_1),(i^{},j_2)\}`$,
3. $`M_i^{(l)}`$ we have $`M=\{(i,j)\}`$ if $`iI^+`$ and $`M=\{(i^{},j)\}`$ if $`iI^{}`$,
4. $`card_i^{(l)}=v_i=dimV_i`$.
We can represents a special collection with an enriched graph whose vertices are the sets $`𝒮_j`$ and completed according with the following rules:
1. we put an arrow from $`𝒮_{j_1}`$ to $`𝒮_{j_2}`$ if there exists $`C=\{(i,j_2),(i^{},j_1)\}𝒞`$,
2. we put an indexed circle box $`_i^l`$ on $`𝒮_j`$ if there exists $`M=\{(i,j)\}_i^{(l)}`$
3. we put an indexed square box $`\mathrm{}_i^l`$ on $`𝒮_j`$ if there exists $`M=\{(i^{},j)\}_i^{(l)}`$
4. if $`E_j`$ is of type $`V_{i,b}`$ or $`V_{i,b^{}}^{}`$ then we add the element $`b`$ or $`b^{}`$ at the left of the corresponding vertex.
5. if $`E_j`$ is of type $`V_{h_0}^{}V_{h_1}`$ then we write $`h`$ at the left of the corresponding vertex.
Observe that a vertex can be marked with a circle and a square but that it cannot be marked with two circles or two squares.
There is a perfect bijection between $`m`$-special collection $`𝔉`$ and graphs as above such that:
1. the cardinality of vertexes marked with $`_i^l`$ is $`v_i`$ for each $`iI^+`$ and $`1lm_i`$,
2. the cardinality of vertexes marked with $`\mathrm{}_i^l`$ is $`v_i`$ for each $`iI^{}`$ and $`1lm_i`$.
We will use the same letter $`𝔉`$ to indicate the collection or the graph.
To a special collection $`𝔉`$ as above we attach a function $`\varphi _𝔉`$ on $`E`$. We define it by the formula
$$\varphi _𝔉(e_1\mathrm{}e_n)=\underset{C𝒞}{}\varphi _C\underset{iI^+}{}\underset{l=1}{\overset{m_i}{}}<o_i^{},_i^{(l)}>\underset{iI^{}}{}\underset{l=1}{\overset{m_i}{}}<o_i,_i^{(l)}>$$
where $`o_i`$ is a non zero element in $`^{v_i}V_i`$, $`o_i^{}`$ is a non zero element in $`^{v_i}V_i^{}`$ and
$`e_j`$ $`=\{\begin{array}{cc}x_j^{}V_i^{}\hfill & \text{ if }E_j=V_{i,b}^{},\hfill \\ y_jV_i\hfill & \text{ if }E_j=V_{i,b},\hfill \\ x_j^{}y_jV_{h_0}^{}V_{h_1}\hfill & \text{ if }E_j=V_{h_0}^{}V_{h_1},\hfill \end{array}`$
$`\varphi _C`$ $`=<x_{j_1}^{},v_{j_2}>\text{ if }C=\{(i^{},j_1),(i,j_2)\}`$
$`{\displaystyle _i^{(l)}}`$ $`=y_{j_1}\mathrm{}y_{j_{v_i}}\text{ if }_i^{(l)}=\{\{(i,j_1)\},\mathrm{},\{(i,j_{v_i})\}\}\text{ and }iI^+`$
$`{\displaystyle _i^{(l)}}`$ $`=x_{j_1}^{}\mathrm{}x_{j_{v_i}}^{}\text{ if }_i^{(l)}=\{\{(i,j_1)\},\mathrm{},\{(i,j_{v_i})\}\}\text{ and }iI^{}`$
Finally we extend $`\varphi _𝔉`$ to all $`E`$ by linearity. By the lemma above and the discussion in 2.1 we deduce easily the following lemma:
###### Lemma 30.
$`E_\chi ^{}`$ is generated by the functions $`\varphi _𝔉`$.
Proposition 20 now follows from lemma 30 and the following claim: claim: for any special collection $`𝔉`$ the function $`\mathrm{}(\varphi _𝔉)`$ is a $`[S]^{G_v}`$-linear combination of the functions $`f_\mathrm{\Delta }`$ described in 2.0.2.
We consider the connected components of the graph. There are only five possible types of paths:
1. closed paths,
2. straight paths leaving from a non boxed vertex and arriving in a non boxed vertex,
3. straight paths leaving from a non boxed vertex and arriving in a circle boxed vertex,
4. straight paths leaving from a square boxed vertex and arriving in a non boxed vertex,
5. straight paths leaving from a square boxed vertex and arriving in a circle boxed vertex.
Let now $`𝔉_0`$ be the union of the connected components of the first two types and $`𝔉_1`$ be the union of the remaing components. Observe that
$$\mathrm{}(\varphi _𝔉)=\mathrm{}(\varphi _{𝔉_0})\mathrm{}(\varphi _{𝔉_1}).$$
Observe also that $`\varphi _{𝔉_0}`$ is an invariant function (indeed this part of the graph corresponds to the situation studied by Lusztig in \[Lu:Q4\]). Since we are interested in generators of $`[S]_{\chi ,1}`$as a $`[S]^G`$-module, we can suppose for simplicity $`𝔉=𝔉_1`$.
Observe now that each connected component $`\mathrm{\Gamma }`$ of the graph of the third type and with a circle $`_{i_1}^l`$ at the end, has an initial vertex which is an $`𝒮_j=\{i_0^{}\}`$ and that is marked with $`b_{i_0}^{}`$ on the left. All the other vertexes of the connected component are of type $`𝒮_j=\{h_0^{},h_1\}`$ and they define a path $`\alpha ^\mathrm{\Gamma }`$ such that $`\alpha _0^\mathrm{\Gamma }=i_0`$ and $`\alpha _1^\mathrm{\Gamma }=i_1`$. We call $`b=b(\mathrm{\Gamma })`$ and $`l=L_1(\mathrm{\Gamma })`$.
In the same way we see that:
1. each connected component $`\mathrm{\Gamma }`$ of the fourth type determines a path $`\alpha ^\mathrm{\Gamma }`$, $`b^{}=b^{}(\mathrm{\Gamma })_{\alpha _1^\mathrm{\Gamma }}^{}`$ and $`l=L_0(\mathrm{\Gamma })`$ such that $`1lm_{\alpha _0^\mathrm{\Gamma }}`$,
2. each connected component $`\mathrm{\Gamma }`$ of the fifth type determines a path $`\alpha ^\mathrm{\Gamma }`$, $`l_0=L_0(\mathrm{\Gamma })`$ and $`l_1=L_0(\mathrm{\Gamma })`$ such that $`1l_0m_{\alpha _0^\mathrm{\Gamma }}`$ and $`1l_1m_{\alpha _1^\mathrm{\Gamma }}`$.
Now we prove the claim in the following way, we construct $`X_j`$ and $`Y_i`$ as in 2.2, a groups homomorphism $`\sigma :G_vG_{XY}`$ such that $`\sigma ^{}c=\chi _m`$, a $`G_v`$ equivariant map $`\rho :SH`$, and a $`G_{XY}`$ $`c`$-covariant function $`f`$ on $`H`$ such that:
1. for all $`\phi ,\alpha ,\beta `$ there exists a $`\chi _m`$-good data such that $`det(\mathrm{\Phi }_{\varphi ,\alpha ,\beta })\rho =f_\mathrm{\Delta }`$,
2. $`\mathrm{}(\varphi _𝔉)=f\rho `$
The claim will clearly follow.
Set
$`J^{}`$ $`=\{(i,l):iI^{}\text{ and }1lm_i\},`$
$`J^+`$ $`=\{(i,l):iI^+\text{ and }1lm_i\}.`$
For all $`(i,l)J^{}`$ choose $`X_{(i,l)}=V_i`$ and for each $`(i,l)J^+`$ choose $`Y_{(i,l)}=V_i`$. For each $`(i_0,l_0)J^{}`$ and for each $`(i_1,l_1)J^+`$ define:
$`r_{(i_0,l_0)(i_1,l_1)}`$ $`=card\{\text{connected component }\mathrm{\Gamma }\text{ of the fifth type such}`$
$`\text{that }\alpha _0^\mathrm{\Gamma }=i_0,\alpha _1^\mathrm{\Gamma }=i_1,L_0(\mathrm{\Gamma })=l_0\text{ and }L_1(\mathrm{\Gamma })=l_1\}`$
We the connected component $`\mathrm{\Gamma }`$ of the set in the left handside as a basis $`e_\mathrm{\Gamma }`$ of the vector space $`^{r_{(i_0,l_0)(i_1,l_1)}}`$. This basis plays the role of the basis $`e_m^{ij}`$ we used to give the identification in (9).
For each $`\mathrm{\Gamma }`$ of the third type choose a one dimensional vector space $`_{b(\mathrm{\Gamma })}`$ and fix a generator $`b_\mathrm{\Gamma }`$. For each $`\mathrm{\Gamma }`$ of the fourth type choose a one dimensional vector space $`_{b^{}(\mathrm{\Gamma })}`$ and fix a generator $`b_\mathrm{\Gamma }^{}`$.
$`X_0`$ $`={\displaystyle \underset{\mathrm{\Gamma }\text{ of the third type}}{}}_{b(\mathrm{\Gamma })}={\displaystyle \underset{\mathrm{\Gamma }\text{ of the third type}}{}}b_\mathrm{\Gamma }`$
$`Y_0`$ $`={\displaystyle \underset{\mathrm{\Gamma }\text{ of the fourth type}}{}}_{b^{}(\mathrm{\Gamma })}={\displaystyle \underset{\mathrm{\Gamma }\text{ of the fourth type}}{}}b_\mathrm{\Gamma }^{}.`$
Now for each connected component $`\mathrm{\Gamma }`$ of the third type define $`\rho ^\mathrm{\Gamma }:S\text{Hom}(_{b(\mathrm{\Gamma })},Y_{(\alpha _1^\mathrm{\Gamma },L_1(\mathrm{\Gamma }))})`$ by
$$s\{\lambda \alpha ^\mathrm{\Gamma }(s)\gamma _{\alpha _0^\mathrm{\Gamma }}(b(\mathrm{\Gamma }))\lambda \}$$
In a similar way define $`\rho ^\mathrm{\Gamma }`$ if $`\mathrm{\Gamma }`$ is the fourth or of the fifth type. Finally define
$$\rho :SH\text{ by }\rho =\underset{\mathrm{\Gamma }}{}\rho ^\mathrm{\Gamma }.$$
Define also a group homomorphism $`\sigma :G_vG_{XY}`$ by $`(\sigma (g_i))_{X_{(i_0,l_0)}}=g_{i_0}`$ and $`(\sigma (g_i))_{Y_{(i_1,l_1)}}=g_{i_1}`$, and observe that $`\rho `$ is $`G_v`$ equivariant.
Now we describe $`\varphi (H^{\stackrel{~}{n}})_c^{}`$ (in general $`\stackrel{~}{n}`$ is less or equal to $`n`$) such that
$$\mathrm{}(\varphi _𝔉)(s)=\mathrm{}(\varphi )(\rho (s)).$$
(18)
We describe $`\varphi `$ by giving a summunds $`\stackrel{~}{E}`$ of $`H^{\stackrel{~}{n}}`$ as in (14) and $`\varphi \stackrel{~}{E}_c^{}`$ as in (15).To define $`\stackrel{~}{E}`$ we have to define $`s_{(i_1,l_1)(i_0,l_0)}`$, $`s_{(i_1,l_1)0}`$ and $`s_{0(i_0,l_0)}`$ for all $`(i_1,l_1)J^+`$ and for all $`(i_0,l_0)J^{}`$. We set
$`s_{(i_1,l_1)(i_0,l_0)}`$ $`=r_{(i_1,l_1)(i_0,l_0)}`$
$`s_{(i_1,l_1)0}`$ $`=card\{\text{ connected component }\mathrm{\Gamma }\text{ of the third type}`$
$`\text{ such that }\alpha _1^\mathrm{\Gamma }=i_1\text{ and }L_1(\mathrm{\Gamma })=l_1\}`$
$`s_{0(i_0,l_0)}`$ $`=card\{\text{ connected component }\mathrm{\Gamma }\text{ of the fourth type}`$
$`\text{ such that }\alpha _0^\mathrm{\Gamma }=i_0\text{ and }L_0(\mathrm{\Gamma })=l_0\}`$
Observe that we can choose a bijection $`q\mathrm{\Gamma }_q`$ between $`\{1,\mathrm{},s_{(i_1,l_1)(i_0,l_0)}\}`$ and the set of connected component $`\mathrm{\Gamma }`$ of the fifth type such that $`\alpha _1^\mathrm{\Gamma }=i_1`$,$`L_1(\mathrm{\Gamma })=l_1`$, $`\alpha _0^\mathrm{\Gamma }=i_0`$ and $`L_0(\mathrm{\Gamma })=l_0`$. So we can define $`\varphi ^{(i_1,l_1),(i_0,l_0),q}`$ by
$$\varphi ^{(i_1,l_1),(i_0,l_0),q}(e_\mathrm{\Gamma })=\delta _{\mathrm{\Gamma },\mathrm{\Gamma }_q}.$$
Observe also that we can choose a bijection $`q\mathrm{\Gamma }_q`$ between $`\{1,\mathrm{},s_{(i_1,l_1)0}\}`$ and the set of connected component $`\mathrm{\Gamma }`$ of the third type such that $`\alpha _1^\mathrm{\Gamma }=i_1`$,$`L_1(\mathrm{\Gamma })=l_1`$. So we can define $`\varphi ^{(i_1,l_1),0,q}`$ by
$$\varphi ^{(i_1,l_1),0,q}(b_\mathrm{\Gamma })=\delta _{\mathrm{\Gamma },\mathrm{\Gamma }_q}.$$
In a similar way define $`\varphi ^{0,(i_0,l_0),q}`$.
Up to a sign which depends on our choices and ordering equation (18) is a tautologically satisfied.
Observe now that by lemma 27 and linearity $`[H]_c`$ is generated by functions $`sdet(\mathrm{\Phi }_{\phi ,\alpha ,\beta }(s))`$ where $`\stackrel{~}{A},\stackrel{~}{B},\phi ,\alpha ,\beta `$ are as in (10) and moreover there exists a basis $`e_1,\mathrm{},e_{r_A}`$ of $`\stackrel{~}{A}`$ and a basis $`\stackrel{~}{e}_1,\mathrm{},\stackrel{~}{e}_{r_B}`$ of $`\stackrel{~}{B}^{}`$ such that for all $`i`$ there exist a connected component of the third type $`\mathrm{\Gamma }_i^A`$ such that $`\alpha (e_i)=b_{\mathrm{\Gamma }_i^A}`$ and for all $`i`$ there exists a connected component of the fourth type $`\mathrm{\Gamma }_i^B`$ such that $`\stackrel{~}{e}_i(\beta (b_\mathrm{\Gamma }^{}))=\delta _{\mathrm{\Gamma }\mathrm{\Gamma }_i^B}`$. So it is enough to prove that if $`\stackrel{~}{A},\stackrel{~}{B},\phi ,\alpha ,\beta `$ are as above then there exists a $`\chi _m`$-good $`\mathrm{\Delta }`$ such that
$$det(\mathrm{\Phi }_{\phi ,\alpha ,\beta })\rho =f_\mathrm{\Delta }.$$
We define
$`A`$ $`=(b_{\mathrm{\Gamma }_i^A})_{i=1,\mathrm{},r_A}\left({\displaystyle D_i}\right)^{dim\stackrel{~}{A}}`$
$`B`$ $`=(b_{\mathrm{\Gamma }_i^B}^{})_{i=1,\mathrm{},r_B}\left({\displaystyle D_i^{}}\right)^{dim\stackrel{~}{B}}`$
$`\alpha _{jh}^{ik}`$ $`={\displaystyle \underset{1qs_{(j,h)(i,k)}}{}}\varphi ^{(j,h),(i,k),q}(e_\mathrm{\Gamma })\alpha ^\mathrm{\Gamma }`$
$`\alpha _l^{ik}`$ $`=\delta _{i,(\alpha ^{\mathrm{\Gamma }_l^B})_0}\delta _{k,L_0(\mathrm{\Gamma }_l^B)}\alpha ^{\mathrm{\Gamma }_l^B}`$
$`\alpha _{jh}^l`$ $`=\delta _{j,(\alpha ^{\mathrm{\Gamma }_l^A})_1}\delta _{h,L_1(\mathrm{\Gamma }_l^A)}\alpha ^{\mathrm{\Gamma }_l^A}`$
The equation (18) follows now by the very definition.
## 3. The action of the Weyl group
For any $`mP`$ and for any $`\lambda Z`$ we defined a variety $`M_{m,\lambda }(d,v)`$. Observe that on both $`m,\lambda `$ there is a natural action of the Weyl group $`W`$. We define an action of the Weyl group also on $`(d,v)`$. We have already described $`d`$ as an element of $`X`$ and $`v`$ as an element of $`Q`$. We can now define
$$\sigma (d,v)=(d,\sigma (vd)+d).$$
Observe that $`\sigma (vd)+dQ`$ so the definition is well given. So it make sense to consider the variety $`M_{\sigma m,\sigma \lambda }(\sigma (d,v))`$ or the variety $`𝔐_{\sigma \zeta }(\sigma (d,v))`$ for $`\zeta `$.
In \[Na1\] Nakajima used analytic methods to prove, in the case of a finite Dynkin diagram, that if $`\zeta `$ is generic then there exists a diffeomorphism of differentiable manifolds
$$\mathrm{\Phi }_{\sigma ,\zeta }:𝔐_\zeta (d,v)𝔐_{\sigma \zeta }(\sigma (d,v))$$
and moreover that $`\mathrm{\Phi }_{\sigma ^{},\sigma \zeta }\mathrm{\Phi }_{\sigma ,\zeta }=\mathrm{\Phi }_{\sigma ^{}\sigma ,\zeta }`$. In the same paper he also asserted that a similar construction could be obtained in the general case using reflection functors as indeed we are going to do.
In \[Lu:Q4\] Lusztig gave a purely algebraic construction of an isomorphism
$$M_{0,\lambda }(d,v)M_{0,s_i\lambda }(s_i(d,v))$$
whenever $`\lambda _i0`$. In this paper we will give a generalization of Lusztig construction.
###### Definition 31.
If $`u^n=Q^{}`$ and $`AQ^{}`$ we define
$$H_u=\{(m,\lambda )PZ:<u,^{},\lambda >=<u^{},m>=0\}\text{ and }H_A=\underset{aA}{}H_a$$
Let $`K=\mathrm{max}\{1,a_{ij}^2:i,jI\}`$. If $`v^n`$ we define
$$\stackrel{~}{U}_v=\{u^I:\mathrm{\hspace{0.17em}0}u_iKv_i\}\text{ and }\stackrel{~}{H}^v=H_{\stackrel{~}{U}_v}.$$
We define also
$$U_{\mathrm{}}=\underset{iI}{}W\alpha _i^{}\text{ and }H^{\mathrm{}}=H_U_{\mathrm{}}.$$
Finally we set $`𝒢_v=\{(m,\lambda )P\times Z_G:\sigma (m,\lambda )H^{\sigma v}\text{ for all }\sigma W\}.`$ Both of the following definition of the set $`𝒢`$ will be fine for us:
$`𝒢`$ $`=\{(v,m,\lambda )Q\times P\times Z_G:(m,\lambda )𝒢_v\}\text{ or}`$
$`𝒢`$ $`=\{(v,m,\lambda )Q\times P\times Z_G:(m,\lambda )H^{\mathrm{}}\}.`$
We observe that in any case $`𝒢`$ is $`W`$-stable.
###### Proposition 32.
For all $`d,v`$, for all $`\sigma W`$ and for all $`(m,\lambda )`$ such that $`(m,\lambda ,v)𝒢_v`$ there exists an algebraic isomorphism:
$$\mathrm{\Phi }_{\sigma ,m,\lambda }^{d,v}:M_{m,\lambda }(d,v)M_{\sigma m,\mathrm{\sigma }\mathrm{\lambda }}\left(\sigma (d,v)\right).$$
Moreover this isomorphims satisfies
$$\mathrm{\Phi }_{\tau ,\sigma m,\mathrm{\sigma }\mathrm{\lambda }}^{\sigma (d,v)}\mathrm{\Phi }_{\sigma ,m,\lambda }^{d,v}=\mathrm{\Phi }_{\tau \sigma ,m,\lambda }^{d,v}.$$
(19)
### 3.1. Generators
In this section we define the actions of the generators $`s_i`$ of $`W`$ following \[Lu:Q4\]. We fix $`iI`$ and $`(d,v)`$, $`\lambda Z`$ and $`mP`$. We call $`(d,v^{})=s_i(d,v)`$, $`\lambda ^{}=s_i\lambda `$ and $`m^{}=s_im`$. Through all this section we assume $`v,v^{}0`$. For the convenience of the reader we write explicit formula in this case:
$`\lambda _j^{}`$ $`=\lambda _jc_{ij}\lambda _i`$ $`m_j^{}`$ $`=m_jc_{ij}m_i\text{ for all }j`$
$`v_i^{}`$ $`=d_iv_i+{\displaystyle \underset{ji}{}}a_{ij}v_j`$ $`v_j^{}=v_j\text{ for all }ji`$
Observe that we can choose
$$D_j^{}=D_j\text{ for all }j\text{ and }V_j^{}=V_j\text{ for all }ji.$$
In particular we have
$$T_i=D_i\underset{h_1=i}{}V_{h_0}=T_i^{}$$
since we suppose that our quiver has not simple loops.
###### Definition 33 (Lusztig \[Lu:Q4\]).
Fix $`\lambda Z_G`$ and define $`Z_i^\lambda (d,v)`$ to be the subvariety of $`S_i(d,v)\times S_i(d,v^{})`$ of pairs $`(s,s^{})=((B,\gamma ,\delta ),(B^{},\gamma ^{},\delta ^{}))`$ such that the following conditions hold:
1. $`B_h(s)=B_h^{}(s^{})`$ for all $`h`$ such that $`h_0,h_1i`$,
2. $`\gamma _j(s)=\gamma _j^{}(s^{})`$ for all $`ji`$ ,
3. $`\delta _j(s)=\delta _j^{}(s^{})`$ for all $`ji`$,
4. the following sequence is exact:
$$\begin{array}{ccccccccc}0& & V_i^{}& \stackrel{a_i^{}}{}& T_i& \stackrel{b_i}{}& V_i& & 0,\end{array}$$
(20)
5. $`a_i^{}(s^{})b_i^{}(s^{})=a_i(s)b_i(s)\lambda _i\mathrm{Id}_{T_i}`$,
6. $`s\mathrm{\Lambda }_\lambda (d,v)`$ and $`s^{}\mathrm{\Lambda }_\lambda ^{}(d,v^{})`$.
###### Lemma 34.
Let $`(s,s^{})S_i(d,v)\times S_i(d,v^{})`$ and suppose that it satisfies conditions 1), 2), 3), 4) , 5) above then:
1. $`s\mathrm{\Lambda }_\lambda (d,v)s^{}\mathrm{\Lambda }_\lambda ^{}(d,v^{})`$,
2. if $`\mu _j(s)\lambda _j\mathrm{Id}_{V_j}=0`$ for all $`ji`$ then $`s\mathrm{\Lambda }_\lambda (d,v)`$ ,
3. if $`\mu _j(s^{})=\lambda _j\mathrm{Id}_{V_j^{}}`$ for all $`ji`$ then $`s^{}\mathrm{\Lambda }_\lambda ^{}(d,v^{})`$.
###### Proof.
2) We have to prove $`b_ia_i\lambda _i\mathrm{Id}_{V_i}=0`$ and by condition 4) it is enough to prove $`b_ia_ib_i=\lambda _ib_i`$. So $`b_ia_ib_i=b_i(a_i^{}b_i^{}\lambda _i)=\lambda _ib_i`$ by conditions 4) and 5).
The proof of 3) is equal to the proof of 2). We prove the implication $``$ in 1). By 2) and 3) it is enough to prove that $`b_j^{}a_j^{}=\lambda _j^{}`$ for $`ji`$.
$`b_j^{}a_j^{}`$ $`={\displaystyle \underset{h_1=j}{}}\epsilon (h)B_h^{}B_{\overline{h}}^{}+\gamma _j^{}\delta _j^{}=`$
$`={\displaystyle \underset{h_1=j,h_0i}{}}\epsilon (h)B_hB_{\overline{h}}+\gamma _j\delta _j+{\displaystyle \underset{h_1=j,h_0=i}{}}\epsilon (h)B_h^{}B_{\overline{h}}^{}=`$
$`=b_ja_j+{\displaystyle \underset{h_1=j,h_0=i}{}}\epsilon (h)\left(B_h^{}B_{\overline{h}}^{}B_hB_{\overline{h}}\right)`$
$`=b_ja_j+{\displaystyle \underset{h_0=j,h_1=i}{}}\left(B_{\overline{h}}\epsilon (h)B_hB_{\overline{h}}^{}\epsilon (h)B_h^{}\right)`$
$`=b_ja_j+{\displaystyle \underset{h_0=j,h_1=i}{}}\left([a_ib_i]_{V_{h_0}}^{V_{h_0}}[a_i^{}b_i^{}]_{V_{h_0}}^{V_{h_0}}\right)`$
$`=\lambda _j+{\displaystyle \underset{h_0=j,h_1=i}{}}\lambda _i=\lambda _j^{}`$
The proof of the converse is completely analougous. ∎
###### Lemma 35.
Let $`\lambda Z_G`$, $`(s,s^{})Z_i^\lambda (d,v)`$ and $`\alpha `$ be an element of the path algebra algebra of type $`(\alpha _0,\alpha _1)`$ then
1. if $`\alpha _0,\alpha _1i`$ there exists an element $`\alpha ^{}`$ of the $`b`$-path algebra of type $`(\alpha _0,\alpha _1)`$ such that $`\alpha ^{}(s^{})=\alpha (s)`$,
2. if $`\alpha _1i`$ there exists an element $`\alpha ^{}`$ of the $`b`$-path algebra of type $`(\alpha _0,\alpha _1)`$ such that $`\alpha ^{}(s^{})\gamma _{\alpha _0}^{}=\alpha (s)\gamma _{\alpha _0}`$,
3. if $`\alpha _0i`$ there exists an element $`\alpha ^{}`$ of the $`b`$-path algebra of type $`(\alpha _0,\alpha _1)`$ such that $`\delta _{\alpha _1}^{}\alpha ^{}(s^{})=\delta _{\alpha _1}\alpha (s)`$,
4. there exists an element $`\alpha ^{}`$ of the $`b`$-path algebra of type $`(\alpha _0,\alpha _1)`$ such that $`\delta _{\alpha _1}^{}\alpha ^{}(s^{})\gamma _{\alpha _0}^{}=\delta _{\alpha _1}\alpha (s)\gamma _{\alpha _0}`$,
###### Proof.
By induction on the length of $`\alpha `$ we can reduce the proof of this lemma to the following identities that are a consequence of condition 5) in definition 33:
$`B_h^{}B_k^{}`$ $`=\{\begin{array}{cc}B_hB_k\hfill & \text{ if }h\overline{k}\hfill \\ B_hB_k\lambda _i\hfill & \text{ if }k=\overline{h}\hfill \end{array}`$
$`\delta _i^{}B_k^{}`$ $`=\delta _iB_k`$
$`B_h^{}\gamma _i^{}`$ $`=B_h\gamma _i`$
$`\delta _i^{}\gamma _i^{}`$ $`=\delta _i\gamma _i\lambda _i`$
for $`h,k`$ such that $`h_0=i=k_1`$. ∎
###### Lemma 36.
Let $`(s,s^{})Z_i^\lambda (d,v)`$ and suppose $`m_i0`$ or $`\lambda _i0`$ then
$$s\text{ is }\chi _m\text{ semistable }s^{}\text{ is }\chi _m^{}\text{ semistable }$$
###### Proof.
We prove only $``$. Let’s do first the case $`m_i0`$. If $`s`$ is $`\chi _m`$ semistable, then there exists $`\mathrm{\Delta }=\{A,B,\alpha _{}^{}\}`$ $`m`$-good such that $`f_\mathrm{\Delta }(s)0`$. Using the notation in 2.0.2 we have $`\phi _\mathrm{\Delta }=det\mathrm{\Psi }_\mathrm{\Delta }`$ where $`\mathrm{\Psi }_\mathrm{\Delta }:YZ`$ is a linear map. In our case we can write $`Z`$ as $`^{m_i}V_i\stackrel{~}{Z}`$ and we obseve that no $`V_i`$ summunds appear in $`Y`$ or $`\stackrel{~}{Z}`$.
Now we construct a new data $`\mathrm{\Delta }^{}=\{A^{},B^{},\alpha _{}^{}{}_{}{}^{}\}`$ such that $`f_\mathrm{\Delta }^{}(s^{})0`$ and $`f_\mathrm{\Delta }^{}`$ a $`\chi ^{}`$-covariant polynomial. Our strategy will be the following: we substitute each $`V_i`$ with the space $`T_i`$ in the space $`Z`$ and we add $`m_i`$ copies of $`V_i^{}`$ to $`Y`$. Let’s do it more precise: first of all the new data will not be $`m^{}`$ good so we have to define $`m_{}^{}{}_{j}{}^{+}`$ and $`m_{}^{}{}_{j}{}^{}`$:
1. $`m_{}^{}{}_{i}{}^{+}=0`$ and $`m_{}^{}{}_{i}{}^{}=m_i=m_i^+`$,
2. $`m_{j}^{}{}_{}{}^{}=m_j^{}`$ and $`m_{j}^{}{}_{}{}^{+}=m_j^++a_{ij}m_i^+`$ for all $`ji`$,
3. $`m_{}^{}{}_{}{}^{}=m^{}`$ and $`m_{}^{}{}_{}{}^{+}=m^++d_im_i^+`$.
Observe that $`m_{j}^{}{}_{}{}^{+}m_{j}^{}{}_{}{}^{}=m_j^{}`$ for all $`j`$ so our data will furnish a $`\chi ^{}`$ equivariant function. Moreover if we define
$$Z^{}=^{m_i}V_i\stackrel{~}{Z}\text{ and }Y^{}=^{m_i}T_iY$$
we observe that they have the numbers of $`V_j^{}`$, $`_a`$, $`_b`$ factors specified by $`m^{}`$. Now we construct the new data $`\mathrm{\Delta }^{}`$ in such a way that with respect to the decompositions above we have:
$`[\mathrm{\Psi }_\mathrm{\Delta }(s)]_{^{m_i}V_i\stackrel{~}{Z}}^Y`$ $`=\left(\begin{array}{c}(\mathrm{Id}b_i)\pi \\ \mathrm{\Phi }\end{array}\right),`$
$`[\mathrm{\Psi }_\mathrm{\Delta }^{}(s^{})]_{^{m_i}T_i\stackrel{~}{Z}}^{^{m_i}V_i^{}Y}`$ $`=\left(\begin{array}{cc}\mathrm{Id}a_i^{}& \pi \\ 0& \mathrm{\Phi }\end{array}\right).`$
If we construct a data with this property we observe that $`\mathrm{\Psi }_\mathrm{\Delta }(s)`$ is an isomorphism if and only if $`\mathrm{\Psi }_\mathrm{\Delta }^{}(s^{})`$ is an isomorphism. Hence $`f_\mathrm{\Delta }(s)0`$ implies $`f_\mathrm{\Delta }^{}(s^{})0`$ and the lemma is proved.
To construct the new data we choose a basis $`e_1,\mathrm{},e_{d_i}`$ of $`D_i`$ and we define the other elements of the data according to the following rules
1. $`A^{}=A`$,
2. if $`B=(b_1,\mathrm{},b_{m^+})`$ we set $`B^{}=(b_1,\mathrm{},b_{m^+},\underset{m_1times}{\underset{}{e_1,\mathrm{},e_1}},\mathrm{},\underset{m_1times}{\underset{}{e_{d_i},\mathrm{}e_{d_i}}})`$,
3. $`\alpha _{}^{}{}_{j_2,h_2}{}^{j_1,h_1}`$ for $`j_1i`$ and $`h_2m_{j_2}^+`$ is an element constructed according to case 1) in the previous lemma,
4. $`\alpha _{}^{}{}_{j_2,h_2}{}^{a}`$ for $`h_2m_{j_2}^+`$ is an element constructed according to case 2) in the previous lemma,
5. $`\alpha _{}^{}{}_{b_l}{}^{j_1,h_1}`$ for $`j_1i`$ and $`lm^+`$ is an element constructed according to case 3) in the previous lemma.
6. $`\alpha _{}^{}{}_{b_l}{}^{i,h}=\alpha _{}^{}{}_{j,k}{}^{i,h}=0`$ if $`lm^+`$ and $`km_j^+`$.
In this way we garantee that the projection of $`\mathrm{\Psi }_\mathrm{\Delta }^{}(s^{})`$ onto $`\stackrel{~}{Z}`$ is equal to $`\left(0\mathrm{\Phi }\right)`$. To define the remaining part of the new data we do not give details on the indexes, but we explain how to construct it. It is clear that we can choose $`\alpha _{}^{}{}_{}{}^{i,h}`$ for the remaining indeces * in such a way that the projection of $`\mathrm{\Psi }_\mathrm{\Delta }^{}(s^{})|_{^{m_i}V_i^{}}`$ on $`^{m_i}T_i`$ is equal to $`\mathrm{Id}a_i^{}`$. Finally we observe that a path $`\beta `$ from $`V_j`$ to $`V_i`$ with $`ji`$ has to go through a summand of $`T_i`$ so there exists a path $`\alpha `$ such that $`\beta (s)=b_i\alpha (s)`$. Now we use the previous lemma to change $`\alpha `$ with a $`\alpha ^{}`$ such that $`\beta (s)=b_i\alpha ^{}(s^{})`$. More generally if $`\beta `$ is an element of the path algebra of type $`(j,i)`$ with $`ji`$ then there exists an element of the $`b`$-path algebra $`\alpha ^{}`$ such that $`\beta (s)=b_i\alpha ^{}(s)`$. In this way we define the elements of the $`b`$-path algebra connecting summunds of $`Y`$ and summunds of $`^{m_i}T_i`$.
In the case $`m_i<0`$ we proceed in a similar way: we choose $`\mathrm{\Delta }`$ $`m`$-good and we have
$$Y=^{m_i}V_i\stackrel{~}{Y},Y^{}=^{m_i}T_i\stackrel{~}{Y},Z^{}=^{m_i}V_i^{}Z.$$
As in the previous case we can find a new data $`\mathrm{\Delta }^{}`$ such that:
$`[\mathrm{\Psi }_\mathrm{\Delta }(s)]_Z^{^{m_i}V_i\stackrel{~}{Y}}`$ $`=\left(\begin{array}{cc}\pi (\mathrm{Id}a_i)& \mathrm{\Phi }\end{array}\right),`$
$`[\mathrm{\Psi }_\mathrm{\Delta }^{}(s^{})]_{^{m_i}V_i^{}Z}^{^{m_i}T_i\stackrel{~}{Y}}`$ $`=\left(\begin{array}{cc}\mathrm{Id}b_i^{}& 0\\ \pi & \mathrm{\Phi }\end{array}\right).`$
Now to conclude that $`\mathrm{\Psi }_\mathrm{\Delta }^{}(s^{})`$ is an isomorphism if $`\mathrm{\Psi }_\mathrm{\Delta }(s)`$ is we need to know that $`b_i^{}`$ is an epimorphism and this is not garantee by $`(s,s^{})Z_i^\lambda (d,v)`$. But if $`\lambda _i0`$ then, since $`b_i^{}a_i^{}=\lambda _i`$, we have that $`b_i^{}`$ is surjective. ∎
###### Definition 37.
Let $`p`$ (resp. $`p^{}`$) be the projections of $`Z_i^\lambda (d,v)`$ on $`\mathrm{\Lambda }_\lambda (d,v)S(d,v)`$ (resp. $`\mathrm{\Lambda }_\lambda ^{}(d,v^{})S(d,v^{})`$). Suppose that $`m_i>0`$ or $`\lambda _i0`$ then we define
$$Z_i^{m,\lambda }=p^1\left(\mathrm{\Lambda }_{m,\lambda }(d,v)\right)=p_{}^{}{}_{}{}^{1}\left(\mathrm{\Lambda }_{m^{},\lambda ^{}}(d,v^{})\right).$$
We define also
$$G_{i,v}=\underset{ji}{}GL(V_j)\times GL(V_i)\times GL(V_i^{}).$$
Observe that there are natural projections from $`G_{i,v}`$ to $`G_v`$ and $`G_v^{}`$, therefore ther are natural actions of $`G_{i,v}`$ on $`S_i(d,v)`$, $`S_i(d,v^{})`$. Observe that here is a natural action of $`G_{i,v}`$ on $`Z_i^\lambda `$ and $`Z_i^{m,\lambda }`$ such that the projections $`p`$, $`p^{}`$ are equivariant.
###### Lemma 38.
Let $`s\mathrm{\Lambda }_{\lambda ,m}(d,v)`$ then
1. if $`\lambda _i0`$ then $`b_i`$ is epi and $`a_i`$ is mono,
2. if $`m_i>0`$ then $`b_i`$ is epi,
3. if $`m_i<0`$ then $`a_i`$ is mono.
###### Proof.
If $`\lambda _i0`$ then the result is clear by $`b_ia_i=\lambda _i`$. Suppose now that $`\lambda _i=0`$ and $`m_i>0`$. Let $`U_i=\mathrm{Im}b_i`$ and let $`V_i=U_iW_i`$. Define now a one parameter subgroup $`g(t)`$ of $`G_V`$ in the following way:
$$[g_i(t)]_{U_iW_i}^{U_iW_i}=\left(\begin{array}{cc}1& 0\\ 0& t^1\end{array}\right)\text{ and }g_j1\text{ for }ji$$
Since $`\mathrm{Im}b_iU_i`$ we have that there exists the limit $`lim_{t0}g(t)s=s_0`$. Let now $`n>0`$ and $`f`$ a $`\chi ^n`$-covariant function on $`S`$ such that $`f(s)0`$. Then
$$f(s_0)=\underset{t0}{lim}f(g(t)s)=\underset{t0}{lim}\underset{GL(V_i)}{\overset{nm_i}{det}}f(s)=\underset{t0}{lim}t^{nm_idimW_i}f(s)$$
So we must have $`dimW_i=0`$. The proof of the third case is completely similar to this one. ∎
###### Lemma 39 (see also Lusztig \[Lu:Q4\]).
If $`m_i>0`$ or $`\lambda _i0`$ then
1. $`p:Z_i^{m,\lambda }(d,v)\mathrm{\Lambda }_{m,\lambda }(d,v)`$ is a principal $`GL(V_i^{})`$ bundle,
2. $`p^{}:Z_i^{m,\lambda }(d,v)\mathrm{\Lambda }_{m^{},\lambda ^{}}(d,v^{})`$ is a principal $`GL(V_i)`$ bundle.
###### Proof.
Lusztig’s proof extend to this case without changes. Let’s prove for example 1. We have to prove: $`i.`$ that the action on the fiber is free, $`ii.`$ that it is transitive. First of all we observe that by the previous lemma if $`s\mathrm{\Lambda }_{m,\lambda }`$ then $`b_i(s)`$ is epi. In particular there exists $`a_i^{}:V_i^{}T_i`$ such that sequence (20) is exact, and clearly $`a_i^{}`$ is univoquely determined up to the action of $`GL(V_i^{})`$, moreover this action is free. So $`i.`$ and $`ii.`$ reduce to the following fact: if $`s\mathrm{\Lambda }_{m,\lambda }`$ and $`a_i^{}`$ is such that sequence (20) is exact, then there exists a unique, $`b_i^{}`$ such that $`a_i^{}b_i^{}=a_ib_i\lambda _i`$. Since $`a_i^{}`$ is mono the unicity is clear. To prove the existence we observe that it is equivalent to $`\mathrm{Im}a_i^{}\mathrm{Im}(a_ib_i\lambda _i)`$. But the last statement is clear since we have: $`\mathrm{Im}a_i^{}=\mathrm{ker}b_i`$ and $`b_i(a_ib_i\lambda _i)=0`$. ∎
###### Proposition 40.
If $`m_i>0`$ or $`\lambda _i0`$ then the projections $`p`$, $`p^{}`$ induces algebraic isomorphisms $`\overline{p}`$ , $`\overline{p}^{}`$:
$$\begin{array}{ccccc}\mathrm{\Lambda }_{m,\lambda }(d,v)//G_v& \underset{\overline{p}}{\overset{}{}}& Z_i^{m,\lambda }(d,v)//G_{i,v}& \underset{\overline{p}^{}}{\overset{}{}}& \mathrm{\Lambda }_{m^{},\lambda ^{}}(d,v^{})//G_v^{}\end{array}$$
###### Proof.
This proposition is a straightforward consequence of the previous lemma and the following general fact (see for example \[GIT\] Proposition 0.2): let $`G`$ be an algebraic groups over $``$ and $`X`$, $`Y`$ two irreducible algebraic variety over $``$; if $`G`$ acts on $`X`$ and $`\phi :XY`$ is such that for all $`yY`$ the fiber $`X_y`$ contains exactly one $`G`$-orbit then $`\phi `$ is a categorical quotient. If we apply this lemma to the projection $`p`$, (resp. $`p^{}`$) and to the group $`GL(V_i^{})`$ (resp. $`GL(V_i)`$) we obtain the required result. ∎
We can use this proposition to define the action of the generators of the Weyl group.
###### Definition 41.
Let $`i,\lambda ,m,d,v,\lambda ^{},m^{},v^{}`$ be as above, and suppose $`d_j0`$, $`v_j,v_j^{}0`$ for all $`j`$ then we define an isomorphism of algebraic variety
$$\mathrm{\Phi }_{s_i,\lambda ,m}^{d,v}:M_{m,\lambda }(d,v)M_{m^{},\lambda ^{}}(d,v^{})$$
in the following way:
1. if $`m_i>0`$ or $`\lambda _i0`$ then we set $`\mathrm{\Phi }_{s_i,\lambda ,m}^{d,v}=\overline{p}^{}\overline{p}^1`$,
2. if $`m_i<0`$ then we exchange the role of $`v,v^{}`$ in the previous construction: more precisely we observe that $`m_i^{}>0`$ so we can define $`\mathrm{\Phi }_{s_i,\lambda ^{},m^{}}^{d,v^{}}:M_{m^{},\lambda ^{}}(d,v^{})M_{m,\lambda }(d,v)`$ and we define $`\mathrm{\Phi }_{s_i,\lambda ,m}^{d,v}=\left(\mathrm{\Phi }_{s_i,\lambda ^{},m^{}}^{d,v^{}}\right)^1`$.
###### Remark 42.
To see that $`\mathrm{\Phi }_{s_i,\lambda ,m}^{d,v}`$ is univoquely defined we have to verify that if $`\lambda _i0`$ and $`m_i<0`$ the two definitions above coincide. This fact reduces easily to the following remark: if $`\lambda _i0`$ then
$$(s,s^{})Z_i^\lambda (d,v)(s^{},s)Z_i^\lambda ^{}(d,v^{}).$$
Let us prove, for example, the $``$ part. Since $`a_ib_i=a_i^{}b_i^{}+\lambda _i=a_i^{}b_i^{}\lambda _i^{}`$ the only thing we have to verify is that the sequence
$$\begin{array}{ccccccccc}0& & V_i& \stackrel{a_i}{}& T_i& \stackrel{b_i^{}}{}& V_i^{}& & 0\end{array}$$
is exact. The surjectivity of $`b_i^{}`$ and the injectivity of $`a_i`$ are a consequence of $`\lambda _i0`$. Since $`dimT_i=dimV_i+dimV_i^{}`$ we need only to prove that $`b_i^{}a_i=0`$. Observe that $`b_i^{}a_i=0`$ if and only if $`a_i^{}b_i^{}a_i=0`$ since also $`a_i^{}`$ is injective. Finally $`a_i^{}b_i^{}a_i=(a_ib_i\lambda _i)a_i=0`$.
### 3.2. Preliminaries
We saw how to define
$$\mathrm{\Phi }_{s_i,m,\lambda }^{d,v}:M_{m,\lambda }(d,v)M_{s_i(m),s_i(\lambda )}\left(s_i(d,v)\right)$$
in the case that $`(\lambda _i,m_i)0`$ and $`d,v,s_iv0`$. To define an action of the Weyl group we have now to garantee that coxeter relations hold. We will prove these relations in the next paragraph. Before doing it we observe that we have to garantee some conditions on $`m,\lambda `$ such that we will be able to define $`\mathrm{\Phi }_{s_i,\sigma m,\sigma \lambda }^{\sigma (d,v)}`$ for any element $`\sigma W`$: this condition will be $`(m,\lambda )𝒢_v`$ (31). We have also to say something about the case $`d_i<0`$ or $`v_i<0`$ for some $`iI`$.
In the case that $`d_i<0`$ for some $`i`$ then $`M_{m,\lambda }\left(\sigma (d,v)\right)=\mathrm{}`$ for all $`\sigma ,m,\lambda `$ by the very definition, so there is nothing to define.
The second trivial case is $`d=v=0`$. Indeed in this case we have $`M_{\sigma m,\sigma \lambda }\left(\sigma (d,v)\right)=\{0\}`$ so the definition is trivial.
The other two cases are threated in the two lemmas below.
In the following we fix $`d`$ such that $`d_i0`$ for all $`i`$. It will be convenient to define an affine action of $`W`$ on $`Q`$ by $`\sigma v=\sigma (vd)+d`$.
###### Lemma 43.
Let $`d0`$ and $`(m,\lambda )𝒢_v`$ if there exists $`\sigma `$ such that $`\sigma v0`$ then $`M_{m,\lambda }(d,v)=\mathrm{}`$.
###### Proof.
Suppose that $`\sigma `$ is an element of minimal length such that $`\sigma v0`$ and let $`l=\mathrm{}(\sigma )`$. We prove the lemma by induction on $`l`$. The case $`l=0`$ is trivial.
*Initial step: $`l=1`$.* If $`s_iv0`$ then we have $`0d_i+a_{ij}v_j<v_i`$. Hence $`dimT_i<dimV_i`$, $`u=(0,\mathrm{},\stackrel{𝑖}{1},0,\mathrm{})\stackrel{~}{U}_v`$ and $`(\lambda _i,m_i)0`$. So $`M_{m,\lambda }(d,v)=\mathrm{}`$ by lemma 38.
*Inductive step: if $`l2`$ then $`l1l`$.* Let $`\sigma =\tau s_i`$ with $`\mathrm{}(\tau )=l1`$ and $`v^{}=s_iv`$, $`\lambda ^{}=s_i\lambda `$, $`m^{}=s_im`$. By induction $`M_{m^{},\lambda ^{}}(d,v^{})=\mathrm{}`$ and, since $`l2`$, $`v^{}0`$. If $`(m_i,\lambda _i)0`$ then we can apply Proposition 40 and we obtain $`M_{m,\lambda }(d,v)M_{m^{},\lambda ^{}}(d,v^{})=\mathrm{}`$. If $`(m_i,\lambda _i)=0`$ then $`u=(0,\mathrm{},\stackrel{𝑖}{1},0,\mathrm{})\stackrel{~}{U}_v`$, hence $`v_i=0`$. Moreover $`\lambda ^{}=\lambda `$ and $`m^{}=m`$ so $`(m_i^{},\lambda _i^{})=0`$ and $`u=(0,\mathrm{},\stackrel{𝑖}{1},0,\mathrm{})U_v^{}`$. Hence $`v_i^{}=0`$ so $`v^{}=v`$ and $`\tau v0`$ against the minimality of $`\sigma `$. ∎
###### Lemma 44.
Let $`(I,H)`$ be connected, $`(m,\lambda )𝒢_v`$ and suppose $`d0`$ and $`\sigma v0`$ for all $`\sigma W`$. If there exists $`iI,\sigma W`$ such that $`\sigma (m,\lambda )=(m^{},\lambda ^{})`$ and $`(m_i^{},\lambda _i^{})=0`$ then $`d=v=0`$.
###### Proof.
Without loss of generality we can assume $`\sigma =1`$.
*First step*: $`v_i=0`$. This is clear since otherwise $`u=(0,\mathrm{},\stackrel{𝑖}{1},0,\mathrm{})U_v`$.
*Second step*: $`d_i=0`$ and $`v_j=0`$ for all $`j`$ such that $`a_{ij}`$. Let $`v^{}=s_iv`$ and observe that $`s_i\lambda =\lambda `$ and $`s_im=m`$. Then as in first step we have $`0=v_i^{}=d_i+_ja_{ij}v_j`$ from which the claim follows.
Let now $`W^{}=<\{s_j:a_{ij}0\text{ and }ji\}>`$. If $`(d,v)0`$ then there exists $`jI`$ and $`\sigma W^{}`$ such that $`a_{ij}0`$ and
$$n=d_j+\underset{hI}{}a_{jh}\stackrel{~}{v}_h>0.$$
where $`\stackrel{~}{v}=\sigma v`$. Since $`(\sigma \lambda )_i=\lambda _i=0=m_i=(\sigma m)_i`$ we can assume $`\sigma =1`$. Let now $`v^{}=s_is_jv`$, $`\lambda ^{}=s_is_j\lambda `$ and $`m^{}=s_is_jm`$, we have:
$`v_i^{}`$ $`=a_{ij}n`$ $`\lambda _i^{}`$ $`=a_{ij}\lambda _j`$ $`m_i`$ $`=a_{ij}m_j`$
$`v_j^{}`$ $`=n`$ $`\lambda _j^{}`$ $`=(a_{ij}^21)\lambda _j`$ $`m_j`$ $`=(a_{ij}^21)m_j.`$
Hence $`u=(0,\mathrm{},\stackrel{𝑗}{a_{ij}},0,\mathrm{},\stackrel{𝑖}{a_{ij}^21},0,\mathrm{})U_v^{}`$ and $`<u^{},\lambda ^{}>=<u^{},m^{}>=0`$ against $`(m,\lambda )𝒢_v`$. ∎
###### Remark 45.
the analougous lemma in the case of $`𝒢=\{(m,\lambda ,v):(m,grl)H^{\mathrm{}}\}`$ are more simple.
### 3.3. Relations
In this section we define an isomorphism of algebraic variety
$$\mathrm{\Phi }_{\sigma ,m,\lambda }^{d,v}:M_{m,\lambda }(d,v)M_{\sigma m,\sigma \lambda }\left(\sigma (d,v)\right).$$
in the case $`(m,\lambda )𝒢_v`$ or $`(m,\lambda )H^{\mathrm{}}`$. In the case $`d0`$ or in the case in which there exists $`\sigma W`$ such that $`\sigma v0`$ or in the case $`d=v=0`$ we have seen in the previous section that there is nothing to define or that the definition is trivial. In the remaing cases we observe that for all $`\tau ,i`$ we have $`(\tau (m)_i,\tau (\lambda )_i)0`$ by lemma 44. Hence we can define $`\mathrm{\Phi }_{\sigma ,m,\lambda }^{d,v}`$ by induction on $`\mathrm{}(\sigma )`$ by the formula
$$\mathrm{\Phi }_{\sigma ,m,\lambda }^{d,v}=\mathrm{\Phi }_{s_i,\tau m,\tau \lambda }^{\tau (d,v)}\mathrm{\Phi }_{\tau ,m,\lambda }^{d,v}.$$
(21)
Of course we have to prove that this definition is well given by checking Coxeter relations:
$$s_i^2=\mathrm{Id},s_is_j=s_js_i\text{ if }a_{ij}=0\text{ and }s_is_js_i=s_js_is_j\text{ if }a_{ij}=1$$
which in our situation take the following form:
$$\mathrm{\Phi }_{s_i,s_i\lambda ,s_im}^{s_i(d,v)}\mathrm{\Phi }_{s_i,\lambda ,m}^{d,v}=\mathrm{Id}$$
(22a)
$$\mathrm{\Phi }_{s_i,s_j\lambda ,s_jm}^{s_j(d,v)}\mathrm{\Phi }_{s_j,\lambda ,m}^{d,v}=\mathrm{\Phi }_{s_j,s_i\lambda ,s_im}^{s_i(d,v)}\mathrm{\Phi }_{s_i,\lambda ,m}^{d,v}$$
(22b)
$$\mathrm{\Phi }_{s_i,s_js_i\lambda ,s_js_im}^{s_js_i(d,v)}\mathrm{\Phi }_{s_j,s_i\lambda ,s_im}^{s_i(d,v)}\mathrm{\Phi }_{s_i,\lambda ,m}^{d,v}=\mathrm{\Phi }_{s_j,s_is_j\lambda ,s_is_jm}^{s_is_j(d,v)}\mathrm{\Phi }_{s_i,s_j\lambda ,s_jm}^{s_j(d,v)}\mathrm{\Phi }_{s_j,\lambda ,m}^{d,v}.$$
(22c)
The first of the two equations is clear by the very definition and remark 42. The second equation is trivial. We need to prove the third equation. We will need the following two simple lemmas of linear algebra which proofs are trivial.
###### Lemma 46.
Let $`V,W,X,Y,Z`$ be finite dimensional vector spaces and $`\alpha ,\beta ,\gamma ,\delta ,\epsilon ,\phi `$ linear maps between them as in the diagrams below. The diagram
$$\begin{array}{ccccccccc}0& & V& \stackrel{\left(\begin{array}{c}\alpha \\ \beta \\ \gamma \end{array}\right)}{}& WXY& \stackrel{\left(\begin{array}{ccc}\delta & 0& 1\\ 0& \epsilon & \phi \end{array}\right)}{}& YZ& & 0\end{array}$$
is exact if and only if the diagram
$$\begin{array}{ccccccccc}0& & V& \stackrel{\left(\begin{array}{c}\alpha \\ \beta \end{array}\right)}{}& WX& \stackrel{\left(\begin{array}{cc}\mathrm{\phi }\mathrm{\delta }& \epsilon \end{array}\right)}{}& Z& & 0\end{array}$$
is exact and $`\gamma =\delta \alpha `$.
###### Lemma 47.
Let $`U,V,W,X,Y,Z`$ be finite dimensional vector spaces and $`\alpha ,\beta ,\gamma ,\delta ,\epsilon ,\phi ,\psi ,\rho ,\sigma `$ linear maps between them as in the diagrams below such that $`\psi \rho :WXZ`$ is an epimorphism. Then the diagram
$$\begin{array}{ccccccc}0U& \stackrel{\left(\begin{array}{c}\alpha \\ \beta \\ \gamma \end{array}\right)}{}& VWX& \stackrel{\left(\begin{array}{ccc}\delta & 0& 1\\ \epsilon & \phi & 0\\ 0& \psi & \rho \end{array}\right)}{}& XYZ& \stackrel{\left(\begin{array}{ccc}\rho & \sigma & 1\end{array}\right)}{}& Z0\end{array}$$
is exact if and only if $`\gamma =\delta \alpha `$ , $`\psi =\sigma \varphi `$ , $`\rho \delta +\sigma \epsilon =0`$ and the diagram
$$\begin{array}{ccccccccc}0& & U& \stackrel{\left(\begin{array}{c}\alpha \\ \beta \end{array}\right)}{}& VW& \stackrel{\left(\begin{array}{cc}\epsilon & \phi \end{array}\right)}{}& Y& & 0\end{array}$$
is exact.
We fix now and $`i,j`$ such that $`a_{ij}=1`$ and we verifies (22c). Let
$`\lambda ^{}`$ $`=s_i\lambda `$ $`m^{}`$ $`=s_im`$ $`v^{}`$ $`=s_iv`$
$`\lambda ^{\prime \prime }`$ $`=s_j\lambda ^{}`$ $`m^{\prime \prime }`$ $`=s_jm^{}`$ $`v^{\prime \prime }`$ $`=s_jv^{}`$
$`\lambda ^{\prime \prime \prime }`$ $`=s_i\lambda ^{\prime \prime }`$ $`m^{\prime \prime \prime }`$ $`=s_im^{\prime \prime }`$ $`v^{\prime \prime \prime }`$ $`=s_iv^{\prime \prime }`$
$`\stackrel{~}{\lambda }`$ $`=s_j\lambda `$ $`\stackrel{~}{m}`$ $`=s_jm`$ $`\stackrel{~}{v}`$ $`=s_jv`$
$`\stackrel{~}{\stackrel{~}{\lambda }}`$ $`=s_i\stackrel{~}{\lambda }`$ $`\stackrel{~}{\stackrel{~}{m}}`$ $`=s_i\stackrel{~}{m}`$ $`\stackrel{~}{\stackrel{~}{v}}`$ $`=s_i\stackrel{~}{v}`$
First of all we observe that since relation (22a) holds we can assume that:
1. $`\lambda _i0`$ or $`m_i>0`$ and $`\lambda _j0`$ or $`m_j>0`$,
2. $`\lambda _j^{}0`$ or $`m_j^{}>0`$ and $`\stackrel{~}{\lambda }_i0`$ or $`\stackrel{~}{m}_i>0`$,
3. $`\lambda _i^{\prime \prime }0`$ or $`m_i^{\prime \prime }>0`$ and $`\stackrel{~}{\stackrel{~}{\lambda }}_j0`$ or $`\stackrel{~}{\stackrel{~}{m}}_j>0`$.
Define
$`Z_{iji}`$ $`=\{(s^{\prime \prime \prime },s)\mathrm{\Lambda }_{m^{\prime \prime \prime },\lambda ^{\prime \prime \prime }}(d,v^{\prime \prime \prime })\times \mathrm{\Lambda }_{m,\lambda }(d,v):s^{\prime \prime }S(d,v^{\prime \prime }),`$
$`\text{ and }s^{}S(d,v^{})\text{ such that }(s^{\prime \prime \prime },s^{\prime \prime })Z_i^{m^{\prime \prime },\lambda ^{\prime \prime }}(d,v^{\prime \prime }),`$
$`(s^{\prime \prime },s^{})Z_j^{m^{},\lambda ^{}}(d,v^{})\text{ and }(s^{},s)Z_i^{m,\lambda }(d,v)\}`$
$`Z_{jij}`$ $`=\{(s^{\prime \prime \prime },s)\mathrm{\Lambda }_{m^{\prime \prime \prime },\lambda ^{\prime \prime \prime }}(d,v^{\prime \prime \prime })\times \mathrm{\Lambda }_{m,\lambda }(d,v):\stackrel{~}{\stackrel{~}{s}}S(d,\stackrel{~}{\stackrel{~}{v}}),`$
$`\text{ and }\stackrel{~}{s}S(d,\stackrel{~}{v})\text{ such that }(s^{\prime \prime \prime },\stackrel{~}{\stackrel{~}{s}})Z_j^{\stackrel{~}{\stackrel{~}{m}},\stackrel{~}{\stackrel{~}{\lambda }}}(d,\stackrel{~}{\stackrel{~}{v}}),`$
$`(\stackrel{~}{\stackrel{~}{s}},\stackrel{~}{s})Z_i^{\stackrel{~}{m},\stackrel{~}{\lambda }}(d,\stackrel{~}{v})\text{ and }(\stackrel{~}{s},s)Z_j^{m,\lambda }(d,v)\}`$
Observe that $`(s^{\prime \prime \prime },s)Z_{iji}p_{m^{\prime \prime \prime }.\lambda ^{\prime \prime \prime }}^{d,v^{\prime \prime \prime }}(s^{\prime \prime \prime })=\mathrm{\Phi }_{s_i}\mathrm{\Phi }_{s_j}\mathrm{\Phi }_{s_i}(p_{m,\lambda }^{d,v}(s))`$ and that $`(s^{\prime \prime \prime },s)Z_{jij}p_{m^{\prime \prime \prime }.\lambda ^{\prime \prime \prime }}^{d,v^{\prime \prime \prime }}(s^{\prime \prime \prime })=\mathrm{\Phi }_{s_j}\mathrm{\Phi }_{s_i}\mathrm{\Phi }_{s_j}(p_{m,\lambda }^{d,v}(s))`$. So relation (22c) is equivalent to $`Z_{iji}=Z_{jij}`$.
Let now $`R_i=D_i_{h:h_1=i,h_0j}V_{h_0}`$, $`R_j=D_i_{h:h_1=i,h_0j}V_{h_0}`$ and observe that $`T_i=R_iV_j`$ and $`T_j=R_jV_i`$. Let $`k`$ be the only element of $`H`$ such that $`k_0=j`$ and $`k_1=i`$. Let $`\epsilon =\epsilon (k)`$. Define also $`A=A(s)=B_k(s)`$, $`B=B(s)=B_{\overline{k}}(s)`$ and for $`l=i,j`$ and $`\{l^{},l\}=\{i,j\}`$ set $`c_l=c_l(s)=\pi _{R_l}^{R_lV_l^{}}a_l(s)`$ and $`d_l=d_l(s)=b_l(s)|_{R_l}`$.
Let now $`(s,s^{\prime \prime \prime })\mathrm{\Lambda }_\lambda (d,v)\times \mathrm{\Lambda }_{\lambda ^{\prime \prime \prime }}(d,v^{\prime \prime \prime })`$ and set $`A^{}=A(s^{})`$, $`B^{}=B(s^{})`$, $`c_l^{}=c_l(s^{})`$ and $`d_l^{}=d_l(s^{})`$ for $`l\{i,j\}`$ and $`\{,^{\prime \prime \prime }\}`$.
If we apply lemmas 46 and 47 to our situation we obtain the following result: $`(s,s^{\prime \prime \prime })Z_{iji}`$ of and only there exists vector spaces $`V_i^{},V_j^{},V_i^{\prime \prime },V_j^{\prime \prime }`$ and linear maps $`A^{},B^{},c_i^{},d_i^{},c_j^{},d_j^{},A^{\prime \prime },B^{\prime \prime },c_i^{\prime \prime },d_i^{\prime \prime },c_j^{\prime \prime },d_j^{\prime \prime }`$ sucht that:
1. $`dimV_l^{}=v_l^{}`$ for $`l\{i,j\}`$ and $`\{^{},^{\prime \prime }\}`$,
2. for each $`\{^{},^{\prime \prime }\}`$ and $`l\{i,j\}`$ $`A^{}\text{Hom}(V_i^{},V_j^{})`$, $`B^{}\text{Hom}(V_j^{},V_i^{})`$, $`c_l\text{Hom}(V_l^{},R_l^{})`$ and $`d_l\text{Hom}(R_l^{},V_l^{})`$,
3. $`V_j^{\prime \prime \prime }=V_j^{\prime \prime }`$, $`c_j^{\prime \prime \prime }=c_j^{\prime \prime }`$, $`d_j^{\prime \prime \prime }=d_j^{\prime \prime }`$ and
$`c_i^{\prime \prime \prime }d_i^{\prime \prime \prime }`$ $`=c_id_i\lambda _i\lambda _j`$ $`c_i^{\prime \prime \prime }B^{\prime \prime \prime }`$ $`=c_i^{}B^{\prime \prime }`$
$`A^{\prime \prime \prime }d_i^{\prime \prime \prime }`$ $`=A^{\prime \prime }d_i^{}`$ $`\epsilon A^{\prime \prime \prime }B^{\prime \prime \prime }`$ $`=\epsilon A^{\prime \prime }B^{\prime \prime }\lambda _j`$
4. $`V_i^{\prime \prime }=V_i^{}`$, $`c_i^{\prime \prime }=c_i^{}`$, $`d_i^{\prime \prime }=d_i^{}`$ and
$`c_j^{\prime \prime }d_j^{\prime \prime }`$ $`=c_jd_j\lambda _i\lambda _j`$ $`c_j^{\prime \prime }A^{\prime \prime }`$ $`=c_jA^{}`$
$`B^{\prime \prime }d_j^{\prime \prime }`$ $`=B^{}d_j`$ $`\epsilon A^{\prime \prime }B^{\prime \prime }`$ $`=\epsilon A^{}B^{}+\lambda _i+\lambda _j`$
5. $`V_j^{}=V_j`$, $`c_j^{}=c_j`$, $`d_j^{}=d_j`$ and
$`c_i^{}d_i^{}`$ $`=c_id_i\lambda _i`$ $`c_i^{}B^{}`$ $`=c_iB`$
$`A^{}d_i^{}`$ $`=Ad_i`$ $`\epsilon A^{}B^{}`$ $`=\epsilon AB\lambda _i`$
6. $`\epsilon c_i^{}B^{\prime \prime }A^{\prime \prime \prime }+c_i^{}d_i^{}c_i^{\prime \prime \prime }=0`$ and $`\epsilon A^{}B^{\prime \prime }=d_ja_j^{\prime \prime }`$,
7. the following diagrams are exact
$$\begin{array}{ccccccccc}0& & V_i^{\prime \prime \prime }& \stackrel{\left(\begin{array}{c}c_i^{\prime \prime \prime }\\ c_j^{\prime \prime }A^{\prime \prime \prime }\end{array}\right)}{}& R_iR_j& \stackrel{\left(\begin{array}{cc}Ad_i& d_j\end{array}\right)}{}& V_j& & 0\\ 0& & V_j^{\prime \prime }& \stackrel{\left(\begin{array}{c}c_j^{\prime \prime }\\ c_i^{}B^{\prime \prime }\end{array}\right)}{}& R_jR_i& \stackrel{\left(\begin{array}{cc}Bd_j& d_i\end{array}\right)}{}& V_i& & 0\\ 0& & V_i^{}& \stackrel{\left(\begin{array}{c}c_i^{}\\ A^{}\end{array}\right)}{}& R_iV_j& \stackrel{\left(\begin{array}{cc}d_i& \epsilon B\end{array}\right)}{}& V_i& & 0\end{array}$$
###### Remark 48.
The first condition in point 6) is equivalent to $`\epsilon B^{\prime \prime }A^{\prime \prime \prime }=d_i^{}c_i^{\prime \prime \prime }`$. Indeed this condition is certainly sufficient. To prove the necessity observe that by the injectivity of $`a_i^{}=(c_i^{}A^{})^t`$ it is enough to prove $`\epsilon c_i^{}B^{\prime \prime }A^{\prime \prime \prime }+c_i^{}d_i^{}c_i^{\prime \prime \prime }=0`$ and $`\epsilon A^{}B^{\prime \prime }A^{\prime \prime \prime }+A^{}d_i^{}c_i^{\prime \prime \prime }=0`$. The first equation is the first condition in point 6) and the second one is a consequence of $`\epsilon A^{}B^{\prime \prime }=d_jc_i^{\prime \prime \prime }`$, $`A^{}d_i^{}=Ad_i`$ and the exactness of the first sequence.
###### Remark 49.
The condition $`(s,s^{\prime \prime \prime })Z_{jij}`$ can be expressed in a similar way. In the prevoius conditions we have only to change $`i`$ with $`j`$ and $`\epsilon `$ with $`\epsilon `$.
We will prove now $`Z_{iji}Z_{jij}`$. To do it we supose that $`A^{},\mathrm{},d_j^{\prime \prime }`$ are given as above and we construct $`\stackrel{~}{A},\stackrel{~}{B},\stackrel{~}{c}_i,\stackrel{~}{d}_i,\stackrel{~}{c}_j,\stackrel{~}{d}_j,\stackrel{~}{\stackrel{~}{A}},\stackrel{~}{\stackrel{~}{B}},\stackrel{~}{\stackrel{~}{c}}_i,\stackrel{~}{\stackrel{~}{d}}_i,\stackrel{~}{\stackrel{~}{c}}_j,\stackrel{~}{\stackrel{~}{d}}_j`$ such that they satisfy the conditions. for $`(s,s^{\prime \prime \prime })Z_{jij}`$.
First step: construction of $`\stackrel{~}{A},\stackrel{~}{B},\stackrel{~}{c}_i,\stackrel{~}{c}_j,\stackrel{~}{d}_i,\stackrel{~}{d}_j`$. Choose $`\stackrel{~}{s}`$ such that $`(\stackrel{~}{s},s)Z_j^{\chi ,\lambda }`$ and define $`\stackrel{~}{A}=A(\stackrel{~}{s})`$, $`\stackrel{~}{B}=B(\stackrel{~}{s})`$, $`\stackrel{~}{c}_l=c_l(\stackrel{~}{s})`$ and $`\stackrel{~}{d}_l=d_l(\stackrel{~}{s})`$ for $`l\{i,j\}`$ .
Now I claim that there exists unique $`\stackrel{~}{\stackrel{~}{A}}:V_i^{\prime \prime \prime }\stackrel{~}{V}_j`$ and $`\stackrel{~}{\stackrel{~}{B}}:\stackrel{~}{V}_jV_i^{\prime \prime \prime }`$ such that:
$$\{\begin{array}{cc}\stackrel{~}{c}_j\stackrel{~}{\stackrel{~}{A}}=c_j^{\prime \prime }A^{\prime \prime \prime }\hfill & \\ \stackrel{~}{B}\stackrel{~}{\stackrel{~}{A}}=\epsilon d_ic_i^{\prime \prime \prime }\hfill & \end{array}\text{ and }\{\begin{array}{cc}\stackrel{~}{\stackrel{~}{A}}\stackrel{~}{\stackrel{~}{B}}=\stackrel{~}{A}\stackrel{~}{B}\epsilon \lambda _i\epsilon \lambda _j\hfill & \\ c_i^{\prime \prime \prime }\stackrel{~}{\stackrel{~}{B}}=c_i\stackrel{~}{B}\hfill & \end{array}$$
*Unicity of $`\stackrel{~}{\stackrel{~}{A}}`$*: since the map $`\stackrel{~}{a}_j=(\stackrel{~}{c}_j\epsilon \stackrel{~}{B})^t`$ is injective the unicity is clear.
*Existence of $`\stackrel{~}{\stackrel{~}{A}}`$*: to prove the existence of $`\stackrel{~}{\stackrel{~}{A}}`$ is enough to prove:
$$\mathrm{Im}\left(\begin{array}{c}c_j^{\prime \prime }A^{\prime \prime \prime }\\ \epsilon d_ic_i^{\prime \prime \prime }\end{array}\right)\mathrm{Im}\left(\begin{array}{c}\stackrel{~}{c}_j\\ \stackrel{~}{B}\end{array}\right)=\mathrm{ker}\left(\begin{array}{cc}d_j& \epsilon A\end{array}\right).$$
So the thesis follows from $`d_jc_j^{\prime \prime }A^{\prime \prime \prime }+Ad_ic_i^{\prime \prime \prime }=0`$.
Let now $`\stackrel{~}{\stackrel{~}{a}}_i=(c_i^{\prime \prime \prime }\stackrel{~}{\stackrel{~}{A}})^t`$. I claim that $`\stackrel{~}{\stackrel{~}{a}}_i`$ is injective and that $`\mathrm{Im}\stackrel{~}{\stackrel{~}{a}}_i=\mathrm{ker}(d_i\epsilon \stackrel{~}{B})=\mathrm{ker}\stackrel{~}{b}_i`$. First of all observe that since $`\stackrel{~}{m}_i>0`$ or $`\lambda _i0`$, $`\stackrel{~}{b}_i`$ is surjective. Observe also that
$$\left(\begin{array}{cc}\stackrel{~}{c}_j& 0\\ 0& \mathrm{Id}_{V_i^{\prime \prime \prime }}\end{array}\right)\left(\begin{array}{c}\stackrel{~}{\stackrel{~}{A}}\\ c_i^{\prime \prime \prime }\end{array}\right)=\left(\begin{array}{c}c_j^{\prime \prime }A^{\prime \prime \prime }\\ c_i^{\prime \prime \prime }\end{array}\right).$$
So $`\stackrel{~}{\stackrel{~}{a}}_i`$ is injective as claimed. Now since $`dimR_i+dim\stackrel{~}{V}_j=dimV_i^{\prime \prime \prime }+dimV_i`$ to prove the last part of the claim it is enough to check that $`\stackrel{~}{b}_i\stackrel{~}{\stackrel{~}{a}}_i=0`$. Indeed
$$\stackrel{~}{b}_i\stackrel{~}{\stackrel{~}{a}}_i=d_ic_i^{\prime \prime \prime }+\epsilon \stackrel{~}{B}\stackrel{~}{\stackrel{~}{A}}=0.$$
*Unicity of $`\stackrel{~}{\stackrel{~}{B}}`$*: this is a consequence of $`\stackrel{~}{\stackrel{~}{a}}_i`$ injective.
*Existence of $`\stackrel{~}{\stackrel{~}{B}}`$*: As for the existence of $`\stackrel{~}{\stackrel{~}{A}}`$ this is equivalent to
$$\mathrm{Im}\left(\begin{array}{c}c_i\stackrel{~}{B}\\ \stackrel{~}{A}\stackrel{~}{B}\epsilon \lambda _i\epsilon \lambda _j\end{array}\right)\mathrm{Im}\left(\begin{array}{c}c_i^{\prime \prime \prime }\\ \stackrel{~}{\stackrel{~}{A}}\end{array}\right)=\mathrm{ker}\left(\begin{array}{cc}d_i& \epsilon \stackrel{~}{B}\end{array}\right).$$
So the thesis follows from $`\epsilon \stackrel{~}{B}\stackrel{~}{A}\stackrel{~}{B}\lambda _i\stackrel{~}{B}\lambda _j\stackrel{~}{B}+d_ic_i\stackrel{~}{B}=0`$.
Finally we set
$`\stackrel{~}{\stackrel{~}{V}}_i`$ $`=V_i^{\prime \prime \prime }`$ $`\stackrel{~}{\stackrel{~}{c}}_i`$ $`=c_i^{\prime \prime \prime }`$ $`\stackrel{~}{\stackrel{~}{d}}_i`$ $`=d_i^{\prime \prime \prime }`$
$`\stackrel{~}{\stackrel{~}{V}}_i`$ $`=\stackrel{~}{V}_j`$ $`\stackrel{~}{\stackrel{~}{c}}_j`$ $`=\stackrel{~}{c}_j`$ $`\stackrel{~}{\stackrel{~}{d}}_j`$ $`=\stackrel{~}{d}_j.`$
The verification of all the conditions is now straightforward.
The inclusion $`Z_{jij}Z_{iji}`$ can be proved similarly and equation (19) is clear by definition. So Proposition 32 is proved.
## 4. A representation of the Weyl group
In this section, following Nakajima \[Na1\], we show how to use the above action to construct an action of the Weyl group on the homology of quiver varieties. Maybe this action is related with the one constructed by Slodowy in the case of flag varieties (\[Slodowy\], ch.4).
First we recall some general about the action of the Weyl group. Let $`Z^{}=Q^{}_{}`$ and $`Z=Q_{}P`$. On $`Z`$, $`Z^{}`$ there is a natural action of $`W`$.
###### Lemma 50.
For all $`uZ^{}`$ the set $`Wu`$ is discrete.
###### Lemma 51.
Consider the action of $`W`$ on $`(Z^{})`$. If $`p(Z^{})`$ then
$$\overline{Wp}\text{ is countable}.$$
If $`p(Z^{})`$ we define $`H_p=\{xP_{}:<x,p>=0\}`$
###### Lemma 52.
If $`p(Z^{})`$ then
$$\overline{WH_p}=\underset{q\overline{Wp}}{}H_q.$$
We define $`=\overline{W_v_U}`$ and $`=𝔷`$. By the previous lemmas $``$ is the union of a countable number of real codimension $`3`$ subspaces in $``$ and in particular $``$ is simply connected. We need also the following definition
$`K`$ $`=\{u^I:{\displaystyle u_i\alpha _i\text{ is dominant and supp}u\text{ is connected }}\}`$
$`P_0`$ $`=\{pP:p\text{ is dominant and }<u^{},p>2\text{ for all }uK\}.`$
Now we choose $`d,v`$ such that $`\overline{d}=\overline{v}`$.
###### Lemma 53.
If $`\overline{d}P_0`$ then $`\stackrel{~}{\mu }`$ is surjective and is a locally trivial bundle over $``$.
###### Proof.
By Proposition 10.5 and Corollary 10.6 in \[Na2\] there exists a closed orbit $`Gs`$ in $`\mathrm{\Lambda }_0(d,v)`$ with trivial stabilizer. Then by Proposition 9 there exists $`tGs`$ such that $`\stackrel{~}{\mu }(t)=0`$ and by lemma 5 and 3 $`d\stackrel{~}{\mu }_t`$ is surjective. Now the surjectivity follows by homogeneity.
The local triviality over $``$ follows also from lemma 5 and 3. ∎
Now consider
$`R`$ $`=\{\lambda Z:(0,\lambda )\},`$
$`\mathrm{\Lambda }(d,v)`$ $`=\{(\lambda ,s)Z\times S:s\mathrm{\Lambda }_\lambda \},`$
$`M(d,v)`$ $`=\mathrm{\Lambda }(d,v)//G_v\text{ and }p:M(d,v)Z\text{ the projection}`$
$`𝔏(d,v)`$ $`=\{(\zeta ,s)\times S:s𝔏_\zeta \},`$
$`𝔐(d,v)`$ $`=𝔏(d,v)/U(V)\text{ and }\stackrel{~}{p}:𝔐(d,v)\text{ the projection}`$
We have the following commutative diagram
By Proposition 10 the diagram is a pull back and by lemma 53 $`p`$ and $`\stackrel{~}{p}`$ are locally trivial over $`R`$, $``$. We call $`M_R=p^1(R)`$ and $`𝔐_{}=\stackrel{~}{p}^1()`$.
Now consider the complex $`=Rp_{}(_{M_R})=ı_M^1R\stackrel{~}{p}_{}(_𝔐_{})`$ which is cohomologically a locally constant complex. We observe now that $`\mathrm{\Pi }_1()`$ is trivial so $`R\stackrel{~}{p}_{}(_𝔐_{})`$ is isomorphic to cohomologically constant complex on $``$ so it is $``$ on $`R`$. In particular for any $`x,yR`$ we have a canonically isomorphism
$$\psi _{x,y}^i:H^i(_x)H^i(_y).$$
Now observe that by Proposition 32 there is an action on $`W`$ on $`R,M_R`$ and that $`p`$ is equivariant with respect to this action. So we can define a $`W`$ action on $`H^i(M_{0,\lambda }(d,v),)`$ by
$$\sigma (c)=\psi _{\sigma \lambda ,\lambda }^iH^i(\mathrm{\Phi }_{\sigma ,0,\lambda }^{d,v})(c)$$
for any $`\sigma W`$. To verify that this is an action we have only to verify that
$$\psi _{\sigma \lambda ^2,\lambda ^1}^iH^i(\mathrm{\Phi }_{\sigma ,0,\lambda ^2}^{d,v})\psi _{\lambda ^1,\lambda ^2}^i(c)=\psi _{\sigma \lambda ^1,\lambda ^1}^iH^i(\mathrm{\Phi }_{\sigma ,0,\lambda ^1}^{d,v})(c).$$
Since $`R`$ is connected and $`H^i(M,)`$ is discrete this is clear. S o we have proved the following corollary.
###### Corollary 54.
If $`d=v`$ and $`(0,\lambda )`$ then there is an action of $`W`$ on $`H^i(M_{0,\lambda }(d,v),)`$.
###### Remark 55.
If $`m_+=(1,\mathrm{},1)`$ and $`\lambda =0`$ it is easy to see that $`d\mu _s`$ is surjetive for all $`s\mathrm{\Lambda }_{m_+,0}(d,v)`$. Then by lemma 53 there is a canonical isomorphism $`H_{}(M_{m_+,0}(d,v))H_{}(M_{0,\lambda }(d,v))`$ if $`(0,\lambda )`$. So by Nakajima’s Theorem (Theorem 10.2 \[Na2\]) it is natural to make the following conjecture:
###### Conjecture 56.
Let $`top=\frac{1}{2}dimH_{}(M_{0,\lambda }(d,v))`$ then
$$H^{top}(M_{(0,\lambda )}(d,v),)\left(L_d\right)_0$$
where $`\left(L_d\right)_0`$ is the $`0`$-weight space of the Kac-Moody algebra associated to the quiver of heighest weight $`_id_i\overline{\omega }_i`$.
## 5. Reduction to the dominant case
As a consequence of Proposition 32 we see that if $`(m,\lambda )𝒢_v`$ then there exists $`\sigma W`$ and $`v^{}=\sigma v`$ such that $`dv^{}`$ is dominant and $`M_{\sigma m,\sigma \lambda }(d,v^{})M_{m,\lambda }(d,v)`$. We generalize now this result to arbitrary $`\lambda `$.
On $`Q`$ we consider the following order: $`v^{}v`$ if and only if $`v_i^{}v_i`$.
We consider now the following construction: let $`v^{}v`$ and fix an embedding $`V_i^{}V_i`$ and a complement $`U_i`$ of $`V^{}`$ in $`V_i`$, then we can define a map $`\stackrel{~}{ȷ}:S(d,v^{})S(d,v)`$ through:
$$\stackrel{~}{ȷ}(B^{},\gamma ^{},\delta ^{})=(\left(\begin{array}{cc}B^{}& 0\\ 0& 0\end{array}\right),\left(\begin{array}{c}\gamma ^{}\\ 0\end{array}\right),\left(\begin{array}{cc}\delta ^{}& 0\end{array}\right))$$
(23)
where the matrices of the new triple represents the maps described through the decomposition $`V_i=V_i^{}U_i`$.
Suppose now that $`(m_i,\lambda _i)=0`$ for all $`i`$ such that $`v_i^{}v_i`$. Then it is easy to see that this map restrict to a map $`ȷ_r:\mathrm{\Lambda }_{m,\lambda }(d,v^{})\mathrm{\Lambda }_{m,\lambda }(d,v)`$ and so enduces a map $`ȷ_v^v^{}=ȷ:M_{0,\lambda }(d,v^{})M_{0,\lambda }(d,v)`$.
###### Lemma 57.
$`ȷ`$ is a closed immersion
###### Proof.
We prove that the map $`ȷ^{\mathrm{}}:[\mathrm{\Lambda }_\lambda (d,v)]^{G(v)}[\mathrm{\Lambda }_\lambda (d,v^{})]^{G(v^{})}`$ is surjective. By proposition 18 this follows by the following two identities:
$$\mathrm{Tr}\left(\alpha \left(ȷ(s)\right)\right)=\mathrm{Tr}\left(\alpha (s)\right)\text{ and }\beta \left(ȷ(s)\right)=\beta (s)$$
for each $`B`$-path $`\alpha `$ and for each admissible path $`\beta `$. ∎
###### Lemma 58.
If $`2v_i>d_i+_{jI}a_{ij}v_j`$ and $`v^{}=v\alpha _i`$ then $`ȷ`$ is an isomorphism of algebraic varieties
###### Proof.
It’s enough to prove that $`ȷ`$ is surjective. Let $`s=(B,\gamma ,\delta )\mathrm{\Lambda }_0(d,v)`$ and consider the sequence (see (4) for the notation) :
$$\begin{array}{ccccc}T_i& \stackrel{b_i}{}& V_i& \stackrel{a_i}{}& T_i.\end{array}$$
Since $`b_ia_i=0`$ and $`2dimV_i>dimT_i`$ we have that $`b_i`$ is not surjective or that $`a_i`$ is not injective.
Suppose that $`b_i`$ is not surjective, then up to the action of $`G_v`$ we can assume that $`\mathrm{Im}b_iv_i^{}`$. Then, for $`t^{}`$ consider $`g_t=(g_{j,t})G_v`$ with
$$g_i=\left(\begin{array}{cc}\mathrm{Id}_{v_i^{}}& 0\\ 0& t^1\end{array}\right)\text{ and }g_j=\mathrm{Id}_{V_j}\text{ for }ji.$$
Then
1. $`g_{i,t}B_h=B_h`$ if $`h_1=i`$ and $`g_i\gamma _i=\gamma _i`$, since $`\mathrm{Im}B_h,\mathrm{Im}\gamma _i\mathrm{Im}b_iv_i^{}`$,
2. $`lim_{t0}B_hg_{i,t}^1=B_h`$ if $`h_0=i`$ and $`\delta _ig_i^1=\delta _i`$
So $`lim_{t0}g_ts=s^{}`$ and it is clear that $`s^{}\stackrel{~}{ȷ}(\mathrm{\Lambda }_0(d,v^{}))`$ and that $`p_0(s)=p_0(s^{})\mathrm{Im}ȷ`$.
If $`b_j`$ is surjective and $`a_i`$ is not injective the argument is similar. ∎
###### Proposition 59.
For all $`\lambda `$ and for all $`d0,v0`$ there exists $`v^{}`$ and $`\sigma W`$ such that $`\overline{d}\overline{v}^{}`$ is dominant and
$$M_{0,\mathrm{\sigma }\mathrm{\lambda }}(d,v^{})M_{0,\lambda }(d,v).$$
###### Proof.
We prove this proposition by induction on the order $``$ on $`Q`$.
First step: $`v=0`$. If $`v=0`$ we can take $`v^{}=v`$ and $`\sigma =1`$.
Inductive step. If $`dv`$ is not dominant then there exists $`i`$ such that $`2v_i>d_i+a_{ij}v_j`$.
If $`\lambda _i0`$ we observe that $`s_iv=v^{}<v`$ (that is $`v^{}v`$ and $`v^{}v`$) and that $`M_{s_im,s_i\lambda }(d,v^{})M_{m,\lambda }(d,v)`$ and so we can apply the inductive hypothesis.
If $`\lambda _i=0`$ we apply the previous lemma and the inductive hypothesis. ∎
## 6. On normality and connectdness of quiver variety in the finite type case
In this section we resctrict our attention to the case of quiver varieties of finite type and to the case $`m=(1,\mathrm{},1)`$ and $`\lambda =0`$ and we fix $`d,v`$. By remark 2 we can assume without loss of generality that $`v_i>0`$ for all $`i`$. We would like to prove the following conjecture:
###### Conjecture 60.
$`M_{m,0}(d,v)`$ is connected and $`M_{0,0}(d,v)`$ is normal.
###### Remark 61.
If $`\stackrel{~}{m}=(m_1,\mathrm{},m_n)_+^I`$ it is easy to see that $`\mathrm{\Lambda }_{\stackrel{~}{m},0}=\mathrm{\Lambda }_{m,0}`$ ao in particular $`M_{m,0}(d,v)`$ is smooth. Instead $`M_{0,0}`$ is a cone so it is clearly connected.
By proposition 59it is enough to prove the theorem in the case $`dv`$ dominant. Unfortunately I’m not able to prove the conjecture only in the case $`dv`$ regular: $`<dv,\alpha _i>>0`$ for all $`i`$. To prove the conjecture in this case we will use the following stratification introduced by Lusztig in \[Lu:Q4\].
###### Definition 62.
For any $`sS`$ and $`iI`$ let
$$V_i^+=V_i^+(s)=\underset{\alpha \text{ a }B\text{path}:\alpha _1=i}{}\mathrm{Im}(\alpha (s)\gamma _{\alpha _0})$$
If $`v^{}=(v_1^{},\mathrm{},v_n^{})^n`$ we define
$$\mathrm{\Lambda }^v^{}=\{s\mathrm{\Lambda }_0(d,v):dimV_i^+(s)=v_i^{}\}.$$
Observe that $`\mathrm{\Lambda }^v=\mathrm{\Lambda }_{m_+,0}(d,v)`$. To prove our result we will use the following lemma of Lusztig.
###### Lemma 63 (Lusztig: \[Lu:Q4\] Proposition 4.5 and Proposition 5.3).
If $`0v_i^{}v_i`$ for each $`i`$ then
$$dim\mathrm{\Lambda }^v^{}(d,v)=dimS\underset{iI}{}dimgl(V_i)<(vv^{})^{},dv>\frac{1}{2}<(vv^{})^{},vv^{}>$$
Our result follows trivially from the following lemma.
###### Lemma 64.
1) If $`dv`$ is dominant then $`\mathrm{\Lambda }_{0,0}(d,v)`$ is a complete intersection.
2) If $`dv`$ is regular then $`\mathrm{\Lambda }_{0,0}`$ is normal and irreducible and $`\mathrm{\Lambda }_{m_+,0}(d,v)`$ is connected.
###### Proof.
Observe that $`\mathrm{\Lambda }_{0,0}(d,v)=\mu ^1(0)`$ so each irreducible component of $`\mathrm{\Lambda }_{0,0}(d,v)`$ must have dimension at least $`dimS_idimgl(V_i)=\delta _V`$.
Suppose now that $`dv`$ is dominant. By Nakajima’s theorem (\[Na2\] Theorem 10.2) $`M_{m_+,0}`$ is not empty. Observe also that by Proposition 14 $`\mathrm{\Lambda }_{m_+,0}(d,v)`$ is a smooth subset of $`\mathrm{\Lambda }_{0,0}(d,v)`$ of dimension $`\delta _V`$.
It is well known that $`\mathrm{\Lambda }_{m_+,0}(d,v)=\mathrm{\Lambda }^v`$. Hence
$$\mathrm{\Lambda }_{0,0}(d,v)\mathrm{\Lambda }_{m_+,0}(d,v)=\underset{v^{}v\text{ and }v^{}v}{}\mathrm{\Lambda }^v^{}.$$
By the lemma above we have that if $`v^{}v`$ and $`v^{}v`$ $`dim\mathrm{\Lambda }^v^{}<\delta _V`$. So $`\mathrm{\Lambda }_{m_+,0}(d,v)`$ must be dense in $`\mathrm{\Lambda }_{0,0}(d,v)`$ and $`\mathrm{\Lambda }_{0,0}(d,v)`$ is a complete intersection. Moreover if $`dv`$ is regular we have that $`dim\mathrm{\Lambda }^v^{}<\delta _V1`$ so the singular locus has codimension at least two and normality and irreducibility follows. Finally by our discussion it is clear that if $`\mathrm{\Lambda }_{m_+,0}(d,v)`$ is disconnected then $`\mathrm{\Lambda }_{0,0}(d,v)`$ is not irreducible. ∎
###### Remark 65.
In the lemma we can substitute $`\mathrm{\Lambda }_{m_+,0}(d,v)`$ with any other subset $`Reg`$ of regular points in $`\mathrm{\Lambda }(d,v)`$. In this way is indeed possible to improve a little bit the theorem but Crawley-Boevey explained me that this strategy cannot work in general becouse there are cases where $`dv`$ is dominant and $`\mathrm{\Lambda }_{0,0}(d,v)`$ is not normal. It should be also pointed out that Crawley-Boevey proved the connectdness in complete generality (\[CrBo\]). He said me that is also able to prove normality for a much bigger class of quiver varieties.
|
warning/0003/math0003070.html
|
ar5iv
|
text
|
# Untitled Document
Statement ≥of ≥the
Alexandru ≥Conjecture
The purpose of this text is to add (at least conjecturally) some more items to the list of analogies between the category $``$ of Harish-Chandra modules and the category $`𝒪`$ of Bernstein-Gelfand-Gelfand which has been established by Bernstein-Gelfand-Gelfand, Vogan, Beilinson, Ginzburg, Soergel and others. These analogies have been suggested by confronting some observations about $`p`$-integrable harmonic forms on real hyperbolic space with results of the people mentioned above about the category $`𝒪`$. To the reader more familiar with $`L^p`$ harmonic forms than with the category $`𝒪`$ my advice is to read in parallel Parts A and B of section 1 (Part B being a detailed example).
1. The ≥main ≥statements
Part A. The Weak Alexandru Conjecture
(1.1) Setting.
$`G`$ is a connected semisimple Lie group $`G`$ with finite center,
$`KG`$ is a maximal compact subgroup,
$`𝔤𝔨`$ are the complexified Lie algebras of $`G`$ and $`K`$,
$`𝔟𝔤`$ is a Borel subalgebra,
$`𝔥`$ a Cartan subalgebra of $`𝔤`$ contained in $`𝔟`$.
For any pair $`𝔪𝔩`$ of (complex) Lie algebras and any $`𝔩`$-module $`V`$, say that $`V`$ is $`𝔪`$-finite if it is a sum of finite dimensional sub-$`𝔪`$-modules, and that $`V`$ is an $`(𝔩,𝔪)`$-module if it is $`𝔪`$-finite and $`𝔪`$-semisimple. The category $`𝒪`$ of BGG-modules is the full subcategory of $`𝔤`$-mod whose objects are the $`𝔟`$-finite $`(𝔤,𝔥)`$-modules of finite length ; whereas the category $``$ of Harish-Chandra modules is the full subcategory of $`𝔤`$-mod whose objects are those $`(𝔤,𝔨)`$-modules of finite length $`V`$ such that for any finite dimensional $`𝔨`$-invariant subspace $`FV`$ the action of $`𝔨`$ on $`F`$ exponentiates to $`K`$. The categories $`𝒪`$ and $``$ are $`\mathrm{}`$-categories in the sense of Bass \[B\] page 57.
For any $`\mathrm{}`$-category $`𝒞`$ let
$$=(𝒞)$$
be the set of isomorphism classes of simple objects of $`𝒞`$ \[assume it is a set\], for each $`i`$ choose a representative
$$L_ii$$
and let
$$\mathrm{}(i)$$
be the projective dimension of $`L_i`$ \[i.e. the supremum in $`\mathrm{}\{+\mathrm{}\}`$ of the set $`\{n\mathrm{}|Ext^n(L_i,)0\}`$\].
(1.2) Definition. The $`𝒞`$-ordering is the smallest partial ordering $``$ on $``$ satisfying
$$\begin{array}{c}i,j\\ \\ \mathrm{}(j)=\mathrm{}(i)+1<\mathrm{}\\ \\ Ext^1(L_j,L_i)0\end{array}\}ij.$$
(1.3) Definition. The subcategory generated by the subset $`𝒥`$ of $``$ is the full sub-$`\mathrm{}`$-category $`𝒥_𝒞`$ of $`𝒞`$ characterized by the condition that an object $`V`$ of $`𝒞`$ belongs to $`𝒥_𝒞`$ iff each simple subquotient of $`V`$ is isomorphic to $`L_j`$ for some $`j𝒥`$.
(1.4) Definition. If $`𝒞`$ is a $`\mathrm{}`$-category and $``$ a full sub-$`\mathrm{}`$-category, say that $``$ is Ext-full in $`𝒞`$ if for all $`V,W`$ the natural morphism
$$Ext_{}^{}(V,W)Ext_𝒞^{}(V,W)$$
is an isomorphism.
\[If all objects of $``$ have finite length it suffices to check the above isomorphism for $`V`$ and $`W`$ simple (because of the long exact sequences and the five-lemma).\] Recall that a subset $`𝒥`$ of $``$ is an initial segment iff
$$\begin{array}{c}i,j\\ \\ j𝒥\\ \\ ij\end{array}\}V𝒥.$$
(1.5) Definition. In the above notation $`𝒞`$ is a Guichardet category if the subcategory generated by any initial segment is Ext-full in $`𝒞`$.
Let Setting (1.1) be in force, denote by $`I`$ the annihilator of the trivial module in the center of $`U(𝔤)`$ and for any sub-$`\mathrm{}`$-category $`𝒞`$ of $`𝔤`$-mod let
$$𝒞_\rho $$
be the full sub-$`\mathrm{}`$-category of $`𝒞`$ whose objects are annihilated by some power of $`I`$.
(1.6) Weak Alexandru Conjecture. The category $`_\rho `$ is a Guichardet category.
Part (a) of the theorem below is due to Cline, Parshall and Scott, and part (b) to Fuser.
(1.7) Theorem.
(a) The category $`𝒪_\rho `$ is a Guichardet category,
(b) the Weak Alexandru Conjecture holds for $`_\rho `$ with $`G=SL(3,\mathrm{})`$, $`\text{Spin}(n,1)`$ or $`SU(n,1)`$.
The next item on the agenda is the Strong Alexandru Conjecture (SAC). Hoping to make it more digest I first give a set of statements of a somewhat geometric flavor which imply the SAC for $`\text{Spin}(2n+1,1)`$; more precisely the setting of the SAC is dual to the one described here. I’ll use horizontal lines to set off this part of the text \[which is merely motivational\].
Part B. A detailed example
Let $`\mathrm{\Omega }_k^p`$ be the space of those $`p`$-forms on hyperbolic $`(2n+1)`$-space which are killed by $`\mathrm{\Delta }^k`$ and $`\mathrm{\Omega }^p`$ the union of the $`\mathrm{\Omega }_k^p`$ ; put
$$:=\{p\mathrm{}|0pn\},$$
$$\mathrm{\Omega }:=\underset{p}{}\mathrm{\Omega }^p.$$
As a general notation if $`V`$ is a vector space, $`W`$ a vector subspace and $`\phi _1,\mathrm{},\phi _k`$ endomorphisms of $`V`$, set
$$W(\phi _1,\mathrm{},\phi _k):=W\left(\underset{i=1}{\overset{k}{}}Ker\phi _i\right).$$
Let $`d`$ be the differential and $`d^{}`$ be the codifferential, and put for $`p`$
$$L_p:=\mathrm{\Omega }^p(d,d^{}),$$
$$M_p:=\mathrm{\Omega }^p(d^{}),$$
$$\overline{M}_p:=\{\begin{array}{ccc}\mathrm{\Omega }^p(\mathrm{\Delta },d^{})& \text{ if }& p<n,\\ \\ L_n& \text{ if }& p=n.\end{array}$$
Equip $`\mathrm{\Omega }_k^p`$ with the $`C^0`$-topology and $`\mathrm{\Omega }^p`$ with the inductive limit topology ; let $`G`$ be the group of orientation preserving hyperbolic isometries ; denote by $`𝒞`$ the category defined by the rule that an object of $`𝒞`$ is a topological $`G`$-module $`V`$ which is isomorphic to a close subspace of some $`\mathrm{\Omega }^{p_1}\mathrm{}\mathrm{\Omega }^{p_r}`$, and a morphism in $`𝒞`$ is a $`G`$-equivariant continuous linear map. The following facts are known :
$``$The category $`𝒞`$ is a $`\mathrm{}`$-category.
$``$The modules introduced above belong to it ; any simple object of $`𝒞`$ is isomorphic to $`L_p`$ for a unique $`p`$ in $``$.
$``$The $`𝒞`$-ordering on $``$ is opposite to the natural ordering ; the projective dimension of $`L_p`$ is $`2n+1p`$.
$``$The category $`_\rho `$ is equivalent to the subcategory of $`𝒞`$ whose objects have finite length, or equivalently are annihilated by some power of $`\mathrm{\Delta }`$.
$``$This subcategory — which I abusively denote by $`_\rho `$ for a short while — is Ext-full in $`𝒞`$ and contains $`\overline{M}_p`$.
$``$The objects of $`𝒞`$ have injective hulls and therefore minimal injective resolutions.
\[In this parenthesis I give a reminder of what’s meant here by minimal resolution and offer to the reader unfamiliar with homological algebra a cheap definition of Ext-groups in this context. Given $`V`$ in $`𝒞`$ there are elements $`p_1,\mathrm{},p_k`$ of $``$ and an isomorphism from the socle of $`V`$ onto $`_{i=1}^kL_{p_i};`$ this isomorphism extends to an embedding $`\phi :VI^0:=_{i=1}^k\mathrm{\Omega }^{p_i}`$ (this statement is sometimes called Frobenius reciprocity). Since the cokernel of $`\phi `$ is again in $`𝒞`$ this process can be iterated, giving rise to an injective resolution $`VI^0I^1\mathrm{}`$, which is clearly minimal. Moreover for each pair of integers $`(j,p)`$ with $`j0`$ and $`p`$ there is a finite dimensional vector space $`Ext^j(L_p,V)`$ acted on trivially by $`G`$ such that $`I^j_pExt^j(L_p,V)_{\mathrm{}}\mathrm{\Omega }^p`$.\]
$``$The injective hull of $`L_p`$ is $`\mathrm{\Omega }^p`$.
$``$The vector spaces $`Ext^{}(V,W)`$ are finite dimensional for $`V,W𝒞`$.
Let me digress a tiny bit by stating the Weak Alexandru Conjecture in this setup. During this parenthetical comment $`p`$ shall be a fixed “number” satisfying $`1p\mathrm{}`$ \[not an element of $``$\]. Let $`K`$ be the stabilizer in $`G`$ of a given point of hyperbolic space, say that an object $`V`$ of $`_\rho `$ is $`p`$-integrable if its $`K`$-finite vectors are, and let $`^p`$ be the full subcategory of $`p`$-integrable objects of $`_\rho `$ ; then $`^p`$ is Ext-full in $`_\rho `$ \[more precisely a subcategory of $`_\rho `$ is generated by an initial segment iff it is of the form $`^p`$ — the number $`p`$ being of course in general nonunique.\]
Going back to the SAC, for $`q`$ consider the filtration $`F_0^q:=0F_1^q:=M_qF_2^q:=\mathrm{\Omega }^q`$ ; then $`d^{}`$ induces an isomorphism $`F_2^q/F_1^qM_{q1}`$ \[with the convention $`M_1=0`$\] ; this filtration is analogous to the filtration of projective modules by Verma modules in the category $`𝒪`$ and is encoded in Axiom (1.10) below. If $`E^{}`$ is a graded vector space let $`E^{}(t)`$ be its Poincar series. Define the $``$ by $``$ matrix $`a`$ with entries in $`\mathrm{}[t]`$ by
$$a_{pq}=Ext^{}(L_p,M_q)(t).$$
To compute this series note that the augmented complex
$$\begin{array}{ccccccccc}M_q& \mathit{}& \mathrm{\Omega }^q& \stackrel{d^{}}{}& \mathrm{\Omega }^{q1}& \stackrel{d^{}}{}& \mathrm{}& \stackrel{d^{}}{}& \mathrm{\Omega }^0\end{array}$$
is “the” minimal injective resolution of $`M_p`$, whence
$$a_{pq}=\{\begin{array}{ccc}t^{qp}& \text{ if }& pq,\\ \\ 0& \text{ if }& q<p.\end{array}$$
Letting $`[V]`$ be the class of $`V`$ in the Grothendieck group $`[𝒞]`$ of $`𝒞`$ we have
$$\left[\overline{M}_p\right]=[L_p]+[L_{p+1}]$$
\[with $`L_{n+1}=0`$\] and
$$[L_p]=\underset{q=p}{\overset{n}{}}(1)^{qp}\left[\overline{M}_q\right]=\underset{q}{}a_{pq}(1)\left[\overline{M}_q\right];$$
this corresponds to the Delorme formula in the category $`𝒪`$ and is encapsulated in Axiom (1.14). The inverse $`a^1=(a_{pq}^1)`$ of $`a`$ being given by
$$a_{pq}^1=\{\begin{array}{cc}(t)^{qp}& \text{if }pqp+1,\\ \\ 0& \text{otherwise}\end{array}$$
the number of occurrences of $`M_p`$ in the filtration $`F_{}^q`$ is $`a_{pq}^1(1)`$, which is an analog of the BGG duality in the category $`𝒪`$, and gives rise to Axiom (1.15) below. Let $`\overline{M}_{p,\text{soc}}^{}`$ the graded object of $`𝒞`$ associated to the socle filtration of $`\overline{M}_q`$ and $`\left[\overline{M}_{q,\text{soc}}^{}\right](t)`$ be its image in $`\mathrm{}[t]_{\mathrm{}}[𝒞]`$, then
$$\left[\overline{M}_{p,\text{soc}}^{}\right](t)=[L_p]+t[L_{p+1}]=\underset{q}{}a_{pq}^1[L_q].$$
The corresponding theorem in the realm of the category $`𝒪`$ is due to Beilinson, Ginzburg and Soergel, and the above formula suggests Axiom (1.16) below. Finally let’s compute $`Ext^{}(L_p,L_q)`$. To simplify the notation, if $`n<p2n+1`$ identify $`\mathrm{\Omega }^p`$ to $`\mathrm{\Omega }^{2n+1p}`$ via the star operator, and decree that $`\mathrm{\Omega }^p=0`$ for $`p<0`$ or $`p>2n+1`$. “The” minimal resolution of $`L_q`$ being
$$\begin{array}{ccccccccc}L_q& & \mathrm{\Omega }^q& \stackrel{\left(\genfrac{}{}{0pt}{}{d}{d^{}}\right)}{}& \mathrm{\Omega }^{q+1}\mathrm{\Omega }^{q1}& \stackrel{\left(\genfrac{}{}{0pt}{}{d0}{0d^{}}\right)}{}& \mathrm{\Omega }^{q+2}\mathrm{\Omega }^{q2}& \stackrel{\left(\genfrac{}{}{0pt}{}{d0}{0d^{}}\right)}{}& \mathrm{}\end{array}$$
we have
$$Ext^{}(L_p,L_q)(t)=t^{|qp|}+t^{2n+1pq}.$$
$`()`$
Setting $`L:=_pL_p`$ and letting $`{}_{}{}^{t}a`$ be the transpose of $`a`$ and $`\delta `$ the diagonal matrix defined by
$$\delta _p=\{\begin{array}{cc}1+t& \text{if }p=n,\\ \\ 1t^2& \text{if }pn,\end{array}$$
$`()`$ reads
$$Ext^{}(L,L)(t)=a\delta {}_{}{}^{t}a.$$
The category $`𝒪`$ analog is the Beilinson-Ginzburg formula and the corresponding Axiom below is (1.18).
Part C. The Strong Alexandru Conjecture
Here are some preliminaries to state the Strong Alexandru Conjecture. Set
$$Z:=\mathrm{}[[z_1,\mathrm{},z_m]],$$
\[where $`z_1,\mathrm{},z_m`$ are indeterminates\] and let $`A`$ be a $`Z`$-algebra which is finitely generated over $`Z`$. Then there is a semisimple subalgebra $`A_0`$ of $`A`$ satisfying $`A=A_0\text{rad}(A)`$. Assume there is a finite set $`F`$ such that $`A_0`$ can — and will — be identified to the algebra $`\mathrm{}^F`$ of functions on $`F`$. For each $`iF`$ define $`e_iA_0`$ by $`e_i(j)=\delta _{ij}`$ \[Kronecker delta\]. As a general notation put
$$A\text{-fd}:=\text{ the category of finite dimensional }A\text{-modules}.$$
I’ll make free use of the facts that by a theorem of Casselman $`A`$-fd is Ext-full in $`A`$-mod and that by a result of BGG the categories $`𝒪_\rho `$ and $`_\rho `$ are equivalent to $`A`$-fd for some algebra $`A`$ as above — the $`\mathrm{}`$-algebra isomorphism class of $`A`$ being unique. It will be tacitly assumed that the category $`𝒞`$ of interest has been set to be $`𝒪_\rho `$ or $`_\rho `$, and that an algebra $`A`$ as above and an equivalence $`𝒞A`$-fd have been chosen, providing in particular an identification $`=F`$ ; the symbol $`L_i`$ denotes at the same time an object of $`𝒞`$ and “the” corresponding object in $`A`$-fd ; more generally I’ll allow myself to navigate rather freely between $`𝒞`$ and $`A`$-fd. Fix a family $`M=(M_i)_i`$ of $`A`$-modules.
(1.8) Definition. An $`M`$-filtration $`F_{}`$ of an $`A`$-module $`V`$ is a sequence
$$0=F_0F_1F_2\mathrm{}F_r=V$$
of submodules such that there exists an $`r`$-tuple $`(i_1,\mathrm{},i_r)`$ of elements of $``$ satisfying $`F_p/F_{p1}M_{i_p}`$ for $`1pr`$. An $`A`$-module is $`M`$-filtrable if it admits an $`M`$-filtration.
Let $``$ be the $`𝒞`$-ordering on $``$ \[see Definition (1.2)\] ; form the small Verma module
$$M_i^{}(A):=Ae_i/\underset{ji}{}Ae_jAe_i$$
and the large Verma module
$$M_i^+(A):=Ae_i/\underset{j>i}{}Ae_jAe_i;$$
letting $`M_i`$ be as in Definition (1.8) set
$$\overline{M}_i:=M_i/\text{rad}(End_AM_i)M_i,$$
$$E^{}(t):=\text{Poincar series of the graded vector space }E^{}\text{,}$$
$$[V]:=\text{class of }V\text{ in the Grothendieck group.}$$
Let $`a`$ be the $``$ by $``$ matrix with entries in $`\mathrm{}[[t]]`$ defined by
$$a_{ij}(t):=Ext_A^{}(M_i,L_j)(t)$$
and if $`a`$ happens to be invertible let $`a_{ij}^1`$ be the $`(i,j)`$ entry of $`a^1`$. Let $`V`$ be in $`A`$-fd ; as a general notation set
$$V_{\text{rad}}^{}:=\left(\text{rad}(A)^iV/\text{rad}(A)^{i+1}V\right);$$
suppose $`V0`$ ; let $`S_1:=0S_1=`$ soc $`V\mathrm{}S_p=V`$ be the socle filtration of $`V`$, with $`S_{p1}V`$ ; say that the radical and socle filtrations coincide if
$$\text{rad}(A)^iV=S_{pi}i.$$
Let $`{}_{}{}^{t}b`$ denote the transpose of any matrix $`b`$ ; and consider the following conditions
(1.9) for each $`i`$ we have $`M_i^+(A)=M_i^{}(A)=M_i`$;
(1.10) for each $`i`$ the module $`Ae_i`$ is $`M`$-filtrable \[see Definition (1.8)\] ;
(1.11) for each $`i`$ in $``$ the module $`M_i`$ is flat over $`End_AM_i;`$
(1.12) for any $`i`$ there is a $`ji`$ such that
$$\mathrm{}(i)=sup\{n\mathrm{}|Ext^n(L_i,L_j)0\};$$
(1.13) there are polynomials $`p_{ij}`$ such that
$$a_{ij}(t)=t^{\mathrm{}(j)\mathrm{}(i)}p_{ij}(t^2),$$
$$p_{ij}0ijp_{ij}(0)=1,$$
$$p_{ii}=1,$$
$$degP_{x,y}<\frac{\mathrm{}(y)\mathrm{}(x)}{2}\text{ if }x<y;$$
in particular $`a`$ is invertible ;
(1.14) $`[L_j]=_ia_{ij}(1)\left[\overline{M}_i\right]`$ ;
(1.15) things can be arranged so that the number of occurrences of $`M_j`$ in the $`M`$-filtration of $`Ae_i`$ in (1.10) is $`a_{ij}^1(1)`$ ;
(1.16) the radical and socle filtrations of $`\overline{M}_j`$ coincide and we have
$$\left(e_i\overline{M}_{j,\text{rad}}^{}\right)(t)=a_{ij}^1(t);$$
(1.17) $`End_A\left(\overline{M}_i\right)=\mathrm{}`$ ≥;
(1.18) there is a diagonal $``$ by $``$ matrix $`d`$ such that
$$Ext_A^{}(A_0,A_0)(t)={}_{}{}^{t}ada.$$
The proposition below is essentially due to Cline, Parshall and Scott ; an elementary proof is given in section 1 \[recall that $`A`$-fd is the category of finite dimensional $`A`$-modules\].
(1.19) Proposition. In the above setting if Conditions (1.9)—(1.12) are satisfied then $`A`$-fd is a Guichardet category.
(1.20) Definitions. If $`A`$ satisfies Conditions (1.9)—(1.18) above, then $`A`$ is a BGG algebra. A beegeegee is a category which is equivalent to $`A`$-fd for some BBG algebra $`A`$.
The theorem below is due to BGG (see \[BGG\]), Beilinson and Ginzburg (see \[BGS\]) and Cline-Parshall-Scott (see statements (3.3.c), (3.5.a) and (3.9.a) in \[CPS\]).
(1.21) Theorem. The category $`𝒪_\rho `$ is a beegeegee.
(1.22) Strong Alexandru Conjecture. The category $`_\rho `$ is a beegeegee.
(1.23) Theorem (Fuser). The above conjecture holds for $`\text{Spin}(n,1),SU(n,1)`$ and $`SL(3,\mathrm{})`$.
The drawback \[at least one of them\] of all this stuff is that it’s almost never computable ! Here is a statement which, although as conjectural as the previous ones, can be submitted to numerical tests. It consists in a computable variant of \[a particular case of\] Condition (1.18), that is in a formula which would express, when $`𝔤`$ and $`𝔨`$ have the same rank, the Poincar series $`Ext_A^{}(A_0,A_0)(t)`$ in terms of computable things. In the Langlands classification $`L_i`$ occurs as the unique simple quotient of a module induced from some parabolic subgroup $`P_i`$ ; let $`𝔭_i=𝔪_i𝔞_i𝔫_i`$ be a Langlands decomposition of $`\mathrm{}_{\mathrm{}}\text{Lie}(P_i)`$ ; put
$$\stackrel{~}{d}_i:=\left(1t^2\right)^{dim𝔞_i};$$
let $`\stackrel{~}{\mathrm{}}(i)`$ be the dimension of the $`K_{\mathrm{}}`$-orbit attached to $`i`$ and $`(\stackrel{~}{p}_{ij})`$ the family of Kazhdan-Lusztig-Vogan polynomials, that is the one denoted $`(P_{\gamma ,\delta })`$ in \[V3\], section 6 ; and set
$$\stackrel{~}{a}_{ij}(t)=t^{\stackrel{~}{\mathrm{}}(j)\stackrel{~}{\mathrm{}}(i)}\stackrel{~}{p}_{ij}(t^2).$$
(1.24) Conjecture. If $`𝔤`$ and $`𝔨`$ have the same rank then
$$Ext_A^{}(A_0,A_0)(t)={}_{}{}^{t}\stackrel{~}{a}\stackrel{~}{d}\stackrel{~}{a}.$$
For a given group $`G`$ one can run the following numerical test.
(a) Compute $`{}_{}{}^{t}\stackrel{~}{a}\stackrel{~}{d}\stackrel{~}{a}`$ ;
(b) look if this is compatible with what’s known of $`Ext_A^{}(A_0,A_0)(t)`$ \[in particular with the $`(𝔤,K)`$-cohomology, as computed by Vogan’s $`U_\alpha `$-algorithm\] ;
(c) pretend Conjecture (1.24) holds and use it to compute the $`_\rho `$-ordering ;
(d) check if there are matrices $`a`$ and $`d`$ such that $`a`$ is upper triangular with ones on the diagonal, $`d`$ is diagonal and we have $`{}_{}{}^{t}ada={}_{}{}^{t}\stackrel{~}{a}\stackrel{~}{d}\stackrel{~}{a}`$.
If the test is successful no conclusion can be drawn ; if it fails then at least one of the involved conjectures is wrong. — In the real rank one case evidence suggests that the classical objects “with tildes” coincide with the nonclassical ones \[“without tildes”\] ; in the case of $`PSp(2,\mathrm{})`$ using results of Vogan \[V\], pp 251-255, one sees such is not the case — but the test is still successful. I think that in the case where the complex ranks of $`𝔤`$ and $`𝔨`$ are different there is a similar formula with $`\stackrel{~}{a}`$ as above and $`\stackrel{~}{d}`$ a certain polynomial valued diagonal matrix.
1. ≥Proof ≥of ≥Proposition ≥(1.19)
In the setting of Proposition (1.19) assume there are at least two elements in $``$ \[otherwise there is nothing to prove\], and consider the following setting
$$i\text{ is a maximal element of },$$
$$e:=e_i,$$
$$I:=AeA,$$
$$B:=A/I,$$
$$𝒥:=\backslash \{i\}.$$
Proposition (1.19) follows from Lemma (2.1) below, which will be proved at the end of the section.
(2.1) Lemma.
(a) The category $`B`$-mod is Ext-full in $`A`$-mod,
(b) The $`𝒞`$-ordering coincides on $`𝒥`$ with the $`𝒥_𝒞`$-ordering \[see (1.2) and (1.3)\],
(c) $`B`$ satisfies Conditions (1.9)—(1.12).
(2.2) Lemma. If $`ji`$ then $`M_jB`$-fd.
Proof of Lemma (2.2). The statement is an immediate consequence of the following observation. For any pair $`(𝒦,k)`$ with $`j𝒦`$ put
$$P(𝒦,k):=Ae_k/\underset{\mathrm{}𝒦}{}Ae_{\mathrm{}}Ae_k.$$
Then $`P(𝒦,k)`$ is a projective cover of $`L_k`$ in $`𝒦_{A\text{-fd.}}`$ If $`k𝒦`$ and $`P(𝒦,k)`$ is in $`_{A\text{-fd}}`$ then $`P(𝒦,k)=P(,\mathrm{})`$. QED
The following lemma is obvious.
(2.3) Lemma. The three conditions (a) $`j`$ is maximal,
(b) $`M_j=Ae_j`$,
(c) $`M_j`$ is projective
are equivalent. In particular $`M_i=Ae_i`$ is projective. QED
Let $``$ be the class of those $`A`$-modules which are isomorphic to a direct sum of finitely many copies of $`Ae_i`$.
(2.4) Lemma. Let $`V`$ be an $`M`$-filtrable $`A`$-module. Then there is an $`M`$-filtration $`F_{}`$ of $`V`$ and a nonnegative integer $`n`$ such that $`F_n`$ and $`e(V/F_n)=0`$.
Proof. The statement results from the fact easy to check that $`eM_j=0`$ for all $`ji`$ and from the observation that each occurrence of a projective into a given module as a subquotient is in fact an occurrence as a submodule. QED
In particular the above lemma provides for each $`j`$ an $`M`$-filtration $`F_{}^j`$ of $`Ae_j`$ by left subideals and a nonnegative integer $`n_j`$ such that $`F_{n_j}^j`$ and $`e(Ae_j/F_{n_j}^j)=0`$. Put
$$J:=\underset{j}{}F_{n_j}^j;$$
then there is a positive integer $`n`$ such that $`J`$ is isomorphic to the direct sum of $`n`$ copies of $`Ae_i`$ ; write this direct sum in the nonsensical form $`_{j=1}^nAe_i`$ ; choose an isomorphism
$$\phi :\underset{j=1}{\overset{n}{}}Ae_i\stackrel{}{}J;$$
and set
$$A_i:=End_AAe,$$
$$C:=eAe=(A_i)^{op}.$$
We may — and will — view any left module over $`A_i`$ as a right module over $`C`$. Put
$$:=\underset{C}{}.$$
(2.5) Lemma.
(a) For $`ji`$ the multiplication map $`AeeAe_jAe_j`$ is an isomorphism,
(b) the natural map $`AeeAA`$ is one-to-one,
(c) we have $`J=I`$ and $`F_{n_j}^j=Ie_j`$. In particular $`Be_j:=Ae_j/Ie_j=Ae_j/F_{n_j}^j`$ is $`M`$-filtrable.
Proof. To prove (a) set $`f:=ee_j`$ and note that the canonical isomorphism $`AeeAeAe`$ is the direct sum of the multiplication maps $`AeeAe_jAe_j`$ and $`AeeAfAf`$. To check (b) consider the diagram
$$\begin{array}{ccc}AeeA& \stackrel{\alpha _1}{}& A\\ \alpha _2& & \phi \\ AeeJ& & \underset{j=1}{\overset{n}{}}Ae_i\\ \alpha _3& & \alpha _5\\ Ae\underset{j=1}{\overset{n}{}}eAe_i& \underset{\alpha _4}{}& \underset{j=1}{\overset{n}{}}AeAe_i,\end{array}$$
where $`\alpha _1,\mathrm{},\alpha _5`$ are defined as follows
$`\alpha _1`$ is the multiplication map,
$`\alpha _2`$ is the identity \[indeed $`e(A/J)=0eA=eJ`$\],
$`\alpha _3`$ is $`1\phi ^1`$,
$`\alpha _4`$ is the canonical map,
$`\alpha _5`$ is the multiplication map.
The maps $`\alpha _2,\alpha _3`$ and $`\alpha _4`$ are clearly isomorphisms and by (a) so is $`\alpha _5`$, whereas $`\phi `$ is one-to-one. Since the diagram commutes $`\alpha _1`$ is one-to-one and (b) is verified. To prove (c) observe $`I=Im\alpha _1=Im\phi =J`$. QED
(2.6) Lemma. Let $`A`$ be an algebra and $`e`$ an element of $`A`$ satisfying $`e^2=e1`$. Put $`I:=AeA`$, $`B:=A/I`$ and $`C:=eAe`$, let $`V`$ and $`W`$ be $`B`$-modules, and consider the following conditions
(a) the natural map $`Ae_CeAIA`$ is one-to-one and $`Ae`$ is right $`C`$-flat,
(b) $`I`$ is right $`A`$-flat,
(c) $`Tor_q^A(B,V)=0`$ for all $`q>0`$,
(d) $`Ext_B^n(V,W)Ext_A^n(V,W)`$ for all $`n`$ (natural isomorphism). Then (a) $``$ (b) $``$ (c) $``$ (d).
Proof. The implication (a) $``$ (b) is clear. The implication (b) $``$ (c) follows from the long (?) exact sequence obtained by applying “$`_AV`$” to the short exact sequence $`IAB`$. The implication (c) $``$ (d) follows from Proposition VI.4.1.3 of Cartan-Eilenberg \[CE\]. QED
Proof of Lemma (2.1). Part (a) \[the Ext-fullness of $`B`$-mod in $`A`$-mod\] follows from Lemma (2.6). In view of Condition (1.12) \[about projective dimensions\] part (b) \[the coincidence of the orderings\] results from (a). Let me prove part (c), claiming that $`B`$ satisfies Conditions (1.9)—(1.12). Condition (1.9) \[the coincidence of the Verma modules\] and Condition (1.11) \[the $`(End_BM_i)`$-flatness of $`M_i`$\] are consequences of part (b) and Lemma (2.2) \[$`M_jB`$-mod\] ; Condition (1.10) \[the $`M`$-filtrability\] follows from (2.5.c) ; Condition (1.12) results from (a). QED
* ≥\* ≥\*
\[B\] Bass H., Algebraic K-theory, Benjamin, New York 1968.
\[BGG\] Bernstein I.N., Gelfand I.M., Gelfand S.I., Category of g-modules, Funct. Anal. Appl. 10, 87-92 (1976).
\[BGS\] Beilinson A., Ginzburg V., Soergel W., Koszul duality patterns in representation theory, J. Am. Math. Soc. 9, No.2, 473-527 (1996).
\[CE\] Cartan H., Eilenberg S., Homological algebra, Princeton University Press, 1956.
\[CPS\] Cline E., Parshall B., Scott L., Finite dimensional algebras and highest weight categories, J. Reine Angew. Math. 391, 85-99 (1988).
\[D\] Delorme P., Extensions dans la catégorie $`𝒪`$ de Bernstein-Gelfand-Gelfand. Applications, manuscript, October 1978.
\[V\] Vogan D., The Kazhdan-Lusztig Conjecture for real reductive groups, in Representation theory of reductive groups, Proceedings of the University of Utah conference 1982, Trombi Peter C. (Ed.), Birkh user, 1983, Progress in mathematics Number 40.
\[V3\] Vogan, D., Irreducible characters of semisimple Lie groups. III : Proof of Kazhdan-Lusztig conjecture in the integral case, Invent. Math. 71 (1983) 381-417.
This text and others are available at http://www.iecn.u-nancy.fr/$``$gaillard
|
warning/0003/math-ph0003023.html
|
ar5iv
|
text
|
# Lifschitz tail in a magnetic field: coexistence of classical and quantum behavior in the borderline case.
## 1 Introduction
The magnetic Lifschitz tail is the asymptotic behavior of the integrated density of states (IDS), $`N(E)`$, at energy $`E`$ near the bottom of the spectrum of the two dimensional random Schrödinger operator with a constant magnetic field $`B`$. The random potential, $`V_\omega `$, represents repulsive impurities that are modelled by a single-site potential profile $`V^{(0)}0`$ convolved with a homogeneous Poisson point process.
The low energy asymptotics of the IDS exhibits two qualitatively different behaviors. For long range $`V^{(0)}`$, the asymptotics of $`N(E)`$ is solely determined by the potential, i.e., by classical effects, hence it is called classical asymptotics or classical regime. In this regime, the low energy behavior of $`N(E)`$ is sensitive to the details of the tail of $`V^{(0)}`$ and it is insensitive to the strength of the magnetic field. For short range potentials, the asymptotics of $`N(E)`$ is determined by the quantum kinetic energy (quantum asymptotics or quantum regime) and it is universal; it depends only on the strength of the magnetic field, but it is insensitive to the potential profile $`V^{(0)}`$.
It has been established in that potentials $`V^{(0)}`$ with algebraically decaying tail of any finite order belong to the classical regime if $`B0`$. The strength of the magnetic field does not appear in the leading term asymptotics of $`N(E)`$. This result is in contrast to the nonmagnetic case, where it has been shown (, , ) that an algebraic decay, $`V^{(0)}(x)|x|^{(d+2)}`$, discriminates between the classical and quantum regimes in $`d`$ dimensions.
Nevertheless, quantum regime does appear in the magnetic case as well, but the discriminating potential decay is much faster than algebraic; in fact it is Gaussian. The existence of the quantum regime for compactly supported potentials was proven in and the Gaussian threshold was conjectured. This threshold has been verified in relying on for the most involved technical part. More precisely, it has been proven that stretched-Gaussian decay leads to the classical asymptotics, while super-Gaussian decay leads to the quantum asymptotics. Later, some refined results were obtained in .
The borderline case, when $`V^{(0)}`$ is asymptotically Gaussian, has not been settled conclusively, only two-sided estimates were given in . The lower bound indicates a coexistence of the classical and quantum effects, as it is determined by $`2\mathrm{}_B^2+\lambda ^2`$. Here $`\mathrm{}_B:=B^{1/2}`$ is the magnetic lengthscale representing the kinetic energy contribution, and $`\lambda `$ is the lengthscale of the Gaussian potential. The upper bound is determined by $`\mathrm{max}\{2\mathrm{}_B^2,\lambda ^2\}`$, indicating no coexistence of the two regimes. The conjecture of was that the lower bound is the true asymptotics.
The purpose of this paper is to show this conjecture. We emphasize that in the Gaussian borderline case both classical and quantum effects are important, hence none of them can be neglected along the proof. This is the main novelty of the present paper, which is a extension of our earlier work .
### 1.1 Definitions
We consider a nonnegative potential function
$$V^{(0)}L_{loc}^2(𝐑^2),V^{(0)}0,$$
(1)
that is strictly positive on a non-empty open set, i.e.,
$$V^{(0)}(x)v\mathrm{𝟏}(|xx_0|a)$$
(2)
for some $`v,a>0`$ and $`x_0𝐑^2`$. Here $`\mathrm{𝟏}()`$ denotes the characteristic function. Let
$$V(x)=V_\omega (x):=\underset{i}{}V^{(0)}(xx_i(\omega ))$$
(3)
be a random potential, where $`x_i(\omega )`$ is the realization of the Poisson point process on $`𝐑^2`$ with a constant intensity $`\nu `$ (here $`\omega `$ refers to the randomness, but we shall usually omit it from the notations). The expectation with respect to this process is denoted by $``$.
We consider the following magnetic Schrödinger operator with a random potential $`V_\omega `$
$$H(B,V_\omega )=H_\omega =\frac{1}{2}\left[(iA)^2B\right]+V_\omega \text{on}L^2(𝐑^2),$$
(4)
where $`A:𝐑^2𝐑^2`$ is a deterministic vector potential (gauge) generating the constant $`B>0`$ magnetic field, i.e., $`\text{curl}A=B`$. The properties we are interested in are independent of the actual gauge choice, so, conveniently, we choose the standard gauge $`A(x):=\frac{B}{2}\left(\genfrac{}{}{0pt}{}{x_2}{x_1}\right)`$. Here $`x=(x_1,x_2)𝐑^2`$. We subtracted the constant $`B/2`$ term in the kinetic energy both for physical reasons (spin coupling) and for mathematical convenience. The spectrum of the free operator $`H(B,V0)`$ is $`\{nB:n=0,1,2,\mathrm{}\}`$.
We also define $`H_{Q,\omega }=H_Q(B,V_\omega )`$ as the restriction of $`H_\omega `$ onto a domain $`Q𝐑^2`$ (with Dirichlet boundary conditions). In this paper, by domain we mean an open, bounded subset of $`𝐑^2`$ with regular (piecewise $`C^1`$) boundary, which is not necessarily connected.
We shall always assume that $`V^{(0)}`$ has sufficient decay so that $`V_\omega L_{loc}^2`$ with probability one, i.e., these operators are almost surely selfadjoint. Moreover, in all cases we consider it is easy to show that
$$inf\text{Spec}H(B,V_\omega )=0\text{almost surely.}$$
(5)
We define the integrated density of states (IDS) as
$$N(E):=\underset{Q𝐑^2}{lim}\frac{1}{|Q|}\text{Tr}P_E(H_{Q,\omega }),$$
(6)
where $`P_E`$ is the spectral projection onto the half line $`(\mathrm{},E]`$, and $`Q𝐑^2`$ is an increasing sequence of nested regular domains, say, squares or disks. The trace is over $`L^2(Q)`$. For the existence of this limit and equivalent definitions we refer to , and references therein.
Following , we assume that $`V^{(0)}`$ has one of the following behaviors at infinity:
Sub-Gaussian decay:
$$\underset{|x|\mathrm{}}{lim}\frac{|x|^2}{\mathrm{log}V^{(0)}(x)}=\mathrm{};$$
Gaussian decay:
$$\underset{|x|\mathrm{}}{lim}\frac{\mathrm{log}V^{(0)}(x)}{|x|^2}=\frac{1}{\lambda ^2}$$
(7)
for some $`0<\lambda <\mathrm{}`$;
Super-Gaussian decay:
$$\underset{|x|\mathrm{}}{lim}\frac{\mathrm{log}V^{(0)}(x)}{|x|^2}=\mathrm{}.$$
The sub-Gaussian decay leads to the classical regime where the potential determines the Lifschitz tail. Hence precise results require more definite tail behavior in this case. The following definition is taken from .
Regular $`(F,\alpha )`$-decay:
$$\underset{|x|\mathrm{}}{lim}\frac{F(1/V^{(0)}(x))}{|x|}=1$$
for some positive function $`F`$, which is regularly varying of index $`1/\alpha [0,\mathrm{}]`$ and is strictly increasing towards infinity. Recall that a positive measurable function $`F`$ is said to be regularly varying of index $`\gamma `$ if $`lim_t\mathrm{}F(ct)/F(t)=c^\gamma `$ for all $`c>0`$. Such class of functions is denoted by $`R_\gamma `$. Two important cases are:
Algebraic decay: $`lim_{|x|\mathrm{}}|x|^\alpha V^{(0)}(x)=\mu `$ with some exponent $`\alpha >2`$ and constant $`0<\mu <\mathrm{}`$. This corresponds to $`F(t)(\mu t)^{1/\alpha }`$.
Stretched Gaussian decay: $`lim_{|x|\mathrm{}}|x|^\alpha \mathrm{log}V^{(0)}(x)=\lambda ^\alpha `$ for some $`0<\lambda <\mathrm{}`$ and $`0<\alpha <2`$. This corresponds to $`F(t)\lambda (\mathrm{log}t)^{1/\alpha }`$.
### 1.2 Results
The result of for algebraically decaying potential $`V^{(0)}`$ is
$$\underset{E0}{lim}E^{2/(\alpha 2)}\mathrm{log}N(E)=C(\alpha ,\mu ,\nu )$$
with an explicitly computed constant $`C(\alpha ,\mu ,\nu )`$.
The general regular $`(F,\alpha )`$-decaying sub-Gaussian case was discussed in details in . The Lifschitz tail is given by the de Bruijn conjugate $`f^\mathrm{\#}`$ of the function $`tf(t)=[t^{1/\alpha }F(t)]^{2\alpha /(2\alpha )}`$. Recall that the de Bruijn conjugate of a slowly varying function $`fR_0`$ is $`f^\mathrm{\#}R_0`$ such that $`f(t)f^\mathrm{\#}(tf(t))1`$ and $`f^\mathrm{\#}(t)f(tf^\mathrm{\#}(t))1`$ as $`t\mathrm{}`$. With this definition
$$\underset{E0}{lim}\frac{E^{2/(\alpha 2)}\mathrm{log}N(E)}{f^\mathrm{\#}(E^{\alpha /(2\alpha )})}=C(\alpha ,\nu )$$
with an explicit constant. In particular, for stretched-Gaussian potential $`V^{(0)}`$ this asymptotics is explicitly given as ()
$$\underset{E0}{lim}\frac{\mathrm{log}N(E)}{|\mathrm{log}E|^{2/\alpha }}=\pi \nu \lambda ^2.$$
For the super-Gaussian case it is proven in that
$$\underset{E0}{lim}\frac{\mathrm{log}N(E)}{|\mathrm{log}E|}=2\pi \nu \mathrm{}_B^2=\frac{2\pi \nu }{B},$$
(8)
with the additional assumption (2). In particular, the super-Gaussian decay includes compactly supported potentials; the case for which (8) was proven in . In a slightly more general definition of super-Gaussian decay was introduced:
$$\underset{R>0}{inf}\text{ess}\underset{|x|>R}{sup}\frac{\mathrm{log}V^{(0)}(x)}{|x|^2}=\mathrm{},$$
(9)
and (8) was proven for such potentials (in addition to the condition (2)).
Finally, the following estimates were given in for the Gaussian case (7)
$$\pi \nu (\lambda ^2+2\mathrm{}_B^2)\underset{E0}{lim\; inf}\frac{\mathrm{log}N(E)}{|\mathrm{log}E|}\underset{E0}{lim\; sup}\frac{\mathrm{log}N(E)}{|\mathrm{log}E|}\pi \nu \mathrm{max}\{\lambda ^2,2\mathrm{}_B^2\}$$
(10)
and the upper bound was weakened to $`2\pi \nu \mathrm{}_B^2`$ in if a more general definition of Gaussian decay is used that is analogous to (9). Our goal is to prove that the lower bound in (10) is the correct one as conjectured in .
###### Theorem 1.1
Suppose that $`V^{(0)}`$ satisfies (1) and (7). Then
$$\underset{E0}{lim}\frac{\mathrm{log}N(E)}{|\mathrm{log}E|}=\pi \nu (\lambda ^2+2\mathrm{}_B^2),\mathrm{}_B:=B^{1/2}.$$
(11)
Since the lower bound (10) has been proven in , we focus only on the upper bound. As usual, we define the Laplace transform of $`N(E)`$ as
$$L(t):=_0^{\mathrm{}}e^{Et}𝑑N(E)=e^{tH_\omega }(x,x).$$
Recall that the diagonal element of the averaged heat kernel is independent of $`x`$. For more details, see . Using a standard Tauberian argument (see for example Appendix of for details), the upper bound in (11) follows from
$$\underset{t\mathrm{}}{lim\; sup}\frac{\mathrm{log}L(t)}{\mathrm{log}t}\pi \nu (\lambda ^2+2\mathrm{}_B^2).$$
(12)
In the rest of the paper we prove (12). Several steps will be similar to , these will not be repeated in details. We give detailed proofs only for the new parts of the argument.
## 2 Localization
We use a two-step localization as in . The first localization is identical to the upper bound in Proposition 3.2 of and the proof is the same.
###### Proposition 2.1
Let $`M:=[m,m]^2`$ be a square box, then
$$L(t)\underset{m\mathrm{}}{lim\; inf}\frac{1}{|M|}\text{Tr}e^{tH_{M,\omega }}.\mathrm{}$$
For the second, more refined localization we cannot neglect the tail of the impurity potentials. We will define effective boundary potentials that estimate the potential tails inside a domain $`\mathrm{\Omega }`$ that come from impurities located outside of $`\mathrm{\Omega }`$.
To prove (12), it is enough to show that
$$\underset{t\mathrm{}}{lim\; sup}\frac{\mathrm{log}L(t)}{\mathrm{log}t}\pi \nu (L^2+2\mathrm{}_B^2)$$
(13)
for any $`L<\lambda `$. We fix two numbers, $`0<L<\overline{L}<\lambda `$, for the rest of the proof and we omit the dependence on $`L`$ and $`\overline{L}`$ of various quantities in the notation.
Using (7), there exists $`R1`$ such that
$$V^{(0)}(x)e^{|x|^2/\overline{L}^2}\text{for all}|x|R.$$
(14)
We also choose $`R=R(\overline{L},B)`$ so large that $`e^{(2R)^2/\overline{L}^2}B`$. For any domain $`\mathrm{\Omega }`$ we define the following boundary potentials ($`\mathrm{\Omega }`$ stands for the boundary of $`\mathrm{\Omega }`$):
$$\overline{V}_\mathrm{\Omega }(x):=\mathrm{exp}\left[(\text{dist}(x,\mathrm{\Omega }))^2/\overline{L}^2\right]\mathrm{𝟏}(x\mathrm{\Omega })\mathrm{𝟏}(\text{dist}(x,\mathrm{\Omega })R),$$
(15)
$$V_\mathrm{\Omega }(x):=\mathrm{exp}\left[(\text{dist}(x,\mathrm{\Omega }))^2/L^2\right]\mathrm{𝟏}(x\mathrm{\Omega }).$$
(16)
Similarly to Section 6 of we fix parameters $`0<\beta <B/2`$, $`1sm`$ and let $`M:=[m,m]^2`$, $`\stackrel{~}{M}:=[ms,m+s]^2`$, $`S:=[s,s]^2`$, $`\stackrel{~}{S}:=[\frac{s}{2},\frac{s}{2}]^2`$ and $`Q_z:=Q+z`$ for any square $`Q𝐑^2`$ and $`z𝐑^2`$. Finally, let $`\lambda _{S_z,\omega }^{(B+2\beta )}`$ be the lowest eigenvalue of
$$\overline{H}_{S_z,\omega }^{(B+2\beta )}:=\frac{1}{2}\left\{\left[i\frac{B+2\beta }{2}\left(\begin{array}{c}x_2\\ x_1\end{array}\right)\right]^2(B+2\beta )\right\}+V_\omega +\overline{V}_{S_z}$$
with Dirichlet boundary conditions on $`S_z`$. The magnetic field of $`\overline{H}_{S_z,\omega }^{(B+2\beta )}`$ is $`B+2\beta `$. Notice that this operator differs from its counterpart in Section 6 of by the additional boundary potential $`\overline{V}_{S_z}`$. We have
###### Proposition 2.2
Assume that $`\beta <1/(2\overline{L}^2)`$, $`\beta s^2128`$ and $`s4R`$. For any $`z𝐑^2`$ there exists a function $`\eta _z`$ supported on $`S_z`$ such that for any $`fH_0^1(M)`$
$$f,H_{M,\omega }f\frac{\beta }{2\pi }_{\stackrel{~}{M}}dzf\eta _z,\overline{H}_{S_z,\omega }^{(B+2\beta )}f\eta _z65s^2e^{\beta s^2/8}f_{L^2(M)}.$$
Using this result, we obtain the following theorem from Proposition 2.1 exactly as Theorem 6.3 was proven in :
###### Proposition 2.3
Let $`\mathrm{}(t):=10(\mathrm{log}t/B)^{1/2}`$, $`s:=n_0\mathrm{}(t)`$ and $`S=[s,s]^2`$. For any fixed $`0<\beta <1/(2\overline{L}^2)`$ and $`n_0(B/\beta )^{1/2}`$, $`n_0𝐙`$
$$\underset{t\mathrm{}}{lim\; sup}\frac{\mathrm{log}L(t)}{\mathrm{log}t}\underset{t\mathrm{}}{lim\; sup}(\mathrm{log}t)^1\mathrm{log}\mathrm{exp}(t\lambda _{S,\omega }^{(B+2\beta )}).\mathrm{}$$
(17)
Proof of Proposition 2.2. Similarly to the proof of Proposition 6.1 in , we define
$$\phi _z(x):=e^{\beta (xz)^2/2}e^{i\beta (x_2z_1x_1z_2)}$$
and $`T_\beta :=i_1+_2+(B/2+\beta )x_2i(B/2+\beta )x_1`$. We use the following identity to localize the kinetic energy for $`fH_0^1(M)`$
$$f,H_{M,\omega }f=\frac{\beta }{\pi }dz\left\{\frac{1}{2}|T_\beta (\phi _zf)|^2+V_\omega |\phi _zf|^2\right\},$$
where we let $``$ denote $`_{𝐑^2}`$. This magnetic localization principle was first used in .
Fix a smooth function $`\theta (x)`$ such that $`\theta 1`$ on $`\stackrel{~}{S}`$, $`\theta 0`$ on $`𝐑^2[\frac{3s}{4},\frac{3s}{4}]^2`$, $`0\theta 1`$ and $`\theta _{\mathrm{}}8s^1`$. Let $`\theta _z(x):=\theta (xz)`$ and $`\eta _z:=\theta _z\phi _z`$. The function $`\theta _z`$ can be commuted with $`T_\beta `$ at the expense of an error of size $`\theta _z^2`$ on the support of $`\theta _z`$. The result is the analogue of (6.13) in
$$f,H_{M,\omega }f\frac{\beta }{2\pi }\left[_{\stackrel{~}{M}_z}\left\{\frac{1}{2}|T_\beta (\eta _zf)|^2+V_\omega |\eta _zf|^2\right\}dz128\pi \beta ^1s^2e^{\frac{\beta }{8}s^2}\right].$$
(18)
Finally, we estimate the boundary term $`dz\overline{V}_{S_z}|\eta _zf|^2`$ in $`\overline{H}_{S_z,\omega }^{(B+2\beta )}`$ using $`\overline{V}_{S_z}|\eta _zf|^2e^{(s/4\overline{L})^2}|\phi _z|^2|f|^2`$;
$$\frac{\beta }{2\pi }_{\stackrel{~}{M}_z}dz\overline{V}_{S_z}|\eta _zf|^2e^{s^2/(4\overline{L})^2}|f|^2,$$
which can be included into the error term in (18). $`\mathrm{}`$.
## 3 Enlargement of obstacles
We follow the basic strategy of Sznitman and its magnetic version from to estimate the lowest eigenvalue $`\lambda _{S,\omega }^{(B+2\beta )}`$ of $`\overline{H}_{S,\omega }^{(B+2\beta )}`$ by the lowest eigenvalue of a Hamiltonian with enlarged, hard-core obstacles. We need this argument for $`\overline{H}_{S,\omega }^{(B+2\beta )}`$, but the actual field does not play much role in this section, so for brevity we consider $`\overline{H}_{S,\omega }^{(B)}`$ and let $`\lambda _{S,\omega }^{(B)}`$ be its smallest eigenvalue.
The main novelty is that we cannot simply use Dirichlet boundary conditions for the ”enlarged obstacle” Hamiltonian, since the potential tail penetrating into the clearing regimes does influence the lowest eigenvalue. We add an appropriate Gaussian boundary potential to the hard-core Dirichlet wall, and we also keep the Gaussian tail of the original potentials. The obstacle configuration $`\omega `$ is fixed throughout this section.
The ”enlarged obstacle” Hamiltonian requires several definitions that were listed in Section 7.1 of . Here we recall only that four parameters, $`\mathrm{}`$, $`b`$, $`\epsilon >0`$ and $`r>0`$ have to be fixed. With these parameters, one defines (Section 7.1 of ) the set of ”good” points (their indices denoted by $`𝒢`$), clearing boxes and the set $`A^1`$, which is the $`\mathrm{}`$-neighborhood of clearing boxes. Recall that a point $`x_i`$ is ”good” if it is not isolated from other points in a certain hierarchical sense. Clearing boxes are squares of size $`\mathrm{}`$ that contain a large regular set (”clearing”) free of good points. Finally we define, for $`s>b`$,
$$\mathrm{\Omega }:=S\underset{i𝒢}{}\left[\overline{B}(x_i,2R)B(x_i,R)\right],\mathrm{\Omega }_+^b:=\left([s+b,sb]^2A^1\right)\underset{i𝒢}{}\overline{B}(x_i,b),$$
(19)
where $`B(x,\rho )`$ denotes the open ball of radius $`\rho `$ about $`x`$. We choose $`\delta =\frac{R}{100}`$. Notice that $`\mathrm{\Omega }`$ is defined by removing annuli around the good points, unlike in , where balls were removed. $`\mathrm{\Omega }_+^b`$ is the ”clearing set”, where the ”enlarged obstacle” Hamiltonian will be defined.
We let
$$U(x):=e^{|x|^2/\overline{L}^2}\mathrm{𝟏}(|x|R),\stackrel{~}{V}_\omega (x):=\underset{i𝒢}{}U(xx_i(\omega )),$$
then we clearly have $`V_\omega \stackrel{~}{V}_\omega `$. This definition of $`\stackrel{~}{V}_\omega `$ is different from (7.3) of . The role of $`v`$ in will be played by the constant $`e^{(2R)^2/\overline{L}^2}`$; this is a lower bound on the potential $`\stackrel{~}{V}_\omega `$ in the annuli $`\{R|xx_i|2R\}`$ around the good points. The role of $`a`$ in is played by $`2R`$. The specific upper bound $`a1`$ imposed in will not be important.
We will estimate the lowest eigenvalue, $`\stackrel{~}{\lambda }`$, of the Hamiltonian with potential $`\stackrel{~}{V}_\omega `$ by $`\stackrel{~}{\lambda }_b`$, the lowest eigenvalue with hard core potential on $`\mathrm{\Omega }_+^b`$. We add boundary potentials to both Hamiltonians. Since we work on multiply connected domains, we must take the gauge freedom into account as in Section 7.2 of . Hence both eigenvalues are defined as the infimum over all gauges on the complementary domain of the obstacles. We recall that for any $`\underset{¯}{\alpha }=\{\alpha _i\}_{i𝒢}[0,2\pi )^𝒢`$ we defined $`B_{\underset{¯}{\alpha }}(x):=B+_{i𝒢}\alpha _iB^{}(xx_i)`$ and its radial gauge $`A_{\underset{¯}{\alpha }}`$, $`\text{curl}A_{\underset{¯}{\alpha }}=B_{\underset{¯}{\alpha }}`$, where $`B^{}:=(4/\pi )\mathrm{𝟏}_{B(0,1/2)}(x)`$ (the definition of $`B^{}`$ in missed a $`\frac{1}{2\pi }`$ factor). The magnetic field $`B_{\underset{¯}{\alpha }}`$ includes flux tubes of strength $`\alpha _i`$ around the good points. We define
$$\stackrel{~}{\lambda }=\stackrel{~}{\lambda }(B):=\underset{\underset{¯}{\alpha }}{inf}\lambda _{\underset{¯}{\alpha }},\lambda _{\underset{¯}{\alpha }}:=inf\text{Spec}\left(\frac{1}{2}\left[(iA_{\underset{¯}{\alpha }})^2B_{\underset{¯}{\alpha }}\right]+\stackrel{~}{V}_\omega +\overline{V}_S\right)_S,$$
(20)
where the subscript refers to Dirichlet boundary conditions on $`S`$. Clearly $`\lambda _{S,\omega }^{(B)}\stackrel{~}{\lambda }`$. Similarly,
$$\stackrel{~}{\lambda }_b=\stackrel{~}{\lambda }_b(B):=\underset{\underset{¯}{\alpha }}{inf}\lambda _{b,\underset{¯}{\alpha }},\lambda _{b,\underset{¯}{\alpha }}:=inf\text{Spec}\left(\frac{1}{2}\left[(iA_{\underset{¯}{\alpha }})^2B_{\underset{¯}{\alpha }}\right]+V_{\mathrm{\Omega }_+^b}\right)_{\mathrm{\Omega }_+^b},$$
(21)
again with Dirichlet boundary conditions on $`\mathrm{\Omega }_+^b`$. Notice that the decay of the boundary potential $`V_{\mathrm{\Omega }_b^+}`$ is slightly stronger than that of $`\stackrel{~}{V}_\omega +\overline{V}_S`$ since $`L<\overline{L}`$.
Let $`g_U`$ denote the Green’s function of any domain $`U`$, i.e., the solution to $`\mathrm{\Delta }g_U=1`$ on $`U`$ and $`g_U=0`$ on $`U`$. We let $`G_U:=\mathrm{max}_{x\overline{U}}g_U(x)`$.
The importance of these functions is that the lowest magnetic Dirichlet eigenvalue of a large domain $`U`$ is essentially $`e^{2BG_U}`$ (a factor 2 was missing on page 349 of ), and the eigenfunction is roughly $`e^{Bg_U}`$ with some cutoff near the boundary. Moreover, for ”round” domains, $`g_U`$ is roughly quadratic in the distance from the boundary. Hence, roughly,
$$\stackrel{~}{\lambda }_b\mathrm{exp}(2BG_{\mathrm{\Omega }_+^b})+_{\mathrm{\Omega }_+^b}\mathrm{exp}\left(\left[\frac{\text{dist}(x,\mathrm{\Omega }_+^b)}{L}\right]^2\right)\left|\frac{\mathrm{exp}Bg_{\mathrm{\Omega }_+^b}(x)}{\mathrm{exp}Bg_{\mathrm{\Omega }_+^b}}\right|^2dx.$$
Here the first term represents the kinetic energy due to localization in the clearing. The second term is the interaction of the ”quantum” wavefunction with the ”classical” effect; the effective contribution of the potential tails. It turns out that the second term dominates. The main contribution comes from the interplay between the Gaussian character of the magnetic eigenfunction $`\mathrm{exp}Bg_{\mathrm{\Omega }_+^b}`$ and the Gaussian potential.
The basic comparison result is the analogue of Corollary 7.3 in (there are two misprints in (7.16) in ; $`r\mathrm{}`$ should be $`r0`$ and a minus sign is missing in front of $`\mathrm{log}K`$).
###### Proposition 3.1
For any fixed positive integer $`n_0`$ we let $`s:=n_0\mathrm{}`$ and let $`\varrho >0`$ be a positive number. For small enough $`r`$, $`\epsilon `$, there exist $`K=K(b,B,r,\mathrm{},s,L,\overline{L},\epsilon ,\varrho )`$ and $`w(r)`$ with $`lim_{r0}w(r)=1`$ such that $`\stackrel{~}{\lambda }_b^{w(r)}\stackrel{~}{\lambda }/K`$ if $`\stackrel{~}{\lambda }\mathrm{min}\{4K,e^{\varrho BG_\mathrm{\Omega }}\}`$, and $`K`$ satisfies
$$\underset{r0}{lim\; sup}\underset{\genfrac{}{}{0pt}{}{b\mathrm{}}{\epsilon 0}}{lim\; sup}\underset{\mathrm{}\mathrm{}}{lim\; sup}\frac{\mathrm{log}K}{\mathrm{}^2}=0.$$
(22)
The basic intuition behind this comparison is that if the lowest eigenvalue of $`H_{S,\omega }`$ is very small, then there must be a big clearing in the obstacle configuration, and the lowest eigenfunction is essentially supported in this clearing. Hence this eigenvalue can be estimated by the Dirichlet eigenvalue within the clearing even with enlarged obstacles. The inclusion of the boundary potential does not change this mechanism, but it changes both eigenvalues. The threshold for such eigenvalues is controlled by two different functions. The control given by $`K`$ is analogous to . The control $`\stackrel{~}{\lambda }e^{\varrho BG_\mathrm{\Omega }}`$ is new.
The proof of Proposition 3.1 is similar to that of Theorem 7.2 in , but we have to include the boundary potential. We first show that the increase of the eigenvalue due to the enlargement is given by the size of the eigenfunction near the boundary (Lemma 3.2). Then, by applying a probabilistic argument, we show that $`g_\mathrm{\Omega }(x)G_\mathrm{\Omega }`$ if $`\mathrm{\Omega }`$ is large and $`x`$ is close to the boundary (Lemma 3.3). In other words, the eigenvalue increases by at most a factor $`e^{o(BG_\mathrm{\Omega })}`$. For technical reasons we give these estimates for a slightly enlarged domain
$$\mathrm{\Theta }:=\mathrm{\Omega }+B(0,2\delta ).$$
In (Lemma 7.7), we finally estimated $`G_\mathrm{\Omega }`$ by the logarithm of the magnetic Dirichlet eigenvalue of $`\mathrm{\Omega }`$ to show that $`e^{o(BG_\mathrm{\Omega })}\stackrel{~}{\lambda }^{o(1)}`$ and therefore $`\stackrel{~}{\lambda }_b\stackrel{~}{\lambda }^{1o(1)}`$. The analogue of this estimate with a boundary potential is more complicated because it requires a control on $`g_\mathrm{\Omega }`$ not only near the boundary. But fortunately we do not need this estimate with a precise constant since it is used only in the error factor $`e^{o(BG_\mathrm{\Omega })}`$. So we choose an alternative method that estimates $`G_\mathrm{\Omega }`$ by the logarithm of the magnetic Dirichlet eigenvalue without boundary potential, exactly as in . The new control $`\stackrel{~}{\lambda }e^{\varrho BG_\mathrm{\Omega }}`$ stems from this modification.
We will state these lemmas precisely, but we give details of the proof only for the modifications compared with .
###### Lemma 3.2
There exist positive constants, $`c_1,c_2`$, depending only on $`B,L,\overline{L},R,b`$, such that
$$\stackrel{~}{\lambda }_bc_1\stackrel{~}{\lambda }s^2e^{2B\eta },$$
(23)
whenever $`\stackrel{~}{\lambda }c_2s^2e^{2B\eta }`$, $`s>2b40R`$, where
$$\eta =\mathrm{max}\left\{g_\mathrm{\Theta }(z):z\overline{\mathrm{\Theta }}\mathrm{\Omega }_+^{2b}\right\},\mathrm{\Theta }:=\mathrm{\Omega }+B(0,2\delta ),\delta :=\frac{R}{100}.$$
Proof. We fix $`\underset{¯}{\alpha }[0,2\pi )^𝒢`$ and let $`\phi _{\underset{¯}{\alpha }}`$ be the normalized eigenfunction belonging to $`\lambda _{\underset{¯}{\alpha }}`$. We can assume that $`\lambda _{\underset{¯}{\alpha }}c_2s^2e^{2B\eta }`$.
Let $`T_{\underset{¯}{\alpha }}:=i_1+_2(A_{\underset{¯}{\alpha }})_1i(A_{\underset{¯}{\alpha }})_2`$, then by variational principle and integration by parts
$$\lambda _{b,\underset{¯}{\alpha }}=\underset{\psi H_0^1(\mathrm{\Omega }_+^b)}{inf}\frac{_{\mathrm{\Omega }_+^b}\frac{1}{2}|T_{\underset{¯}{\alpha }}\psi |^2+V_{\mathrm{\Omega }_+^b}|\psi |^2}{_{\mathrm{\Omega }_+^b}|\psi |^2}.$$
Let $`\theta `$ be a cutoff function such that $`\theta 1`$ on $`\mathrm{\Omega }_+^{2b}`$, $`\theta 0`$ on $`𝐑^2\mathrm{\Omega }_+^b`$, $`0\theta 1`$ and $`|\theta |4b^1`$. Then
$$\lambda _{\underset{¯}{\alpha }}=\frac{1}{2}_S|T_{\underset{¯}{\alpha }}\phi _{\underset{¯}{\alpha }}|^2+_S(\stackrel{~}{V}_\omega +\overline{V}_S)|\phi _{\underset{¯}{\alpha }}|^2$$
$$\frac{1}{4}_S|T_{\underset{¯}{\alpha }}(\theta \phi _{\underset{¯}{\alpha }})|^2\theta _{\mathrm{}}^2_{Ssupp(\theta )}|\phi _{\underset{¯}{\alpha }}|^2+c_3^1_SV_{\mathrm{\Omega }_+^b}|\phi _{\underset{¯}{\alpha }}|^2,$$
using the pointwise inequality $`V_{\mathrm{\Omega }_+^b}(x)c_3\left[\stackrel{~}{V}_\omega (x)+\overline{V}_S(x)\right]`$ with some $`c_3=c_3(B,L,\overline{L},b)1`$. We use $`\psi :=\theta \phi _{\underset{¯}{\alpha }}`$ as a trial function to obtain
$$\lambda _{b,\underset{¯}{\alpha }}\frac{2c_3\lambda _{\underset{¯}{\alpha }}+16b^2_{S\mathrm{\Omega }_+^{2b}}|\phi _{\underset{¯}{\alpha }}|^2}{1_{S\mathrm{\Omega }_+^{2b}}|\phi _{\underset{¯}{\alpha }}|^2}$$
similarly to (7.19) in .
To complete the proof of (23), we need the upper estimate
$$_{S\mathrm{\Omega }_+^{2b}}|\phi _{\underset{¯}{\alpha }}|^2c_4\lambda _{\underset{¯}{\alpha }}s^2e^{2B\eta }$$
(with some $`c_4=c_4(B,\overline{L},R,b)`$) whose derivation is identical to the rest of the proof of Lemma 7.4 . The only difference is that the balls $`\overline{B}(x_i,a)`$ are replaced with the annuli $`\overline{B}(x_i,2R)B(x_i,R)`$, according to the new definition of $`\mathrm{\Omega }`$ in (19), and the constant $`v`$ is replaced with $`e^{(2R)^2/\overline{L}^2}`$. In particular
$$S\mathrm{\Omega }_+^{2b}\left(\mathrm{\Omega }\mathrm{\Omega }_+^{2b}\right)\underset{i𝒢}{}\left[\overline{B}(x_i,2R)B(x_i,R)\right],$$
and we estimate
$$_{_{i𝒢}\left[\overline{B}(x_i,2R)B(x_i,R)\right]}|\phi _{\underset{¯}{\alpha }}|^2\lambda _{\underset{¯}{\alpha }}e^{(2R)^2/\overline{L}^2}$$
instead of the second inequality in (7.20) in . The details are omitted. $`\mathrm{}`$
The next lemma is an analogue of Lemma 7.5 in . It states that near the boundary of the enlarged obstacles the Green’s function $`g_\mathrm{\Theta }`$ is much smaller than its maximum $`G_\mathrm{\Theta }`$.
###### Lemma 3.3
Let $`20Rb`$, $`40b\mathrm{}s`$, $`r<1/4`$ and $`G_\mathrm{\Theta }c_5\mathrm{}`$ with some $`c_5=c_5(R)`$.
(i) For small enough $`\epsilon `$, there exist $`\mathrm{}_0(\epsilon ,b)>0`$ and $`0<k=k(\epsilon ,b)1/4`$ such that
$$\underset{x\mathrm{\Sigma }\mathrm{\Theta }}{sup}g_\mathrm{\Theta }(x)\left(\frac{b}{\mathrm{}}\right)^kG_\mathrm{\Theta }\text{for}\mathrm{}\mathrm{}_0(\epsilon ,b)$$
(24)
with $`\mathrm{\Sigma }:=\left(\overline{S}_{}^{2\delta }S_+^{2b}\right)_{i𝒢}\overline{B}(x_i,2b)`$, where $`S_\pm ^c:=(s\pm c,sc)^2𝐑^2`$.
(ii) There exists a positive number $`c_0`$ such that
$$g_\mathrm{\Theta }(x)\left[(1c_0)^{1/r}G_\mathrm{\Theta }+c_0^1r^2\mathrm{}^2\right]+\underset{y\mathrm{\Sigma }\mathrm{\Theta }}{sup}g_\mathrm{\Theta }(y)$$
(25)
for all $`x\mathrm{\Theta }`$, $`xA^1S`$.
Proof. The proof is almost identical to that of Lemma 7.5 . The only difference is that $`\mathrm{\Omega }`$ is defined by removing annuli, hence the case $`xB(x_i,R+2\delta )`$, $`i𝒢`$, needs a separate estimate. For such $`x`$, the exit time from $`\mathrm{\Theta }`$ is at most the hitting time of the circle $`B(x_i,R+2\delta )`$, hence its expected value depends only on $`R`$. This estimate is taken into account in Lemma 3.3 by the extra requirement $`G_\mathrm{\Theta }c_5\mathrm{}`$. $`\mathrm{}`$
We use the following lemma to establish that the maximal expected hitting time is essentially the same for $`\mathrm{\Omega }`$ and $`\mathrm{\Theta }=\mathrm{\Omega }+B(0,2\delta )`$. The proof is identical to that of Lemma 7.6 ; the geometric condition used in is satisfied for the new definition of $`\mathrm{\Omega }`$ (19) as well.
###### Lemma 3.4
There exists two constant $`c_6,c_7`$ depending on $`R`$ such that
$$G_\mathrm{\Omega }G_\mathrm{\Theta }c_6G_\mathrm{\Omega }+c_7.\mathrm{}$$
(26)
Now we are ready to prove Proposition 3.1. We note that $`\overline{\mathrm{\Theta }}\mathrm{\Omega }_+^{2b}\mathrm{\Sigma }(A^1)^c`$. Hence for large enough $`\mathrm{}`$, the combination of (24), (25) and (26) gives
$$\eta c_8\left[(1c_0)^{1/r}G_\mathrm{\Omega }+c_0^1r^2\mathrm{}^2\right]$$
(27)
with some $`c_8=c_8(R)`$, similarly to (7.43) in . This means that $`\eta G_\mathrm{\Omega }`$ for small $`r`$ if $`G_\mathrm{\Omega }r^2\mathrm{}^2`$. Since $`\stackrel{~}{\lambda }e^{\varrho BG_\mathrm{\Omega }}`$, we see that $`e^{2B\eta }`$ is bounded by a small inverse power of $`\stackrel{~}{\lambda }`$, so from (23) we get that $`\stackrel{~}{\lambda }_b\stackrel{~}{\lambda }^{1o(1)}`$ as $`r0`$. The case $`G_\mathrm{\Omega }=O(r^2\mathrm{}^2)`$ can be included in the error factor $`K`$. The details are very similar to and are left to the reader. $`\mathrm{}`$.
## 4 Proof of the upper bound in Theorem 1.1
We recall the definition of $`\stackrel{~}{\lambda }(B)`$ and $`\stackrel{~}{\lambda }_b(B)`$ from (20) and (21). Using the notations and results of Section 3 for $`B`$ replaced with $`B+2\beta `$, we have $`\stackrel{~}{\lambda }(B+2\beta )\lambda _{S,\omega }^{(B+2\beta )}`$ with $`0\beta B/2`$, $`s=n_0\mathrm{}(t)`$ and $`\mathrm{}=\mathrm{}(t)=10\sqrt{\frac{\mathrm{log}t}{B}}`$. We will need $`\mathrm{\Omega }`$ defined in (19) and we note that $`\mathrm{\Omega }`$ depends on $`B`$, $`R`$, $`t`$, $`b`$, $`\epsilon `$, $`r`$ and $`n_0`$.
Combining Proposition 2.3 with Proposition 3.1, we see that
$$\underset{t\mathrm{}}{lim\; sup}\frac{\mathrm{log}L(t)}{\mathrm{log}t}$$
(28)
$$\underset{t\mathrm{}}{lim\; sup}(\mathrm{log}t)^1\mathrm{log}\left[\mathrm{exp}(tK\stackrel{~}{\lambda }_b^{w(r)})+\mathrm{exp}\left(te^{\varrho (B+2\beta )G_\mathrm{\Omega }}\right)+\mathrm{exp}(4Kt)\right]$$
with $`\stackrel{~}{\lambda }_b=\stackrel{~}{\lambda }_b(B+2\beta )`$. Since the $`n_0\mathrm{}`$ limit will always be taken before $`\beta 0`$, the condition $`n_0(B/\beta )^{1/2}`$ of Proposition 2.3 is satisfied. The last term in (28) is negligible for small enough $`\epsilon ,r`$ and large enough $`b`$, using (22). The estimates on the other two terms are given in the following propositions:
###### Proposition 4.1
For any magnetic field $`B>0`$
$$\underset{n_0\mathrm{}}{lim\; sup}\underset{r0}{lim\; sup}\underset{\genfrac{}{}{0pt}{}{b\mathrm{}}{\epsilon 0}}{lim\; sup}\underset{t\mathrm{}}{lim\; sup}(\mathrm{log}t)^1\mathrm{log}\mathrm{exp}\left(tK\left[\stackrel{~}{\lambda }_b(B)\right]^{w(r)}\right)=\pi \nu (L^2+2B^1).$$
(29)
###### Proposition 4.2
For any magnetic field $`B>0`$
$$\underset{n_0\mathrm{}}{lim\; sup}\underset{r0}{lim\; sup}\underset{\genfrac{}{}{0pt}{}{b\mathrm{}}{\epsilon 0}}{lim\; sup}\underset{t\mathrm{}}{lim\; sup}(\mathrm{log}t)^1\mathrm{log}\mathrm{exp}\left(te^{BG_\mathrm{\Omega }}\right)\frac{2\pi \nu }{B}.$$
(30)
Theorem 1.1 follows from these propositions via (12) just by choosing $`\varrho <2/(2+BL^2)`$, using Proposition 4.1 with a magnetic field $`B+2\beta `$ and Proposition 4.2 with a magnetic field $`\varrho (B+2\beta )`$, and finally letting $`\beta 0`$. $`\mathrm{}`$
Proof of Proposition 4.1. Let $`\mathrm{\Omega }`$ be an arbitrary domain and let $`B>0`$ fixed. Let
$$\widehat{\lambda }^{(B)}(\mathrm{\Omega }):=inf\{inf\text{Spec}\left(\frac{1}{2}\left[(i\widehat{A})^2B\right]+V_\mathrm{\Omega }\right)_\mathrm{\Omega }:\widehat{A}𝒜(\mathrm{\Omega })C^{\mathrm{}}(\overline{\mathrm{\Omega }}),\text{curl}\widehat{A}=B\text{on}\mathrm{\Omega }\}$$
be the smallest eigenvalue of the magnetic Hamiltonian with boundary potential and Dirichlet boundary conditions on $`\mathrm{\Omega }`$. We also took the infimum over all possible gauges, which is unnecessary for simply connected $`\mathrm{\Omega }`$. Here $`𝒜(\mathrm{\Omega })`$ is the set of real analytic vectorfields on $`\mathrm{\Omega }`$.
In Section 5 we show the following estimate for $`\widehat{\lambda }^{(B)}(\mathrm{\Omega })`$:
###### Proposition 4.3
For any $`\kappa >0`$, $`L>0`$, $`B>0`$ and any domain $`\mathrm{\Omega }`$ with volume $`|\mathrm{\Omega }|C(\kappa ,L,B)`$, we have
$$\widehat{\lambda }^{(B)}(\mathrm{\Omega })\mathrm{exp}\left[\frac{|\mathrm{\Omega }|}{\pi (L^2+2B^1)}(1+\kappa )\right].$$
(31)
Using this estimate and that $`\widehat{\lambda }^{(B)}`$ is a monotone function of the domain, the proof of Proposition 4.1 is identical to the argument in Section 8 of . $`\mathrm{}`$
Proof of Proposition 4.2. We consider
$$\mathrm{\Omega }^{}:=S\underset{i𝒢}{}\overline{B}(x_i,a)$$
with some fixed $`0<a<R`$, we let $`\stackrel{~}{V}^{}:=v\mathrm{𝟏}\left(x_{i𝒢}\overline{B}(x_i,a)\right)`$ with $`v=e^{(2R)^2/L^2}`$ and we let $`\stackrel{~}{\lambda }^{}=\stackrel{~}{\lambda }^{}(B)`$ be the infimum over $`\underset{¯}{\alpha }`$ of the lowest eigenvalue of $`\frac{1}{2}[(iA_{\underset{¯}{\alpha }})^2B_{\underset{¯}{\alpha }}]+\stackrel{~}{V}^{}`$. These are exactly the set $`\mathrm{\Omega }`$, the potential $`\stackrel{~}{V}`$ and the eigenvalue $`\stackrel{~}{\lambda }`$ in , but here we use the star superscript to distinguish them from their counterparts used in the present paper.
We claim that for any fixed $`B>0`$, $`n_0,R,L,r,\epsilon ,b`$
$$\underset{t\mathrm{}}{lim\; sup}(\mathrm{log}t)^1\mathrm{log}\mathrm{exp}\left(te^{BG_\mathrm{\Omega }}\right)\underset{t\mathrm{}}{lim\; sup}(\mathrm{log}t)^1\mathrm{log}\mathrm{exp}\left(te^{BG_\mathrm{\Omega }^{}}\right).$$
(32)
For the proof, we define
$$\mathrm{\Omega }^\mathrm{\#}:=S\underset{i𝒢}{}\overline{B}(x_i,2R).$$
Clearly $`g_\mathrm{\Omega }(x)=g_{\mathrm{\Omega }^\mathrm{\#}}(x)`$ for any $`x_{i𝒢}\overline{B}(x_i,R)`$, while $`g_\mathrm{\Omega }(x)c_9(R)`$ if $`x\overline{B}(x_i,R)`$ for some $`i𝒢`$. Hence $`G_\mathrm{\Omega }G_{\mathrm{\Omega }^\mathrm{\#}}+c_9(R)G_\mathrm{\Omega }^{}+c_9(R)`$, where the second inequality follows from $`\mathrm{\Omega }^\mathrm{\#}\mathrm{\Omega }^{}`$.
We then recall that Section 8 of , from (8.1) through (8.8) actually gave the following bound (there $`B`$ was replaced by $`B+2\beta `$):
$$\underset{n_0\mathrm{}}{lim\; sup}\underset{r0}{lim\; sup}\underset{\genfrac{}{}{0pt}{}{b\mathrm{}}{\epsilon 0}}{lim\; sup}\underset{t\mathrm{}}{lim\; sup}(\mathrm{log}t)^1\mathrm{log}e^{tN\stackrel{~}{\lambda }^{}(B)}\frac{2\pi \nu }{B}.$$
(33)
for any function $`N`$ satisfying
$$\underset{r0}{lim\; sup}\underset{\genfrac{}{}{0pt}{}{b\mathrm{}}{\epsilon 0}}{lim\; sup}\underset{t\mathrm{}}{lim\; sup}\frac{\mathrm{log}N}{\mathrm{log}t}=0.$$
Using Lemma 7.7 of , stating that $`\stackrel{~}{\lambda }^{}`$ is smaller than $`e^{BG_\mathrm{\Omega }^{}}`$ modulo negligible factors, we easily obtain (30) from (32) and (33). $`\mathrm{}`$
## 5 Estimate on the magnetic eigenvalue with a boundary potential
In this section we prove Proposition 4.3. Let $`D`$ be the disk of radius $`R_\mathrm{\Omega }:=\pi ^{1/2}|\mathrm{\Omega }|^{1/2}`$ centered at the origin. For any function $`a(r)`$ with
$$02\pi a(r)rB\pi r^2,\text{for all}0rR_\mathrm{\Omega },$$
(34)
we define a radial gauge $`A_{rad}(x)=a(r)\left(\begin{array}{c}\mathrm{sin}\theta \\ \mathrm{cos}\theta \end{array}\right)`$ (in polar coordinates, $`x=re^{i\theta }`$) that generates the radial magnetic field $`B_{rad}(x)=\text{curl}A_{rad}(x)=a^{}(r)+r^1a(r)`$. Condition (34) requires the flux of the magnetic field $`B_{rad}`$ to be not smaller than that of the constant $`B`$ field on all concentric disks $`B(0,r)`$.
Let $``$ and $`_D`$ be the Hilbert spaces of radially symmetric $`H^1(𝐑^2)`$ and $`H_0^1(D)`$ functions, respectively. Let
$$H(a):=\frac{1}{2}\left[(iA_{rad})^2B_{rad}\right]$$
be defined on $`_D`$, and let $`\lambda (a)`$ be its lowest eigenvalue. It is easy to see that the corresponding eigenfunction can be chosen nonnegative. In fact
$$f,H(a)f|f|,H(a)|f|=|f|,\frac{1}{2}(\mathrm{\Delta }+a^2B_{rad})|f|f_D.$$
(35)
Proposition 2.1 of states that the lowest magnetic Dirichlet eigenvalue on $`\mathrm{\Omega }`$ with a constant magnetic field $`B`$ is minorized by $`\lambda (a)`$ for some $`a(r)`$ that satisfies (34).
For any nonnegative function $`\psi `$ we denote its symmetric rearrangement by $`\psi ^{}`$, i.e., $`\psi ^{}`$ is the unique radial function with the property that $`|\{\psi c\}|=|\{\psi ^{}c\}|`$ for any $`c`$. It is not stated explicitly in , but actually the proof of Proposition 2.1 in gives the following result from which the comparison of the eigenvalues has been derived.
###### Proposition 5.1
Let $`\widehat{A}𝒜(\mathrm{\Omega })C^{\mathrm{}}(\overline{\mathrm{\Omega }})`$, $`\text{curl}\widehat{A}=B`$ on $`\mathrm{\Omega }`$ and let $`\widehat{H}=\frac{1}{2}\left[(i\widehat{A})^2B\right]`$. Then for any function $`fH_0^1(\mathrm{\Omega })`$ there exists a function $`a(r)`$ that satisfies (34) such that
$$f,\widehat{H}f|f|^{},H(a)|f|^{}.$$
For the proof one only has to notice that the radial trial function $`q(r)`$ defined on p. 289 of as $`q(r)=\mathrm{\Lambda }^1(h^{}(r))`$ is actually the symmetric rearrangement of $`\psi =|f|`$ since $`h=\mathrm{\Lambda }(\psi )`$ and $`\mathrm{\Lambda }`$ is strictly monotone. $`\mathrm{}`$
Now we include the boundary potential $`V_\mathrm{\Omega }`$ (see (16)). We replace $`V_\mathrm{\Omega }`$ by
$$\widehat{V}_\mathrm{\Omega }:=\frac{1}{\pi L^2}\mathrm{𝟏}_{\mathrm{\Omega }^c}e^{||^2/L^2}$$
where $`\mathrm{𝟏}_{\mathrm{\Omega }^c}`$ is the characteristic function of $`\mathrm{\Omega }^c`$ and $``$ denotes the convolution. It is straightforward that $`\widehat{V}_\mathrm{\Omega }V_\mathrm{\Omega }`$ on $`\mathrm{\Omega }`$ and $`\widehat{V}_\mathrm{\Omega }1`$ everywhere.
By the Riesz rearrangement inequality,
$$f,(1\widehat{V}_\mathrm{\Omega })f=\frac{1}{\pi L^2}|f(x)|^2e^{|xy|^2/L^2}\mathrm{𝟏}_\mathrm{\Omega }(y)dxdy\frac{1}{\pi L^2}\left(|f|^2\right)^{}(x)e^{|xy|^2/L^2}\mathrm{𝟏}_D(y)dxdy,$$
where the disk $`D`$ is the symmetric rearrangement of the set $`\mathrm{\Omega }`$. A simple estimate yields
$$f,\widehat{V}_\mathrm{\Omega }f|f|^{},W_\eta |f|^{}$$
for any $`fL^2(\mathrm{\Omega })`$ and any $`\eta >0`$ with
$$W_\eta (x):=\frac{1}{2}\mathrm{exp}\left((1+\eta ^1)\frac{(R_\mathrm{\Omega }|x|)^2}{L^2}(1+\eta )\right),|x|R_\mathrm{\Omega }.$$
From these estimates and Proposition 5.1 we conclude that there exists a radial function $`a(r)`$, satisfying (34), such that
$$\widehat{\lambda }^{(B)}(\mathrm{\Omega })inf\text{Spec}\left(H(a)+W_\eta \right)\text{on}_D.$$
Next we claim that if $`a_1(r)a_2(r)`$ satisfy (34), then
$$inf\text{Spec}\left(H(a_1)+W_\eta \right)inf\text{Spec}\left(H(a_2)+W_\eta \right)\text{on}_D.$$
This is proven exactly as Lemma 3.1 in . It is easy to check that the inclusion of a bounded nonnegative radial potential $`W_\eta `$ does not alter the trial function argument.
Therefore $`H(a)+W_\eta `$ has the lowest eigenvalue if $`a(r)=Br/2`$, i.e., in case of the constant field. Using (35), this eigenvalue is the same as the lowest eigenvalue, $`\lambda _\eta `$, of
$$H_\eta :=H_{osc}+W_\eta ,\text{with}H_{osc}:=\frac{1}{2}\left[\mathrm{\Delta }+\frac{Bx^2}{4}B\right]\text{on}_D.$$
Let $`\phi _0(x)=\sqrt{\frac{B}{2\pi }}e^{Bx^2/4}`$ span the kernel of the harmonic oscillator $`H_{osc}`$ on $``$, and let $`P:=|\phi _0\phi _0|`$ be the projection onto this kernel. It is well known that $`H_{osc}`$ has a gap of size $`B`$ above zero on $``$. We can estimate $`\lambda _\eta `$ by decomposing the eigenfunction $`f_D`$ as $`f=Pf+(IP)f`$:
$$\lambda _\eta =(IP)f,H_{osc}(IP)f+f,W_\eta fB(IP)f^2+f,W_\eta f.$$
(36)
Furthermore,
$$\lambda _\eta _DW_\eta |f|^2\frac{1}{2}_DW_\eta |Pf|^22_DW_\eta |(IP)f|^2\frac{1}{2}_DW_\eta |Pf|^2\lambda _\eta B^1,$$
using (36) and $`W_\eta 1/2`$. Hence
$$\lambda _\eta \frac{B}{2(B+1)}_DW_\eta |Pf|^2.$$
(37)
Since $`Pf^2+(IP)f^2=f^2=1`$ and $`(IP)f^2\lambda _\eta B^1`$ from (36), we have $`Pf^2=|f,\phi _0|^21\lambda _\eta B^1`$. We can assume that $`\lambda _\eta <B/2`$, otherwise Proposition 4.3 is trivial. Hence
$$\lambda _\eta \frac{B}{4(B+1)}_DW_\eta |\phi _0|^2=\frac{B^2}{16\pi (B+1)}e^{(1+\eta ^1)}_De^{(1+\eta )(R_\mathrm{\Omega }|x|)^2/L^2}e^{Bx^2/2}dx$$
(38)
$$C(B,L)\mathrm{exp}\left[(1+\eta ^1)\frac{R_\mathrm{\Omega }^2}{L^2+2B^1}(1+\eta )\right].$$
From this bound, Proposition 4.3 easily follows. $`\mathrm{}`$
Remark: From the integration (38) one can see the interplay between the Gaussian eigenfunction $`\phi _0`$ and the Gaussian potential $`W_\eta `$. In particular, the main contribution comes from the intermediate regime around $`|x|\frac{2}{BL^2+2}R_\mathrm{\Omega }`$.
|
warning/0003/hep-ex0003014.html
|
ar5iv
|
text
|
# 1 Outline
## 1 Outline
It is an exciting time to be working at the interface between physics and astronomy. Experiments built to detect astrophysical and atmospheric sources of neutrinos, such as Super-Kamiokande, are providing what may be the first definitive evidence for neutrino mass. Cosmology experiments, such as those measuring the the redshifts of type Ia supernovae and those determining the anisotropy of the microwave background, may have an important impact on particle physics as they confront our general view of the origin of the Universe. There have been similarly exciting developments at the high energy frontier of astronomy. Among other things:
1. we have detected a gamma-ray burst that is probably the most powerful explosion recorded since the Big Bang,
2. we have discovered extragalactic astrophysical sources that beam intense fluxes of TeV $`\gamma `$-radiation to us, and
3. we have observed individual particles (possibly protons) arriving from outer space with energies exceeding 25 Joules.
These new developments are clearly of prime importance for the field of high energy astrophysics, but, even more, research in this area in the future may make important contributions to particle physics depending, of course, on the nature of the discoveries.
In this paper, I provide a summary of the field of high-energy astronomy using photons, cosmic rays, and neutrinos. This field is rapidly developing because of recent discoveries and because of new experimental techniques that derive largely from accelerator-based particle physics detectors. I start with a broad overview of the field and discuss the general scientific motivations. Then, the three exciting developments listed above are described in more detail, followed by a discussion of the theoretical considerations. A selective review of the field in terms of the experimental techniques, results, and future prospects makes up the last part of the paper. I conclude with a summary of the prospects for the future.
## 2 Broad Overview
We learn about the Universe outside the immediate neighborhood of the Solar System by studying the arrival of four distinct messengers: 1) photons, 2) cosmic rays, 3) neutrinos, and 4) gravity waves . Here, cosmic rays are defined as nuclei (p, n, He<sup>++</sup>, etc.) electrons, and their antiparticles. In the future, we may discover other stable particles that convey information across interstellar space. If we do, we can add them to the previous list.
To date, the detection the photons over a wide range of energies has been the basis for the vast majority of astronomical discoveries. Cosmic rays provide important information about high energy processes occurring in our galaxy. Neutrinos and gravity waves both offer great astronomical potential, but are difficult to detect. So far, the neutrino source list is limited to our Sun and the supernova SN1987A.
Photon, cosmic ray, and neutrino astronomy are closely related at high energies (E$`>1`$GeV), both in terms of their astrophysical production mechanisms and in terms of their detection techniques. High energy particles are produced astrophysically by acceleration processes rather than by thermal processes which dominate at lower energies. The high energy particle fluxes typically exhibit rapidly falling power-law spectra, which leads to the requirement of very large detectors. The high energy domain spans a wide dynamic range (from 1 GeV to $`10^{11}`$ GeV), and we cannot expect a single detection technique to work at all energies.
In the future, we imagine studying astrophysical sources with multiple messengers that provide complementary information. A prototypical source would be gamma-ray bursts (GRBs). We know a great deal about GRBs from the electromagnetic radiation they produce, but there is also speculation that they are the source of the highest energy cosmic rays and that they produce a detectable neutrino flux . Thus, we anticipate that astronomy using these different messengers will become interrelated as the experiments become more powerful and as more detections are made.
The general scientific motivation for particle astronomy is multi-faceted. On the physical side, we use high energy radiation to probe extreme conditions of magnetic or gravitational potential. Our understanding of such astrophysical situations is still far from complete. The copious flux of high energy cosmic rays argues for efficient astrophysical accelerators with beam energies well beyond what we can achieve on Earth. High energy particle astronomy may also shed light on aspects of particle physics or cosmology beyond their respective standard models. On the astronomical side, the movement towards high energies continues the historical expansion of astronomy from optical light into new wavebands (radio, infrared, X-ray, and $`\gamma `$-ray). It is important to emphasize, however, that when discussing the many aspects of high energy astronomy, there is not yet a generally applicable “Standard Model”. For the experimentalist, this situation is ideal in that there are few constraints. New experimental results often change the general paradigm.
## 3 Three Recent Exciting Results
Here I discuss three of the most exciting results in this field in the last few years. The choices are, of course, very subjective.
### 3.1 Gamma Ray Bursts
Even after thirty years of research, the nature of gamma ray bursts (GRBs) remains one of the most important mysteries of astrophysics. We know from the BATSE detector on the Compton Gamma Ray Observatory that the arrival directions of GRBs are consistent with isotropy . However, the relatively poor angular resolution of BATSE of a few degrees has made the detection of GRB counterparts (i.e. sources at optical or radio wavelengths) difficult. The lack of counterparts has hindered efforts to pin down the distance scale of GRBs.
In 1997, a major breakthrough was achieved with the first detection of optical counterparts. This work was accomplished by the Beppo/SAX satellite detector in conjunction with powerful optical telescopes, both on the ground (e.g. Keck) and in space (e.g. HST). By using a combination of wide and narrow-field X-ray telescopes, Beppo/SAX determines the positions of some bursts with an accuracy of several arc minutes. This excellent localization allowed the detection of fading optical counterparts for several dozen bursts . Redshift values have been determined for approximately twenty counterparts; they indicate that the correlated bursts are cosmological in origin (typical $`z1`$). The typical inferred energy outputs of the bursts range between $`10^{51}`$ and $`10^{53}`$ ergs (assuming isotropic emission).
The detection of optical counterparts to GRBs is clearly a landmark discovery. Until 1999, however, optical observations were limited to time periods well after the actual bursts because of the long delay (typically 8 hours) required by Beppo/SAX to achieve accurate position localization. What we really want to do is to carry out out simultaneous $`\gamma `$-ray and optical observations of GRBs. Such observations require a fast slewing telescope with a wide field of view to cover the large BATSE error box.
Several wide-field optical telescopes have been constructed or are currently under development. One of the pioneering experiments, known as ROTSE-I, consists of an array of 35 mm telephoto lenses coupled to large format CCD detectors. With an overall field of view of $`16^{}`$x$`\mathrm{\hspace{0.17em}16}^{}`$, and by slewing in an automated fashion upon receiving an alert from the GRB Coordinates Network (GCN), ROTSE is ideally suited for rapid optical follow-up of GRBs. The careful design of the experiment paid off on January 23, 1999 when ROTSE made the first detection of contemporaneous optical radiation from a GRB .
From the perspective of BATSE, GRB990123 was relatively ordinary, but the optical signature detected by ROTSE, was truly remarkable. As shown in Figure 1, the optical brightness increased by 3 magnitudes in 25 seconds and then waned by 5 magnitudes over a period of 8 minutes. At its peak brightness, the optical magnitude was 8.95 which, when combined with the measured redshift value of 1.60, meant that this burst was the most luminous object ever detected. GRB990123 was brighter than the brightest quasar by several orders of magnitude. Assuming isotropic emission, the inferred energy release of GRB990123 exceeds $`10^{54}`$ergs. Theoretical models, already grappling to explain the wide variety of GRB phenomena, are further strained to deal with this remarkable energy output.
A wide range of theoretical models have been proposed to explain gamma ray bursts. There are almost as many models as bursts! The basic difficulty is to construct a physical mechanism that can produce and extract the intense high energy emission we observe. The general picture calls for a cataclysmic event which produces a relativistic fireball of material with Lorentz factors approaching 1000 . The relativistic material escapes from the region of high energy density along a collimated jet. High energy radiation results when the jet collides with nearby ambient material. The nature of the original cataclysmic event is not fully understood. Generally favored pictures include colliding neutron stars and hypernovae (“failed supernovae” of heavy stars) .
The recent discoveries in the area of gamma ray bursts are clearly profound, but the overall puzzle is far from being solved. Many additional questions remain. For example, when considering the time durations and spectral shapes of GRBs, we know from the BATSE data that there are at least two classes of bursts, (if not more) . For instrumental reasons, Beppo/SAX is sensitive only to longer bursts (duration greater than 1 second). Thus, so far, we can confidently ascribe a cosmological origin only to a portion of the burst population. The medium and short length bursts may have a different origin. We can expect GRB research to remain exciting for years to come.
### 3.2 TeV $`\gamma `$-rays from Extragalactic Sources
The field of very high energy (VHE) $`\gamma `$-ray astronomy has come of age in the last ten years . Ground-based telescopes using the atmospheric Cherenkov technique have made strong detections of TeV photons from galactic and extragalactic sources. The most exciting development in this area has been the discovery of strong VHE emission from active galactic nuclei (AGN), which make up a broad class of extragalactic objects including quasars. AGN are powerful sources at all wavelengths, but the importance of their $`\gamma `$-ray output has only recently been fully appreciated. We now know that in some AGN, those known as blazars, the bulk of their power is emitted at high energies. Explaining this observational fact is an important problem facing theorists.
In 1997, the blazar source Markarian 501 (Mrk 501) entered a period of dramatic activity. The TeV $`\gamma `$-ray emission from the source increased by up to a factor of fifty from earlier epochs. Figure 2 shows the VHE $`\gamma `$-ray flux detected from this source over a four year period by the Whipple experiment. The difference between the average flux level seen in 1997 and in the other years is striking. At certain times, the $`\gamma `$-ray flux from Mrk 501 exceeded the brightest known source in our galaxy, the Crab Nebula, by factors of three to four. This happened in spite of the fact that Mrk 501 is more distant from us than the Crab by a factor of 10,000. At its maximum brightness, Mrk 501 beamed 10<sup>11</sup> VHE photons per second to the Earth’s surface. The VHE spectrum extends to maximum energies above $`20`$TeV.
Another impressive feature of the 1997 emission from Mrk 501 was the high degree of variability observed. The $`\gamma `$-ray flux varied by as much as an order of magnitude from night to night and by factors of two on hourly time scales. These variations imply that the acceleration region where the $`\gamma `$-rays are produced must be very compact, presumably several light hours across (times any relativistic Doppler factors). The flux variations of Mrk 501 were studied by several state-of-the-art Cherenkov telescopes around the globe and there was good agreement between the flux levels recorded by the different experiments.
AGN are such luminous objects that the only feasible source for their power is in the intense gravitational field near black holes. The general picture of an AGN is a supermassive ($`10^810^9`$ M) black hole surrounded by a rotating disk. Accreting material provides the power for the broad-band emission observed. It also powers relativistic jets directed along the angular momentum axis of the black hole. As shown in Figure 3, shock acceleration in the jets produces high energy beams of electrons or protons which interact with radiation fields to produce secondary beams of $`\gamma `$-rays and neutrinos. A key aspect of this model is that we observe blazars when their beams are directed into our line of sight.
In the general blazar picture, the high energy $`\gamma `$-rays probe acceleration processes in the jet. The fact that very rapid $`\gamma `$-ray variability has been observed indicates that the acceleration region relevant to high energy particle production is deep within the jet, and possibly quite close to the black hole itself. Thus, understanding the VHE $`\gamma `$-ray emission is of crucial interest.
### 3.3 Cosmic Rays with Energies Exceeding $`10^{20}`$eV
We have known about the existence of extremely high energy ($`E>10^{20}`$eV) cosmic rays for more than forty years. In 1966, it was realized by Greisen, Zatsepin, and Kuz’min (GZK) that cosmic rays above the energy of $`6\times 10^{19}`$eV (if they existed!) would interact with the 3 K cosmic microwave background radiation (CMBR). This GZK cutoff limits the mean free path for the highest energy cosmic rays to less than 100 Mpc, a distance that is quite small in comparison to typical extragalactic scales. Until recently, conventional wisdom held that cosmic rays above the GZK cutoff would not be detected on Earth because: 1) it was hard to construct astrophysical sources capable of accelerating particles to these energies, and 2) if such source existed, they would almost certainly be located at great distances and the cosmic rays they produced would be absorbed by the CMBR. This wisdom proved to be wrong. It now seems clear that the cosmic ray spectrum continues to energies of $`10^{20}`$eV, and beyond.
The most compelling evidence for particles beyond the GZK cutoff comes from the AGASA experiment . AGASA is a large surface array covering an area of approximately $`100`$km<sup>2</sup>, and located in the central portion of the Japanese island of Honshu. The array samples the particle cascade in giant air showers produced from extremely high energy cosmic ray interactions in the atmosphere. The cosmic ray energy is estimated from the particle density determined a fixed distance from the core of the shower. Uncertainties in the energy measurement are estimated by several techniques, including simulation methods. The typical energy resolution is $`30`$%. More importantly, the proportion of events with a 50% or more overestimation in energy is less than 3%.
The latest cosmic ray energy spectrum released by AGASA is shown in Figure 4 . Based on a data sample collected between 1993 and 1998, the spectrum shows no evidence for a cutoff above $`6\times 10^{19}`$eV. Instead, seven events are detected above $`10^{20}`$eV. The measured spectrum from AGASA is inconsistent with that expected from a population of extragalactic sources distributed uniformly in the Universe. Barring an unforeseen loophole to the GZK cutoff, the AGASA results, in conjunction with the detection by the Fly’s Eye experiment of a cosmic ray at $`3\times 10^{20}`$eV , strongly support the idea the sources of the highest energy cosmic ray events are “local”.
Considering the enormous rigidity of $`10^{20}`$eV particles and the weakness of the intergalactic magnetic field ($`B<10^{10}`$g), we do not expect substantial deflection of such particles over path lengths of 50 Mpc. Thus, it is reasonable to consider doing astronomy with the extremely high energy cosmic rays. The arrival directions of the cosmic rays above $`10^{19}`$eV have been examined. No obvious concentration in the sky can be discerned. Also, there is no significant correlation between the cosmic ray positions and those of known astronomical sources (e.g. quasars, radio galaxies, etc.). Thus, we are faced with a quandary similar to that of gamma ray bursts. Since the sources of the highest energy cosmic rays cannot be correlated with known astronomical objects, they must represent something new. To achieve their extreme energies, the cosmic rays must be produced in a truly remarkable astrophysical accelerator. They could, in fact, come from physics at a higher mass scale. It is perhaps this latter possibility that makes these particles so intriguing.
## 4 Theoretical Considerations
### 4.1 Extreme Astrophysics
As introduced in Section 2, we expect high energy $`\gamma `$-rays, cosmic rays, and neutrinos to be produced via particle interactions at sites of powerful acceleration. When discussing such production, we must distinguish between the power source and the acceleration mechanism. For sources using conventional (i.e. known) physics, the power source will make use of extreme electromagnetic or gravitational potentials. Pulsars, such as the Crab, are examples of the former case. Supernova remnants and AGN are examples of the latter.
High energy $`\gamma `$-radiation has been detected from a number of known radio pulsars. The radiation can be categorized as being pulsed (i.e. with the same period as that detected in the radio) or unpulsed. For both categories, the ultimate power source derives from the spin-down of a highly magnetized neutron star. At the highest energies, the radiation is believed to originate from inverse-Compton scattering of soft photons by a relativistic wind of electrons . In the case of the Crab, from which $`\gamma `$-rays up to 50 TeV have been detected, the electron spectrum is believed to extend above $`1`$PeV, making the Crab the high-energy accelerator known in the Universe.
Supernova remnants (SNRs) are attractive candidates for particle acceleration because of the enormous power contained in the original explosion. Supernova explosions typically release $`10^{51}`$ergs of kinetic energy and they occur every few decades somewhere in the galaxy. Therefore, the average power injected into the interstellar medium by these explosions is $`10^{42}`$ergs s<sup>-1</sup>. Since the power required to replenish the cosmic rays is $`10^{40}10^{41}`$ergs s<sup>-1</sup>, supernova remnants can explain the origin of the cosmic rays if $`110`$% of the kinetic energy released in the explosion goes into the acceleration of protons and nuclei.
Particle acceleration in SNRs is expected to occur at the shock front produced as the remnant material traverses the interstellar medium. Charged particles scatter off magnetic field irregularities as they diffuse back and forth across the shock front. Because the velocity distributions of the scattering centers on either side of the shock are isotropic, the particles see a converging flow of scattering centers on both sides of the front. They thus gain energy with each round trip passage. This particle acceleration process, first proposed by Fermi in a different context , continues as long as the particles are contained in the vicinity of the shock front. The containment time is proportional to energy, and so diffusive shock acceleration naturally leads to a power law energy spectrum for the accelerated particles, (dN/dE $``$E<sup>α</sup>). For strong (highly supersonic) shocks, the power law spectral index is $`\alpha 2`$.
The maximum obtainable particle energy in SNRs is determined by the lifetime of the shock itself. As the shock wave expands, it slows downs and weakens. For typical shock lifetimes ($`5000`$yr), the maximum particle energy is calculated to be Z$`\times 10^{14}`$eV, where Z is the particle charge . Thus, SNRs offer a plausible explanation for the origin of cosmic rays up to an energy of $`10^{15}`$eV (and possibly up to $`10^{16}`$eV), but a new source is required to explain higher energy cosmic rays.
There is a great deal of speculation about the origin of particles above $`10^{16}`$eV. Galactic objects do not in general have the combination of size and magnetic field strength to contain a particle at these energies . Therefore, it is generally believed that the sources must be extragalactic. AGN are attractive possibilities because they are known to be powerful emitters of $`\gamma `$-rays. They may also produce very high energy cosmic rays and neutrinos.
The general model for AGN was outlined in Section 3.2 (see Figure 3). Charged particles are accelerated to high energies in the jet via the Fermi shock mechanism. The accelerated particles are likely to be electrons or protons (or a combination of the two). Electrons are certainly needed to explain the broad-band synchrotron emission observed. If electrons dominate, $`\gamma `$-rays could naturally be produced by inverse-Compton scattering off low energy photons . The soft photons may originate from synchrotron radiation, from disk emission, or from reprocessed emission from clouds or dust . If protons dominate, they would interact with radiation fields to produce electromagnetic cascades that would ultimately produce high energy $`\gamma `$-rays . Protons may also produce neutrinos via a beam dump mechanism (i.e. by impinging on material to produce charged pions which then decay).
Extragalactic sources, such as AGN, may explain the origin of cosmic rays between $`10^{16}`$eV and $`10^{19}`$eV. However, as discussed in Section 3.3, the highest energy particles cannot come from large extragalactic distances. Their origin is a true mystery. Given our inability to explain their origin via known astrophysics, we are forced to seriously consider other possibilities, such as new astrophysics or particle physics beyond the Standard Model.
### 4.2 Non-Standard Astrophysics or Particle Physics
When explaining the existence of particles arriving at Earth with very high energies, we must naturally consider possible sources outside known astrophysics. In this case, particles would acquire their energies in a “top down” picture,i.e. as coming from physics at a higher mass scale, as opposed to the “bottom up” picture of acceleration. Cosmology must come into play here, because the early Universe offers the most natural conditions needed to generate interactions or particles beyond the Standard Model. Here we give a few examples of non-standard sources of very high energy particles.
There is speculation that the highest energy cosmic rays result from the collapse of topological defects produced in the early Universe . Defects such as cosmic strings resulting from phase transitions could produce grand unified (GUT) scale particles at with masses of $`10^{14}10^{15}`$GeV. GUT particles would decay to leptons and quarks, producing electromagnetic and hadronic cascades whose eventual tertiary products would cosmic rays or $`\gamma `$-rays. Under certain assumptions, the cosmic rays would be produced close enough to Earth to avoid the GZK cutoff.
Another interesting possibility is that of primordial black holes. Black holes are not truly black, but emit a spectrum of radiation . As the black hole loses energy, it gets hotter, and this results in more available quantum states for the emission. Eventually, the black hole evaporates in a final explosion that releases a burst of radiation. Small black holes created in the early Universe with masses near $`10^{14}`$g would be in the process of evaporation now. Such primordial black holes would produce bursts of VHE $`\gamma `$-rays on time scales of microseconds to seconds .
The search for evidence of supersymmetry is a very important frontier area of particle physics. Given the possibility that supersymmetry may manifest itself at TeV energies, it is natural to consider VHE astrophysical signatures. One possible signature would be the direct detection of neutral supersymmetric particles. For example, it has been proposed that the highest energy cosmic rays are supersymmetric strongly interacting particles ($`S^0`$) having a much greater path length through the CMBR than ordinary nucleons . A second possible signature for supersymmetry would be the detection of secondary decay products from a supersymmetric particle. The most interesting possibility here is that of the neutralino, X̃<sup>0</sup>. Neutralinos may be the lightest supersymmetric particles and they may also comprise a large fraction of the dark matter in the galaxy . Since neutralinos must couple to ordinary matter through weak interactions, the weak scale cross-section determines the temperature at which the neutralino would freeze out, and hence its relic abundance. Neutralinos concentrating in the galactic center would annihilate via loop processes, X̃<sup>0</sup>$`{}_{}{}^{0}\gamma \gamma ,Z\gamma `$, yielding a very high energy $`\gamma `$-ray signature . Similarly, neutrinos would be produced from the annihilations of gravitationally trapped neutralinos in the center of the Sun or the Earth. Estimates for the $`\gamma `$-ray or neutrino fluxes are sensitive to parameters of the given model of supersymmetry, to the mass of the neutralino, and to the density profile of galactic dark matter. In spite of large uncertainties in expected flux levels, the search for supersymmetry by these techniques is important for several reasons: 1) the signature, a highly directional, mono-energetic flux of $`\gamma `$-rays or neutrinos, would be unambiguous, (i.e. it would provide “smoking gun” evidence for supersymmetry), 2) a statistically significant detection would allow measurement of the neutralino mass, and 3) the next generation of $`\gamma `$-ray and neutrino telescopes are sensitive to neutralino masses up to several TeV, which may be out of reach of accelerator experiments for some time to come.
## 5 The Experimental Situation
### 5.1 Detection Techniques: $`\gamma `$-rays and Cosmic Rays
As shown in Figure 5, $`\gamma `$-rays and cosmic rays span an enormous range of energies, from $`1`$MeV to $`10^{20}`$ eV. Given this range, a single detection technique will not suffice. Satellite and balloon experiments above the Earth’s atmosphere operate at MeV and GeV energies. At high energies, the particle flux is small enough so that space-borne instruments become flux limited. For $`\gamma `$-rays, the practical upper limit of sensitivity for the instruments on the Compton Gamma Ray Observatory is $`20`$GeV. For cosmic rays, the upper range of space-borne detectors is higher by several orders of magnitude due to the greater cosmic ray flux.
At energies above the reach of balloon and satellite experiments, we use the Earth’s atmosphere as the target and absorber of high energy particles. Primary $`\gamma `$-rays and cosmic rays interact in the atmosphere to create extensive air showers that propagate to the ground. Ground-based detectors sample the Cherenkov radiation, the charged particles, or the blue/UV fluorescence from nitrogen excitations to determine the arrival directions and energies of the incoming primary particles. As shown in Figure 5, atmospheric Cherenkov telescopes, air shower arrays, and Nitrogen fluorescence detectors operate in various energy ranges from $`200`$GeV to $`10^{20}`$eV.
### 5.2 $`\gamma `$-ray Experiments and Results
There are a number of ground-based $`\gamma `$-ray telescopes in operation around the world. This paper does not attempt to be comprehensive, and several recent articles have reviewed the experimental situation . The state-of-the-art Cherenkov telescopes include Whipple in Arizona, USA, HEGRA on the island of La Palma, CAT at Themis, France, and CANGAROO at Woomera, Australia. Operating air shower arrays include HEGRA and the Tibet Array in Yanbajing, Tibet.
$`\gamma `$-ray astronomy is a quickly changing field with a variety of sources and new phenomena discovered during the last decade. The EGRET experiment on the Compton Gamma Ray Observatory has detected $`150`$ sources at energies between 30 MeV and 20 GeV. Many of the EGRET sources remain unidentified, but others have been associated with pulsars and AGN. At energies above 200 GeV, the current status of the field is shown in Figure 6. Atmospheric Cherenkov telescopes have detected at least seven VHE sources and five additional sources have been tentatively identified.
The sources detected at very high energies can be categorized as pulsar nebulae, AGN, and SNRs. The detection of each different source category was a significant advance which brought us new information about high energy astrophysical phenomena . These detections have also raised numerous questions. For example, for AGN, remarkable emission has been detected from a few nearby sources, but many questions remain about how these sources work and why we cannot see more distant objects.
For the future, there are a variety of important areas that need exploration:
1. To date, no sensitive experiments have operated in the energy range between 20 and 200 GeV. New instruments, and possibly new techniques, are needed to explore this energy band where the likelihood of exciting new astrophysics is high. For example, it is in this energy range that we expect to measure spectral features resulting from the absorption of $`\gamma `$-rays from by the cosmic infrared radiation .
2. So far, it appears that we have detected very high energy radiation from the most luminous astrophysical sources. Significant increases in flux sensitivity should lead to the detection of fainter objects and more types of sources.
3. The most important $`\gamma `$-ray results at energies above 200 GeV have come from atmospheric Cherenkov telescopes which have limited fields of view and operate only on dark, clear nights. We need instruments which have wide fields of view and operate with close to 100% duty cycle. Such instruments would be able to detect dramatic transient phenomena, such as gamma ray bursts.
To address the major goals for the future, there are a number of new high energy $`\gamma `$-ray telescopes soon to come on line or being proposed for the future. To extend the reach of ground-based telescopes to lower energies, the STACEE and CELESTE experiments each use large arrays of solar heliostat mirrors to detect fainter Cherenkov showers and thus to observe $`\gamma `$-rays at lower energies. Both experiments are in the final stages of construction and have reported promising initial results . The MAGIC imaging Cherenkov telescope , which consists of a large 17 m reflector and a state-of-the-art camera, is now under construction. MAGIC hopes to achieve an effective energy threshold as low as 20 GeV.
A new air shower experiment that has recently come on line is MILAGRO . Consisting of a large man-made pond of purified water viewed by $`1000`$ photomultiplier tubes, MILAGRO detects air showers at a median energy threshold of 1 TeV. By using the air shower technique, MILAGRO will carry out the first all-sky survey at TeV energies. The encouraging initial results from a prototype of MILAGRO include the tentative detection of Mrk 501 and possible correlated TeV $`\gamma `$-ray emission from a gamma ray burst .
In order to substantially improve the sensitivity, and to extend the energy range, of the atmospheric Cherenkov technique, several groups are developing new telescopes that employ arrays of imaging reflectors / HESS, initially consisting of four 12 m reflectors in Namibia, and Super-CANGAROO, consisting of four 10 m reflectors located in Australia, will be the premier ground-based $`\gamma `$-ray telescopes in the Southern Hemisphere. In the Northern Hemisphere, VERITAS , is being proposed as an array of seven 10 m reflectors. As shown in Figure 7, the dramatic improvement in sensitivity of VERITAS over Whipple (which represents the current generation of atmospheric Cherenkov telescopes) will allow for much more detailed study of rapid very high energy phenomena.
The single most important new $`\gamma `$-ray telescope will be flown in space. GLAST , will be a state-of-the-art detector using many techniques of experimental particle physics, such as Si-strip tracking and CsI calorimetry. With a very wide field of view, and a suitable pointing strategy, GLAST will scan the entire sky on every orbit, offering unparalleled coverage of transient $`\gamma `$-ray phenomena, such as AGN flares and gamma ray bursts. GLAST will have substantially improved characteristics (angular resolution, energy resolution, energy range, etc.) relative to its predecessor, EGRET. The resulting improvement in sensitivity of GLAST will enable the detection of up to two orders of magnitude more sources (e.g. approximately 3000-4000 AGN). An artist’s conception of GLAST is shown in Figure 8.
### 5.3 Cosmic Ray Experiments and Results
The variety of experimental techniques for the detection of high energy cosmic rays is summarized in Figure 5. Satellite instruments have measured the elemental composition of the cosmic rays at energies up to $`10^{14}`$ eV. An important outstanding issue concerns the level of antimatter, e.g. positrons, antiprotons, antihelium, etc. A primordial source of antimatter would be an important discovery. Searching for this source is one of the main goals of the Alpha Magnetic Spectrometer (AMS) satellite detector. In June 1998, AMS took initial cosmic ray data during a flight on the space shuttle Discovery. Initial results based on the AMS data have been presented . AMS is now in the process of reconfiguration for an extended physics program on the International Space Station.
At energies above $`10^{14}`$ eV, ground-based air shower detectors provide the only effective way to carry out cosmic ray measurements. Determining the cosmic ray composition by indirect detection techniques has been notoriously difficult. No clear understanding of the composition exists in the energy region near the bend in the energy spectrum (the “knee” at $`3\times 10^{15}`$eV). It is in this region, however, where we expect to learn something about the origin of cosmic rays . Future ground-based and satellite experiments may shed light on this difficult, but important, problem.
For the highest energy cosmic rays above the GZK cutoff, the current experimental situation is summarized in Figure 4. There is good evidence that the cosmic ray spectrum extends to $`10^{20}`$ eV and beyond, but the statistics are not overwhelming. It is important to keep in mind that the flux of particles above $`10^{20}`$ eV is approximately 1 per km<sup>2</sup> per century! The AGASA experiment, which has detected the majority of the events beyond the GZK cutoff, has a collection factor of approximately 120 km$`{}_{}{}^{2}`$sr. The newly commissioned Fly’s Eye HiRes experiment improves upon this factor by an order of magnitude. Fly’s Eye HiRes consists of two fluorescence detector sites located 13 km apart. Each detector employs spherical mirrors to reflect air shower fluorescence light onto photomultiplier tube cameras. The construction for the experiment has been completed, and preliminary results based on data-taking with an initial configuration have been reported . The early Fly’s Eye results provide clear confirmation of the AGASA result that the cosmic ray spectrum continues to $`10^{20}`$eV, and beyond.
The exciting results relating to the highest energy cosmic rays have prompted the development of new, and even larger, cosmic ray experiments. The Auger Project plans to construct two giant air shower arrays, one each in the Southern and Northern Hemispheres. The arrays incorporates both particle and nitrogen fluorescence detectors . Each array consists of 1600 particle detection stations on a grid covering approximately 3000 km<sup>2</sup>. Three fluorescence detectors are used to improve the energy calibration of the ground array. As shown in Figure 9, a detector station employs a large water tank viewed by self-contained instrumentation to record the arrival of the air shower particles. The Southern Hemisphere array is currently under construction in Mendoza, Argentina.
The Northern Hemisphere array of Auger is planned for construction in Millard County, Utah. A very large fluorescence detector, the Telescope Array, is also being considered for construction in the same general region as Auger. A possible future satellite experiment, OWL/Airwatch , would detect giant air showers using a downward-looking fluorescence detector. This ambitious instrument is currently under development by groups in the United States and Italy.
### 5.4 Very High Energy Neutrino Astrophysics
As pointed out by Greisen forty years ago , there is an obvious connection between $`\gamma `$-ray and neutrino astronomy. Both fields involve the detection of neutral particles produced as secondaries in high energy astrophysical accelerators. In many cases, the same sources give rise to fluxes of high energy $`\gamma `$-rays and neutrinos.
There are important differences, however, between photon and neutrino astronomy. Neutrinos have the great advantage over $`\gamma `$-rays in that their very small interaction cross section allows them to travel unimpeded from their production sites to Earth. Thus, in principle, neutrino astronomy allows us to probe the dense central regions of objects such as gamma ray bursts, AGN, and supernovae. We can also search for neutrino signals coming from the annihilation of dark matter concentrated at the center of the Earth and Sun. Neutrino astronomy may also shed light on important aspects of particle physics. A clear demonstration of this possibility is the evidence for the oscillation of neutrinos produced by cosmic ray interactions in our atmosphere.
The most probable source of high energy astrophysical neutrinos is from the decay of charged pions produced in a hadronic beam. Thus, unlike $`\gamma `$-rays, which can arise from electromagnetic processes such as synchrotron radiation and inverse-Compton scattering, neutrinos will likely arise from astrophysical sources containing proton beams. Although we have not yet conclusively identified sources such as these, we know from the flux of cosmic rays that energetic hadrons are produced somewhere in large quantities.
The fact that neutrinos have a very small interaction cross section also poses a big difficulty for the experimentalist in that very large detectors are required. So far, this difficulty has limited the astrophysical information learned from neutrino telescopes. It is difficult to predict for certain how large experiments must be in order to detect neutrinos from AGN or gamma ray bursts. The appropriate detector collection area is probably 1 km<sup>2</sup>, or perhaps even larger .
For some time, it has been recognized that clear water (liquid or ice) would make a very suitable medium for neutrino detection. The basic idea is to instrument a large volume of water with photomultiplier tubes arranged on long strings. Upward-going neutrinos (coming from astrophysical sources on the other side of the Earth) interact in the large volume of Earth below the detector. The high energy muons produced from these interactions are detected via their Cherenkov radiation in the instrumented volume. The general method is shown in Figure 10. The upward-going signature is required to reduce the very large ($`>10^5`$) background of downward-going muons produced in cosmic ray air showers. The cosmic rays also produce a flux of atmospheric neutrinos on the other side of the globe that can be used as a calibration source for the detector.
Early pioneering work in building very large neutrino telescopes was carried out by the DUMAND and BAIKAL collaborations. Indeed, BAIKAL was the first detector to observe very high energy atmospheric neutrinos . Presently, the AMANDA installation at the South Pole represents the current state-of-the-art in high energy neutrino telescopes. AMANDA has grown in stages during the last five years. The current instrument, AMANDA-II, comprises 675 optical modules on 19 strings located at ice depths of 1450 to 2050 meters. AMANDA-II has a collection area of approximately 30,000 m<sup>2</sup> ( 1 TeV). An earlier version of the detector, AMANDA-B, operated between 1997 and 1999 and recorded a large data set. As shown in Figure 11, by means of stringent quality cuts on the track fitting, a clean sample of 17 muon tracks due to atmospheric neutrino events has been isolated . This result demonstrates that AMANDA can detect a neutrino signal.
To increase the likelihood of detecting astrophysical neutrinos, several groups have proposed water Cherenkov telescopes on much larger scales than AMANDA. There are three projects being considered for deployment in the deep Mediterranean water: ANTARES, NESTOR, and NEMO. Of these various efforts, the French ANTARES project is perhaps the furthest along. The ANTARES group has secured approval to build a telescope consisting of 1000 optical modules arranged on 13 vertical strings. The modules will deployed at water depths of 2100 to 2400 m, 30 km off shore. The baseline design for ANTARES calls for an effective collection area of approximately 0.1 km<sup>2</sup>. The neutrino energy threshold should be below 100 GeV which will permit the study of atmospheric neutrino oscillations, in addition to the search for neutrinos from astrophysical sources.
A very large km<sup>3</sup> size detector is being considered for deployment in the South Pole ice. The IceCube collaboration is proposing an array of 81 detector strings arranged on a square grid of 125 m string spacing . Each string would hold 60 optical modules with a vertical module spacing of 16 m. The effective neutrino energy for IceCube of approximately 400 GeV would be somewhat higher than for ANTARES, but IceCube will have a very large collection area ($`1`$km<sup>2</sup> 10 TeV) for neutrinos from astrophysical sources.
## 6 Summary
Astronomy using very high energy particles ($`\gamma `$-rays, cosmic rays, and neutrinos) is a diverse and rapidly developing field. It is currently a field largely driven by experimental results where significant scientific progress can be made on a relatively short time scale. For example, in the last few years there have been several exciting discoveries. Using $`\gamma `$-rays, we are probing remarkable and unexpected phenomena in objects such as active galaxies and gamma ray bursts. We are also searching for the origins of the cosmic radiation. At the very highest energies, we are discovering cosmic ray particles that probably should not be there, but are. Discoveries in this field often raise as many questions as they answer.
For the future, there will be an expanding interest in this field, both to understand astrophysics under extreme conditions and to search for evidence of physics beyond the standard models of elementary particles and cosmology. Future next-generation experiments in space and on the ground will greatly expand our discovery horizon. For $`\gamma `$-rays, the new projects include more powerful atmospheric Cherenkov telescopes and a new $`\gamma `$-ray satellite, GLAST. Larger air shower detectors, such as Fly’s Eye HiRes and Auger, will explore questions relating to the highest energy cosmic rays. In the world of neutrinos, the AMANDA experiment has demonstrated the ability to detect neutrinos produced in the atmosphere. Future experiments, such as IceCube and ANTARES, will greatly expand the possibility of detecting high energy astrophysical neutrino sources. We can only hope that our overall knowledge of high energy particles from the Universe advances at the same rate as new instruments are being constructed.
I wish to acknowledge the help and encouragement of many people in the particle physics and astrophysics communities. The contributions of the following people were particularly important: Katsushi Arisaka, Steve Barwick, Michael Catanese, Corbin Covault, Jim Cronin, Francis Halzen, Charles Jui, Tadashi Kifune, Peter Leonard, Eckart Lorenz, John Matthews, Masaki Mori, Gus Sinnis, David Smith, Simon Swordy, Steve Ritz, Pierre Sokolsky, Masahiro Takeda, Masahiro Teshima, Trevor Weekes, and Heinz Völk. I also thank the organizers of the Lepton-Photon Symposium (especially John Jaros, Helen Quinn, and Michael Peskin) for their encouragement and patience. Any inaccuracies are my fault alone. This research is supported in part by the National Science Foundation.
|
warning/0003/cond-mat0003482.html
|
ar5iv
|
text
|
# Anomalous diffusion associated with nonlinear fractional derivative Fokker-Planck-like equation: Exact time-dependent solutions
## Abstract
We consider the $`d=1`$ nonlinear Fokker-Planck-like equation with fractional derivatives $`\frac{}{t}P(x,t)=D\frac{^\gamma }{x^\gamma }\left[P(x,t)\right]^\nu `$. Exact time-dependent solutions are found for $`\nu =\frac{2\gamma }{1+\gamma }`$ ($`\mathrm{}<\gamma 2`$). By considering the long-distance asymptotic behavior of these solutions, a connection is established, namely $`q=\frac{\gamma +3}{\gamma +1}`$ ($`0<\gamma 2`$), with the solutions optimizing the nonextensive entropy characterized by index $`q`$ . Interestingly enough, this relation coincides with the one already known for Lévy-like superdiffusion (i.e., $`\nu =1`$ and $`0<\gamma 2`$). Finally, for $`(\gamma ,\nu )=(2,0)`$ we obtain $`q=5/3`$ which differs from the value $`q=2`$ corresponding to the $`\gamma =2`$ solutions available in the literature ($`\nu <1`$ porous medium equation), thus exhibiting nonuniform convergence.
PACS: 05.60.+w, 05.20.-y, 05.40.+j, 66.10.Cb
A great variety of diffusive problems in nature, namely those referred to as normal diffusion, are satisfactorily described by the Fokker-Planck linear equation
$$\frac{}{t}P(𝐱,t)=D^2P(𝐱,t)$$
(1)
where $`P(𝐱,t)`$ is the density of probability in the $`𝐱\{x_1,x_2,\mathrm{},x_d\}`$ space and $`D>0`$ is the diffusion coefficient. Such processes are currently characterized by the fact that $`𝐱^2t`$, as shown by Einstein in his celebrated 1905 paper on Brownian motion.
More recently, several works have focused on the same type of linear equation but with fractional derivatives. More precisely
$$\frac{}{t}P(𝐱,t)=D^\gamma P(𝐱,t)(\mathrm{}<\gamma 2)$$
(2)
where $`^\gamma _{i=1}^d\frac{^\gamma }{x_i^\gamma }`$ . Also, the nonlinear equation with ordinary derivatives has been focused on as well. More precisely
$$\frac{}{t}P(𝐱,t)=D^2\left[P(𝐱,t)\right]^\nu (\nu >1)$$
(3)
(no solutions are known for $`\nu 1`$ which are integrable ).
These two generalized Fokker-Planck equations have been used to study anomalous Lévy-like diffusion as well as correlated-like diffusive processes in porous media . The present paper addresses the unification of both equations as follows:
$$\frac{}{t}P(x,t)=D^\gamma \left[P(x,t)\right]^\nu (\mathrm{}<\gamma 2)$$
(4)
We will restrict ourselves to the $`d=1`$ case. More specifically, we are interested in normalized scaled solutions of the type
$$P(x,t)=\frac{1}{\varphi \left(t\right)}F\left[\frac{x}{\varphi \left(t\right)}\right]$$
(5)
Inserting this form into Eq.(4) (and, without loss of generality, setting $`D=1`$) we obtain:
$$\frac{\dot{\varphi }\left(t\right)}{\varphi \left(t\right)^2}\left[\frac{d}{dz}F\left(z\right)+zF\left(z\right)\right]=\frac{1}{\varphi \left(t\right)^{\nu +\gamma }}\frac{d^\gamma }{dz^\gamma }\left[F\left(z\right)^\nu \right]$$
(6)
where we have used the generic property
$$\frac{d^\delta }{dx^\delta }F\left(ax\right)=a^\delta \frac{d^\delta }{dz^\delta }F\left(z\right)(\delta )$$
(7)
with $`z=ax`$. This basic property holds not only for the ordinary derivative but also for all fractional operators we are aware of. By choosing the ansatz
$$\frac{\dot{\varphi }\left(t\right)}{\varphi \left(t\right)^{2\nu \gamma }}=k$$
(8)
where $`k`$ is an arbitrary constant, we obtain
$$\varphi \left(t\right)=\frac{1}{\left(k_1t+k_2\right)^{\frac{1}{\nu +\gamma 1}}}$$
(9)
$$\frac{d^\gamma }{dz^\gamma }\left[F\left(z\right)\right]^\nu =k\frac{d}{dz}\left[zF\left(z\right)\right]$$
(10)
with $`k_1\left(\gamma +\nu 1\right)k`$, $`k_2`$ being another arbitrary constant. Finally making an integration we obtain
$$\frac{d^{\gamma 1}}{dz^{\gamma 1}}\left[F\left(z\right)\right]^\nu =kzF\left(z\right)+c$$
(11)
where $`c`$ is another arbitrary constant.
Thus far it has not been necessary to specify the fractional operator we refer to. Indeed, several fractional generalizations exist for the ordinary derivative, namely the Riemann-Liouville (based on Laplace transform), Weyl (based on Fourier transform) and Caputo (also based on Laplace transform) ones. From now on we will use the Riemann-Liouville operator, since it is for this one that it has been possible to find new exact solutions. In this case we will work with the positive $`x`$ axis and, later on, we will use symmetry to extend the results to the entire real axis (we are working, in other words, with $`\frac{^\gamma }{|x|^\gamma }`$). Also, we will use the following generic result (see Appendix):
$$D_x^\delta \left[x^\alpha \left(a+bx\right)^\beta \right]=a^\delta \frac{\mathrm{\Gamma }\left[\alpha +1\right]}{\mathrm{\Gamma }\left[\alpha +1\delta \right]}x^{\alpha \delta }\left(a+bx\right)^{\beta \delta }$$
(12)
with $`D_x^\delta \frac{d^\delta }{dx^\delta }`$ and $`\delta \alpha +\beta +1`$. By defining $`g(x)x^{\frac{\alpha }{\nu }}\left(a+bx\right)^{\frac{\beta }{\nu }}`$ and $`\lambda \alpha \left(1\frac{1}{\nu }\right)\delta `$, and rearranging the indices, Eq.(12) can be rewritten as follows:
$$D_x^\delta \left[g\left(x\right)\right]^\nu =\frac{\mathrm{\Gamma }\left[\alpha +1\right]}{\mathrm{\Gamma }\left[\alpha +1\delta \right]}a^\delta x^\lambda g\left(x\right).$$
(13)
Using this property in Eq. (11) and, for simplicity, choosing $`c=0`$, we find
$`\alpha `$ $`=`$ $`{\displaystyle \frac{\left(2\gamma \right)\gamma }{12\gamma }}`$ (14)
$`\beta `$ $`=`$ $`{\displaystyle \frac{\gamma ^23\gamma +2}{12\gamma }}`$ (15)
$`\nu `$ $`=`$ $`{\displaystyle \frac{2\gamma }{1+\gamma }}.`$ (16)
These results allow us to write the solution in the form
$$P(x,t)=\frac{A}{\left(|k_1|t\right)^{\frac{\gamma +1}{\gamma ^2\gamma +1}}}\left[\frac{z^{\gamma \left(\gamma +1\right)}}{\left(1+bz\right)^{1\gamma ^2}}\right]^{\frac{1}{12\gamma }}$$
(17)
$$A=\left[k\frac{\mathrm{\Gamma }\left(\beta \right)}{\mathrm{\Gamma }(\alpha +1))}\right]^{\frac{1+\gamma }{12\gamma }}$$
(18)
$$z\frac{x}{\left(|k_1|t\right)^{\frac{\gamma +1}{\gamma ^2\gamma +1}}}$$
(19)
where $`b`$ is an arbitrary constant (to be taken, later on, as $`\pm 1`$ according to the specific solutions that are studied) and where, without loss of generality, we have set $`k_2=0`$ and $`a=1`$. Indeed, the $`k_2`$ constant can be incorporated into a shift of the origin of time, and $`a`$ can be incorporated into the normalization constant $`A`$. We also mention the exact solution
$$F\left(z\right)z^{\frac{\gamma }{\nu 1}}$$
(20)
that is not normalizable. Several regions will have to be considered, namely
$$\mathrm{}<\gamma <1,1<\gamma <0,\mathrm{\hspace{0.17em}0}<\gamma <\frac{1}{2},\frac{1}{2}<\gamma <1,\mathrm{\hspace{0.17em}1}<\gamma <2$$
(21)
We start with the region $`\mathrm{}<\gamma <1`$ for which, again without loss of generality, we can choose $`b=1`$. The normalization condition implies
$$A\underset{1}{\overset{1}{}}\left[\frac{z^{\gamma \left(\gamma +1\right)}}{\left(1z\right)^{1\gamma ^2}}\right]^{\frac{1}{12\gamma }}𝑑z=2A\frac{\mathrm{\Gamma }\left[\frac{\gamma ^2\gamma +1}{12\gamma }\right]\mathrm{\Gamma }\left[\frac{\gamma \left(\gamma 2\right)}{2\gamma 1}\right]}{\mathrm{\Gamma }\left[1\gamma \right]}=1$$
(22)
See Fig. 1. Also, as we can see from the limits of the different regions (21), we will have to consider different particular cases, namely $`\gamma =1,\gamma =0,\gamma =1/2,\gamma =1`$ and $`\gamma =2`$.
Let us start with $`\gamma =1`$ and arbitrary $`\nu `$. The corresponding equation is
$$\frac{}{t}P(x,t)=\underset{0}{\overset{x}{}}\left[P(y,t)\right]^\nu 𝑑y$$
(23)
To solve it let us go back to Eq. (10); after derivation with respect to $`z`$, we obtain
$$k\frac{d^2zF(z)}{dz^2}=\left[F\left(z\right)\right]^\nu $$
(24)
We are not going to treat this equation in detail; we rather limit ourselves to remark that the value $`\gamma =1`$ corresponds to $`\nu \pm \mathrm{}`$ in the curve (16). In the region $`1<\gamma <0`$ the probability density has a compact support, on the edges of which it diverges. See Fig. 1.
Let us now address the $`0<\gamma <1/2`$ region (where $`b=1`$). In the limit $`\gamma =0`$ ($`\gamma =1/2`$) integrability fails at infinity (at the origin); no such problems exist for $`0<\gamma <1/2`$. Normalization implies
$$A=\left[k\frac{\mathrm{\Gamma }\left(\frac{23\gamma +\gamma ^2}{12\gamma }\right)}{\mathrm{\Gamma }\left(\frac{1\gamma ^2}{12\gamma }\right)}\right]^{\frac{1+\gamma }{12\gamma }}=\frac{\mathrm{\Gamma }\left(\frac{1\gamma ^2}{12\gamma }\right)}{2\mathrm{\Gamma }\left(\gamma \right)\mathrm{\Gamma }\left(\frac{\gamma ^2\gamma +1}{12\gamma }\right)}$$
(25)
It is easy to show that $`P(x,t)`$ achieves a maximum (see Fig. 2) at
$$z=\frac{\gamma }{\left(12\gamma \right)}$$
(26)
Let us now address the $`\gamma =0`$ limit. It corresponds to the equation
$$\frac{}{t}P(x,t)=\left[P(x,t)\right]^\nu $$
(27)
that can be resolved analytically for arbitrary $`\nu `$. To obtain this solution it is convenient to go back to Eq.(10). It follows that
$$F\left(z\right)=\frac{B}{z}\left[1+cz^{1\nu }\right]^{\frac{1}{1\nu }}$$
(28)
Therefore, in the $`z\mathrm{}`$ (or, equivalently $`x\mathrm{}`$) limit, we have that $`F(z)1/z`$ if $`\nu >1`$ and $`F(z)`$ is a constant if $`\nu <1`$. The $`\nu =1`$ case needs specific discussion and we obtain $`F(z)z^{\frac{1}{k}1}`$. It is worthy reminding that the $`\gamma =0`$ solutions cannot be considered as distributions of probabilities since they are not normalizable.
The $`\gamma =1/2`$ limiting case corresponds to the following linear equation:
$$\frac{}{t}P(x,t)=\frac{^{\frac{1}{2}}}{x^{\frac{1}{2}}}P(x,t)$$
(29)
This equation can be solved by using the Laplace transform on both $`t`$ and $`x`$. It can also be solved by taking the limit $`\gamma 1/2`$. We have followed this procedure and, after tedious though straightforward calculations, we obtain $`k\frac{8}{3}(12\gamma )`$ and $`A[\frac{8\pi }{3}(12\gamma )]^{1/2}/2`$, which, replaced into Eq. (17), yield
$$P_{1/2}(x,t)\underset{\gamma \frac{1}{2}}{lim}P_\gamma (x,t)=\frac{1}{4\sqrt{\pi }t^2}\frac{\mathrm{exp}(\frac{t^2}{4x})}{(x/t^2)^{3/2}},$$
(30)
which is a distribution of the Poisson type. Let us stress that the above distribution can indistinctively be obtained by solving the fractional differential equation for $`\gamma =1/2`$, or by taking the $`\gamma >1/2`$ and the $`\gamma <1/2`$ solutions and then considering the $`\gamma 1/2`$ limit.
Let us now focus on the $`1/2<\gamma <2`$ region (where $`b=1`$). The solutions strictly vanish inside an interval which contains the origin. Outside this interval, the solutions are everywhere finite if $`1/2<\gamma <1`$, whereas they diverge if $`1<\gamma <2`$ (see Fig. 2).
The solutions corresponding to the $`1/2<\gamma <1`$ region are as follows:
$$P(x,t)=\frac{A}{\left(kt\right)^{\frac{\gamma +1}{\gamma ^2\gamma +1}}}\left[\frac{z^{\gamma \left(\gamma +1\right)}}{\left(1+z\right)^{1\gamma ^2}}\right]^{\frac{1}{12\gamma }}$$
(31)
Normalization implies
$$2A\underset{1}{\overset{\mathrm{}}{}}\left[\frac{z^{\gamma \left(\gamma +1\right)}}{\left(1+z\right)^{1\gamma ^2}}\right]^{\frac{1}{12\gamma }}𝑑z=2A\frac{\mathrm{\Gamma }\left[\gamma \right]\mathrm{\Gamma }\left[\frac{\gamma \left(\gamma 2\right)}{2\gamma 1}\right]}{\mathrm{\Gamma }\left[\frac{\gamma \left(\gamma +1\right)}{2\gamma 1}\right]}=1$$
(32)
from which $`A`$ is uniquely determined; finally, $`k`$ is obtained from $`A`$ by using Eq. (18). See Fig. 2.
Let us now focus the special case $`\gamma =1`$. In this case the equation becomes
$$\frac{}{t}P(x,t)=\frac{}{x}P(x,t)^\nu .$$
(33)
Its generic solution of the form indicated in Eq. (5)) is:
$$[F\left(z\right)]^\nu =kzF(x,t)+c,$$
(34)
which implicitly determines $`F(z)`$. The solution corresponding to $`c=0`$ is
$$F\left(z\right)z^{1/(\nu 1)}.$$
(35)
In the region $`1<\gamma <2`$ we have the same analytic solution that we had in the region $`1/2<\gamma <1`$ (i.e., Eq. (31)); however, at the point $`z=1`$, a divergence is now present (see Fig. 2). It is clear that this solution cannot be used without appropriate asymptotic considerations for $`\gamma =2`$, since for $`\gamma 2`$, $`\nu 0`$.
Let us now specifically focus on the possible $`x\mathrm{}`$ tail of $`P(x,t)`$ for arbitrary $`t`$. For $`\mathrm{}<\gamma <0`$ the support is compact, hence there is no tail. For $`0<\gamma <2`$, we obtain, using either Eq. (17) (for $`0<\gamma <1/2)`$ or Eq. (31) (for $`1/2<\gamma <2`$), the following asymptotic behavior:
$$P(x,t)\frac{1}{t^{\frac{\gamma +1}{\gamma ^2\gamma +1}}}z^{\frac{\gamma \left(\gamma +1\right)+\gamma ^21}{12\gamma }}\frac{t^{\frac{\gamma \left(\gamma +1\right)}{\gamma ^2\gamma +1}}}{x^{1+\gamma }}$$
(36)
We can easily verify that the exponent $`[\gamma \left(\gamma +1\right)]/[\gamma ^2\gamma +1]`$ monotonically increases from zero to $`\frac{2}{\sqrt{3}}+1`$ when $`\gamma `$ increases from zero to $`\frac{1+\sqrt{3}}{2}`$. Also, we verify that, for $`0<\gamma <1`$, both $`|x|`$ and $`x^2`$ diverge; for $`1<\gamma <2`$, only the latter does. For $`\gamma <0`$, all momenta are finite since the support is compact. Finally, in all cases, $`x`$ vanishes because of symmetry.
Let us now address the solutions for $`\nu =1`$ and arbitrary $`\gamma `$. The equation to be solved becomes
$$\frac{}{t}P(x,t)=\frac{^\gamma }{x^\gamma }P(x,t)$$
(37)
This equation can be solved using Laplace transform and the solutions are discussed in . Moreover, we can see (in Fig. 3), that the point $`(\gamma ,\nu )=(1/2,1)`$ is at the crossing of two solvable lines. Since solution (30) has been found as the limit of either one of these lines, it seems reasonable to conjecture that the same solution is found as a limit along any curve through that point.
Let us summarize the present work. We have addressed a generic Fokker-Planck-like diffusive equation, namely the one-dimensional case of Eq. (3), and have looked for exact scaled solutions of the type in Eq. (5). In the $`(\gamma ,\nu )`$ parameter space, the solutions corresponding to $`\gamma =2`$ and $`\nu >1`$ (porous medium equation) as well as to $`0<\gamma <2`$ and $`\nu =1`$ are available in the literature, as already mentioned. We are now exhibiting exact solutions along two new lines, namely the line $`\gamma =1`$ and arbitrary $`\nu `$, and the line indicated in Eq. (16) (i.e., $`\nu =(2\gamma )/(1+\gamma )`$ with $`\mathrm{}<\gamma <2`$). For the latter, we observe on Eq. (36) that the spatial asymptotic behavior is characterized by the exponent $`1+\gamma `$, which exactly coincides with that corresponding to Lévy superdiffusion. This is a remarkable result, since the present solutions concern a nonlinear Riemann-Liouville-fractional differential equation, ! and not the usual linear Fourier-fractional one, whose solutions are known to essentially be Lévy distributions. It would no doubt be interesting to know whether the same behavior is obtained no matter the value of $`\nu `$ ($`0<\gamma <2`$).
Let us finally mention a connection between the present problem and the solutions obtained from the optimization, under appropriate constraints (normalization and finite $`q`$-expectation value of $`x^2`$ in the interval $`(\mathrm{},\mathrm{})`$), of the nonextensive entropy $`S_q[1𝑑xp(x)^q]/[q1]`$. It has been shown that these optimizing distributions precisely coincide with the solutions of the present diffusive problem for $`\gamma =2`$. It comes out that $`q=2\nu `$ ($`\nu >1`$) . Along the line indicated in Eq. (16), the exact solutions of the entropic optimization problem and the present diffusive one do not coincide for arbitrary value of $`x`$. However, comparison of the $`x\mathrm{}`$ asymptotic behaviors is possible . Indeed, by identifying the behavior exhibited in (36) with the behavior $`1/|x|^{2/(q1)}`$ obtained for the entropic problem, we obtain
$$q=\frac{\gamma +3}{\gamma +1}(0<\gamma <2),$$
(38)
which, as commented above, precisely reproduces the connection established for Lévy distributions . By using Eq. (16), this relation can be rewritten as follows:
$$q=\frac{5+2\nu }{3}(0<\nu <2).$$
(39)
(We remind that the distributions for $`\nu >2`$ and $`\nu <1`$, i.e., $`\gamma <0`$, have compact support). The present nontrivial solution provides, for $`(\gamma ,\nu )=(2,0)`$, $`q=5/3`$, whereas the porous medium equation solution $`q=2\nu `$ provides $`q=2`$. This discrepancy exhibits that the point $`(\gamma ,\nu )=(2,0)`$ is a singular one, at least within the fractional derivative that we have adopted in this work.
Another point worth to be mentioned is that we have compared the present solutions with those optimizing $`S_q`$ defined in the interval $`(\mathrm{},\mathrm{})`$ and using finite $`q`$-expectation for $`x^2`$ (whereas the $`q`$-expectation value of $`x`$ vanishes). This is appropriate since, through symmetry, we have extended the solutions that we have found in the interval $`(0,\mathrm{})`$ to the entire real axis. Another possibility would of course be to compare the present results in the positive real semi-axis with those optimizing the entropy $`S_q`$ defined in the same semi-axis and using finite $`q`$-expectation value of $`x`$. If we did so, relation (39) would be replaced by $`q=(\gamma +2)/(\gamma +1)`$ (see also ).
Finally, last but not least, now that we have finished the presentation of the various exact solutions that emerged within the present work, it is worthy to mention that it would be very welcome the discussion of the stability of such solutions. More precisely, if we start at $`t=0`$ with an arbitrary distribution $`P(x,0)`$ and make it evolve through the present differential equation, what would be the $`t\mathrm{}`$ asymptotic distribution $`P_a(x,t)`$? ($`a`$ stands for attractor, in the space of the distributions). For instance, if the evolution is determined through the convolution product (i.e., a linear fractional-derivative Fokker-Planck equation, using a Fourier-based definition of fractional derivative), then the standard and the Lévy-Gnedenko central limit theorems apply, and consequently the attractor $`P_a(x,t)`$ is either a Gaussian or a Lévy distribution (respectively when the second cumulant is finite or infinite). If the evolution is instead det! ermined through a nonlinear integer-derivative Fokker-Planck equation like the one considered in , then $`P_a(x,t)`$ is given (as numerical verifications have shown) by the distributions which optimize the nonextensive entropy, where $`x`$ scales with a simple function of $`t`$. If the time evolution is obtained, as sometimes done, through recursive use of maps , $`P_a(x,t)`$ can present a variety of shapes depending on the specific map which is used. Finally, in our present case (nonlinear fractional-derivative Fokker-Planck equation using a Laplace-based definition of fractional derivative), the solutions we have found might well be $`P_a(x,t)`$. This point, however, deserves analysis on its own.
APPENDIX
We give a short review of the property of the fractional operator used to solve the nonlinear equation (4). The Riemann - Liouville operator is used in many applications of fractional calculus. The usual integral representation for this operator is:
$$D_x^\alpha f\left(x\right)=\frac{1}{\mathrm{\Gamma }\left(n\alpha \right)}\frac{d^n}{dx^n}\underset{0}{\overset{x}{}}\frac{dtf\left(t\right)}{\left(xt\right)^{\alpha +1}}(n1<\alpha <n).$$
(A-1)
For the calculations in this paper we have instead used the following equivalent form:
$$D_x^\alpha x^\rho =\frac{\mathrm{\Gamma }\left(\rho +1\right)}{\mathrm{\Gamma }\left(\rho +1\alpha \right)}x^{\rho \alpha }$$
(A-2)
We have also used the generalized Leibnitz formula for this kind of fractional derivative, namely
$$D_x^\alpha \left[f\left(x\right)g\left(x\right)\right]=\underset{n=0}{\overset{\mathrm{}}{}}\left(\genfrac{}{}{0pt}{}{\alpha }{n}\right)D_x^{\alpha n}\left[f\left(x\right)\right]D_x^n\left[g\left(x\right)\right]$$
(A-3)
Let us also mention that, for this operator, the following property holds under Laplace transform:
$$L\left[D_x^\alpha f\left(x\right)\right]=s^\alpha F\left(s\right)f^{\left(\alpha 1\right)}\left(0\right)$$
(A-4)
for $`\alpha 1`$ and where $`f^{\left(\alpha 1\right)}\left(0\right)`$ means fractional derivative calculated in $`x=0`$.
Let us now show how formula (A-3) leads to Eq. (12), that is used in the text. By assuming no restrictions on the parameters $`\alpha `$, $`\beta `$ and $`\gamma `$ and applying the generalized Leibnitz rule to the function $`x^\alpha \left(a+bx\right)^\beta `$ we obtain
$$D_x^\delta \left[x^\alpha \left(a+bx\right)^\beta \right]=\underset{n=0}{\overset{\mathrm{}}{}}\left(\genfrac{}{}{0pt}{}{\delta }{n}\right)D_x^{\delta n}\left[x^\alpha \right]D_x^n\left[\left(a+bx\right)^\beta \right].$$
(A-5)
After some algebra we obtain
$$D_x^\delta \left[x^\alpha \left(a+bx\right)^\beta \right]=\underset{n=0}{\overset{\mathrm{}}{}}\left(\genfrac{}{}{0pt}{}{\delta }{n}\right)\frac{\mathrm{\Gamma }\left(\alpha +1\right)}{\mathrm{\Gamma }\left(\alpha +1\delta +n\right)}x^{\alpha \delta +n}\frac{\mathrm{\Gamma }\left(\beta +1\right)}{\mathrm{\Gamma }\left(\beta +1n\right)}\left(1\right)^nb^n\left(a+bx\right)^{\beta n}.$$
(A-6)
A closed form for this series can be achieved if $`\alpha +\beta +1=\delta `$. Indeed, by using the Gamma function property $`\mathrm{\Gamma }\left(z\right)\mathrm{\Gamma }\left(1z\right)=\frac{\pi }{\mathrm{sin}\pi z}`$ we obtain
$$D_x^\delta \left[x^\alpha \left(a+bx\right)^\beta \right]=a^\delta \frac{\mathrm{\Gamma }\left(\alpha +1\right)}{\mathrm{\Gamma }\left(\alpha +1\delta \right)}x^{\alpha \delta }\left(a+bx\right)^{\beta \delta },$$
(A-7)
which essentially is Eq. (12).
For completeness, it is worthy to briefly mention here a recent variation of Riemann-Liouville operator that we have mentioned in the text, namely, the Caputo derivative. Its definition is:
$$^CD_x^\alpha f\left(x\right)\frac{1}{\mathrm{\Gamma }\left(m\alpha \right)}\underset{0}{\overset{x}{}}\frac{dtf^{\left(m\right)}\left(t\right)}{\left(xt\right)^{\alpha +1m}}.$$
(A-8)
($`C`$ stands for Caputo). The main advantage with respect to the Riemann-Liouville operator is that Caputo derivative of a constant is zero, which is not the case of the Riemann-Liouville one. Substantially, this kind of fractional derivative is a formal generalization of the integer derivative under Laplace transform. As disadvantage, it exhibits the fact that, whenever the derivation index is an integer number, it recovers the usual derivative excepting for an additive constant, whereas the Riemann-Liouville operator has no such disagreable property.
Finally, let us also mention the definition of Weyl fractional derivative. It is based on the properties of Fourier transform, and it is defined as follows:
$$^WD_x^\alpha f\left(x\right)=\underset{k=\mathrm{}}{\overset{+\mathrm{}}{}}(ık)^\alpha c_k\mathrm{exp}(ıkx),$$
(A-9)
its continuum version being
$$^WD_x^\alpha f\left(x\right)=\frac{1}{2\pi }_{\mathrm{}}^+\mathrm{}d\omega (ı\omega )^\alpha \widehat{f}\left(\omega \right)\mathrm{exp}(ı\omega x).$$
(A-10)
($`W`$ stands for Weyl).
|
warning/0003/astro-ph0003145.html
|
ar5iv
|
text
|
# Near and mid-IR sub-arcsecond structure of the dusty symbiotic star R Aqr
## 1 Introduction
R Aqr is a mass-losing long-period variable (LPV) star believed to be in a symbiotic system with an obscured hot companion whose presence is betrayed by nebular emission lines. Long noted for its peculiar elliptical visible nebulosity extending up to 2’ (first reported in Lampland 1923), this star has recently been subject to intense scrutiny with the latest generation of astronomical instrumentation both on the ground and in space. Much of this attention can be traced to the detection in the optical (Wallerstein (1978); Herbig (1980)) and radio (Sopka et al. (1982)) of what is believed to be the nearest astrophysical jet to the earth ($`200`$ pc; Van Leeuwen Feast Whitelock & Yudin (1997); Hollis, Pedelty & Lyon (1997)). Thought to originate in the accretion disk around an unseen hot sub-dwarf, this jet has been studied extensively at UV, optical and radio wavelengths (e.g. Burgarella, Vogel & Paresce (1992); Paresce & Hack (1994); Lehto & Johnson (1992)).
A difference in location of the peak intensity of the $`\nu =1,J=10`$ SiO maser line and the nearby lineless continuum from 7 mm VLA maps was reported by Hollis, Pedelty and Lyon (1997). These authors interpreted the SiO peak as the location of the LPV and the continuum peak marking the accretion disk of the companion, and were thus able to derive a binary separation of $`55\pm 2`$ mas with a position angle of $`18^{}\pm 2^{}`$.
With its ability to penetrate layers of obscuring dust, the infrared might be thought an ideal wavelength to image the inner regions of this system, from photospheres of the stars to dust and jets in their immediate surroundings. However the high angular resolutions required and the presence of the enormously luminous LPV have limited the efficacy of this approach. Infrared interferometric measurements with the IOTA array (Van Belle et al. (1996)) yielded photospheric diameter measurements of $``$14–15 mas for the LPV but no evidence for an additional component. Speckle observations at J band (Karovska MCCarthy & Christou (1994)) suggestive of an elongated image were interpreted in terms of a distorted inner dust shell, possibly due to the presence of the companion. Anandarao & Pottasch (1986) were able to match the near- to far-IR fluxes using a two-temperature dust shell model, while more recently Le Sidaner & Le Bertre (1996) were able to fit similar spectral energy distributions with a simple uniform outflow model.
In the following section of this paper, we briefly describe our observational methods and apparatus for obtaining near- and mid-IR interferometric measurements. In Section 3 the experimental results are given, together with discussion of model fitting and physical interpretation. A brief summary of our conclusions is found in Section 4.
## 2 Observations
A brief synopsis of the observing techniques used to secure the near- and mid-IR interferometric measurements is given below. Near-IR data was taken at the Keck-I telescope in 1998 July and 1999 January covering 4 narrowband wavelengths between 1.25 and 3.08 $`\mathrm{\mu m}`$. Visibility data at 11.15 $`\mathrm{\mu m}`$ were taken with the Infrared Spatial Interferometer (ISI) over the period 1992 to 1999 with various baselines sampling the visibility function.
### 2.1 Near-Infrared Interferometry
Near-IR observations at the Keck-I telescope have utilized the technique of aperture masking interferometry, in which a mask is placed over the telescope pupil so as to only pass light from selected regions, in effect transforming the telescope into an array of small subapertures. Starlight passing through the mask, in this case a 21-hole non-redundant configuration, was brought to a focus in the Near InfraRed Camera (Matthews & Soifer (1994); Matthews et al. (1996)), a $`256\times 256`$ InSb array with a pixel scale of 20.57 milli-arcsec/pixel. At focus, an interference pattern is formed containing sets of fringes at various spatial frequencies corresponding to baselines in the pupil. Subsequent analysis of datasets consisting of 100 rapid-exposure ($`140`$ msec) frames using Fourier techniques enables recovery of the visibility amplitudes and closure phases corresponding to the complex visibility function of the object. With the exception of the sparse telescope pupil which has been shown to confer signal-to-noise advantage for bright targets, the observational and data processing techniques used were very similar to those utilized in speckle interferometry (for a review, see Roddier (1988)).
After calibration utilizing nearly-contemporaneous observations of nearby point-source stars, the visibility data could be interpreted in a number of ways, including the fitting of model brightness distributions. The recovery of diffraction-limited maps was also possible with the help of self-calibration methods such as the Maximum-Entropy Method (Gull & Skilling (1984); Sivia (1987)) or CLEAN algorithm (Högbom (1974)). The signal-to-noise advantages of sparse-aperture observations for bright sources have been exploited previously (e.g., Baldwin et al. (1986); Haniff et al. (1987); Roddier (1988)), while a more detailed description of this particular apparatus and technique may be found in Monnier et al. (1999) and Tuthill et al. (1999a). An observing log of near-IR observations is given in Table 1 showing the dates, filters and stellar phases (Mattei (1999)) pertaining to our two sets of measurements in 1998 June and 1999 January.
### 2.2 Mid-Infrared Interferometry
Our mid-IR visibility data were obtained at 11.15 $`\mathrm{\mu m}`$ with the U.C. Berkeley ISI, a two-element heterodyne stellar interferometer located on Mt. Wilson, CA. Both telescopes are mounted within movable semi-trailers which allowed periodic reconfiguration of the baseline from 4 to 32 m over the course of these measurements. Detailed descriptions of the apparatus and recent upgrades can be found in Bester et al. (1990; 1994), Lipman (1998) and Hale et al. (1999). A journal of R Aqr observations taken with the ISI is provided in Table 2, which shows observations broken into seven separate observing epochs labeled for convenience ‘A’ through ‘G’. As individual visibility measurements taken with the ISI can have low signal-to-noise, we have averaged together measurements taken over these periods with the result that the observing parameters (such as sky position angle or stellar variability phase) pertinent to any single datum can be traced back to the range of values given in Table 2. Observations of K giant stars $`\alpha `$ Tau and $`\alpha `$ Boo were utilized to monitor system visibility drifts which could be as large as 15% from year to year, thus ensuring reliable calibration (to better than 5%) of science targets as is discussed at some length in Danchi et al. (1990; 1994)
## 3 Results and Discussion
### 3.1 Near-Infrared Visibility Measurements
The principal aim of the near-IR program was to attempt to detect the presence of hot circumstellar material, either in the immediate environment of the LPV or associated with the symbiotic companion. The full two-dimensional Fourier coverage out to 10 m baselines afforded by the masking experiment has proved highly successful at recovering asymmetric circumstellar structures in other systems (e.g., Monnier et al. (1999); Tuthill et al. 1999b ). However, in the case of R Aqr, the optical depth to the star is relatively low in the near-IR ($`\tau =`$ 0.01 at 2.2 $`\mu `$m), so that the star greatly outshines the contribution from the dust shell. No circumstellar features were detected at any wavelength in the near-IR from J through K bands, and we are able to place an upper limit of around $`\mathrm{\Delta }M`$$`\stackrel{>}{_{}}`$ 5 magnitudes for the relative brightness in the near-IR of any such companion.
Failure to detect a second star may seem surprising in view of previous indications of a hot companion. Resolution of the interferometry is about 15 mas (FWHM) and the projected distance to a second star might possibly be less than this. However, Hollis et al. (1997) reported indications in 1996 of a companion separated by 55 mas, or 11 AU assuming the distance to R Aqr is 200 pc. The orbital period was estimated to be 44 years, so that its distance during the present measurements would have to be comparable with 55 mas and easily resolved. However, a near-IR intensity less than 1% of that of the R Aqr Mira is in fact not unreasonable. The LPV has a diameter of 15 mas, or 3 AU, whereas a hot star even as large as our Sun would have a diameter at least 100 times smaller, and even allowing for its higher irradiance, a factor of $`10^3`$ less flux should result in the near-IR.
Although the stellar disk of R Aqr itself was not well resolved with the baselines available, its size was determined by fitting the available data. Figure 1 displays visibility curves for four different colors and two observing epochs (see Table 1). In order to increase the signal-to-noise in the plotted data, the two-dimensional visibility function for each wavelength has been azimuthally averaged, allowing visibilities to be plotted as a function of baseline length. On this occasion, additional correction had to be made for the size of the calibrator star 30 Psc, which was itself a late-type (M3III) variable expected from effective temperature and photometric arguments to present an angular diameter of 7.52 mas (Ochsenbein & Halbwachs (1982)). Also overplotted on Figure 1 are the best-fitting circular uniform disk model profiles with diameters given in the figure key. As short baseline visibilities can suffer from poor calibration due to seeing changes (the well-known ‘seeing spike’ problem c.f. Tuthill et al. 1999a), we have chosen to fit data only at spatial frequencies higher than $`2\times 10^5`$ rad<sup>-1</sup>.
It should be noted that over the wavelength range 1.236 – 2.260 $`\mu `$m (upper four sets) the stellar diameter is really at the limit of our detection. Visibility curves have only dropped some $`1020`$% as compared to an unresolved source. Although the formal diameter errors to the plotted data are small, the larger errors given in Figure 1 reflect systematic sources of uncertainty due to seeing-related mismatches between the source and point calibrator measurements. The likely spread of this dominant term in the uncertainty was determined by examination of large volumes of additional calibration data taken with the same observing parameters. Asymmetries in the errors are a consequence of the nonlinear relationship between visibilities and angular diameters – symmetric error bars on a visibility datum will result in asymmetric diameter errors for such barely-resolved targets. Despite these large uncertainties, we find good agreement with the published K band angular diameters of van Belle et al. (1996) of 14.95 and 14.06 mas in 1995 Jul 11 and 1995 Oct 07. We note also, in passing, that our 1999 January 2.26 $`\mu `$m diameter taken after minimum light is larger than the 1998 June observation near to maximum. Such cycle-dependent size behavior has been studied by Burns et al. (1997) for the case of R Leo, and although our enlargement is in rough accord with this earlier work, our large errors and poor phase sampling cast doubt on the tentative detection of such dynamical changes here.
The masking experiment at the Keck telescope recorded full two-dimensional Fourier information out to the maximum possible $``$10 m baselines, and in the preceding discussions we have restricted our attention only to the azimuthally averaged visibilities. Departures from circular symmetry would be manifest in the visibilities as a modulation with position angle, while departures from inversion symmetry are indicated by non-zero closure phase signals. As R Aqr is barely resolved, second-order terms in the visibility function were hard to extract given the random and systematic errors on the visibility data. However at all near-IR wavelengths of observation, the visibilities were consistent with a circular disk, with the upper limits on any ellipticity given by the error bars to the circular disk fits given in Figure 1. With the exception of the 3.08 $`\mu `$m observations mentioned below, closure phase signals were within the error bars of zero, implying no significant departures from inversion symmetry in the data.
Images of the R Aqr system in the J band were recently reported by Karovska, McCarthy and Christou (1994) from 1991 November speckle observations with the MMT (effective aperture 6.86 m). These authors found R Aqr to be dramatically elongated, presenting a size of 60 mas along a position angle of 140 while being unresolved in the orthogonal direction. Unfortunately, there is no evidence for such a signal in our data. Our 1.236 $`\mu `$m data rule out elongations greater than 15 mas in any direction, and we note that the expected signal from a 60 mas elongation would cause our Keck visibility function to pass beyond the first null – a radical departure from the observations. Unless profound secular changes in the brightness profile occurred between 1991 and 1999, such elongations appear to be ruled out and we can see of no way to reconcile these two datasets. Full determination of source structure must await the coming generation of separated-element imaging interferometers working at substantially higher angular resolutions.
It is apparent, from Figure 1, that the 3.08 $`\mu `$m data are qualitatively quite different from the shorter wavelength sets, with significantly larger ($``$34 mas) best-fit uniform disk diameters. Such dramatic enlargement has been reported before, and seems to be most pronounced for objects of extremely late spectral type (Tuthill et al. 1999c ). Furthermore, at both 1998 June and 1999 January epochs, the visibilities can be seen to be a poor match to the uniform disk profile, implying instead that a significant proportion of the flux originates in an extended halo. Simple geometrical models, consisting of a uniform disk plus an extended Gaussian shell, were investigated to see if the 3.08 $`\mu `$m visibility functions could be fit. Good fits to both epochs of data were found for models in which $``$20% of the flux came from a Gaussian halo of $``$100 mas FWHM. Following earlier workers in the optical (Labeyrie et al. (1977); Scholz & Takeda (1987); Haniff, Scholz & Tuthill (1995)), changes of angular diameter and radial profile with wavelength can be attributed to the reprocessing of radiation at different levels in the atmosphere by molecular layers whose opacity characteristics are strongly wavelength-dependent. In contrast to the abundance of prominent molecular features such as TiO and VO in optical spectra, recent ISO spectra of a number of late M-type stars (Tsuji et al. (1997)) did not show strong spectral structures near 3.1 $`\mu `$m, although there was some unaccounted absorption in this region. Water ice shows an extremely strong absorption at 3.1 $`\mu `$m, however this would seem to be ruled out on physical grounds and further investigation of the possible blanketing effects of common molecules such as CO and H<sub>2</sub>O is needed. An alternative scenario is that this extended halo arises from emission from the innermost hot circumstellar dust, just beginning to make its presence felt as we move towards longer wavelength.
As a final note concerning the 3.08 $`\mu `$m data, the closure phase signals did exhibit a small (few degrees) departure from zero as would be expected if the source had a non centro-symmetric brightness distribution. One way the phase signals could be modeled was to allow the $``$100 mas FWHM Gaussian halo (mentioned above) to move 10 mas to the southwest with respect to the star. Although this identification of non-centro symmetric elements must be labeled as tentative, it is hoped that with the longer baselines available to separated-element interferometers, any such complexity in this source will be subject to careful scrutiny.
### 3.2 Mid-Infrared Visibility Measurements
Visibilities at 11.15 $`\mu `$m, recorded over the seven observing epochs (‘A’ through ‘G’ from Table 2), are plotted in Figure 2. Although an interferometer baseline is a vector quantity, visibilities are shown as a function of the (scalar) baseline length. This approach is often adopted in cases such as this where coverage of the UV plane is sparse so as to make full two-dimensional modeling of the visibility function unwarranted.
Data are separated into three categories according to the stellar phase $`\mathrm{\Phi }`$ of the LPV at the time of observation (Mattei (1999)); $`\mathrm{\Phi }>0.875`$ or $`\mathrm{\Phi }<0.125=`$ ‘Maximum’; $`0.375<\mathrm{\Phi }<0.625=`$ ‘Minimum’; with all other data considered ‘Intermediate’. Although the spatial frequency coverage is far from uniform for all stellar phases, the general shapes of the visibility curves are broadly similar at all phases, and they exhibit two clear components. The relatively rapid drop in visibility at low spatial frequencies argues for a resolved dust shell component, while the flattening off towards high spatial frequencies is indicative of a compact (stellar) component contributing some 30$``$40% of the flux.
The physical parameters of the R Aqr circumstellar environment were modeled using radiative transfer computations of simulated stellar dust shells attempting to match the visibility data, together with published spectrophotometry. The modeling code used for this purpose, based on the work of Wolfire & Cassinelli (1986), calculates the equilibrium temperature of the dust shell as a function of the distance from the star (assumed to be a blackbody). A wide range of optical properties of dust grains and density distributions of dust shells can be modeled, as is described in more detail in Danchi et al. (1994) and Monnier et al. (1997). As R Aqr has an oxygen-rich atmosphere, it is appropriate to use optical constants for astronomical silicates while the grain size distribution follows that of Mathis, Rumpl & Nordsieck (1977) with dust opacities calculated assuming spheroidal Mie scattering using a method developed by Toon & Ackerman (1981). In the absence of detailed multi-wavelength polarization measurements, this standard model was favored over more complex scenarios such as those with elongated grains, which should have little effect on the main conclusions of the paper except perhaps to produce somewhat different optical depths in model fits. Radiative transfer calculations at 67 wavelengths allow the wavelength-dependent visibility curves, the broadband spectral energy distribution, and the mid-infrared spectrum all to be computed for comparison with observations. Dense sampling of the silicate band in the 10 $`\mathrm{\mu m}`$ region is particularly useful for comparison with published spectrophotometry. Figure 3 gives near- and mid-IR spectrophotometric data taken from various literature sources for comparison with model spectra. The optical constants of Draine & Lee (DL 1984), Ossenkopf, Henning, & Mathis (OHM 1992), and David & Papoular (DP 1990), were all tested in modeling the dust shell, with the result that OHM (used hereafter) and DP constants both provided high-quality fits to the shape of the silicate feature, while DL optical constants were not preferred.
For model calculations, the simplifying assumption of uniform isotropic outflow has been made, resulting in a spherical dust shell with a $`\rho r^2`$ density distribution beginning at a discrete dust condensation radius. Although there is considerable evidence for departures from spherical symmetry at larger scales, the primary constraints on the models from the mid-IR visibilities and spectra were not extensive enough to warrant more complex models.
Two models were constructed, one at stellar maximum and one at minimum and their computed mid-IR visibility functions are given in Figure 2, while simulated spectra and physical parameters are in Figure 3. The diameter of the central star was fixed at 16 mas (or 3.2 AU assuming a distance of 200 pc) from near-IR observations (Section 3.1; Van Belle et al. (1996)) and the temperature fixed to match the mid-IR luminosity. As can be seen from Figure 2, the visibility data at maximum and minimum light are well fit by the model curves, and we note that all parameters of the dust shell model are identical for the two models, with the increase in stellar flux being simulated by an elevation in the effective stellar temperature of the star by 500 K. This in turn leads to an elevation of the temperature of the dust at the inner radius from 850 K to 1170 K.
Serious inadequacies in this simplest-case model scenario are revealed upon comparison of model predictions with spectro-photometry and intermediate-phase visibility data. As can be seen from Figure 3, the simulated spectral energy distributions are not well matched to the observations. In the mid-IR, although the overall flux level and the shape of the model silicate feature are in fair accord with the measurements of Monnier et al. (1998), the spectral slope towards longer wavelengths is too steep. This lack of long-wavelength flux argues for the presence of more cool dust, or a cooler underlying stellar spectrum. Near-IR photometric measurements (Le Bertre (1993); Kamath & Ashok (1999)) fall well below the models implying increased near-IR opacity, or again a cooler underlying star; our model overestimates the bolometric luminosity. However, such decreases in stellar temperature and increases in density of the circumstellar dust shell were not possible without seriously perturbing the fits to the visibility data, in particular the relatively high visibilities recorded at long baselines.
Put simply, the strong point-source (30$``$40%) evident from the mid-IR visibility curves drives models to low optical depths and in turn high stellar luminosities to account for the mid-IR flux (which could otherwise come from the dust if the optical depth were higher). Although a wide range of model parameter space was explored, including variations on the density-radius law, the inner and outer radii, the dust opacity properties, and the dust densities, no spherically symmetric model could be found which overcame the shortcomings with the spectral fitting. The reader is therefore cautioned against taking the model parameters from Figure 3 as an accurate description of this star; rather they describe the best fit to our interferometric dataset which can be provided by a uniform-outflow spherical dust shell.
There are many imperfections in the modeling process, mostly due to oversimplifications or limited knowledge, which could lead to misfitting of the data. Although the use of a blackbody spectrum for the star is common practice, evolved stars exhibit significant departures (see for example, Lobel Doyle & Bagnulo (1999)). The optical properties of the dust can also have a significant impact on the radiative transfer results, and detailed knowledge of crucial aspects such as the true grain size distribution are lacking (the dust of Mathis, Rumpl & Nordsieck (1977) is more appropriate for the ISM). However, it may be that not all difficulties in fitting may stem from the fact that the simple dust radiative transfer models, with available silicate opacity functions do not give a realistic picture of actual circumstellar dust shells.
Examination of the intermediate-phase visibilities (bottom panel, Figure 2) points to an additional possible explanation for the poor spectral performance of our models. While visibilities from data set ‘A’ fall between the Max and Min model curves as might be expected, those from data set ‘F’ are higher than either, especially at low spatial frequencies well covered in the other plots. A likely explanation for this change in visibility can be found by examination of Table 2: most of the short baseline data underpinning the models (sets ‘C’, ‘D’ and ‘G’) along with set ‘A’ were taken at a position angle on the sky of around $`70100^{}`$, while the baselines constituting data set ‘F’ were taken at a very different orientation ($`145^{}`$ for the short-baselines). Hence the difference in these visibilities taken at angles differing by $`4575^{}`$ argues for substantial departures from sphericity. In such an anisotropic circumstellar environment, the spherical model assumptions break down and failure to fit to the spectral energy distribution may be the result of complications such as patchy extinction of the star or radiation leaking through holes in the cloud. Indeed, the extension of asymmetric structure from the well-known nebula down to successively smaller scales is not surprising, and has already been suggested from findings of time-variable polarization (Deshpande et al. (1987)).
The visibility data set ‘F’ has rather dense sampling over a range of spatial frequencies together with small probable errors, and as is seen from Figure 2, it exhibits wiggles or bumps (e.g. at $`6\times 10^5`$ rad<sup>-1</sup>) which indicate departure from a simple smooth radial intensity distribution. In order to assess this further evidence for anisotropy in the R Aqr dust shell, model brightness distributions (not, in this case, based on radiative transfer computations) have been fit, the result of which is overplotted as a dotted line in panel 3 of Figure 2. The stellar component of our best-fit model contributed 36% of the flux, with 61% coming from a spherical circumstellar dust shell assumed to have a simple Gaussian profile with a FWHM of 350 mas. A third component, modeled as a localized point-source contributing 3% of the total flux and located 700 mas from the star at a position angle of $`100^{}`$ (or $`80^{}`$), produces the wiggles noted in the data. Although such a feature might be some localized concentration of dust, or possibly a local warming, serious discussion should be deferred until more complete Fourier coverage can be obtained. With the sparse coverage afforded by the single interferometer baseline, a host of different models could be devised in which some additional flux is originating at some distance from the star and thereby generating the high-frequency signal we see. Therefore the likely significance to be placed upon our model parameters is low. This discussion is included simply as additional evidence for departure from a uniform outflow, and as an encouragement to further efforts at full high angular resolution imaging of this system in the mid-IR.
When non-spherical dust halos occur, it is rarely possible to discriminate between a number of viable outflow models when only relatively limited Fourier data are available, even though reasonable models may be chosen (e.g. Lopez et al. (1997)). The simple isotropic outflow model given here provides a reasonable fit to most of our mid-IR visibility data. However, deviations from the model are clear, and the more involved task of mapping and fitting of spectral data await the recovery of the full two-dimensional complex visibility function. For this reason, a detailed comparison with other models such as the two-shell model of Anandarao & Pottasch (1986) or the uniform outflow model of Le Sidaner & Le Bertre (1996) is not given here. However, our inner dust shell radius is in agreement with results of these earlier workers (for Anandarao & Pottasch (1986), we consider only their hot inner shell which should dominate over the cool outer shell ($`T_2=87.1`$ K) in this wavelength region). On the other hand, significant differences include lower optical depths and hotter stellar temperatures.
All spherically symmetric models investigated have been found to be inadequate in fitting the expanded dataset encompassing near- and mid-IR interferometry and the spectral energy distribution from the near- to the far-IR. This finding is in agreement with studies of $`o`$ Ceti, another symbiotic system observed interferometrically in the mid-IR Lopez et al. (1997). With further work at high resolutions, particularly directed towards full imaging of these systems, the effects of the companion’s presence, both gravitational and radiative, in the shaping of the circumstellar dust shell may be elucidated.
## 4 Conclusions
High-resolution interferometric studies of R Aqr at narrow bandwidths within the J, H and K bands find no evidence for any significant departure from the best-fitting model of a marginally-resolved stellar disk, which we identify with the LPV component of the system. Any companion present and separated by more than $`15`$ mas must exhibit a magnitude difference in excess of $`\mathrm{\Delta }M`$$`\stackrel{>}{_{}}`$ 5 mag in the near infrared. These results are consistent with those of Van Belle et al. (1996), but not with Karovska et al. (1994). An enlargement, by approximately a factor of 2, is reported for the apparent size at a wavelength of 3.08 $`\mathrm{\mu m}`$ as compared with shorter near-infrared bands, which is attributed to molecular blanketing by an unidentified species in the atmosphere, or thermal emission from material in an extended halo. Tentative evidence for asymmetry at this wavelength is also reported. In the mid-IR, visibility data obtained with the ISI have been used to constrain simple uniform spherical outflow models. Using self-consistent radiative transfer calculations, good fits to most of the visibility data have been obtained. However, serious shortcomings were found when comparing the synthetic spectral energy distributions to measurements, many of which may be attributed to inadequate knowledge of the detailed dust shell parameters. Furthermore, mid-IR interferometry obtained at a different position angle on the sky points towards substantial departures both from spherical symmetry and from a simple radial distribution. These indicate that the shortcomings of the present models are probably in large part be due to the oversimplistic assumption of an isotropic uniform outflow. Further high-resolution studies, both in near-infrared with higher spatial resolution and dynamic range, and in the mid-infrared with more complete Fourier coverage, are important in order to fill in the gaps in understanding of this star.
We thank Everett Lipman for help in securing many of the ISI measurements published herein, and for developmental work on the instrument. Long-baseline interferometry in the mid-infrared at U.C. Berkeley is supported by the National Science Foundation (Grants AST-9315485, AST-9321289, AST-9500525, & AST-9731625) and by the Office of Naval Research (OCNR N00014-89-J-1583 & FDN0014-96-1-0737). Some of the measurements herein were obtained at the W.M. Keck Observatory, made possible by the generous support of the W.M. Keck Foundation, and operated as a scientific partnership among the California Institute of Technology, the University of California, and NASA.
|
warning/0003/math0003145.html
|
ar5iv
|
text
|
# Study on the family of 𝐾3 surfaces induced from the lattice (𝐷₄)³⊕⟨-2⟩⊕⟨2⟩
## 0 Introduction
In 1992 K.Matsumoto, T.Sasaki and M. Yoshida studied the period mapping for a family of $`K3`$ surfaces of type $`(3,6)`$, that is the family of double sextic surface over $`𝐏^2`$ ramified along 6 lines in general position, and Matsumoto gave the description of the inverse mapping in terms of theta constants. It gives the modular mapping for the 4 dimensional Shimura variety in the Siegel upper half space $`_4`$ derived from the family of 4-dimensional abelian varieties with generalized complex multiplication by $`\sqrt{1}`$. So we call it MSY modular mapping. Shiga showed an arithmetic application of MSY modular mapping in . This story is the consequence of the eventual coincidence of 2 different bounded symmetric domains between $`D_{IV}^4`$ and $`I_{2,2}`$. There are a few (finite) such exceptional coincidences. The highest one is the (analytic) equivalence between $`D_{IV}^6`$ and $`D_{II}(4)`$ ( in terms of Lie algebra $`so(2,6;𝐑)so(4,𝐇)`$) and it contains the above coincidence of MSY case.
Our present study is the first step to get the extended model of the MSY modular mapping using this equivalence.
Let us consider the rank 14 lattice $`P=D_4^322`$. We define a $`K3`$ surface $`S`$ of type $`P`$ with the property that $`P\mathrm{Pic}(S)`$ (see Definition 2.1), where $`\mathrm{Pic}(S)`$ indicates the Picard lattice of $`S`$. In this article we study the family of $`K3`$ surfaces of type $`P`$ with a certain fixed multi-polarization. And we do not discuss the representation of the inverse of the period mapping.
We mention that our family is already appeared in the work ( in section 7) and they obtained the differential equation coming from this family standing on a different view point.
Throughout this article we work on the field $`𝐂`$.
In Section 1 we study the family $``$ of double covering surfaces over $`𝐏^1\times 𝐏^1`$ branching along 4 bidegree $`(1,1)`$ curves imposed with a certain generality condition. The element of $``$ is called a double $`4H`$ surface ( see Definition 2.2). Such a surface is given in the affine form
$$S=S(x):w^2=\underset{k=1}{\overset{4}{}}(x_1^{(k)}st+x_2^{(k)}s+x_3^{(k)}t+x_4^{(k)}),$$
where we use the notation
$$x_k=\left(\begin{array}{cc}x_1^{(k)}& x_2^{(k)}\\ x_3^{(k)}& x_4^{(k)}\end{array}\right)M(2,𝐂).$$
We set up the view point that a general member $`S`$ of $``$ to be an elliptic fibred surface (Proposition 1.1). It becomes to be a $`K3`$ surface. Then we construct the basis of the transcendental lattice $`\mathrm{Tra}(S)`$ of a general member $`S`$. Fortiori we know the structures of the Picard lattice $`\mathrm{Pic}(S)`$:
### Theorem 1
For a general member $`S`$ of $``$ it holds
$$\mathrm{Tra}(S)U(2)^22^4.$$
By making its orthogonal complement in the $`K3`$ lattice $`L`$ we get the Picard lattice
$$\mathrm{Pic}(S)D_4^322.$$
In Section 2 we show that a $`K3`$ surface $`X`$ of type $`P`$ can be realized as a double covering surface studied in Section 1 provided a certain fixed polarization (Theorem 2).
In Section 3 we define a fixed marking of $`S`$ of a $`K3`$ surface of type $`P`$. First we study the period domain for the family of such marked surfaces. That is a 6 dimensional domain given in the form
$$D^+=\{\eta =[\eta _1,\mathrm{},\eta _8]𝐏^7:{}_{}{}^{t}\eta A\eta =0,{}_{}{}^{t}\overline{\eta }A\eta >0,\mathrm{}(\eta _3/\eta _1)>0\},$$
where $`A=U(2)U(2)(2I_4)`$. It is a bounded symmetric domain of type IV. Next we determine the modular group for the isomorphism classes of the marked surfaces. It is given as the principal congruence subgroup $`G^+(2)`$ of level 2 in the full group
$$G^+=\{gPGL(8,𝐙):{}_{}{}^{t}gAg=A,g(D^+)=D^+\}$$
of the positive isometries for the lattice $`U(2)^22^4`$. The exact statement is given in Theorem 3.
In Section 4 we consider a general form of the period
$$u(x)=_C\mathrm{\Omega }=_C\underset{k=1}{\overset{4}{}}(x_1^{(k)}st+x_2^{(k)}s+x_3^{(k)}t+x_4^{(k)})^{1/2}dsdt,$$
where $`C`$ indicates an element of $`H_2(S(x),𝐙)`$. We describe the differential equation for $`_C\mathrm{\Omega }`$ as a multi-valued analytic function of $`16`$ variables $`x=(x_1,\mathrm{},x_4)`$, where we use the notation
$$x_k=\left(\begin{array}{cc}x_1^{(k)}& x_2^{(k)}\\ x_3^{(k)}& x_4^{(k)}\end{array}\right)M(2,𝐂).$$
That is a certain type of GKZ hypergeometric differential equation ( Proposition 4.1 and 4.2) that is not so called Aomoto-Gelfond type ( compare with ). We show the regularity of the system outside some divisor in the parameter space as stated in Theorem 4. We show the regular holonomicity of the system and make the calculation of the rank based on the theory of . The exact statement is given in Theorem 5. We determine the monodromy group for our system in Theorem 6.
In Section 5 we construct the family $`\{KS(\eta ):\eta D^+\}`$ of Kuga-Satake varieties corresponding to $``$. That becomes to be equivalent with the family of $`8`$ dimensional Abelian varieties with the Hamilton quaternion endomorphism structure (Theorem 7).
## 1 Lattice structure of a double $`4H`$ surface
### 1.1 Setting up the situation
We consider an algebraic surface $`S^{}`$ obtained as a double cover over $`𝐏^1\times 𝐏^1`$ ramifying along 4 different rational curves $`H_1,H_2,H_3,H_4`$ of bidegree $`(1,1)`$. Here we suppose the following generality condition:
(g1) $`H_i(i=1,2,3,4)`$ is irreducible,
(g2) $`H_iH_j(ij)`$ consists of 2 different points,
(g3) For any different 3 indices $`i,j,k`$ we have $`H_iH_jH_k=\mathrm{}`$.
We denote $`\pi `$ the projection $`S^{}𝐏^1\times 𝐏^1`$. Set $`L_i=\pi ^{}H_i`$. The surface $`S^{}`$ has 12 singular points of type $`A_1`$ corresponding to the intersections $`L_iL_j(ij)`$. By the desingularization procedure we get a $`K3`$ surface $`S`$.
If we have the condition
(e1) the algebraic variety $`S^{}`$ has at most simple singularities
instead of $`(\mathrm{g1}),(\mathrm{g2}),(\mathrm{g3})`$, we obtain a $`K3`$ surface by the same procedure. Henceforth we describe the curve $`H_k=H(x^k)(k=1,2,3,4)`$ in the form
$$(s,1)\left(\begin{array}{cc}x_1^{(k)}& x_2^{(k)}\\ x_3^{(k)}& x_4^{(k)}\end{array}\right)\left(\begin{array}{c}t\\ 1\end{array}\right)=0$$
with a matrix
$$x^k=\left(\begin{array}{cc}x_1^{(k)}& x_2^{(k)}\\ x_3^{(k)}& x_4^{(k)}\end{array}\right)\mathrm{M}(2,𝐂)$$
and an affine coordinate $`(s,t)`$ of $`𝐏^1\times 𝐏^1`$. Let $`S^{}=S^{}(x)`$ be the algebraic variety obtained from $`H(x^1)\mathrm{}H(x^4)`$ by the above way, and let $`S(x)`$ be its desingularization. Set
$$X^{}=\{x=(x^1,x^2,x^3,x^4)\mathrm{M}(2,𝐂)^4:H(x^1),\mathrm{},H(x^4)\text{satisfy}(g1),(g2),(g3)\},$$
and set
$$X^{}=\{x=(x^1,x^2,x^3,x^4)\mathrm{M}(2,𝐂)^4:S^{}\text{has at most simple singularities}\}.$$
###### Definition 1.1.
We call a double $`4H`$ surface the $`K3`$ surface $`S`$ obtained as $`S=S(x),xX^{}`$. An extended double $`4H`$ surface is a $`K3`$ surface obtained as $`S=S(x),xX^{}`$. Let $``$ denote the totality of double $`4H`$ surfaces:
$$\{S(x):xX^{}\}$$
We use the following notations.
L: the $`K3`$ lattice $`E_8^2U^3`$,where $`E_8`$ denotes the negative definite even unimodular lattice of rank 8 and $`U`$ denotes the hyperbolic lattice $`\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)`$. The basis of $`L`$ is fixed once and always.
$`\mathrm{Pic}(S)`$: the Picard lattice of $`S`$,
$`\mathrm{Tra}(S)`$: the transcendental lattice of $`S`$. That is defined as the orthogonal complement of $`\mathrm{Pic}(S)`$ in $`H^2(S,𝐙)`$
###### Remark 1.1.
We consider $`X^{}`$ to be the parameter space of $``$, but it contains many abundant parameters. Our family $``$ contains 6 essential parameters. We shall discuss the problem of abundance in Section 2 and Section 3.
### 1.2 Elliptic fibration
In this section we determine $`\mathrm{Tra}(S)`$ and $`\mathrm{Pic}(S)`$ for a general member of $``$. For this work we always consider the problem in the dual lattice $`H_2(S,𝐙)`$ and we observe a special member of $``$.
Let $`s,t`$ be the affine coordinates of $`𝐏^1\times 𝐏^1`$. Henceforth we denote the $`s`$-space ($`t`$-space ) by $`𝐏^1(s)(𝐏^1(t))`$, respectively. Set
$$f_1(s)=s,f_2(s)=\frac{16}{s},f_3(s)=\frac{5s+2}{2s+5},f_4(s)=\frac{3s+64}{4s+27}$$
Consider the following four $`(1,1)`$ curves :
$$H_i:t=f_i(s)(i=1,2,3,4),$$
and make the double $`4H`$ surface :
$$S_0:w^2=\underset{i=1}{\overset{4}{}}(tf_i(s))$$
that is derived from the above system $`\{H_1,\mathrm{},H_4\}`$. Let $`\pi _1`$ be the projection from $`S_0`$ to $`𝐏^1(s)`$. For a generic point on $`𝐏^1(s)`$ we obtain an elliptic curve $`\pi ^1(s)`$. So we get an elliptic surface $`(S_0,\pi _1,𝐏^1)`$.
It is easy to observe the following.
###### Proposition 1.1.
The elliptic surface $`(S_0,\pi _1,𝐏^1)`$ has 12 singular fibres corresponding to the intersections $`H_iH_j(ij)`$. Every singular fibre is of type $`I_2`$ ( according to the Kodaira classification of the singular fibres). These are situated over real points :
$`s=\pm 4(=H_1H_2),\pm 1(=H_1H_3),2\mathrm{and}8(=H_1H_4),`$
$`2\mathrm{and}8(=H_2H_3),\pm 12(=H_3H_4),\pm \sqrt{19}(=H_3H_4).`$
We denote these 12 points by $`s_j(j=1,\mathrm{},12)`$ according to the ascending order and set $`\mathrm{Sing}=\{s_1,\mathrm{},s_{12}\}`$.
We fix a base point $`s_0=4\sqrt{1}`$ and set $`E_0=\pi _1^1(s_0)`$. It is given by
$$w^2=(tf_1(4i))(tf_2(4i))(tf_3(4i))(tf_4(4i)).$$
So this is a double cover over $`𝐏^1(t)`$ branched at 4 points $`t=f_k(4i)(k=1,2,3,4)`$. On $`E_0`$ we take a basis $`\alpha _1,\alpha _2`$ of $`H_1(E_0,𝐙)`$ with the intersection multiplicity $`\alpha _1\alpha _2=1`$ as in Fig.1.
The fundamental group $`\pi _1(𝐏^1\mathrm{Sing},)`$ acts on $`H_1(E_0,𝐙)`$ as the monodromy group. We describe the generator system for it. Let $`\gamma _i(i=1,\mathrm{},12)`$ be a loop starting from $`s_0`$ and goes around $`s_i`$ in the positive sense on the upper half plain except the circuit around $`s_i`$ (see Fig.2 ).
By a direct observation we can calculate the monodromy transformation $`T_j=T(\gamma _j)`$ of $`H_1(E_0,𝐙)`$ along $`\gamma _j`$ with respect to the basis $`\{\alpha _1,\alpha _2\}`$.
###### Proposition 1.2.
The circuit matrices for $`T_j=T(\gamma _j)`$ are given by the following table. Here $`T_j`$ acts from right to the system $`\{\alpha _1,\alpha _2\}`$.
| $`s_j`$ | $`\pm 4,\pm \sqrt{19}`$ | $`\pm 1,\pm 12`$ | $`8,2`$ | $`2,8`$ |
| --- | --- | --- | --- | --- |
| $`T_j`$ | $`\left(\begin{array}{cc}1& 2\\ 0& 1\end{array}\right)`$ | $`\left(\begin{array}{cc}1& 0\\ 2& 1\end{array}\right)`$ | $`\left(\begin{array}{cc}1& 2\\ 2& 3\end{array}\right)`$ | $`\left(\begin{array}{cc}3& 2\\ 2& 1\end{array}\right)`$ |
| vanishing cycle | $`\alpha _1`$ | $`\alpha _2`$ | $`\alpha _1+\alpha _2`$ | $`\alpha _1\alpha _2`$ |
We choose a fixed point $`\underset{¯}{s}`$ in the lower half plane of $`𝐏^1(s)`$. Make cut lines $`c_i`$ by the line segment between $`\underset{¯}{s}`$ and $`s_i(i=1,\mathrm{},12)`$. By restricting $`(S_0,\pi _1,𝐏^1)`$ on $`𝐏^1_{i=1}^{12}c_i`$ we have a topologically trivial fibration. So we can determine the cycles $`\alpha _1,\alpha _2`$ of $`H_1(\pi ^1(s),𝐙)`$ for any $`s𝐏^1_{i=1}^{12}c_i`$ using this trivialization. If we make the continuation of the system $`\{\alpha _1,\alpha _2\}`$ passing through the line $`c_i`$ from left to right, It is transformed according to Table 1.
### 1.3 Two systems $`\{\mathrm{\Gamma }_i\}`$ and $`\{C_i\}`$ of $`H_2(S_0,𝐙)`$
Let $`r`$ be an oriented arc on $`𝐏^1(s)`$ starting from $`s_0`$. We make a 2-chain on $`S_0`$ by the continuation of a starting 1-cycle $`\alpha H_1(E_0,𝐙)`$ along $`r`$, and denote it by
$$r\times \alpha .$$
Here we define its orientation as the ordered pair of the ones of $`r`$ and $`\alpha `$. Using this notation we make a system $`\{\mathrm{\Gamma }_1,\mathrm{},\mathrm{\Gamma }_8\}`$ of 8 elements in $`H_2(S_0,𝐙)`$ (see Fig.3, Fig.4, Fig.5, Fig.6):
$`\mathrm{\Gamma }_1=\gamma (7)^1\gamma (8)^1\gamma (9)^1\gamma (11)^1\gamma (12)^1\times \alpha _1,`$
$`\mathrm{\Gamma }_2=\gamma (4)^1\gamma (5)^1\gamma (6)^1\gamma (7)^1\gamma (8)^1\gamma (9)^1\times \alpha _2,`$
$`\mathrm{\Gamma }_3=\gamma (5)^1\gamma (6)^1\gamma (7)^1\gamma (8)^1\times \alpha _1,`$
$`\mathrm{\Gamma }_4=\gamma (2)^1\gamma (3)^1\gamma (4)^1\gamma (5)^1\times \alpha _2,`$
$`\mathrm{\Gamma }_5=\gamma (12)\times (\alpha _1+\alpha _2)+\gamma (10)^1\gamma (11)^1\gamma (12)^1\times (\alpha _2),`$
$`\mathrm{\Gamma }_6=\gamma (6)\times (\alpha _1+\alpha _2)+\gamma (4)^1\gamma (5)^1\gamma (6)^1\times (\alpha _2),`$
$`\mathrm{\Gamma }_7=\gamma (9)\times (\alpha _1\alpha _2)+\gamma (7)^1\gamma (8)^1\gamma (9)^1\times (\alpha _1),`$
$`\mathrm{\Gamma }_8=\gamma (3)\times (\alpha _1\alpha _2)+\gamma (1)^1\gamma (2)^1\gamma (3)^1\times (\alpha _1).`$
where the composite arcs are performed from right to left.
We note that any of $`\{\mathrm{\Gamma }_1,\mathrm{},\mathrm{\Gamma }_8\}`$ is a 2-cycle according to the monodromy action in Table 1. So we regard them a system in $`H_2(S_0,𝐙)`$. Let $`L(\mathrm{\Gamma })`$ denote the sublattice $`𝐙\mathrm{\Gamma }_i`$.
Next we construct another system that is dual to $`L(\mathrm{\Gamma })`$. Let $`r(i)`$ be the oriented line segment from $`s_0`$ to $`s_i`$. Using this notation we construct 8 elements in $`H_2(S_0,𝐙)`$ ( see Fig.7):
$`C_1=(r(6)r(7))\times \alpha _2,`$
$`C_2=(r(9)r(10))\times \alpha _1,`$
$`C_3=(r(7)r(12))\times \alpha _2,`$
$`C_4=(r(4)r(9))\times \alpha _1,`$
$`C_5=r(10)\times \alpha _1r(11)\times (\alpha _1+\alpha _2)r(12)\times \alpha _2,`$
$`C_6=r(4)\times \alpha _1r(5)\times (\alpha _1+\alpha _2)r(6)\times \alpha _2,`$
$`C_7=r(7)\times \alpha _2r(8)\times (\alpha _1\alpha _2)r(9)\times \alpha _1,`$
$`C_8=r(1)\times \alpha _2r(2)\times (\alpha _1\alpha _2)r(3)\times \alpha _1.`$
We note that any element of $`\{C_1,\mathrm{},C_8\}`$ ends with a vanishing cycle at the terminal point of the base arc. So it is a 2-cycle on $`S_0`$. Let $`L(C)`$ denote the sublattice $`𝐙C_i`$.
###### Proposition 1.3.
We have the the matrices of the intersection numbers $`\mathrm{\Gamma }_i\mathrm{\Gamma }_j`$ and $`\mathrm{\Gamma }_iC_j`$ as follows :
$`(\mathrm{\Gamma }_i\mathrm{\Gamma }_j)=\left(\begin{array}{cccccccc}0& 2& 0& 0& 0& 0& 0& 0\\ 2& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 2& 0& 0& 0& 0\\ 0& 0& 2& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 2& 0& 0& 0\\ 0& 0& 0& 0& 0& 2& 0& 0\\ 0& 0& 0& 0& 0& 0& 2& 0\\ 0& 0& 0& 0& 0& 0& 0& 2\end{array}\right).`$
$`(\mathrm{\Gamma }_iC_j)=\left(\begin{array}{cccccccc}1& 0& 0& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 0\\ 0& 1& 1& 0& 1& 0& 0& 0\\ 1& 0& 0& 0& 0& 1& 0& 0\\ 1& 1& 1& 1& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\end{array}\right).`$
\[Proof\]. We can calculate all the intersection numbers by a direct observation of the cycles. For example we consider $`\mathrm{\Gamma }_1\mathrm{\Gamma }_2`$. They have 2 intersections on a vertical line between $`4`$ and $`\sqrt{19}`$ of $`\mathrm{\Gamma }_2`$ saying $`p_1,p_2`$ (see Fig.8).
At the upper intersection $`p_1`$ the base arcs intersect with a negative sign, and the fibre cycles $`\alpha _1,\alpha _2`$ intersect with a positive sign. By considering the definition of the orientation on $`\mathrm{\Gamma }_i`$ we get there intersection number $`+1`$. The situation is quite the same for the intersection in the lower half plane. So we have $`\mathrm{\Gamma }_1\mathrm{\Gamma }_2=2`$. We can proceed this type of calculation until we get the full intersection matrices.
q.e.d.
###### Remark 1.2.
1) By using the conventional notations $`U=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)`$ and $`U(k)=\left(\begin{array}{cc}0& k\\ k& 0\end{array}\right)`$, we have
$$(\mathrm{\Gamma }_i\mathrm{\Gamma }_j)=U(2)U(2)2^4.$$
2) We note that $`(\mathrm{\Gamma }_iC_j)`$ is a unimodular matrix. So we can find a generator system $`\{\stackrel{ˇ}{\mathrm{\Gamma }}_1,\mathrm{},\stackrel{ˇ}{\mathrm{\Gamma }}_8\}`$ of $`L(C)`$ with $`\mathrm{\Gamma }_i\stackrel{ˇ}{\mathrm{\Gamma }}_j=\delta _{ij}`$.
###### Lemma 1.1.
$`L(\mathrm{\Gamma })`$ and $`L(C)`$ are nondegenerate rank 8 lattices. $`L(\mathrm{\Gamma })`$ is a primitive sublattice in $`H_2(S_0,𝐙)L`$.
Proof\]. The first statement is a direct consequence of the fact that $`(\mathrm{\Gamma }_iC_j)`$ is a unimodular matrix.
Suppose $`L(\mathrm{\Gamma })`$ is not primitive. It means that there is an element $`\lambda LL(\mathrm{\Gamma })`$ with $`k\lambda L(\mathrm{\Gamma })`$ for some integer $`k`$.Put
$$k\lambda =\underset{i=1}{\overset{8}{}}k_i\mathrm{\Gamma }_i.$$
Then we have
$$k(\lambda \stackrel{ˇ}{\mathrm{\Gamma }}_i)=(k\lambda )\stackrel{ˇ}{\mathrm{\Gamma }}_i=k_i.$$
It indicates $`kk_i(i=1,\mathrm{},8)`$. So it is deduced $`\lambda L(\mathrm{\Gamma })`$. This is a contradiction.
q.e.d
### 1.4 The algebraic sublattice
For the moment we study the divisors on the reference surface $`S_0`$. Let $`P_{ij}^\pm `$ denote two intersections $`H_iH_j(1i<j4)`$, here we distinguish them by the signature of the $`s`$-coordinate ( see Proposition 1.1). The surface $`S_0`$ has an exceptional curve corresponding to every intersection $`P_{ij}^\pm (1i<j4)`$. We denote them by $`E_{ij}^\pm `$. Let us consider the projection $`\pi :S𝐏^1\times 𝐏^1`$. Let $`F_s`$ denote the pull back $`\pi ^{}(𝐏^1(s))\times \{t\})`$ of the generic $`s`$-coordinate line, and let $`F_t`$ denote the pull back of the generic $`t`$-coordinate line. Set
$$G_i=\frac{1}{2}(\pi ^{}H_i\underset{ji}{}(E_{ij}^++E_{ij}^{})),(i=1,2,3,4).$$
This is the reduced divisor coming from the 4 sections $`H_i`$ of the elliptic fibration $`(S_0,\pi _1,𝐏^1(s))`$.By an easy observation we have the following.
###### Lemma 1.2.
We have intersection numbers among $`E_{ij}^\pm `$, $`G_i`$, $`F_s`$, $`F_t`$ :
$`G_iG_i=2,F_sF_t=2,G_iF_s=G_iF_t=1,`$
$`G_iE_{ij}^+=G_iE_{ij}^{}=1,E_{ij}^+E_{ij}^+=E_{ij}^{}E_{ij}^{}=2,`$
and all the other intersections are 0. We let $`L(div)`$ denote the sublattice of $`H_2(S_0,𝐙)=L`$ generated by $`E_{ij}^\pm ,G_i,F_s,F_t`$.
###### Remark 1.3.
If we observe the construction of $`\mathrm{\Gamma }_i`$, we can show that $`L(div)L(\mathrm{\Gamma })^{}`$. Where $``$ indicates the orthogonal complement in the full lattice $`H_2(S_0,𝐙)L`$.
###### Proposition 1.4.
We have
$`L(\mathrm{\Gamma })^{}=L(div)=E_{ij}^\pm ,G_i,F_s,F_t`$
$`(D_4)^322`$
on the reference surface $`S_0`$.
Proof\]. Set the sublattices $`L_i(i=1,2,3)`$ as follows:
$`L_1=G_1,E_{14}^+,E_{14}^{},F_tE_{23}^+,`$
$`L_2=G_2,E_{24}^+,E_{24}^{},F_tE_{13}^+,`$
$`L_3=G_3,E_{34}^+,E_{34}^{},F_tE_{12}^+.`$
They are isometric with the lattice
$$D_4=\left(\begin{array}{cccc}2& 1& 1& 1\\ 1& 2& 0& 0\\ 1& 0& 2& 0\\ 1& 0& 0& 2\end{array}\right),$$
and they are perpendicular each others. Set
$`\mathrm{\Delta }_1=F_s+F_tE_{12}^+E_{13}^+E_{23}^+,`$
$`\mathrm{\Delta }_2=G_1+G_2+G_3G_4+F_s+3F_tE_{12}^+E_{13}^+E_{23}^+,`$
then we have
$$\mathrm{\Delta }_1\mathrm{\Delta }_1=2,\mathrm{\Delta }_1\mathrm{\Delta }_2=0,\mathrm{\Delta }_2\mathrm{\Delta }_2=2.$$
Now we put
$$P=L_1L_2L_3\mathrm{\Delta }_1\mathrm{\Delta }_2.$$
Then it holds
$`PL(div)L(\mathrm{\Gamma })^{},`$
$`\mathrm{rank}P=14,\mathrm{discr}P=2^8.`$
On the other hand we have $`\mathrm{discr}L(\mathrm{\Gamma })=2^8`$. By considering the fact that $`L`$ is unimodular and $`L(\mathrm{\Gamma })`$ is a primitive sublattice, we get the conclusion that $`P=L(div)=L(\mathrm{\Gamma })^{}`$.
q.e.d.
Now we extend Proposition 1.4 to the general situation.
###### Theorem 1.
Let $`S`$ be a double $`4H`$ surface. If we have $`\mathrm{rank}\mathrm{Pic}(S)=14`$, then it holds
$`\mathrm{Tra}(S)U(2)U(2)(2I_4),`$
$`\mathrm{Pic}(S)D_4D_4D_422.`$
\[Proof\]. We note that the family $`=\{S(x)\}`$ is fibred over the parameter space of $`\mathrm{M}(2,𝐂)^4`$ . This fibration is locally topological trivial on a Zariski open subset in $`X^{}\mathrm{M}(2,𝐂)^4`$. Using this trivialization we can proceed the same argument for a general member as for the specialized element $`S_0`$.In fact always $`L(div)`$ and $`P`$ are defined as a sublattice in $`\mathrm{Pic}(S(x))`$. We have the systems $`L(\mathrm{\Gamma })`$ and $`L(C)`$ also. So Proposition 1.4 is valid for every $`S`$ of $``$. If the condition is satisfied, then $`L(\mathrm{\Gamma })`$ cannot contain any divisor class. Hence we obtain the required conclusion.
q.e.d
###### Remark 1.4.
As we will see in Section 2, we have $`\mathrm{rank}\mathrm{Pic}(S)=14`$ for a general member of $``$.
## 2 Family $``$ as that of lattice $`K3`$ surfaces
### 2.1 General properties
In this section we show the converse of Theorem 1. Before starting this argument we state some general properties of a $`K3`$ surface.
###### Lemma 2.1.
Let $`S`$ be a $`K3`$ surface that is given by a minimal nonsingular model.
1) For an irreducible curve $`C`$ on $`S`$, $`C^20`$ or $`C^2=2`$.
2) If $`D\mathrm{Pic}(\mathrm{S})`$ satisfies $`D^22`$, then either $`D`$ or $`D`$ is an effective class ( that is given by an effective divisor).
This is deduced from the Riemann-Roch theorem by a routine argument.
###### Definition 2.1.
We call an element $`D\mathrm{Pic}(S)`$ is nef when it holds $`DC0`$ for any effective class $`C`$.
###### Proposition 2.1 (Pjateckiĭ-Šapiro and Šafarevič ).
Let $`S`$ be a $`K3`$ surface, then we have the following:
1) Suppose $`x\mathrm{Pic}(S)`$ satisfies $`x0,x^2=0`$. Then there exists an isometry $`\gamma `$ of $`\mathrm{Pic}(S)`$ such that $`\gamma (x)`$ becomes to be effective and nef.
2) Suppose $`x\mathrm{Pic}(S)`$ is effective, nef and $`x^2=0`$ then $`x`$ is a multiple class of an elliptic curve (i.e. $`x=m[E]`$ for a certain $`m𝐍`$ and an elliptic curve $`E`$ on $`S`$).
3) The linear system of an elliptic curve $`C`$ on $`X`$ determines an elliptic fibration $`S𝐏^1`$.
### 2.2 A lattice $`K3`$ surface and its realization as a double $`4H`$ surface
Let us consider the lattice
$`P=D_4^322,`$ (2.1)
and we let $`(,)`$ denote the bilinear form on an abstract lattice $`P`$.
###### Definition 2.2.
We say a $`K3`$ surfaces $`S`$ is $`\mathrm{of}\mathrm{type}P`$ if
$`(\mathrm{a})S`$ admits an embedding $`P\mathrm{Pic}(S)`$,
and $`S`$ is $`\mathrm{of}\mathrm{exact}\mathrm{type}P`$ if
$`(\mathrm{a}^{})S`$ admits an isomorphism $`P\mathrm{Pic}(S)`$.
By the argument in Section 1 we can find the generator system of $`P`$:
$`f_s,f_t,g_i(i=1,\mathrm{},4),e_{ij}^+,e_{ij}^{}(ij,1i,j4),`$ (2.2)
with the following properties among them :
(p1)
$`(g_i,g_i)=2,(f_s,f_t)=2,(g_i,f_s)=(g_i,f_t)=1`$
$`(g_i,e_{ij}^+)=(g_i,e_{ij}^{})=1,(e_{ij}^+,e_{ij}^+)=(e_{ij}^{},e_{ij}^{})=2,`$
and all other intersections equal $`0`$,
(p2)
$`f_s+f_t=2g_i+{\displaystyle \underset{ji}{}}(e_{ij}^++e_{ij}^{}).`$
We note that this system is determined uniquely up to isometries of $`P`$.
###### Theorem 2.
Let $`S`$ be a $`K3`$ surface of type $`P`$ with the property:
(n) $`f_s,f_t`$ are nef classes under the identification in $`(\mathrm{a})`$.
Then $`S`$ has a representation as an extended double $`4H`$ surface.
When $`S`$ is of exact type $`P`$ with the above conditioin $`(\mathrm{n})`$, $`S`$ has a representation as a double $`4H`$ surface.
\[Proof\]. Always we regard the element of $`P`$ under the identification in $`(\mathrm{a})`$ or $`(\mathrm{a}^{})`$. We show the first assertion.
(Claim 1) We can assume that every $`e_{ij}^\pm `$ is effective.
Let $`e`$ be one of the twelve elments $`e_{ij}^\pm `$. Since $`ee=2`$, either $`e`$ or $`e`$ is effective. If $`e`$ is effective then we perform the reflection $`\gamma `$ on $`\mathrm{Pic}(S)`$ determined by $`e`$. And we can use
$$\gamma (f_s),\gamma (f_t),\gamma (g_i),\gamma (e_{ij}^\pm )$$
instead of the system $`(2.2)`$. According to the orthogonality of the system $`\{e_{ij}^\pm \}`$ we can iterate this procedure until we get the required effective system. Since we have $`\gamma (f_s)=f_s,\gamma (f_t)=f_t`$, the nef property for $`f_s,f_t`$ is always satisfied throughout this process.
(Claim 2) We can find a double covering $`S^{}`$, that is birationally equivalent with $`S`$,
$$\pi =(\pi _1,\pi _2):S^{}𝐏^1\times 𝐏^1,$$
ramified along a bidegree $`(4,4)`$ curve $`B`$. $`S^{}`$ has at most simple singularities coming from the singular points of $`B`$.
By Proposition 2.1 $`f_s,f_t`$ are multiple classes of elliptic curve. Observing the assumption $`g_if_s=g_if_t=1`$ we know that they are reduced classes. Again by Proposition 2.1 $`f_s`$ and $`f_t`$ determine two different elliptic fibrations $`\pi _1:X𝐏^1`$ and $`\pi _2:X𝐏^1`$, respectively. Set
$$\pi =(\pi _1,\pi _2):S𝐏^1\times 𝐏^1.$$
This map is surjective and of degree 2, because we have $`\pi _1^1(x)\pi _1^1(y)=f_sf_t=2`$ for any $`(x,y)𝐏^1\times 𝐏^1`$. Let $`L_s,L_t`$ be two lines $`𝐏^1\times \{\mathrm{}\}`$ and $`\{\mathrm{}\}\times 𝐏^1`$, respectively. So we obtain a double covering $`\pi =(\pi _1,\pi _2):S𝐏^1\times 𝐏^1`$. Generic fibers $`\pi _1^{}(x)`$ and $`\pi _2^{}(y)`$ are elliptic curves realized as double coverings over $`𝐏^1`$. So each of them has four branch points. Consequently, the branch locus $`B`$ of $`\pi `$ is a curve of degree $`(4,4)`$. Since $`S`$ is a $`K3`$ surface, the canonical class $`K_S=0`$. Hence $`B`$ has at most simple singularities.
(Claim 3) Let $`G_i`$ be the effective divisor representing $`g_i`$.Then $`H_i=\pi _{}G_i`$ is a $`(1,1)`$-curve $`(i=1,2,3,4)`$.
Since $`g_ig_i=2`$, either $`g_i`$ or $`g_i`$ is effective. But $`f_s`$ is nef and has intersection $`f_sg_i=1`$, then we know that $`g_i`$ is effective. By the projection formula, we obtain
$$\pi _{}G_iL_s=G_i\pi ^{}L_s=g_if_s=1.$$
By the same argument, we get $`\pi _{}G_iL_t=1`$ also. Hence we obtain the required property.
(Claim 4) The effective class $`e_{ij}^\pm `$ is obtained by an exceptional divisor coming from the singularity of $`B`$.
Let $`E_{ij}^\pm `$ be the effective divisor representing $`e_{ij}^\pm `$. Since $`(e_{ij}^\pm ,f_s)=(e_{ij}^\pm ,f_t)=0`$ we get $`\pi _{}E_{ij}^\pm =0`$. It indicates that every $`E_{ij}^\pm `$ is an exceptional divisor derived from the singularity of $`B`$.
(Claim 5). We have $`B=H_1+H_2+H_3+H_4`$
We consider the $`(1,1)`$ curve $`H_i=\pi _{}G_i`$. As a divisor class $`\pi ^{}H_i`$ equals to $`f_s+f_t`$. By the starting assumption $`(\mathrm{p2})`$, $`\pi ^{}H_i`$ and $`2G_i+_{ji}(E_{ij}^++E_{ij}^{})`$ are linearly equivalent. So there exist a principal divisor $`D`$ such that
$$\pi ^{}H_i=2G_i+\underset{ji}{}(E_{ij}^++E_{ij}^{})+D.$$
Then
$$2H_i=\pi _{}\pi ^{}H_i=2\pi _{}G_i+\pi _{}D=2H_i+\pi _{}D.$$
This is an equality as divisors itsselves (not as divisor classes), so we get $`\pi _{}D=0`$. This implies that $`D`$ is an sum of exceptional divisors. But $`D`$ is principal, so we must have $`D=0`$. By observing the equality
$$\pi ^{}H_i=2G_i+\underset{ji}{}(E_{ij}^++E_{ij}^{})$$
we know that $`\pi ^{}H_i`$ has six components of the form $`E_{ij}^\pm `$, note that we don’t have any cancellation with an effective divisor $`G_i`$. Hence $`H_i`$ meets six different singular points on $`B`$. The sum of intersection numbers at these double points exceeds the number $`(B,H_i)=8`$. In case $`H_i`$ is irreducible it must be an irreducible component of $`B`$. In case we have $`H_i=\mathrm{}_1+\mathrm{}_2`$ with a $`(0,1)`$ curve and a $`(1,0)`$ curve, at least one component, saying $`\mathrm{}_1`$, is contained in $`B`$. If $`\mathrm{}_2`$ is not contained in $`B`$, we have $`\mathrm{}_2(B\mathrm{}_1)=3`$. So $`\mathrm{}_2`$ contains only one $`\pi (E_{ij}^\pm )`$. Namely $`\mathrm{}_1`$ contains 5 others. Observing $`\mathrm{}_1(B\mathrm{}_1)=4`$ we know that is too many. So we have $`H_iB`$. Hence we obtained the claim.
According to the above arguments we obtained the double covering $`\pi :S𝐏^1\times 𝐏^1`$. realizing the extended double $`4H`$ surface.
Next let us prove the second assertion using the condition $`(\mathrm{a}^{})`$. It is enough to show
(Claim 6) $`S^{}`$ has exactly 12 singular points of type $`A_1`$ on $`\pi ^{}B`$.
Let $`L(exc)`$ be the sublattice generated by the system $`\{e_{ij}^\pm \}`$ in $`P=\mathrm{Pic}(S)`$, this is isometric with $`2^{12}`$. Suppose $`xP`$ is a class of irreducible $`(2)`$-curve with $`\pi _{}x=0`$. We can describe $`x`$ in the form
$$x=m_sf_s+m_tf_t+e+\underset{i=1}{\overset{4}{}}m_ig_i,eL(exc)$$
with the coefficients $`m_s,m_t,m_i𝐙`$. By assumption it holds $`\pi _{}xL_s=\pi _{}xL_t=0`$, so we obtain
$$m_s=m_t,2m_s+m_1+m_2+m_3+m_4=0.$$
According to the starting condition $`(\mathrm{p1}),(\mathrm{p2})`$, we have
$`2(xe)`$ $`=`$ $`2(m_s(f_s+f_t)+m_1g_i+m_2g_2+m_3g_3+m_4g_4)`$
$`=`$ $`{\displaystyle \underset{i=1}{\overset{4}{}}}m_i(2g_if_sf_t)`$
$`=`$ $`{\displaystyle \underset{i=1}{\overset{4}{}}}m_i{\displaystyle \underset{ji}{}}(e_{ij}^++e_{ij}^{})`$
So we have $`2xL(excep)`$. By assumption we have $`(2x,2x)=8`$. But such an element in $`L(exc)`$ should be the form $`\pm 2e_{ij}^\pm `$ or $`\pm e_1\pm e_2\pm e_3\pm e_4`$. Where $`\{e_i(i=1,2,3,4)\}`$ are distinct four elements of $`\{e_{ij}^\pm \}`$. In the case $`2x=\pm 2e_{ij}^\pm `$, we conclude $`xL(excep)`$.
In the later case, we obtain
$$\pm e_1\pm e_2\pm e_3\pm e_42e=\underset{i=1}{\overset{4}{}}m_i\underset{ji}{}(e_{ij}^++e_{ij}^{}).$$
Observe the right hand side, the number of the odd coefficients of $`e_{ij}^\pm `$ should be one of $`0,6,8`$. It does not attained in the left hand. So this case does not happen. Hence we obtained the claim.
q.e.d
## 3 Modular group and marking with a multi-polarization
### 3.1 Congruence subgroup $`G(2)`$
Let $`L`$ be the K3 lattice with the fixed basis as defined in Section 1 and $`(,)`$ be the corresponding bilinear form. We consider the sublattice $`P=(\mathrm{D}_4)^322`$ with the fixed generator system $`\{f_s,f_t,e_{ij}^\pm ,g_i\}`$ as in Section 2. Set $`T=P^{}=\mathrm{U}(2)^22^4`$ with the fixed basis so that $`A=U(2)U(2)(2I_4)`$ be the intersection form. We denote the group of isometries of a lattice $`M`$ by $`\mathrm{O}(M)`$. Set $`G=\mathrm{O}(T)`$ and
$$\mathrm{O}(L,P)=\{g\mathrm{O}(L):g(x)=x\text{for}xP\}.$$
Note that we have $`g(T)=T`$ for $`g\mathrm{O}(L,P)`$. So we can regard $`\mathrm{O}(L,P)`$ as a subgroup of $`G`$. Now G is given as
$$G=\{gM(8,𝐙):{}_{}{}^{t}gAg=A\}.$$
Set
$$G(2)=\{gG:gI\text{mod}2\}.$$
Let $`\alpha T`$ be a $`(2)`$element, that is $`(\alpha ,\alpha )=2`$. It determines a reflection $`\gamma _\alpha \mathrm{O}(L,P)`$ by putting
$$\gamma _\alpha (x)=x+(x,\alpha )\alpha .$$
###### Remark 3.1.
According to and , we know that $`G(2)`$ is generated by the reflections of the above form $`\gamma _\alpha `$. So it holds $`G(2)\mathrm{O}(L,P)`$.
###### Proposition 3.1.
We have
$$G(2)=\mathrm{O}(L,P)$$
.
\[Proof\]. It is enough to show $`\mathrm{O}(L,P)G(2)`$. Here we use the notations and results in . Let $`q_T:A_T𝐐`$ (resp. $`q_P`$) be the discriminant form of $`T`$ (resp. $`P`$), where $`A_T=T^{}/T`$.
(Fact): $`H=L/TP`$ is a maximal isotropy subgroup for $`q_Tq_P`$ in $`A_TA_P`$. Namely it holds $`(q_Tq_P)_H0`$ and $`H`$ is maximal with this property. Moreover the canonical projections $`HA_T`$ and $`HA_P`$ induce the isomorphism
$$A_THA_P.$$
Hence we have
$$\mathrm{O}(L,P)Ker\{G\mathrm{O}(q_T)\}$$
, where $`G\mathrm{O}(q_T)`$ indicates the canonical map. We can easily show that
$$G(2)=Ker\{G\mathrm{O}(q_T)\}.$$
q.e.d
### 3.2 Marking with a multi-polarization
###### Definition 3.1.
Let $`S`$ be a $`K3`$ surface with an embedding $`P\mathrm{Pic}(S)`$. The triple $`(S,\phi ,P)`$ is a $`P`$-marking of $`S`$ provided
1) $`\phi :H_2(S,𝐙)L`$ is an isometry,
2) $`\phi ^1(f_s),\phi ^1(f_t),\phi ^1(e_{ij}^\pm ),\phi ^1(g_i)`$ are effective and $`\phi ^1(f_s),\phi ^1(f_t)`$ are nef.
According to Theorem 2 we have the double covering representation $`\pi :S𝐏^1\times 𝐏^1`$ and the class
$$\pi _{}(\phi ^1(g_1)+\mathrm{}+\phi ^1(g_4))$$
gives the class of the ramification divisor.
###### Definition 3.2.
Let $`(S,\phi ,P)`$ and $`(S^{},\phi ^{},P)`$ be $`P`$-markings of $`S`$ and $`S^{}`$, respectively. An isomorphism $`\rho :SS^{}`$ is an isomorphism between these markings provided $`\phi =\phi ^{}\rho _{}`$.
###### Remark 3.2.
Because we have $`\rho _{}\phi ^1(f_s)=(\phi ^{})^1(f_s)`$ and $`\rho _{}\phi ^1(f_t)=(\phi ^{})^1(f_t)`$ , such an isomorphism $`\rho `$ preserves the covering structure over $`𝐏^1\times 𝐏^1`$.
Namely we have the following commutative diagram with an element $`\sigma \mathrm{PGL}_2(𝐂)\times \mathrm{PGL}_\mathrm{𝟐}(𝐂)`$ :
$`\begin{array}{ccc}S& \stackrel{\rho }{}& S^{}\\ \pi & & \pi ^{}& & \\ 𝐏^1\times 𝐏^1& \stackrel{\sigma }{}& 𝐏^1\times 𝐏^1\end{array}`$
###### Definition 3.3.
Let $`H=(H_1,H_2,H_3,H_4)`$, $`H^{}=(H_1^{},H_2^{},H_3^{},H_4^{})`$ be ordered sets of four curves of bidegree $`(1,1)`$. We say that $`H`$ and $`H^{}`$ are equivalent if there exists $`\sigma \mathrm{PGL}_2(𝐂)\times \mathrm{PGL}_\mathrm{𝟐}(𝐂)`$ such that $`\sigma (H_i)=H_i^{}(i=1,2,3,4)`$.
###### Remark 3.3.
Let $`(S,\phi ,P)`$ and $`(S^{},\phi ^{},P)`$ be $`P`$-markings of $`S`$ and $`S^{}`$, respectively.
(1) If these two $`P`$-markings are isomorphic then we have the same equivalent class of the ordered sets $`H=(H_1,H_2,H_3,H_4)`$.
(2) But we don’t have the converse of (1).
### 3.3 Modular group
Let $`\{\mathrm{\Gamma }_1,\mathrm{},\mathrm{\Gamma }_8\}`$ be the basis of $`T`$ such that $`(\mathrm{\Gamma }_i,\mathrm{\Gamma }_j)=A`$. We have elements $`\{\stackrel{ˇ}{\mathrm{\Gamma }}_1,\mathrm{},\stackrel{ˇ}{\mathrm{\Gamma }}_8\}`$ in $`L`$ such that $`(\stackrel{ˇ}{\mathrm{\Gamma }}_i,\mathrm{\Gamma }_j)=\delta _{ij}`$. The system $`\{\stackrel{ˇ}{\mathrm{\Gamma }}_i\}`$ is uniquely determined modulo $`P`$.
Let $`(S,\phi ,P)`$ be a $`P`$-marking of a surface $`S`$. We consider the period
$$[_{\phi ^1(\stackrel{ˇ}{\mathrm{\Gamma }}_1)}\mathrm{\Omega }:\mathrm{}:_{\phi ^1(\stackrel{ˇ}{\mathrm{\Gamma }}_8)}\mathrm{\Omega }]𝐏^7$$
where $`\mathrm{\Omega }`$ is a holomorphic 2-form on $`S`$. The bilinear relation
$$_S\mathrm{\Omega }\mathrm{\Omega }=0,_S\overline{\mathrm{\Omega }}\mathrm{\Omega }>0$$
implies that the period belongs to the domain
$$D=\{\eta 𝐏^7:{}_{}{}^{t}\eta A\eta =0,{}_{}{}^{t}\overline{\eta }A\eta >0\}.$$
The domain $`D`$ has two connected components
$$D=D^+D^{},D^\pm =\{(\eta _1,\mathrm{},\eta _8)D:\pm \mathrm{Im}(\eta _3/\eta _1)>0\}$$
and we can take $`D^+`$ as the period domain for the family of isomorphism classes of marked surfaces $`\{(S(x),\phi ,P):xX^{}\}`$ .
###### Remark 3.4.
Two domains $`D^\pm `$ are complex conjugate. That is
$$[_{\phi ^1(\stackrel{ˇ}{\mathrm{\Gamma }}_1)}\mathrm{\Omega }:\mathrm{}:_{\phi ^1(\stackrel{ˇ}{\mathrm{\Gamma }}_8)}\mathrm{\Omega }]D^+[_{\phi ^1(\stackrel{ˇ}{\mathrm{\Gamma }}_1)}\overline{\mathrm{\Omega }}:\mathrm{}:_{\phi ^1(\stackrel{ˇ}{\mathrm{\Gamma }}_8)}\overline{\mathrm{\Omega }}]D^{}.$$
The group $`G`$ acts on the domain $`D`$. Set
$$G(2)^+=\{gG(2):g(D^+)=D^+\}.$$
It is a subgroup of $`G(2)`$ with index 2.
###### Theorem 3.
Let $`(S,\phi ,P)`$ and $`(S^{},\phi ^{},l_s^{},P)`$ be $`P`$-markings of extended double $`4H`$ surfaces $`S`$ and $`S^{}`$ , respectively. Let $`\eta ,\eta ^{}D^+`$ be corresponding periods. Then these markings are isomorphic if and only if
$$g(\eta )=\eta ^{}$$
for some $`gG(2)^+(=\mathrm{O}^+(L,P):=\{g\mathrm{O}(L,P):g(D^+)=D^+\})`$.
###### Proof.
Assume $`g(\eta )=\eta ^{}`$ for some $`gG(2)^+`$. According to Proposition 3.1 there exist an element $`\widehat{g}\mathrm{O}(L,P)`$ such that $`\widehat{g}|_T=g`$. Then we have
$$[_{\phi ^1\widehat{g}(\stackrel{ˇ}{\mathrm{\Gamma }}_1)}\mathrm{\Omega }:\mathrm{}:_{\phi ^1\widehat{g}(\stackrel{ˇ}{\mathrm{\Gamma }}_8)}\mathrm{\Omega }]=[_{(\phi ^{})^1(\stackrel{ˇ}{\mathrm{\Gamma }}_1)}\mathrm{\Omega }^{}:\mathrm{}:_{(\phi ^{})^1(\stackrel{ˇ}{\mathrm{\Gamma }}_8)}\mathrm{\Omega }^{}],$$
where $`\mathrm{\Omega }^{}`$ indicates the holomorphic form on $`S^{}`$. Consider the composition
$$f=(\phi ^{})^1\widehat{g}^1\phi :\mathrm{H}_2(S,𝐙)LL\mathrm{H}_2(S^{},𝐙).$$
The above composite isomorphism induces the dual map
$$f^{}:\mathrm{H}^2(S^{},𝐙)\mathrm{H}^2(S,𝐙)$$
with $`f^{}(\mathrm{H}^{2,0}(S^{}))=\mathrm{H}^{2,0}(S)`$. Moreover, $`f^{}`$ preserves ample classes. Hence the Torelli theorem for K3 surfaces assures that there exists the unique isomorphism $`\rho :SS^{}`$ such that $`\rho _{}=f`$. It is obvious that $`\rho `$ is an isomorphism of marked surfaces.
The converse is derived by the same argument. ∎
## 4 The hypergeometric differential equation for the periods
Let us take an element $`CH_2(S_0,𝐙)`$, and let $`C(x)H_2(S(x),𝐙)`$ denote its continuation to $`x=(x^1,x^2,x^3,x^4)X^0`$ ( that is multivalued and depends on the paths to $`x`$ in $`X^{}`$). Now we investigate the differential equation for the period
$`u(x)={\displaystyle _{C(x)}}\mathrm{\Omega }={\displaystyle _{C(x)}\{\underset{p=1}{\overset{4}{}}(x_{11}^p\xi _1\eta _1+x_{12}^p\xi _1\eta _2+x_{21}^p\xi _2\eta _1+x_{22}^p\xi _2\eta _2)^{1/2}\}\omega }`$ (4.1)
$`\omega =(\xi _1d\xi _2\xi _2d\xi _1)(\eta _1d\eta _2\eta _2d\eta _1),`$ (4.2)
where $`[\xi _1,\xi _2]`$ and $`[\eta _1,\eta _2]`$ denote the homogeneous coordinates of $`𝐏^1(s)`$ and $`𝐏^1(t)`$, respectively. Note that this integral does not depend on the affine representatives of the homogeneous coordinate of $`𝐏^1(s)`$ and $`𝐏^1(t)`$.
There are left and right actions of $`\mathrm{GL}(2,𝐂)`$ and a multiplicative action of $`(𝐂^{})^4`$ on $`X^0`$. They induce the following actions on the period $`u(x)`$ :
$`u(gx)=u(gx^1,\mathrm{},gx^4)g\mathrm{GL}(2,𝐂),`$
$`u(xh)=u(x^1{}_{}{}^{t}h,\mathrm{},x^4{}_{}{}^{t}h)h\mathrm{GL}(2,𝐂),`$
$`u(\lambda x)=u(\lambda _1x^1,\mathrm{},\lambda _4x^4)\lambda =(\lambda _1,\mathrm{},\lambda _4)(𝐂^{})^4`$
###### Lemma 4.1.
We have the following equalities.
$`(1)`$ $`u(\lambda x)={\displaystyle \underset{p=1}{\overset{4}{}}}\lambda _p^{a_p1}u(x).`$ (4.3)
$`(2)`$ $`u(gx)={\displaystyle \frac{1}{\mathrm{det}(g)}}u(x),u(xh)={\displaystyle \frac{1}{\mathrm{det}(h)}}u(x).`$ (4.5)
###### Proposition 4.1.
The integral $`u(x)`$ satisfies the following systems :
$`{\displaystyle \underset{1j,k2}{}}x_{jk}^p{\displaystyle \frac{u}{x_{jk}^p}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}u(x)(p=1,2,3,4),`$ (4.6)
$`E_1:{\displaystyle \underset{p=1}{\overset{4}{}}}{\displaystyle \underset{j=1}{\overset{2}{}}}x_{lj}^p{\displaystyle \frac{u}{x_{mj}^p}}`$ $`=`$ $`\delta _{lm}u(x)(\mathrm{},m\{1,2\}),`$ (4.8)
$`{\displaystyle \underset{p=1}{\overset{4}{}}}{\displaystyle \underset{j=1}{\overset{2}{}}}x_{jl}^p{\displaystyle \frac{u}{x_{jm}^p}}`$ $`=`$ $`\delta _{lm}u(x)(\mathrm{},m\{1,2\}),`$ (4.10)
$`E_2:{\displaystyle \frac{^2u}{x_{ij}^qx_k\mathrm{}^p}}`$ $`=`$ $`{\displaystyle \frac{^2u}{x_k\mathrm{}^qx_{ij}^p}}(i,j,k,\mathrm{}\{1,2\},p,q\{1,2,3,4\}),`$ (4.11)
$`{\displaystyle \frac{^2u}{x_{11}^qx_{22}^p}}`$ $`=`$ $`{\displaystyle \frac{^2u}{x_{21}^qx_{12}^p}}(p,q\{1,2,3,4\}).`$ (4.12)
\[Proof\] Differentiate the first equality (4.2) in Lemma 4.1 with respect to $`\lambda _i`$ and put $`\lambda =(\lambda _1,\lambda _2,\lambda _3,\lambda _4)=(1,1,1,1)`$. Then we get the first equality (4.4).
Differentiate the second equality (4.3) in Lemma 3.1 for the left action with respect to the $`ij`$-component $`g_{ij}`$ of $`g`$ and put $`g=I_2`$. Then we get the second equality (4.5). We get the third equality by using the equality for the right action of $`\mathrm{PGL}(2,𝐂)`$ with the same method. The system $`E_2`$ is deduced from the direct computation and the commutativity of the partial differentiations and the integral.
q.e.d.
Now we show the system $`E_1+E_2`$ is a holonomic system on $`X^0`$. We consider the variety $`B`$ in the cotangent bundle $`T^{}(X^{})`$ defined by
$`{\displaystyle \underset{1j,k2}{}}x_{jk}^p\xi _{jk}^p`$ $`=`$ $`0(p=1,2,3,4),`$ (4.13)
$`{\displaystyle \underset{p=1}{\overset{4}{}}}{\displaystyle \underset{j=1}{\overset{2}{}}}x_\mathrm{}j^p\xi _{mj}^p`$ $`=`$ $`0(\mathrm{},m\{1,2\}),`$ (4.14)
$`{\displaystyle \underset{p=1}{\overset{4}{}}}{\displaystyle \underset{j=1}{\overset{2}{}}}x_j\mathrm{}^p\xi _{jm}^p`$ $`=`$ $`0(\mathrm{},m\{1,2\}),`$ (4.15)
$`\xi _{ij}^q\xi _k\mathrm{}^p\xi _k\mathrm{}^q\xi _{ij}^p`$ $`=`$ $`0(i,j,k,\mathrm{}\{1,2\},p,q\{1,2,3,4\}),`$ (4.17)
$`\xi _{11}^q\xi _{22}^p\xi _{21}^q\xi _{12}^p`$ $`=`$ $`0(p,q\{1,2,3,4\}),`$ (4.19)
where $`\xi _{jk}^p`$ stands for $`\frac{u}{x_{ij}^k}`$. We regard $`B`$ as a fiber space over $`X^{}`$ in the product of $`X^{}`$ and a space of symbols $`\xi =(\xi ^1,\xi ^2,\xi ^3,\xi ^4)`$ for
$$\xi ^p=\left(\begin{array}{cc}\xi _{11}^p& \xi _{12}^p\\ \xi _{21}^p& \xi _{22}^p\end{array}\right).$$
The characteristic variety itself does not necessarily coincide with $`B`$, but it contains $`B`$, which we will call the fake characteristic variety. We will prove that the fake characteristic variety consists of only $`\{0\}`$ (zero-section) on $`X^0`$. It implies that any solution of $`E_1+E_2`$ on a simply connected domain in $`X^0`$ is holomorphic.
Let us fix an arbitrary point
$`x=(x^1,x^2,x^3,x^4)X^0,x^p=\left(\begin{array}{cc}x_{11}^p& x_{12}^p\\ x_{21}^p& x_{22}^p\end{array}\right)(p=1,2,3,4).`$
Put
$`H_p:f_p=x_{11}^p\xi _1\eta _1+x_{12}^p\xi _1\eta _2+x_{21}^p\xi _2\eta _1+x_{22}^p\xi _2\eta _2=0,`$
and set
$`D(pq)=\mathrm{Det}(x^p)\mathrm{Trace}((x^p)^1x^q)=x_{11}^px_{22}^q+x_{22}^px_{11}^qx_{12}^px_{21}^qx_{21}^px_{12}^q.`$
We set
$$M_{pqr}=\left(\begin{array}{cccc}x_{11}^p& x_{12}^p& x_{21}^p& x_{22}^p\\ x_{11}^q& x_{12}^q& x_{21}^q& x_{22}^q\\ x_{11}^r& x_{12}^r& x_{21}^r& x_{22}^r\end{array}\right),p,q,r\{1,2,3,4\},$$
and let $`(ijk)`$ denote the $`3\times 3`$ minor determinant of $`M_{pqr}`$ induced from the $`i,j,k`$-th column vectors. Put
$$D(pqr)=(234)(123)(134)(124).$$
###### Lemma 4.2.
We have the following :
(1) $`D(pq)=D(qp),D(pp)=2\mathrm{d}\mathrm{e}\mathrm{t}(x^p)`$
(2) $`H_p`$ is irreducible if and only if
$$D(pp)0$$
(3) Let $`H_p`$ and $`H_q`$ be different and both irreducible . Then they have different intersection points if and only if
$$D(pq)^2D(pp)D(qq)0.$$
(4) We have $`H_pH_qH_r=\mathrm{}`$ if and only if $`D(pqr)0`$.
\[Proof\].
The first two claims are obvious. So we consider the third statement.
$`H_p`$ is expressed in the form
$$\frac{\xi _2}{\xi _1}=\frac{x_{11}^p\eta _1+x_{12}^p\eta _2}{x_{21}^p\eta _1+x_{22}^p\eta _2}.$$
So we obtain the intersections $`H_pH_q`$ from
$$\frac{x_{11}^p\eta _1+x_{12}^p\eta _2}{x_{21}^p\eta _1+x_{22}^p\eta _2}=\frac{x_{11}^q\eta _1+x_{12}^q\eta _2}{x_{21}^q\eta _1+x_{22}^q\eta _2}.$$
Then the intersection comes from the eigen vector of $`(x_p)^1x^q`$. Hence we get the required condition.
Next we examine the last statement. Generally we have
$$M_{pqr}\left(\begin{array}{c}(234)\\ (134)\\ (124)\\ (123)\end{array}\right)=\left(\begin{array}{c}0\\ 0\\ 0\end{array}\right).$$
Note that we have $`H_pH_qH_r\mathrm{}`$ if and only if there exists a solution $`\zeta ={}_{}{}^{t}(\zeta _1,\zeta _2,\zeta _3,\zeta _4)`$ with $`\zeta _1\zeta _4\zeta _2\zeta _3=0`$ for
$$M_{pqr}\left(\begin{array}{c}\zeta _1\\ \zeta _2\\ \zeta _3\\ \zeta _4\end{array}\right)=\left(\begin{array}{c}0\\ 0\\ 0\end{array}\right)$$
If $`\mathrm{rank}M_{pqr}<3`$ we can find easily such a solution $`\zeta `$. In case $`\mathrm{rank}M_{pqr}=3`$ we have 1 dimensional solution space for $`M_{pqr}\zeta =0`$. So $`{}_{}{}^{t}((234),(134),(124),(123))`$ becomes a required solution only when $`D(pqr)=0`$.
q.e.d
###### Theorem 4.
Let $`x=(x^1,x^2,x^3,x^4)\mathrm{M}(2,𝐂)^4`$ be a point on $`X^0`$, namely $`x`$ satisfies the conditions:
(g1). Any $`H_p`$ is irreducible i.e. $`D(pp)0`$,
(g2). $`H_p`$ and $`H_q`$ have 2 different intersections i.e. $`D(pq)^2D(pp)D(qq)0`$ for any $`ij`$
(g3). $`H_pH_qH_r=\mathrm{}`$ i.e $`D(pqr)0`$ for any $`(pq)(qr)(rp)0`$
Then any local solution of the system $`E_1+E_2`$ around the point $`xX^0`$ is locally holomorphic.
\[Proof\]. We show $`B_x=\{0\}`$ in several steps.
(Step1 ). Let us consider the Segre embedding
$$\psi :𝐏^1\times 𝐏^1𝐏^3,$$
that is defined by
$$[\xi _1,\xi _2]\times [\eta _1,\eta _2]\left(\begin{array}{cc}\xi _1\eta _1& \xi _1\eta _2\\ \xi _2\eta _1& \xi _2\eta _2\end{array}\right).$$
Note that the point $`\left(\begin{array}{cc}\xi _{11}& \xi _{12}\\ \xi _{21}& \xi _{22}\end{array}\right)𝐏^3`$ belongs to $`\mathrm{Im}(\psi )`$ if and only if $`\xi _{11}\xi _{22}\xi _{12}\xi _{21}=0`$. Let $`\mathrm{\Delta }`$ denote the diagonal map
$$\mathrm{\Delta }:𝐏^3𝐏^3\times 𝐏^3\times 𝐏^3\times 𝐏^3,P(P,P,P,P).$$
Set
$$\xi ^p=\left(\begin{array}{cc}\xi _{11}^p& \xi _{12}^p\\ \xi _{21}^p& \xi _{22}^p\end{array}\right),$$
and regard $`\xi =(\xi ^1,\xi ^2,\xi ^3,\xi ^4)`$ as a homogeneous coordinate on $`𝐏^3\times 𝐏^3\times 𝐏^3\times 𝐏^3`$. Then the system $`E_2`$ in Proposition 4.1 determines exactly the image of $`\mathrm{\Delta }\psi `$.
(Step 2). Let $`\xi =(\xi ^1,\xi ^2,\xi ^3,\xi ^4)`$ be a point on $`B_x`$. Then at least two of $`\{\xi ^p\}`$ should be $`O`$. Suppose in contrary three of them , saying $`\xi ^1,\xi ^2,\xi ^3`$, are not equal to $`O`$. According to the argument in Step 1 they have the pull backs in $`𝐏^1\times 𝐏^1`$, and these are the same point $`P=([\xi _1,\xi _2],[\eta _1,\eta _2])`$. Replace $`\xi ^p`$ by $`\psi (P)`$ in (4.9), which yields $`PH_p`$. Hence we obtain $`PH_1H_2H_3`$. This contradicts $`(\mathrm{g3})`$.
(Step 3). Let $`\xi =(\xi ^1,\xi ^2,\xi ^3,\xi ^4)`$ be a point on $`B_x`$. If three of $`\{\xi ^p\}`$ are $`O`$, then all of them equal $`O`$.
Suppose $`\xi ^2=\xi ^3=\xi ^4=0`$. Then the equation (4.10) reduces to
$$x^1{}_{}{}^{t}\xi _{}^{1}=O.$$
Because $`x^1`$ is assumed to be invertible, we obtain $`\xi ^1=O`$.
(Step 4). We don’t have the case $`\xi ^10,\xi ^20,\xi ^3=\xi ^4=0`$.
Suppose it happens, then we have $`\xi ^2=c\xi ^1`$ for some constant $`c(0)`$. The two equations (4.10) and (4.11) are expressed in the form
$$x^1{}_{}{}^{t}\xi _{}^{1}+x^2{}_{}{}^{t}\xi _{}^{2}=0,{}_{}{}^{t}\xi _{}^{1}x^1+{}_{}{}^{t}\xi _{}^{2}x^2=0.$$
Because we supposed $`x^2`$ to be invertible, we have
$$(x^2)^1x^1{}_{}{}^{t}\xi _{}^{1}={}_{}{}^{t}\xi _{}^{1}x^1(x^2)^1.$$
So we get 4 linear equations for 4 unknowns $`\xi _{11}^1,\xi _{12}^1,\xi _{21}^1,\xi _{22}^1`$. Together with the 2 equations coming from the first equation (4.9) for $`p=1,2`$ we obtain the system of linear equations
$$M\left(\begin{array}{c}\xi _{11}^1\\ \xi _{12}^1\\ \xi _{21}^1\\ \xi _{22}^1\end{array}\right)=\left(\begin{array}{c}0\\ 0\\ 0\\ 0\\ 0\\ 0\end{array}\right)().$$
By a direct calculation we have
$$M=\left(\begin{array}{cccc}x_{12}^1x_{21}^2x_{21}^1x_{12}^2& x_{12}^1x_{22}^2x_{22}^1x_{12}^2& x_{22}^1x_{21}^2x_{21}^1x_{22}^2& 0\\ x_{11}^1x_{12}^2x_{12}^1x_{11}^2& 0& x_{11}^1x_{22}^2x_{22}^1x_{11}^2& x_{12}^1x_{22}^2x_{22}^1x_{12}^2\\ x_{21}^1x_{11}^2x_{11}^1x_{21}^2& x_{22}^1x_{11}^2x_{11}^1x_{22}^2& 0& x_{22}^1x_{21}^2x_{21}^1x_{22}^2\\ 0& x_{11}^1x_{12}^2x_{12}^1x_{11}^2& x_{21}^1x_{11}^2x_{11}^1x_{21}^2& x_{21}^1x_{12}^2x_{12}^1x_{21}^2\\ x_{11}^1& x_{12}^1& x_{21}^1& x_{22}^1\\ x_{11}^2& x_{12}^2& x_{21}^2& x_{22}^2\end{array}\right).$$
Let $`(i,j,k,\mathrm{})`$ denote the $`4\times 4`$ minor of $`M`$ obtained by taking $`i,j,k,\mathrm{}`$-th row vectors. By assumption the system ($``$) has a nontrivial solution $`\xi ^1`$, so every $`(i,j,k,\mathrm{})`$ should be $`0`$. By a direct calculation we have the following:
$`(1256)=(x_{12}^1x_{22}^2x_{22}^1x_{12}^2)(D(12)^2D(11)D(22)),`$
$`(1356)=(x_{22}^1x_{21}^2x_{21}^1x_{22}^2)(D(12)^2D(11)D(22)),`$
$`(1456)=(x_{21}^1x_{12}^2x_{12}^1x_{21}^2)(D(12)^2D(11)D(22)),`$
$`(2356)=(x_{11}^1x_{22}^2x_{22}^1x_{11}^2)(D(12)^2D(11)D(22)),`$
$`(2456)=(x_{12}^1x_{11}^2x_{11}^1x_{12}^2)(D(12)^2D(11)D(22)),`$
$`(3456)=(x_{11}^1x_{21}^2x_{21}^1x_{11}^2)(D(12)^2D(11)D(22)),`$
and other $`4\times 4`$ minors are $`0`$. Here recall the 2nd assumption for $`x`$. So we have $`(D(12)^2D(11)D(22))`$ is not $`0`$. Hence we obtain
$$x_{ij}^1x_k\mathrm{}^2=x_{ij}^2x_k\mathrm{}^1$$
for any indices $`i,j,k,\mathrm{}\{1,2\}`$. It means $`x^2=cx^1`$ and $`H_1=H_2`$. This is a contradiction.
By the combination of the above arguments we obtain $`B_x=\{0\}`$.
q.e.d.
The system of partial differential equations $`E_1+E_2`$ is closely related to the GKZ hypergeometric system introduced by Gel’fand, Kapranov and Zelevinski. Let us explain the relation and evaluate the dimension of the solution space of $`E_1+E_2`$.
Define a set of operators $`E_1(\mathrm{},m)`$ and $`E_1^{}(\mathrm{},m)`$ by
$`E(\mathrm{},m)`$ $`:`$ $`{\displaystyle \underset{p,j}{}}x_\mathrm{}j^p_{mj}^p+\delta _\mathrm{}m,`$
$`E^{}(\mathrm{},m)`$ $`:`$ $`{\displaystyle \underset{p,j}{}}x_j\mathrm{}^p_{jm}^p+\delta _\mathrm{}m.`$
Let $`D`$ be the Weyl algebra
$$𝐂x_{11}^1,x_{12}^1,x_{21}^1,x_{22}^1,x_{11}^2,\mathrm{},x_{22}^4,_{11}^1,_{12}^1,_{21}^1,_{22}^1,_{11}^2,\mathrm{},_{22}^4.$$
Consider the GKZ-hypergeometric ideal $`H_A(\beta )`$ in $`D`$ associated to the matrix
$$A=\left(\begin{array}{cccccccccccccccc}1& 1& 1& 1& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 1& 1& 1& 1& 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0& 1& 1& 1& 1& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 1& 1& 1& 1\\ 1& 1& 0& 0& 1& 1& 0& 0& 1& 1& 0& 0& 1& 1& 0& 0\\ 1& 0& 1& 0& 1& 0& 1& 0& 1& 0& 1& 0& 1& 0& 1& 0\end{array}\right)$$
and $`\beta =(1/2,1/2,1/2,1/2,1,1)`$.
###### Proposition 4.2.
Our system of partial differential equations $`E_1`$ and $`E_2`$ consists of three groups of operators
$$A\theta {}_{}{}^{t}\beta ,$$
$$I_A=\{^u^vAu=Av,u,v𝐍_0^{16}\},$$
$$E_1(\mathrm{},m)\mathrm{and}E_1^{}(\mathrm{},m)(\mathrm{}m).$$
Here, we denote by $`\theta `$ the column vector $`(x_{jk}^i_{jk}^i)`$ of Euler operators of length 16.
The GKZ-hypergeometric ideal $`H_A(\beta )`$ is the left ideal in $`D`$ generated by $`A\theta {}_{}{}^{t}\beta `$ and $`I_A`$. Since the toric ideal $`I_A`$ is homogeneous, the $`D`$-module $`D/H_A(\beta )`$ is regular holonomic on $`X=𝐂^{16}`$ (Hotta’s theorem, see, e.g., \[14, p.82\]). Our toric ideal $`I_A`$ satisfies the following properties:
1. The initial ideal of $`I_A`$ with respect to the reverse lexicographic order is generated by square free monomials.
2. The toric ideal is Cohen-Macaulay.
3. The multiplicity of $`I_A`$ is $`20`$.
4. The variety $`V(I_A)`$ admits a natural action of $`(𝐂^{})^4`$ and $`V(I_A)/(𝐂^{})^4`$ is isomorphic to $`\mathrm{Im}(\mathrm{\Delta }\psi )`$ in Theorem 4.
The first statement can be easily checked by Buchberger’s criterion. The second fact follows from the first (see, e.g., \[14, p.153\]). The third fact can be shown by computing the Hilbert polynomial of $`I_A`$ on computers. Therefore, by theorems due to Gel’fand, Kapranov and Zelevinsky, the rank of the solution space of $`H_A(\beta )`$ is $`20`$ and the singular locus agrees with the zero set of the principal $`A`$-determinant (see, e.g., \[14, p.173\]).
We denote by E the left ideal in $`D`$ generated by first order operators $`E_1(\mathrm{},m)`$, $`(\mathrm{}m)`$ and the GKZ-hypergeometric ideal $`H_A(\beta )`$.
###### Theorem 5.
The $`D`$-module $`D/E`$ is regular holonomic on $`X`$, and the rank of the solution space of $`E=E_1+E_2`$ is equal to 8.
\[Proof\]. Since $`H_A(\beta )E`$ holds and the GKZ system is regular holonomic, the $`D`$-module $`D/E`$ is also regular holonomic on $`X`$. The differential operators $`E_1(\mathrm{},m),(\mathrm{}m)`$ are used to extract the space of the period maps from the $`20`$ dimensional solution space of the GKZ system.
Next let us show the second statement. Already in Section 3 we constructed 8 independent periods. So it is enough to say that ”at most 8 dimensional”. Although it is possible in principle to evaluate the rank by computer and Oaku’s algorithm (see, e.g., \[14, p.31\]), we could not evaluate it because of an exhaustion of memory. So, we try to find sufficiently many initial terms for the left ideal $`E`$ and a suitable weight. For this system, we chose a weight $`w^1=1w,w^2=4w,w^3=9w,w^4=16w`$, $`w=\left(\begin{array}{cc}1& 2\\ 0& 4\end{array}\right)`$. Put $`W=(w^1,w^2,w^3,w^4)`$. The initial ideal generated by $`\mathrm{in}_{(W,W)}(E_1)`$ and $`\mathrm{in}_{(W,W)}(E_2)`$ has the rank $`20`$, which means that we do not have sufficiently many initial terms. We computed a partial Gröbner basis of $`E_1`$ and $`E_2`$ with the weight $`(W,W)`$ up to the degree 7 in the homogenized Weyl algebra. The ideal generated by the initial terms of the partial Gröbner basis has rank $`8`$. It follows from Theorem 2.5.1 of and the argument of the regular holonomicity that the rank is bounded by $`8`$.
q.e.d.
According to Theorem 4 and Theorem 5 the solution space of $`E_1+E_2`$ looks like a vector bundle over $`X^{}`$ of rank 8. Let us take a $`P`$-marking $`(S(x_0),\phi ,P)`$ , then we can choose the basis for the solution space of $`E_1+E_2`$
$$\{_{\phi ^1(\stackrel{ˇ}{\mathrm{\Gamma }}_1)}\mathrm{\Omega },\mathrm{},_{\phi ^1(\stackrel{ˇ}{\mathrm{\Gamma }}_8)}\mathrm{\Omega }\}$$
at $`x_0X^{}`$. So the system $`E_1+E_2`$ induces a representation of $`\pi _1(X^{},x_0)`$ over $`GL(8,𝐙)`$.
###### Definition 4.1.
The monodromy group Mono for $`(,X^{})`$ is the image of this representation.
###### Remark 4.1.
Set $`\mathrm{\Sigma }=(GL(2,𝐂)X^{}/GL(2,𝐂))/(𝐂^{})^4`$. According to Lemma 4.1 we have the same period along the orbit of the actions of $`GL(2,𝐂)`$ and $`(𝐂^{})^4`$. So the above monodromy representation reduces to that of $`\pi _1(\mathrm{\Sigma },)`$.
Let $`B_r`$ be the sublattice of $`L`$ generated by the elements $`f_s,f_t,g_i(i=1,\mathrm{},4)`$. Put
$$O(L,Br)=\{gO(L):g(P)=P,g_{Br}=\mathrm{id}\}.$$
We define
$$O^+(L,Br)=\{gO(L,Br):g_T(D^+)=D^+\}.$$
The group $`\pi _1(X^{},x_0)`$ acts on $`LH_2(S(x_0),𝐙)`$ by the natural way, and the branch locus is fixed under this action. So $`\mathrm{Mono}`$ is identified with a certain subgroup of $`O^+(L,Br)`$.
###### Lemma 4.3.
Let $`(S_1,\phi _1,P)`$ and $`(S_2,\phi _2,P)`$ be $`P`$-markings, and let $`x_1,x_2`$ be corresponding points on $`X^{}`$. Suppose there is an isomorphism $`\rho :S_1S_2`$ such that we have $`\phi _2\rho _{}\phi _1^1O^+(L,Br)`$. Then the points $`x_1`$ and $`x_2`$ determine the same point on $`\mathrm{\Sigma }=(GL(2,𝐂)X^{}/GL(2,𝐂))/(𝐂^{})^4`$.
\[Proof\].
Because we have $`\rho _{}\phi _1^1(f_s)=\phi _2^1(f_s)`$ and $`\rho _{}\phi _1^1(f_s)=\phi _2^1(f_s)`$, the isomorphism $`\rho `$ preserves the covering structure over $`𝐏^1\times 𝐏^1`$. Moreover we have $`\rho _{}\phi _1^1(g_i)=\phi _2^1(g_i)(i=1,2,3,4)`$. Hence $`\rho `$ preserves the branch locus with the same numbering order. So $`x_1`$ and $`x_2`$ have the required property.
q.e.d.
###### Theorem 6.
We have
$$O^+(L,Br)_T=\mathrm{Mono}.$$
\[Proof\]. It is enough to show that $`O^+(L,Br)_T\mathrm{Mono}`$. Let us fix a $`P`$-marking $`(S_0,\phi _0,P)`$ corresponding to the initial point $`x_0X^{}`$. Let $`g`$ be an arbitrary element of $`O^+(L,Br)`$, and put $`g^{}=g_T`$. Let $`\eta _0D^+`$ be the period point determined by $`(S_0,\phi _0,P)`$, and set $`\eta _1=g^{}(\eta _0)`$. Let us take an oriented arc $`r`$ in $`D^+`$ that starts from $`\eta _0`$ and ends at $`\eta _1`$. Because of the surjectivity of the period map, we can find a $`P`$-marking $`(S(\eta ),\phi (\eta ),P)`$ for every point $`\eta `$ on $`r`$. Set $`(S_1,\phi _1,P)`$ be the terminal marking and let $`x_1`$ be the corresponding point on $`X^{}`$. Because of the injectivity of the period map we have $`(S_1,\phi _1,P)=(S_1,\phi _1g^1,P)`$. Then by virtue of the above Lemma $`x_0`$ and $`x_1`$ determine the same point on $`\mathrm{\Sigma }`$. Hence we get unique closed arc $`\gamma `$ in $`\mathrm{\Sigma }`$ corresponding to $`r`$. It means that $`g^{}`$ is the monodromy transformation coming from the arc $`\gamma `$.
q.e.d.
## 5 Periods and Kuga-Satake varieties
In this section we construct the abelian variety attached to the extended double $`4H`$ surface starting from its period. The reader will find the method in and . The detailed calculation and argument is exposed in also.
Let us consider the lattice $`T`$ defined by the intersection matrix $`A=U(2)U(2)(2I_4)`$ and put $`V_k=Tk(k=𝐑\mathrm{or}𝐐)`$. Let $`Q(x)`$ denote the quadratic form on $`T`$ and at the same time on $`V_k`$. Let $`Tens(T)`$ and $`Tens(V_k)`$ be the corresponding tensor algebras. And we let $`Tens^+(T)`$ and $`Tens^+(V_k)`$ denote the subalgebras composed of the parts with even degree in $`Tens(T)`$ and $`Tens(V_k)`$ respectively. We consider the two sided ideal $`I`$ in $`Tens^+(V_k)`$ generated by the elements of the form $`xxQ(x)`$ for $`xV_k`$, and the ideal $`I_𝐙`$ in $`Tens(T)`$ by the same manner. The corresponding even Clifford algebra is defined by
$$C^+(V_k,Q)=Tens^+(V_k)/I.$$
By the same manner we define the even Clifford algebra over $`𝐙`$ by
$$C^+(T,Q)=Tens^+(T)/I_𝐙.$$
We note that $`C^+(V_𝐑,Q)`$ is a $`128`$ dimensional real vector space and $`C^+(T,Q)`$ is a lattice in it. So we obtain a real torus
$$𝒯_𝐑=C^+(V_𝐑,Q)/C^+(T,Q).$$
Let $`𝐅`$ denote the quaternion algebra
$$𝐐𝐐i𝐐j𝐐ij$$
with $`i^2=j^2=1`$. By some routine calculations of Clifford algebra we obtain the following.
###### Proposition 5.1.
We have the isomorphism of algebras $`C^+(V_𝐐,Q)\mathrm{M}(4,𝐅)\mathrm{M}(4,𝐅)`$.
Let a complex vector $`\underset{¯}{\eta }=(\eta _1,\mathrm{},\eta _8)`$ be a representative of a point $`\eta =[\eta _1,\mathrm{},\eta _8]D^+`$. So it has an ambiguity of the multiplication by a non zero complex number. Put $`\underset{¯}{\eta }=s+it(s,t𝐑^8)`$. If we impose the condition $`(st)^2=1`$ in $`C^+(V_𝐑,Q)`$, the representative is uniquely determined up to a multiplication by a complex unit. We denote it by
$$\underset{¯}{\eta }=m_1(\eta )+im_2(\eta ),$$
and put
$$m(\eta )=m_1(\eta )m_2(\eta ).$$
It is uniquely determined by $`\eta `$ without any ambiguity. According to the imposed condition the element $`m(\eta )C^+(V_𝐑,Q)`$ defines a complex structure on $`C^+(V_𝐑,Q)`$ by the left action. It induces a complex structure on the real torus $`𝒯_𝐑`$ also. We denote this complex torus by $`(T,m(\eta ))`$.
Let $`\{\epsilon _1,\mathrm{},\epsilon _8\}`$ be the basis of $`T`$ with the intersection matrix $`U(2)U(2)(2I_4)`$. And let $`\{e_1,\mathrm{},e_8\}`$ be a orthonormal basis of $`V`$ given by
$$(e_1,\mathrm{},e_8)=(\epsilon _1,\mathrm{},\epsilon )(\left(\begin{array}{cccc}\frac{1}{2}& \frac{1}{2}& 0& 0\\ 0& 0& \frac{1}{2}& \frac{1}{2}\\ \frac{1}{2}& \frac{1}{2}& 0& 0\\ 0& 0& \frac{1}{2}& \frac{1}{2}\end{array}\right)(I_4)).$$
Then the corresponding intersection matrix takes the form $`I_2(I_2)(2I_4)`$.
Let $`\iota `$ be an involution on $`C^+(V,Q)`$ induced from the transformation
$$\iota :e_{i_1}e_{i_2}\mathrm{}e_{i_k}e_{i_k}\mathrm{}e_{i_2}e_{i_1}$$
for the basis. Set $`\alpha =4e_2e_1`$. According to the method in \[St\] we know that
$$E(x,y)=tr(\alpha x^\iota y)$$
determines a Riemann form on $`(L,m(\eta ))`$. We denote this abelian variety by $`A^+(\eta )`$, that is so called the Kuga-Satake variety attached to the $`K3`$ surface corresponding to the period $`\eta `$. In this way we can construct a family of abelian varieties
$$𝒜^+=\{A^+(\eta ):\eta D^+\}$$
induced from the lattice $`T`$ parameterized by the domain $`D^+`$. We can construct the ”conjugate family”
$$𝒜^{}=\{A^{}(\eta ):\eta D^{}\}$$
parameterized by
$$D^{}=\{\eta =[\eta _1,\mathrm{},\eta _8]:{}_{}{}^{t}\eta A\eta =0,{}_{}{}^{t}\overline{\eta }A\eta >0,\mathrm{}(\eta _3/\eta _1)<0\}$$
by the same procedure with the Riemann form $`E^{}(x,y)=\mathrm{tr}(\alpha x^\iota y)`$. The right action of $`C^+(V_𝐐,Q)`$ on $`(V,T,m(\eta ))`$ commutes with the left action of $`\alpha (\eta )`$. So we have
$$C^+(T_𝐐)\mathrm{End}(\mathrm{A}^\pm (\eta ))𝐐$$
for any $`A^\pm (\eta )`$. For a general member $`\eta 𝒟^+`$, the endmorphism ring is given by
$$\mathrm{End}_𝐐(\mathrm{A}(\eta ))=\mathrm{End}(\mathrm{A}(\eta ))𝐐\mathrm{C}^+(\mathrm{V}_𝐐).$$
According to Proposition 5.1 we obtain
###### Theorem 7.
For a general member $`\eta 𝒟^+`$, $`A^+(\eta )`$ is isogeneous to a product of abelian varieties $`(A_1(\eta )\times A_2(\eta ))^4`$ where $`A_1(\eta )`$ and $`A_2(\eta )`$ are 8-dimensional simple abelian varieties with $`\mathrm{End}_𝐐(\mathrm{A}_\mathrm{i}(\eta ))=𝐅(\mathrm{i}=1,2)`$.
###### Remark 5.1.
Here we describe the relation between $`A_1(\eta )`$ and $`A_2(\eta )`$. Now we define the linear involution $``$ on $`V_𝐑`$ by
$$e_1^{}=e_1\mathrm{and}e_i^{}=e_i(i=2,\mathrm{},8).$$
It can be extended on $`C^+(V_𝐑,Q)`$ as an autmorphism of algebra. We define an involution $`\sigma `$ on $`𝒟`$ :
$$\sigma :𝒟𝒟,(\eta _1,\mathrm{},\eta _8)(\eta _2,\eta _1,\eta _3,\mathrm{}\eta _8).$$
So we have $`𝒟_+^\sigma =𝒟_{}`$. It is easy to check that we have
$$A_2(\eta )A_1(\eta ^\sigma ),A_1(\eta )A_2(\eta ^\sigma ),$$
where $``$ indicates the isogenous relation.
T. Tsutsui
Department of Mathematics and Informatics
Faculty of Science
Chiba University
Yayoi-cho Inage-ku Chiba 263-8522, JAPAN
K. Koike, H. Shiga
Graduate School of Science
Chiba University
Yayoi-cho Inage-ku Chiba 263-8522, JAPAN
N. Takayama
Department of Mathematics
Faculty of Science
Kobe University
Rokko Kobe 657-8501, JAPAN
|
warning/0003/gr-qc0003023.html
|
ar5iv
|
text
|
# Entropy using Path Integrals for Quantum Black Hole Models
## Abstract
Several eigenvalue equations that could describe quantum black holes have been proposed in the canonical quantum gravity approach. In this paper, we choose one of the simplest of these quantum equations to show how the usual Feynman’s path integral method can be applied to obtain the corresponding statistical properties. We get a logarithmic correction to the Bekenstein-Hawking entropy as already obtained by other authors by other means.
In the early seventies, in his insightful work Bekenstein proposed the quantization of a black hole. He suggested that its surface gravity is proportional to its temperature and that the area of its event horizon is proportional to its entropy. In his remarkable work he conjectured that the horizon area of non-extremal black holes plays the role of a classical adiabatic invariant. He concluded that the horizon area should have a discrete spectrum with uniformly spaced eigenvalues. Another result at that same time, was the discovery of the mechanical laws in the framework of general relativity, that govern non-extremal black holes by . These laws have a striking analogy with those of thermodynamics. This similarity was not well understood until Hawking, in his seminal work a year later, discovered that black holes evaporate, they radiate as black bodies with a temperature proportional to the surface gravity. One can then argue that the laws of black hole mechanics have a real thermodynamical meaning and that the entropy of a black hole is one fourth of its horizon area.
These remarkable results of nearly three decades ago, point out to a deeper relation among classical gravitation, quantum mechanics and statistical properties. Hawking’s semi-classical calculations allowed to interpret the relation between classical black hole mechanics and its thermodynamics leading to the celebrated entropy area expression.
During these years, different approaches have emerged to try to understand the interplay between the quantum and classical descriptions, and the statistical properties of black holes. We briefly mention the main quantum gravity non-perturbative formalisms. String theory, whose building blocks are essentially D-branes , loop quantum gravity , and canonical quantum gravity, having as elementary constituents the quantum excitations of the geometry itself . These proposals emerge essentially from different principles. String theory is sought to be fundamental, it provides the precise expression for the temperature and entropy. Loop and canonical gravity do not provide the corresponding proportionality constants unambiguously. These, however, deal directly with the curved black hole geometry. In string theory, one carries out, for example, the calculation of Hawking´s radiation in flat space. Moreover, in string theory these calculations have been up to now, only possible for extremal and nearly extremal black holes .
In this paper, we are interested in canonical quantum gravity treatments where a Hamiltonian quantum theory of spherically symmetric vacuum spacetime can be defined (for our purposes we could also consider a ball or shell of dust collapsing to a black hole ). On the other hand, supposing a uniformly spaced area spectrum , the Schwarzschild black hole has been treated as a microcanonical ensemble . Also considering a mini-superspace approach and by means of statistical techniques previously used in hadron physics, the usual area of black hole thermodynamics is recovered . Moreover, assuming that the area spectrum of the black hole is uniformly spaced, a grand canonical ensemble has been considered with the ADM mass (the Hamiltonian) and the horizon area as separately observables . It is argued that in this way the partition function is not divergent. However, the result in that work, a logarithmic correction to the Bekenstein-Hawking entropy, is the same already obtained by Kastrup by means of an analytic continuation approach.
We will take into account the results mentioned in the previous paragraph, and that the quantum Hamiltonians defined to study spherically symmetric vacuum spacetime should describe, for an observer at rest far away, the quantum mechanics of a black hole . We could have taken also models of a shell or ball collapsing on its way to a black hole formation . In our procedure it is not absolutely necessary that the resulting area spectrum is uniformly spaced. However we need and eigenvalue equation for the Arnowitt-Deser-Misner (ADM) mass of the hole or its horizon area, which we do not consider independent variables. We will proceed, with one of the most simple examples, by taking the Hamiltonian operator defined by Mäkela and Repo , taking $`Q=0`$, for the Reissner- Nordström black hole Hamiltonian. Our approach is straightforward, an eigenvalue equation has been defined for a particular parametrization of the variables of interest, where for a Schwarzschild black hole a certain factor ordering has been chosen. The quantum-mechanical model is represented by an eigenvalue equation for a linear harmonic oscillator. For such quantum equations, we already know how to obtain the corresponding statistical mechanics . For $`Q0`$ and other variable parametrizations and factor orderings, the same method outlined here can be followed. We show how to proceed, with a simple example.
With Hamiltonians involving only momenta and coordinates, we can use path integrals to get the corresponding statistical properties. The equation of interest, in this paper, can be obtained from that in Mäkela and Repo , for $`Q=0`$ is
$$\frac{\mathrm{}^2G^2}{c^6}a^{s1}\frac{d}{da}\left(a^s\frac{d}{da}\mathrm{\Psi }(a)\right)=(a\frac{2GM}{c^2})\mathrm{\Psi }(a),$$
(1)
where $`a`$ and $`P_a^2=`$ -$`\frac{\mathrm{}^2G^2}{c^6}a^s\frac{d}{da}\left(a^s\frac{d}{da}\mathrm{\Psi }(a)\right)`$ are phase coordinates obtained from the phase space coordinates $`m`$ and $`P_m`$, by means of an appropriate canonical transformation and $`m(t):=M(t,r)`$ and $`p_m(t):=_{\mathrm{}}^{\mathrm{}}𝑑rP_M(t,r)`$. The variable $`m`$ can be defined as the mass $`M`$ of the hole when Einstein’s equations are satisfied , $`s`$ is a factor ordering parameter. In particular, if we choose $`s=2`$ and identifying $`R_s=\frac{2GM}{c^2}`$, one gets
$$\frac{\mathrm{}^2G^2}{c^6}\frac{1}{a}\left(\frac{d^2}{da^2}+\frac{2}{a}\frac{d}{da}\right)\mathrm{\Psi }(a)=(aR_s)\mathrm{\Psi }(a).$$
(2)
We now transform
$$\mathrm{\Psi }(a)=\frac{1}{a}U(a)\mathrm{and}x=aR_s,$$
(3)
where the variable $`x`$ describes the gravitational degrees of freedom of the Schwarzschild black hole. We also introduce the appropriate constants and consider the fact that the energy of “excitation” associated with the variable $`a`$ is not positive; the physical reason for this is simply that the total energy of the black hole is included and the ADM energy is equal to zero. Then, the quantum equation (2) transforms into
$$\left(\frac{1}{2}\mathrm{}_p\mathrm{}^2E_p\mathrm{}\frac{d^2}{dx^2}+\frac{E_p\mathrm{}}{2\mathrm{}_p\mathrm{}^2}x^2\right)U(x)=\frac{R_s}{4\mathrm{}_p\mathrm{}}E_sU(x),$$
(4)
where $`E_s=Mc^2`$ is the black hole ADM energy and $`\mathrm{}_p\mathrm{}=\sqrt{\frac{G\mathrm{}}{c^3}}`$ and $`E_p\mathrm{}=\sqrt{\frac{c^5\mathrm{}}{G}}`$ are Planck’s length and energy respectively. Being (4) a quantum linear oscillator, one immediately gets
$$\frac{R_s(n)}{4\mathrm{}_p\mathrm{}}E_s(n)=\left(n+\frac{1}{2}\right)E_p\mathrm{},$$
(5)
this result coincides with Beckenstein’s proposal .
Our procedure can be summarized as follows. We start with the Wheeler-DeWitt equation for the Schwarzschild black hole (4), we then use Feynman´s path integrals approach to statistical mechanics. This will allow us to calculate the free energy and the partition function. So that the statistical properties of the black hole, temperature and entropy will be deduced.
The classical partition function for the harmonic oscillator is given by
$$Z_{class}=\frac{1}{\beta \mathrm{}\omega }.$$
(6)
Applying the path integral formalism, the changes in the potential function due to quantum mechanical effects can be included by means of the use of . The “corrected” potential for the oscillator is
$$V=\frac{m\omega ^2}{2}\left(x^2+\frac{\beta \mathrm{}^2}{12m}\right),$$
(7)
which gives the partition function
$$Z_{approx}=\frac{e^{\frac{(\beta \mathrm{}\omega )^2}{24}}}{\beta \mathrm{}\omega }.$$
(8)
It should be pointed out that this result is for the Euclidean space.
We now utilize this method, that has been so successful in quantum mechanics, in the quantum gravity model of interest. Obviously, we do not have a quantum theory of gravity, but we have a quantum harmonic oscillator eigenvalue equation for the Schwarzschild black hole (4), that provides the expected energy eigenvalues (5) . We should note that the Schrödinger equation (4) for black holes, corresponds to a quantum oscillator in Lorentzian space. Consequently, we need to modify the partition function defined in (8) to go from the Lorentzian to the Euclidean world. After a rotation $`\beta i\beta `$, which is equivalent to a Wick rotation $`ti\tau `$, the partition function has the following form
$$Z_{approx}^{}=i\frac{e^{\frac{(\beta \mathrm{}\omega )^2}{24}}}{\beta \mathrm{}\omega }.$$
(9)
As already stated by Kastrup , this analytic continuation should be used to get the thermodynamical properties of interest, so we have for the internal energy
$$\overline{E}=\frac{\mathrm{ln}(ImZ_{approx}^{})}{\beta }=Mc^2,$$
(10)
which should be the internal gravitational energy of the black hole. According to equation (4) the frequency of our harmonic oscillator results in $`\mathrm{}\omega =\sqrt{\frac{3}{2\pi }}E_p\mathrm{}`$, so from equations (9) and (10), we get
$$\frac{E_p\mathrm{}^2}{8\pi }\beta ^2Mc^2\beta 1=0,$$
(11)
the positive solution for this equation for the case $`E_p\mathrm{}Mc^2`$ leads to
$$\beta =\frac{8\pi Mc^2}{E_p\mathrm{}^2}\left[1+\frac{1}{8\pi }\left(\frac{E_p\mathrm{}}{Mc^2}\right)^2\right]=\beta _H\left[1+\frac{1}{\beta _HMc^2}\right],$$
(12)
which is Hawking´s temperature $`\beta _H=\frac{1}{kT_H}`$, plus a small correction, that is the same as that obtained by Kastrup . This calculation is straightforward, that is because the quantum equation corresponds to a linear oscillator, but even if we would have a more complicated equation, we would only need to calculate a more difficult path integral. This would be the case, for example, for the Reissner-Nordström black hole $`Q0`$ .
The entropy for the model, can be obtained again, by means of the rotation $`(\beta i\beta `$, $`kik)`$, we get
$$Im\left(\frac{S}{ik}\right)=\left[\mathrm{ln}\left(ImZ_{approx}^{}\right)+\beta \overline{E}\right],$$
(13)
at $`\beta \mathrm{}\omega =\beta \sqrt{\frac{3}{2\pi }}E_{pl}`$. By putting the expression in terms of the partition function
$$Im\left(\frac{S}{ik}\right)=\left[\mathrm{ln}\left(ImZ_{approx}^{}\right)\beta \frac{\mathrm{ln}(ImZ_{approx}^{})}{\beta }\right],$$
(14)
and substituting equation (8) we get
$$\frac{S}{k}=\frac{A_s}{4l_p\mathrm{}^2}\left[1+\frac{1}{8\pi }\frac{E_p\mathrm{}^2}{\left(Mc^2\right)^2}\right]^2+\frac{1}{2}\mathrm{ln}\left(\frac{A_s}{4l_p\mathrm{}^2}\left[1+\frac{1}{8\pi }\frac{E_p\mathrm{}^2}{\left(Mc^2\right)^2}\right]^2\right)+\frac{1}{2}\mathrm{ln}(24)1,$$
(15)
where $`A_s=4\pi R_s^2`$ is the area of the horizon.
In terms of the Bekenstein-Hawking relation $`\frac{S_{BH}}{k}=\frac{A_s}{4l_p\mathrm{}^2}`$ and ignoring terms of higher order, we finally obtain the logarithmic correction as previously obtained by Kastrup and recently by Gour using different procedures.
$$\frac{S}{k}=\frac{S_{BH}}{k}+\frac{1}{2}\mathrm{ln}\left(\frac{S_{BH}}{k}\right).$$
(16)
We should mention that it has already been shown in that the black hole thermal fluctuations are small (and the same results follow in our case). Consequently the thermal interactions with the heat bath are small and, a canonical statistical treatment seems justified. The above previous results would then be understood as to provide the conditions on the heat bath if it is to be in thermal equilibrium with the back hole.
In this paper, we have shown how the Feynman method to obtain the statistical mechanics of a system from its fundamental quantum mechanics can be applied to some quantum black hole models. For this purpose, as an example, we have chosen one of the simplest Wheeler-DeWitt equations proposed in the literature to describe the quantum mechanics of a black hole. For any Hamiltonian in terms only of momenta and coordinates, the path integral method provides us with a powerful tool to describe physical properties of the system of interest. So, the procedure can be generalized to a wide variety of Wheeler-DeWitt equations that have been or could be proposed for black holes. One can consider models describing collapsing dust or a shell on the process to form a black hole , or those concerned with spherically symmetric vacuum spacetime . The proposals can be for black hole metrics including charge, angular momentum and even other possible physical properties. Even in these more general cases, the minisuperspace approach leaves out, necessarily, degrees of freedom that are not considered in the quantum model and for some physical purposes it will probably been unable to provide some particular desired physical results that may require to take into account, for example, particle interactions and (or) higher angular momentum modes . We have chosen an equation in minisuperspace (4) that has been parametrized in such a way that it provides Bekenstein area spectrum (5) . These area eigenvalues have been used directly, by other authors to get Hawking´s entropy and the logarithmic correction. For the purpose of getting the temperature (12) and the black hole entropy and its logarithmic correction (16), we have shown here that the minisuperspace approach, by means of Feynman´s method, gives the expected physical properties already obtained by other means .Some other Hamiltonians for more general metrics, with and without matter and their supersymmetric generalizations, are under study and will be reported elsewhere.
Acknowledgments
We thank Carlos Martinez and Pedro Ludwig for several remarks. This work was supported in part by CONACyT grant 28454E, and CONACyT-NFS grant.
|
warning/0003/hep-ph0003005.html
|
ar5iv
|
text
|
# Transverse QCD Dynamics Near the Light Cone
## 1 Introduction
The solution of Quantum Chromodynamics on the light cone is still an unsolved theoretical task the present status of which is reviewed in ref. . In a recent paper a formulation of QCD in coordinates near the light cone has been proposed which has the advantage of keeping a direct link to equal time theories. The Cauchy problem is well defined in near light cone coordinates, since the initial data are given on a space like surface. This formulation avoids the solution of constraint equations which, on the quantum field theoretic level, may be very complicated. The problem of a nontrivial vacuum appears in a solvable form related to the transverse dynamics. This is physically very appealing, since in high energy reactions the incoming particles propagate near the light cone and interact mainly exchanging particles with transverse momenta. We would like to connect successful models for the soft transverse nonperturbative dynamics of high energy reactions (cf. refs. ) to the underlying QCD Hamiltonian. In diffractive reactions, fast Lorentz contracted hadrons experience the confining forces of QCD in their transverse extensions and interact with other fast moving hadrons via soft interactions. In a theory of total cross sections the nonperturbative infrared dynamics in the transverse plane is essential. The objective of this paper is to investigate the effective transverse Hamiltonian in QCD near the light cone for small light cone momenta.
Near light cone QCD has a nontrivial vacuum which we claim cannot be neglected even in the light cone limit. We will demonstrate the existence of massless excitations in the zero mode theory. These excitations do not decouple in the light cone limit. Genuine nonperturbative techniques must be used to investigate the behavior of this limit. In principle the additional parameter which labels the coordinate system can be chosen arbitrarily. We will show that the zero mode Hamiltonian depends on an effective coupling constant containing this parameter and evolves towards an infrared fixed point. Therefore, we propose a sophisticated choice of the frame dependence which facilitates calculations in the near light cone frame. We use the infrared fixed point of the zero mode Hamiltonian to follow a trajectory in the space of couplings, where high resolution is synchronized with the light cone limit. Note, we consider the parameter associated with the frame dependence as yet another coupling constant. As shown in Ref. the zero mode sector is also relevant to the definition of M-theory in light cone coordinates.
We choose the following near light cone coordinates which smoothly interpolate between the Lorentz and light front coordinates:
$`x^t=x^+`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left\{\left(1+{\displaystyle \frac{\eta ^2}{2}}\right)x^0+\left(1{\displaystyle \frac{\eta ^2}{2}}\right)x^3\right\},`$
$`x^{}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(x^0x^3\right).`$ (1)
The transverse coordinates $`x^1`$, $`x^2`$ are unchanged; $`x^t=x^+`$ is the new time coordinate, $`x^{}`$ is a spatial coordinate. As finite quantization volume we will take a torus and its extension in “-”, as well as in “1, 2” direction is $`L`$. The scalar product of two 4-vectors $`x`$ and $`y`$ is given with $`\stackrel{}{x}_{}\stackrel{}{y}_{}=x^1y^1+x^2y^2`$ as
$`x_\mu y^\mu `$ $`=`$ $`x^{}y^++x^+y^{}\eta ^2x^{}y^{}\stackrel{}{x}_{}\stackrel{}{y}_{}`$ (2)
$`=`$ $`x_{}y_++x_+y_{}+\eta ^2x_+y_+\stackrel{}{x}_{}\stackrel{}{y}_{}.`$
Obviously, the light-cone is approached as the parameter $`\eta `$ goes to zero. For non-zero $`\eta `$, the transition to the coordinates introduced above can be formally identified as a Lorentz-boost combined with a linear transformation, which avoids time dependent boundary conditions . The boost parameter $`\beta =v_3`$ is given by
$$\beta =\frac{1\eta ^2/2}{1+\eta ^2/2},$$
(3)
indicating that for $`\eta ^20`$ the relative velocity $`v_3c(1)`$. Of course, this is connected to the well-known interpretation of the ’tilted’ light-cone frame in terms of the infinite momentum frame.
The choice of near light cone coordinates allows to quantize the theory on a space-like finite interval of length $`L`$ in $`x^{}`$ at equal times, i.e. $`\mathrm{\Delta }x^+=0`$. The invariant length squared of this interval for $`\mathrm{\Delta }x_{}^2=0`$ is related to the length of the compact $`x^{}`$ dimension.
$`\mathrm{\Delta }s^2`$ $`=`$ $`\mathrm{\Delta }x^{}\mathrm{\Delta }x^++\mathrm{\Delta }x^+\mathrm{\Delta }x^{}\eta ^2(\mathrm{\Delta }x^{})^2\mathrm{\Delta }x_{}^2`$ (4)
$`=`$ $`\eta ^2L^2.`$
For simplicity we consider also the transverse dimensions periodic in $`L`$. Previous work of the St. Petersburg and Erlangen groups assumed a fixed tilted coordinate system with fixed transverse ultraviolet cut off. Our main purpose in this paper is to consider the zero mode fields on a transverse lattice with varying transverse lattice spacing $`a`$. We propose to approach light-cone dynamics by synchronizing the continuum limit $`\mathrm{\Lambda }=\pi /a\mathrm{}`$ with the light-cone limit $`x^+\frac{1}{\sqrt{2}}(x^0+x^3)`$. The vacuum fluctuations in the transverse directions induce a second order phase transition which allows to have a continuum limit of the lattice theory. Thereby we can eliminate the cut off in a controlled way, preserving the nontrivial vacuum structure.
The gauge fixing procedure in the modified light-cone gauge $`_{}A_{}=0`$ involves zero modes dependent on the transverse coordinates. These zero mode fields carry zero linear momentum $`p_{}`$ in near light cone coordinates, but finite amount of $`p_0+p_3`$. They correspond to ”wee” partons in the language of the original parton model of Feynman. In $`SU(2)`$ the zero mode fields $`a_{}(x_{})`$ are proportional to $`\tau ^3`$, i.e. they can be chosen color diagonal. The use of an axial gauge is very natural for the light-cone Hamiltonian even more so than in the equal-time Hamiltonian. The asymmetry of the background zero mode naturally coincides with the asymmetry of the space coordinates on the light cone. The zero mode fields describe disorder fields. Depending on the effective coupling the zero mode transverse system will be in the massive or massless phase. Its second order phase transition allows to perform the continuum limit in the zero mode Hamiltonian. Our main conjecture is that the continuum limit and light cone limit can be realized simultaneously at this critical point. The evolution of the coupling determines the approach to the light cone with transverse resolution approaching zero. The resulting relation is reminiscent of the simple behavior one gets from considering naive scaling relations and dynamics in the infinite momentum frame. We will also calculate the contribution of the two dimensional zero modes to the ground state energy and show that it scales with the three dimensional volume. This estimate demonstrates very simply the relevance of modes with lower dimensionality to the full problem in the case of critical second order behavior.
The naive considerations for a simultaneous light cone and continuum limit go as follows. We characterize the ’infinite’ momentum frame by giving the momenta of the proton and photon in usual coordinates. The fast moving proton carries $`(P,0_{},P)`$ and the photon $`q=(\frac{\nu M}{2P},\sqrt{Q^2},\frac{\nu M}{2P})`$, where $`\nu `$ is the energy transfer in the laboratory. The dominant contribution comes from energy conserving transitions, where the energy of the quark in the final state equals the sum of initial quark and photon energies. For collinear quarks with momentum $`p=(x_BP,0_{},x_BP)`$, where $`x_B`$ denotes the momentum fraction, one explicitly finds
$$\sqrt{Q^2+\left(x_BP\frac{\nu M}{2P}\right)^2}=x_BP+\frac{\nu M}{2P},$$
(5)
which yields the “scaling variable” $`x_B=\frac{Q^2}{2M\nu }`$ as long as $`PQ/x_B`$, i.e. the “infinite” momentum has to be large enough:
$$\gamma M=PQ=1/a.$$
(6)
At the same time, the photon has a wavelength $`1/Q`$ which resolves transverse details of size $`a`$ in the hadronic wavefunction. By eq. (3), the $`\gamma `$ factor of the “infinite” momentum is related to $`\eta `$. Therefore we expect that the parameter $`\eta `$ approaches zero with the transverse cutoff $`\mathrm{\Lambda }=\pi /a\mathrm{}`$
$$\eta =\frac{1}{\sqrt{2}\gamma }\frac{1}{\sqrt{2}}aM0.$$
(7)
One of the objectives of this paper is to derive the precise relation between $`\eta `$ and $`a`$. We will simulate the zero mode dynamics on a $`(2+1)`$ dimensional lattice to find the fixed point of the effective coupling. The light-like limit is governed by an infrared fixed point (cf. ref. ). Most previous discussions of the zero mode problem have been on the classical equations of motion level, here we will work on the fully nonperturbative quantum level. Our work is not directly related to M-theory, but some of the consequences may be relevant for other studies with light like compactification.
## 2 Near Light Cone QCD Hamiltonian
In Ref. the near light cone Hamiltonian has been derived. A similar Hamiltonian has been obtained in equal time coordinates . Here we will sketch the derivation. We restrict ourselves to the color gauge group $`SU(2)`$ and dynamical gluons; only an external (fermionic) charge density $`\rho _m`$ is considered here.
Since the $`A_+^a`$ coordinates have no momenta conjugate to them, the Weyl gauge $`A_+^a=0`$ is the starting point for a canonical formulation. The canonical momenta of the dynamical fields $`A_{}^a,A_i^a`$ are given by
$`\mathrm{\Pi }_{}^a`$ $`=`$ $`{\displaystyle \frac{}{F_+^a}}=F_+^a,`$
$`\mathrm{\Pi }_i^a`$ $`=`$ $`{\displaystyle \frac{}{F_{+i}^a}}=F_i^a+\eta ^2F_{+i}^a.`$ (8)
¿From this, we get the Weyl gauge Hamiltonian density
$$_W=\frac{1}{2}\mathrm{\Pi }_{}^a\mathrm{\Pi }_{}^a+\frac{1}{2}F_{12}^aF_{12}^a+\frac{1}{2\eta ^2}\underset{i=1,2}{}\left(\mathrm{\Pi }_i^aF_i^a\right)^2.$$
(9)
The Hamiltonian has to be supplemented by the original Euler–Lagrange equation for $`A_+`$ as constraints on the physical states (Gauss’ Law constraints)
$`G^a(x_{},x^{})|\mathrm{\Phi }`$ $`=`$ $`\left(D_{}^{ab}\mathrm{\Pi }_{}^b+D_{}^{ab}\mathrm{\Pi }_{}^b+g\rho _m^a\right)|\mathrm{\Phi }`$ (10)
$`=`$ $`\left(D_{}^{ab}\mathrm{\Pi }_{}^b+G_{}^a\right)|\mathrm{\Phi }=0.`$
In order to obtain a Hamiltonian formulated in terms of unconstrained variables, one has to resolve the Gauss’ Law constraint. Via unitary gauge fixing transformations a solution of Gauss’ Law with respect to components of the chromo–electric field $`\mathrm{\Pi }_{}`$ can be accomplished. This gives a Hamiltonian independent of the conjugate gauge fields $`A_{}`$, i.e. the latter become cyclic variables. Classically this would correspond to the light front gauge $`A_{}=0`$. However, this choice is not legitimate if we want to consider the theory in a finite box. Instead, the (classical) Coulomb light front gauge $`_{}A_{}=0`$ is compatible with gauge invariance and periodic boundary conditions. The reason is that $`A_{}`$ carries information on the (gauge invariant) eigenvalues of the spatial Polyakov line matrix
$$\widehat{𝒫}(x_{})=P\mathrm{exp}\left[ig𝑑x^{}A_{}(x_{},x^{})\right],$$
(11)
which can be written in terms of a diagonal matrix with the zero mode field $`a_{}^3(x_{})\frac{\tau _3}{2}=a_{}(x_{})`$
$$\widehat{𝒫}(x_{})=V\mathrm{exp}\left[igLa_{}(x_{})\right]V^{}.$$
(12)
Obviously we have to keep these ‘zero modes’ $`a_{}(x_{})`$ as dynamical variables, while the other components of $`A_{}`$ are eliminated. The zero mode degrees of freedom are independent of $`x^{}`$ and, therefore, correspond to quantities with zero longitudinal momentum $`p_{}`$.
In order to eliminate the momentum $`\mathrm{\Pi }_{}`$, conjugate to $`A_{}`$, by means of Gauss’ Law, one needs to ‘invert’ the covariant derivative $`D_{}`$ which simplifies to $`d_{}=_{}ig[a_{},..]`$. On the space of physical states one can simply make the replacement<sup>2</sup><sup>2</sup>2The inversion of $`d_{}`$ can be explicitly performed in terms of its eigenfunctions, cf. (10).
$$\mathrm{\Pi }_{}(x_{},x_{})p_{}(x_{})\left(d_{}^1\right)G_{}(x_{},y^{}).$$
(13)
The operator $`p_{}(x_{})`$ is also diagonal and $`p_{}^3(x_{})`$ is the momentum conjugate to the zero mode $`a_{}^3(x_{})`$. It has eigenvalue zero with respect to $`d_{}`$, i.e. $`d_{}p_{}=0`$, and is therefore not constrained.
The appearance of the zero modes implies a residual Gauss’ Law, which arises from the $`x_{}`$ integration over eq. (10) for the $`a=3`$ component using periodic boundary conditions in the $`x_{}`$ direction. This constraint on two dimensional fields can be handled in full analogy to QED, since it concerns the diagonal part of color space. A further Coulomb gauge fixing in the $`SU(2)`$ 3–direction eliminates the color neutral, $`x^{}`$–independent, two–dimensional longitudinal gauge field in favor of a neutral chromo–electric field
$$e_{}(x_{})=g_{}𝑑y^{}𝑑y_{}d(x_{}y_{})\left\{f^{3ab}A_{}^a(y_{},y^{})\mathrm{\Pi }_{}^b(y_{},y^{})+\rho _m^3(y_{},y^{})\right\}\frac{\tau ^3}{2}.$$
(14)
Here we use the periodic Greens function of the two dimensional Laplace operator
$$d(z_{})=\frac{1}{L^2}\underset{\stackrel{}{n}\stackrel{}{0}}{}\frac{1}{p_n^2}e^{ip_nz_{}},p_n=\frac{2\pi }{L}\stackrel{}{n},$$
(15)
where $`\stackrel{}{n}=(n_1,n_2)`$ and $`n_1,n_2`$ are integers.
As a remnant of the local Gauss’ Law constraints, a global condition
$$Q^3|\mathrm{\Phi }^{}=𝑑y^{}𝑑y_{}\left\{f^{3ab}A_{}^a(y_{},y^{})\mathrm{\Pi }_{}^b(y_{},y^{})+\rho _m^3(y_{},y^{})\right\}|\mathrm{\Phi }^{}=0$$
(16)
emerges. The physical meaning of this equation is that the neutral component of the total color charge, including external matter as well as gluonic contributions, must vanish in the sector of physical states. It can be shown (see ) that there must be $`\stackrel{~}{Q}_{12}(x_{})=0`$ everywhere in transverse space in order to avoid an infinite Coulomb energy. The two conditions together suggest that physical states have to be color singlets.
The final Hamiltonian density in the physical sector explicitly reads (in terms of unconstrained $`A_{}`$ and $`\mathrm{\Pi }_{}`$ obtained after a shift is made by subtracting averages)
$``$ $`=`$ $`\text{tr}\left[_1A_2_2A_1ig[A_1,A_2]\right]^2+{\displaystyle \frac{1}{\eta ^2}}\text{tr}\left[\mathrm{\Pi }_{}\left(_{}A_{}ig[a_{},A_{}]\right)\right]^2`$ (17)
$`+`$ $`{\displaystyle \frac{1}{\eta ^2}}\text{tr}\left[{\displaystyle \frac{1}{L}}e_{}_{}a_{}\right]^2+{\displaystyle \frac{1}{2L^2}}p_{}^3(x_{})p_{}^3(x_{})`$
$`+`$ $`{\displaystyle \frac{1}{L^2}}{\displaystyle _0^L}𝑑z^{}{\displaystyle _0^L}𝑑y^{}{\displaystyle \underset{p,q,n}{\overset{^{}}{}}}{\displaystyle \frac{G_{qp}(x_{},z^{})G_{pq}(x_{},y^{})}{\left[\frac{2\pi n}{L}+g(a_q(x_{})a_p(x_{}))\right]^2}}e^{i2\pi n(z^{}y^{})/L},`$
where $`p`$ and $`q`$ are matrix labels for rows and columns, $`a_q=(a_{})_{qq}`$ and the prime indicates that the summation is restricted to $`n0`$ if $`p=q`$. The operator $`G_{}(x_{},x^{})`$ is defined as
$$G_{}=_{}\mathrm{\Pi }_{}+gf^{abc}\frac{\tau ^a}{2}A_{}^b\left(\mathrm{\Pi }_{}^c\frac{1}{L}e_{}^c\right)+g\rho _m.$$
(18)
The last two terms of the Weyl gauge Hamiltonian come from the original term $`\mathrm{\Pi }_{}^2`$ with squared electric field strengths in $`x^{}`$ direction. After elimination of $`\mathrm{\Pi }_{}`$ the zero mode part and the light cone Coulomb energy in the axial gauge remain. In the Coulomb term one sees the role of the zero mode fields as infrared regulators of the spatial momenta $`p_{}=2\pi n/L`$ which are quantized due to the compact interval $`L`$. Since the two dimensional theory has been Coulomb gauge fixed, the electric field $`\stackrel{}{e}_{}`$ replaces the canonical momentum of the longitudinal, neutral gauge field which has been eliminated.
We note that the terms containing $`\stackrel{}{e}_{}`$ and the momentum $`\stackrel{}{\mathrm{\Pi }}_{}`$, have the pre-factor $`1/\eta ^2`$. Physically this pre-factor signals the increase of transverse electric energies with the boost factor $`\gamma =(\sqrt{2}\eta )^1`$. The boost also couples transverse electric fields with transverse magnetic fields. In the light-cone limit the pre-factor diverges and the adjacent brackets become constraint equations. This reflects the corresponding reduction of the number of degrees of freedom if one goes exactly on the light cone. We do not follow this procedure, but keep a finite $`\eta `$ as a kind of Lagrange parameter.
A characteristic feature of an exact light-cone formulation is the triviality of the ground state. This may simplify explicit calculations, e.g. of the hadron spectrum. However, the light cone vacuum is definitely not trivial in the zero mode sector. In fact, already in ref we have shown strong and weak coupling solutions of the zero mode Hamiltonian. Massless modes influence the dynamics in the light cone limit, whereas massive modes decouple when $`\eta 0`$. In contrast to earlier work, we solve the zero mode Hamiltonian in this paper numerically showing the transition from the massive phase to the massless phase.
The zero mode degrees of freedom $`a_{}`$ couple to the transverse three dimensional gluon fields $`A_i`$ via the second magnetic term in $``$, the Coulomb term and directly via the electric field $`\stackrel{}{e}_{}`$. We remark that the transverse electric fields $`\mathrm{\Pi }_{}`$ and $`\stackrel{}{e}_{}`$ are dual to the magnetic fields $`_{}A_{}`$ and $`_{}a_{}`$. This duality is typical for the light cone and is absent in the equal time case. Since duality plays an important role in supersymmetric QCD its role in light cone theories should deserves to be investigated in greater detail.
In the following we neglect the couplings between the three and two dimensional fields and consider the pure zero mode Hamiltonian
$$h=d^2x\left[\frac{1}{2L}p_{}^3(\stackrel{}{x}_{})p_{}^3(\stackrel{}{x}_{})+\frac{L}{2\eta ^2}(_{}a_{}^3)^2\right].$$
(19)
The global constraint, eq. (16), does not contain $`a_{}`$ and $`p_{}`$ and, consequently, is irrelevant for the time being. Even at this level of severe approximations the zero mode Hamiltonian differs from the corresponding one in QED. The reason is the hermiticity defect of the canonical momentum $`p^{}`$. In the Schrödinger representation eq. (19)
$$h=d^2x\left[\frac{1}{2L}\frac{1}{J(a_{}^3(\stackrel{}{x}_{}))}\frac{}{a_{}^3(\stackrel{}{x}_{})}J\left(a_{}^3(\stackrel{}{x}_{})\right)\frac{}{a_{}^3(\stackrel{}{x}_{})}+\frac{L^2}{2\eta ^2}(_{}a_{}^3)^2\right],$$
(20)
contains the Jacobian $`J(a_{})`$ which equals the Haar measure of $`SU(2)`$
$$J\left(a_{}^3(\stackrel{}{x}_{})\right)=\mathrm{sin}^2\left(\frac{gL}{2}a_{}^3(\stackrel{}{x}_{})\right).$$
(21)
It stems from the gauge fixing procedure, effectively introducing curvilinear coordinates. It also appears in the functional integration volume element for calculating matrix elements. It is convenient to introduce dimensionless variables
$$\phi (\stackrel{}{x}_{})=\frac{gLa_{}^3(\stackrel{}{x}_{})}{2},$$
(22)
which vary in a compact domain $`0\phi \pi `$. We regularize the above Hamiltonian $`h`$ by introducing a lattice spacing $`a`$ between transversal lattice points $`\stackrel{}{b}`$. Next we appeal to the physics of the infinite momentum frame and factorize the reduced true energy from the Lorentz boost factor $`\gamma =\sqrt{2}/\eta `$ and the cut off by defining $`h_{\mathrm{red}}`$
$$h=\frac{1}{2\eta a}h_{\mathrm{red}}$$
(23)
In the continuum limit of the transverse lattice theory we let $`a`$ go to zero. For small lattice spacing we obtain the reduced Hamiltonian
$$h_{\mathrm{red}}=\underset{\stackrel{}{b}}{}\left\{g_{\mathrm{eff}}^2\frac{1}{J}\frac{}{\phi (\stackrel{}{b})}J\frac{}{\phi (\stackrel{}{b})}+\frac{1}{g_{\mathrm{eff}}^2}\underset{i=1,2}{}\left(\phi (\stackrel{}{b})\phi (\stackrel{}{b}+\stackrel{}{e}_i)\right)^2\right\}.$$
(24)
with the effective coupling constant
$$g_{\mathrm{eff}}^2=\frac{g^2L\eta }{4a}$$
(25)
The first part of the Hamiltonian contains the kinetic (electric) energy of the $`SU(2)`$ rotators on a half circle at each lattice point and the second part gives the potential (magnetic) energy of these rotators due to the differences of angles at nearest neighbor sites. For further discussion we define these electric and magnetic parts as:
$`h_{\mathrm{red}}`$ $`=`$ $`h_\mathrm{e}+h_\mathrm{m}`$ (26)
$`h_\mathrm{e}`$ $`=`$ $`{\displaystyle \underset{\stackrel{}{b}}{}}h_\stackrel{}{b}^\mathrm{e}`$
$`h_\mathrm{m}`$ $`=`$ $`{\displaystyle \underset{\stackrel{}{b}}{}}{\displaystyle \underset{i=1,2}{}}h_{\stackrel{}{b},\stackrel{}{b}+\stackrel{}{e}_i}^\mathrm{m}`$
In the strong coupling domain where we have analytical solutions of the zero mode Hamiltonian, the numerical Hamiltonian lattice theory agrees with the analytical solutions for the mass gap. The real question concerns the continuum limit of the zero mode system which occurs outside of the strong coupling region. For this purpose we have to find the region of vanishing mass gap for the lattice Hamiltonian.
## 3 Lattice Calculation of the Zero Mode Hamiltonian
We solve the equivalent lattice theory in a Euclidean formulation. The zero mode Hamiltonian represents a $`2+1`$ dimensional theory in two spatial and one time direction. The lattice has $`N_xa=N_ya`$ extensions in transverse space and $`N_T\mathrm{\Delta }\tau =T`$ extension in near light cone time. To set up the density matrix one has to write down the Trotter formula for the given Hamiltonian. Using $`h_{\mathrm{red}}=h_\mathrm{e}+h_\mathrm{m}`$ we have
$$\mathrm{exp}(Th_{\mathrm{red}})=\underset{N_T\mathrm{}}{lim}\left[\mathrm{exp}\left(\mathrm{\Delta }\tau h_\mathrm{m}/2\right)\mathrm{exp}\left(\mathrm{\Delta }\tau h_\mathrm{e}\right)\mathrm{exp}\left(\mathrm{\Delta }\tau h_\mathrm{m}/2\right)\right]^{N_T},$$
(27)
where each time evolution step $`\mathrm{\Delta }\tau `$ can be separately done for the electric and magnetic part of the Hamiltonian. For definiteness, we choose $`\mathrm{\Delta }\tau =a/2`$ in all the following. In the Appendix we show that this choice of $`\mathrm{\Delta }\tau `$ optimizes the updating procedure since it generates approximately equal widths for the weight functions resulting from the kinetic and potential energies. The different time slices will labeled with the index $`l`$. The electric Hamiltonian can be evaluated by inserting products of complete set single site eigenfunctions $`C_{n_l}`$ (cf. ). Practically a maximal number $`N_{\mathrm{max}}=100`$ of eigenfunctions is fully sufficient to reach convergence in the interval of couplings we need.
$$\phi _{l+1}|h^\mathrm{e}|\phi _l=\underset{n_{l+1},n_l}{}\phi _{l+1}|n_{l+1}n_{l+1}|h^\mathrm{e}|n_ln_l|\phi _l\underset{n_l=0}{\overset{N_{\mathrm{max}}}{}}C_{n_l}(\phi _{l+1})g_{\mathrm{eff}}^2n_l(n_l+2)C_{n_l}(\phi _l)$$
(28)
with the (single site) eigenfunctions and eigenvalues given as:
$`C_{n_l}(\phi _l)`$ $`=`$ $`\sqrt{{\displaystyle \frac{2}{\pi }}}\left\{{\displaystyle \frac{\mathrm{sin}\left((n_l+1)\phi _l\right)}{\mathrm{sin}\phi _l}}\right\},`$ (29)
$`h^\mathrm{e}C_{n_l}(\phi _l)`$ $`=`$ $`g_{\mathrm{eff}}^2n_l(n_l+2)C_{n_l}(\phi _l).`$ (30)
The eigenfunctions form an orthonormal set with respect to a scalar product which contains the Jacobian in the measure. The Jacobian $`J\left(\phi (\stackrel{}{b})\right)`$ has been defined above (21). The magnetic part is diagonal in $`\{\phi (\stackrel{}{b})\}`$ i.e. local in time:
$$\{\phi _l(\stackrel{}{b})\}|h_\mathrm{m}|\{\phi _l(\stackrel{}{b})\}=h_\mathrm{m}(\{\phi _l(\stackrel{}{b})\})=\underset{\stackrel{}{b}}{}\underset{i=1,2}{}h_{\stackrel{}{b},\stackrel{}{b}+\stackrel{}{e}_i}^\mathrm{m}\left(\phi _l(\stackrel{}{b})\phi _l(\stackrel{}{b}+\stackrel{}{e}_i)\right).$$
(31)
The full partition function is given by an integral over all time slices
$`Z`$ $`=`$ $`\text{tr}\mathrm{exp}\left[Th_{\mathrm{red}}\right]={\displaystyle \underset{\stackrel{}{b}}{}\left(J(\phi (\stackrel{}{b}))d\phi (\stackrel{}{b})\right)\{\phi (\stackrel{}{b})\}|\mathrm{exp}\left[Th_{\mathrm{red}}\right]|\{\phi (\stackrel{}{b})\}}`$
$`=`$ $`{\displaystyle \underset{\stackrel{}{b},l}{}\left(J(\phi (\stackrel{}{b}))d\phi (\stackrel{}{b})\underset{n_l=0}{\overset{N_{\mathrm{max}}}{}}C_{n_l}(\phi _{l+1}(\stackrel{}{b}))\mathrm{exp}\left[g_{\mathrm{eff}}^2n_l(n_l+2)\mathrm{\Delta }\tau \right]C_{n_l}(\phi _l(\stackrel{}{b}))\right)}`$
$`\times `$ $`{\displaystyle \underset{l}{}}\mathrm{exp}\left[h_\mathrm{m}(\{\phi _l(\stackrel{}{b})\})\mathrm{\Delta }\tau \right].`$
Because of the Jacobian the dominant contributions to the partition function come from $`\phi _l`$ values around $`\frac{\pi }{2}`$. The Hamiltonian is invariant under reflections $`\phi _l\frac{\pi }{2}\frac{\pi }{2}\phi _l`$. In the strong coupling limit, where $`g_{\mathrm{eff}}^2`$ is large, the half rotators act almost independently on each lattice site producing a large mass gap. The system is disordered. For decreasing coupling constant $`g_{\mathrm{eff}}^2`$ the movement of the individual rotators becomes locked from one site to the next. Long range correlations develop. The order parameter (or “magnetization”) of the system is the expectation value of the trace of the Polyakov line (11), which on the lattice has the form
$$P=\overline{\frac{1}{2}tr\widehat{𝒫}}=\frac{1}{N^3}\underset{\stackrel{}{b},\tau }{}\mathrm{cos}\phi (\stackrel{}{b},\tau ).$$
(33)
The order parameter is odd under the above symmetry operation. Operator expectation values are evaluated with the density matrix defined by $`h_{\mathrm{red}}`$ and $`T`$:
$$O=\frac{1}{Z}\text{tr}O\mathrm{exp}[Th_{\mathrm{red}}].$$
(34)
At each particular $`\beta _g`$
$$\beta _g1/g_{\mathrm{eff}}^2$$
(35)
we simulate lattice sizes with equal number of sites in all directions $`N=N_x=N_y=N_T`$ and $`N=4,6,8,12,16`$. A Metropolis algorithm is used for updating with the steps size and the number of hits adapted to $`\beta _g`$. In order to tabulate the Boltzmann weights related to the timelike links we also discretized the continuous angle variables to a system of $`N_\phi =60`$ orientations: $`\phi _j=j\pi /N_\phi `$, $`j=[0,N_\phi ]`$.
In this explorative investigation we generated between $`5000`$ and $`50000`$ uncorrelated configurations depending on the $`\beta _g`$ and lattice size. All our calculations were done on a cluster of AlphaStations with the single site Metropolis updating algorithm sketched above, so the accuracy can be improved using a more powerful algorithms and/or more computing resources.
We calculate the $`\beta _g`$ dependence of the following quantities: average electric energy $`\epsilon _\mathrm{e}`$ and magnetic energy $`\epsilon _\mathrm{m}`$ per site, the average of the absolute value of the order parameter $`|P|`$, the susceptibility $`\chi `$ and the normalized fourth cumulant $`g_r`$ (which gives the deviation of the moments of the Polyakov expectation value from a pure Gaussian behavior):
$`\epsilon _{\mathrm{e},\mathrm{m}}`$ $`=`$ $`{\displaystyle \frac{1}{N^2}}h_{\mathrm{e},\mathrm{m}}`$ (36)
$`\chi `$ $`=`$ $`N^3\left(P^2P^2\right)`$ (37)
$`g_r`$ $`=`$ $`{\displaystyle \frac{P^4}{P^2^2}}3.`$ (38)
Firstly we calculated the ground state energy for strong couplings $`g_{\mathrm{eff}}^2`$ and compared with the exact calculation . The obtained results from the lattice agree with the analytical result. For higher values of $`\beta _g`$ the lattice calculation agrees within $`1020\%`$ with the previous effective double site calculation for the energy per site. The ultraviolet regularization with $`\mathrm{\Delta }\tau `$ influences the final lattice result. The energies are measured with high accuracy and practically do not depend on the infrared cutoff, i.e. on the total lattice size for $`N8`$. The situation for other variables is different.
In order to see the emergence of a massless phase, we investigated the (connected) time correlation functions of the (trace of the) Polyakov line operators:
$`K(0,\tau )={\displaystyle \frac{1}{N^4}}{\displaystyle \underset{\stackrel{}{b},\stackrel{}{b^{}}}{}}\mathrm{cos}(\phi (\stackrel{}{b},0))\mathrm{cos}(\phi (\stackrel{}{b}^{},\tau ))|P|^2.`$ (39)
The correlation masses are obtained from a fit of the time correlation functions $`K(0,\tau )|_{\tau =n\mathrm{\Delta }\tau }`$ to a parameterization taking the periodicity in time into account
$$K(0,n\mathrm{\Delta }\tau )=c\text{cosh}\left[m\mathrm{\Delta }\tau \left(n\frac{N_T}{2}\right)\right].$$
(40)
The obtained masses for different lattice sizes are shown in Fig.1a. For each fixed $`\beta _g`$ we extrapolated the mass to the limit $`N\mathrm{}`$ using a linear $`1/N`$ parameterization
$$m(\beta _g,N)=m_0(\beta _g)+\frac{C(\beta _g)}{N}.$$
(41)
Our data for $`m(\beta _g,N)`$ can be fitted with this linear form (41) for $`\beta _g>3`$. The resulting dependence of $`m_0`$ on $`\beta _g`$ is presented in Fig.1b. The mass vanishes in the infinite volume limit near $`\beta 5`$, which is a first guess for the critical coupling $`\beta _g^{}`$.
In order to extract more exact information on the critical behavior
$$m(\beta _g)\left(\frac{1}{\beta _g}\frac{1}{\beta _g^{}}\right)^\nu .$$
(42)
from our data on small systems we have done a Finite Size Scaling (FSS) analysis of the variables $`|P|`$, $`\chi `$ and $`g_r`$, searching for the critical coupling $`\beta _g^{}`$ in the interval $`3\beta _g9`$. We could determine the critical indices $`\beta `$, $`\gamma `$ and $`\nu `$ using an approach employed for $`SU(2)`$ gauge theory at the finite temperature transition by Engels et al. . The general form of the scaling relations for a variable $`V`$ ($`V`$ = $`|P|`$, $`\chi `$, $`g_r`$) is
$$V(t,N)=N^{\rho /\nu }F_V(tN^{1/\nu },g_iN^{y_i}),$$
(43)
where $`\rho `$ is the corresponding critical index ($`\beta `$, $`\gamma `$, 0) for the respective quantities and $`t`$ is the reduced inverse coupling $`\beta _g`$:
$$t=\frac{\beta _g^{}\beta _g}{\beta _g}.$$
(44)
In practice eq. (43) is computed near $`t=0`$. Expanding up to first order in $`t`$ and taking into account only the largest irrelevant exponent $`y_1\omega `$ one obtains
$$V(t,N)=N^{\rho /\nu }\left[c_0+\left(c_1+c_2N^\omega \right)tN^{1/\nu }+c_3N^\omega \right].$$
(45)
With this parameterization we analyze our data for lattice sizes $`N=4,6,8,12,16`$. The $`\beta _g`$ dependence in eq. (45) is valid near the critical point. The quality of the fit to the data in different intervals on $`\beta _g`$ can serve as a guide to localize the critical point $`\beta _g^{}`$ . We find that in our case this parameterization can be applied for $`|P|`$ and $`\chi `$ only for $`\beta _g>5`$: for smaller $`\beta _g`$ the $`\chi ^2/D.F.`$ becomes too large. Finally we use the parameterization of eq. (45) in the interval $`5\beta _g9`$ (see Fig.2) where $`|t|<0.3`$ and an expansion of first order in $`t`$ still is applicable. In this $`\beta _g`$ interval the fits for $`|P|`$ and $`\chi `$ yield a $`\chi ^2/D.F2`$. To have a reliable error estimate for the extracted parameters we use the jackknife method. The results for $`\beta _c^{}`$ and critical indices $`\beta `$, $`\gamma `$ and $`\nu `$ are presented in Table 1. We use $`\omega =1`$, but in fact the results do not depend on $`\omega `$ in the interval $`0.8\omega 1.2`$.
It should be noticed that the system frequently flips between the two ordered states on a finite lattice. Therefore the expectation value $`P`$ in equation (37) vanishes with good accuracy. Analogously to the treatment of the magnetization in the 3-dimensional Ising model, we should - instead of the expression for $`\chi `$ in eq. (37) - define the susceptibility in the broken phase as $`\chi _{\mathrm{broken}}=N^3\left(P^2|P|^2\right)`$. For the Ising case this susceptibility is supposed to converge towards the correct infinite volume limit. Due to lack of statistics, however, we did not use separate expressions for $`\chi `$ in the two phases. Our data for $`g_r`$ are even much less accurate which excludes the possibility to use them in the FSS analysis.
The values found for the critical indices $`\beta `$, $`\gamma `$ and $`\nu `$ are not far from those of the Ising model ($`D=3`$) (see Table 1). They are much less accurate but also in agreement with a high statistics analysis of the $`SU(2)`$ deconfinement transition at finite temperature. We interprete the behavior of the light cone zero mode theory as a consequence of the underlying $`Z(2)`$ symmetry of the Hamiltonian. $`Z(2)`$ transformations correspond to reflections of the angle variables $`\phi `$ around $`\frac{\pi }{2}`$. The Hamiltonian and the measure of eq. (17) are invariant under these transformations. The resulting critical behavior is common to Ising like models. The critical coupling itself is a nonuniversal quantity and we have found a rough value. More extended work with the lattice Hamiltonian near the light cone is needed to clarify the continuous transition further.
## 4 Conclusions
The scaling analysis gives indications that there is a second order transition between a phase with massive excitations at strong coupling and a phase with massless excitations in weak coupling. To reach a higher accuracy large scale simulations of the Hamiltonian zero mode system are needed. In this context it is advisable to treat the coupled system of three and two dimensional modes together. A calculation in the epsilon expansion gives the zero of the $`\beta `$-function as an infrared stable fixed point. This shows that the limit of large longitudinal dimensions $`L`$ is well defined. Using the running coupling constant $`g_{\mathrm{eff}}^2`$ of the zero mode system we have (cf. (42))
$$ma=\frac{1}{\zeta _0g_{\mathrm{eff}}^{\mathrm{\hspace{0.17em}2}}}\left(g_{\mathrm{eff}}^2g_{\mathrm{eff}}^{\mathrm{\hspace{0.17em}2}}\right)^\nu $$
(46)
with $`g_{\mathrm{eff}}^{\mathrm{\hspace{0.17em}2}}=0.17\pm 0.03`$ and $`\nu =0.56\pm 0.05`$ from our lattice calculations (cf. Table 1). The mass can be interpreted as an inverse correlation length $`ma=a/\xi `$. When the correlation length $`\xi `$ approaches the temporal size of the lattice $`\xi N_Ta/2`$ in the infrared limit we can identify $`ma`$ with $`a/L`$, where $`L`$ is our infrared length scale. Therefore the effective coupling $`g_{\mathrm{eff}}^2=g^2\eta \frac{L}{4a}`$ (cf. (25)) runs with $`a/L`$ in the following way
$$g_{\mathrm{eff}}^2=g_{\mathrm{eff}}^{\mathrm{\hspace{0.17em}2}}+\left(\zeta _0g_{\mathrm{eff}}^{\mathrm{\hspace{0.17em}2}}\right)^{1/\nu }\left(\frac{a}{L}\right)^{1/\nu }.$$
(47)
The coupling to the three dimensional modes produces the usual evolution of $`g^2`$ in $`SU(2)`$ QCD, where the coupling $`g^2(a)`$ at the lattice scale $`a`$ is related to the coupling $`g_0^2`$ defined at the infrared scale $`L`$ as follows
$$g^2(a)=\frac{g_0^2}{1+\frac{g_0^2}{4\pi ^2}\frac{22}{3}\mathrm{log}\frac{L}{a}}.$$
(48)
Combining the eqs. (47) and (48) we can synchronize the approach to the light cone, i.e. the limit $`\eta 0`$ with the continuum limit $`a0`$. The condition that the three-dimensional evolution of $`g^2`$ has to be compatible with the two-dimensional evolution of $`g_{\mathrm{eff}}^2`$ towards $`g_{\mathrm{eff}}^{\mathrm{\hspace{0.17em}2}}`$ yields that for $`a0`$ the light cone parameter $`\eta `$ approaches zero as
$$\eta (a)\frac{g_{\mathrm{eff}}^{\mathrm{\hspace{0.17em}2}}}{\pi ^2}\frac{22}{3}\frac{a}{L}\mathrm{log}\frac{L}{a}.$$
(49)
This relation is similar to the naive scaling result deduced in section $`1`$ from physics arguments besides logarithmic modification. Although at the outset we had as parameters $`N_L=L/a`$ and $`\eta `$ in the Hamiltonian we reduce this multiparameter problem to a problem with one single coupling constant which can be chosen as $`g`$.
We can use the energy and the critical behavior of the Euclidean $`2+1`$ dimensional system to obtain its contribution to the ground state energy of the system in $`3+1`$ dimensions. In our calculations the energy of the zero mode Hamiltonian $`h`$ near the critical coupling is proportional to $`N_{}^2/\eta `$. Using the scaling of $`\eta `$ determined from the compatibility of 2-dimensional and 3-dimensional dynamics, we can eliminate $`\eta `$ in favor of $`N_L=L/a`$ and get a zero mode energy which grows linearly with the 3-dimensional volume. That means the zero mode dynamics becomes relevant for the full problem, i.e.
$$haN_LN_{}^2.$$
(50)
In reality the zero mode system is coupled to the $`3+1`$ dimensional system the extension of which is large. So we expect that the couplings of the $`2+1`$ dimensional system can get modified. Because of the universal dynamics near the infrared fixed point this will not change the qualitative form of the volume dependence obtained in eq. (50). The behavior of the zero mode Hamiltonian has to be taken into account when we try to solve for the infinite momentum frame solutions of the complete Hamiltonian. The discussion until now often centers on the order of limits. In order to guarantee a meaningful limit of the near light cone theory, one first has to take $`L\mathrm{}`$ and then one can take $`\eta 0`$. This is a workable procedure in solvable $`1+1`$ dimensional models. In QCD in $`3+1`$ dimensions it hardly can serve as a prescription. Numerical solutions have to be obtained with a finite cutoff and the continuum limit can only be reached approximately via scaling relations. It is in this case that the role of zero mass excitations will crucially enter the physics. As one sees above the synchronization of the continuum limit with the light cone limit leads to results which are independent of the near light cone parameter $`\eta `$. The independence of the ground state energy near the critical coupling gives us back Lorentz invariance which states that the physics should not depend on the reference frame. The vacuum energy is not allowed to depend on $`\eta `$. There are dynamical zero modes in near light cone QCD which contain the physics of the nonperturbative QCD vacuum, the physics of the gluon and quark condensates. The nontrivial structure of the nonabelian gauge theory enters in our calculation through the kinetic operator of the rotators which at finite resolution have a finite mass gap vanishing only in the continuum limit.
The near light cone description of QCD based on a modified axial gauge $`_{}A_{}=0`$ yields a very natural formulation of high energy scattering. Polyakov variables near the light cone have already entered the calculations of high energy cross sections in two different ways: For the geometrical size of cross sections in ref. the correlations between Polyakov lines along the projectile and target directions are important. For this situation an interpolating gauge may be useful. For the energy dependence of cross sections an evolution equation of Polyakov line correlators approaching the light cone has been discussed in ref. which generalizes the BFKL- equation. This situation is close to the treatment in this paper. Now the decisive step is to back up our theoretical work on effective transverse QCD dynamics near the light cone with phenomenological consequences for high energy scattering.
## Appendix
A choice of $`\mathrm{\Delta }\tau `$ is proposed by comparing Monte-Carlo weights $`W_t`$ and $`W_s`$ which link nearest neighbors in time and space respectively
$`W_t(\phi _1,\phi _2)`$ $`=`$ $`\phi _1(x,y,\tau +1)|e^{\mathrm{\Delta }\tau h^\mathrm{e}}|\phi _2(x,y,\tau )\sqrt{J(\phi _1)J(\phi _2)},`$ (51)
$`W_s(\phi _1,\phi _2)`$ $`=`$ $`\phi _1(x+1,y,\tau )|e^{\mathrm{\Delta }\tau h^\mathrm{m}}|\phi _2(x,y,\tau ).`$ (52)
The probability distribution for the update of the field $`\phi (x,y,\tau )`$ is given by the product of weights for all links connecting it to neighbors in time and space. So the updating procedure for $`\phi `$ is more effective when $`W_t`$ and $`W_s`$ have similar distributions near the angle $`\phi =\pi /2`$ where the Jacobian has its maximum. For $`\mathrm{\Delta }\tau =a/10`$ (cf. Fig. 3a) the spatial weights $`W_s`$ have a much wider distribution than the timelike weights $`W_t`$, whereas for $`\mathrm{\Delta }\tau =a`$ (cf. Fig. 3c) the situation is reverse. For the intermediate case $`\mathrm{\Delta }\tau =a/2`$ (cf. Fig. 3b) the widths of the distributions are very similar, therefore the updating procedure is optimal.
|
warning/0003/gr-qc0003100.html
|
ar5iv
|
text
|
# Gravitational Energy-Momentum Density in Teleparallel Gravity
\[
## Abstract
In the context of a gauge theory for the translation group, a conserved energy-momentum gauge current for the gravitational field is obtained. It is a true spacetime and gauge tensor, and transforms covariantly under global Lorentz transformations. By rewriting the gauge gravitational field equation in a purely spacetime form, it becomes the teleparallel equivalent of Einstein’s equation, and the gauge current reduces to the Møller’s canonical energy-momentum density of the gravitational field.
\]
The definition of an energy-momentum density for the gravitational field is one of the oldest and most controversial problems of gravitation. As a true field, it would be natural to expect that gravity should have its own local energy-momentum density. However, it is usually asserted that such a density can not be locally defined because of the equivalence principle . As a consequence, any attempt to identify an energy-momentum density for the gravitational field leads to complexes that are not true tensors. The first of such attempt was made by Einstein who proposed an expression for the energy-momentum density of the gravitational field which was nothing but the canonical expression obtained from Noether’s theorem . Indeed, this quantity is a pseudotensor, an object that depends on the coordinate system. Several other attempts have been made, leading to different expressions for the energy-momentum pseudotensor for the gravitational field .
Despite the existence of some controversial points related to the formulation of the equivalence principle , it seems true that, in the context of general relativity, no tensorial expression for the gravitational energy-momentum density can exist. However, as our results show, in the gauge context, the existence of an expression for the gravitational energy-momentum density which is a true spacetime and gauge tensor turns out to be possible. Accordingly, the absence of such expression should be attributed to the general relativity description of gravitation, which seems to be not the appropriate framework to deal with this problem .
In spite of some skepticism , there has been a continuous interest in this problem . In particular, a quasilocal approach has been proposed recently which is highly clarifying . According to this approach, for each gravitational energy-momentum pseudotensor, there is an associated superpotential which is a hamiltonian boundary term. The energy-momentum defined by such a pseudotensor does not really depend on the local value of the reference frame, but only on the value of the reference frame on the boundary of a region — then its quasilocal character. As the relevant boundary conditions are physically acceptable, this approach validates the pseudotensor approach to the gravitational energy-momentum problem. It should be mentioned that these results were obtained in the context of the general relativity description of gravitation.
In the present work a different approach will be used to re-examine the gravitational energy-momentum problem. Due to the fundamental character of the geometric structure underlying gauge theories, the concept of currents, and in particular the concepts of energy and momentum, are much more transparent when considered from the gauge point of view . Accordingly, we are going to consider gravity as described by a gauge theory . Our basic interest will be concentrated on the gauge theories for the translation group , and in particular on the so called teleparallel equivalent of general relativity . It is important to remark that this equivalence is true only in the absence of spinor matter fields .
Let us start by reviewing the fundamentals of the teleparallel equivalent of general relativity. We use the Greek alphabet $`(\mu ,\nu ,\rho ,\mathrm{}=0,1,2,3)`$ to denote indices related to spacetime, and the Latin alphabet $`(a,b,c,\mathrm{}=0,1,2,3)`$ to denote indices related to the tangent space (fiber), assumed to be a Minkowski space with the metric $`\eta _{ab}=\text{diag}(+1,1,1,1)`$. A gauge transformation is defined as a local translation of the tangent-space coordinates,
$$\delta x^a=\delta \alpha ^bP_bx^a,$$
(1)
with $`P_a=/x^a`$ the translation generators, and $`\delta \alpha ^a`$ the corresponding infinitesimal parameters. Denoting the gauge potentials by $`A^a_\mu `$, the gauge covariant derivative of a general matter field $`\mathrm{\Psi }`$ is
$$𝒟_\mu \mathrm{\Psi }=h^a{}_{\mu }{}^{}_{a}^{}\mathrm{\Psi },$$
(2)
where
$$h^a{}_{\mu }{}^{}=_\mu x^a+c^2A^a_\mu $$
(3)
is a nontrivial tetrad field, with $`c`$ the speed of light. From the covariance of $`D_\mu \mathrm{\Psi }`$, we obtain the transformation of the gauge potentials:
$$A^a^{}{}_{\mu }{}^{}=A^a{}_{\mu }{}^{}c^2_\mu \delta \alpha ^a.$$
(4)
As usual in abelian gauge theories, the field strength is given by
$$F^a{}_{\mu \nu }{}^{}=_\mu A^a{}_{\nu }{}^{}_\nu A^a{}_{\mu }{}^{},$$
(5)
which satisfies the relation
$$[𝒟_\mu ,𝒟_\nu ]\mathrm{\Psi }=c^2F^a{}_{\mu \nu }{}^{}P_{a}^{}\mathrm{\Psi }.$$
(6)
It is important to remark that, whereas the tangent space indices are raised and lowered with the metric $`\eta _{ab}`$, the spacetime indices are raised and lowered with the riemannian metric
$$g_{\mu \nu }=\eta _{ab}h^a{}_{\mu }{}^{}h_{}^{b}{}_{\nu }{}^{}.$$
(7)
A nontrivial tetrad field induces on spacetime a teleparallel structure which is directly related to the presence of the gravitational field. In fact, given a nontrivial tetrad $`h^a_\mu `$, it is possible to define a Cartan connection
$$\mathrm{\Gamma }^\rho {}_{\mu \nu }{}^{}=h_a{}_{}{}^{\rho }_{\nu }^{}h^a{}_{\mu }{}^{},$$
(8)
which is a connection presenting torsion, but no curvature . As a natural consequence of this definition, the Cartan covariant derivative of the tetrad field vanishes identically:
$$_\nu h^a{}_{\mu }{}^{}_\nu h^a{}_{\mu }{}^{}\mathrm{\Gamma }^\theta {}_{\mu \nu }{}^{}h_{}^{a}{}_{\theta }{}^{}=0.$$
(9)
This is the absolute parallelism condition. The torsion of the Cartan connection is
$$T^\rho {}_{\mu \nu }{}^{}=\mathrm{\Gamma }^\rho {}_{\nu \mu }{}^{}\mathrm{\Gamma }^\rho {}_{\mu \nu }{}^{},$$
(10)
from which we see that the gravitational field strength is nothing but torsion written in the tetrad basis:
$$F^a{}_{\mu \nu }{}^{}\text{ = }c^2h^a{}_{\rho }{}^{}T_{}^{\rho }{}_{\mu \nu }{}^{}.$$
(11)
The Cartan connection $`\mathrm{\Gamma }^\rho _{\mu \nu }`$ and the Levi-Civita connection of the metric (7), denoted by $`\stackrel{}{\mathrm{\Gamma }}{}_{}{}^{\rho }_{\mu \nu }`$, are related by
$$\mathrm{\Gamma }^\rho {}_{\mu \nu }{}^{}=\stackrel{}{\mathrm{\Gamma }}{}_{}{}^{\rho }{}_{\mu \nu }{}^{}+K^\rho {}_{\mu \nu }{}^{},$$
(12)
with
$$K^\rho {}_{\mu \nu }{}^{}=\frac{1}{2}(T_\mu {}_{}{}^{\rho }{}_{\nu }{}^{}+T_\nu {}_{}{}^{\rho }{}_{\mu }{}^{}T^\rho {}_{\mu \nu }{}^{})$$
(13)
the contorsion tensor.
The gauge gravitational field Lagrangian is given by
$$_G=\frac{hc^4}{16\pi G}S^{\rho \mu \nu }T_{\rho \mu \nu },$$
(14)
where $`h=\mathrm{det}(h^a{}_{\mu }{}^{})`$, and
$`S^{\rho \mu \nu }=S^{\rho \nu \mu }\frac{1}{2}[K^{\mu \nu \rho }g^{\rho \nu }T^{\theta \mu }{}_{\theta }{}^{}+g^{\rho \mu }T^{\theta \nu }{}_{\theta }{}^{}]`$
is a tensor written in terms of the Cartan connection only. As usual in gauge theories, it is quadratic in the field strength. By using relation (12), this lagrangian can be rewritten in terms of the Levi-Civita connection. Up to a total divergence, the result is the Hilbert-Einstein Lagrangian of general relativity
$$=\frac{c^4}{16\pi G}\sqrt{g}\stackrel{}{R},$$
(15)
where the identification $`h=\sqrt{g}`$ has been made.
By performing variations in relation to the gauge field $`A_a^\rho `$, we obtain from the gauge lagrangian $`_G`$ the teleparallel version of the gravitational field equation,
$$_\sigma (hS_a{}_{}{}^{\sigma \rho })\frac{4\pi G}{c^4}(hj_a{}_{}{}^{\rho })=0,$$
(16)
where $`S_a{}_{}{}^{\sigma \rho }h_a{}_{}{}^{\lambda }S_{\lambda }^{}^{\sigma \rho }`$. Analogously to the Yang-Mills theories ,
$$hj_a{}_{}{}^{\rho }\frac{_G}{h^a_\rho }=\frac{c^4}{4\pi G}hh_a{}_{}{}^{\lambda }S_{\mu }^{}{}_{}{}^{\nu \rho }T_{}^{\mu }{}_{\nu \lambda }{}^{}h_a{}_{}{}^{\rho }_{G}^{}$$
(17)
stands for the gravitational gauge current, which in this case represents the energy and momentum of the gravitational field. The term $`(hS_a{}_{}{}^{\sigma \rho })`$ is called superpotential in the sense that its derivative yields the gauge current $`(hj_a{}_{}{}^{\rho })`$. Because of the anti-symmetry of $`S_a^{\sigma \rho }`$ in the last two indices, $`(hj_a{}_{}{}^{\rho })`$ is conserved as a consequence of the field equation:
$$_\rho (hj_a{}_{}{}^{\rho })=0.$$
(18)
Making use of the identity
$$_\rho hh\mathrm{\Gamma }^\nu {}_{\nu \rho }{}^{}=h(\mathrm{\Gamma }^\nu {}_{\rho \nu }{}^{}K^\nu {}_{\rho \nu }{}^{}),$$
(19)
this conservation law can be rewritten as
$$D_\rho j_a{}_{}{}^{\rho }_\rho j_a{}_{}{}^{\rho }+(\mathrm{\Gamma }^\rho {}_{\lambda \rho }{}^{}K^\rho {}_{\lambda \rho }{}^{})j_a{}_{}{}^{\lambda }=0,$$
(20)
where $`D_\rho `$ is the teleparallel version of the covariant derivative, which is nothing but the Levi-Civita covariant derivative of general relativity rephrased in terms of the Cartan connection . As can be easily checked, $`j_a^\rho `$ transforms covariantly under a general spacetime coordinate transformation, and is invariant under local (gauge) translation of the tangent-space coordinates. This means that $`j_a^\rho `$ is a true spacetime and gauge tensor. However, it transforms covariantly only under a global tangent-space Lorentz transformation.
Let us now proceed further and find out the relation between the above gauge approach and general relativity. By using Eq. (8) to express $`_\rho h_a^\lambda `$, the field equation (16) can be rewritten in a purely spacetime form,
$$_\sigma (hS_\lambda {}_{}{}^{\sigma \rho })\frac{4\pi G}{c^4}(ht_\lambda {}_{}{}^{\rho })=0,$$
(21)
where now
$$ht_\lambda {}_{}{}^{\rho }=\frac{c^4}{4\pi G}h\mathrm{\Gamma }^\mu {}_{\nu \lambda }{}^{}S_{\mu }^{}{}_{}{}^{\nu \rho }+\delta _\lambda {}_{}{}^{\rho }_{G}^{}$$
(22)
stands for the teleparallel version of the canonical energy-momentum pseudotensor of the gravitational field. Despite not explicitly apparent, as a consequence of the local Lorentz invariance of the gauge Lagrangian $`_G`$, the field equation (21) is symmetric in $`(\lambda \rho )`$. Furthermore, by using Eq. (12), it can be rewritten in terms of the Levi-Civita connection only. As expected, due to the equivalence between the corresponding Lagrangians, it is the same as Einstein’s equation:
$$\frac{h}{2}[\stackrel{}{R}_{\mu \nu }\frac{1}{2}g_{\mu \nu }\stackrel{}{R}]=0.$$
(23)
The canonical energy-momentum pseudotensor $`t_\lambda ^\rho `$ is not simply the gauge current $`j_a^\rho `$ with the algebraic index “$`a`$” changed to the spacetime index “$`\lambda `$”. It incorporates also an extra term coming from the derivative term of Eq. (16):
$$t_\lambda {}_{}{}^{\rho }=h^a{}_{\lambda }{}^{}j_{a}^{}{}_{}{}^{\rho }+\frac{c^4}{4\pi G}\mathrm{\Gamma }^\mu {}_{\lambda \nu }{}^{}S_{\mu }^{}{}_{}{}^{\nu \rho }.$$
(24)
We see thus clearly the origin of the connection-term which transforms the gauge current $`j_a^\rho `$ into the energy-momentum pseudotensor $`t_\lambda ^\rho `$. Through the same mechanism, it is possible to appropriately exchange further terms between the derivative and the current terms of the field equation (21), giving rise to different definitions for the energy-momentum pseudotensor, each one connected to a different superpotential $`(hS_\lambda {}_{}{}^{\rho \sigma })`$. Like the gauge current $`(hj_a{}_{}{}^{\rho })`$, the pseudotensor $`(ht_\lambda {}_{}{}^{\rho })`$ is conserved as a consequence of the field equation:
$$_\rho (ht_\lambda {}_{}{}^{\rho })=0.$$
(25)
However, in contrast to what occurs with $`j_a^\rho `$, due to the pseudotensor character of $`t_\lambda ^\rho `$, this conservation law can not be rewritten with a covariant derivative.
Because of its simplicity and transparency, the teleparallel approach to gravitation seems to be much more appropriate than general relativity to deal with the energy problem of the gravitational field. In fact, Møller already noticed a long time ago that a satisfactory solution to the problem of the energy distribution in a gravitational field could be obtained in the framework of a tetrad theory. In our notation, his expression for the gravitational energy-momentum density is
$$ht_\lambda {}_{}{}^{\rho }=\frac{}{_\rho h^a_\mu }_\lambda h^a{}_{\mu }{}^{}+\delta _\lambda {}_{}{}^{\rho },$$
(26)
which is nothing but the usual canonical energy-momentum density yielded by Noether’s theorem. Using for $``$ the gauge Lagrangian (14), it is an easy task to verify that Møller’s expression coincides exactly with the teleparallel energy-momentum density appearing in the field equation (21-22). Since $`j_a^\rho `$ is a true spacetime tensor, whereas $`t_\lambda ^\rho `$ is not, we can say that the gauge current $`j_a^\rho `$ is an improved version of the Møller’s energy-momentum density $`t_\lambda ^\rho `$. Mathematically, they can be obtained from each other by Eq. (24). It should be remarked, however, that both of them transform covariantly only under global tangent-space Lorentz transformations. This is, we believe, the farthest one can go in the direction of a tensorial definition for the energy and momentum of the gravitational field. The lack of a local Lorentz covariance can be considered as the teleparallel manifestation of the pseudotensor character of the gravitational energy-momentum density in general relativity. Accordingly, we can say that, if it were possible to define a local Lorentz covariant gauge current in the teleparallel gravity, the corresponding general relativity energy-momentum density would be represented by a true spacetime tensor.
The results can be summarized as follows. In the context of a gauge theory for the translation group, we have obtained an energy-momentum gauge current $`j_a^\rho `$ for the gravitational field which transforms covariantly under spacetime general coordinate transformations, and is invariant under local (gauge) translations of the tangent-space coordinates. This means essentially that $`j_a^\rho `$ is a true spacetime and gauge tensor. By rewriting the gauge field equation in a purely spacetime form, it becomes equivalent to Einstein’s equation of general relativity, and the gauge current $`j_a^\rho `$ reduces to the canonical energy-momentum pseudotensor of the gravitational field, which coincides with Møller’s well-known expression. In the ordinary context of general relativity, therefore, the energy-momentum density for the gravitational field will always be represented by a pseudotensor.
According to the quasilocal approach, to any energy-momentum pseudotensor there is an associated superpotential which is a hamiltonian boundary term . On the other hand, the teleparallel field equations explicitly exhibit both the superpotential and the gravitational energy-momentum complex. We see then that, in fact, by appropriately exchanging terms between the superpotential and the current terms of the field equation (21), it is possible to obtain different gravitational energy-momentum pseudotensors with their associated superpotentials. In this context, our results can be rephrased according to the following scheme. First, notice that the left-hand side of the field equation (21) as a whole is a true tensor, though each one of its two terms is not. Then if we extract the spurious part from the first term — so that it becomes a true spacetime and gauge tensor — and add this part to the second term — the energy-momentum density — it becomes also a true spacetime and gauge tensor. We thus arrive at the gauge-type field equation (16), with $`(hS_a{}_{}{}^{\sigma \rho })`$ as the superpotential, whose corresponding expression for the conserved energy-momentum density for the gravitational field, given by $`j_a^\rho `$, though transforming covariantly only under a global tangent-space Lorentz transformation, is a true spacetime and gauge tensor.
The authors would like to thank FAPESP-Brazil, CAPES-Brazil and CNPq-Brazil for financial support.
|
warning/0003/hep-ph0003003.html
|
ar5iv
|
text
|
# Inclusive jet production on the nucleus in the perturbative QCD with 𝑁_𝑐→∞
## 1 Introduction
In the colour dipole approach of A.H.Mueller with $`N_c\mathrm{}`$ the interaction with a heavy nucleus is exactly described by the sum of fan diagrams constructed of BFKL pomerons, each of them splitting into two . The equation for this sum \[3-6\] has recently been studied perturbatively in , by asymptotic estmates in and finally solved numerically in . The results indicate that the cross-section of a particle on the nucleus saturates at high energies to its geometrical value $`2\pi R_A^2`$ corresponding to the scattering on a black disk. In itself this is not a very surprising result: since the days of the old Regge-Gribov theory it has been known that the supercritical pomeron fan diagrams with a triple pomeron vertex lead to a constant cross-section . The new elements in the recent developments is that it has been understood that in the high-colour limit no splitting of the pomeron in more than two occurs and that the fan diagrams seem to be unitary by themselves so that they indeed represent the full scattering amplitude on the nucleus.
An interesting byproduct of the numerical solution in is the found gluon density of the nucleus in the combined momentum-impact parameter space
$$\frac{xG(x,q,b)}{^2q^2b}=\frac{N_c}{2\alpha _s\pi ^2}h(y,q,b)$$
(1)
where the rapidity $`y`$ is related to $`x`$ by
$$y=\mathrm{ln}\frac{1}{x}$$
Function $`h`$ defined via the sum of BFKL fan diagrams below (Section 2) and found in numerically is a soliton wave in $`y\mathrm{ln}q`$ space moving towards higher rapidities with a constant velocity and preserving its nearly Gaussian shape. A fit to numerical data gives
$$h(y,q,b)=h_0e^{a(\xi \xi _0(y,b))^2}$$
(2)
where $`\xi =\mathrm{ln}q`$, $`h_00.3`$ and $`a0.3`$ are universal and $`\xi _0`$ is nearly linear in $`y`$:
$$\xi _0=c(b)+\mathrm{\Delta }_0y,\mathrm{\Delta }_0=2.3\overline{\alpha }$$
(3)
where we standardly define $`\overline{\alpha }=\alpha _sN_c/\pi `$. The term $`c(b)`$ actually depends on a dimensionless parameter $`B`$:
$$B=\pi \alpha _s^2AT(b)R_N^2$$
where $`R_N`$ is the nucleon radius and $`T(b)`$ is the nucleus profile function. Numerical results for $`c(b)`$ are well described by
$$c(b)=3.11+(2/3)\mathrm{ln}B$$
Actually the gluon density does not really show itself in the total scattering cross-section at high rapidities. As mentioned, the cross-section becomes purely geometrical and insenstive to any dynamics. One expects however that the gluon density will be directly felt in production processes, where it is responsible for hard collisions giving rise to observable jets.
In this paper we study simplest of the production processes: the single and double inclusive jet production in hA collisions. It has to be noted that these observables were also studied for the fan diagrams in the old Regge-Gribov theory. The most important conclusion which was drawn from this study is that the multiplicities are strongly damped, as compared to the naive eikonal result. For hA scattering one obtains a multiplicity independent of $`A`$ (instead of $`A^{1/3}`$). As we shall see this result does not hold for the BFKL fan diagrams: the jet production rate grows with $`A`$ although somewhat slowlier than in the eikonal approach, approximately as $`A^{2/9}`$. The energy dependence remains esentially the same as in older studies: the production rate grows with rapidity $`Y`$ essentially as $`\mathrm{exp}\mathrm{\Delta }Y`$, where $`\mathrm{\Delta }`$ is the BFKL intercept. With $`\alpha _s0.2`$ its value is close to 0.5, so that the multiplicities reach very high values at rapidities in the region $`20÷30`$.
## 2 Single jet inclusive production
In the BFKL framework, at fixed impact parameter $`b`$, a single scattering contribution to the forward ampitude for the interaction of the projectile particle with a nucleus has a form
$$𝒜_1(Y,b)=isg^4AT(b)d^2r_1d^2r_2\rho (r)G(Y,r_1,r_2)\rho _N(r_2)2isd^2r_1\rho (r_1)\mathrm{\Phi }_1(y,r_1,b)$$
(4)
Here $`\rho `$ and $`\rho _N`$ are the colour densities of the projectile and the target nucleon respectively. Function $`G`$ is the forward BFKL Green function taken at an overall rapidity $`Y`$
$$G(Y,r,r^{})=\frac{rr^{}}{32\pi ^2}\underset{n=\mathrm{}}{\overset{+\mathrm{}}{}}e^{in(\varphi \varphi ^{})}_{\mathrm{}}^{\mathrm{}}\frac{d\nu e^{Y\omega (\nu )}}{[\nu ^2+(n1)^2/4][\nu ^2+(n+1)^2/4]}(r/r^{})^{2i\nu },$$
(5)
where $`\varphi `$ and $`\varphi ^{}`$ are the azimuthal angles and
$$\omega (\nu )=\overline{\alpha }(\psi (1)\mathrm{Re}\psi (1/2+i\nu ))$$
(6)
are the BFKL levels. Due to the azimuthal symmetry of the projectile colour density one may retain only the term with zero orbital momenta $`n=0`$ in (5). The corresponding inclusive jet production rate
$$I(y,k,b)\frac{(2\pi )^2\sigma }{y^2k^2b}$$
(7)
is obtained from the imaginary part of (4) divided by $`s`$ after the substitution
$$G(Y,r_1,r_2)d^2rG(y_1,r_1,r)V_k(r)G(y,r,r_2)$$
(8)
where $`r`$’s are relative distances between the gluons and $`V_k(r)`$ is a vertex for the emission
$$V_k(r)=\frac{4N_c\alpha _s}{k^2}\stackrel{}{\mathrm{\Delta }}e^{ikr}\stackrel{}{\mathrm{\Delta }}$$
(9)
local in $`r`$ (a differential operator, the arrows shows the direction of its action). The rapidity of the produced jet relative to the projectile is $`y_1=Yy`$.
As mentioned, in the limit $`N_c\mathrm{}`$ the total scattering amplitude which includes multiple interactions in the target is given by the sum of fan diagrams shown in Fig 1. To obtain the corresponding inclusive cross-section, one has to make the substitition (8) in one of the pomerons. However due to the well-known AGK rules all such substitutions in the pomerons below the upper splitting point (Fig. 2a) cancel, since the other branch, counting from the upper point, can be both ”cut” and ”uncut” (i.e. represent both real and virtual processes). The only left contribution is the one with the substitution (8) in the uppermost pomeron (Fig. 2b). It is clearly seen from Fig. 2b that the vertex $`V_k`$ is coupled from below to exactly the full amplitude $`\mathrm{\Phi }(y,r,b)`$ which multiplies the projectile density in Fig. 1 and substitutes the single pomeron exchange contribution $`\mathrm{\Phi }_1(Y,r_1,b)`$ in (4). Thus one obtains the inclusive cross-section for jet production on the nucleus at fixed $`b`$ as
$$I(y,k,b)=2d^2r^{}d^2r\rho (r^{})G(y_1,r^{},r)V_k(r)\mathrm{\Phi }(y,r,b)$$
(10)
The total inclusive production rate as a function of $`y`$ and $`k`$ is obtained from (10) after the integration over $`b`$
Function $`\varphi (y,r,b)=\mathrm{\Phi }(y,r,b)/(2\pi r^2)`$, in the momentum space, satisfies a nonlinear equation
$$\frac{\varphi (y,q,b)}{\stackrel{~}{y}}=H\varphi (y,q,b)\varphi ^2(y,q,b),$$
(11)
where $`\stackrel{~}{y}=\overline{\alpha }y`$ and $`H`$ is the BFKL Hamiltonian for the so-called semi-amputated function . As mentioned in the Introduction, this equation was solved numerically in .
The Laplace operator in (9) acting on $`\mathrm{\Phi }(y,r,b)`$ gives
$$_r^2\mathrm{\Phi }(y,r,b)=2\pi _r^2r^2\frac{d^2q}{(2\pi )^2}e^{iqr}\varphi (y,q,b)=2\pi _r^2\frac{d^2q}{(2\pi )^2}e^{iqr}_q^2\varphi (y,q,b)$$
$$=2\pi \frac{d^2q}{(2\pi )^2}e^{iqr}q^2_q^2\varphi (y,q,b)=2\pi h(y,r,b)$$
(12)
Note that in the integration by parts used to obtain the second expression one has to be careful, since function $`\varphi (y,q,b)`$ has a logarithmic singularity at $`q=0`$ of the form $`\mathrm{ln}q`$. As a result, action of the Laplacian operator in $`q`$ generates a $`\delta ^2(q)`$ term which in its turn leads to a constant in $`\mathrm{\Phi }`$ (actually unity). However this term is annihilated by the Laplacian operator in $`r`$, so that no extra term arises from the intergration by parts.
Function $`h(y,q,b)`$, which is the gluon density, up to a trivial factor (Eq. (1)) is just a Fourier transform of $`h(y,r,b)`$:
$$h(y,q,b)=q^2_q^2\varphi (y,q,b)$$
(13)
So we find the inclusive cross-section in the form
$$I(y,k,b)=\frac{16\pi \alpha _sN_c}{k^2}d^2r^{}d^2r\rho (r^{})_r^2G(y_1,r^{},r)e^{irk}h(y,r,b)$$
(14)
We observe that the definition of $`h(y,q,b)`$ as the gluon density introduced in is supported by its role in the jet production on the nucleus.
Applying the Laplace operator to the Green function (5) (with only $`n=0`$ terms retained) we obtain
$$_r^2G(y_1,r^{},r)=\frac{r^{}}{8\pi ^2r}_{\mathrm{}}^{\mathrm{}}\frac{d\nu e^{y_1\omega (\nu )}}{(\nu +i/2)^2}(r^{}/r)^{2i\nu },$$
(15)
We are going to study the production rate at $`y_1>>1`$, that is far from the projectile fragmentation region. Then we can use the asymptotics of (15) at high $`y_1`$, which comes from $`\nu 0`$. Presenting $`\omega (\nu )`$ in the standard manner
$$\omega (\nu )=\mathrm{\Delta }y\beta \nu ^2$$
(16)
where
$$\mathrm{\Delta }=\overline{\alpha }\mathrm{\hspace{0.17em}4}\mathrm{ln}2,\mathrm{and}\beta =\overline{\alpha }\mathrm{\hspace{0.17em}14}\zeta (3)$$
(17)
we get for (15)
$$_r^2G(y_1,r^{},r)\frac{r^{}}{2\pi ^2r}e^{\mathrm{\Delta }y_1}\sqrt{\frac{\pi }{\beta y_1}}$$
(18)
With (18), integration over $`r^{}`$ converts $`r^{}`$ into the average transverse dimension of the projectile $`r_P=R_P`$. Doing the integration over $`r`$ we obtain
$$I(y,k,b)=\frac{8\overline{\alpha }}{k^2}R_Pe^{\mathrm{\Delta }y_1}\sqrt{\frac{\pi }{\beta y_1}}F(y,k,b)$$
(19)
where
$$F(y,k,b)=\frac{d^2r}{r}e^{ikr}h(y,r,b)=\frac{d^2q}{2\pi }\frac{h(y,q,b}{|k+q|}$$
(20)
As mentioned in the Introduction, the gluon density $`h(y,q,b)`$ has a sharp peak at $`q=q_0(y,b)=\mathrm{exp}\xi _0(y,b)`$ and is practically zero everywhere else. With growing $`y`$ $`q_0(y,b)`$ grows exponentially, so that at high $`y`$ we can safely neglect $`k`$ in the denominator of (20) as compared to $`q`$. Then $`F`$ becomes independent of $`k`$ and given by
$$F(y,b)=\frac{d^2q}{2\pi q}h(y,q,b)=_{\mathrm{}}^+\mathrm{}𝑑\xi qh(y,q,b)$$
(21)
where $`\xi =\mathrm{ln}q`$. One easily finds that
$$𝑑\xi h(y,q,b)=1$$
(22)
Indeed Eq. (13) can be rewritten as
$$h(y,q,b)=\frac{d^2}{d\xi ^2}\varphi (y,q,b)$$
(23)
Putting (23) into (21) and taking into account that at $`\xi +\mathrm{}`$ $`\varphi =const\xi `$ and at $`\xi +\mathrm{}`$ $`\varphi 0`$ we obtain (22). So (21) in fact has a meaning of the average gluon momentum in the nucleus $`q`$. Since $`h`$ has a sharp maximun at $`\xi =\xi _0(y,b)`$ given by (3) one expects that $`F(y,b)`$ will approximately be given just by $`q_0(y,b)`$. Numerical integration gives more accurate values which are well fitted by the formula similar to (3) but with slightly different parameters:
$$F(y,b)=q_1(y,b)=1.2B^{2/3}\frac{e^{\mathrm{\Delta }_1y}}{R_N\sqrt{y}},\mathrm{\Delta }_1=2.38\overline{\alpha }$$
(24)
With this approximate form of $`F`$ we find the final expression for the inclusive rate production (in (GeV/c)<sup>-2</sup>)
$$I(y,k,b)=1.2\frac{8\overline{\alpha }}{k^2}\frac{R_P}{R_N}[\pi \alpha _s^2AT(b)R_N^2]^{2/3}e^{\mathrm{\Delta }Yϵy}\sqrt{\frac{\pi }{\beta y(Yy)}}$$
(25)
where
$$ϵ=\mathrm{\Delta }\mathrm{\Delta }_1=0.39\overline{\alpha }$$
and is relatively small. As we observe, the inclusive jet production is maximal in nucleus fragmentation region and slowly falls as one moves to higher $`y`$.
Since the total hA cross-section saturates at large $`Y`$ an immediate consequence of (25) is that the multiplicity grows like $`\mathrm{exp}\mathrm{\Delta }Y`$, that is, as the cross-section generated by a single pomeron exchange. This conclusion is also true for the inclusive production from the sum of fan diagrams in the old Regge-Gribov theory with a local pomeron. However this theory predicts a much stronger $`y`$-dependence, corresponding to (25) with $`ϵ=\mathrm{\Delta }`$.
As to the $`A`$-dependence, a novel element in the BFKL fan diagrams is factor $`[\pi \alpha _s^2AT(b)R_N^2]^{2/3}`$. For a nucleus with constant density inside a sphere of radius $`R_A=A^{1/3}R_0`$ its integration over all $`b`$ leads to a factor
$$2\pi R_A^2A^{2/9}\frac{3}{8}\gamma ,\gamma =\left(\frac{3}{2}\frac{R_N^2}{R_0^2}\right)^{2/3}$$
(26)
Since the total hA cross-section at large $`Y`$ saturates to its black-disc limit $`2\pi R_A^2`$ the jet multiplicity results proportional to $`A^{2/9}`$:
$$\mu (y,k)=\frac{1}{\sigma _{hA}^{tot}}d^2bI(y,k,b)=1.2A^{2/9}\frac{3\overline{\alpha }}{k^2}\frac{R_P}{R_N}\gamma e^{\mathrm{\Delta }Yϵy}\sqrt{\frac{\pi }{\beta y(Yy)}}$$
(27)
So unlike the fan diagrams in the old Regge-Gribov theory, the BFKL model predicts rising of the multiplicities with $`A`$, although somewhat weaker than according to the eikonal model ($`A^{1/3}`$).
Although the gluon distribution in the nucleus results radically changed as compared to the single nucleon, in particular, strongly damped at small momenta, this seems to have no effect on the $`k`$-dependence of the jet production rate at moderate $`k<<q_1(y,b)`$: it remains proportional to $`1/k^2`$ and growing as $`k0`$. The inclusive cross-section integrated over all $`k`$ remains infinite, as for a single BFKL pomeron exchange. Thus discussing the integrated cross-section we have to cutoff the spectrum from below by some $`k_{min}`$. The formal expression for this cross-section can be trivially found integrating (19) over all $`k>k_{min}`$. A simple estimate of its magnitude can be made in the logarithmic approximation in the large parameter $`q_1(y,b)/k_{min}`$. Evidently
$$\frac{d^2k}{(2\pi )^2k^2}F(y,k,b)\theta (kk_{min})=\frac{d^2q}{2\pi }h(y,q,b)\frac{d^2k}{(2\pi )^2k^2}\frac{1}{|k+q|}\theta (kk_{min})$$
$$\frac{1}{2\pi }\frac{d^2q}{2\pi q}h(y,q,b)\mathrm{ln}(q/k_{min})=\frac{1}{2\pi }q\mathrm{ln}(q/k_{min})$$
(28)
So, up to factor $`1/2\pi `$, integration over $`k`$ leads to the average value of $`q\mathrm{ln}(q/k_{min})`$ in the nucleus. Approximating it by $`q_1(y,b)\mathrm{ln}(q_1(y,b)/k_{min})`$, so that for fixed $`b`$ one finds
$$I(y,b)\frac{d^2k}{(2\pi )^2}\theta (kk_{min})=\frac{4\overline{\alpha }}{\pi }R_Pe^{\mathrm{\Delta }(Yy)}\sqrt{\frac{\pi }{\beta (Yy)}}q_1(y,b)\mathrm{ln}\frac{q_1(y,b)}{k_{min}}$$
(29)
with $`q_1(y,b)`$ given by (24).
At high $`y`$ the logarithmic factor reduces to $`\mathrm{\Delta }_1y`$. Its effect, opposite to the exponential, is to flatten the distribution $`I(y)`$ as compared to $`I(y,k)`$. The $`A`$ dependence evidently becomes somewhat stronger. Numerical integration of $`I(k,y,b)`$ with function $`h`$ found in gives results which more or less agree with this approximate estimates. At $`\overline{\alpha }Y=6`$ for lead ($`A=207`$) we obtain multiplicities which in the region $`2\overline{\alpha }y5`$ very slowly grow from $`\mathrm{1.13\hspace{0.17em}10}^5`$ to $`\mathrm{1.26\hspace{0.17em}10}^5`$. The found $`A`$ dependence can be fitted by $`A^{0.26}`$. The large magnitude of the multiplicity is due to the exponential factor $`\mathrm{exp}\mathrm{\Delta }Y=\mathrm{exp}(\overline{\alpha }Y4\mathrm{ln}2)`$.
## 3 Double inclusive cross-sections
Passing to double inclusive cross-sections and using the AGK rules we find two contributions schematically shown in Figs. 3 a,b.
A simpler contribution Fig. 3a corresponds to emission of both gluons from the upper pomeron. To obtain it we have to make two insertions (8) in the upper BFKL Green function. We get
$$I_a(y_1,k_1,y_2,k_2,b)\left[\frac{(2\pi )^4\sigma }{y_1^2k_1y_2d^2k_2^2b}\right]_{Fig.3a}$$
$$=2d^2r^{}d^2r_1d^2r_2\rho (r^{})G(Yy_1,r^{},r_1)V_{k_1}(r_1)G(y_1y_2,r_1,r_2)V_{k_2}(r_2)\mathrm{\Phi }(y_2,r_2,b)$$
(30)
where we assumed that $`Y>>y_1>>y_2>>1`$ in the lab. system.
The only new element in this equation is the second BFKL Green function with Laplacians applied in both coordinates. Using (5) we find
$$\mathrm{\Delta }_1\mathrm{\Delta }_2G(y,r_1,r_2)=\frac{1}{2\pi ^2r_1r_2}𝑑\nu e^{y\omega (\nu )}\left(\frac{r_1}{r_2}\right)^{2i\nu }$$
(31)
Combining this with (15) for the first Green function (with the integration variable $`\nu _1`$) we find an integral over $`r_1`$:
$$d^2r_1e^{ik_1r_1}r_1^{22i\nu +2i\nu _1}=\pi (k_1/2)^{2i(\nu \nu _1)}\frac{1}{i(\nu _1\nu i0)}\frac{\mathrm{\Gamma }(1+i\nu _1i\nu )}{\mathrm{\Gamma }(1i\nu _1+i\nu )}$$
(32)
As before we limit ourserlves with the leading contribution as both $`Yy_1`$ and $`y_1y_2`$ become large, which correspond to small values of both $`\nu `$ and $`\nu _1`$. Correspondingly we put $`\nu =\nu _1=0`$ everywhere except in the exponents $`\omega (\nu _1)(Yy_1)`$ and $`\omega (\nu )(y_1y_2)`$ and in the denominator of (32). We find that only the real part of (32) contributes, which is equivalent to taking
$$\frac{1}{i(\nu _1\nu i0)}\pi \delta (\nu _1\nu )$$
Using (16) we finally get for the double inclusive cross-section (30)
$$I_a(y_1,k_1,y_2,k_2,b)=16\pi R_P\frac{\overline{\alpha }^2}{k_1^2k_2^2}e^{\mathrm{\Delta }(Yy_2)}\sqrt{\frac{\pi }{\beta (Yy_2)}}F(y_2,k_2,b)$$
(33)
where function $`F(y,k,b)`$ was defined in Section 2 (Eq.(20)). A remarkable feature of this cross-section is that it depends only on the rapidity of the slowlier jet $`y_2`$ and does not depend on $`y_1`$ altogether. In the approximation (24) the dependence on $`y_2`$ results the same as for the single inclusive cross-section, namely $`\mathrm{exp}(ϵy_2)`$, so that the cross-section slowly diminishes as $`y_2`$ increases.
Using approximation (24), integrating over $`b`$ and dividing by $`\sigma _{hA}^{tot}`$ we get the corresponding multiplicity distribution
$$\mu _a(y_1,k_1,y_2,k_2)=1.2A^{2/9}6\pi \frac{R_P}{R_N}\gamma \frac{\overline{\alpha }^2}{k_1^2k_2^2}e^{\mathrm{\Delta }Yϵy_2}\sqrt{\frac{\pi }{\beta y_2(Yy_2)}}$$
(34)
where $`\gamma `$ was defined in (26).
The second contribution (Fig. 3b) corresponds to emission of both jets just after the first splitting of the upper pomeron. To write it we first take the contribution from all non-trivial fan diagrams to $`\mathrm{\Phi }`$
$$\mathrm{\Phi }(Y,r^{},b)=\frac{g^2N_c}{8\pi ^3}_0^Y𝑑y^{}\underset{i=1}{\overset{3}{}}d^2x_i\delta ^2(x_1+x_2+x_3)$$
$$\frac{x_3^2_{3}^{}{}_{}{}^{4}}{x_1^2x_2^2}G(Yy^{},r^{},x_3)\mathrm{\Phi }(y^{},x_1,b)\mathrm{\Phi }(y^{},x_2,b).$$
(35)
where we denoted $`x_i,i=1,2,3`$ the transverse coordinates of the triple pomeron vertex. According to the AGK rules the corresponding contribution to the absorptive part with both lower branches cut is twice (35) with the opposite sign. To finally obtain the double inclusive cross-section we have to insert into both $`\mathrm{\Phi }`$’s on the right-hand side verteces $`V_k(r)`$ as was done in (10). In this way we obtain the second part of the double inclusive cross-section as
$$I_b(y_1,k_1,y_2,k_2,b)=\frac{g^2N_c}{2\pi ^3}d^2r^{}\rho (r^{})d^2r_1d^2r_2_{y_1}^Y𝑑y^{}\underset{i=1}{\overset{3}{}}d^2x_i\delta ^2(x_1+x_2+x_3)$$
$$\frac{x_3^2_{3}^{}{}_{}{}^{4}}{x_1^2x_2^2}G(Yy^{},r^{},x_3)G(y^{}y_1,x_1,r_1)V_{k_1}(r_1)\mathrm{\Phi }(y_1,r_1,b)G(yy_2,x_2,r_2)V_{k_2}\mathrm{\Phi }(y_2,r_2,b).$$
(36)
We again assume $`Y>>y_1>>y_2>>1`$ so that the rapidity of the splitting point $`y^{}`$ varies from $`y_1`$ to $`Y`$.
We easily find
$$x_3^2_3^4G(y,r^{},x_3)=\frac{r^{}}{2\pi ^2x_3}𝑑\nu _3e^{y\omega (\nu _3)}\left(\frac{r^{}}{x_3}\right)^{2i\nu _3}$$
(37)
Using this and representation (15) for the other two BFKL Green functions we find an integral over $`x_i`$, $`i=1,2,3`$
$$J(\nu _i)=\underset{i=1}{\overset{3}{}}d^2x_i\delta ^2(x_1+x_2+x_3)\frac{1}{x_1x_2x_3}x_1^{2i\nu _1}x_2^{2i\nu _2}x_3^{2i\nu _3}$$
(38)
This integral diverges at large $`x_i`$’s for real $`\nu `$’s. In fact, presenting the $`\delta `$ function as an integral over momentum $`q`$ we find
$$J=2\pi \underset{l=1}{\overset{3}{}}2^{2i\nu _l}\frac{\mathrm{\Gamma }(1/2i\nu _l)}{\mathrm{\Gamma }(1/2+i\nu _l)}\frac{d^2q}{q^3}q^{2i(\nu _1+\nu _2+\nu _3)}$$
(39)
which does not exist for real $`\nu `$ due to the divergence at $`q=0`$. To overcome this difficulty we have to shift the integration contours in $`\nu `$’s into the lower half-plane, using the fact that integrands are analytic in the strip $`i/2<\mathrm{Im}\nu <+i/2`$. We choose to do it in a symmetric manner putting
$$\nu _i=i/6+\overline{\nu }_i,i=1,2,3$$
(40)
With such a choice we get an extra factor $`q`$ in the integral over $`q`$, which then gives a $`\delta `$ function:
$$J=\pi ^2\underset{l=1}{\overset{3}{}}\frac{\mathrm{\Gamma }(1/2i\nu _l)}{\mathrm{\Gamma }(1/2+i\nu _l)}\delta (\overline{\nu }_1+\overline{\nu }_2+\overline{\nu }_3)$$
(41)
Now we pass to the integration over the rapidity of the splitting point. The integral is done trivially:
$$_{y_1}^Y𝑑y^{}e^{y(\omega _1+\omega _2\omega _3)}=\frac{e^{Y(\omega _1+\omega _2\omega _3)}e^{y_1(\omega _1+\omega _2\omega _3)}}{\omega _1+\omega _2\omega _3}$$
(42)
where we denote $`\omega _1=\omega (\nu _1)`$ etc. Multiplying by the rest exponential factors we find the overall rapidity factor as
$$\frac{e^{(Yy_1)\omega _1+(Yy_2)\omega _2}e^{(Yy_1)\omega _3+(y_1y_2)\omega _2}}{\omega _1+\omega _2\omega _3}$$
(43)
Note that each of the two terms depends only on two of the three $`\nu `$ variables. Therefore integrating over $`\nu `$’s we can always take as independent variables just those which are present in the exponents in (43).
Having this in mind we again use the fact that all rapidity differences in the exponents are large and use the saddle point method. Athough the initial contours of integration in the independent $`\nu `$’s lie along the lines $`\mathrm{Im}\nu =1/6`$, the saddle points remain at $`\nu =0`$, that is at $`\overline{\nu }=(1/6)i`$. This means that independent $`\nu `$’s will be close to zero, as usual. The only price we have to pay for the shift in the contours of integration will be the value of the dependent $`\nu `$ at the saddle point.
Take the first term in (43), for which $`\nu _1`$ and $`\nu _2`$ are independent variables and $`\nu _3`$ is the dependent one. At the saddle point we shall have $`\nu _1=\nu _2=0`$ and $`\nu _3=(1/6)i\overline{\nu _1}\overline{\nu _2}=(1/2)i`$. Accordingly we put $`\nu _1=\nu _2=0`$ and $`\nu _3=(1/2)i`$ in all terms except in the exponent and the singular $`\mathrm{\Gamma }`$-function:
$$\mathrm{\Gamma }(1/2i\nu _3)=\mathrm{\Gamma }(i\nu _1+i\nu _2)=\frac{\mathrm{\Gamma }(1+i\nu _1+i\nu _2)}{i(\nu _1+\nu _2i0)}\frac{1}{i(\nu _1+\nu _2i0)}$$
(44)
Again only the real part contributes, so that (43) actually gives
$$\pi \delta (\nu _1+\nu _2)$$
Using (16) and integrating over $`\nu _1`$ and $`\nu _2`$ we obtain a factor
$$\frac{1}{\mathrm{\Delta }}e^{\mathrm{\Delta }(2Yy_1y_2)}\sqrt{\frac{\pi }{\beta (2Yy_1y_2)}}$$
(45)
and an extra factor $`r^{}`$ due to the fact that $`\nu _3=(1/2)i`$ at the saddle point.
For the second term in (43) the independent variables are $`\nu _2`$ and $`\nu _3`$ which are zero at the saddle point. The dependent $`\nu _1`$ is equal to $`(1/2)i`$ at the saddle point. The $`\mathrm{\Gamma }`$ functions in (41) will give $`\pi \delta (\nu _2+\nu _3)`$, so that the integration over $`\nu _2`$ and $`\nu _3`$ will give a factor
$$\frac{1}{\mathrm{\Delta }}e^{\mathrm{\Delta }(Yy_2)}\sqrt{\frac{\pi }{\beta (Yy_2)}}$$
(46)
and we shall have an extra factor $`r_1`$ due to $`\nu _1=(1/2)i`$ at the saddle point.
Other factors in (36) are straightforward. Combining all of them we find the expression for the double inclusive cross-section corresponding to Fig. 2b in the form
$$I_b(y_1,k_1,y_2,k_2,b)=\frac{4}{\mathrm{ln}2}\frac{\overline{\alpha }^2}{k_1^2k_2^2}e^{\mathrm{\Delta }(2Yy_1y_2)}\sqrt{\frac{\pi }{\beta (2Yy_1y_2)}}F(y_2,k_2,r)$$
$$\left[r^2_PF(y_1,k_1,b)R_Pe^{\mathrm{\Delta }(Yy_1)}\sqrt{\frac{2Yy_1y_2}{Yy_2}}h(y_1,k_1,b)\right]$$
(47)
where
$$r^2_P=d^2rr^2\rho (r)$$
(48)
The second term is evidently exponentially small as compared to the first and we can safely neglect it, having in mind all approximation already made. So our final expression is
$$I_b(y_1,k_1,y_2,k_2,b)=$$
$$\frac{4}{\mathrm{ln}2}r^2_P\frac{\overline{\alpha }^2}{k_1^2k_2^2}e^{\mathrm{\Delta }(2Yy_1y_2)}\sqrt{\frac{\pi }{\beta (2Yy_1y_2)}}F(y_1,k_1,b)F(y_2,k_2,r)$$
(49)
Recalling our expression (19) for the single inclusive cross-section we can rewrite it as
$$I_b(y_1,k_1,y_2,k_2,b)=\frac{1}{16\mathrm{ln}2}\frac{r^2_P}{r_P^2}\sqrt{\frac{\beta }{\pi }\frac{(Yy_1)(Yy_2)}{2Yy_1y_2}}I(y_1,k_1,b)I(y_2,k_2,b)$$
(50)
After the integration over $`b`$ with the approximation (24) we get the corresponding multiplicity distribution:
$$\mu _b(y_1,k_1,y_2,k_2)=1.44A^{4/9}\frac{6}{5\mathrm{ln}2}\frac{r^2_P}{R_N^2}\gamma ^2\frac{\overline{\alpha }^2}{k_1^2k_2^2}e^{2\mathrm{\Delta }Yϵ(y_1y_2)}\sqrt{\frac{\pi }{\beta y_1y_2(2Yy_1y_2)}}$$
(51)
Evidently at high $`Y`$ the part $`\mu _b`$ dominates due to the factor $`\mathrm{exp}2\mathrm{\Delta }Y`$ as compared with only $`\mathrm{exp}\mathrm{\Delta }Y`$ in $`\mu _a`$
Correlations are determined by the ratio
$$R(Y,y_1,y_2)=\frac{\mu (y_1,k_1,y_2,k_2)}{\mu (y_1,k_1)\mu (y_2,k_2)}$$
(52)
(evidently it does not depend on $`k_1`$ and $`k_2`$). Neglecting the part $`\mu _a`$ we find from (27) and (51)
$$R(Y,y_1,y_2)=\frac{2}{15\mathrm{ln}2}\frac{r^2_P}{R_P^2}\sqrt{\frac{\beta }{\pi }\frac{(Yy_1)(Yy_2)}{2Yy_1y_2}}$$
(53)
If the two jets, forward and backward in the c.m. system, are taken symmetric, that is,
$$y_1=\frac{1}{2}Y+y,y_2=\frac{1}{2}Yy$$
(54)
Eq. (53) becomes
$$R(Y,y)=\frac{1}{15\mathrm{ln}2}\frac{r^2_P}{R_P^2}\sqrt{\frac{\beta }{\pi }\left(Y\frac{4y^2}{Y}\right)}$$
(55)
With the rapidity gap between the jet fixed and growing energy the term $`4y^2/Y`$ can be neglected, so that the right-hand side of (55) becomes independent of the gap $`2y`$ and proportional to $`\sqrt{Y}`$. This means that eventually it inevitably becomes greater than unity, so that the correlations become positive. The concrete threshold for this depends on the values of $`\overline{\alpha }`$ and ratio $`r^2/R_P^2`$. With the Yukawa distribution of colour in the projectile the last ratio is equal 2. Taking $`\alpha _s=0.2`$ we get for (55)
$$R0.195\sqrt{Y}$$
(56)
So with these values correlations remain negative up to quite high energies corresponding to $`Y<25`$.
Note that in the framework of the old Regge-Gribov theory, due to the assumed relations between the parameters and smallness of the intercept $`\mathrm{\Delta }`$, it was argued that the ratio $`R`$ for the part $`\mu _b`$ is identically equal to unity, so that correlations are entirely generated by the part $`\mu _a`$ . In the present approach this does not hold. Although the factors exponential in rapidities are indeed the same in $`\mu _b`$ and the product $`\mu (y_1,k_1)\mu (y_2,k_2)`$, the power factors are not, so that their final $`Y`$-dependence is different. Also the relation between the coefficients is not fixed but depends on $`\alpha _s`$. As a result the ratio $`R`$ for $`\mu _b`$ is different from unity and depends on energy and $`\alpha _s`$, so that the correlations become mainly determined by $`\mu _b`$ and thus much stronger than assumed in earlier studies.
## 4 Conclusions
We have studied the inclusive jet production off the nucleus in the perturbative QCD with $`N_c`$ infinity, in which the total hA amplitude is exactly given by a sum of fan diagrams constructed of BFKL pomerons. We have used the numerical results obtained for this sum in our earlier paper . Our main results are the following.
Our formulas confirm interpretation of function $`h(y,q,b)`$ defined by Eq. (13) as the gluon density in the nucleus. In this was done by comparison with the perturbative expression for the structure function, the validity of which for large values of the gluon density can be questioned.
The inclusive jet production rate is found to rapidly grow with energy in the same manner as the cross-section generated by a single pomeron exchange. This result is an immediate consequence of the saturation of the total hA cross-section, and was also found in earlier studies of the fan diagrams with local pomerons.
However in contrast to these older studies, our jet production rate is $`A`$ dependent: jet multiplicities grow roughly as $`A^{2/9}`$. So the $`A`$ dependence is somewhat reduced as compared to the eikonal result $`A^{1/3}`$, but not so strongly as thought before. This result may have direct consequences for particle production in nucleus-nucleus collisions, where we may expect the multiplicities grow as $`A^{10/9}`$ for identical nuclei, less than in the eikonal approach ($`A^{4/3}`$) but greater than in the older fan diagram estimates ($`A^{2/3}`$).
The double inclusive jet production is dominated by the contribution from Fig. 2b. Unlike older studies, this contribution is not cancelled in the correlation function and leads to correlations which start negative at lower energies and then change to positive at a certain (rather high) energy. Their magnitude corresponds to a double pomeron exchange and so grows very large at high energies.
In conclusion we want to remark that since the gluon distribution in the nucleus is strongly shifted to higher momenta, the notorious difficulty of the diffusion towards very small momenta seems to be overcome. This problem always raised doubts of the validity of the perturbative treatment of hadronic amplitudes: even if the initial wave function is centered at high-momenta, its evolution in rapidity inevitably introduces low momentum contributions, for which the perturbative treatment is not valid. The large nucleus seems to damp these small momenta contributions, so that the use of the pertubation theory becomes justified.
## 5 References
1. A.Mueller, Nucl. Phys.,B415 (1994) 373.
2. A.Mueller and B.Patel, Nucl. Phys.,B425 (1994) 471.
3. M.A.Braun and G.P.Vacca, Eur. Phys. J C6 (1999) 147.
4. Yu. Kovchegov, Phys. Rev D60 (1999) 034008.
5. I.Balitsky, hep-ph/9706411; Nucl. Phys. B463 (1996) 99.
6. M.A.Braun, hep-ph/0001268
7. Yu. Kovchegov, preprint CERN-TH/99-166 (hep-ph/9905214).
8. E.Levin and K.Tuchin, preprint DESY 99-108, TAUP 2592-99 (hep-ph/9908317).
9. A.Schwimmer, Nucl. Phys. B94 (1975)445.
10. L.N.Lipatov in: ”Perturbative QCD”, Ed. A.H.Mueller, World Sci., Singapore (1989) 411.
## 6 Figure captions
Fig. 1. The amplitude for hA scattering. Lines denote BFKL pomerons.
Fig. 2. Inclusive jet production from the lower (a) and upper (b) pomerons.
Fig. 3. Double inclusive jet production from the upper pomeron (a) and immediately after the first spliting (b).
|
warning/0003/cond-mat0003353.html
|
ar5iv
|
text
|
# Universality of the Kondo Effect in a Quantum Dot out of Equilibrium
## I Introduction
The Kondo effect results from exchange interaction of itinerant electrons with a localized spin state. This interaction leads to local spin polarization of the electron gas. The polarization becomes significant only at low temperatures, due to the existence of collective states with small excitation energies. Simultaneously with the modification of the spin susceptibility, the scattering properties are modified significantly from those of the original localized spin state. The modification is especially striking in the case of antiferromagnetic exchange interaction, when the “spin cloud” formed out of free electron gas tends to screen the magnetic moment of the localized state. In this case, the scattering cross-section grows as the temperature is being lowered, and reaches the unitary limit at $`T0`$. This phenomenon is responsible for the non-monotonic temperature behavior of the resistivity of metals with magnetic impurities at low temperatures, which was the first experimentally observed manifestation of the Kondo resonance. However, the system of magnetic impurities embedded in a metal sample does not offer much control over the parameters even at the stage of the fabrication of the sample, not speaking of the dynamic variation of the parameters in the course of an experiment.
Another class of systems, whose transport properties can also be affected by the Kondo effect, but which offer a much higher level of control over the system parameters, is provided by quantum dots. A quantum dot in a semiconductor planar heterostructure is a confined few-electron system contacted by sheets of two-dimensional electron gas (leads). If the total number of electrons on the dot is odd, then the dot is similar to a magnetic impurity. Junctions between the dot and the leads produce an overlap of the states in the dot and in the 2D leads. This overlap leads to the exchange interaction between the spin of the dot and spins of the itinerant 2D electrons. At sufficiently low temperatures, the “spin cloud” is formed of the electrons in the leads. Like in bulk metals, the scattering off the resulting many-body state is being enhanced as the temperature is being lowered, and reaches unitary limit at $`T0`$. Different are only the manifestations of the Kondo resonance in bulk metals and in the quantum dot systems. In bulk metals, the enhancement of scattering by the Kondo resonance increases the resistivity. In quantum dot systems, on the contrary, the scattering facilitates transport through the dot. Kondo effect results in a specific temperature dependence of the linear conductance across the dot. If temperature is lowered, the conductance first drops due to the Coulomb blockade phenomenon, and then increases again due to the Kondo effect. At $`T0`$, the conductance $`G`$ reaches its maximum, which corresponds to the unitary limit of tunneling. Recently, the Kondo effect in a quantum dot was observed experimentally.
The quantum dot devices are highly controllable, and can be operated in regimes inaccessible in the conventional magnetic impurity systems, that were used previously for studying the Kondo effect. Irradiation of a quantum dot with an ac field offers a new, clever way of affecting its dynamics, which enables one to study the Kondo anomaly in essentially non-equilibrium conditions. The ac field can be applied to the gate, thus modulating the dot’s potential with respect to the leads; alternatively, one may apply ac bias to the leads. In any case, driving the system out of equilibrium affects the dc conductance discussed above. Measuring the dc $`IV`$ characteristics, one can investigate the effects of the irradiation on the Kondo anomaly.
A generic theoretical description of a quantum dot operates with a significant number of parameters and energy scales describing the system. Nevertheless, in the case of no ac field, it turns out that the low-energy properties of the quantum dot system which are related to the Kondo effect are controlled by only one relevant energy scale, which is the Kondo temperature $`T_K`$. The Kondo temperature, in turn, depends on the microscopic parameters of the system, e.g., on the gate voltage and conductances of the dot-lead junctions. Such universality allows for easier understanding and description of the problem.
The ac field introduces new parameters to the problem, thus apparently breaking the universal description, which is valid in the static case. This re-emerging abundance of parameters makes it difficult to build a consistent description of the effects of the ac field on the Kondo conductance. The theoretical work performed up to the date, was concentrated on some specific regimes. Goldin and Avishai considered the case of very strong ac bias with the help of the third-order perturbation theory in the dot-lead coupling. Nordlander et al. analyzed the effects of ac field of sufficiently high frequency, which ionizes the dot. They have conjectured that even at the temperature of the thermal bath $`T=0`$, the finite rate of the dot ionization results in a finite effective temperature “seen” by the Kondo state. This way, irradiation provides the cut-off for the Kondo singularity and reduces the conductance. Later, we have demonstrated that even in the absence of the dot ionization, irradiation is able to flip the spin of the dot, thus bringing decoherence into the Kondo state and diminishing the Kondo effect. In addition to analytical methods, a number of numerical approaches have been used to study the conductance of a Kondo system out of equilibrium at certain sets of values of the bare parameters of the system. Because of the large number of the parameters involved, the results of such calculations are hard to analyze. At any rate, such a limited consideration could not reveal the universalities of the problem. In our view, it also does not provide an insight into the regimes which do not allow for a perturbative treatment.
In our earlier paper we have learned how to apply the Renormalization Group (RG) technique for a Kondo system out of equilibrium. It allowed us to sum up the infinite series of the perturbation theory in the dot lead coupling. This treatment, valid at relatively high decoherence rates, yielded the expression for the conductance as a function of only one parameter, which is the ratio of the decoherence rate $`\mathrm{}/\tau `$ to the Kondo temperature. Thus we established that the Kondo temperature $`T_K`$ remains a meaningful parameter of the theory, even if the ac field strongly suppresses the manifestations of the Kondo effect. We do not see a way, however, of measuring the decoherence rate directly. Moreover, the definition of this quantity in Ref. makes sense only at a sufficiently high frequency of the ac field, $`\mathrm{}\omega T_K`$. Therefore, the dimensionless ratio $`\mathrm{}/\tau T_K`$ can not be the only parameter describing the effect of ac field on the Kondo system.
In this paper, we find the correct dimensionless variables that characterize the amplitude and frequency of the ac field perturbing the Kondo system. If the ac field biases the dot, then the proper variables are $`eV_{\mathrm{ac}}/T_K`$ and $`\mathrm{}\omega /T_K`$, where $`V_{\mathrm{ac}}`$ is the amplitude of the ac bias. In the case of ac modulation of the gate voltage, the perturbation introduced by the ac field is characterized by the corresponding dimensionless variation of the Kondo temperature $`\delta T_K/T_K`$, which in principle can be independently measured. In terms of the proper pair of variables, the behavior of the Kondo conductance is universal. We find analytically the asymptotes of the universal dependence by further developing the RG treatment (valid in the case of strong suppression of the Kondo effect), and by generalizing the Nozières’ Fermi-liquid theory onto the non-equilibrium case (which adequately describes the limit of weak perturbation by a low-frequency field).
Within this picture, we are able to describe in a consistent way the effect of irradiation in a wide range of frequencies of the ac perturbation - from zero to the dot ionization threshold; this includes the most interesting for the current experiments region $`\mathrm{}\omega T_K`$. The remarkable feature of the presented description is that the Kondo temperature remains the only relevant energy scale, despite the essentially non-equilibrium character of the problem.
The outline of our paper is as follows. In Sec. II we introduce the description of the system by means of the time-dependent Kondo Hamiltonian. Then we derive the expression for the Kondo conductance of the dot in the absence of ac field. The purpose of this derivation is to present the formalism which later would be suitable to describe the non-equilibrium states produces by the irradiation.
In Sections IIIVI we consider the effect of ac modulation of the gate voltage on the Kondo conductance. At higher frequencies, $`\mathrm{}\omega /T_K1`$, the decoherence of the Kondo many-body state is the principal channel via which the ac field influences the Kondo anomaly. This frequency domain is considered in Sections IIIIV. Depending on the strength of modulation, $`\delta T_K/T_K`$, the suppression of the Kondo conductance is significant (Section III), or relatively weak (Sections IV). In Sec. V we consider the limit of very small frequencies of the ac field. The decoherence probability in these conditions is exponentially small. However, the adiabatic evolution of the collective state, caused by the ac modulation, affects its scattering properties. It results in small deviation of the conductance from the unitary limit. The results of Sections IIIV are summarized in Sec. VI.
In Sec. VII we consider the effect of ac bias on the Kondo conductance. It turns out that a strong suppression of the Kondo effect is possible only if $`eV_{\mathrm{ac}}/T_K1`$. Upon the increase of the frequency, the suppression diminishes. This frequency dependence is opposite to the one in the case of gate voltage modulation.
In Sec. VIII we consider the satellite peaks which can be created by irradiation in the $`I`$$`V`$ characteristic of a quantum dot.
Finally in Sec. IX we compare experimental results with our theory.
## II The Kondo effect in a quantum dot
### A The model
The system we study is a quantum dot attached to two leads by high-resistance junctions, so that the charge of the dot is nearly quantized. The Kondo effect emerges in a quantum dot occupied by an odd number of electrons at temperatures below the mean level spacing in the dot. Under such conditions, the topmost occupied level is special, since it is filled by only one electron. It is this level which produces the Kondo effect. The other levels, occupied by two or no electrons are unimportant in our discussion (similarly to the inner shells of a magnetic impurity in the conventional Kondo effect). Therefore the model of the dot attached to two leads can be truncated to the Anderson single-level impurity model,
$`\widehat{H}`$ $`=`$ $`{\displaystyle \underset{k,\sigma ,\alpha }{}}(\xi _k+eV_\alpha )c_{k\sigma \alpha }^{}c_{k\sigma \alpha }^{}+{\displaystyle \underset{k,\sigma ,\alpha }{}}(v_\alpha c_{k\sigma \alpha }^{}d_\sigma ^{}+\mathrm{H}.\mathrm{c}.)`$ (2)
$`+`$ $`{\displaystyle \underset{\sigma }{}}(E_d+eV_{\mathrm{dot}}\mathrm{cos}\omega t)d_\sigma ^{}d_\sigma ^{}+Ud_{}^{}d_{}^{}d_{}^{}d_{}^{};`$ (5)
$`V_{L,R}=\pm \frac{1}{2}\left[V_{\mathrm{dc}}+V_{ac}\mathrm{cos}(\omega ^{}t+\varphi _0)\right],`$
$`\mathrm{\Gamma }_\alpha 2\pi \nu |v_\alpha |^2.`$
Here the first two terms correspond to non-interacting electrons in the two leads ($`\alpha =L,R`$), and tunneling of free electrons between the dot and the leads, respectively. The dot is described by the third and fourth terms of the Hamiltonian, $`E_d`$ and $`UE_d`$ are the ionization and the electron addition energy, respectively. The tunneling matrix elements $`v_\alpha `$ are related to the widths $`\mathrm{\Gamma }_\alpha `$ by Eq. (5), where $`\nu `$ is the density of states in a lead. The ac field can be applied to the gate, which is coupled to the dot capacitively, and thus can modulate the energy of the electron localized in the dot with the amplitude $`eV_{\mathrm{dot}}`$. We assume that the leads can be either only dc-biased, or an additional ac bias can be applied.
### B Time-dependent Schrieffer-Wolff transformation
In the present paper we consider the dot in the Kondo regime, $`UE_d,E_d\mathrm{\Gamma }_{L,R}`$. Under such conditions, the number of electrons on the dot is a well-defined quantity. In the limit of infinitely small tunneling, the ground state of the system described by the Hamiltonian (II A) is doubly degenerate due to the spin of the (single) electron which occupies the level $`d`$. The states with two or no electrons on the dot are higher in energy by $`UE_d`$ or $`E_d`$ respectively, and are not important for the low-energy dynamics of the system. In our paper, we study the irradiation effects when the applied fields do not drive the dot out of the Kondo regime,
$$eV_{\mathrm{dc}},eV_{\mathrm{dot}},eV_{\mathrm{ac}}<E_d,UE_d.$$
(6)
Therefore, the excited states with two or no electrons on the dot are to be projected out. This can be achieved by the Schrieffer-Wolff transformation, modified to account for the time dependence of the parameters of the Hamiltonian (II A). In the present subsection we perform this transformation, which finally yields the description of the quantum dot system by means of the Kondo Hamiltonian with time-dependent parameters.
First we move all dependence of the Hamiltonian on the applied voltages $`V_{\mathrm{dot}}`$, $`V_{\mathrm{dc}}`$, $`V_{\mathrm{ac}}`$ to the off-diagonal terms. It is achieved by the unitary transformation
$`𝖴`$ $`=`$ $`\mathrm{exp}\{{\displaystyle \frac{ie}{\mathrm{}}}{\displaystyle ^t}dt^{}[{\displaystyle \underset{k,\sigma ,\alpha }{}}V_\alpha (t^{})c_{k\sigma \alpha }^{}c_{k\sigma \alpha }^{}`$ (8)
$`+V_{\mathrm{dot}}(t^{})d_\sigma ^{}d_\sigma ^{}]\}.`$
After this transformation, the Hamiltonian has the form
$`\widehat{H}^{}`$ $`=`$ $`𝖴\widehat{H}𝖴^{}i\mathrm{}{\displaystyle \frac{𝖴}{t}}𝖴^{}`$ (9)
$`=`$ $`{\displaystyle \underset{k,\sigma ,\alpha }{}}\xi _k^{}c_{k\sigma \alpha }^{}c_{k\sigma \alpha }^{}+{\displaystyle \underset{\sigma }{}}(E_d)d_\sigma ^{}d_\sigma ^{}+Ud_{}^{}d_{}^{}d_{}^{}d_{}^{}`$ (10)
$`+{\displaystyle \underset{k,\sigma ,\alpha }{}}[\stackrel{~}{v}_\alpha ^{}(t)c_{k\sigma \alpha }^{}d_\sigma ^{}+\mathrm{H}.\mathrm{c}.],`$ (11)
with
$$v_\alpha (t)v_\alpha \mathrm{exp}\left\{\frac{ie}{\mathrm{}}^t𝑑t^{}\left[V_\alpha (t^{})V_{\mathrm{dot}}(t^{})\right]\right\}.$$
(12)
Now we can make the time-dependent Schrieffer-Wolff transformation, which is defined by the unitary operator
$$𝖶=\mathrm{exp}(𝖵)$$
(13)
with
$`𝖵={\displaystyle \underset{k,\sigma ,\alpha }{}}\{[w_{k\alpha }^{(1)}(t)(1n_\sigma )+w_{k\alpha }^{(2)}(t)n_\sigma ]d_\sigma ^{}c_{k\sigma \alpha }^{}`$ (14)
$`\mathrm{H}.\mathrm{c}.\}.`$ (15)
The functions $`w_{k\alpha }^{(j)}(t)`$ are to be found from the condition
$$0=\widehat{H}_v+[𝖵,\widehat{H}_1]i\mathrm{}\frac{𝖵}{t},$$
(16)
where $`\widehat{H}_v`$ is the part of the Hamiltonian (11) responsible for mixing of electron states in the leads and in the dot \[the last term of Eq. (11)\], and $`\widehat{H}_1`$ describes the uncoupled dot and leads \[the first three terms in Eq. (11)\]. The condition (16) ensures that the resulting Hamiltonian $`𝖶\widehat{H}^{}𝖶^{}`$ has no linear-in-$`v_\alpha `$ terms, which account for the variations of the number of electrons in the dot.
The only difference of the transformation (13)–(15) from the conventional Schrieffer-Wolff transformation is the time dependence of $`w_{k\alpha }^{(j)}`$. For the static Anderson Hamiltonian, these factors are constant; in our case they are functions of time because of the time variations of the Hamiltonian (II A). Solving Eq. (16) for $`w_{k\alpha }^{(1)}(t)`$, we obtain
$$w_{k\alpha }^{(1)}(t)=\left[i^t𝑑t^{}e^{i(\xi _kE_d)t^{}/\mathrm{}}v_\alpha (t^{})\right]e^{i(\xi _kE_d)t/\mathrm{}}.$$
(17)
When the applied ac fields are slow enough, $`\mathrm{}\omega /E_d,\mathrm{}\omega ^{}/E_d1`$, one can solve Eq. (16) in the adiabatic approximation, neglecting the third term in it. This approach yields a simplified expression for $`w_{k\alpha }^{(1)}(t)`$:
$$w_{k\alpha }^{(1)}(t)\frac{v_\alpha (t)}{E_deV_\alpha (t)+eV_{\mathrm{dot}}(t)}.$$
(18)
Here we have also neglected the single-electron energies $`\xi _k`$ because the Kondo effect is produced by the states close to the Fermi level, whose energy is small in comparison to $`E_d`$. The formulas for $`w_{k\alpha }^{(2)}(t)`$ are analogous to Eqs. (17)–(18), only $`E_d`$ must be replaced by $`E_dU`$.
Applying transformation (13)–(15) to the Hamiltonian (11), we come to the Kondo Hamiltonian
$`\widehat{H}_K`$ $`=`$ $`\widehat{H}_0+\widehat{H}_𝒥,\widehat{H}_0={\displaystyle \underset{k,\sigma ,\alpha }{}}\xi _k^{}c_{k\sigma \alpha }^{}c_{k\sigma \alpha }^{},`$ (20)
$`\widehat{H}_𝒥`$ $`=`$
k,σ,α
k,σ,α
𝒥αα(t)(14δσσ+S^lsσσl)ckσαckσα,subscript
k,σ,α
k,σ,α
subscript𝒥𝛼superscript𝛼𝑡14subscript𝛿𝜎superscript𝜎subscript^𝑆𝑙subscriptsuperscript𝑠𝑙𝜎superscript𝜎subscriptsuperscript𝑐𝑘𝜎𝛼subscriptsuperscript𝑐absentsuperscript𝑘superscript𝜎superscript𝛼\displaystyle\!\sum_{\parbox{21.68121pt}{$k,\sigma,\alpha$\\
$k^{\prime}\!,\sigma^{\prime}\!\!,\alpha^{\prime}$}}\!{\cal J}_{\alpha\alpha^{\prime}}(t)\left(\textstyle\frac{1}{4}\delta_{\sigma\sigma^{\prime}}+\hat{S}_{l}s^{l}_{\sigma\sigma^{\prime}}\right)c^{\dagger}_{k\sigma\alpha}c^{\phantom{\dagger}}_{k^{\prime}\sigma^{\prime}\alpha^{\prime}}\;, (21)
where $`\widehat{𝐬}`$ and $`\widehat{𝐒}`$ are the spin operators of the electrons in the leads and of the electron on the isolated level, respectively; we assume summation over the repeating indices $`l=x,y,z`$. In the framework of the Hamiltonian (II B), the state of the dot is fully described by its spin. The terms of the Anderson Hamiltonian (II A) that are responsible for the electron tunneling to/from the dot, and for the Coulomb interaction of electrons in the dot have been transformed to the term $`\widehat{H}_𝒥`$ of the Kondo Hamiltonian (II B). This term represents exchange interaction between the spin of the dot and the electrons in the leads. The Hamiltonian (II B) operates within the band $`E_d<\xi _k<UE_d`$, see Ref. . The coupling parameters $`𝒥_{\alpha \alpha ^{}}(t)`$ are given by
$$𝒥_{\alpha \alpha ^{}}(t)=\stackrel{~}{v}_\alpha (t)\left[w_{k^{}\alpha ^{}}^{(1)}(t)+w_{k^{}\alpha ^{}}^{(2)}(t)\right]^{}+\text{C.c.}$$
(22)
The applied bias is accounted for by the time dependence of $`𝒥_{\alpha \alpha ^{}}(t)`$ with $`\alpha \alpha ^{}`$. The dependence of the right-hand side of Eq. (22) on the indices $`k`$ is negligible near the Fermi level; therefore we can disregard this dependence in the coupling constants $`𝒥`$.
In this paper, we are primarily interested in the irradiation effects when the applied fields are unable to ionize the dot,
$`eV_{\mathrm{dc}},eV_{\mathrm{dot}},eV_{\mathrm{ac}}E_d,UE_d,`$
$`\mathrm{}\omega ,\mathrm{}\omega ^{}E_d,UE_d.`$
Under these conditions, one can use the approximate solution (18) of the equation (16), expanding it in powers of small parameters $`eV_{\mathrm{dot}}/E_d`$, $`eV_{\mathrm{dot}}/(UE_d)`$, etc. For simplicity we will consider the cases when the system is affected by only one kind of ac field: either ac voltage applied to the gate, or the ac bias.
In the former case, $`eV_{\mathrm{dot}}0`$, $`eV_{\mathrm{ac}}=0`$ , we obtain the following expression for the coupling parameters:
$$𝒥_{\alpha \alpha ^{}}(t)=𝒥_{\alpha \alpha ^{}}^{(0)}\left[1+\gamma \mathrm{cos}\omega t\right]\mathrm{exp}\left[\frac{ie}{\mathrm{}}(V_{\mathrm{dc},\alpha }V_{\mathrm{dc},\alpha ^{}})t\right],$$
(23)
where the exchange constants $`𝒥_{\alpha \alpha ^{}}^{(0)}`$ are given by
$$𝒥_{\alpha \alpha ^{}}^{(0)}\frac{\sqrt{\mathrm{\Gamma }_\alpha \mathrm{\Gamma }_\alpha ^{}}}{\pi \nu \stackrel{~}{E}_d},\stackrel{~}{E}_d\frac{(UE_d)E_d}{U}.$$
(24)
The exponential factor in Eq. (23) is due to the dc bias, which produces the phase difference between the electrons in the left and right lead. The cosine term accounts for the applied ac field and stems from the adiabatic variation of the electron energy in the dot, $`E_d+eV_{\mathrm{dot}}(t)`$, see Eq. (18). The strength of the applied ac field is characterized by the dimensionless parameter
$$\gamma eV_{\mathrm{dot}}\frac{2E_dU}{(UE_d)E_d}1.$$
(25)
If the ac field is applied to the leads rather than to the gate, $`eV_{\mathrm{dot}}=0`$, $`eV_{\mathrm{ac}}0`$, the expressions for $`𝒥_{\alpha \alpha ^{}}(t)`$ read
$`𝒥_{\alpha \alpha }(t)`$ $`=`$ $`𝒥_{\alpha \alpha }^{(0)},`$ (26)
$`𝒥_{LR}(t)`$ $`=`$ $`𝒥_{LR}^{(0)}\mathrm{exp}\left[{\displaystyle \frac{ieV_{\mathrm{dc}}t}{\mathrm{}}}+i\gamma ^{}\mathrm{sin}(\omega ^{}t+\varphi _0)\right],`$ (27)
where $`𝒥_{\alpha \alpha ^{}}^{(0)}`$ is given by Eq. (24). The ac bias creates the phase difference between the electrons in the left and right lead, and therefore enters the exponent in (27) together with the dc bias. The relevant parameter characterizing the strength of the ac perturbation here is
$$\gamma ^{}\frac{eV_{\mathrm{ac}}}{\mathrm{}\omega ^{}}.$$
(28)
The variation of the electron energy in the dot with respect to the leads, $`E_d\pm eV_{\mathrm{ac}}(t)`$, see Eq. (18), generates a term smaller by a factor of $`\mathrm{}\omega ^{}/\stackrel{~}{E}_d`$, and is neglected in Eq. (27).
In the limit of small amplitude of ac bias, $`\gamma ^{}1`$, the expression (27) for $`𝒥_{\alpha \alpha ^{}}(t)`$ may be further simplified by dropping the terms of high orders in $`\gamma ^{}`$. Expanding the right-hand side of (27) in powers of $`\gamma ^{}`$ up to the first power, we arrive at
$`𝒥_{\alpha \alpha }(t)`$ $`=`$ $`𝒥_{\alpha \alpha }^{(0)},`$ (29)
$`𝒥_{LR}(t)`$ $`=`$ $`𝒥_{LR}^{(0)}\mathrm{exp}\left[{\displaystyle \frac{ieV_{\mathrm{dc}}t}{\mathrm{}}}\right]\left[1+i\gamma ^{}\mathrm{sin}(\omega ^{}t+\varphi _0)\right].`$ (30)
### C Kondo conductance in equilibrium
In the framework of the Kondo Hamiltonian (II B)-(22), two types of tunneling between the left and right leads are possible: regular elastic cotunneling \[the first term in parentheses in Eq. (21)\], and “exchange cotunneling” \[the second term\]. In an act of “exchange cotunneling,” the simultaneous flip of the spins of the tunneling electron and of the dot can occur. In the case of weak coupling ($`\nu |𝒥_{\alpha \alpha ^{}}^{(0)}|u1`$), one may apply the perturbation theory to evaluate the conductance through the dot. It turns out that at $`T0`$, the higher-order terms of the perturbation theory series grow, finally making the series diverge, signaling the Kondo anomaly. This phenomenon was extensively studied for the magnetic impurities in metals. In the subsection II C we demonstrate how a similar behavior emerges in the tunneling through a quantum dot. The main purpose of this subsection is to present the formalism which is suitable for the treatment of a non-equilibrium case at hand. For simplicity, we first consider the case of no ac field. Effects of the ac field are included into consideration in subsequent sections.
Unlike the conventional treatment of the Kondo problem, we have to consider the Kondo anomaly directly in the conductance, rather than in the scattering amplitude. This need emerges from the kinetic nature of the problem at $`\gamma ,\gamma ^{}0`$. To calculate the differential dc conductance $`G(V_{\mathrm{dc}})`$, we employ the non-equilibrium Keldysh technique in the time representation. In this formalism
$$G(V_{\mathrm{dc}})=\frac{}{V_{\mathrm{dc}}}𝖲(\mathrm{},0)\widehat{I}(0)𝖲(0,\mathrm{})_0,$$
(31)
where
$`\widehat{I}(t)`$ $`=`$ ie
k,σ
k,σ
[𝒥LR(t)(14δσσ+S^lsσσl)ckσLckσR\displaystyle\frac{ie}{\hbar}\sum_{\parbox{14.45377pt}{$k$,$\sigma$\\
$k^{\prime}\!,\sigma^{\prime}$}}\left[{\cal J}_{LR}(t)\left(\textstyle\frac{1}{4}\delta_{\sigma\sigma^{\prime}}+\hat{S}_{l}s^{l}_{\sigma\sigma^{\prime}}\right)c^{\dagger}_{k\sigma L}c^{\phantom{\dagger}}_{k^{\prime}\sigma^{\prime}R}\right. (33)
$`\mathrm{H}.\mathrm{c}.]`$
is the current operator, and $`𝖲(t_2,t_1)`$ is the evolution matrix determined by $`\widehat{H}_𝒥`$.
In the lowest non-vanishing (second) order of the perturbation theory in the coupling constant $`𝒥_{\alpha \alpha ^{}}^{(0)}`$, the conductance of the dot is given by the expression
$$G^{(2)}=\pi ^2\frac{e^2}{\pi \mathrm{}}\nu ^2\left[𝒥_{LR}^{(0)}\right]^2.$$
(34)
The logarithmic divergences appear starting from the terms of the third order in $`𝒥_{\alpha \alpha ^{}}^{(0)}`$. A representative term has the following structure:
$`{\displaystyle \frac{e^2}{\pi \mathrm{}}}{\displaystyle \frac{\left[𝒥_{LR}^{(0)}\right]^2𝒥_{RR}^{(0)}}{\mathrm{}^3}}{\displaystyle _{\mathrm{}}^0}𝑑t_1{\displaystyle _{t_1}^0}𝑑t_2\widehat{S}_j(0)\widehat{S}_k(t_1)\widehat{S}_l(t_2)\epsilon ^{jkl}`$ (35)
$`\times \left[t_1\mathrm{cos}(eVt_1/\mathrm{})+t_2\mathrm{cos}(eVt_2/\mathrm{})\right]`$ (36)
$`\times {\displaystyle \underset{k_1,k_2,k_3}{}}G_{k_1}(t_2)G_{k_2}(t_2t_1)\overline{G}_{k_3}(t_1),`$ (37)
$`\epsilon ^{jkl}`$ is the antisymmetric unit tensor and $`G_k(t)`$ and $`\overline{G}_k(t)`$ are the time-ordered and anti-time-ordered Green functions of free electrons in the leads, given by
$$G_k(t)=\{\begin{array}{cc}i[1f(\xi _k)],\hfill & \text{if}t>0,\hfill \\ if(\xi _k)],\hfill & \text{if}t<0,\hfill \end{array}$$
(38)
with $`f(\xi )`$ being the Fermi distribution function. This and other terms of the same structure yield the Kondo divergence in the conductance.
If there is no external ac field, the averages $`\widehat{S}_j(t_1)\widehat{S}_k(t_2)\widehat{S}_l(t_3)`$ are independent of time and equal $`(i/4)\epsilon _{jkl}`$. After adding up all the cubic in $`𝒥_{\alpha \alpha ^{}}^{(0)}`$ terms in the expression for the conductance $`G`$ \[one of them is given by Eq. (37)\], summing over the electron states $`k_i`$, and performing the integration over $`t_2`$ \[see Eq. (37)\], we arrive at
$`G^{(3)}(T,V_{\mathrm{dc}})={\displaystyle \frac{3\pi ^2}{2}}{\displaystyle \frac{e^2}{\pi \mathrm{}}}\nu ^3\left[𝒥_{LR}^{(0)}\right]^2\left[𝒥_{RR}^{(0)}+𝒥_{LL}^{(0)}\right]`$ (39)
$`\times {\displaystyle _{\mathrm{}}^0}dt{\displaystyle \frac{(t)\mathrm{cos}(eV_{\mathrm{dc}}t/\mathrm{})}{\mathrm{sinh}^2(\pi Tt/\mathrm{})+(T/D_0)^2}}\left({\displaystyle \frac{\pi T}{\mathrm{}}}\right)^2`$ (40)
Here
$$D_0\sqrt{E_d(UE_d)}$$
(41)
is the effective bandwidth. For the sake of simplicity, further we will mostly consider the zero-bias conductance $`G_{\mathrm{peak}}`$. In this case, Eq. (40) yields
$$G_{\mathrm{peak}}^{(3)}(T)=\frac{3\pi ^2}{2}\frac{e^2}{\pi \mathrm{}}\nu ^3\left[𝒥_{LR}^{(0)}\right]^2\left[𝒥_{RR}^{(0)}+𝒥_{LL}^{(0)}\right]\mathrm{ln}\frac{D_0}{T}.$$
(42)
The results for the finite-bias conductance $`G^{(3)}(V_{\mathrm{dc}})`$ with $`eV_{\mathrm{dc}}>T`$ can be obtained from Eq. (42) by replacing $`T`$ with $`eV_{\mathrm{dc}}`$.
Thus the second \[Eq. (34)\] and third \[Eq. (42)\] orders of the perturbation theory in the coupling constant $`𝒥_{\alpha \alpha ^{}}^{(0)}`$ yield the following expression for the dot conductance
$`G_{\mathrm{peak}}`$ $`=`$ $`{\displaystyle \frac{3\pi ^2}{4}}{\displaystyle \frac{e^2}{\pi \mathrm{}}}\nu ^2\left[𝒥_{LR}^{(0)}\right]^2\left[1+2\nu \left(𝒥_{RR}^{(0)}+𝒥_{LL}^{(0)}\right)\mathrm{ln}{\displaystyle \frac{D_0}{T}}\right]`$ (43)
$`+`$ $`{\displaystyle \frac{\pi ^2}{4}}{\displaystyle \frac{e^2}{\pi \mathrm{}}}\nu ^2\left[𝒥_{LR}^{(0)}\right]^2.`$ (44)
Here we have split the quadratic in $`𝒥_{\alpha \alpha ^{}}^{(0)}`$ contribution (34) in two: the one due to the “exchange cotunneling”, which entered the first term in Eq. (44), and the one due to regular cotunneling, which became the last term in Eq. (44). The cubic in $`𝒥_{\alpha \alpha ^{}}^{(0)}`$ term in Eq. (44) grows as the temperature is being lowered, demonstrating the Kondo anomaly. The regular cotunneling does not produce terms growing at low temperatures and bias, and does not contribute to the Kondo effect. Equation (44) is valid while $`TT_KD_0\mathrm{exp}[1/\nu (𝒥_{LL}^{(0)}+𝒥_{RR}^{(0)})]`$.
If this condition is not satisfied, then the expansion up to the cubic in $`𝒥_{\alpha \alpha ^{}}^{(0)}`$ terms is insufficient. At $`TT_K`$, the conductance can be derived in the leading logarithmic approximation. The latter consists in the summation of the most diverging terms in each order in $`𝒥_{\alpha \alpha ^{}}^{(0)}`$, i.e. terms, proportional to $`[𝒥_{LR}^{(0)}]^2[𝒥_{\alpha \alpha ^{}}^{(0)}\mathrm{ln}(D_0/T)]^n`$, in the series for $`G`$. To perform this summation, we modify the “poor man’s scaling” technique. In the framework of this technique, the electron bandwidth $`D`$ is gradually reduced, and the exchange constants in the Kondo Hamiltonian (II B) are renormalized to compensate for this band reduction, i.e. $`𝒥_{\alpha \alpha ^{}}^{(0)}`$ is replaced with some $`𝒥_{\alpha \alpha ^{}}(D)`$ . The proper dependence of $`𝒥_{\alpha \alpha ^{}}`$ on $`D`$ should be derived from the condition of invariance of physical quantities with respect to the RG transformation. Finally, the renormalized Hamiltonian with the reduced band width will allow for the calculation of the conductance in the second order of the perturbation theory in the renormalized exchange constants $`𝒥_{\alpha \alpha ^{}}`$; the resulting expression will be equal to the sum of the dominant terms of all orders of the perturbation theory in the initial, bare exchange constants $`𝒥_{\alpha \alpha ^{}}^{(0)}(D=D_0)`$.
For the non-equilibrium system we consider, the Renormalization Group (RG) equations for the exchange constants should be derived from the condition of the invariance of the linear conductance (or current) under the RG transformation, rather than the invariance of the scattering amplitudes. In the main logarithmic approximation which we are going to employ, the (invariant) conductance must be evaluated in the two lowest non-vanishing orders of the perturbation theory, namely in the second and third ones, see Eqs. (34), (40). The Kondo divergence (and, therefore, the renormalization of $`𝒥_{\alpha \alpha ^{}}`$) occur due to exchange scattering \[the second term in braces in Eq. (21)\] only. Therefore we single out this contribution in the term of the second order in $`𝒥_{\alpha \alpha ^{}}`$,
$$G_{\mathrm{exch}}^{(2)}(D)=\frac{3\pi ^2}{4}\frac{e^2}{\pi \mathrm{}}\nu ^2\left[𝒥_{LR}^{(0)}(D)\right]^2.$$
(45)
In the third order in the exchange constants, the conductance is given by Eq. (40). The resulting condition of invariance of $`G`$ under the transformation, which corresponds to the “poor man‘s scaling”, has the following form:
$`{\displaystyle \frac{}{D}}`$ $`\{G_{\mathrm{exch}}^{(2)}(D)`$ (47)
$`+{\displaystyle \frac{3\pi ^2}{2}}{\displaystyle \frac{e^2}{\pi \mathrm{}}}\nu ^3\left[𝒥_{LR}\right]^2[𝒥_{RR}+𝒥_{LL}]\mathrm{ln}{\displaystyle \frac{D}{T}}\}=0.`$
Within the accuracy of this equation, when differentiating the second term, we should neglect any implicit dependence on $`D`$ through the parameters $`𝒥_{\alpha \alpha ^{}}(D)`$.
Equation (47), together with Eq. (45), yields the equation for the evolution of $`𝒥_{LR}`$:
$$\frac{d𝒥_{LR}}{dD}=\nu \frac{𝒥_{LR}(𝒥_{RR}+𝒥_{LL})}{D},$$
(48)
The corresponding equations for $`𝒥_{RR}`$ and $`𝒥_{LL}`$ can be derived from the condition of invariance of other physical quantities under the RG transformation. For this purpose, we pick the spin current from the left and right lead,
$`I_\alpha ^{(s)}`$ $`=`$ $`𝖲(\mathrm{},0)\widehat{I}_\alpha ^{(s)}(0)𝖲(0,\mathrm{})_0,`$ (49)
$`\widehat{I}_\alpha ^{(s)}(t)`$ $``$ $`i[\widehat{H}_𝒥,{\displaystyle \underset{k}{}}\left(c_{k\alpha }^{}c_{k\alpha }^{}c_{k\alpha }^{}c_{k\alpha }^{}\right)],`$ (50)
which is induced by applying infinitely small magnetic field to the leads. The resulting two equations will be independent, in contrast to the corresponding equations for the charge, because the spin of the dot can vary while the charge cannot.
Evaluating $`I_\alpha ^{(s)}`$ in the second and third orders of the perturbation theory in $`𝒥_{\alpha \alpha ^{}}^{(0)}`$, similarly to Eq. (44), and differentiating it by $`D`$, we arrive at
$`{\displaystyle \frac{d𝒥_{RR}}{dD}}=\nu {\displaystyle \frac{𝒥_{RR}^2+𝒥_{LR}^2}{D}},`$ (51)
$`{\displaystyle \frac{d𝒥_{LL}}{dD}}=\nu {\displaystyle \frac{𝒥_{LL}^2+𝒥_{LR}^2}{D}},`$ (52)
Equations (48), (51), and (52) make a complete system, which, with the initial conditions
$$𝒥_{\alpha \alpha ^{}}(D_0)=𝒥_{\alpha \alpha ^{}}^{(0)}\frac{\sqrt{\mathrm{\Gamma }_\alpha \mathrm{\Gamma }_\alpha ^{}}}{\pi \nu \stackrel{~}{E}_d},$$
(53)
see Eq. (24), yields
$$𝒥_{LR}(D)=\frac{2\sqrt{\mathrm{\Gamma }_L\mathrm{\Gamma }_R}}{\mathrm{\Gamma }_L+\mathrm{\Gamma }_R}\frac{1}{2\nu \mathrm{ln}(D/T_K)}.$$
(54)
The Kondo temperature $`T_K`$ here is given by
$$T_K=\mu \sqrt{\frac{(\mathrm{\Gamma }_L+\mathrm{\Gamma }_R)U}{\pi }}\mathrm{exp}\left[\frac{\pi \stackrel{~}{E}_d}{(\mathrm{\Gamma }_L+\mathrm{\Gamma }_R)}\right],$$
(55)
with $`\mu 1`$. To obtain the pre-exponential factor $`\sqrt{(\mathrm{\Gamma }_L+\mathrm{\Gamma }_R)U/\pi }`$ in the equation for $`T_K`$ one has in fact to include the next in $`𝒥_{\alpha \alpha ^{}}(D)`$ order in the RG equations, see Ref. .
The renormalization should proceed until the band width is reduced to $`T`$. After that, the current and conductance can be calculated in the Born approximation (45) in the renormalized exchange constant $`𝒥_{LR}`$ given by Eq. (54) with $`D=T`$. The resulting expression for the conductance in the domain $`TT_K`$ is
$$G_{\mathrm{peak}}=\frac{3\pi ^2}{16}\frac{1}{[\mathrm{ln}(T/T_K)]^2}G_U,$$
(56)
where
$$G_U\frac{e^2}{\pi \mathrm{}}\frac{4\mathrm{\Gamma }_L\mathrm{\Gamma }_R}{(\mathrm{\Gamma }_L+\mathrm{\Gamma }_R)^2}$$
(57)
is the conductance of the dot in the unitary limit of tunneling.
At $`TT_K`$, one can expand Eq. (56) into the series of powers of $`𝒥_{\alpha \alpha ^{}}^{(0)}\mathrm{ln}(D_0/T)`$. The first term of the series is the conductance calculated in the Born approximation \[see Eq. (45)\], the second term yields the lowest order Kondo correction given by Eq. (42).
The RG technique can be also used to derive the dependence of Kondo conductance on the applied dc bias in the domain $`eV_{\mathrm{dc}}T_K`$, $`eV_{\mathrm{dc}}>T`$. Starting from Eqs. (40) and (45) and proceeding along the lines of Eqs. (47)–(54), we arrive at
$$G(V_{\mathrm{dc}})=\frac{3\pi ^2}{16}\frac{1}{[\mathrm{ln}(eV_{\mathrm{dc}}/T_K)]^2}G_U.$$
(58)
Thus the Renormalization Group technique (47)–(56) allows one to perform summation of infinite series of the perturbation theory in the exchange constants $`𝒥_{\alpha \alpha ^{}}^{(0)}`$. The results obtained in this way are valid in a wider domain of parameters as compared to the results of the finite-order perturbation theory. The RG technique reveals the meaning of the energy scale $`T_K`$. The resulting expressions (56), (58) for physical quantities contain the single relevant characteristic of the system, $`T_K`$, rather than numerous parameters of the Anderson Hamiltonian \[Eq. (5)\]. For example, in Eq. (56) the dependence of the differential conductance on the applied bias is expressed in terms of the dimensionless variable $`T/T_K`$. The dependence of $`G/G_U`$ on this variable is given by some universal function of any value of $`T/T_K`$; its high-temperature asymptote (56) is established with the help of RG technique. Similarly, the frequency and magnitude of the applied ac field may enter into some new universal formulas for $`G/G_U`$ in the form of dimensionless variables, being normalized by $`T_K`$. The generalization of the RG technique which we presented in this section, will allow us to check the validity of this conjecture and to establish the asymptotes of these new universal dependences.
## III Spin decoherence by ac gate voltage
Now we include into consideration the effects of an ac field. As we have shown in our earlier paper, the ac field can bring decoherence into the dynamics of the dot’s spin, thus affecting the Kondo conductance. We start our study of the irradiation-modified Kondo anomaly from the consideration of the decoherence.
### A Mechanisms of spin decoherence
In terms of the Anderson Hamiltonian (5) the loss of coherence by the dot spin occurs when an electron leaves the dot and another electron, with the opposite spin, enters it. If the frequency of the applied ac field is large enough, $`\mathrm{}\omega >E_d,UE_d`$, this process can consist of two real processes: the dot gets ionized by the ac field, and then an electron from a lead enters the dot to fill the vacancy. Alternatively, an extra electron can be put in the dot, and then the electron which was initially present in the dot leaves it.
In the present paper we deal with a more subtle case, when the applied ac field is unable to ionize the dot. In this case the dot can still change its spin, even at zero bias, by means of the “spin-flip cotunneling,” which is shown schematically on Fig. 1. In the course of this process, an electron, which interacts with the dot spin \[see Eq. (II B)\], absorbs a photon and hops to a state above the Fermi level, while the spin of the dot flips. In terms in the Anderson Hamiltonian (5), this process cannot be described as two separate real processes. Instead, the change of the dot spin occurs as a single process when a state with two or no electrons in the dot appears only as a virtual intermediate state.
The rate of the spin-flip cotunneling can be calculated with the help of the Kondo Hamiltonian given by Eqs. (II B), (22), and (18). In the case of weak modulation, $`\gamma 1`$, see Eq. (27), it is sufficient to account for the single-photon processes only, and use the reduced form of the Hamiltonian, given by Eqs. (II B), (23)-(25). An infinitely small dc bias, needed for measurements of the linear conductance, does not affect the rate of spin-flip cotunneling. Therefore in this subsection we set $`V_{\mathrm{dc}}=0`$ for the sake of simplicity.
Applying the Fermi Golden Rule we obtain
$$\frac{\mathrm{}}{\tau }=\frac{1}{8\pi }\mathrm{}\omega \left[\frac{\mathrm{\Gamma }_L+\mathrm{\Gamma }_R}{\stackrel{~}{E}_d}\right]^2\gamma ^2,$$
(59)
where $`\gamma `$ is given by Eq. (25).
The amplitude of inelastic transitions yielding Eq. (59) was evaluated in the lowest-order perturbation theory. It corresponds to the first order in the amplitude of the ac perturbation, and the zeroth order in the time-independent (at $`V_{\mathrm{dc}}=0`$) part of the exchange interaction (21). Accounting for the terms of higher-order in this time-independent part renormalizes the amplitude of the inelastic transition (which is still linear in the amplitude of the ac field). Similarly to the calculation of the conductance, we intend to collect the leading logarithmic terms in the renormalization of the inelastic transition amplitude. This can be accomplished by the renormalization group (RG) transformation, described in subsection II C. The transformation reduces the electron band width $`D`$ and renormalizes the matrix elements $`𝒥_{\alpha \alpha ^{}}`$ of the Kondo Hamiltonian (II B) to account for this band reduction. Finally one can calculate the decoherence rate in the second order perturbation theory in renormalized $`𝒥_{\alpha \alpha ^{}}`$; the result given by such a treatment equals the sum of infinite perturbation theory series in the initial Hamiltonian.
The RG transformation starts from the bandwidth $`D=D_0`$ given by Eq. (41), and the initial matrix elements
$$𝒥_{\alpha \alpha ^{}}(D)|_{D=D_0}=𝒥_{\alpha \alpha ^{}}^{(0)}\left[1+\gamma \mathrm{cos}\omega t\right],$$
(60)
cf. Eq. (23). While the width $`D`$ of the band exceeds $`\mathrm{}\omega `$, the time dependence of the Hamiltonian matrix elements (60) can be treated adiabatically, i.e. time $`t`$ in the right-hand side of Eq. (60) can be considered as just a parameter. The RG equations, derived from the condition of invariance of physical quantities under the transformation, have the same form (48), (51), and (52). The transformation must be stopped when the bandwidth is reduced to the values of the order of the frequency $`\mathrm{}\omega `$ of the applied ac field. Expanding the solution of the RG equations (48), (51), and (52) with the initial condition (60) in powers of $`\gamma `$ up to the first power, we obtain
$`𝒥_{\alpha \alpha ^{}}(D)|_{D\mathrm{}\omega }`$ $`=`$ $`{\displaystyle \frac{2\sqrt{\mathrm{\Gamma }_\alpha \mathrm{\Gamma }_\alpha ^{}}}{\mathrm{\Gamma }_\alpha +\mathrm{\Gamma }_\alpha ^{}}}{\displaystyle \frac{1}{2\nu \mathrm{ln}(\mathrm{}\omega /T_K)}}`$ (61)
$`\times `$ $`\left[1+\gamma {\displaystyle \frac{\pi \stackrel{~}{E}_d}{\mathrm{\Gamma }_L+\mathrm{\Gamma }_R}}{\displaystyle \frac{1}{\mathrm{ln}(\mathrm{}\omega /T_K)}}\mathrm{cos}\omega t\right].`$ (62)
The Fermi Golden rule applied to the Hamiltonian (II B) with $`𝒥_{\alpha \alpha ^{}}`$ given by Eq. (62) yields the following expression for the decoherence rate:
$$\frac{\mathrm{}}{\tau T_K}=\frac{3\pi }{32}\frac{\mathrm{}\omega }{T_K}\frac{1}{\left[\mathrm{ln}(\mathrm{}\omega /T_K)\right]^4}\left[\frac{\delta T_K}{T_K}\right]^2.$$
(63)
Here we introduced the relative amplitude
$$\frac{\delta T_K}{T_K}\gamma \frac{\pi \stackrel{~}{E}_d}{\mathrm{\Gamma }_L+\mathrm{\Gamma }_R}$$
(64)
of adiabatic variations of the “time-dependent Kondo temperature”. The latter is defined by
$$T_K(t)\mu \sqrt{\frac{(\mathrm{\Gamma }_L+\mathrm{\Gamma }_R)U}{\pi }}\mathrm{exp}\left[\frac{\pi \stackrel{~}{E}_d(t)}{(\mathrm{\Gamma }_L+\mathrm{\Gamma }_R)}\right],$$
(65)
with
$$\stackrel{~}{E}_d(t)=\stackrel{~}{E}_d(1+\gamma \mathrm{cos}\omega t),$$
(66)
cf. Eqs. (55) and (24).
One can see that the amplitude of the ac field enters Eq. (63) through the dimensionless parameter $`\delta T_K/T_K`$. The value of this parameter, in principle, can be directly measured. Representation of $`\mathrm{}/\tau T_K`$ in terms of $`\delta T_K/T_K`$ will allow us to build a universal description of the ac field effect on the Kondo conductance.
As we mentioned before, the spin-flip cotunneling is essentially different from the dot ionization with its subsequent filling. During a process of spin-flip cotunneling, the ionized dot exists only as a virtual state. Therefore, the spin-flip cotunneling persists at frequencies $`\mathrm{min}\{E_d,UE_d\}>\mathrm{}\omega `$, leading to the decoherence of the dot spin state without ionization.
### B Effects of spin decoherence on Kondo conductance
As we have just shown, the external ac field is able to flip the dot’s spin. Therefore, in the presence of the ac field, the averages of the type $`S_j(t_1)S_k(t_2)S_l(t_3)`$ no longer equal $`S_j(0)S_k(0)S_l(0)=TrS_jS_kS_l(i/4)\epsilon _{jkl}`$. In the limiting case $`|t_mt_n|\tau `$ ($`mn`$), the orientations of the dot spin at $`t=t_1`$, $`t_2`$, $`t_3`$ are independent of each other, because of the spin-flip cotunneling, and one has
$$S_j(t_1)S_k(t_2)S_l(t_3)=S_j(t_1)S_k(t_2)S_l(t_3)=0.$$
At finite time intervals $`|t_mt_n|`$, the spin correlator decays exponentially, with the spin-flip cotunneling rate being the characteristic decay rate:
$`\widehat{S}_j(t_1)\widehat{S}_k(t_2)\widehat{S}_l(t_3)=(i/4)\epsilon _{jkl}\mathrm{exp}(t_{max}/\tau ),`$ (67)
$`t_{max}\mathrm{max}\{|t_1t_2|,|t_2t_3|,|t_1t_3|\}.`$ (68)
Equation (67) can be derived using the formalism of equations of motion. In the framework of this formalism, Eq. (67) appears as the solution of the equation
$`{\displaystyle \frac{}{t_1}}\widehat{S}_j(t_1)\widehat{S}_k(t_2)\widehat{S}_l(t_3)`$ (69)
$`=𝖲(t_3,t_1)\left(i[\widehat{H}_𝒥,\widehat{S}_j]\right)𝖲(t_1,t_2)\widehat{S}_k𝖲(t_2,t_3)\widehat{S}_l_0`$ (70)
where $`𝖲(t,t^{})`$ is the evolution matrix determined by $`\widehat{H}_𝒥`$. Expanding $`𝖲(t,t^{})`$ in powers of $`\gamma 𝒥_{\alpha \alpha ^{}}^{(0)}`$ up to the first power, we arrive at
$`{\displaystyle \frac{}{t_1}}\widehat{S}_j(t_1)\widehat{S}_k(t_2)\widehat{S}_l(t_3)`$ (71)
$`={\displaystyle \frac{1}{\tau }}\left[\theta (t_3t_1)\theta (t_2t_1)\theta (t_1t_3)\theta (t_1t_2)\right]`$ (72)
$`\times \widehat{S}_j(t_1)\widehat{S}_k(t_2)\widehat{S}_l(t_3),`$ (73)
where $`\tau `$ is given by Eq. (59). Equation (73) with $`\tau `$ given by Eq. (63) can be obtained by expanding the evolution matrix $`𝖲(t,t^{})`$ up to the second power in $`\gamma 𝒥_{\alpha \alpha ^{}}^{(0)}`$, and using the RG technique described in Sec. II C.
The leading effect of the irradiation is in cutting off the logarithmic divergences in the time integrals of type (37). One can easily see that with the time-decaying spin correlation function (67), $`G_{\mathrm{peak}}^{(3)}`$ is finite even at $`T0`$:
$$G^{(3)}=\frac{3\pi ^2}{2}\frac{e^2}{\pi \mathrm{}}\nu ^3\left[𝒥_{LR}^{(0)}\right]^2\left[𝒥_{RR}^{(0)}+𝒥_{LL}^{(0)}\right]\mathrm{ln}\frac{D_0\tau }{\mathrm{}}.$$
(74)
As we have shown in Sec. III A, the spin decoherence by external irradiation does not require ionization of the impurity level, and therefore exists at frequencies below $`E_d`$, $`UE_d`$.
The effect of the irradiation on the Kondo conductance is not analytic in the intensity of the ac field. It cannot be obtained by a finite-order perturbation theory in $`\gamma `$ in the formula (31). To obtain Eq. (74) directly from Eq. (31) using the perturbation theory series in $`𝒥_{\alpha \alpha ^{}}^{(0)}`$, one would need to add up all the terms proportional to $`[𝒥_{\alpha \alpha ^{}}^{(0)}]^3[\gamma 𝒥_{\alpha \alpha ^{}}^{(0)}]^{2n}`$.
The finite-order perturbation theory \[Eqs. (34) and (74)\] can be used to evaluate the Kondo conductance only if the decoherence rate $`\mathrm{}/\tau `$ is much larger than the Kondo temperature $`T_K`$. At lower decoherence rates we have to take into account terms of all orders in $`𝒥_{\alpha \alpha ^{}}^{(0)}`$. It can be done by means of the Renormalization Group technique described in Sec. II C. One RG equation is to be derived from the condition of invariance of the conductance, given by the second and third orders of the perturbation theory in $`𝒥_{\alpha \alpha ^{}}^{(0)}`$ \[Eqs. (45) and (74)\], similarly to Eq. (48). The other two RG equations can be obtained using the requirement of invariance of the spin current (49) under the RG transformation. The resulting set of equations coincides with the one given by Eqs. (48), (51), and (52). When the decoherence rate exceeds the temperature $`T`$, the RG transformation must be stopped when the bandwidth $`D`$ reaches $`\mathrm{}/\tau `$ rather than $`T`$. Then the linear conductance can be evaluated in the second order perturbation theory in the renormalized exchange constants $`𝒥_{LR}`$, given by Eq. (54) with $`D=\mathrm{}/\tau `$:
$$G_{\mathrm{peak}}=\frac{3\pi ^2}{16}\frac{1}{\left[\mathrm{ln}(\mathrm{}/\tau T_K)\right]^2}G_U.$$
(75)
Here $`\mathrm{}/\tau T_K`$ is determined by Eq. (63). Equation (75) is the central formula of this section. It defines the conductance of the quantum dot as a function of two dimensionless parameters: $`\mathrm{}\omega /T_K`$ and $`\delta T_K/T_K`$ \[Eq. (64)\]. The region of validity of Eq. (75) is determined by the condition
$$\frac{\mathrm{}}{\tau T_K}1,$$
(76)
and corresponds to the regime of strong suppression of the Kondo effect by the external ac field. At fixed strength of the ac field the spin-flip rate (59) decreases with the decrease of ac field frequency $`\omega `$. Correspondingly, the peak conductance (75) grows. The crossover from weak to strong \[$`GG_U`$\] Kondo effect occurs when $`\mathrm{}/\tau T_K`$. Equations (59) and (63) show that this value of $`\mathrm{}/\tau `$ is reached while $`\mathrm{}\omega /T_K1`$.
## IV Weak spin decoherence
In this section we consider the regime of “intermediate suppression” of the Kondo effect by the ac radiation. By “intermediate” we mean that the decoherence is relatively weak
$$\mathrm{}/\tau <T_K,$$
(78)
and the formula (75) is no longer valid, but the frequency is still sufficiently high
$$T_K<\mathrm{}\omega ,$$
(79)
so that the RG result (63) for the decoherence rate holds. In this regime, the formation of the many-body state is not suppressed, because of Eq. (78). However, Eq. (79) allows for sudden spin flips. The complicated nature of the many-body state hampers the quantitative consideration of this regime, and we limit ourselves to qualitative analysis.
When the many-body Kondo resonance is fully formed, the conductance of the dot equals $`G_U`$ \[Eq. (57)\] and corresponds to the unitary limit of tunneling through the dot. An act of spin flip destroys the many-body state, and the conductance drops substantially below the value given by Eq. (57). The time necessary for the many-body state to be restored equals approximately $`\mathrm{}/T_K`$. Therefore, the fraction of time which the system spends in the highly-conducting ($`GG_U`$) state equals approximately $`1a\mathrm{}/\tau T_K`$, where $`a1`$. The resulting time-averaged conductance of the dot can be estimated as
$$G_{\mathrm{peak}}=\left[1a\frac{\mathrm{}}{\tau T_K}\right]G_U.$$
(80)
The rate $`\mathrm{}/\tau `$ of the spin-flip processes here is given by Eq. (63). Under conditions (IV), parameter $`a`$ does not depend on the characteristics of the ac field. The value of $`a`$ should be found from the quantum-mechanical problem of evolution, which starts with a state “prepared” by the flip of the impurity spin, and results eventually in the re-formation of a Kondo polaron. Our qualitative treatment of the regime (IV) does not allow us to find the exact value of the universal coefficient $`a`$, which however can be found from a numeric calculation.
At the upper limit of applicability, $`\mathrm{}/\tau T_K`$, the peak conductance given by Eq. (80) matches the result (75).
## V Low-frequency ac field: adiabatic approximation
In Sections II C and III we considered the case when the isolated spin is only weakly screened by the many-electron state formed around it. The complete screening was suppressed either by relatively high temperature, $`T>T_K`$, or by large bias $`eV_{\mathrm{dc}}>T_K`$, or by the decoherence. In the case of Sec. IV, the spin-screening cloud is able to form; however, the spin flips, produced by the irradiation, occasionally destroy this many-body state, thus reducing the conductance.
In this section, we consider the case of low frequencies of the ac field $`\mathrm{}\omega T_K`$, when the energy of a photon is insufficient to flip the dot’s spin in the fully formed many-body Kondo state. For the irradiation to be the leading cause of the deviation of the conductance from the unitary limit, we suppose the temperature and bias to be also low, $`T,eV_{\mathrm{dc}}T_K`$. The RG technique we used before is not applicable in this regime. Therefore we need another approach to evaluate the conductance of the quantum dot system and the effects of the external irradiation on it.
The required approach is provided by the scaling theory of Nozières and Blandin. This theory states that the renormalization-group transformation, whose initial stage was described in Sec. II C, can be continued, and finally leads to a fixed point. At the fixed point, the system exhibits Fermi-liquid behavior, and its Hamiltonian has a relatively simple form. This fixed-point Hamiltonian can be used to study the properties of the Kondo system at low temperatures, $`TT_K`$. Mapping the quantum dot system in the Kondo regime onto the regular one-channel Kondo problem, we can employ the fixed-point Hamiltonian to evaluate the dc current through the dot induced by the applied bias.
The external ac field disturbs the many-particle state formed near the isolated spin, leading to the deviations of the system behavior from that dictated by the (static) fixed-point Hamiltonian. In this section we study the case when the frequency of the field is low ($`\mathrm{}\omega T_K`$), so that the many-body state is not destroyed but rather adiabatically varied by the ac field, as the level in the dot goes up and down \[see Eq. (5)\]. Then the current through the dot can be evaluated with the help of the fixed-point Hamiltonian with time-dependent parameters.
Now we map the problem of transport through the dot onto the regular scattering problem. For this purpose, it is convenient to use the basis of $`s`$ and $`p`$ scattering states rather than that of the left-lead and right-lead states. These two bases are connected by
$`a_{k\sigma }^{(s)}=\xi c_{k\sigma L}+\eta c_{k\sigma R},a_{k\sigma }^{(p)}=\eta c_{k\sigma L}+\xi c_{k\sigma R},`$ (81)
$`\text{where}\xi {\displaystyle \frac{v_L}{\sqrt{v_L^2+v_R^2}}},\eta {\displaystyle \frac{v_R}{\sqrt{v_L^2+v_R^2}}}.`$ (82)
The $`p`$-states are decoupled from the dot, so the dot-lead coupling term in the Anderson Hamiltonian (5) has the form
$$\sqrt{v_L^2+v_R^2}\underset{k,\sigma }{}(a_{k\sigma }^{(s)}d_\sigma ^{}+\mathrm{H}.\mathrm{c}.).$$
The initial basis $`c_{k\sigma \alpha }`$ is composed of the states residing entirely in the left or right lead, which is convenient for the problem of two leads connected by a weak link, when the inter-lead tunneling is to be considered as a perturbation. In terms of incident and reflected/transmitted waves, these states correspond to the waves incident from one of the leads to the dot and completely reflected back to the same lead. Therefore the $`s`$-waves of Eq. (81), which enter the new basis, have the scattering phase equal to $`\pi /2`$.
Making the Schrieffer-Wolff transformation, we arrive to the regular Kondo problem, which at low temperatures can be studied with the help of the fixed-point Hamiltonian. Under these conditions, the $`s`$-wave electrons, interacting via the isolated spin, form the screening cloud. This many-body state still has the Fermi-liquid properties, though its scattering characteristics are different from those of just an isolated spin. One of the principal differences is the shift of the scattering phase by $`\pi /2`$ for the states at the Fermi level. This suggests another change of the basis for the sake of convenience: from $`s`$-waves having the scattering phase equal to $`\pi /2`$, $`a_{k\sigma }^{(s)}`$, to those with the scattering phase $`\pi `$.
The formal relation between the $`a_{k\sigma }^{(s)}`$ and the new basis, which we denote $`b_{k\sigma }`$, is given by
$`b_{k\sigma }`$ $``$ $`{\displaystyle 𝑑xe^{ikx}\widehat{\mathrm{\Psi }}_\sigma (x)},`$ (83)
$`\widehat{\mathrm{\Psi }}_\sigma (x)`$ $`=`$ $`\mathrm{exp}\left[i\pi {\displaystyle _{\mathrm{}}^x}𝑑x^{}g(x^{})\right]\widehat{\psi }_\sigma ^{(s)}(x),`$ (84)
$`\widehat{\psi }_\sigma ^{(s)}(x)`$ $``$ $`{\displaystyle 𝑑ke^{ikx}a_{k\sigma }^{(s)}},`$ (85)
where $`g(x)`$ is an arbitrary function obeying $`_{\mathrm{}}^{\mathrm{}}𝑑x^{}g(x^{})=1`$. The “coordinate” $`x`$ was introduced for convenience to separate the incoming and outgoing parts of the scattering states, which correspond to negative and positive values of $`x`$ respectively.
Before the scattering region \[$`x\mathrm{}`$ in Eq. (85)\], the wave functions of the states $`b_{k\sigma }`$ and $`a_{k\sigma }^{(s)}`$ coincide. Therefore the states
$$c_{k\sigma L}^{(\mathrm{in})}=\xi b_{k\sigma }\eta a_{k\sigma }^{(p)},c_{k\sigma R}^{(\mathrm{in})}=\eta b_{k\sigma }+\xi a_{k\sigma }^{(p)}$$
(86)
represent waves incident from a left or right lead only. Passing the scattering region, the wave function of the state $`b_{k\sigma }`$ acquires an extra phase of $`\pi `$ as compared to that of $`a_{k\sigma }^{(s)}`$. Then one can see that the states
$$c_{k\sigma L}^{(\mathrm{out})}=\xi b_{k\sigma }\eta a_{k\sigma }^{(p)},c_{k\sigma R}^{(\mathrm{out})}=\eta b_{k\sigma }+\xi a_{k\sigma }^{(p)}$$
(87)
have an outgoing wave only in the left or right lead respectively. The current operator in terms of these states equals simply
$$\widehat{I}(V)=\underset{k,\sigma }{}\left(c_{k\sigma L}^{(\mathrm{out})}c_{k\sigma L}^{(\mathrm{out})}c_{k\sigma R}^{(\mathrm{out})}c_{k\sigma R}^{(\mathrm{out})}\right).$$
(88)
The fixed-point Hamiltonian in the basis $`b_{k\sigma }`$, $`a_{k\sigma }^{(s)}`$ has a relatively simple form:
$`\widehat{H}_{\mathrm{fp}}`$ $`=`$ $`v_F{\displaystyle \underset{k\sigma }{}}kb_{k\sigma }^{}b_{k\sigma }+v_F{\displaystyle \underset{k\sigma }{}}ka_{k\sigma }^{(p)}a_{k\sigma }^{(p)}`$ (89)
$``$ $`{\displaystyle \frac{v_F}{\nu T_K}}{\displaystyle \underset{k_1k_2\sigma }{}}(k_1+k_2)b_{k_1\sigma }^{}b_{k_2\sigma }^{}`$ (90)
$`+`$ $`{\displaystyle \frac{1}{\nu ^2T_K}}{\displaystyle \underset{k_1k_2k_3k_4}{}}:b_{k_1}^{}b_{k_2}^{}b_{k_3}^{}b_{k_4}^{}:,`$ (91)
where $`:\mathrm{}:`$ denotes normal ordering. The spectrum of electrons is linearized, $`\epsilon _k=v_Fk`$, since the reduced bandwidth is of the order of $`T_K\epsilon _F`$; The Kondo temperature $`T_K`$ is the only energy scale of the fixed-point Hamiltonian (89).
The third term in Eq. (89) determines the phase shift which a quasiparticle acquires as it passes through the dot. This shift is energy-dependent: it equals $`\pi `$ at the Fermi level, as we discussed before; and $`\pi +\epsilon _k/T_K`$ in general case. In terms of waves incident from the left or right lead, such behavior of the phase shift is analogous to those in tunneling through a resonant state tied to the Fermi level. The fourth term in the Hamiltonian describes the interaction of the quasiparticles of the Fermi liquid at the fixed point. The $`p`$-waves are not affected by the Kondo screening, so the Hamiltonian for them has the same form as the one given by Eqs. (5), (81).
Using the Hamiltonian (89), we can rewrite the current operator (88) in the form more convenient for the further calculations:
$$\widehat{I}(V)=(2\eta \xi )^2\left\{\frac{e^2}{\pi \mathrm{}}V\frac{ie}{\mathrm{}}[\widehat{H}_{\mathrm{fp}},\underset{k,\sigma }{}c_{k\sigma R}^{(\mathrm{out})}c_{k\sigma R}^{(\mathrm{out})}]\right\}.$$
(92)
The first term in Eq. (92) is the current that would flow if all the states were able for resonant tunneling through the dot; the scattering between the left- and right-incident species (which is just backscattering when $`\xi =\eta `$) reduces the magnitude of the current, and is accounted for by the second term.
To evaluate the conductance of the dot, we employ the Keldysh technique (31), treating the last two terms of Hamiltonian (89) as a perturbation. At infinitely small temperature and bias, the current through the dot is transferred by electrons at the Fermi level. The transmission coefficient for these electrons equals $`(2\xi \eta )^2`$, i.e. the second (backscattering) term in the current operator (92) yields zero. Therefore the dot at this conditions has maximum conductance, $`G=G_U(e^2/\pi \mathrm{})(2\eta \xi )^2`$. At finite temperatures, the electrons which facilitate the current are spread within a strip of width $`T`$ near the Fermi level. The departure of the particle energy from the Fermi level in the system (89) leads to the deviation of its scattering phase from $`\pi `$, i.e. from the resonance. Therefore the conductance in this case will be lower than $`G_U`$. Indeed, substituting Eqs. (92) and (89) into Eq. (31) and employing the second order of the perturbation theory in the last two terms of Hamiltonian (89), we arrive at
$`G_{\mathrm{peak}}(T)`$ $`=`$ $`{\displaystyle \frac{e^2}{\pi \mathrm{}}}(2\eta \xi )^2\{1{\displaystyle \frac{v_F}{\nu T_K^2}}{\displaystyle \underset{k}{}}k^2[{\displaystyle \frac{df(v_Fk)}{dk}}]`$ (93)
$`{\displaystyle \frac{2}{v_F\nu ^3T_K^2}}{\displaystyle \underset{k_1k_2k_3}{}}\left[{\displaystyle \frac{df(v_Fk_1)}{dk_1}}\right]f(v_Fk_2)`$ (94)
$`\times [1f(v_Fk_3)]f[v_F(k_1k_2+k_3)]]\}`$ (95)
$`=`$ $`\left[1\pi ^2\left({\displaystyle \frac{T}{T_K}}\right)^2\right]G_U,`$ (96)
where $`f(\epsilon )1/[\mathrm{exp}(\epsilon /T)+1]`$ is the Fermi distribution function. One can see from Eq. (96) that the conductance of the quantum dot system at low temperatures decreases with growth of the temperature. This behavior has been observed experimentally and is analogous to the decrease of the resistivity in a regular Kondo system (bulk metal with magnetic impurities).
The differential conductance of the dot at finite bias $`V_{\mathrm{dc}}`$, with $`TeV_{\mathrm{dc}}T_K`$, can be derived analogously to Eq. (96). The resulting formula
$$G(V_{\mathrm{dc}})=\left[1\frac{3}{8}\left(\frac{eV_{\mathrm{dc}}}{T_K}\right)^2\right]G_U$$
(97)
shows that $`G(V_{\mathrm{dc}})`$ decreases with the growth of the bias applied to the dot.
Slow ($`\mathrm{}\omega T_K`$) ac field results in adiabatic time-dependence of Kondo temperature, see Eq. (65). The time-dependent part of the Hamiltonian (89) with $`1/T_K(t)`$ given by Eq. (65) accounts for the interaction of quasiparticles with ac field. To consider this part of Hamiltonian in the conventional terms of electron-photon interaction, we expand $`1/T_K(t)`$ in Fourier series:
$$\frac{1}{T_K(t)}\underset{n}{}\frac{1}{T_K^{(n)}}e^{in\omega t}$$
(98)
After an act of photon absorption, a quasiparticle facilitating the current is transferred from the Fermi level, i.e. away from the resonance. As a result, at low temperatures the ac field must reduce the conductance of a quantum dot in the Kondo regime. At $`\mathrm{}\omega T_K`$, the conductance can be calculated in the second order of the perturbation theory in the time-dependent part of the Hamiltonian. Substituting Eq. (98) into Eq. (89), and then using the Keldysh formalism (31) to evaluate the conductance we arrive at
$`G_{\mathrm{peak}}`$ $`=`$ $`\{1{\displaystyle \underset{n}{}}\left({\displaystyle \frac{1}{T_K^{(n)}}}\right)^2[{\displaystyle \frac{v_F^2}{\nu }}{\displaystyle \underset{k}{}}k^2\delta (v_Fkn\mathrm{}\omega )`$ (101)
$`+{\displaystyle \frac{2}{\nu ^3}}{\displaystyle \underset{k_1k_2k_3}{}}\delta (v_Fk_1n\mathrm{}\omega )\theta (v_Fk_2)`$
$`\times [1\theta (v_Fk_3)]\theta [v_F(k_1k_2+k_3)]]\}G_U`$
$`=`$ $`\left\{13{\displaystyle \underset{n}{}}\left({\displaystyle \frac{\mathrm{}n\omega }{T_K^{(n)}}}\right)^2\right\}G_U,`$ (102)
where for simplicity we set temperature to zero. Transforming Eq. (102) back from $`1/T_K^{(n)}`$ to $`1/T_K(t)`$ \[Eq. (98)\], we finally obtain
$`G_{\mathrm{peak}}`$ $`=`$ $`\left\{13\overline{\left({\displaystyle \frac{d}{dt}}{\displaystyle \frac{1}{T_K(t)}}\right)^2}\right\}G_U`$ (103)
$``$ $`\left\{1{\displaystyle \frac{3}{2}}\left({\displaystyle \frac{\delta T_K}{T_K}}\right)^2\left({\displaystyle \frac{\mathrm{}\omega }{T_K}}\right)^2\right\}G_U,`$ (104)
where $`\overline{\mathrm{}}`$ denotes averaging over the period of variation of $`T_K(t)`$, and $`\delta T_K/T_K`$ is defined by Eq. (64).
The single-photon decoherence processes described in Sec. III A do not occur in this regime, because the energy necessary to flip the dot’s spin is increased by its interaction with the screening “spin cloud” in the leads, and is of the order of $`T_K\mathrm{}\omega `$. The rate of the spin flip due to many-photon processes is exponentially small in $`T_K/\mathrm{}\omega `$.
## VI Scaling formula for the conductance
In this section we summarize the results obtained in Secs. IIIV for the effect of the periodic modulation of the dot’s potential on the Kondo conductance.
In the absence of ac irradiation, the quantum dot system is described by a number of physical parameters, see Eqs. (II B), (24). However, in the Kondo regime all these parameters combine into a single relevant energy scale, $`T_K`$, see Eq. (55), controlling the behavior of the system, see, e.g., Eqs. (56), and (96). The periodic modulation $`V_{\mathrm{dot}}\mathrm{cos}\omega t`$ of the dot potential adds two more parameters to the initial Hamiltonian (II B), and, most importantly, drives the system into a non-equilibrium state. Surprisingly, such a drastic perturbation does not break down the universal description of the problem, and the Kondo temperature remains the only relevant energy scale. We have shown that the effect of the irradiation is described by two dimensionless parameters $`\mathrm{}\omega /T_K`$ and $`\delta T_K/T_KV_{\mathrm{dot}}`$, where $`\delta T_K`$ has the meaning of the adiabatic variation of the Kondo temperature under the influence of ac modulation, see Eq. (64).
At sufficiently large frequencies $`\omega `$ of the ac field, when
$$\frac{\mathrm{}\omega }{T_K}>\frac{32}{3\pi }\frac{\left[\mathrm{ln}(\delta T_K/T_K)\right]^4}{\left[\delta T_K/T_K\right]^2},$$
(105)
the rate $`\mathrm{}/\tau `$ of the spin-flip cotunneling exceeds the Kondo temperature $`T_K`$. The spin-flip cotunneling brings decoherence into the spin dynamics of the dot, destroying the Kondo resonance. Small lifetime of the Kondo resonance leads to a significant suppression of the Kondo effect, see Sec. III B. The dependence of the zero-bias dc conductance $`G_{\mathrm{peak}}`$ of the dot on the power and frequency of the ac field is given by Eqs. (75) and (63).
Upon lowering the frequency $`\omega `$, condition (105) breaks down, and $`\mathrm{}/\tau `$ becomes smaller than the Kondo temperature. Under such conditions, the strong suppression of the Kondo conductance is not possible. However, the conductance still may deviate from the unitary limit $`G_U`$. The violation of the condition (105) occurs while $`\mathrm{}\omega `$ still exceeds $`T_K`$. The zero-bias conductance in this regime can be estimated by Eq. (80) and (63).
At frequencies below the Kondo temperature, the ac field is unable to flip the spin of the dot, and the spin-flip cotunneling does not occur. In this regime, the ac-driven deviation from the unitary limit is small and can be accounted within the framework of the Fermi-liquid description. The main role of the ac field is to scatter the conduction electrons, transferring them to the energies off the Fermi level. These scattered electrons miss the Kondo resonance, which is tied to the Fermi level. It produces the small deviation of the dc conductance $`G_{\mathrm{peak}}`$ from the unitary limit, see Sec. V, Eq. (103).
The results obtained for these three regimes match each other on the corresponding limits of applicability. It allows us to piece together the dependence of $`G_{\mathrm{peak}}`$ on $`\delta T_K/T_K`$ and $`\mathrm{}\omega `$ in a broad frequency range, see Fig. 2.
This dependence allows us to conjecture that at small $`\delta T_K/T_K`$ the conductance can be cast in the following form:
$$\frac{G_{\mathrm{peak}}}{G_U}=F\left[\left(\frac{\delta T_K}{T_K}\right)^2f\left(\frac{\mathrm{}\omega }{T_K}\right)\right],$$
(106)
with two universal functions $`F(x)`$ and $`f(y)`$. Each of the functions depends on only one variable; they have the following asymptotes
$`F(x)=\{\begin{array}{cc}1ax,\hfill & \text{if}x1,\hfill \\ {\displaystyle \frac{3\pi ^2}{16}}{\displaystyle \frac{1}{(\mathrm{ln}x)^2}},\hfill & \text{if}x1,\hfill \end{array}`$ (109)
and
$`f(y)=\{\begin{array}{cc}{\displaystyle \frac{1}{a}}3y^2\text{},\hfill & \text{if}y1,\hfill \\ {\displaystyle \frac{3\pi }{32}}{\displaystyle \frac{y}{(\mathrm{ln}y)^4}},\hfill & \text{if}y1.\hfill \end{array}`$ (112)
The numerical parameter $`a1`$ is introduced and discussed in Sec. IV. When $`\mathrm{}\omega T_K`$, the argument of function $`F`$ has the meaning of dimensionless decoherence rate $`\mathrm{}/\tau T_K`$.
## VII Conductance suppression by ac bias
In sections IIIVI we considered the effects of modulation of the dot potential on the Kondo conductance. In the present section we study the conductance in the system where the ac field is applied to the leads, thus creating alternating bias $`V_{\mathrm{ac}}`$. The parameters characterizing such a field are the dimensionless amplitude $`eV_{\mathrm{ac}}/T_K`$ and frequency $`\mathrm{}\omega ^{}/T_K`$.
First we consider the case of “fast” ac bias, $`\mathrm{}\omega ^{}\mathrm{max}\{T_K,eV_{\mathrm{ac}}\}`$. Under these conditions, the ac bias affects the Kondo conductance through the decoherence of the dot’s spin. The dependence of the corresponding decoherence rate $`\mathrm{}/\tau ^{}`$ on the amplitude and frequency of the ac bias can be calculated with the help of the Renormalization Group technique which we used in Sec. II C, III A, and III B. The resulting expression reads
$$\frac{\mathrm{}}{\tau ^{}T_K}=\frac{1}{\pi }\frac{G_U}{e^2/\pi \mathrm{}}\left(\frac{eV_{\mathrm{ac}}}{T_K}\right)^2\frac{T_K}{\mathrm{}\omega ^{}}\frac{1}{[\mathrm{ln}(\mathrm{}\omega ^{}/T_K)]^2}.$$
(113)
Note that, in contrast to the ac modulation of the gate voltage (Sec. III A), in the case of ac bias the rate of decoherence decreases with the growth of the field frequency $`\omega ^{}`$. The parameter
$$\frac{G_U}{e^2/\pi \mathrm{}}\frac{4\mathrm{\Gamma }_L\mathrm{\Gamma }_R}{(\mathrm{\Gamma }_L+\mathrm{\Gamma }_R)^2}$$
characterizes the asymmetry in the dot, and emerges in the expressions for quantities associated with electron transfer between the leads.
When $`\mathrm{}/\tau ^{}T_K>1`$, the conductance can be evaluated by means of the perturbation theory, see Sec. II C. The decaying function $`S_j(t_1)S_k(t_2)S_l(t_3)`$, which enters the terms of the perturbation theory, provides the large-time cut-off for the integrals in equations of the type (37). The derivation of the expressions for the conductance is identical to the one given in Sec. III B, cf. Eq. (75). The final formula reads
$$G_{\mathrm{peak}}=\frac{3\pi ^2}{16}\frac{1}{\left[\mathrm{ln}(\mathrm{}/\tau ^{}T_K)\right]^2}G_U,$$
(114)
with $`\mathrm{}/\tau ^{}T_K`$ given by Eq. (113).
At smaller amplitudes, $`\mathrm{}/\tau ^{}T_K<1`$, the ac bias is unable to suppress the formation of the Kondo many-electron state. For this case we may repeat the reasoning of Sec. IV. As the result, we obtain
$$G_{\mathrm{peak}}=\left[1a\frac{\mathrm{}}{\tau ^{}T_K}\right]G_U,$$
(115)
i.e. the Kondo conductance is only weakly suppressed.
In the opposite limit of slow variations of bias, $`\mathrm{}\omega ^{}\mathrm{max}\{T_K,eV_{\mathrm{ac}}\}`$, one can use the adiabatic approximation,
$$G_{\mathrm{peak}}=\overline{G(V_{\mathrm{ac}}\mathrm{cos}\omega ^{}t)}.$$
(116)
Here $`G(V)`$ is the differential dc conductance at finite bias $`V`$, and $`\overline{\mathrm{}}`$ denotes averaging over the period of variation of ac bias. For $`eV/T_K1`$, the conductance is given by Eq. (97). Substituting it into Eq. (116), we obtain
$$G_{\mathrm{peak}}=\left\{1\frac{3}{16}\left(\frac{eV_{\mathrm{ac}}}{T_K}\right)^2\right\}G_U.$$
(117)
In the opposite case, $`eV_{\mathrm{ac}}/T_K1`$, we obtain using Eqs. (58) and (116):
$$G_{\mathrm{peak}}=\frac{3\pi ^2}{16}\frac{1}{[\mathrm{ln}(eV_{\mathrm{ac}}/T_K)]^2}G_U.$$
(118)
Figure 3 shows the possible regimes of the ac bias effect on the Kondo conductance of a dot. At
$$\mathrm{}\omega ^{}\mathrm{max}\{eV_{\mathrm{ac}},T_K\},$$
(120)
the peak conductance depends only on $`eV_{\mathrm{ac}}/T_K`$, see Eqs. (117), (118). In the opposite case of high frequencies,
$$\mathrm{}\omega ^{}\mathrm{max}\{eV_{\mathrm{ac}},T_K\},$$
(121)
the peak conductance depends only on $`\mathrm{}/\tau ^{}T_K`$, see Eqs. (114), (115). Thus in both regions (120) and (121), $`G_{\mathrm{peak}}`$ is a function of a single variable. However, the corresponding variables are different in the two regions. Therefore, at the crossover between these two frequency domains, the peak conductance $`G_{\mathrm{peak}}(\mathrm{}\omega ^{}/T_K,eV_{\mathrm{ac}}/T_K)`$ can not be cast into a simple form of a single-variable function.
It is instructive to consider the peak conductance as a function of the frequency $`\omega ^{}`$ at a fixed field amplitude $`eV_{\mathrm{ac}}`$. At small amplitudes, $`eV_{\mathrm{ac}}/T_K1`$, the suppression of the Kondo effect is weak at any frequency. At stronger fields, $`eV_{\mathrm{ac}}/T_K1`$, the Kondo effect is suppressed far below the unitary limit at low frequencies. The height of the zero-bias peak grows with the increase of the field frequency $`\omega ^{}`$, see Fig. 4. The description of the crossover between the regimes (120) and (121) is developed in Appendix A.
## VIII Satellite conductance peaks
In Sections IIIVII of this paper we mostly concentrated on the effects of the ac field on the zero-bias Kondo conductance $`G_{\mathrm{peak}}`$. In this section we study how the ac field modifies the finite-bias differential Kondo conductance $`G(V_{\mathrm{dc}})`$.
Without an ac field, the dependence of the differential Kondo conductance on $`V_{\mathrm{dc}}`$ is given by Eq. (97) for $`eV_{\mathrm{dc}}T_K`$ and by Eq. (58) for $`eV_{\mathrm{dc}}T_K`$. As we have seen from the previous sections, the ac field reduces the height of the zero-bias peak, $`G_{\mathrm{peak}}G(V_{\mathrm{dc}}=0)`$.
Another effect of external irradiation on the differential conductance $`G(V_{\mathrm{dc}})`$ is in producing satellite peaks at $`eV_{\mathrm{dc}}=\pm n\mathrm{}\omega `$. If an external ac field is applied, then, at $`eV_{\mathrm{dc}}=\pm n\mathrm{}\omega `$, a tunneling electron can hop from a state at the Fermi level in one lead to a state at the Fermi level in the other lead, emitting or absorbing $`n`$ photons. Thus at finite bias the external irradiation can effectively put a tunneling electron into zero-bias conditions, and the Kondo anomaly in the conductance is revived. The height of these peaks can be calculated from the formula (31) similarly to Eq. (40). At low enough irradiation level, $`\gamma 1`$, it is sufficient to consider only one-photon processes, accounted for by the Hamiltonian (II B) with the coupling constants $`𝒥_{\alpha \alpha ^{}}(t)`$ given by Eqs. (23)–(25). In this approximation, we will be able to describe the first pair of satellite peaks, which emerge next to the main, zero-bias, peak. The resulting correction to the conductance at $`e|V_{\mathrm{dc}}|`$ close to $`\mathrm{}\omega `$ has the form
$`G_{\mathrm{stl},\pm }^{(3)}(V_{\mathrm{dc}})={\displaystyle \frac{3\pi ^2}{8}}{\displaystyle \frac{e^2}{\pi \mathrm{}}}\nu ^3\left[𝒥_{LR}^{(0)}\right]^2\left[𝒥_{RR}^{(0)}+𝒥_{LL}^{(0)}\right]\gamma ^2`$ (122)
$`\times {\displaystyle _{\mathrm{}}^0}dt(t)\mathrm{exp}(|t|/\tau _{\mathrm{stl}})`$ (123)
$`\times \left({\displaystyle \frac{\pi T}{\mathrm{}}}\right)^2{\displaystyle \frac{\mathrm{cos}[(eV_{\mathrm{dc}}\pm \mathrm{}\omega )t/\mathrm{}]}{\mathrm{sinh}^2(\pi Tt/\mathrm{})+(T/D_0)^2}}.`$ (124)
When $`eV\pm \mathrm{}\omega `$, the cosine function cuts off the logarithmic divergence here. However, when $`eV_{\mathrm{dc}}\pm \mathrm{}\omega `$, the cosine factor becomes essentially constant \[cf. Eq. (40) at $`V_{\mathrm{dc}}0`$\], and the differential conductance has a peak again. At $`T0`$, the height of the satellite conductance peak is determined by the spin decoherence rate $`\mathrm{}/\tau _{\mathrm{stl}}`$. We must mention that $`\tau _{\mathrm{stl}}`$ may be significantly shorter than $`\tau `$ given by Eq. (59). The time $`\tau `$ characterizes the spin decoherence at zero bias, whereas the satellite corresponds to a finite bias $`eV_{\mathrm{dc}}=\pm \mathrm{}\omega `$. In the latter case, the spin decoherence occurs mostly due to the tunneling of electrons through the dot (see Fig. 5, and also Ref. ). The rate of this process at $`\mathrm{}\omega T_K`$ is given by
$$\frac{\mathrm{}}{\tau _{\mathrm{stl}}}=\frac{1}{2\pi }\mathrm{}\omega \frac{\mathrm{\Gamma }_L\mathrm{\Gamma }_R}{\stackrel{~}{E}_d^2}.$$
(125)
To extend the result (125) to lower frequencies $`\mathrm{}\omega T_K`$, we employ the RG technique. As the result, we obtain the decoherence rate as a function of the universal parameter $`\mathrm{}\omega /T_K`$:
$$\frac{\mathrm{}}{\tau _{\mathrm{stl}}T_K}=\frac{3\pi }{32}\frac{G_U}{e^2/\pi \mathrm{}}\frac{\mathrm{}\omega /T_K}{[\mathrm{ln}(\mathrm{}\omega /T_K)]^2}.$$
(126)
Note that Eqs. (125) and (126) imply $`eV_{\mathrm{dc}}\mathrm{}\omega `$.
Using the RG technique, we derive the formula for $`G_{\mathrm{stl},\pm }(V_{\mathrm{dc}})`$, which is the contribution of one-photon processes to the differential dc conductance:
$`G_{\mathrm{stl},\pm }(V_{\mathrm{dc}})`$ $`=`$ $`\left[\mathrm{ln}{\displaystyle \frac{\sqrt{(\mathrm{}/\tau _{\mathrm{stl}})^2+(eV_{\mathrm{dc}}\pm \mathrm{}\omega )^2}}{T_K}}+\mathrm{ln}{\displaystyle \frac{\mathrm{}\omega }{T_K}}\right]^4`$ (127)
$`\times `$ $`{\displaystyle \frac{3\pi ^2}{4}}\left[{\displaystyle \frac{\delta T_K}{T_K}}\right]^2G_U.`$ (128)
One can see from Eqs. (126)–(128) that $`G_{\mathrm{stl}}`$ depends on the parameters of irradiation only through the universal variables $`\delta T_K/T_K`$ and $`\mathrm{}\omega /T_K`$. The details of derivation of Eq. (128) are given in Appendix B.
The full expression for $`G(V_{\mathrm{dc}})`$ reads
$$G(V_{\mathrm{dc}})=G_{\mathrm{main}}(V_{\mathrm{dc}})+G_{\mathrm{stl},+}(V_{\mathrm{dc}})+G_{\mathrm{stl},}(V_{\mathrm{dc}}).$$
(129)
Here $`G_{\mathrm{main}}`$ accounts for the tunneling through the dot without absorption/emission of photons and is responsible for the zero-bias Kondo peak. The other two terms in Eq. (129) describe the satellite peaks in $`G(V_{\mathrm{dc}})`$.
As the criterion for resolution of the satellites, we adopt the requirement that the function $`G(V_{\mathrm{dc}})`$ must be non-monotonic on the sides of the zero-bias peak. This requirement can be reformulated as:
$$\frac{\delta T_K}{T_K}>1,\frac{\mathrm{}\omega }{T_K}1.$$
(130)
In the derivation of the condition (130) we used Eq. (129) with $`G_{\mathrm{main}}(V_{\mathrm{dc}})`$ given by Eq. (58). Such a form of the elastic Kondo conductance $`G_{\mathrm{main}}`$ should be used because at $`e|V_{\mathrm{dc}}|\mathrm{}\omega `$ it is suppressed mainly due to the finite bias rather than due to the decoherence, since $`\mathrm{}\omega \mathrm{}/\tau _{\mathrm{stl}}\mathrm{}/\tau `$.
The above consideration was performed for the ac field applied to the gate. The case of ac bias is can be considered similarly. The third-order perturbation theory result for the shape of satellite peak may be obtained from Eqs. (124) and (128) by replacing $`\gamma ^2`$ with $`(\gamma ^{})^2`$. The RG treatment yields
$`G_{\mathrm{stl},\pm }(V_{\mathrm{dc}})`$ $`=`$ $`\left[\mathrm{ln}{\displaystyle \frac{\sqrt{(\mathrm{}/\tau _{\mathrm{stl}})^2+(eV_{\mathrm{dc}}\pm \mathrm{}\omega )^2}}{T_K}}+\mathrm{ln}{\displaystyle \frac{\mathrm{}\omega }{T_K}}\right]^4`$ (131)
$`\times `$ $`\left[\mathrm{ln}{\displaystyle \frac{\mathrm{}\omega }{T_K}}\right]^4{\displaystyle \frac{3\pi ^2}{4}}\left[{\displaystyle \frac{eV_{\mathrm{ac}}}{\mathrm{}\omega ^{}}}\right]^2G_U.`$ (132)
The condition for the satellite peaks to be clearly visible takes the form
$$\frac{eV_{\mathrm{ac}}}{\mathrm{}\omega ^{}}>1,\frac{\mathrm{}\omega }{T_K}1.$$
(133)
Conditions (130) and (133), together with Eq. (75), demonstrate that upon the increase of the amplitude of the ac field, the zero-bias peak is suppressed first, and only after that the satellite peaks may become distinguishable from the background conductance.
## IX Comparison with experiment
In this section, we discuss possible ways of comparison of the theory presented here with experiments. Ideally, one should measure the dc conductance through a dot while applying the ac bias to the gates or leads in a controllable way, and varying its frequency $`\omega `$ and amplitude in a broad range. Our results predict that the data obtained at various values of $`T_K`$ should be scalable, when using the proper dimensionless variables, Eqs. (109), (112), (114), (115), (117), and (118). We predict also that at a fixed magnitude of ac field, the suppression of the Kondo effect becomes more severe with the increase of $`\omega `$, if the applied field modulates the gate potential; the dependence on $`\omega `$ in the case of ac field applied to the leads is opposite. If $`\omega `$ significantly exceeds $`T_K`$, see Eqs. (130) and (133), observation of “satellites” at $`eV=\mathrm{}\omega `$ of the main Kondo singularity may become possible. The appearance of even small side peaks though is associated with a strong suppression of the zero-bias Kondo singularity.
Presently, there is only one experiment aimed at observation of effects of irradiation on the Kondo conductance in the quantum dot. In the analysis of this experiment below, we will see that the frequencies used were of the order of $`T_K`$, and the range of $`\omega `$ was less than one decade. In addition, it was impossible to calibrate the amplitude of the field applied to the device; the attenuation coefficients were frequency-dependent. Therefore, one could not perform the measurements of $`G_{\mathrm{peak}}(\omega )`$ at a fixed field amplitude. However, it was possible to measure the dependence of the peak conductance on the amplitude of the applied field at a discrete set of fixed frequencies. The ac source was powerful enough to allow the authors of Ref. to completely suppress the Kondo anomaly. Nominally, the ac field was applied to the gate, but it was apparently hard to exclude “leaking” of the field to other electrodes of the device. It creates further ambiguity in the interpretation of the experimental results.
The Kondo temperature of the system can be found from the dependencies $`G_{\mathrm{peak}}(T)`$ and $`G(V_{\mathrm{dc}})`$ measured in the absence of irradiation. Comparing the experimental results with the theory \[Eqs. (96) and (97)\], we obtain the estimate $`T_K3060`$ $`\mu `$eV. The external irradiation was applied through a high-frequency coaxial cable, coupled capacitively to the gate. The frequency of the irradiation ranged from 10 GHz to 50 GHz (i.e. $`\mathrm{}\omega `$ between 40$`\mu `$eV and 200 $`\mu `$eV).
The zero-bias Kondo peak was clearly observed in the $`G(V_{\mathrm{dc}})`$ dependence, when no ac field was applied. With the increase of the amplitude of the ac field, the height of the zero-bias peak decreased for each of the frequencies used. Such a behavior is in agreement with the our conclusion that irradiation must suppress the Kondo effect in a quantum dot even when dot ionization does not occur. Also, the satellites did not appear, in agreement with our theory for $`\mathrm{}\omega T_K`$, see Sec. VIII.
The authors of Ref. attempted the data collapse for the dependence of $`G_{\mathrm{peak}}`$ on the amplitude $`V_{\mathrm{irr}}`$ and frequency $`\omega `$ of the ac field. In this procedure, $`V_{\mathrm{irr}}`$ was scaled by some frequency-dependent parameter, in order to bring to a single curve the dependencies $`G_{\mathrm{peak}}`$ vs. $`V_{\mathrm{irr}}`$ measured at different values of $`\omega `$. Successful data collapse means that a universal dependence exists,
$$G_{\mathrm{peak}}=F[V_{\mathrm{irr}}f(\mathrm{}\omega )].$$
(134)
In the experiment the scaled curves coincided with good precision, see Fig. 7 of Ref. for conductances $`G_{\mathrm{peak}}`$ ranging from the unitary limit down to values small compared with $`G_U`$. Our theory indeed allows for such a single-parameter scaling (134), see Eqs. (109) and (112) and Fig. 2, in the case of ac modulation of the gate voltage. On the contrary, according to Fig. 3 scaling of the type (134) would be possible only in the domain of large or small conductances, if the ac field of frequency $`\omega T_K`$ is applied to the leads.
We must mention here that the experimentally measured Kondo conductance was suppressed by the irradiation uniformly across the Coulomb blockade valley, including its middle point, where $`E_d=UE_d`$. In our theory, however, the suppression is not uniform, and vanishes at $`E_d=U/2`$, if the ac field is applied to the gate \[see Eqs. (25), (63), and (64)\]. We see two possible reasons for the discrepancy. First, our conclusion is valid for a weak modulation only. This condition most probably was not satisfied in the experiment performed at $`\mathrm{}\omega T_K`$; at such frequencies strong modulation of the gate potential, $`\delta T_KT_K`$, is required to achieve significant suppression of the Kondo effect. Second, leakage of the ac field to the leads would result in suppression of the Kondo conductance by ac bias. This effect is not so sensitive to a specific value of $`E_d`$, see Section VII.
## X Conclusion
We considered the Kondo conductance of a quantum dot subjected to ac field. We have shown that, despite the essentially non-equilibrium character of the problem, the Kondo temperature $`T_K`$ \[Eq. (55)\] remains the only relevant energy scale. The dc Kondo conductance depends on the ac field only through two dimensionless variables which are the frequency and the amplitude of the ac perturbation divided by $`T_K`$. In terms of these two variables, conductance is a universal function. The form of this function, and the relation of the perturbation amplitude to the “bare” value of the ac field amplitude, depends on the way the ac field is applied.
If the ac field is applied to the gate, then the strength of the perturbation is characterized by the amplitude $`\delta T_K`$ of adiabatic variations of the Kondo temperature, see Eqs. (64)–(66). At low frequencies, $`\mathrm{}\omega <T_K`$, the conductance is close to the unitary limit \[Sec. V, Eq. (103)\]. At higher frequencies, $`\mathrm{}\omega >T_K`$, the ac field suppresses the Kondo effect by means of the decoherence of the dot’s spin (Sec. III, IV). The value of the zero-bias conductance decreases with the increase of the frequency $`\omega `$ of the ac field. The results we obtained for the modulation of the gate voltage are summarized in Sec. VI.
If the ac field is applied across the dot, then the proper variable is the corresponding dimensionless bias between the leads, $`V_{\mathrm{ac}}/T_K`$. A “slow” field, $`\mathrm{}\omega ^{}<\mathrm{max}\{eV_{\mathrm{ac}},T_K\}`$, suppresses the Kondo effect essentially the same way as a finite dc bias does, see Eqs. (117), (118). A “fast” ac field, $`\mathrm{}\omega ^{}>\mathrm{max}\{eV_{\mathrm{ac}},T_K\}`$, affects the Kondo conductance through the decoherence of the dot spin \[Eqs. (113), (114), (115)\]. At the fixed amplitude of the field, the suppression of the Kondo effect diminishes with the increase of the ac field frequency $`\omega ^{}`$.
The ac field also produces satellite peaks in the dependence of the differential dc conductance on the dc bias. However, the satellite maxima in the conductance are inevitably small, see Sec. VIII.
The analysis of the experiment \[Sec. IX\] demonstrates good agreement between our theoretical results and the results of the experiment.
###### Acknowledgements.
The work at the University of Minnesota was supported by NSF Grant DMR 97-31756. LG acknowledges the hospitality of the Delft University of Technology. LG and AK acknowledge also the hospitality of Institute of Theoretical Physics supported by NSF Grant PHY 94-07194 at University of California at Santa Barbara, where a part of the work was performed. The authors are grateful to L.P. Kouwenhoven, D. Goldhaber-Gordon and Y. Meir for useful discussions.
## A ac bias with $`eV_{\mathrm{ac}}\mathrm{}\omega ^{}`$
In this Appendix we describe the crossover between the regimes of “slow” and “fast” ac bias, which occurs at
$$\mathrm{}\omega ^{}eV_{\mathrm{ac}}T_K,$$
(A1)
see Fig. 3. Throughout this Appendix, we will use the finite-order perturbation theory to evaluate the decoherence rate and conductance. The RG technique is abandoned here, since the finite-order perturbation theory is sufficient in the region defined by condition (A1).
### 1 Decoherence by ac bias
Unlike the case $`\gamma ^{}eV_{\mathrm{ac}}/\mathrm{}\omega ^{}1`$, in the crossover region (A1) the decoherence rate is determined also by many-photon processes. Using the Fermi Golden Rule with the Hamiltonian of Eqs. (II B) and (27), we arrive at
$$\frac{\mathrm{}}{\tau ^{}}=\frac{2}{\pi }\frac{\mathrm{\Gamma }_L\mathrm{\Gamma }_R}{\stackrel{~}{E}_d^2}M\left(\gamma ^{}\right)\mathrm{}\omega ^{},$$
(A2)
where $`M(x)xJ_0(x)J_1(x)+x^2\left[J_0(x)\right]^2+x^2\left[J_1(x)\right]^2`$. In the case of “fast” ac field, $`\gamma ^{}eV_{\mathrm{ac}}/\mathrm{}\omega ^{}1`$, equation (A2) reduces to
$$\frac{\mathrm{}}{\tau ^{}}=\frac{1}{\pi }\frac{\mathrm{\Gamma }_L\mathrm{\Gamma }_R}{\stackrel{~}{E}_d^2}\left(\gamma ^{}\right)^2\mathrm{}\omega ^{}.$$
(A3)
The latter formula is similar to Eq. (59) and accounts for single-photon processes only. Equation (A3) is the first term in the perturbation theory series in $`𝒥_{\alpha \alpha ^{}}`$ for $`\mathrm{}/\tau ^{}`$. The summation of leading terms of all orders in $`𝒥_{\alpha \alpha ^{}}`$ can be performed with the help of the RG technique, and yields Eq. (113) \[cf. Eqs. (59) and (63) respectively\].
In the limiting case of “slow” ac bias, $`\gamma ^{}1`$, Eq. (A2) is reduced to
$$\frac{\mathrm{}}{\tau ^{}}=\frac{4}{\pi ^2}eV_{\mathrm{ac}}\frac{\mathrm{\Gamma }_L\mathrm{\Gamma }_R}{\stackrel{~}{E}_d^2},$$
(A4)
cf. Eq. (125).
### 2 Conductance
The conductance in the crossover region (A1) can be evaluated with the third-order perturbation-theory series in $`𝒥_{\alpha \alpha ^{}}(t)`$. Using Eqs. (II B), (27), (31), and (33), we arrive at
$`G_{\mathrm{peak}}^{(3)}(T,V_{\mathrm{dc}})={\displaystyle \frac{3\pi ^2}{2}}{\displaystyle \frac{e^2}{\pi \mathrm{}}}\nu ^3\left[𝒥_{LR}^{(0)}\right]^2\left[𝒥_{RR}^{(0)}+𝒥_{LL}^{(0)}\right]`$ (A5)
$`\times {\displaystyle _{\mathrm{}}^0}dt{\displaystyle \frac{(t)\mathrm{cos}(eV_{\mathrm{dc}}t/\mathrm{})\mathrm{exp}(|t|/\tau )}{\mathrm{sinh}^2(\pi Tt/\mathrm{})+(T/D_0)^2}}\left({\displaystyle \frac{\pi T}{\mathrm{}}}\right)^2`$ (A6)
$`\times \mathrm{cos}\left[\gamma ^{}\mathrm{sin}(\omega ^{}t+\varphi _0)\gamma ^{}\mathrm{sin}\varphi _0\right]`$ (A7)
\[cf. Eq. (40)\]. The expression (A6) accounts for the current induced by both dc and ac biases. To single out the former contribution, which is the true dc conductance, we must average over the phase $`\varphi _0`$. This averaging in the limit of zero temperature and dc bias yields
$`G_{\mathrm{peak}}^{(3)}`$ $`=`$ $`{\displaystyle \frac{3\pi ^2}{2}}{\displaystyle \frac{e^2}{\pi \mathrm{}}}\nu ^3\left[𝒥_{LR}^{(0)}\right]^2\left[𝒥_{RR}^{(0)}+𝒥_{LL}^{(0)}\right]`$ (A9)
$`\times {\displaystyle _{\mathrm{}}^0}dt{\displaystyle \frac{\mathrm{exp}(|t|/\tau )}{\sqrt{t^2+(\mathrm{}/\pi D_0)^2}}}J_0\left[2\gamma ^{}\mathrm{sin}{\displaystyle \frac{\omega ^{}t}{2}}\right].`$
In the limit $`\omega ^{}0`$, equation (A9) yields
$$G_{\mathrm{peak}}^{(3)}\mathrm{ln}(D_0/eV_{\mathrm{ac}}).$$
(A10)
This result is analogous to Eq. (118) at $`eV_{\mathrm{ac}}T_K`$. At small frequencies, $`\omega ^{}\tau <1`$, the corrections to Eq. (A10) are proportional to $`\mathrm{exp}(1/\omega ^{}\tau )`$. At larger frequencies $`\mathrm{}/\tau <\mathrm{}\omega ^{}eV_{\mathrm{ac}}`$, the two leading terms in the expansion of the right-hand side of Eq. (A6) in powers of $`1/\gamma ^{}`$ are
$`G_{\mathrm{peak}}^{(3)}`$ $`=`$ $`{\displaystyle \frac{3\pi ^2}{32}}{\displaystyle \frac{e^2}{\pi \mathrm{}}}\nu ^3\left[𝒥_{LR}^{(0)}\right]^2\left[𝒥_{RR}^{(0)}+𝒥_{LL}^{(0)}\right]`$ (A12)
$`\times \left[\mathrm{ln}{\displaystyle \frac{D_0}{eV_{\mathrm{ac}}}}+{\displaystyle \frac{1}{\gamma ^{}}}\mathrm{ln}\omega ^{}\tau \right].`$
At even larger frequencies, $`\mathrm{}\omega ^{}eV_{\mathrm{ac}}`$, Eq. (A9) yields
$$G_{\mathrm{peak}}^{(3)}\mathrm{ln}(D_0\tau ^{}/\mathrm{}).$$
(A13)
This is the first logarithmic term of the series that are summed up in Eq. (114).
Note that the results given by Eqs. (A10), (A12), and (A13) match each other at the corresponding applicability limits.
## B RG transformation for satellite peaks
In this Appendix, we describe the RG transformation we used to derive Eq. (128). Unlike the other instances of application of the RG technique in our paper, the RG transformation of this Appendix consists of two stages. The first stage is analogous to the one considered in Sec. III A \[Eqs. (60), (62)\]. It stops when the bandwidth $`D`$ reaches $`eV_{\mathrm{dc}}\mathrm{}\omega `$. Since $`\mathrm{}/\tau _{\mathrm{stl}}<\mathrm{}\omega `$ \[see Eq. (126)\], we need to reduce the band further to account for all possible virtual transitions contributing to the Kondo anomaly in $`G_{\mathrm{stl}}(V_{\mathrm{dc}})`$.
The RG transformation we consider in this Appendix is aimed at evaluation of the conductance $`G_{\mathrm{stl}}`$. In processes which contribute to the singularity in $`G_{\mathrm{stl}}`$, a tunneling electron jumps from the Fermi level in one lead to the Fermi level in the other, emitting or absorbing a photon. Therefore, the reduced band in each lead must be centered at its Fermi level. When $`D`$ is below $`\mathrm{}\omega `$, transitions of only two types are possible within such a band. First, there can be transitions within a lead without absorption/emission of a photon. Second, transitions from the higher-potential lead to the lower-potential lead with emission of a photon, and reverse transition with absorption of a photon are also possible. The other types of transitions bring electrons out of the reduced band and should be excluded from the consideration. Such a treatment yields the RG equations
$`{\displaystyle \frac{d𝒥_{\alpha \alpha }}{dD}}`$ $`=`$ $`\nu {\displaystyle \frac{𝒥_{\alpha \alpha }^2}{D}},`$ (B1)
$`{\displaystyle \frac{d𝒥_{LR}}{dD}}`$ $`=`$ $`\nu {\displaystyle \frac{𝒥_{LR}(𝒥_{LL}+𝒥_{RR})}{D}},`$ (B2)
with the initial conditions
$`𝒥_{\alpha \alpha }(D)|_{D\mathrm{}\omega }`$ $`=`$ $`{\displaystyle \frac{2\mathrm{\Gamma }_\alpha }{\mathrm{\Gamma }_L+\mathrm{\Gamma }_R}}{\displaystyle \frac{1}{2\nu \mathrm{ln}(\mathrm{}\omega /T_K)}}.`$ (B3)
$`𝒥_{LR}(D)|_{D\mathrm{}\omega }`$ $`=`$ $`{\displaystyle \frac{2\sqrt{\mathrm{\Gamma }_L\mathrm{\Gamma }_R}}{\mathrm{\Gamma }_L+\mathrm{\Gamma }_R}}{\displaystyle \frac{1}{4\nu [\mathrm{ln}(\mathrm{}\omega /T_K)]^2}}{\displaystyle \frac{\delta T_K}{T_K}}.`$ (B4)
cf. Eq. (62). The second stage of the transformation must be stopped at $`DD^{}\sqrt{(\mathrm{}/\tau _{\mathrm{stl}})^2+(eV_{\mathrm{dc}}\pm \mathrm{}\omega )^2}`$. Expanding the solution for $`𝒥_{LR}`$ in powers of $`\delta T_K/T_K`$ up to the first power, we obtain
$`𝒥_{LR}|_{DD^{}}={\displaystyle \frac{2\sqrt{\mathrm{\Gamma }_L\mathrm{\Gamma }_R}}{\mathrm{\Gamma }_L+\mathrm{\Gamma }_R}}{\displaystyle \frac{\delta T_K}{T_K}}`$ (B5)
$`\times \left[\mathrm{ln}{\displaystyle \frac{\sqrt{(\mathrm{}/\tau _{\mathrm{stl}})^2+(eV_{\mathrm{dc}}\pm \mathrm{}\omega )^2}}{T_K}}+\mathrm{ln}{\displaystyle \frac{\mathrm{}\omega }{T_K}}\right]^2.`$ (B6)
The conductance $`G_{\mathrm{stl}}`$ must be calculated in the second-order perturbation theory in $`𝒥_{LR}(DD^{})`$, given by Eq. (B6). This calculation finally yields Eq. (128).
|
warning/0003/cond-mat0003120.html
|
ar5iv
|
text
|
# Conductivity activation energy for bilayer heterostructures at integer filling factors
## Hamiltonian of 2DEG bilayer
The electronic Hamiltonain of a 2DEG in a confining potential $`V(\stackrel{}{\rho })`$ and in an external magnetic field $`H`$ consists of a one-particle part as well as a Coulomb interaction part:
$$\begin{array}{c}H=\psi _\alpha ^+(\stackrel{}{\rho })\left(\frac{1}{2m}\left[i\stackrel{}{}+\stackrel{}{A}(\stackrel{}{\rho })\right]^2+V(\stackrel{}{\rho })|g|\mu _BH\sigma _{\alpha \beta }^z\right)\psi _\beta (\stackrel{}{\rho })d^3\stackrel{}{\rho }+\\ \frac{1}{2}\frac{e^2}{|\stackrel{}{\rho }\stackrel{}{\rho ^{}}|}\psi _\alpha ^+(\stackrel{}{\rho })\psi _\beta ^+(\stackrel{}{\rho ^{}})\psi _\beta (\stackrel{}{\rho ^{}})\psi _\alpha (\stackrel{}{\rho })d^3\stackrel{}{\rho }d^3\stackrel{}{\rho ^{}},\end{array}$$
where $`\alpha ,\beta =\pm `$ are spin indices and thereafter a sum over repeated indices is implied. We use such units that $`\mathrm{}=1`$, $`e=c`$ and $`H=B=1`$. The latter implies that all distances can be expressed in terms of the so-called magnetic length: $`l_H=\sqrt{c\mathrm{}/eH}=1`$. We split three coordinates $`\stackrel{}{\rho }`$ into a perpendicular to the layer coordinate $`\xi `$ and two in-plane coordinates $`\stackrel{}{r}=(x,y)=(z,\overline{z})`$. We assume that the confining potential is uniform over the plane: $`V(\stackrel{}{\rho })=V(\xi )`$, and represents a double well structure in the transverse direction as shown on the Fig.1, with the two wells being separated by the distance $`d`$. We use only two eigen functions: the lowest energy symmetric $`\chi _S(\xi )`$ and antisymmetric $`\chi _A(\xi )`$, from a set of one-electron eigen functions in the confining potential $`V(\xi )`$ and we expand an electron second-quantized operator in terms of these two eigen functions:
$$\psi _\alpha (\stackrel{}{\rho })=\underset{\tau ,n,p}{}\chi _\tau (\xi )\varphi _{n,p}(\stackrel{}{r})c_{n\alpha \tau p},$$
(1)
where $`c_{\alpha \tau p}^+`$ and $`c_{\alpha \tau p}`$ are electron creation and anhilation operators, $`\varphi _{n,p}(z\overline{z})`$ is an electron wave function number $`p`$ in the Landau gauge in the $`n`$’s Landau Level, the index $`\tau =1,2`$ being the layer index and the layer wave functions read:
$$\chi _{1,2}(\xi )=\frac{\chi _S(\xi )\pm \chi _A(\xi )}{\sqrt{2}}.$$
(2)
We restrict our model to the case of a sufficiently strong magnetic field, such that the cyclotron energy $`\mathrm{}\omega _0`$ dominates over the Coulomb, Zeeman and the level splitting: $`E_AE_S`$, energies. Thus, we specialize to the lowest Landau Level and retain only the term $`n=0`$ in (1).
Plugging the wave function (1) into the Hamiltonian (Hamiltonian of 2DEG bilayer) we find a 2DEG Hamiltonian as:
$$\begin{array}{c}H=\frac{1}{2m}c_{\alpha \tau p}^+c_{\alpha \tau p}c_{\alpha \tau _1p}^+\left(t\tau _{\tau _1\tau _2}^x+\mu ^z\tau _{\tau _1\tau _2}^z\right)c_{\alpha \tau _2p}|g|\mu _BHc_{\alpha \tau p}^+\sigma _{\alpha \beta }^zc_{\beta \tau p}+\\ +\frac{1}{2}\underset{p_1..p_4}{}d^2\stackrel{}{r}d^2\stackrel{}{r^{}}V_{\tau _2\tau _3}^{\tau _1\tau _4}(\stackrel{}{r}\stackrel{}{r^{}})\varphi _{p_1}^{}(\stackrel{}{r})\varphi _{p_2}^{}(\stackrel{}{r^{}})\varphi _{p_3}(\stackrel{}{r^{}})\varphi _{p_4}(\stackrel{}{r})\\ c_{\alpha \tau _1p_1}^+c_{\beta \tau _2p_2}^+c_{\beta \tau _3p_3}c_{\alpha \tau _4p_4},\end{array}$$
(3)
where we have defined a hopping constant:
$$t=\frac{1}{2}d^2\stackrel{}{r}𝑑\xi \varphi _p^{}(\stackrel{}{r})\chi _{\tau _1}(\xi )\tau _{\tau _1\tau _2}^xV(\xi )\chi _{\tau _2}(\xi )\varphi _p(\stackrel{}{r}),$$
(4)
as well as an external electrostatic potential created by an asymmetric gate charge:
$$\mu ^z=\frac{1}{2}d^2\stackrel{}{r}𝑑\xi \varphi _p^{}(\stackrel{}{r})\chi _{\tau _1}(\xi )\tau _{\tau _1\tau _2}^zV(\xi )\chi _{\tau _2}(\xi )\varphi _p(\stackrel{}{r}),$$
(5)
whereas the Coulomb interaction matrix reads:
$$V_{\tau _2\tau _3}^{\tau _1\tau _4}(\stackrel{}{r}\stackrel{}{r^{}})=\frac{\chi _{\tau _1}(\xi )\chi _{\tau _2}(\xi ^{})\chi _{\tau _3}(\xi ^{})\chi _{\tau _4}(\xi )}{\sqrt{(\xi \xi ^{})^2+\left(\stackrel{}{r}\stackrel{}{r^{}}\right)^2}}𝑑\xi 𝑑\xi ^{}.$$
(6)
We use notations: $`\tau ^x`$, $`\tau ^y`$ and $`\tau ^z`$, for the Pauli matrices in the layer space whereas we use notations: $`\sigma ^x`$, $`\sigma ^y`$ and $`\sigma ^z`$, for the Pauli matrices in the spin space. The hopping constant can be related to the splitting of the symmetric and the antisymmetric levels: $`t=E_AE_S`$. The electrostatic potential $`\mu ^z`$, which can be viewed as a difference between the chemical potentials in the two layers, breaks down the symmetry between the two wells of $`V(\xi )`$ potential. This term appears naturally when a single gate is fabricated to control the electron density in the bilayer. In the limit $`d0`$, $`\mu ^z`$ vanishes too, whereas in the limit of large layer separation: $`d\mathrm{}`$, $`\mu ^z\mathrm{}`$ and electrons reside only on the layer adjacent to the gate. We assume that the energy of a capacitor formed by the two layers is much lower than the characteristic Coulomb energy $`e^2/\kappa l_H^3`$, per area, where $`\kappa `$ is the dielectric constant. We note the invariance of the Coulomb energy: (6) under the following transformations: $`\tau _1\tau _4`$, $`\tau _2\tau _3`$ as well as $`(\tau _1\tau _4)(\tau _2\tau _3)`$. To fully exploit these symmetries we cast the Eq.6 into a more suitable representation:
$$V_{\tau _2\tau _3}^{\tau _1\tau _4}(\stackrel{}{r}\stackrel{}{r^{}})=V^{\mu \nu }(\stackrel{}{r}\stackrel{}{r^{}})\tau _{\tau _1\tau _4}^\mu \tau _{\tau _2\tau _3}^\nu ,$$
(7)
where $`\tau ^0`$ is the unit matrix, $`V^{\mu \nu }`$ is a $`3\times 3`$ symmetric interaction matrix with indices $`\mu `$, $`\nu `$ running over a set $`(0,z,x)`$. If there is a symmetry of the Coulomb interaction under an exchange of layers: $`(\xi \xi ^{})(\xi \xi ^{})`$ and $`12`$ then it restricts further values of the interaction matrix: $`V^{0z}=0`$ and $`V^{zx}=0`$. But in the presence of a gate asymmetry we shall keep the matrix element: $`V^{0z}`$. Therefore, the Coulomb interaction matrix for symmetric bilayer 2DEG depends on four parameters: $`V^{00}>0`$, $`V^{0x}`$, $`V^{xx}>0`$, $`V^{zz}>0`$. We note also, that $`V^{0x}\chi `$, $`V^{xx}\chi ^2`$, whereas $`V^{zz}d^2/|z|^3`$ as $`|z|\mathrm{}`$. In the following we shall neglect $`V^{xx}`$ matrix element.
Next, we split the total bilayer Hamiltonian (3) into two parts: the first one contains a dominant Coulomb energy term:
$$\begin{array}{c}H^{sym}=\frac{1}{2m}c_{\alpha \tau p}^+c_{\alpha \tau p}+\frac{1}{2}\frac{d^2\stackrel{}{q}}{(2\pi )^2}V^{00}(\stackrel{}{q})N(\stackrel{}{q})N(\stackrel{}{q}),\end{array}$$
(8)
where $`V^{\mu \nu }(\stackrel{}{q})`$ is the Fourier transform of $`V^{\mu \nu }`$ $`(\stackrel{}{r})`$ multiplied by a factor $`\mathrm{exp}(\stackrel{}{q}^2/4)`$ and the electron density operator reads:
$$N(\stackrel{}{q})=\underset{p}{}c_{\alpha \tau p}^+c_{\alpha \tau pq_y}\mathrm{exp}iq_x(p\frac{q_y}{2})$$
(9)
This part of the Hamiltonian is invariant under uniform rotations from the $`SU(4)`$ Lee group in the combined spin and layer space. Every of its eigen energy is hugely degenerate. Given any eigen state $`|\mathrm{\Psi }_0`$ a set of related eigen states can be generated by applying rotations: $`|\mathrm{\Psi }=U|\mathrm{\Psi }_0`$, where $`USU(4)`$. For Landau level filling factor $`\nu =1`$, $`\nu =2`$ and $`\nu =3`$ we assume that the the bilayer ground state is uniform over $`p`$-orbitals:
$$\mathrm{\Psi }=\underset{i=1}{\overset{\nu }{}}\underset{p}{}c_{\alpha _i\tau _ip}^+|empty,$$
(10)
and we prove in the next section that this state is stable with respect to long-range spatial perturbations. One can easily check by inspection that any such wave-function (10) represents an eigen function of the $`H^{sym}`$ (8). The remaining few terms in the Hamiltonian (3) are treated like perturbations:
$$\begin{array}{c}H^{anis}=c_{\alpha \tau _1p}^+\left(t\tau _{\tau _1\tau _2}^x+\mu ^z\tau _{\tau _1\tau _2}^z\right)c_{\alpha \tau _2p}|g|\mu _BHc_{\alpha \tau p}^+\sigma _{\alpha \beta }^zc_{\beta \tau p}+\\ +\frac{1}{2}\frac{d^2\stackrel{}{q}}{(2\pi )^2}V^{\mu \nu }(\stackrel{}{q})T^\mu (\stackrel{}{q})T^\nu (\stackrel{}{q}),\end{array}$$
where (see e.g. )
$$T^\mu (\stackrel{}{q})=\underset{p}{}c_{\alpha \tau _1p}^+\tau _{\tau _1\tau _2}^\mu c_{\alpha \tau _2pq_y}\mathrm{exp}iq_x(p\frac{q_y}{2})$$
(11)
with $`(\mu \nu )(00)`$. The Hamiltonian (Hamiltonian of 2DEG bilayer) breaks down the $`SU(4)`$ symmetry but it is still invariant under separate rotations in the spin and layer space: $`SU(2)SU(2)`$. We shall call this part of the Hamiltonian the anisotropy Hamiltonian. It lifts the degeneracy of eigen states of the $`SU(4)`$-symmetric Hamiltonian (8). An important point to note here is that a splitting of energy levels is determined by matrix elements of weak anisotropy Hamiltonian (Hamiltonian of 2DEG bilayer) truncated to a linear space of the symmetric Hamiltonian (8) level degeneracy. There are no Fermi-liquid type renormalizations of the constants of the anisotropy Hamiltonian (Hamiltonian of 2DEG bilayer) due to the $`SU(4)`$-symmetric Hamiltonian (8). In other word the mean-field Hartree-Fock approach is perfect for the $`\nu =1`$, $`\nu =2`$ and $`\nu =3`$ cases.
Our guiding analogy in treating the total bilayer Hamiltonian (8, Hamiltonian of 2DEG bilayer) lies in the theory of magnetism. We will see below that there exists a local order parameter: $`Q`$, very much like magnetization. And we aim to express the total bilayer Hamiltonian (8,Hamiltonian of 2DEG bilayer) in terms of this order parameter $`Q`$. The exchange-like Hamiltonian (8) has to be expanded in powers of spatial variations of order parameter $`Q(\stackrel{}{r})`$ with the second power of gradients being the important contribution, whereas only locally homogeneous $`Q`$ has to be retain in the anisotropy Hamiltonian (Hamiltonian of 2DEG bilayer). In the next section we carry out the first step whereas in the next-to-next section we transform the anisotropy energy.
## SU(4) Symmetric Case
In this section we specialize to the $`SU(4)`$-symmetric part of the bilayer 2DEG Hamiltonian (8) which is invariant under the global rotations of a four component electron spinor by the $`4\times 4`$ matrix $`U`$ from the $`SU(4)`$ Lee group. We find it useful the ground state of operators $`c`$ to be the reference state. Any non-homogeneous state is then generated by rotation: $`U(t,\stackrel{}{r})`$. And the action of bilayer is some functional of it:
$$S[U(t,\stackrel{}{r})]=i\text{tr}\mathrm{log}𝒟c^+(t)𝒟c(t)\mathrm{exp}\left(i𝑑t\right)$$
(12)
where the symmetric Lagrangian of bilayer 2DEG reads:
$$\begin{array}{c}=c_{\alpha \tau _1}^+\left[i\frac{}{t}\frac{1}{2m}\left(i\stackrel{}{}+\stackrel{}{A}_0+\stackrel{}{\mathrm{\Omega }}_{\alpha \beta \tau _1\tau _2}\right)^2\right]c_{\beta \tau _2}d^2\stackrel{}{r}+\\ +c_{\alpha \tau _1}^+\mathrm{\Omega }_{\alpha \beta \tau _1\tau _2}^tc_{\beta \tau _2}d^2\stackrel{}{r}\frac{1}{2}\frac{d^2\stackrel{}{q}}{(2\pi )^2}V^{00}(\stackrel{}{q})N(\stackrel{}{q})N(\stackrel{}{q}),\end{array}$$
(13)
where $`\stackrel{}{A}_0=(0,A^y)`$ is the vector potential of the external magnetic field and
$$\mathrm{\Omega }^t=iU^+\frac{}{t}U,\stackrel{}{\mathrm{\Omega }}=iU^+\stackrel{}{}U.$$
(14)
In the limit of slow spatial variations of rotation $`U(t,\stackrel{}{r})`$ this functional has an expansion in powers of $`\mathrm{\Omega }`$. In this context the functional (12) is called an effective low-energy Goldstone Action, and we are going to find it in this section. All calculations follows step in step those done in Ref. for the case of a single layer, and here we emphasized only the points of difference.
Calculations of the Goldstone Action can be carried through for three filling factors of the bilayer 2DEG: $`\nu =1`$, $`\nu =2`$ and $`\nu =3`$, at once. We start with choosing the reference state: i) $`n=1`$, one electron with spin up fills every orbital $`p`$ of the lowest Landau Level of the first layer: $`|+_1`$; ii) $`n=2`$, two electrons - one with spin up on the first layer and the other with spin down on the second layer - fill every orbital $`p`$ the lowest Landau Level: $`|+_1|_2`$; iii) three electrons \- one with spin down on the first layer and the two other with spin up and down on the second layer - fill every orbital $`p`$ the lowest Landau Level: $`|_1|+_2|_2`$. The case iii) reduces to the case i) if one makes the electron-hole transformation.
The one-electron Green functions defined for the reference state of the Hamiltonian (13) in the homogeneous limit: $`\stackrel{}{\mathrm{\Omega }}=0`$, reads:
$$\begin{array}{c}G_{\alpha \tau _1,\beta \tau _2}^0(ϵ)=\frac{1}{2}\left(\mathrm{\Sigma }_{\alpha \beta ,\tau _1\tau _2}^0+\mathrm{\Sigma }_{\alpha \beta ,\tau _1\tau _2}^z\right)\frac{1}{ϵ+E_0\mu i0}+\\ +\frac{1}{2}\left(\mathrm{\Sigma }_{\alpha \beta ,\tau _1\tau _2}^0\mathrm{\Sigma }_{\alpha \beta ,\tau _1\tau _2}^z\right)\frac{1}{ϵ\mu +i0},\end{array}$$
(15)
where $`\mu `$ is the chemical potential,
$$\mathrm{\Sigma }_{\alpha \beta ,\tau _1\tau _2}^0=\sigma _{\alpha \beta }^0\tau _{\tau _1\tau _2}^0=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& 1\end{array}\right),$$
(16)
and in the case $`n=1,3`$:
$$\mathrm{\Sigma }_{\alpha \beta ,\tau _1\tau _2}^z=\pm \left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& 1\end{array}\right),$$
(17)
whereas in the case $`n=2`$:
$$\mathrm{\Sigma }_{\alpha \beta ,\tau _1\tau _2}^z=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& 1\end{array}\right).$$
(18)
The effective action can be written as:
$$S=S_0+S_2,$$
(19)
where
$$S_0=i\text{tr}\mathrm{log}\frac{G}{G_0}$$
(20)
and $`S_2`$ is conveniently represented in terms of diagrams in the Fig.2. As it is explained in the Ref. the first order perturbation correction to the Green function: $`\delta G=GG_0`$, contains an electron propagation on the first excited Landau Level. Thus, we also need the bare Green function of an excited electron on this level:
$$\begin{array}{c}G_{\alpha \tau _1,\beta \tau _2}^1(ϵ)=\frac{1}{2}\left(\mathrm{\Sigma }_{\alpha \beta ,\tau _1\tau _2}^0+\mathrm{\Sigma }_{\alpha \beta ,\tau _1\tau _2}^z\right)\frac{1}{ϵ1/m+E_1\mu +i0}+\\ +\frac{1}{2}\left(\mathrm{\Sigma }_{\alpha \beta ,\tau _1\tau _2}^0\mathrm{\Sigma }_{\alpha \beta ,\tau _1\tau _2}^z\right)\frac{1}{ϵ1/m\mu +i0},\end{array}$$
(21)
where
$$E_0=2E_1=\sqrt{\frac{\pi }{2}}\frac{e^2}{\kappa l_H}.$$
(22)
As it is explained in the Ref. the gradient vector field enters the following terms in the Hamiltonian
$$H_1=\frac{1}{2m}c_{\alpha \tau _1}^+\left(\mathrm{\Omega }_{\alpha \beta ,\tau _1\tau _2}^+\widehat{\mathrm{\Pi }}_{}+\mathrm{\Omega }_{\alpha \beta ,\tau _1\tau _2}^{}\widehat{\mathrm{\Pi }}_+\right)c_{\beta \tau _2}d^2\stackrel{}{r},$$
(23)
$$H_2=\frac{1}{2m}c_{\alpha \tau _1}^+\left[\left(\stackrel{}{\mathrm{\Omega }}^2\right)_{\alpha \beta ,\tau _1\tau _2}i\left(_\mu \mathrm{\Omega }_\mu \right)_{\alpha \beta ,\tau _1\tau _2}\right]c_{\beta \tau _2}d^2\stackrel{}{r},$$
(24)
where the operators $`\widehat{\mathrm{\Pi }}_\pm `$ shift an electron between the adjacent Landau Levels:
$$\widehat{\mathrm{\Pi }}_{}\varphi _{np}(\stackrel{}{r})=\sqrt{2n}\varphi _{n1p}(\stackrel{}{r}),\widehat{\mathrm{\Pi }}_+\varphi _{n1p}(\stackrel{}{r})=\sqrt{2n}\varphi _{np}(\stackrel{}{r}),$$
(25)
though only the $`n=0`$ and $`n=1`$ Landau Level states are relevant for our problem. The two new gradient vector fields in the Hamiltonian (23) are defined as follows:
$$\mathrm{\Omega }_\pm =iU^+\left(_yi_x\right)U,$$
(26)
These two components are real: $`\mathrm{\Omega }_{}=\mathrm{\Omega }_+^{}`$. An expansion of the 2DEG action up to the second power of the Hamiltonian (23) reads:
$$\delta S_0=i\text{tr}\left(H_1G_0\right)+\frac{i}{2}\text{tr}\left(H_1G_0H_1G_0\right)+i\text{tr}\left(H_2G_0\right).$$
(27)
Combining this expansion with diagrams of the Fig.2 we find the low-energy Goldstone Hamiltonian as follows:
$$H_G=\frac{E_1}{8}\frac{d^2\stackrel{}{r}}{2\pi }\text{tr}\left((\mathrm{\Sigma }^0\mathrm{\Sigma }^z)\mathrm{\Omega }_{}(\stackrel{}{r})(\mathrm{\Sigma }^0+\mathrm{\Sigma }^z)\mathrm{\Omega }_+(\stackrel{}{r})\right)$$
(28)
The insertion matrices in (28) are non-negative diagonal ones and they represent the occupation number for the electron states:
$$N=\frac{1}{2}\left(\mathrm{\Sigma }^0+\mathrm{\Sigma }^z\right)=\left(\begin{array}{cc}1& 0\\ & \\ 0& 0\end{array}\right)$$
(29)
and
$$\mathrm{\Sigma }^0N=\frac{1}{2}\left(\mathrm{\Sigma }^0\mathrm{\Sigma }^z\right)=\left(\begin{array}{cc}0& 0\\ & \\ 0& 1\end{array}\right),$$
(30)
where blocks are $`1\times 1`$ and $`3\times 3`$ in the case $`\nu =1,3`$ and $`2\times 2`$ in the case $`\nu =2`$. It follows immediately that $`H_G0`$.
The matrices $`N`$ and $`1N`$ can be viewed as projector operators that allow only those rotations into the Hamiltonian (28) that do change the ground state. It is useful to understand what particular sub-set of $`SU(4)`$ Lee group these physical rotations form. The vector field $`\mathrm{\Omega }_\mu `$ can be expanded in the basis of fifteen generators of $`SU(4)`$ Lee group: $`\{\mathrm{\Sigma }^l\}`$, with $`l=\mathrm{1..15}`$. And we subdivide them into two complementary sets: the first one includes those generators that do commute with occupation number matrices (29, 30), and we shall called it an $`(even)`$ set, whereas the second one includes the remaining generators, and we shall call it the $`(odd)`$ set. Generators of the $`(even)`$ set constitute an algebra itself. This algebra has a normal Abelian subalgebra formed by a single generator: $`\mathrm{\Sigma }^z\text{tr}\mathrm{\Sigma }^z/4`$. A Lee group built around the $`(even)`$ set of generators is called a stabilizer sub-group $`S`$ of $`SU(4)`$ Lee group. The $`(odd)`$ set always contains an even number of generators. Specifically, eight in the case of $`\nu =2`$ and six in the case $`\nu =1,3`$.
The Hamiltonian (28) must be invariant under the time reversal symmetry. The time reversal operator can be chosen as a complex conjugate operator: $`UU^{}`$. It follows that $`\stackrel{}{\mathrm{\Omega }}\stackrel{}{\mathrm{\Omega }}^T`$. Now it is evident that $`\mathrm{\Omega }_{}`$ and $`\mathrm{\Omega }_+`$ get interchanged under the time reversal in (28), and this changes the Hamiltonian. But we remember that the time reversal is always accompanied by inverting the magnetic field $`B^z`$ and, thus, we can restore the time reversal symmetry by multiplying the antisymmetric part of the Hamiltonian $`H_G`$ by the sign of the magnetic field:
$$H_G=\frac{E_1}{2}\frac{d^2\stackrel{}{r}}{2\pi }\left[\text{tr}\left((\mathrm{\Sigma }^0N)\mathrm{\Omega }_\mu N\mathrm{\Omega }_\mu \right)+i\text{sgn}(B^z)ϵ_{\mu \nu }\text{tr}\left(\mathrm{\Omega }_\mu \mathrm{\Omega }_\nu N\right)\right]$$
(31)
Now if we define a useful gradient vector field:
$$\mathrm{\Omega }_\mu ^z(\stackrel{}{r})=\frac{i}{2}\text{tr}\left(U^+(\stackrel{}{r})\mathrm{\Sigma }^z_\mu U(\stackrel{}{r})\right),$$
(32)
then it is straightforward to rewrite the Eq.(31):
$$H_G=\frac{E_1}{2}\frac{d^2\stackrel{}{r}}{2\pi }\left[\text{tr}\left((\mathrm{\Sigma }^0N)\mathrm{\Omega }_\mu N\mathrm{\Omega }_\mu \right)+\text{sgn}(B^z)\text{curl}\mathrm{\Omega }^z\right],$$
(33)
The first term in $`H_G`$ (33) is the gradient energy whereas the second term is proportional to the topological index of an excited state:
$$𝒬=\text{curl}\mathrm{\Omega }^z\frac{d^2\stackrel{}{r}}{2\pi }=Z,$$
(34)
where $`Z`$ is the set of integer numbers. The case $`𝒬=\pm 1`$ corresponds to the simplest spin skyrmion in the first layer being rotated by a $`SU(4)`$ matrix to become a general bilayer skyrmion. The energy constant in $`H_G`$ (33) coincide identically with that of the one-layer case , which means that the bilayer skyrmion energy is the same as found for one layer. But there is an important difference between a bilayer skyrmion and a spin skyrmion in a single ferromagnetic layer. Namely, as was shown in the Ref. any skyrmion carries a charge density in the core:
$$n(\stackrel{}{r})=\frac{\text{curl}\mathrm{\Omega }^z(\stackrel{}{r})}{2\pi }=\frac{R^2}{\pi \left(R^2+\stackrel{}{r}^2\right)^2},$$
(35)
where $`R`$ is the radius of the skyrmion’s core, contrary to that in the bilayer a charge of the skyrmion’s core can be delocalized over the two layers with a much longer tails:
$$n_{1,2}(\stackrel{}{r})\pm \frac{1}{\left(R^2+\stackrel{}{r}^2\right)}.$$
(36)
The total charge of two layers does converge according to (35) at large distances from the skyrmion’s core.
The Goldstone Hamiltonian (33) can be cast in a special order parameter representation. To do this we define a non-homogeneous order parameter matrix $`Q`$ as follows:
$$Q(\stackrel{}{r})=U(\stackrel{}{r})NU^+(\stackrel{}{r}).$$
(37)
This electronic order parameter has an important property:
$$A=\text{tr}(AQ),$$
(38)
where $`A`$ is any operator. Inspecting the particular difinition of $`N`$ (29) it becomes evident that rotations from the denominator sub-group $`S`$ leaves the order parameter intact. Thus, rotations in (37) can be restricted to a coset or, in other word, a physical space of the bilayer 2DEG:
$$\frac{U(4)}{U(\nu )U(4\nu )}.$$
(39)
Now it a straightforward calculation to rewrite $`H_G`$ (33) in terms of the order parameter matrix:
$$H=\frac{E_1}{4}\text{tr}\left(\stackrel{}{}Q\stackrel{}{}Q\right)\frac{d^2\stackrel{}{r}}{2\pi }+\text{sgn}(B^z)\frac{E_1}{2}ϵ_{\mu \nu }\text{tr}\left(Q_\mu Q_\nu Q\right)\frac{d^2\stackrel{}{r}}{2\pi }.$$
(40)
In this representation the topological index appears as an index of a map of the order parameter coset space into a 2D plane. The index selection rule (34) is a consequence of a well known homotopy group identity:
$$𝒬=\pi _2\left(\frac{U(4)}{U(\nu )U(4\nu )}\right)=Z.$$
(41)
In the end we have to include the Coulomb energy of charge distribution inside the skyrmion core:
$$\delta H_G=\frac{1}{2}d^2\stackrel{}{r}d^2\stackrel{}{r^{}}\frac{\text{curl}\mathrm{\Omega }^z(\stackrel{}{r})}{2\pi }\frac{e^2}{|\stackrel{}{r}\stackrel{}{r^{}}|}\frac{\text{curl}\mathrm{\Omega }^z(\stackrel{}{r}^{})}{2\pi }.$$
(42)
## Anisotropic Part of Coulomb Energy. Phase Diagram.
In this section we cast the anisotropic part of the bilayer Hamiltonian (Hamiltonian of 2DEG bilayer) in terms of the order parameter matrix $`Q`$. It can be conveniently done by the following Harteree-Fock average of $`c`$-operator product in (Hamiltonian of 2DEG bilayer):
$$\begin{array}{c}\tau _{\tau _1\tau _4}^\mu \tau _{\tau _2\tau _3}^\nu <c_{\alpha \tau _1p_1}^+c_{\beta \tau _2p_2}^+c_{\beta \tau _3p_3}c_{\alpha \tau _4p_4}>=\\ =\delta _{p_1p_4}\delta _{p_2p_3}\text{tr}(Q\tau ^\mu )\text{tr}(Q\tau ^\nu )\delta _{p_1p_3}\delta _{p_2p_4}\text{tr}(Q\tau ^\mu Q\tau ^\nu ),\end{array}$$
where $`\tau ^\mu `$ acts on four-spinor as $`\tau ^\mu \sigma ^0`$. Next, we define the following Coulomb anisotropy constants:
$$E^{ab}=\frac{dzd\overline{z}}{2\pi l_H^2}V^{ab}(|z|)\mathrm{exp}\frac{|z|^2}{2l_H^2}\frac{dzd\overline{z}}{2\pi l_H^2}V^{ab}(|z|),$$
(43)
where the last approximation holds for $`(ab)(00)`$ in the limit $`dl_H`$. And, finally, we rewrite the anisotropy Hamiltonian (Hamiltonian of 2DEG bilayer) in terms of order parameter matrix $`Q`$:
$$\begin{array}{c}H^{anis}/𝒩=\left(t+(\nu 1)E^{0x}\right)\text{tr}\left(Q\tau ^x\right)\left(\mu ^z+(\nu 1)E^{0z}\right)\text{tr}\left(Q\tau ^z\right)\\ |g|\mu _BH\text{tr}\left(Q\sigma ^z\right)+\frac{1}{2}E^{zz}\left[\text{tr}\left(Q\tau ^z\right)\text{tr}\left(Q\tau ^z\right)\text{tr}\left(Q\tau ^zQ\tau ^z\right)\right],\end{array}$$
(44)
where $`𝒩`$ is the number of degeneracy of the Landau Level. The Eqs.(40, 44) defines the effective long-range Hamiltonian of a bilayer at integer filling factors. At non-zero temperatures thermal fluctuations of the order parameter soften the anisotropy constants in the Hamiltonian (44). The relevant calculation can be found in eg. and the result reads:
$$\begin{array}{c}\left(t+(\nu 1)E^{0x}\right)_R=\left(t+(\nu 1)E^{0x}\right)\left(\frac{l_h}{R^{}}\right)^{8T/E_1}\\ \left(t+(\nu 1)E^{0z}\right)_R=\left(t+(\nu 1)E^{0z}\right)\left(\frac{l_h}{R^{}}\right)^{8T/E_1}\\ (|g|\mu _BH)_R=|g|\mu _BH\left(\frac{l_h}{R^{}}\right)^{8T/E_1}\\ E_R^{zz}=E^{zz}\left(\frac{l_h}{R^{}}\right)^{24T/E_1},\end{array}$$
(45)
where the spatial scale $`R_{}^2=l_H^2E_1/\text{max}(t,\mu ^z,|g|\mu _BH,E^{zz})`$ indicates the excitation wavelength where its anisotropy energy starts to compete with its exchange energy. Note that the three first constants renormalize as an external field whereas $`E^{zz}`$ constant renormalizes as an easy-axis anisotropy. Although the Coulomb energy $`E_1100KT1K`$ in most experiments, the specific number: $`24=3\times 8`$, which is related to the order of anisotropy and to the eight degree of freedom for thermal fluctuations in the case of $`SU(4)`$ symmetry, makes the renormalization of the constant $`E^{zz}`$ noticeable.
As we have seen in the previous section the order parameter can be parameterized by six or eight angles in the case $`\nu =1,3`$ or $`\nu =2`$. Actually, not every of those rotations corresponds to a physically distinct eigen state. We restrict the calculation of the total bilayer energy up to a first order in powers of the anisotropy Hamiltonian, which means that we shall need only its diagonal matrix elements. But, these are real matrix elements of course, despite the fact that in an external magnetic field there is no time reversal symmetry. Hence, if the Hamiltonian is real one so real has to be its ground state. One generates all real eigen states from a reference state by rotations from the $`SO(4)`$ sub-group of the $`SU(4)`$ group. This group has 6 parameters with two of them falling into the denominator sub-group. Thus, the ground state differs from the reference state by just four rotations. One can view locally the 8D manifold of order parameter as a composition of four unit vectors: magnetization of the first and the second layers and the two hopping-tau vectors which represent the distribution of spin-up and spin-down electron density over the two layers. Now the first two term in the Hamiltonian (Hamiltonian of 2DEG bilayer) are external fields acting on these four vectors. On the other hand the Coulomb energy couples pairs of tau vectors via an exchange interaction. This instructive picture allows us to identify only three special global rotations that do change the total bilayer energy. We start with the case $`\nu =2`$ and we use a set of trial many electron wave functions parameterized by the three angles of relevant in our case rotations: $`\theta _\pm `$ and $`\vartheta `$,:
$$\underset{p}{}U(\vartheta ,\vartheta )R(\theta _+,\theta _{})c_{+1p}^+c_{2p}^+|empty,$$
(46)
where spins in the layer 1,2 are first rotated by $`\pm \vartheta `$:
$$U_{\beta \tau _2}^{\alpha \tau _1}(\vartheta ,\vartheta )=\left(\frac{\tau ^0+\tau ^z}{2}\right)_{\tau _1\tau _2}\mathrm{exp}(i\frac{\vartheta }{2}\sigma ^y)_{\alpha \beta }+\left(\frac{\tau ^0\tau ^z}{2}\right)_{\tau _1\tau _2}\mathrm{exp}(i\frac{\vartheta }{2}\sigma ^y)_{\alpha \beta },$$
(47)
and then wave functions of electrons with spin $`\pm `$ spill over the two layers, the process described by two distinct angles: $`\theta _\pm `$,:
$$R_{\beta \tau _2}^{\alpha \tau _1}(\theta _+,\theta _{})=\left(\frac{\sigma ^0+\sigma ^z}{2}\right)_{\alpha \beta }\mathrm{exp}(i\frac{\theta _+}{2}\tau ^y)_{\tau _1\tau _2}+\left(\frac{\sigma ^0\sigma ^z}{2}\right)_{\alpha \beta }\mathrm{exp}(i\frac{\theta _{}}{2}\tau ^y)_{\tau _1\tau _2}.$$
(48)
This set includes the singlet-liquid state at $`\theta _\pm =\pi /2`$ and $`\vartheta =0`$ and the canted antiferromagnetic state at $`\theta _\pm =0`$. The order parameter reads:
$$Q=URNR^+U^+,$$
(49)
with $`N`$ being the electron density calculated with the reference state of the previous section (see (29)). Now we substitute (49) into the anisotropic Hamiltonian (44) to find the the total anisotropy bilayer energy as:
$$\begin{array}{c}E^{anis}=E^{zz}\mathrm{cos}\theta _+\mathrm{cos}\theta _{}(t+E^{0x})\mathrm{cos}\vartheta (\mathrm{sin}\theta _++\mathrm{sin}\theta _{})\\ (\mu ^z+E^{0z})(\mathrm{cos}\theta _+\mathrm{cos}\theta _{})|g|\mu _BH\mathrm{sin}\vartheta (\mathrm{cos}\theta _++\mathrm{cos}\theta _{}),\end{array}$$
(50)
The minimum of this energy corresponds to three phases: a) ferromagnetic $`\vartheta =\pi /2`$, $`\theta _+=\theta _{}=0`$; b) spin singlet $`\vartheta =0`$, $`\theta _+=\pi \theta _{}=\theta `$; and c) canted antiferromagnetic state otherwise, as it is shown on the Fig.3. It is identical to that found in the Ref.. A line of continuous phase transitions between the ferromagnetic phase and the canted antiferromagnetic phase is given by the following equation:
$$\left[\left(E^{zz}+|g|\mu _BH\right)^2\left(\mu ^z+E^{0z}\right)^2\right]|g|\mu _BH=\left(t+E^{0x}\right)^2\left(E^{zz}+|g|\mu _BH\right)$$
(51)
In the spin singlet phase the interlayer mixing phase: $`\theta `$, is determined by the equation:
$$\left(E^{zz}\mathrm{sin}\theta +t+E^{0x}\right)\mathrm{cos}\theta =\left(\mu ^z+E^{0z}\right)\mathrm{sin}\theta .$$
(52)
A phase transition line that separate the spin singlet phase from the canted antiferromagnetic phase is given parametrically by the equation:
$$\left((t+E^{0x})\mathrm{sin}\theta E^{zz}+\left(\mu ^z+E^{0z}\right)\mathrm{cos}\theta \right)\left(t+E^{0x}\right)=\left(|g|\mu _BH\right)^2\mathrm{sin}\theta ,$$
(53)
where $`\theta `$ is determined from (52). This phase transition is a continuous one also.
In the case $`\nu =1`$ there is no Coulomb interaction energy and the total bilayer energy reads:
$$E^{anis}=t\mathrm{sin}\theta \mu ^z\mathrm{cos}\theta |g|\mu _BH\mathrm{cos}\vartheta ,$$
(54)
The minimum of this energy is given by electron spin being directed along the magnetic field: $`\vartheta =0`$, whereas $`\theta =\mathrm{tan}^1t/\mu ^z`$. There is no phase transition in the case $`\nu =1`$ and the only phase can be characterized as ferromagnetic in both the spin and the layer spaces.
The case $`\nu =3`$ formally reduces to the case $`\nu =1`$ although here the Coulomb interaction energy does not vanish identically. We find renormalizations to the one-particle electron Hamiltonian whereas the total energy being similar to the case $`\nu =1`$:
$$E^{anis}=(t+2E^{0x})\mathrm{sin}\theta (\mu ^z+2E^{0z})\mathrm{cos}\theta |g|\mu _BH\mathrm{cos}\vartheta ,$$
(55)
There is no phase transition in this case either.
## Anisotropic energy gap of one Skyrmion
In this section we find an anisotropic part of the total skyrmion gap energy. The non-homogeneous order parameter that represents one skyrmion is given by the Belavin-Polyakov (BP) skyrmion solution for $`|𝒬|=1`$ :
$$Q_{BP}(z\overline{z})=\frac{R^2}{R^2+|z|^2}\left(\begin{array}{cc}|z|^2& zR\\ & \\ \overline{z}R& R^2\end{array}\right)$$
(56)
where only the shown above four matrix elements differs from those in the electron density matrix $`N`$ (29). The skyrmion order parameter has to be rotated by a homogeneous matrix $`RU`$ calculated in the previous section in such a way that the order parameter far away from the skyrmion center minimizes the anisotropy energy. In addition to this rotation we have to allow all homogeneous rotations from the denominator sub-group $`S`$ that actually transform the BP skyrmion order parameter (56): $`W`$. Thus, a general skyrmion order parameter reads:
$$Q(\stackrel{}{r})=RUWQ_{BP}(z\overline{z})W^+U^+R^+.$$
(57)
First, we consider the case $`\nu =2`$. Here the matrix $`W`$ is parameterized by seven angles:
$$\left(\begin{array}{cccc}\mathrm{cos}\beta _fe^{i(\gamma _f+\alpha _f)}& \mathrm{sin}\beta _fe^{i(\gamma _f\alpha _f)}& 0& 0\\ \mathrm{sin}\beta _fe^{i(\gamma _f+\alpha _f)}& \mathrm{cos}\beta _fe^{i(\gamma _f\alpha _f)}& 0& 0\\ 0& 0& \mathrm{cos}\beta _ee^{i(\gamma _e+\alpha _e)}& \mathrm{sin}\beta _ee^{i(\gamma _e\alpha _e)}\\ 0& 0& \mathrm{sin}\beta _ee^{i(\gamma _e+\alpha _e)}& \mathrm{cos}\beta _ee^{i(\gamma _e\alpha _e)}\end{array}\right)$$
(58)
The additional seventh parameter angle of the denominator group just rotates the coordinates: $`ze^{i\gamma _7}z`$. We find by explicit calculation that the skyrmion anisotropic energy does not depend on parameters $`\gamma _e`$, $`\gamma _f`$ and $`\gamma _7`$ whereas $`\alpha _e=0`$ and $`\alpha _f=\pi `$ correspond to the skyrmion energy minimum. Thus we shall express the anisotropic skyrmion energy in terms of the two relevant parameters $`\beta _e`$ and $`\beta _f`$ that rotates the core of skyrmion in the empty-electron space and filled-electron space correspondingly. Besides these angles the skyrmion energy depends on free parameters entering the BP solution. The calculation here has been performed only for skyrmion topological index: $`|𝒬|=1`$ (56) and in this case there is only one such parameter, namely, the radius $`R`$ of the core of skyrmion. There are few different spatial integrals that we encounter in calculation. And only one of them is logarithmically divergent. We calculate the skyrmion anisotropy energy with the logarithmic accuracy. It means that parts of the anisotropy energy coming from different anisotropy sources are all multiplied by the same spatial integral, which we denote by a constant $`A`$:
$$A=\left(\frac{R}{l_H}\right)^2\mathrm{log}\frac{R^{}}{R}.$$
(59)
where $`R^{}`$ is the inverse mass of Goldstone excitations in the model (33,50). Also, the Zeeman energy of skyrmion reads:
$$\begin{array}{c}E_Z^{skyr}=A_Z=A|g|\mu _BH(\frac{1}{2}\mathrm{sin}\vartheta (\mathrm{cos}\beta _f\mathrm{cos}\beta _e)(\mathrm{cos}\theta _+\mathrm{cos}\theta _{})+\\ +\mathrm{sin}\vartheta (\mathrm{cos}\theta _++\mathrm{cos}\theta _{})\mathrm{cos}\vartheta (\mathrm{sin}\beta _f+\mathrm{sin}\beta _e)\mathrm{sin}\frac{\theta _++\theta _{}}{2})\end{array}$$
(60)
The hopping energy of skyrmion reads:
$$\begin{array}{c}E_H^{skyr}=A_H=A(t+E^{0x})(\frac{1}{2}\mathrm{cos}\vartheta (\mathrm{cos}\beta _f\mathrm{cos}\beta _e)(\mathrm{sin}\theta _+\mathrm{sin}\theta _{})+\\ +\mathrm{cos}\vartheta (\mathrm{sin}\theta _++\mathrm{sin}\theta _{})\mathrm{sin}\vartheta (\mathrm{sin}\beta _f+\mathrm{sin}\beta _e)\mathrm{cos}\frac{\theta _++\theta _{}}{2})\end{array}$$
(61)
The gate asymmetry energy of skyrmion reads:
$$\begin{array}{c}E_G^{skyr}=A_G=A(\mu _z+E^{0z})((\mathrm{cos}\theta _+\mathrm{cos}\theta _{})\\ \frac{1}{2}(\mathrm{cos}\beta _f\mathrm{cos}\beta _e)(\mathrm{cos}\theta _++\mathrm{cos}\theta _{}))\end{array}$$
(62)
And finally the anisotropic Coulomb energy of skyrmion reads:
$$\begin{array}{c}E_{zz}^{skyr}=A_{zz}=AE^{zz}(\frac{1}{2}(1+\mathrm{cos}\beta _f\mathrm{cos}\beta _e)(1+\mathrm{cos}\theta _+\mathrm{cos}\theta _{})+\\ +2\mathrm{cos}\theta _+\mathrm{cos}\theta _{}\frac{1}{2}\mathrm{sin}\beta _f\mathrm{sin}\beta _e\mathrm{sin}\theta _+\mathrm{sin}\theta _{}).\end{array}$$
(63)
There is also a contribution to the skyrmion energy coming from an non-homogeneous BP electric charge distribution inside the skyrmion core (42):
$$E_C^{skyr}=\frac{1}{2}d^2\stackrel{}{r}d^2\stackrel{}{r^{}}\frac{R^2}{\pi (r^2+R^2)^2}\frac{e^2}{\kappa |\stackrel{}{r}\stackrel{}{r^{}}|}\frac{R^2}{\pi (r^2+R^2)^2}=\frac{3\pi ^2}{64}\frac{e^2}{\kappa R}.$$
(64)
The minimum of the total anisotropic skyrmion energy (60,61,62, 63): $`^{sky}=_Z+_H+_G+_{zz}`$, over the two parameters $`\beta _e`$ and $`\beta _f`$ was found numerically and is denoted as: $`_{min}^{skyr}`$. Next, we find a minimum the total skyrmion energy including (64): $`E^{skyr}=A_{min}^{skyr}+E_C^{skyr}`$, with respect to the skyrmion radius $`R`$:
$$\mathrm{\Delta }=\frac{𝒬+|𝒬|}{2}E_1+3\left(_{min}^{skyr}\left(\frac{3\pi ^2e^2}{128\kappa l_H}\right)^2\mathrm{log}\frac{e^2}{\kappa l_H_{min}^{skyr}}\right)^{1/3}.$$
(65)
This formula holds in the limit $`^{skyr}E_1`$ and, thus, the second term is much smaller then the first term in the Eq.(65) as they are calculated. But, it is important in the case of antiskyrmion $`𝒬=|𝒬|`$, the gap is only a relatively small anisotropic energy. The resulting anisotropic part of a skyrmion gap is shown on the Fig.4 in the case $`\mu ^z+E^{0z}=0`$. Note the two prominent cusp-like lines in the skyrmion gap sheet coincide with the two phase transition lines from the Fig.3. A view of another cross-section of the skyrmion gap: $`|g|\mu _BH=0`$, is shown on the Fig.5. A minima here also coincide with the canted-antiferromagnetic phase. A skyrmion in the ferromagnetic state is a spin-skyrmion with the spin rotations being localized inside one of the two layer, whereas a skyrmion in the spin-singlet state is a layer-skyrmion with the electron density being distributed over the two layers.
In the experimental setup $`^{skyr}E_1`$ and our formulas can be compared with the experimental results only qualitatively. But they found a profound disappearence of the thermal activation gap maximum at some interval on the $`\nu =2`$ line. In our theory this would correspond to a minimum in the skyrmion gap and we suggest that this indicate the canted antiferromagnetic phase.
In the case $`\nu =1`$ we parameterize general rotations from the denominator sub-group by four angles in such a way that the electron density matrix becomes:
$$WQ_{BP}(z\overline{z})W^+=\frac{1}{R^2+|z|^2}\left(\begin{array}{cc}|z|^2& Rz|f\\ & \\ R\overline{z}f|& R^2|ff|\end{array}\right)$$
(66)
where
$$|f=(\mathrm{cos}\frac{\beta }{2},\mathrm{sin}\frac{\beta }{2}\mathrm{cos}\alpha e^{i\lambda _1},\mathrm{sin}\frac{\beta }{2}\mathrm{sin}\alpha e^{i\lambda _2})$$
(67)
Once again we find that the skyrmion energy does not depend on the parameters $`\lambda _1`$ and $`\lambda _2`$. A straightforward calculation shows then that the Zeeman skyrmion energy is
$$\mathrm{\Delta }=2A|g|\mu _BH\left(1\mathrm{sin}^2\frac{\beta }{2}\mathrm{cos}^2\alpha \right),$$
(68)
the Hopping skyrmion energy is
$$E_H^{skyr}=At\left(\mathrm{sin}\beta \mathrm{sin}\alpha +\mathrm{sin}\theta \left[1+\mathrm{cos}^2\alpha \mathrm{sin}^2\frac{\beta }{2}\right]\right)$$
(69)
whereas the gate asymmetry skyrmion energy is:
$$E_G^{skyr}=A\mu _z\left(\mathrm{cos}\beta \left[1\mathrm{sin}^2\frac{\beta }{2}\mathrm{cos}^2\alpha \right]\mathrm{cos}^2\alpha \mathrm{sin}^2\frac{\theta }{2}+\mathrm{cos}\theta \right)$$
(70)
Note that there is no Coulomb energy in the case $`\nu =1`$. Searching for minima of Eq.(68,69,70) varying parameters $`\beta `$ and $`\alpha `$ and fixing $`\theta =\mathrm{tan}^1t/\mu ^z`$ we find the minimal anisotropy energy of skyrmion to be:
$$\mathrm{\Delta }=\frac{𝒬+|𝒬|}{2}E_1+\text{min}(2\sqrt{t^2+\mu _z^2},2|g|\mu _BH).$$
(71)
The skyrmion gap is given by the same formula as in the case $`\nu =2`$: (65). Generally, there is no prominent minima in the skyrmion gap here. Nevetherless in the limit $`t|g|\mu _BH`$ such a minimum occurs.
The case $`\nu =3`$ reduces to the case $`\nu =1`$ if one simply to renormalize the constant in (71) as it was explained in the previous section.
## Conductivity Activation Energy
The conductivity activation energy was measured in experiments for bilayer system and in one layer system . It was known for quite some time that this energy is considerably less than a typical exchange constant $`e^2/\kappa l_H`$. This fact was the main motivation for the experimental search of topological excitations. If one considers an act of creation of skyrmion and anti-skyrmion pair at large separation then one easily gets the pair energy: $`E_1=1/2E_0`$. This is definitely less than the creation of electron-hole pair at large separation: $`E_0`$, but still of the order of $`e^2/\kappa l_H`$, for one layer. In a bilayer system the prominent minimum in the minimal activation energy for skyrmion antiskyrmion pair is still of the order of $`E_1`$, in spite of the fact that experimental conductivity activation energy goes to zero very sharply .
This controversy can be overcome by considering the creation of neutral antyskyrmion and electron pair at large separation. First, we consider a one layer case. The energy of the additional electron with reversed spin does not contain exchange Coulomb energy and consist from anisotropy energy only. The energy of anti-skyrmion also has only anisotropic energy (65). The anisotropic energy of electron can be neglected in the limit of large ratio: $`E_1/|g|\mu _BH`$, which holds in most experiments. The total anti-skyrmion energy is:
$$|g|\mu _BB(1\mathrm{cos}\beta (r))\frac{d^2\stackrel{}{r}}{2\pi l_H^2}+\frac{1}{2}\frac{e^2}{\kappa |\stackrel{}{r}\stackrel{}{r}^{}|}\text{curl}\mathrm{\Omega }^z(\stackrel{}{r})\text{curl}\mathrm{\Omega }^z(\stackrel{}{r}^{})d^2\stackrel{}{r}d^2\stackrel{}{r}^{}.$$
(72)
. Plugging the BP solution we calculate this energy to be:
$$2|g|\mu _BBR^2\mathrm{log}\frac{R^{}}{R}+\frac{3\pi ^2}{64}\frac{e^2}{\kappa R}.$$
(73)
Minimizing it further with respect to radius of the anti-skyrmion core $`R`$ we find the core radius:
$$R=\frac{1}{2}l_H\left(\frac{3\pi ^2e^2}{32\kappa l_H|g|\mu _B\mathrm{log}(R^{}/l_H)}\right)^{1/3}$$
(74)
and the activation energy:
$$\mathrm{\Delta }=\frac{3}{2}\left(|g|\mu _BH\mathrm{log}(R^{}/l_H)\left(\frac{3\pi ^2e^2}{128\kappa l_H}\right)^2\right)^{1/3}.$$
(75)
The upper limit under the logarithm $`R^{}`$ is defined by the validity of BP solution. Essentially, it is a distance where gradient energy becomes of the order of Zeeman energy:
$$R^{}=l_H\sqrt{\frac{E_1}{|g|\mu _BH}}.$$
(76)
Therefore the logarithmic factor is of the order of unity. If it were rather large one would need to compare the energy of $`𝒬=1`$ and $`𝒬=2`$ antiskyrmions, with the latter being logarithmic free. In any case an important point is that the activation energy is magnetic field dependent:
$$\mathrm{\Delta }=KH^{2/3},$$
(77)
which in agreement with the various experimental results . The absolute value of the constant $`K`$ (75)in this relation also conforms experiment . Note that the Zeeman energy of electron is much less that $`\mathrm{\Delta }`$ (77). In the case of double layer system the situation is more complicated due to large number of parameters defining the $`H^{anis}`$. Nevertheless, assuming anisotropy energy is small and neglecting for the same reason the anisotropy energy of an electron we get the deep minimum in the canted antiferromagnetic phase in accordance with the experimental result . It should be noted that in the process of creation of anti-skyrmion electron pair the total topological charge of the system is not conserved as opposed to the skyrmion antiskyrmion pair. Therefore, this process goes tentatively on the sample boundary. Also the existence of magnetic field makes it possible to violate electron-hole symmetry usually assumed in theories used the basis of projected on the lowest Landau Level wave functions. The violation of the electron-hole symmetry is related to the topological charge and gives rise to the skyrmion vs antiskyrmion energy difference due to the topological term in the action (33).
## Acknowledgement
Valuable discussions with V.T.Dolgopolov and A.Maltsev are gratefully acknowledged. This work was supported by grant INTAS N97-31980.
|
warning/0003/hep-ph0003076.html
|
ar5iv
|
text
|
# Spin content of the Λ hyperon
## I Introduction
The spin and flavor content of the nucleon has been an extensively studied issue for well over a decade, since the EMC announced that the strange quark is strongly polarized opposite to the valence quarks, indicating that the quarks carry only a small fraction of the nucleon spin. In contrast to the Ellis-Jaffe sum rule , the polarization of the strange quark turned out to be nonnegligible. A series of following experiments confirmed the EMC result . Hence, it is also natural to investigate the structure of other baryons. In particular, the $`\mathrm{\Lambda }`$ hyperon is interesting, since $`\mathrm{\Delta }q_\mathrm{\Lambda }`$ are related to the fragmentation functions which can be measured experimentally . According to the naive quark model the spin of $`\mathrm{\Lambda }`$ comes solely from the strange quark, while the up and down quarks form the spin singlet, so that they make no contribution. However, an analysis based on the data of hyperon semileptonic decays and lepton-nucleon deep-inelastic scattering predicts $`\mathrm{\Delta }s_\mathrm{\Lambda }0.6`$ and $`\mathrm{\Delta }u_\mathrm{\Lambda }=\mathrm{\Delta }d_\mathrm{\Lambda }0.2`$. This analysis assumes, however, the flavor SU(3) symmetry, i.e. it uses the relations: $`\mathrm{\Delta }u=\mathrm{\Delta }d=\frac{1}{3}(\mathrm{\Delta }\mathrm{\Sigma }D)`$ and $`\mathrm{\Delta }s=\frac{1}{3}(\mathrm{\Delta }\mathrm{\Sigma }+2D)`$, with $`\mathrm{\Delta }\mathrm{\Sigma }`$ identical for all octet baryons. It implies that the $`\mathrm{\Lambda }`$ spin is not completely carried by the strange quark. While the above analysis shows a discrepancy with the naive quark model, one should note that the effect of SU(3) symmetry breaking was not taken into account.
Recently, we have investigated hyperon semileptonic decays and the spin content of the nucleon with SU(3) flavor symmetry breaking taken into account . The SU(3) symmetry breaking was implemented via the Chiral Quark-Soliton Model ($`\chi `$QSM) in such a way that all dynamical variables in the model were fixed by the experimental data for the semileptonic decay constants. The results for the proton were: $`\mathrm{\Delta }u_\text{p}=0.72\pm 0.07`$, $`\mathrm{\Delta }d_\text{p}=0.54\pm 0.07`$, $`\mathrm{\Delta }s_\text{p}=0.33\pm 0.51`$, and $`\mathrm{\Delta }\mathrm{\Sigma }_\text{p}=0.51\pm 0.41`$. The large errors in $`\mathrm{\Delta }s_\text{p}`$ and $`\mathrm{\Delta }\mathrm{\Sigma }_\text{p}`$ are due to the large experimental uncertainties of the $`\mathrm{\Xi }^{}`$ decay constants. The conclusion in that work was as follows: First, statements concerning $`\mathrm{\Delta }s_\text{p}`$ and $`\mathrm{\Delta }\mathrm{\Sigma }_\text{p}`$ based on SU(3) flavor symmetry are premature. Second, accurate results can be obtained only by reducing the experimental uncertainty for $`\mathrm{\Xi }`$ decays. It is of great interest to study the spin content of the $`\mathrm{\Lambda }`$, using the same framework as in Ref. . The aim of this paper is thus to find out what we can reliably conclude on the spin structure of the $`\mathrm{\Lambda }`$, based on hyperon semileptonic decays.
Let us first briefly recall how the standard analysis is carried out. Three diagonal axial-vector coupling constants define integrated polarized quark densities for a given baryon B:
$`g_\mathrm{A}^{(3)}(\text{B})`$ $`=`$ $`\mathrm{\Delta }u_\text{B}\mathrm{\Delta }d_\text{B},`$ (1)
$`\sqrt{3}g_\mathrm{A}^{(8)}(\text{B})`$ $`=`$ $`\mathrm{\Delta }u_\text{B}+\mathrm{\Delta }d_\text{B}2\mathrm{\Delta }s_\text{B},`$ (2)
$`g_\mathrm{A}^{(0)}(\text{B})`$ $`=`$ $`\mathrm{\Delta }u_\text{B}+\mathrm{\Delta }d_\text{B}+\mathrm{\Delta }s_\text{B}.`$ (3)
Note that in our normalization $`g_\mathrm{A}^{(0)}(`$B$`)=`$ $`\mathrm{\Delta }\mathrm{\Sigma }_\text{B}`$.
Assuming SU(3) symmetry, one can calculate $`g_\text{A}^{(3,8)}(`$B$`)`$ in terms of the reduced matrix elements $`F`$ and $`D`$. For the proton and $`\mathrm{\Lambda }`$ one gets:
$`g_\mathrm{A}^{(3)}(\text{p})=F+D,`$ $`\sqrt{3}g_\mathrm{A}^{(8)}(\text{p})=3FD,`$ (4)
$`g_\mathrm{A}^{(3)}(\mathrm{\Lambda })=0,`$ $`\sqrt{3}g_\mathrm{A}^{(8)}(\mathrm{\Lambda })=2D.`$ (5)
The constants $`F`$ and $`D`$ can be in principle extracted from the hyperon semileptonic decays. For example:
$$A_1=\left(g_1/f_1\right)^{(\mathrm{n}\mathrm{p})}=F+D,A_4=\left(g_1/f_1\right)^{(\mathrm{\Sigma }^{}\mathrm{n})}=FD.$$
(6)
For convenience, we denote the ratios of axial-vector to vector decay constants by $`A_i`$ (see Table I). Taking for these decays experimental values (see Table I), one gets $`F=0.46`$ and $`D=0.80`$.
Since $`g_\mathrm{A}^{(0)}(`$B$`)`$ does not correspond to the SU(3) current, it cannot be expressed in terms of $`F`$ and $`D`$ without further assumptions. Thus, in order to extract all $`\mathrm{\Delta }q_\text{p}`$ separately, one needs some additional information. Either another experimental input is needed, or a model which predicts $`g_\mathrm{A}^{(0)}(`$B$`)`$ in terms of the $`F`$ and $`D`$.
The first possibility can be realized by taking the experimental result for the first moment of the spin structure function $`g_1^\text{p}(x)`$ of the proton:
$$I_\text{p}=\underset{0}{\overset{1}{}}𝑑xg_1^\text{p}(x)=\frac{1}{18}\left(4\mathrm{\Delta }u_\text{p}+\mathrm{\Delta }d_\text{p}+\mathrm{\Delta }s_\text{p}\right)\left(1\frac{\alpha _\text{s}}{\pi }+\mathrm{}\right).$$
(7)
Recent analysis implies $`I_\text{p}=0.124\pm 0.011`$ which translates into:
$$\mathrm{\Gamma }_\text{p}4\mathrm{\Delta }u_\text{p}+\mathrm{\Delta }d_\text{p}+\mathrm{\Delta }s_\text{p}=2.56\pm 0.23.$$
(8)
if $`\alpha _\mathrm{s}(Q^2=3(\mathrm{GeV}/c)^2)=0.4`$ is assumed. Taking for $`F=0.46`$ and for $`D=0.80`$ together with Eq.(8), one gets for the nucleon: $`\mathrm{\Delta }u_\text{p}=0.79`$, $`\mathrm{\Delta }d_\text{p}=0.47`$ and $`\mathrm{\Delta }s_\text{p}=0.13`$, which implies $`\mathrm{\Delta }\mathrm{\Sigma }_\text{p}=0.19`$, a fairly small number as compared with the naive expectation from the quark model: $`\mathrm{\Delta }\mathrm{\Sigma }_\text{p}=1`$.
It is important to realize that $`\mathrm{\Delta }\mathrm{\Sigma }_\text{p}`$ is not directly measured; it is extracted from the data through some theoretical model. In the above example we have assumed the SU(3) symmetry and used two arbitrarily chosen hyperon decays (6). One could, however, use any two $`A_i`$’s out of 6 known hyperon decays to extract $`F`$ and $`D.`$ The number of combinations which one can form to extract $`F`$ and $`D`$ is 14 (actually 15, but two conditions are linearly dependent). Taking these 14 combinations into account, one gets:
$$F=0.40÷0.55,D=0.70÷0.89.$$
(9)
These are the uncertainties of the central values due to the theoretical error caused by using the exact SU(3) symmetry to describe the hyperon semileptonic decays. These uncertainties are further increased by the experimental errors of all individual decays.
Looking at Eq.(9), one might get an impression that a typical error associated with using SU(3) symmetry in analyzing the hyperon decays is of the order of 15 % or so. While this is true for the hyperon decays, the values of $`\mathrm{\Delta }q_\mathrm{B}`$ and $`\mathrm{\Delta }\mathrm{\Sigma }_\mathrm{B}`$ of the various baryons might be much more affected by the symmetry breaking. Indeed for $`F`$ and $`D`$ corresponding to (9) and $`\mathrm{\Gamma }_\text{p}`$ as given by Eq.(8) $`\mathrm{\Delta }\mathrm{\Sigma }_\text{p}=0.02÷0.30`$.
As will be shown in the following, the $`\chi `$QSM predicts in the chiral limit :
$$g_\mathrm{A}^{(0)}(\text{B})=9F5D$$
(10)
for any baryon B. Here $`g_\mathrm{A}^{(0)}`$ is very sensitive to small variations of $`F`$ and $`D`$, since it is a difference of the two, with relatively large multiplicators. Indeed, for the 14 fits mentioned above the central value for $`g_\mathrm{A}^{(0)}`$ of the nucleon varies between $`0.25`$ to approximately 1 and a similar feature is expected for any baryon, particularly for the $`\mathrm{\Lambda }`$. Thus, despite the fact that hyperon semileptonic decays are relatively well described by the model in the chiral limit, the singlet axial-vector constant is basically undetermined. This is a clear signal of the importance of the symmetry breaking for this quantity.
One could argue that this kind of behavior is just an artifact of the $`\chi `$ QSM. However, the scenario of a rotating soliton (which is by the way used also in the Skyrme-type models) is very plausible and cannot be a priori discarded on the basis of first principles. The $`\chi `$QSM is a particular realization of this scenario and we use it as a tool to investigate the sensitivity of the singlet axial-vector current to the symmetry breaking effects in hyperon semileptonic decays. In fact, conclusions similar to ours have been obtained in chiral perturbation theory in Ref..
In Ref. we have at length discussed the properties of the model formula for $`g_\mathrm{A}^{(0)}`$ in two limiting cases , i.e. large (Skyrme model limit) and small (quark model limit) soliton sizes. In the Skyrme model the ratio $`F/D=5/9`$ and $`\mathrm{\Delta }\mathrm{\Sigma }_\text{p}`$ vanishes. In the quark model $`F/D=2/3`$ and $`F+D=5/3`$, and therefore $`\mathrm{\Delta }\mathrm{\Sigma }_\text{p}=1`$. We also gave numerical arguments in support of our approach: namely releasing the model assumptions concerning $`g_\mathrm{A}^{(0)}`$ and using $`\mathrm{\Gamma }_\text{p}`$ as an additional input one arrives at almost identical numerical results as using (10).
It is virtually impossible to analyze the symmetry breaking in weak decays without resorting to some specific model . In this paper, following Ref. , we will implement the symmetry breaking for the hyperon decays using the $`\chi `$ QSM (see Ref. for review) which satisfactorily describes the axial-vector properties of hyperons .
The model provides a link between the matrix elements of the octet of the axial-vector currents, responsible for hyperon decays, and the matrix elements of the singlet axial-vector current. In the present work we will study the relation between hyperon semileptonic decays and integrated polarized quark distributions for the $`\mathrm{\Lambda }`$ hyperon. We will use the $`\chi `$QSM only to identify the algebraic structure of the symmetry breaking ($`m_\text{s}`$ corrections). The dynamical quantities, so called inertia parameters which are in principle calculable within the model , will be treated as free parameters. By adjusting them to the experimentally known semileptonic decays we allow not only for maximal phenomenological input but also for minimal model dependence. In Ref. we have already studied the magnetic moments of the octet and decuplet in this way.
Such a ”model-independent” approach – used for example by Adkins and Nappi in the context of the Skyrme model – is of interest for at least two reasons. First, it can be considered as a QCD-motivated tool to analyze and classify (in terms of powers of $`m_\mathrm{s}`$ and $`1/N_\mathrm{c}`$) the symmetry-breaking terms for a given observable. For nontrivial operators such as axial-vector form factors a general analysis, without referring to some specific model, is virtually impossible. Second, this ”model-independent” analysis provides an information for the model builders as well. It tells us what are the best predictions the model can ever produce. Indeed, model calculations in the framework of the $`\chi `$QSM are not as unique as one might think: They depend on adopted regularizations, cutoff parameter, or the constituent quark mass. Moreover, in the SU(3) version of the $`\chi `$QSM a quantization ambiguity appears . Therefore, if the “model-independent” analysis would have failed to describe the data, that would mean that the model did not correctly include all necessary physics relevant for a given observable. On the other hand, the success of such an analysis gives a strong hint for the model builders that the model is correct and worth exploring.
As far as the symmetry breaking is concerned, our results are identical to the ones obtained in Refs. within QCD in the large $`N_\mathrm{c}`$ limit. Indeed, the $`\chi `$QSM is a specific realization of the large $`N_\mathrm{c}`$ limit. The truly new ingredient of our analysis is the model formula for the singlet axial-vector constant $`g_\mathrm{A}^{(0)}`$, i.e. Eq.(10), which we use to calculate quantities relevant for polarized high energy experiments.
## II Hyperon decays in the Chrial Quark Soliton Model
The discussion in this section follows closely Ref. . The transition matrix elements of the hadronic axial-vector current $`B_2|A_\mu ^X|B_1`$ can be expressed in terms of three independent form factors:
$$B_2|A_\mu ^X|B_1=\overline{u}_{B_2}(p_2)\left[\left\{g_1^{B_1B_2}(q^2)\gamma _\mu \frac{ig_2^{B_1B_2}(q^2)}{M_1}\sigma _{\mu \nu }q^\nu +\frac{g_3^{B_1B_2}(q^2)}{M_1}q_\mu \right\}\gamma _5\right]u_{B_1}(p_1),$$
(11)
where the axial-vector current is defined as
$$A_\mu ^X=\overline{\psi }(x)\gamma _\mu \gamma _5\lambda _X\psi (x)$$
(12)
with $`X=\frac{1}{2}(1\pm i2)`$ for strangeness conserving $`\mathrm{\Delta }S=0`$ currents and $`X=\frac{1}{2}(4\pm i5)`$ for $`|\mathrm{\Delta }S|=1`$. Similar expressions hold for the hadronic vector current, where the $`g_i`$ are replaced by $`f_i`$ ($`i=1,2,3`$ ) and $`\gamma _5`$ by $`\mathrm{𝟏}`$.
Hadronic matrix elements such as $`B_2|A_\mu ^X|B_1`$ can be easily evaluated within the $`\chi `$QSM . Taking into account the $`1/N_c`$ rotational and $`m_\mathrm{s}`$ corrections, we can write the resulting axial-vector constants $`g_1^{B_1B_2}(0)`$ in the following form<sup>*</sup><sup>*</sup>*In the following we will assume that the baryons involved have $`S_3=+\frac{1}{2}`$.:
$`g_1^{(B_1B_2)}`$ $`=`$ $`a_1B_2|D_{X3}^{(8)}|B_1+a_2d_{pq3}B_2|D_{Xp}^{(8)}\widehat{S}_q|B_1+{\displaystyle \frac{a_3}{\sqrt{3}}}B_2|D_{X8}^{(8)}\widehat{S}_3|B_1`$ (13)
$`+`$ $`m_s[{\displaystyle \frac{a_4}{\sqrt{3}}}d_{pq3}B_2|D_{Xp}^{(8)}D_{8q}^{(8)}|B_1+a_5B_2|(D_{X3}^{(8)}D_{88}^{(8)}+D_{X8}^{(8)}D_{83}^{(8)})|B_1`$ (14)
$`+`$ $`a_6B_2|(D_{X3}^{(8)}D_{88}^{(8)}D_{X8}^{(8)}D_{83}^{(8)})|B_1].`$ (15)
$`\widehat{S}_q`$ ($`\widehat{S}_3`$) stand for the $`q`$-th (third) component of the spin operator of the baryons. The $`D_{ab}^{()}`$ denote the SU(3) Wigner matrices in representation $``$. The $`a_i`$ denote parameters depending on the specific dynamics of the chiral soliton model (see for example Refs. and references therein). Their explicit form in terms of a Goldstone mean field can be found in Ref. . As mentioned already, in the present approach we will not calculate this mean field but treat $`a_i`$ as free parameters to be adjusted to experimentally known semileptonic hyperon decays.
Because of the SU(3) symmetry breaking due to the strange quark mass $`m_\mathrm{s}`$, the collective baryon Hamiltonian is no more SU(3)-symmetric. The octet states are mixed with the higher representations such as antidecuplet $`\overline{\mathrm{𝟏𝟎}}`$ and eikosiheptaplet $`\mathrm{𝟐𝟕}`$ . In the linear order in $`m_\mathrm{s}`$ the wave function of a state $`B=(Y,I,I_3)`$ of spin $`S_3`$ is given as:
$$\psi _{B,S_3}=()^{\frac{1}{2}S_3}\left(\sqrt{8}D_{BS}^{(8)}+c_B^{(\overline{10})}\sqrt{10}D_{BS}^{(\overline{10})}+c_B^{(27)}\sqrt{27}D_{BS}^{(27)}\right),$$
(16)
where $`S=(1,\frac{1}{2},S_3)`$. Mixing parameters $`c_B^{()}`$ can be found for example in Ref. . They are given as products of a numerical constant $`N_B^{()}`$ depending on the quantum numbers of the baryonic state $`B`$ and a dynamical parameter $`c_{}`$ depending linearly on $`m_\mathrm{s}`$ (which we assume to be 180 MeV) and the model parameter $`I_2`$, which is responsible for the splitting between the octet and higher exotic multiplets .
Analogously to Eq.(15), one obtains in the $`\chi `$QSM diagonal axial-vector coupling constants. In that case $`X`$ can take two values: $`X=3`$ and $`X=8`$. For $`X=0`$ (singlet axial-vector current) we have the following expression :
$$\frac{1}{2}g_A^{(0)}(\text{B})=\frac{1}{2}a_3+\sqrt{3}m_\text{s}(a_5a_6)B|D_{83}^{(8)}|B.$$
(17)
This equation is remarkable, since it provides a link between an octet and singlet axial-vector current. It is the most important model input in our analysis. Pure QCD-arguments based the large $`N_c`$ expansion do not provide such a link. Moreover, due to the structure of the matrix element $`B|D_{83}^{(8)}|B`$, the $`g_A^{(0)}(`$B$`)`$ are identical inside the isospin multiplets. We predict much stronger symmetry breaking for the $`\mathrm{\Lambda }`$ than for the proton, since
$$\sqrt{3}p|D_{83}^{(8)}|p=\frac{1}{10},\sqrt{3}\mathrm{\Lambda }|D_{83}^{(8)}|\mathrm{\Lambda }=\frac{3}{10},$$
(18)
for spin $`S_3=+1/2`$.
Instead of calculating 7 dynamical parameters $`a_i(i=1,\mathrm{},6)`$ and $`I_2`$ (which enters into $`c_{\overline{10}}`$ and $`c_{27}`$) within the $`\chi `$QSM, we shall fit them from the hyperon semileptonic decays data. It is convenient to introduce the following set of 7 new parameters:
$`r={\displaystyle \frac{1}{30}}\left(a_1{\displaystyle \frac{1}{2}}a_2\right),s={\displaystyle \frac{1}{60}}a_3,x={\displaystyle \frac{1}{540}}m_\mathrm{s}a_4,y={\displaystyle \frac{1}{90}}m_\mathrm{s}a_5,z={\displaystyle \frac{1}{30}}m_\mathrm{s}a_6,`$
$$p=\frac{1}{6}m_\mathrm{s}c_{\overline{10}}\left(a_1+a_2+\frac{1}{2}a_3\right),q=\frac{1}{90}m_\mathrm{s}c_{27}\left(a_1+2a_2\frac{3}{2}a_3\right).$$
(19)
Employing this new set of parameters, we can express all possible semileptonic decays of the octet baryons:
$`A_1=\left(g_1/f_1\right)^{(\mathrm{n}\mathrm{p})}`$ $`=`$ $`14r+2s44x20y4z4p+8q,`$ (20)
$`A_2=\left(g_1/f_1\right)^{(\mathrm{\Sigma }^+\mathrm{\Lambda })}`$ $`=`$ $`9r3s42x6y3p+15q,`$ (21)
$`A_3=\left(g_1/f_1\right)^{(\mathrm{\Lambda }\mathrm{p})}`$ $`=`$ $`8r+4s+24x2z+2p6q,`$ (22)
$`A_4=\left(g_1/f_1\right)^{(\mathrm{\Sigma }^{}\mathrm{n})}`$ $`=`$ $`4r+8s4x4y+2z+4q,`$ (23)
$`A_5=\left(g_1/f_1\right)^{(\mathrm{\Xi }^{}\mathrm{\Lambda })}`$ $`=`$ $`2r+6s6x+6y2z+6q,`$ (24)
$`A_6=\left(g_1/f_1\right)^{(\mathrm{\Xi }^{}\mathrm{\Sigma }^0)}`$ $`=`$ $`14r+2s+22x+10y+2z+2p4q.`$ (25)
The U(1) and SU(3) axial-vector constants $`g_A^{(0,3,8)}`$ can be also expressed in terms of the new set of parameters (19). In the case of the proton and the $`\mathrm{\Lambda }`$ we have the singlet axial-vector constants:
$$g_A^{(0)}(\mathrm{p})=60s18y+6z,g_A^{(0)}(\mathrm{\Lambda })=60s+54y18z,$$
(26)
for the triplet ones, we writeTriplet $`g_A^{(3)}`$’s are proportional to $`I_3`$, formulae in Eq.(27) correspond to the highest isospin state.:
$$g_A^{(3)}(\mathrm{p})=14r+2s44x20y4z4p+8q,g_A^{(3)}(\mathrm{\Lambda })=0,$$
(27)
and for the octet one, we obtain:
$$g_A^{(8)}(\mathrm{p})=\sqrt{3}(2r+6s+12x+4p+24q),g_A^{(3)}(\mathrm{\Lambda })=\sqrt{3}(6r+2s36x+36q).$$
(28)
Let us finally note that there is certain redundancy in Eq.(25-28), namely by redefinition of $`q`$ and $`x`$ we can get rid of the variable $`p`$:
$$x^{}=x\frac{1}{9}p,q^{}=q\frac{1}{9}p.$$
(29)
## III Spin content of $`\mathrm{\Lambda }`$ hyperon
As shown in the last section there are 6 free parameters which have to be fitted from the data. There are 2 chiral parameters: $`r`$ and $`s`$, related closely to $`F`$ and $`D`$:
$$F=5(sr),D=3(s+3r).$$
(30)
and 4 proportional to $`m_\mathrm{s}`$: $`x^{}`$, $`y`$, $`z`$, and $`q^{}`$. Since there are six known hyperon semileptonic decays, we can express all model parameters as linear combinations of these decay constants, and subsequently all quantities of interest can be expressed in terms of the input amplitudes. In the following we will use the experimental values of Refs. , which are presented in Table I.
Before doing this, let us, however, observe that there exist two linear combinations $`A_i`$’s which within the model are free of the $`m_\mathrm{s}`$ corrections:
$`42r+6s`$ $`=`$ $`A_1+2A_6,`$ (31)
$`90r+90s`$ $`=`$ $`3A_18A_26A_3+6A_4+6A_5.`$ (32)
Solving Eq.(32) for $`r`$ and $`s`$, we obtain the chiral-limit expressions for hyperon semileptonic decays and integrated quark densities (i.e. with $`x^{}=y=z=q^{}=0`$). The corresponding $`F`$ and $`D`$ take the following form:
$`F`$ $`=`$ $`{\displaystyle \frac{1}{12}}(4A_14A_23A_3+3A_4+3A_5+5A_6),`$ (33)
$`D`$ $`=`$ $`{\displaystyle \frac{1}{12}}(4A_2+3A_33A_43A_5+3A_6).`$ (34)
Numerically:
$$F=0.50\pm 0.07,D=0.77\pm 0.04.$$
(35)
With these values for $`F`$ and $`D`$ together with Eq.(8) one gets: $`\mathrm{\Delta }u_\text{p}=0.81`$, $`\mathrm{\Delta }d_\text{p}=0.47`$ and $`\mathrm{\Delta }s_\text{p}=0.20`$, which implies $`\mathrm{\Delta }\mathrm{\Sigma }_\text{p}=0.15`$. The advantage of using Eq.(34) consists in the fact that $`F`$ and $`D`$ do not need to be refitted when $`m_\text{s}`$ corrections are added.
Another important point is, that Eq.(34) is more general than the model considered here. In fact they follow from the large $`N_\text{c}`$ QCD, as discussed in Ref. . The errors come from the experimental errors of the decay amplitudes and are dominated by the errors of the $`\mathrm{\Xi }^{}`$ decays. It is of utmost importance to reduce the errors of these decays in order to get better accuracy for $`F`$ and $`D`$.
In the case of the $`\mathrm{\Lambda }`$ Eq.(5) implies that $`\mathrm{\Delta }u_\mathrm{\Lambda }=\mathrm{\Delta }d_\mathrm{\Lambda }`$ and one has in the chiral limit:
$`\mathrm{\Delta }u_\mathrm{\Lambda }^{(0)}`$ $`=`$ $`3F2D=A_1{\displaystyle \frac{5}{3}}A_2{\displaystyle \frac{5}{4}}A_3+{\displaystyle \frac{5}{4}}A_4+{\displaystyle \frac{5}{4}}A_5+{\displaystyle \frac{3}{4}}A_{6,}`$ (36)
$`\mathrm{\Delta }s_\mathrm{\Lambda }^{(0)}`$ $`=`$ $`3FD=A_1{\displaystyle \frac{4}{3}}A_2A_3+A_4+A_5+A_{6.}`$ (37)
Numerical values $`\mathrm{\Delta }u_\mathrm{\Lambda }^{(0)}=\mathrm{\Delta }d_\mathrm{\Lambda }^{(0)}=0.03\pm 0.14`$ and $`\mathrm{\Delta }s_\mathrm{\Lambda }^{(0)}=0.74\pm 0.17`$ (see Table II) are closer to the naive quark model expectations: $`\mathrm{\Delta }u_\mathrm{\Lambda }=\mathrm{\Delta }d_\mathrm{\Lambda }=0`$ and $`\mathrm{\Delta }s_\mathrm{\Lambda }=1`$, than to the numbers quoted in Ref.: $`\mathrm{\Delta }u_\mathrm{\Lambda }=\mathrm{\Delta }d_\mathrm{\Lambda }=0.23\pm 0.06`$ and $`\mathrm{\Delta }s_\mathrm{\Lambda }=0.58\pm 0.07`$. This is reflected in the fact that
$$\mathrm{\Delta }\mathrm{\Sigma }^{(0)}=9F5D=3A_1\frac{14}{3}A_2\frac{7}{2}A_3+\frac{7}{2}A_4+\frac{7}{2}A_5+\frac{5}{2}A_6.$$
(38)
(which is identical to all hadrons) reads: $`\mathrm{\Delta }\mathrm{\Sigma }^{(0)}=0.68\pm 0.44`$ and is much larger than the value required by using $`\mathrm{\Gamma }_\text{p}`$ as an additional input. Indeed, as explained in Ref. , in the chiral limit one is not able to reproduce the value of $`\mathrm{\Gamma }_\text{p}`$ (see Table II).
The two least known amplitudes $`A_5`$ and $`A_6`$ are almost entirely responsible for the errors of $`\mathrm{\Delta }q_\mathrm{\Lambda }`$. However, since the coefficients which enter into Eqs.(37,38) are not too large, the absolute errors are relatively small.
The full expressions are obtained by solving the remaining 4 equations for $`m_\mathrm{s}`$ dependent parameters $`x^{}`$, $`y`$, $`z`$ and $`q^{}`$. Also in this case we are able to link integrated quark densities $`\mathrm{\Delta }q`$ to the hyperon decays:
$`\mathrm{\Delta }u_\mathrm{\Lambda }`$ $`=`$ $`\mathrm{\Delta }d_\mathrm{\Lambda }={\displaystyle \frac{A_2}{3}}{\displaystyle \frac{A_3}{4}}+{\displaystyle \frac{A_4}{4}}+{\displaystyle \frac{13A_5}{4}}{\displaystyle \frac{A_6}{4}},`$ (39)
$`\mathrm{\Delta }s_\mathrm{\Lambda }`$ $`=`$ $`{\displaystyle \frac{15A_1}{4}}{\displaystyle \frac{13A_2}{2}}{\displaystyle \frac{87A_3}{16}}{\displaystyle \frac{21A_4}{16}}+{\displaystyle \frac{45A_5}{16}}+{\displaystyle \frac{51A_6}{16}}`$ (40)
$`\mathrm{\Delta }\mathrm{\Sigma }_\mathrm{\Lambda }`$ $`=`$ $`{\displaystyle \frac{15A_1}{4}}{\displaystyle \frac{46A_2}{6}}{\displaystyle \frac{95A_3}{16}}{\displaystyle \frac{13A_4}{16}}+{\displaystyle \frac{149A_5}{16}}+{\displaystyle \frac{46A_6}{16}}.`$ (41)
To guide the eye it is convenient to restore the linear $`m_\text{s}`$ dependence for the quark densities in the following way:
$`\mathrm{\Delta }q=\mathrm{\Delta }q^{(0)}+{\displaystyle \frac{m_\text{s}}{180\text{MeV}}}\left(\mathrm{\Delta }q\mathrm{\Delta }q^{(0)}\right),`$and similarly for $`\mathrm{\Delta }\mathrm{\Sigma }`$. This dependence is explicitly shown in Fig.1, where we plot the central values and “experimental” error bars (shaded areas) of $`\mathrm{\Delta }q_\mathrm{\Lambda }`$’s.
Fig.1. $`\mathrm{\Delta }q_\mathrm{\Lambda }`$ as functions of $`m_\mathrm{s}`$.
Fig.2. $`\mathrm{\Delta }\mathrm{\Sigma }_\mathrm{p}`$ and $`\mathrm{\Delta }\mathrm{\Sigma }_\mathrm{\Lambda }`$ as functions of $`m_\mathrm{s}`$.
In Fig.2 we plot $`m_\text{s}`$ dependence of $`\mathrm{\Delta }\mathrm{\Sigma }`$ both for the proton and for the $`\mathrm{\Lambda }`$. In order to make the plot readable we have denoted theoretical errors as error bars around the black dots which correspond to the chiral limit and full theoretical prediction. The splitting between the proton and the $`\mathrm{\Lambda }`$ is caused by the term proportional to $`a_5a_6`$ in Eq.(17). Numerical values can be found in Table II. We see that for $`m_\text{s}=180`$ MeV apart from fitting all hyperon semileptonic decays (which is our input) we reproduce $`\mathrm{\Gamma }_\text{p}`$ with relatively small error. The errors of $`\mathrm{\Delta }\mathrm{\Sigma }`$ and $`\mathrm{\Delta }s`$ are much bigger. The central values, however, differ from the “standard” ones. Interestingly $`\mathrm{\Delta }s_\text{p}`$ in proton is rather large and positive, however, the error bars are so large that the quark model value $`\mathrm{\Delta }s_\text{p}=0`$ is not excluded. In the $`\mathrm{\Lambda }`$ the $`\mathrm{\Delta }s_\mathrm{\Lambda }`$ is larger than 1, but again the errors are large. The errors for $`\mathrm{\Delta }u`$ and $`\mathrm{\Delta }d`$ both in the proton and in the $`\mathrm{\Lambda }`$ are much smaller. For the $`\mathrm{\Lambda }`$ we get that the non-strange quarks almost do not carry spin in surprising accordance with the expectations of the naive quark model.
As already discussed, the errors on for $`\mathrm{\Delta }q`$’s and $`\mathrm{\Delta }\mathrm{\Sigma }`$ come almost entirely from the large errors of the $`\mathrm{\Xi }^{}`$ decays ($`A_5`$ and $`A_6`$). Instead of using these two hyperon semileptonic decays $`A_5`$ and $`A_6`$ as input, we can use the experimental value for $`\mathrm{\Gamma }_\text{p}`$ as given by Eq.(8) and $`\mathrm{\Delta }\mathrm{\Sigma }_\text{N}`$, which we vary in the range from 0 to 1. In Fig.3 we plot our predictions for $`A_5`$ and $`A_6`$ (solid lines), together with the experimental error bands for these two decays. It is clearly seen from Fig.3 that the allowed region for $`\mathrm{\Delta }\mathrm{\Sigma }_\text{N}`$, in which the theoretical prediction falls within the experimental error bars amounts to $`\mathrm{\Delta }\mathrm{\Sigma }_\text{N}=0.20÷0.45`$.
Fig.3. $`A_5`$ (lower line) and $`A_6`$ (upper line) as functions of $`\mathrm{\Delta }\mathrm{\Sigma }_\mathrm{p}`$.
Fig.4. $`\mathrm{\Delta }q_\mathrm{\Lambda }`$’s and $`\mathrm{\Delta }\mathrm{\Sigma }_\mathrm{\Lambda }`$ as functions of $`\mathrm{\Delta }\mathrm{\Sigma }_\mathrm{p}`$.
In Fig.4 we plot the dependence of $`\mathrm{\Delta }q_\mathrm{\Lambda }`$’s and $`\mathrm{\Delta }\mathrm{\Sigma }_\mathrm{\Lambda }`$ upon $`\mathrm{\Delta }\mathrm{\Sigma }_\text{p}`$ (with $`\mathrm{\Gamma }_\text{p}`$ fixed by Eq.(8)). We see rather strong correlation of these quantities with $`\mathrm{\Delta }\mathrm{\Sigma }_\text{p}`$. Within the allowed region $`0.20<\mathrm{\Delta }\mathrm{\Sigma }_\text{p}<0.45`$ the strange quark density $`\mathrm{\Delta }s_\mathrm{\Lambda }`$ varies from 0.84 to 1.10. Interestingly, in the central region around $`\mathrm{\Delta }\mathrm{\Sigma }_\text{p}0.35`$ the strange quark density in $`\mathrm{\Lambda }`$ is close to 1 in accordance with an intuitive assumptions of the naive quark model. Nonstrange quarks contribute to the spin of the Lambda at the level of $`0.04`$, and $`\mathrm{\Delta }\mathrm{\Sigma }_\mathrm{\Lambda }=0.92`$.
## IV Summary
In this paper we studied the influence of the SU(3) symmetry breaking in hyperon semileptonic decays on the determination of the integrated polarized quark densities $`\mathrm{\Delta }q_\mathrm{\Lambda }`$ in the $`\mathrm{\Lambda }`$. Using the Chiral Quark-Soliton Model we have obtained a satisfactory parametrization of all available experimental data on semileptonic decays. In this respect our analysis is identical to QCD analysis in the large $`N_\mathrm{c}`$ of Ref..
The new ingredient of our analysis consists in using the model formula for the singlet axial-vector current in order to make contact with the high energy polarization experiments.
The model contains 6 free parameters which can be fixed by 6 known hyperon decays. Unfortunately $`g_1/f_1`$ for the two known decays of $`\mathrm{\Xi }^{}`$ have large experimental errors, which influence our predictions for $`\mathrm{\Delta }q_\mathrm{\Lambda }`$. Our strategy was very simple: using model parametrization we expressed $`\mathrm{\Delta }q_\mathrm{\Lambda }`$’s and $`\mathrm{\Delta }\mathrm{\Sigma }_\mathrm{\Lambda }`$ in terms of the six known hyperon decays. Errors were added in quadrature.
There are two points which have to be stressed here. Our fit respects chiral symmetry in a sense that the leading order parameters $`r`$ and $`s`$ (or equivalently $`F`$ and $`D`$) are fitted to the linear combinations of the hyperon decays which are free from $`m_\mathrm{s}`$ corrections. As discussed in Ref. it is impossible to use the SU(3) symmetric parametrization as given by Eq.(32) and reproduce $`\mathrm{\Gamma }_\text{p}`$ (as far as the central values are concerned) . With $`m_\mathrm{s}`$ corrections turned on one hits the experimental value for $`\mathrm{\Gamma }_\text{p}`$ (see Table II), however, the value of $`\mathrm{\Delta }\mathrm{\Sigma }_\text{p}`$ is practically undetermined, due the the experimental error of $`\mathrm{\Xi }^{}`$ decays.
The nature of the $`m_\text{s}`$ is such that the central value of $`\mathrm{\Delta }\mathrm{\Sigma }_\text{p}`$ is relatively large, whereas $`\mathrm{\Delta }s_\text{p}`$ is positive, however, still compatible with 0 within large errors. So one can accommodate all existing data with $`\mathrm{\Delta }q_\text{p}`$ much closer to the expectations of the naive quark model than in the standard, SU(3) symmetric approach. This trend is even stronger in the case of the $`\mathrm{\Lambda }`$, where $`\mathrm{\Delta }u_\mathrm{\Lambda }=\mathrm{\Delta }d_\mathrm{\Lambda }0`$ and $`\mathrm{\Delta }s_\mathrm{\Lambda }1`$. SU(3) symmetry breaking effects cause $`\mathrm{\Delta }\mathrm{\Sigma }_\mathrm{\Lambda }>\mathrm{\Delta }\mathrm{\Sigma }_\text{p}`$, so that the $`\mathrm{\Lambda }`$ is in a sense more nonrelativistic than the nucleon.
Our analysis shows clearly that if one wants to link hyperon semileptonic decays with high-energy polarized experiments, one cannot neglect SU(3) symmetry breaking for the former. In this respect our conclusion agrees with Refs.. Similarly to Ref. we see that $`\mathrm{\Delta }s_\text{p}=0`$ is not ruled out by present experiments. Therefore, the results for $`\mathrm{\Delta }s_\mathrm{\Lambda }`$ and $`\mathrm{\Delta }\mathrm{\Sigma }_\mathrm{\Lambda }`$ based on the exact SU(3) symmetry are in our opinion premature.
## Acknowledgments
This work has partly been supported by the BMBF, the DFG and the COSY–Project (Jülich). We are grateful to M.V. Polyakov for fruitful discussions. H.-Ch.K. and M.P. thank P.V. Pobylitsa for critical comments. H.-Ch.K. has been supported by the Korean Physical Society (1999). M.P. has been supported by Polish grant PB 2 P03B 019 17.
|
warning/0003/cond-mat0003231.html
|
ar5iv
|
text
|
# Optical transitions in broken gap heterostructures
## I Introduction
The characteristic feature of Al<sub>x</sub>Ga<sub>1-x</sub>Sb/InAs heterostructures with $`x`$$`<`$0.3 is the overlap of the InAs conduction band with the Al<sub>x</sub>Ga<sub>1-x</sub>Sb valence band . Such systems, often referred to as broken gap heterostructures, exhibit interesting negative persistent photoconductivity , semimetal-semiconductor transition induced by magnetic field , and intrinsic exciton . As the Al concentration increases from $`x`$=0.3, a staggered band alignment appears at the Al<sub>x</sub>Ga<sub>1-x</sub>Sb/InAs interface with the valence band edge of Al<sub>x</sub>Ga<sub>1-x</sub>Sb lying in the band gap of InAs. On the other hand, at the GaSb/Al<sub>x</sub>Ga<sub>1-x</sub>Sb interface the band alignment is of the straddled type. Therefore, in a InAs/Al<sub>x</sub>Ga<sub>1-x</sub>Sb/GaSb heterostructure, carriers can tunnel from the InAs conduction band to the GaSb valence through the Al<sub>x</sub>Ga<sub>1-x</sub>Sb barrier.
InAs, GaSb, and AlSb form a family of semiconductors with sufficient lattice match for epitaxy growth. Interband tunneling devices based on the InAs/AlGaSb/GaSb double barrier resonant tunneling have been fabricated, which exhibit high frequency response and peak-to-valley current ratio at room temperatures . Because of their potential for technological application, resonant tunneling and resonant magnetotunneling have been investigated extensively.
To the contrary of tunneling processes, very little work has been done on the electronic structure and optical properties of broken gap heterostructures . We notice that due to the overlap of InAs conduction band and the GaSb valence band, it is possible to form new type of eigenstates with interesting optical properties. We will elaborate this with Fig. 1 where the conduction band edge $`E_c`$ and the valence band edge $`E_v`$ are marked for a GaSb/InAs heterostructure embedded in the band gap of AlSb. Treating the AlSb at both sides as potential barriers, the system AlSb/GaSb/InAs/AlSb is essentially a new type of quantum well which we name as broken gap quantum well (BGQW). If we diminish the thickness of either the GaSb layer or the InAs layer to zero, the BGQW reduces to a conventional quantum well which, with bipolar doping becomes a quantum well laser.
If in a properly designed BGQW there exist two eigenenergies as shown in Fig. 1, under bipolar doping, the BGQW can lase in the infrared frequency region. Such semiconductor infrared lasers or detectors were proposed earlier . The lower energy state in Fig. 1 is characterized by the degree of hybridization between the GaSb valence band states and the InAs conduction band states. The enhancement of optical transition requires the hybridized wave function having a large amplitude in the InAs layer. On the other hand, to drain out the electrons in the hybridized state in order to acheive the population inversion for lasing, the wave function should have a large component from the GaSb layer. A thorough understanding of this hybridization is important not only for opto-electronic devices, but also for the fundamental theory of transport parallel to interfaces.
The optical property of a BGQW, which is of most theoretical interest and of most importance to opto-electronics, involves the optical transition between these two levels. Hence, relevant issues to be studied are the dependence of optical matrix elements on the degree of conduction-valence hybridization, the tuning of hybridization by varying the width of the GaSb layer, and the sensitivity of quantized levels to this tuning. In this paper we will perform a theoretical calculation of electronic structure and optical matrix elements in order to investigate these items. In Sec. II we present the model system we study and the basic assumptions in our theoretical calculation. Our results for the electronic structure and optical matrix elements are given in Sec. III and Sec. IV respectively, and discussed in Sec. V.
## II Model and theory
Our model system shown in Fig. 1 is a GaSb/InAs heterostructure sandwiched between two thick AlSb layers. The growth direction is which defines the $`z`$ axis. The $`x`$ axis is along and the $`y`$ axis is along . This system was studied earlier with a simple two band model . The essential feature of our model is the existence in BGQW quantized conduction band states in the InAs layer, and hybridized quantum levels in the energy regime where the GaSb valence band overlaps with the InAs conduction band. Therefore, for simplicity we can well assume that the two thick AlSb layers are bulk AlSb. As an illustration, one conduction band level and one hybridized level are plotted in Fig. 1. We will perform band structure calculation for the BGQW shown in Fig. 1. The lattice constants are 6.096 Å for GaSb, 6.058 Å for InAs, and 6.136 Å for AlSb. With such small lattice mismatch, the strain estimated with the deformation potential of these materials is not large enough to modify our conclusion qualitatively. Hence, in our theoretical calculation, strain will be neglected. As will be discussed in Sec. III, the essential feature of our model mentioned above, and hence the conclusion reached by our calculation will remain intact even this weak strain effect is taken into account.
We employ the effective bond orbital model (EBOM) which is a tight-binding-like model defined on the Bravais lattice of the underlying crystal. The basis are constructed as linear combinations of orbitals centered on each lattice site having $`s`$\- or $`p`$-symmetry, together with spin eigenstates. The so-constructed basis functions diagonalize the spin-orbit interaction. Our model Hamiltonian includes local and nearest neighbor matrix elements. The degree of inversion asymmetry of the BGQW heterostructure under study is much higher than that of individual InAs or GaSb constituent layers. Hence, in our calculation the effect of weak inversion symmetry is neglected for simplicity. On the other hand, the strong inversion symmetry of the BGQW and its combined effect with the spin-orbit coupling at finite in-plane wavevector are fully taken into account in our theoretical treatment.
It is worthwhile to point out that the EBOM is not an atomistic model like the tight binding model. While the tight binding model is built from orbitals representing the degree of freedom of each individual atom, the bond-orbitals in the EBOM are confined to a primitive unit cell and represent more than one atom. The symmetry of the EBOM is identical to that of the k$``$p model, and therefore in both models there is no asymmetry between the (110) and (1$`\overline{1}`$0) directions. On the other hand, the difference between the EBOM and the k$``$p model is that the k$``$p model is a continum model without a natural momentum cutoff, but the EBOM is discretized on the Bravais lattice of the crystal.
The bulk parameters of the EBOM, which are required for our calculation, are obtained by fitting the bulk band-structure of EBOM to that of the k$``$p Hamiltonian near the $`\mathrm{\Gamma }`$ point to the second order in $`k`$, as described in Ref . The EBOM can therefore be considered as an effective mass theory discretized on the Bravais lattice. For small k vectors the two models are equivalent. However, since EBOM is discretized on the Bravais lattice, it gives better results at large $`k`$. It is important to point out that usually, as in the present case, it is not even necessary to know the precise form of the effective bond orbitals.
Because of the large overlap between the InAs conduction band and the GaSb valence band, the split-off band must be included in our treatment. Hence, we will deal essentially an eight band Kane model. We consider the BGQW heterostructures grown in the direction, with $`n`$ labelling the (001) planes. Then, the position of a lattice site $`\stackrel{}{R}_n`$ in the $`n`$th plane can be expressed as $`(n\frac{a}{2},\stackrel{}{R}_n)`$, where $`a`$ is the lattice constant of the conventional unit cell. The eight bond orbitals at each lattice site are conventionally labelled as ($`S\frac{1}{2}\frac{1}{2}`$), ($`S\frac{1}{2}\frac{1}{2}`$), ($`P\frac{3}{2}\frac{3}{2}`$), ($`P\frac{3}{2}\frac{1}{2}`$), ($`P\frac{3}{2}\frac{1}{2}`$), ($`P\frac{3}{2}\frac{3}{2}`$), ($`P\frac{1}{2}\frac{1}{2}`$), and ($`P\frac{1}{2}\frac{1}{2}`$). We denote $`|n,\stackrel{}{R}_n,\alpha `$ for the $`\alpha `$th bond orbital at position $`(n\frac{a}{2},\stackrel{}{R}_n)`$. In terms of these bond orbitals, we define a basis set of planar orbitals
$$|n,\stackrel{}{k}_{},\alpha =N_{}^{1/2}\underset{\stackrel{}{R}_n}{}\mathrm{exp}(i\stackrel{}{k}_{}\stackrel{}{R}_n)|n,\stackrel{}{R}_n,\alpha ,$$
(1)
where $`\stackrel{}{k}_{}`$ is an in-plane wave vector, and $`N_{}`$ is the number of sites in each plane.
Using this basis of planar orbitals, the Hamiltonian
$`H={\displaystyle \underset{\stackrel{}{k}_{}}{}}H_\stackrel{}{k}_{}`$
is decomposed into a linear combination of partial Hamiltonians $`H_\stackrel{}{k}_{}`$. The partial Hamiltonian
$`H_\stackrel{}{k}_{}`$ $`=`$ $`{\displaystyle \underset{n\alpha \beta }{}}e_{n,\stackrel{}{k}_{},\alpha \beta }|n,\stackrel{}{k}_{},\alpha n,\stackrel{}{k}_{},\beta |`$ (3)
$`+{\displaystyle \underset{n\alpha \beta }{}}(v_{n,\stackrel{}{k}_{},\alpha \beta }|n+1,\stackrel{}{k}_{},\alpha n,\stackrel{}{k}_{},\beta |+h.c.)`$
is diagonal with respect to the in-plane wave vectors $`\stackrel{}{k}_{}`$, and block-tridiagonal with respect to the remaining quantum numbers. This form of Hamiltonian is suitable for numerical computations.
The matrix elements in $`H_\stackrel{}{k}_{}`$, except for those $`v_{n,\stackrel{}{k}_{},\alpha \beta }`$ which connect the two bond orbital planes forming an interface, are set to the values for the corresponding bulk materials. At an interface, we follow the commonly accepted approach to determine the value of $`v_{n,\stackrel{}{k}_{},\alpha \beta }`$ as the average of the parameter values of the bulk materials on each side of the interface. The materials parameter values from which the effective bond orbital parameters are derived are given in Table I. The valence band offsets are 0.56 eV for GaSb-InAs, 0.18 eV for AlSb-InAs, and -0.38 eV for AlSb-GaSb.
For a given $`\stackrel{}{k}_{}`$, the eight eigenfunctions
$`|\mathrm{\Psi }_{\gamma ,\stackrel{}{k}_{}}={\displaystyle \underset{n,\alpha }{}}F_{n,\stackrel{}{k}_{},\alpha }^\gamma |n,\stackrel{}{k}_{},\alpha ;\gamma =1,2,\mathrm{}`$
and the corresponding band energies $`E_{\gamma ,\stackrel{}{k}_{}}`$ can be readily obtained from Eq. (3) by diagonalizing a finite tridiagonal matrix for a BGQW of finite width. Knowing the eigensolutions, we will calculate the optical matrix elements.
The representation of the momentum operator in the bond-orbital basis contains the leading local and nearest neighbor matrix elements. They are determined by mapping the bulk matrix elements to the k$``$p results up to and including terms linear in wave vector . In the planar orbital basis, its $`\stackrel{}{k}_{}`$-diagonal part can be expressed as
$`P^\nu (\stackrel{}{k}_{})`$ $`=`$ $`{\displaystyle \underset{n\alpha \beta }{}}P_{n,\stackrel{}{k}_{},\alpha \beta }^\nu |n,\stackrel{}{k}_{},\alpha n,\stackrel{}{k}_{},\beta |`$ (5)
$`+{\displaystyle \underset{n\alpha \beta }{}}(Q_{n,\stackrel{}{k}_{},\alpha \beta }^\nu |n+1,\stackrel{}{k}_{},\alpha n,\stackrel{}{k}_{},\beta |+h.c.),`$
where $`\nu `$=$`x,y,z`$ is the polarization direction. The matrix elements are
$`P_{n,\stackrel{}{k}_{},\alpha \beta }^\nu `$ $`=`$ $`\sqrt{{\displaystyle \frac{2}{m_0}}}n,\stackrel{}{k}_{},\alpha |p_\nu |n,\stackrel{}{k}_{},\beta ,`$ (6)
$`Q_{n,\stackrel{}{k}_{},\alpha \beta }^\nu `$ $`=`$ $`\sqrt{{\displaystyle \frac{2}{m_0}}}n+1,\stackrel{}{k}_{},\alpha |p_\nu |n,\stackrel{}{k}_{},\beta ,`$ (7)
where $`m_0`$ is the free electron mass. In the above matrix elements we have included the prefactor $`\sqrt{2/m_0}`$ such that these matrix elements have units $`(eV)^{1/2}`$.
## III Electronic structure
In the following we will examine the electronic structure of the InAs-GaSb BGQW shown in Fig. 1. Before the InAs layers and the GaSb layers are coupled together, we will denote the $`n`$th electronic levels in the conduction band by En, the $`n`$th energy levels in the heavy-hole band by Hn, and the $`n`$th energy levels in the light-hole band by Ln. Such conventional labeling is unambiguous for $`\stackrel{}{k}_{}`$=0.
The symmetry properties of EBOM at $`\stackrel{}{k}_{}`$=0 allows the conduction band states in the InAs layers to hybridize with the light-hole band states in the GaSb layers, but not with the heavy-hole band states. We would like first to locate the region of BGQW structures in which we can fine tune the degree of such hybridization. For this purpose, we perform two energy level calculations; one with zero InAs layer width and the other with zero GaSb layer width. The results are plotted in Fig. 2 as functions of the number of atomic layers. We see that for 58 atomic layers, the GaSb L1 level overlaps with the InAs E1 level. Hence, the region around 60 atomic layers will be the reasonable starting BGQW structure for tuning the system into maximum hybridization. We should be aware of the fact that different confinement boundary conditions also affect the degree of hybridization, as will be seen in our computed results below.
We will then set 60 atomic layers for the InAs constituent and vary the GaSb width from 30 to 80 atomic layers. In the absence of interface coupling, the H1, E1, H2 and L1 levels at $`\stackrel{}{k}_{}`$=0 are shown in Fig. 3 as dashed curves (hole states are the same as in Fig. 2). In the BGQW sample with 80 atomic layers of GaSb, these four levels are ordered as H1$`>`$L1$`>`$H2$`>`$E1. When the interface coupling is turned, these four levels are shown in Fig. 3 as solid curves. For the convenience of description, we use the results at 80 atomic layers of GaSb to label them as H1$`>`$E1$`>`$H2$`>`$L1. Because of the symmetry properties at $`\stackrel{}{k}_{}`$=0, the conduction band states have negligible influence on H1 and H2 levels. The heavy hole levels are only slightly perturbed by the change of boundary conditions. On the other hand, the L1 and E1 levels are strongly hybridized and repel each other. The difference between the solid curves and the dashed curves is also due to the different confinement boundary conditions as mentioned above. In the region of GaSb width between 50 and 60 atomic layers, the separation of the two resulting hybridized levels is about 40-55 meV.
To demonstrate the spatial properties of hybridized states, we have calculated the real space occupation probability
$`𝒪_n^\gamma (\stackrel{}{k}_{}=0){\displaystyle \underset{\alpha }{}}|F_{n,\stackrel{}{k}_{}=0,\alpha }^\gamma |^2;\gamma =1,2,\mathrm{}`$
$`𝒪_n^\gamma (\stackrel{}{k}_{}`$=0) of five eigenstates in a BGQW with 56 GaSb atomic layers and 60 InAs atomic layers are plotted in Fig. 4. From top to bottom, the first state is E2 which is largely confined in InAs layers. The second and fourth state, which localized in GaSb layers, are H1 and H2 respectively. The third and the lowest states are E1-L1 hybridized states. To be consistent with our earlier convention, the third state is E1, and the lowest state is L1. One way to further analyze this issue is to decompose the total occupation probability into partial occupation probability for each $`\alpha `$ band. Such decomposition is shown in Fig. 5 with dotted curves for the E1 state, and the solid curves for the L1 state. In this figure, the two upper plots are for the $`P_{\frac{3}{2},\pm \frac{1}{2}}`$ components and the two lower plots are for the $`S_{\frac{1}{2},\pm \frac{1}{2}}`$ components. We see that the tail of the E1 (or L1) state in the GaSb (or InAs) region has the same component profile as its main part in the InAs (or GaAs) region. Such features were also obtained for InAs-GaSb superlattices in a very recent paper .
In our derivation of the eigensolutions of a BGQW heterostructure, the effect of strain has been neglected. The separations between the so-obtained eigenenergies for $`\stackrel{}{k}_{}`$=0, which are determined by the confinement potential, are substantially larger than the energy shifts produced by the strain. Hence, even the effect of strain is included in an improved calculation, the formation of E1-L1 hybridized states is still dominated by the confinement potential, and the degree of hybridization can still be tuned by changing the width of the GaSb and/or the InAs layers. In this respect, the conclusion reached with our present calculation will remain intact when the effect of strain is taken into account.
An interesting question is, when a charge carrier occupies a E1-L1 hybridized state, will its physical properties electron-like or hole-like? For example, what will be the cyclotron effective mass of such a charge carrier? One relevant quantity is the probability
$`𝒪_{\mathrm{InAs}}^\gamma (\stackrel{}{k}_{}=0){\displaystyle \underset{n\mathrm{InAs}}{}}𝒪_n^\gamma (\stackrel{}{k}_{}=0);\gamma =1,2,\mathrm{}`$
to find the electron in InAs layers when it occupies the $`\gamma `$th eigenstate. For the two hybridized states (labeled as E1 and L1 in our convention) in a InAs-GaSb BGQW with 60 InAs atomic layers, the results are plotted in Fig. 6 as function of the number of GaSb atomic layers. The solid curves are for $`𝒪_{\mathrm{GaSb}}^\gamma (\stackrel{}{k}_{}`$=0), and the dashed curves are for $`𝒪_{\mathrm{InAs}}^\gamma (\stackrel{}{k}_{}`$=0). For thin GaSb layers, the charge carrier in a hybridized state is mostly either electron-like or hole-like. However, as the width of GaSb layers increases towards 80 atomic layers, the characteristic features of the charge carrier remain to be studied.
The above calculations for zone center ($`\stackrel{}{k}_{}`$=0) states are repeated for finite $`\stackrel{}{k}_{}`$, and the dispersion relations of various subbands are shown in Fig. 7 for a BGQW with 60 InAs atomic layers and 60 GaSb atomic layers. We have chosen $`\stackrel{}{k}_{}`$ along the direction, or the $`x`$ axis. At zone center, from top to bottom the levels are E2, H1, E1, H2, L1, and L2. In region of finite $`\stackrel{}{k}_{}`$, the spin-splitting of levels and the anticrossing of levels are clearly seen. We will return to these dispersion relations later for further discussions.
We have mentioned that our BGQW system was studied earlier with a simple two-band model , where only the H1 and the E1 levels are included. Therefore, in Ref. there is no E1-L1 hybridization. The energy bands shown in Fig. 2 of Ref. corresponds to our E1 and H1 bands in Fig. 7, within the region 0$`<`$$`ka`$$`<`$0.12. The two-band calculation does yield an anticrossing at finite $`\stackrel{}{k}_{}`$, but gives no spin-orbit splitting.
## IV Optical matrix elements
Knowing the eigenfunctions of the BGQW heterostructure, the optical matrix elements are readily calculated with Eq. (5). Because for finite $`\stackrel{}{k}_{}`$ each dispersion curve is spin-orbit splited into two curves, we distinguish them as Ena and Enb for conduction band states, Hna and Hnb for heavy-hole band states, and Lna and Lnb for light-hole band states. For convenience, in the region of small $`\stackrel{}{k}_{}`$, the lower spin-orbit splited dispersion curve is assigned with a, and the higher one with b.
There are numerous optical transitions between each pair of states for different polarizations of the electromagnetic waves. The selection rules are complicated by the anticrossing of levels as well as the hybridization between InAs conduction band states and GaSb valence band states. Hence, here our study will focus on the cases of our interest: the optical transitions from the E2a and E2b levels to the E1a, E1b, H1a, and H1b levels. As will be discussed later, such transitions are relevant to the possible infrared lasers based on BGQW heterostructures. We would like to mention that parallel to our work, in a very recent paper , in InAs-GaSb superlattices the optical transition matrix elements at $`\stackrel{}{k}_{}`$=0 with in-plane polarization were studied as functions of the superlattice period.
By analyzing our extensive numerical results, we have reached the following selection rules for transitions from the E2a and E2b levels to the E1a, E1b, H1a, and H1b levels. For $`\stackrel{}{k}_{}`$ in the or the direction, among the two transitions from E2a (or E2b) to E1a or E1b, only one is allowed. Similarly, among the two transitions from E2a (or E2b) to H1a or H1b, only one is allowed. However, if $`\stackrel{}{k}_{}`$ is along a low symmetry direction in the two-dimensional Brillouin zone, such as for example, all transitions are allowed. It is worthwhile to point out that these selection rules for a BGQW, which is asymmetric, happen to be the same as the selection rules for the intersubband transitions in a symmetric conduction band quantum well, derived with a simplified eight band model .
The numerical results to be discussed below are obtained for a BGQW with 60 InAs atomic layers and 60 GaSb atomic layers, the band structure of which is given in Fig. 7. From Fig. 6 we see a significant E1-L1 hybridization in this BGQW heterostructure. The square of the amplitude of optical matrix elements with $`z`$-polarization are shown in Fig. 8. At $`\stackrel{}{k}_{}`$=0, the optical matrix element is zero for the E2$``$H1 transition, but is large for the E2$``$E1 transition, as expected from the symmetry properties. As $`\stackrel{}{k}_{}`$ increases, the anticrossing between the E1 band and the H1 band occurs around $`\stackrel{}{k}_{}a`$$``$0.75. Consequently, the E2$``$E1 transition drops sharply to zero, while the E2$``$H1 transition picks up its strength rapidly.
By changing the polarization of the electromagnetic wave and repeating the calculations of optical matrix elements, the results for $`y`$-polarization is shown in Fig. 9 and for $`x`$-polarization in Fig. 10. The strength of these transitions are substantially smaller than those for the $`z`$ polarization, by about two orders of magnitude. The strong wavevector dependence of the optical matrix elements in Figs. 9 and 10 is due to the combined effect of spin-orbit splitting and anticrossing of energy bands. The drastic difference between the curves in Fig. 9 and the corresponding curves in Fig. 10 indicates the strong anisotropy of the optical matrix elements with respect to in-plane polarizations.
## V Discussion
The very unusual feature of the BGQW heterostructures is the new eigenstates formed by the hybridization of conduction band states at one side of the interface with the valence band states at the other side. We made a thorough investigation on the formation of such states and their influence on the optical properties of a BGQW. However, their influence on the transport properties parallel to interfaces is perhaps a more important issue for fundamental study. If one tunes the the system to have the Fermi energy lying in a conduction-valence hybridized two-dimensional energy band, depending on the degree of hybridization, the parallel transport properties can change from completely electron-like to complete hole-like.
Another relevant issue for fundamental study is the spin-orbit splitting. Recently, there has been much interest in the effect of spin-orbit interactions in two-dimensional electron gases. For example, the Shubnikov-de Haas measurements of spin-orbit splitting have been performed on symmetric InAs/GaSb quantum wells and on symmetric InAs/AlSb quantum wells, in which the origin of spin-orbit interaction is the lack of inversion symmetry in bulk crystal structure. Because of the asymmetry of the BGQW itself, very large spin-orbit splittings appear in the energy bands as shown in Fig. 7. Consequently, BGQW heterostructures are good candidate for investigating phenomena which are related to spin-orbit interaction.
BGQW heterostructures have great potential in technological applications for lasers and detectors, tunnable in infrared wavelengths . Our calculations can serve as the theoretical modeling of such infrared lasers. To make a bipolar laser similar to the conventional quantum lasers, in Fig. 1 the AlSb layers at the InAs side should be $`n`$-doped, and the AlSb layers at the GaSb side should be $`p`$-doped. The number of InAs atomic layers should be around 60, and the number of GaSb atomic layers should be between 70 and 80. Then, from Fig. 6 we see a strong E1-L1 hybridyzation. By adjusting the acceptor concentration, we can set the Fermi level between the E1 level and the H2 level shwon in Fig. 4. When such BGQW heterostructure is connected to external circuit, electrons which are injected into the E2 energy band will relax to the zone center and make radiation transitions to the E1 band due to the large optical matrix elements of $`z`$-polarized light. Then, the strong conduction-valence hybridization of the E1 band states provides the rapid draining of electrons from the E1 band into external circuit via the GaSb valence band states. From Fig. 7 we estimate the energy of emitted photons to be 0.18-0.2 eV, corresponding to wavelength about 6 micrometers. Around this wavelength there exist infrared windows, and consequently such radiation source will be extremely useful.
To close this paper, we must mention the very recent work of Ref. in which the electronic structure and optical matrix elements at the zone center ($`\stackrel{}{k}_{}`$=0) of InAs-GaSb superlattices have been calculated with both a plane wave pseudopotential method and the eight band k$``$p method. In that paper the effects of superlattice period on various physical quantities at zone center were studied in details. In particular, the authors of Ref. found a zone center E1-H1 coupling manifested by band anticrossing at superlattice period $`n`$=28, as well as a zone center L1-H2 coupling and anticrossing around $`n`$=13. Such features are absent in the k$``$p method because it fails to recognize the atomistic details in no-common-atom superlattice. How these features will affect quantitatively the 2D dispersion relation, and hence our results of the formation of E1-L1 hybridization and the in-plane physical properties of a single BGQW heterostructure, as well as their impact on the theoretical modeling of infrared lasers and detectors, remains to be an open question.
###### Acknowledgements.
This work was supported by the Norwegian Research Council (NFR) under grant no. 111071/431.
|
warning/0003/hep-ph0003133.html
|
ar5iv
|
text
|
# 1 Introduction
## 1 Introduction
The success of the standard model (SM) predictions is remarkably high and, indeed, to some extent, even beyond what we theorists would have expected. A common view before LEP started operating was that some new physics related to the electroweak symmetry breaking should show up when precisions at the percent level on some electroweak observable could be reached. As we know, on the contrary, even reaching sensitivities better than the percent has not given rise to any indication of departure from the SM predictions. All that can be summarized in a powerful statement about the “low–energy” limit of any kind of new physics beyond the SM: such new physics has to reproduce the SM with great accuracy when we consider its limit at energy scales of the order of the electroweak scale.
The fact that with the SM we have a knowledge of fundamental interactions up to energies of $`𝒪(100)`$ GeV should not be underestimated: it represents a tremendous and astonishing success of our gauge theory approach in particle physics and it is clear that it represents one of the great achievements in a century of great conquests in physics. Having said that, we are now confronting ourselves with an embarrassing question: if the SM is so extraordinarily good, does it make sense do go beyond it? The answer, in our view, is certainly positive. This “yes” is not only motivated by what we could define “philosophical” reasons, but there are specific motivations pushing us beyond the SM: we will group them in two broad categories, theoretical and “observational” reasons.
### 1.1 Theoretical reasons for new physics
The major theoretical conundra of the SM are related the following issues: flavor problem, unification of the fundamental interactions and gauge hierarchy problem. We briefly remind what they are about.
Flavor problem. All the masses and mixings of fermions are just free (unpredicted) parameters in the SM. To be sure, there is not even any hint in the SM about the number and rationale of fermion families. Leaving aside predictions for individual masses, not even a rough relation among fermion masses within the same generation or among different generations is present.
Unification of forces. At the time of the Fermi theory we had two couplings to describe the electromagnetic and the weak interactions (the electric constant and the Fermi constant, respectively). In the SM we are trading off those two couplings with two new couplings, the gauge couplings of $`SU(2)`$ and $`U(1)`$. Moreover, the gauge coupling of the strong interactions is very different from the other two. We cannot say that the SM represents a true unification of fundamental interactions, even leaving aside the problem that gravity is not considered at all by the model.
Gauge hierarchy. Fermion and vector boson masses are “protected” by symmetries in the SM (i.e., their masses can arise only when we break certain symmetries). On the contrary, the Higgs scalar mass does not enjoy such a symmetry protection. We would expect such mass to naturally jump to some higher scale where new physics sets in (this new energy scale could be some grand unification scale or the Planck mass, for instance). The only way to keep the Higgs mass at the electroweak scale is to perform incredibly accurate fine tunings of the parameters of the scalar sector.
### 1.2 “Observational” reasons for new physics
We have already said that all the experimental particle physics results of these last years have marked one success after the other of the SM. What do we mean then by “observational” difficulties for the SM? It is curious that such difficulties do not arise from observations within the strict particle physics domain, but rather they originate from possible “clashes” of the particle physics SM with the standard model of cosmology (i.e., the Hot Big Bang) or the standard model of the Sun. Explicitly we have in mind the following points.
Dark Matter. Denoting with $`\mathrm{\Omega }`$ the ratio of the energy density to the critical energy density, the problem of dark matter (DM) can be summarized in the following two numbers: $`\mathrm{\Omega }_{DM}=0.3`$ and $`\mathrm{\Omega }_B<0.1`$. The first number denotes the amount of the contribution to $`\mathrm{\Omega }`$ due to DM as inferred from measurements at the level of cluster of galaxies. The upper bound denotes the highest contribution of baryonic matter to $`\mathrm{\Omega }`$ to have compatibility with one of the main pillars of the Big Bang model: nucleosynthesis. The clash between the above two numbers underlines the fact that we definitely need a large amount of non–baryonic DM. In the electroweak SM, no viable non–baryonic candidate exists to fulfill this task (remember that in the SM neutrinos are strictly massless). Hence, the existence of a (large) amount of non–baryonic DM pushes us to introduce new particles in addition to those of the SM.
Baryogenesis. Given that we have strong evidence that the Universe is vastly matter–antimatter asymmetric (i.e. no sizeable amount of primordial antimatter has survived), it is appealing to have a dynamical mechanism to give rise to such large baryon–antibaryon asymmetry starting from a symmetric situation. In the SM it is not possible to have such an efficient mechanism for baryogenesis. In spite of the fact that at the quantum level sphaleronic interactions violate baryon number in the SM, such violation cannot lead to the observed large matter–antimatter asymmetry (both CP violation is too tiny in the SM and, also, the present experimental lower bounds on the Higgs mass do not allow for a conveniently strong electroweak phase transition). Hence, a dynamical baryogenesis calls for the presence of new particles and interactions beyond the SM (successful mechanisms for baryogenesis in the context of new physics beyond the SM are well known).
Inflation. Several serious cosmological problems (flatness, causality, age of the Universe, …) are beautifully solved if the early Universe underwent some period of exponential expansion (inflation). The SM with its Higgs doublet does not succeed to originate such an inflationary stage. Again some extensions of the SM, where in particular new scalar fields are introduced, are able to produce a temporary inflation of the early Universe.
### 1.3 The SM as an effective low–energy theory
The above theoretical and “observational” arguments strongly motivate us to go beyond the SM. On the other hand, the clear success of the SM in reproducing all the known phenomenology up to energies of the order of the electroweak scale is telling us that the SM has to be recovered as the low–energy limit of such new physics. Indeed, it may even well be the case that we have a “tower” of underlying theories which show up at different energy scales.
If we accept the above point of view, we may try to find signals of new physics considering the SM as a truncation to renormalizable operators of an effective low–energy theory which respects the $`SU(3)\times SU(2)\times U(1)`$ symmetry and whose fields are just those of the SM. The renormalizable (i.e. of canonical dimension less or equal to four) operators giving rise to the SM enjoy three crucial properties which have no reason to be shared by generic operators of dimension larger than four. They are the conservation (at any order in perturbation theory) of Baryon (B) and Lepton (L) numbers and an adequate suppression of Flavor Changing Neutral Current (FCNC) processes through the GIM mechanism.
Now consider the new physics (directly above the SM in the “tower” of new physics theories) to have a typical energy scale $`\mathrm{\Lambda }`$. In the low–energy effective Lagrangian, such scale appears with a positive power only in the quadratic scalar term (scalar mass) and in the dimension zero operator which can be considered a cosmological constant. Notice that $`\mathrm{\Lambda }`$ cannot appear in dimension three operators related to fermion masses because chirality forbids direct fermion mass terms in the Lagrangian. Then, in all operators of dimension larger than four, $`\mathrm{\Lambda }`$ will show up in the denominator with powers increasing with the dimension of the corresponding operator.
The crucial question that all of us, theorists and experimentalists, ask ourselves is: where is $`\mathrm{\Lambda }`$? Namely is it close to the electroweak scale (i.e. not much above $`100`$ GeV) or is $`\mathrm{\Lambda }`$ of the order of the grand unification scale or the Planck scale? B– and L–violating processes and FCNC phenomena represent a potentially interesting clue to answer this fundamental question.
Take $`\mathrm{\Lambda }`$ to be close to the electroweak scale. Then we may expect non–renormalizable operators with B, L and flavor violations not to be largely suppressed by the presence of powers of $`\mathrm{\Lambda }`$ in the denominator. Actually this constitutes, in general, a formidable challenge for any model builder who wants to envisage new physics close to $`M_W`$. Theories with dynamical breaking of the electroweak symmetry (technicolour) and low–energy supersymmetry constitute examples of new physics with a “small” $`\mathrm{\Lambda }`$. In these lectures we will see that the above general considerations on potentially large B, L and flavor violations apply to the SUSY case (it is well–known that FCNC represent a major problem also in technicolour schemes).
Alternatively, given the abovementioned potential danger of having a small $`\mathrm{\Lambda }`$, one may feel it safer to send $`\mathrm{\Lambda }`$ to super–large values. Apart from kind of “philosophical” objections related to the unprecedented gap of many orders of magnitude without any new physics, the above discussion points out a typical problem of this approach. Since the quadratic scalar terms have a coefficient in front scaling with $`\mathrm{\Lambda }^2`$, we expect all scalar masses to be of the order of the super–large scale $`\mathrm{\Lambda }`$. This is the gauge hierarchy problem, and it constitutes the main (if not only) reason to believe that SUSY should be a low–energy symmetry.
Notice that the fact that SUSY should be a fundamental symmetry of Nature (something of which we have little doubt given the “beauty” of this symmetry) does not imply by any means that SUSY should be a low–energy symmetry, namely that it should hold unbroken down to the electroweak scale. SUSY may well be present in Nature but be broken at some very large scale (Planck scale or string compactification scale). In that case SUSY would be of no use in tackling the gauge hierarchy problem and its phenomenological relevance would be practically zero. On the other hand, if we invoke SUSY to tame the growth of the scalar mass terms with the scale $`\mathrm{\Lambda }`$, then we are forced to take the view that SUSY should hold as a good symmetry down to a scale $`\mathrm{\Lambda }`$ close to the electroweak scale. Then B, L and FCNC may be useful for us to shed some light on the properties of the underlying theory from which the low–energy SUSY Lagrangian resulted. Let us add that there is an independent argument in favor of this view that SUSY should be a low–energy symmetry. The presence of SUSY partners at low energy creates the conditions to have a correct unification of the strong and electroweak interactions. If they were at $`M_{\mathrm{Planck}}`$ and the SM were all the physics up to super–large scales, the program of achieving such a unification would largely fail, unless one complicates the non–SUSY GUT scheme with a large number of Higgs representations and/or a breaking chain with intermediate mass scales is invoked.
In the above discussion, we stressed that we are not only insisting on the fact that SUSY should be present at some stage in Nature, but we are asking for something much more ambitious: we are asking for SUSY to be a low–energy symmetry, namely it should be broken at an energy scale as low as the electroweak symmetry breaking scale. This fact can never be overestimated. There are indeed several reasons pushing us to introduce SUSY: it is the most general symmetry compatible with a local, relativistic quantum field theory, it softens the degree of divergence of the theory, it looks promising for a consistent quantum description of gravity together with the other fundamental interactions. However, all these reasons are not telling us where we should expect SUSY to be broken. For that matter, we could even envisage the maybe “natural” possibility that SUSY is broken at the Planck scale. What is relevant for phenomenology is that the gauge hierarchy problem and, to some extent, the unification of the gauge couplings are actually forcing us to ask for SUSY to be unbroken down to the electroweak scale, hence implying that the SUSY copy of all the known particles, the so–called s–particles should have a mass in the $`100`$$`1000`$ GeV mass range. If LEP and Tevatron are not going to see any SUSY particle, at least the advent of LHC will be decisive in establishing whether low–energy SUSY actually exists or it is just a fruit of our (ingenious) speculations. Although even after LHC, in case of a negative result for the search of SUSY particles, we will not be able to “mathematically” exclude all the points of the SUSY parameter space, we will certainly be able to very reasonably assess whether the low–energy SUSY proposal makes sense or not.
Before the LHC (and maybe Tevatron) direct searches for SUSY signals we should ask ourselves whether we can hope to have some indirect manifestation of SUSY through virtual effects of the SUSY particles.
We know that, in the past, virtual effects (i.e. effects due to the exchange of yet unseen particles in the loops) were precious in leading us to major discoveries, like the prediction of the existence of the charm quark or the heaviness of the top quark long before its direct experimental observation. Here we focus on the potentialities of SUSY virtual effects in processes which are particluraly suppressed (or sometime even forbidden) in the SM; the flavor changing neutral current phenomena and the processes where CP violation is violated.
## 2 Flavor, CP and New Physics
The generation of fermion masses and mixings (“flavor problem”) gives rise to a first and important distinction among theories of new physics beyond the electroweak standard model.
One may conceive a kind of new physics which is completely “flavor blind”, i.e. new interactions which have nothing to do with the flavor structure. To provide an example of such a situation, consider a scheme where flavor arises at a very large scale (for instance the Planck mass) while new physics is represented by a supersymmetric extension of the SM with supersymmetry broken at a much lower scale and with the SUSY breaking transmitted to the observable sector by flavor–blind gauge interactions. In this case, one may think that the new physics does not cause any major change to the original flavor structure of the SM, namely that the pattern of fermion masses and mixings is compatible with the numerous and demanding tests of flavor changing neutral currents.
Alternatively, one can conceive a new physics which is entangled with the flavor problem. As an example consider a technicolour scheme where fermion masses and mixings arise through the exchange of new gauge bosons which mix together ordinary fermions and technifermions. Here we expect (correctly enough) new physics to have potential problems in accommodating the usual fermion spectrum with the adequate suppression of FCNC. As another example of new physics which is not flavor blind, take a more conventional SUSY model which is derived from a spontaneously broken N=1 supergravity and where the SUSY breaking information is conveyed to the ordinary sector of the theory through gravitational interactions. In this case we may expect that the scale at which flavor arises and the scale of SUSY breaking are not so different and possibly the mechanism itself of SUSY breaking and transmission is flavor–dependent. Under these circumstances, we may expect a potential flavor problem to arise, namely that SUSY contributions to FCNC processes are too large.
### 2.1 The Flavor Problem in SUSY
The potentiality of probing SUSY in FCNC phenomena was readily realized when the era of SUSY phenomenology started in the early 80’s . In particular, the major implication that the scalar partners of quarks of the same electric charge but belonging to different generations had to share a remarkably high mass degeneracy was emphasized.
Throughout the large amount of work in this last decade, it became clearer and clearer that generically talking of the implications of low–energy SUSY on FCNC may be rather misleading. Even in the Minimal SUSY extension of the SM (MSSM) from the point of view of the particle content, we have a host of different situations. The so–called Constrained Minimal Supersymmetric Standard Model (CMSSM) is the simplest possibility, and the FCNC contributions can be computed in terms of a very limited set of unknown new SUSY parameters. Remarkably enough, this minimal model succeeds to pass all the set of FCNC tests unscathed. To be sure, it is possible to severely constrain the SUSY parameter space, for instance using $`bs\gamma `$, in a way which is complementary to what is achieved by direct SUSY searches at colliders.
However, the CMSSM is by no means equivalent to low–energy SUSY. A first sharp distinction concerns the mechanism of SUSY breaking and transmission to the observable sector which is chosen. As we mentioned above, in models with gauge–mediated SUSY breaking (GMSB models ) it may be possible to avoid the FCNC threat “ab initio” (notice that this is not an automatic feature of this class of models, but it depends on the specific choice of the sector which transmits the SUSY breaking information, the so–called messenger sector). The other more “canonical” class of SUSY theories (including also CMSSM) has gravitational messengers and a very large scale at which SUSY breaking occurs. In this talk we will focus only on this class of gravity–mediated SUSY breaking models. Even sticking to this more limited choice, we have a variety of options with very different implications for the flavor problem.
First, there exists an interesting large class of SUSY realizations where the customary R–parity (which is invoked to suppress proton decay) is replaced by other discrete symmetries which allow either baryon or lepton violating terms in the superpotential. But, even sticking to the more orthodox view of imposing R–parity, we are still left with a large variety of extensions of the MSSM at low energy. The point is that low–energy SUSY “feels” the new physics at the super–large scale at which supergravity (i.e., local supersymmetry) broke down. In this last couple of years, we have witnessed an increasing interest in supergravity realizations without the so–called flavor universality of the terms which break SUSY explicitly. Another class of low–energy SUSY realizations, which differ from the MSSM in the FCNC sector, is obtained from SUSY–GUT’s. The interactions involving super–heavy particles in the energy range between the GUT and the Planck scale bear important implications for the amount and kind of FCNC that we expect at low energy.
Even when R–parity is imposed, the FCNC challenge is not over. It is true that in this case, analogously to what happens in the SM, no tree level FCNC contributions arise. However, it is well–known that this is a necessary but not sufficient condition to consider the FCNC problem overcome. The loop contributions to FCNC in the SM exhibit the presence of the GIM mechanism and we have to make sure that in the SUSY case with R parity some analog of the GIM mechanism is active.
To give a qualitative idea of what we mean by an effective super–GIM mechanism, let us consider the following simplified situation where the main features emerge clearly. Consider the SM box diagram responsible for the $`K^0`$$`\overline{K}^0`$ mixing and take only two generations, i.e. only the up and charm quarks run in the loop. In this case, the GIM mechanism yields a suppression factor of $`𝒪((m_c^2m_u^2)/M_W^2)`$. If we replace the W boson and the up quarks in the loop with their SUSY partners and we take, for simplicity, all SUSY masses of the same order, we obtain a super–GIM factor which looks like the GIM one with the masses of the superparticles instead of those of the corresponding particles. The problem is that the up and charm squarks have masses which are much larger than those of the corresponding quarks. Hence the super–GIM factor tends to be of $`𝒪(1)`$ instead of being $`𝒪(10^3)`$ as it is in the SM case. To obtain this small number we would need a high degeneracy between the mass of the charm and up squarks. It is difficult to think that such a degeneracy may be accidental. After all, since we invoked SUSY for a naturalness problem (the gauge hierarchy issue), we should avoid invoking a fine–tuning to solve its problems! Then, one can turn to some symmetry reason. For instance, just sticking to this simple example that we are considering, one may think that the main bulk of the charm and up squark masses is the same, i.e. the mechanism of SUSY breaking should have some universality in providing the mass to these two squarks with the same electric charge. Flavor universality is by no means a prediction of low–energy SUSY. The absence of flavor universality of soft–breaking terms may result from radiative effects at the GUT scale or from effective supergravities derived from string theory. Indeed, from the point of view of these effective supergravity theories, it may appear more natural not to have such flavor universality. To obtain it one has to invoke particular circumstances, like, for instance, strong dilaton over moduli dominance in the breaking of supersymmetry, something which is certainly not expected on general ground.
Another possibility one may envisage is that the masses of the squarks are quite high, say above few TeV’s. Then, even if they are not so degenerate in mass, the overall factor in front of the four–fermion operator responsible for the kaon mixing becomes smaller and smaller (it decreases quadratically with the mass of the squarks) and, consequently, one can respect the observational result. We see from this simple example that the issue of FCNC may be closely linked to the crucial problem of the way we break SUSY.
We now turn to some general remarks about the worries and hopes that CP violation arises in the SUSY context.
### 2.2 CP Violation in SUSY
CP violation has major potentialities to exhibit manifestations of new physics beyond the standard model. Indeed, it is quite a general feature that new physics possesses new CP violating phases in addition to the Cabibbo–Kobayashi–Maskawa (CKM) phase $`\left(\delta _{\mathrm{CKM}}\right)`$ or, even in those cases where this does not occur, $`\delta _{\mathrm{CKM}}`$ shows up in interactions of the new particles, hence with potential departures from the SM expectations. Moreover, although the SM is able to account for the observed CP violation in the kaon system, we cannot say that we have tested so far the SM predictions for CP violation. The detection of CP violation in $`B`$ physics will constitute a crucial test of the standard CKM picture within the SM. Again, on general grounds, we expect new physics to provide departures from the SM CKM scenario for CP violation in $`B`$ physics. A final remark on reasons that make us optimistic in having new physics playing a major role in CP violation concerns the matter–antimatter asymmetry in the universe. Starting from a baryon–antibaryon symmetric universe, the SM is unable to account for the observed baryon asymmetry. The presence of new CP–violating contributions when one goes beyond the SM looks crucial to produce an efficient mechanism for the generation of a satisfactory $`\mathrm{\Delta }`$B asymmetry.
The above considerations apply well to the new physics represented by low–energy supersymmetric extensions of the SM. Indeed, as we will see below, supersymmetry introduces CP violating phases in addition to $`\delta _{\mathrm{CKM}}`$ and, even if one envisages particular situations where such extra–phases vanish, the phase $`\delta _{\mathrm{CKM}}`$ itself leads to new CP–violating contributions in processes where SUSY particles are exchanged. CP violation in $`B`$ decays has all potentialities to exhibit departures from the SM CKM picture in low–energy SUSY extensions, although, as we will discuss, the detectability of such deviations strongly depends on the regions of the SUSY parameter space under consideration.
In any MSSM, at least two new “genuine” SUSY CP–violating phases are present. They originate from the SUSY parameters $`\mu `$, $`M`$, $`A`$ and $`B`$. The first of these parameters is the dimensionful coefficient of the $`H_uH_d`$ term of the superpotential. The remaining three parameters are present in the sector that softly breaks the N=1 global SUSY. $`M`$ denotes the common value of the gaugino masses, $`A`$ is the trilinear scalar coupling, while $`B`$ denotes the bilinear scalar coupling. In our notation, all these three parameters are dimensionful. The simplest way to see which combinations of the phases of these four parameters are physical is to notice that for vanishing values of $`\mu `$, $`M`$, $`A`$ and $`B`$ the theory possesses two additional symmetries . Indeed, letting $`B`$ and $`\mu `$ vanish, a $`U(1)`$ Peccei–Quinn symmetry originates, which in particular rotates $`H_u`$ and $`H_d`$. If $`M`$, $`A`$ and $`B`$ are set to zero, the Lagrangian acquires a continuous $`U(1)`$ $`R`$ symmetry. Then we can consider $`\mu `$, $`M`$, $`A`$ and $`B`$ as spurions which break the $`U(1)_{PQ}`$ and $`U(1)_R`$ symmetries. In this way, the question concerning the number and nature of the meaningful phases translates into the problem of finding the independent combinations of the four parameters which are invariant under $`U(1)_{PQ}`$ and $`U(1)_R`$ and determining their independent phases. There are three such independent combinations, but only two of their phases are independent. We use here the commonly adopted choice:
$$\phi _A=\mathrm{arg}\left(A^{}M\right),\phi _B=\mathrm{arg}\left(B^{}M\right).$$
(1)
where also $`\mathrm{arg}\left(B\mu \right)=0`$, i.e. $`\phi _\mu =\phi _B`$.
The main constraints on $`\phi _A`$ and $`\phi _B`$ come from their contribution to the electric dipole moments of the neutron and of the electron. For instance, the effect of $`\phi _A`$ and $`\phi _B`$ on the electric and chromoelectric dipole moments of the light quarks ($`u`$, $`d`$, $`s`$) lead to a contribution to $`d_N^e`$ of order
$$d_N^e2\left(\frac{100\mathrm{G}\mathrm{e}\mathrm{V}}{\stackrel{~}{m}}\right)^2\mathrm{sin}\phi _{A,B}\times 10^{23}\mathrm{e}\mathrm{cm},$$
(2)
where $`\stackrel{~}{m}`$ here denotes a common mass for squarks and gluinos. The present experimental bound, $`d_N^e<1.1^{25}`$ e cm, implies that $`\phi _{A,B}`$ should be $`<10^2`$, unless one pushes SUSY masses up to $`𝒪`$(1 TeV). A possible caveat to such an argument calling for a fine–tuning of $`\phi _{A,B}`$ is that uncertainties in the estimate of the hadronic matrix elements could relax the severe bound in Eq. (2) .
In view of the previous considerations, most authors dealing with the MSSM prefer to simply put $`\phi _A`$ and $`\phi _B`$ equal to zero. Actually, one may argue in favor of this choice by considering the soft breaking sector of the MSSM as resulting from SUSY breaking mechanisms which force $`\phi _A`$ and $`\phi _B`$ to vanish. For instance, it is conceivable that both $`A`$ and $`M`$ originate from one same source of $`U(1)_R`$ breaking. Since $`\phi _A`$ “measures” the relative phase of $`A`$ and $`M`$, in this case it would “naturally”vanish. In some specific models, it has been shown that through an analogous mechanism also $`\phi _B`$ may vanish.
If $`\phi _A=\phi _B=0`$, then the novelty of SUSY in CP violating contributions merely arises from the presence of the CKM phase in loops where SUSY particles run . The crucial point is that the usual GIM suppression, which plays a major role in evaluating $`\epsilon _K`$ and $`\epsilon ^{}/\epsilon `$ in the SM, in the MSSM case (or more exactly in the CMSSM) is replaced by a super–GIM cancellation which has the same “power” of suppression as the original GIM (see previous section). Again, also in the CMSSM, as it is the case in the SM, the smallness of $`\epsilon _K`$ and $`\epsilon ^{}/\epsilon `$ is guaranteed not by the smallness of $`\delta _{\mathrm{CKM}}`$, but rather by the small CKM angles and/or small Yukawa couplings. By the same token, we do not expect any significant departure of the CMSSM from the SM predictions also concerning CP violation in $`B`$ physics. As a matter of fact, given the large lower bounds on squark and gluino masses, one expects relatively tiny contributions of the SUSY loops in $`\epsilon _K`$ or $`\epsilon ^{}/\epsilon `$ in comparison with the normal $`W`$ loops of the SM. Let us be more detailed on this point.
In the CMSSM, the gluino exchange contribution to FCNC is subleading with respect to chargino ($`\chi ^\pm `$) and charged Higgs ($`H^\pm `$) exchanges. Hence, when dealing with CP violating FCNC processes in the CMSSM with $`\phi _A=\phi _B=0`$, one can confine the analysis to $`\chi ^\pm `$ and $`H^\pm `$ loops. If one takes all squarks to be degenerate in mass and heavier than $`200`$ GeV, then $`\chi ^\pm \stackrel{~}{q}`$ loops are obviously severely penalized with respect to the SM $`W^+`$$`q`$ loops (remember that at the vertices the same CKM angles occur in both cases).
The only chance for the CMSSM to produce some sizeable departure from the SM situation in CP violation is in the particular region of the parameter space where one has light $`\stackrel{~}{q}`$, $`\chi ^\pm `$ and/or $`H^\pm `$. The best candidate (indeed the only one unless $`\mathrm{tan}\beta m_t/m_b`$) for a light squark is the stop. Hence one can ask the following question: can the CMSSM present some novelties in CP–violating phenomena when we consider $`\chi ^+`$$`\stackrel{~}{t}`$ loops with light $`\stackrel{~}{t}`$, $`\chi ^+`$ and/or $`H^+`$?
Several analyses in the literature tackle the above question or, to be more precise, the more general problem of the effect of light $`\stackrel{~}{t}`$ and $`\chi ^+`$ on FCNC processes . A first important observation concerns the relative sign of the $`W^+`$$`t`$ loop with respect to the $`\chi ^+`$$`\stackrel{~}{t}`$ and $`H^+`$$`t`$ contributions. As it is well known, the latter contribution always interferes positively with the SM one. Interestingly enough, in the region of the MSSM parameter space that we consider here, also the $`\chi ^+`$$`\stackrel{~}{t}`$ contribution interferes constructively with the SM contribution. The second point regards the composition of the lightest chargino, i.e. whether the gaugino or higgsino component prevails. This is crucial since the light stop is predominantly $`\stackrel{~}{t}_R`$ and, hence, if the lightest chargino is mainly a wino, it couples to $`\stackrel{~}{t}_R`$ mostly through the $`LR`$ mixing in the stop sector. Consequently, a suppression in the contribution to box diagrams going as $`\mathrm{sin}^4\theta _{LR}`$ is present ($`\theta _{LR}`$ denotes the mixing angle between the lighter and heavier stops). On the other hand, if the lightest chargino is predominantly a higgsino (i.e. $`M_2\mu `$ in the chargino mass matrix), then the $`\chi ^+`$–lighter $`\stackrel{~}{t}`$ contribution grows. In this case, contributions $`\theta _{LR}`$ become negligible and, moreover, it can be shown that they are independent on the sign of $`\mu `$. A detailed study is provided in reference . For instance, for $`M_2/\mu =10`$, they find that the inclusion of the SUSY contribution to the box diagrams doubles the usual SM contribution for values of the lighter $`\stackrel{~}{t}`$ mass up to $`100`$$`120`$ GeV, using $`\mathrm{tan}\beta =1.8`$, $`M_{H^+}=100`$ TeV, $`m_\chi =90`$ GeV and the mass of the heavier $`\stackrel{~}{t}`$ of 250 GeV. However, if $`m_\chi `$ is pushed up to 300 GeV, the $`\chi ^+`$$`\stackrel{~}{t}`$ loop yields a contribution which is roughly 3 times less than in the case $`m_\chi =90`$ GeV, hence leading to negligible departures from the SM expectation. In the cases where the SUSY contributions are sizeable, one obtains relevant restrictions on the $`\rho `$ and $`\eta `$ parameters of the CKM matrix by making a fit of the parameters $`A`$, $`\rho `$ and $`\eta `$ of the CKM matrix and of the total loop contribution to the experimental values of $`\epsilon _K`$ and $`\mathrm{\Delta }M_{B_d}`$. For instance, in the above–mentioned case in which the SUSY loop contribution equals the SM $`W^+`$$`t`$ loop, hence giving a total loop contribution which is twice as large as in the pure SM case, combining the $`\epsilon _K`$ and $`\mathrm{\Delta }M_{B_d}`$ constraints leads to a region in the $`\rho `$$`\eta `$ plane with $`0.15<\rho <0.40`$ and $`0.18<\eta <0.32`$, excluding negative values of $`\rho `$.
In conclusion, the situation concerning CP violation in the MSSM case with $`\phi _A=\phi _B=0`$ and exact universality in the soft–breaking sector can be summarized in the following way: the MSSM does not lead to any significant deviation from the SM expectation for CP–violating phenomena as $`d_N^e`$, $`\epsilon _K`$, $`\epsilon ^{}/\epsilon `$ and CP violation in $`B`$ physics; the only exception to this statement concerns a small portion of the MSSM parameter space where a very light $`\stackrel{~}{t}`$ ($`m_{\stackrel{~}{t}}<100`$ GeV) and $`\chi ^+`$ ($`m_\chi 90`$ GeV) are present. In this latter particular situation, sizeable SUSY contributions to $`\epsilon _K`$ are possible and, consequently, major restrictions in the $`\rho `$$`\eta `$ plane can be inferred. Obviously, CP violation in $`B`$ physics becomes a crucial test for this MSSM case with very light $`\stackrel{~}{t}`$ and $`\chi ^+`$. Interestingly enough, such low values of SUSY masses are at the border of the detectability region at LEP II.
In next Section, we will move to the case where, still keeping the minimality of the model, we switch on the new CP violating phases. Later on we will give up also the strict minimality related to the absence of new flavor structure in the SUSY breaking sector and we will see that, in those more general contexts, we can expect SUSY to significantly depart from the SM predictions in CP violating phenomena.
## 3 Flavor Blind SUSY Breaking and CP Violation
We have seen in the previous section that in any MSSM there are additional phases which can cause deviations from the predictions of the SM in CP violation experiments. In fact, in the CMSSM, there are already two new phases present, Eq.(1), and for most of the MSSM parameter space, the experimental bounds on the electric dipole moments (EDM) of the electron and neutron constrain these phases to be at most $`𝒪(10^2)`$. However, in the last few years, the possibility of having non–zero SUSY phases has again attracted a great deal of attention. Several new mechanisms have been proposed to suppress supersymmetric contributions to EDMs below the experimental bounds while allowing SUSY phases $`𝒪(1)`$. Methods of suppressing the EDMs consist of cancellation of various SUSY contributions among themselves , non universality of the soft breaking parameters at the unification scale and approximately degenerate heavy sfermions for the first two generations . In the presence of one of these mechanisms, large supersymmetric phases are naturally expected and EDMs should be generally close to the experimental bounds. <sup>1</sup><sup>1</sup>1In a more general (and maybe more natural) MSSM there are many other CP violating phases that contribute to CP violating observables.
In this section we will study the effects of these phases in CP violation observables as $`\epsilon _K`$, $`\epsilon ^{}/\epsilon `$ and $`B^0`$ CP asymmetries. In particular we will show that the presence of large susy phases is not enough to produce sizeable supersymmetric contributions to these observables. In fact, in the absence of the CKM phase, a general MSSM with all possible phases in the soft–breaking terms, but no new flavor structure beyond the usual Yukawa matrices, can never give a sizeable contribution to $`\epsilon _K`$, $`\epsilon ^{}/\epsilon `$ or hadronic $`B^0`$ CP asymmetries. However, we will see in the next section, that as soon as one introduces some new flavor structure in the soft Susy–breaking sector, even if the CP violating phases are flavor independent, it is indeed possible to get sizeable CP contribution for large Susy phases and $`\delta _{CKM}=0`$. Then, we can rephrase our sentence above in a different way: A new result in hadronic $`B^0`$ CP asymmetries in the framework of supersymmetry would be a direct proof of the existence of a completely new flavor structure in the soft–breaking terms. This means that $`B`$–factories will probe the flavor structure of the supersymmetry soft–breaking terms even before the direct discovery of the supersymmetric partners .
To prove this we will consider any MSSM, i.e. with the minimal supersymmetric particle content, with general complex soft–breaking terms, but with a flavor structure strictly given by the two familiar Yukawa matrices or any matrix strictly proportional to them. In these conditions, the most general structure of the soft–breaking terms at the large scale, that we call $`M_{GUT}`$, is,
$`(m_Q^2)_{ij}=m_Q^2\delta _{ij}(m_U^2)_{ij}=m_U^2\delta _{ij}(m_D^2)_{ij}=m_D^2\delta _{ij}`$
$`(m_L^2)_{ij}=m_L^2\delta _{ij}(m_E^2)_{ij}=m_E^2\delta _{ij}m_{H_1}^2m_{H_2}^2`$
$`m_{\stackrel{~}{g}}e^{i\phi _3}m_{\stackrel{~}{W}}e^{i\phi _2}m_{\stackrel{~}{B}}e^{i\phi _1}(A_U)_{ij}=A_Ue^{i\phi _{A_U}}(Y_U)_{ij}`$
$`(A_D)_{ij}=A_De^{i\phi _{A_D}}(Y_D)_{ij}(A_E)_{ij}=A_Ee^{i\phi _{A_E}}(Y_E)_{ij}.`$ (3)
where all the allowed phases are explicitly written and one of them can be removed by an R–rotation. All other numbers or matrices in this equation are always real. Notice that this structure covers, not only the CMSSM , but also most of Type I string motivated models considered so far from phenomenology , gauge mediated models , minimal effective supersymmetry models , etc.
Experiments of CP violation in the $`K`$ or $`B`$ systems only involve supersymmetric particles as virtual particles in the loops. This means that the phases in the soft–breaking terms can only appear in these experiments through the mass matrices of the SUSY particles. Then, the key point in our discussion will be the role played by the SUSY phases and the soft–breaking terms flavor structure in the low–energy sparticle mass matrices.
It is important to notice that, even in a model with flavor–universal soft–breaking terms at some high energy scale, as this is the case, some off–diagonality in the squark mass matrices appears at the electroweak scale. Working on the the so–called Super CKM basis (SCKM), where squarks are rotated parallel to the quarks so that Yukawa matrices are diagonalized, the squark mass matrix is not flavor diagonal at $`M_W`$. This is due to the fact that at $`M_{GUT}`$ there are always two non–trivial flavor structures, namely the two Yukawa matrices for the up and down quarks, not simultaneously diagonalizable. This implies that through RGE evolution some flavor mixing leaks into the sfermion mass matrices. In a general Supersymmetric model, the presence of new flavor structures in the soft breaking terms would generate large flavor mixing in the sfermion mass matrices. However, in the CMSSM, the two Yukawa matrices are the only source of flavor change. Always in the SCKM basis, any off–diagonal entry in the sfermion mass matrices at $`M_W`$ will be necessarily proportional to a product of Yukawa couplings. Then, a typical estimate for the element $`(i,j)`$ in the $`L`$$`L`$ down squark mass matrix at the electroweak scale would necessarily be (see for details),
$`(m_{}^{2}{}_{LL}{}^{(D)})_{ij}cm_Q^2Y_{ik}^uY_{jk}^{u}{}_{}{}^{},`$ (4)
with $`c`$ a proportionality factor between $`0.1`$ and 1. This rough estimate provides the order of magnitude of the different entries in the sfermion mass matrices. It is important to notice that if the phases of these elements were $`𝒪(1)`$, due to some of the phases in equation (3), we would be able to give sizeable contributions, or even saturate, the different CP observables . Then, it is clear that the relevant question for CP violation experiments is the presence of imaginary parts in these off–diagonal entries.
As explained in , once we have solved the Yukawa RGEs, the RGE equations of all soft–breaking terms are a set of linear differential equations. Then, they can be solved as a linear function of the initial conditions. For instance the scalar masses are,
$`m_S^2(M_W)={\displaystyle \underset{i}{}}\eta _S^{(\varphi _i)}m_{\varphi _i}^2+{\displaystyle \underset{ij}{}}\left(\eta _S^{(g_ig_j)}e^{i(\phi _i\phi _j)}+\eta _S^{(g_ig_j)T}e^{i(\phi _i\phi _j)}\right)m_{g_i}m_{g_j}`$
$`+{\displaystyle \underset{i}{}}\eta _S^{(g_i)}m_{g_i}^2+{\displaystyle \underset{ij}{}}\left(\eta _S^{(g_iA_j)}e^{i(\phi _i\phi _{A_j})}+\eta _S^{(g_iA_j)T}e^{i(\phi _i\phi _{A_j})}\right)m_{g_i}A_j`$
$`+{\displaystyle \underset{i}{}}\eta _S^{(A_i)}A_i^2+{\displaystyle \underset{ij}{}}\left(\eta _S^{(A_iA_j)}e^{i(\phi _{A_i}\phi _{A_j})}+\eta _S^{(A_iA_j)T}e^{i(\phi _{A_i}\phi _{A_j})}\right)A_iA_j`$
where $`S=Q,U,D`$, $`\varphi _i`$ refers to any scalar, $`g_i`$ to the different gauginos and $`A_i`$ to any tri–linear coupling. In this equation, the different $`\eta `$ matrices are $`3\times 3`$ matrices, strictly real and all the allowed phases have been explicitly written. Regarding the imaginary parts, due to the hermiticity of the sfermion mass matrices, any imaginary part will always be associated to the non–symmetric part of the $`\eta _S^{(g_ig_j)}`$, $`\eta _S^{(A_iA_j)}`$ or $`\eta _S^{(g_iA_j)}`$ matrices. To estimate the size of these anti–symmetric parts, we can go to the RGE equations for the scalar mass matrices, where we use the same conventions and notation as in . Taking advantage of the linearity of these equations, we can directly write the evolution of the anti–symmetric parts, for instance $`\widehat{m}_Q^2=m_Q^2(m_Q^2)^T`$, as,
$`{\displaystyle \frac{d\widehat{m}_Q^2}{dt}}=`$ $`[{\displaystyle \frac{1}{2}}(\stackrel{~}{Y}_U\stackrel{~}{Y}_U^{}+\stackrel{~}{Y}_D\stackrel{~}{Y}_D^{})\widehat{m}_Q^2+{\displaystyle \frac{1}{2}}\widehat{m}_Q^2(\stackrel{~}{Y}_U\stackrel{~}{Y}_U^{}+\stackrel{~}{Y}_D\stackrel{~}{Y}_D^{})+`$ (6)
$`2i\text{Im}\{\stackrel{~}{A}_U\stackrel{~}{A}_U^{}+\stackrel{~}{A}_D\stackrel{~}{A}_D^{}\}+\stackrel{~}{Y}_U\widehat{m}_U^2\stackrel{~}{Y}_U^{}+\stackrel{~}{Y}_D\widehat{m}_D^2\stackrel{~}{Y}_D^{}]`$
where, due to the reality of Yukawa matrices, we have used $`Y^T=Y^{}`$, and following a tilde over the couplings ($`\stackrel{~}{Y}`$, $`\stackrel{~}{A}`$, …) denotes a re–scaling by a factor $`1/(4\pi )`$. The evolution of the $`R`$$`R`$ squark mass matrices, $`m_U^2`$ and $`m_D^2`$, is completely analogous. With the initial conditions in equation (3), $`\widehat{m}_Q^2`$, $`\widehat{m}_U^2`$ and $`\widehat{m}_D^2`$ at $`M_{GUT}`$ are identically zero. Then, we can safely neglect the last two terms in equation (6) because they will only be a second order effect. This means that the only source for $`\widehat{m}_Q^2`$ in equation (6) is necessarily $`\text{Im}\{A_UA_U^{}+A_DA_D^{}\}`$ (also for $`\widehat{m}_{U,D}^2`$) .
The next step is then to analyze the RGE for the tri–linear couplings,
$`{\displaystyle \frac{d\stackrel{~}{A}_U}{dt}}={\displaystyle \frac{1}{2}}\left({\displaystyle \frac{16}{3}}\stackrel{~}{\alpha }_3+3\stackrel{~}{\alpha }_2+{\displaystyle \frac{1}{9}}\stackrel{~}{\alpha }_1\right)\stackrel{~}{A}_U\left({\displaystyle \frac{16}{3}}\stackrel{~}{\alpha }_3M_3+3\stackrel{~}{\alpha }_2M_2+{\displaystyle \frac{1}{9}}\stackrel{~}{\alpha }_1M_1\right)\stackrel{~}{Y}_U`$
$`(2\stackrel{~}{A}_U\stackrel{~}{Y}_U^{}\stackrel{~}{Y}_U+3Tr(\stackrel{~}{A}_U\stackrel{~}{Y}_U^{})\stackrel{~}{Y}_U+{\displaystyle \frac{5}{2}}\stackrel{~}{Y}_U\stackrel{~}{Y}_U^{}\stackrel{~}{A}_U+{\displaystyle \frac{3}{2}}Tr(\stackrel{~}{Y}_U\stackrel{~}{Y}_U^{})\stackrel{~}{A}_U+`$
$`\stackrel{~}{A}_D\stackrel{~}{Y}_D^{}\stackrel{~}{Y}_U+{\displaystyle \frac{1}{2}}\stackrel{~}{Y}_D\stackrel{~}{Y}_D^{}\stackrel{~}{A}_U)`$ (7)
with an equivalent equation for $`A_D`$. With the general initial conditions in equation (3), $`A_U`$ is complex at any scale. However, we are interested in the imaginary parts of $`A_UA_U^{}`$. At $`M_{GUT}`$ this combination is exactly real, but, due to different renormalization of different elements of the matrix, this is not true anymore at a different scale. Nevertheless, a careful analysis of equation (3) is enough to convince ourselves that these imaginary parts are extremely small. Let us, for a moment, neglect the terms involving $`\stackrel{~}{A}_D\stackrel{~}{Y}_D^{}`$ or $`\stackrel{~}{Y}_D\stackrel{~}{Y}_D^{}`$ from the above equation or, strictly speaking, from the complete set of MSSM RGE. Then, the only flavor structure appearing in equation (3) at $`M_{GUT}`$ is $`Y_U`$. We can always go to the basis where $`Y_U`$ is diagonal and then we will have $`A_U`$ exactly diagonal at any scale. In particular this means that $`\text{Im}\{A_UA_U^{}\}`$ would always exactly vanish. A completely parallel reasoning can be applied to $`A_D`$ and $`\text{Im}\{A_DA_D^{}\}`$. Hence, simply taking into account the flavor structure, our conclusion is that, necessarily, any non–vanishing element of $`\text{Im}[A_UA_U^{}+A_DA_D^{}]`$ and hence of $`\widehat{m}_Q^2`$ must be proportional to $`(\stackrel{~}{Y}_D\stackrel{~}{Y}_D^{}\stackrel{~}{Y}_U\stackrel{~}{Y}_U^{}H.C.)`$. So, we can expect them to be,
$`(\widehat{m}_Q^2)_{i<j}K(Y_DY_D^{}Y_UY_U^{}H.C.)_{i<j}`$
$`(\widehat{m}_Q^2)_{12}K\mathrm{cos}^2\beta (h_sh_t\lambda ^5)`$
$`(\widehat{m}_Q^2)_{13}K\mathrm{cos}^2\beta (h_bh_t\lambda ^3)`$
$`(\widehat{m}_Q^2)_{23}K\mathrm{cos}^2\beta (h_bh_t\lambda ^2),`$ (8)
where $`h_i=m_i^2/v^2`$, with $`v=\sqrt{v_1^2+v_2^2}`$ the vacuum expectation value of the Higgs, $`\lambda =\mathrm{sin}\theta _c`$ and $`K`$ is a proportionality constant that includes the effects of the running from $`M_{GUT}`$ to $`M_W`$. To estimate this constant, we have to keep in mind that the imaginary parts of $`A_UA_U^{}`$ are generated through the RGE running and then, these imaginary parts generate $`\widehat{m}_Q^2`$ as a second order effect. This means that roughly $`K𝒪(10^2)`$ times a combination of initial conditions as in equation (3). So, we estimate these matrix elements to be $`(\mathrm{cos}^2\beta \{10^{12},6\times 10^8,3\times 10^7\})`$ times initial conditions. This was exactly the result we found for the $`A`$$`g`$ terms in in the framework of the CMSSM. In fact, now it is clear that this is the same for all the terms in equation (3), $`g_i`$$`A_j`$, $`g_i`$$`g_j`$ and $`A_i`$$`A_j`$, irrespectively of the presence of an arbitrary number of new phases. This discussion can be directly applied for the $`R`$$`R`$ matrices.
Hence, so far, we have shown that the $`L`$$`L`$ or $`R`$$`R`$ squark mass matrices are still essentially real. The only complex matrices, then, will still be the $`L`$$`R`$ matrices that include, from the very beginning, the phases $`\phi _{A_i}`$ and $`\phi _\mu `$. Once more, the size of these entries is determined by the Yukawa elements with these two phases providing the complex structure. In fact, we can follow the same reasoning used after Eq.(3). In the absence of the terms involving $`\stackrel{~}{A}_D\stackrel{~}{Y}_D^{}`$ or $`\stackrel{~}{Y}_D\stackrel{~}{Y}_D^{}`$, the $`\stackrel{~}{A}_U`$ matrix would be exactly diagonalized when we diagonalize the Yukawa matrices. So, any off–diagonal element in the SCKM basis will be proportional to three Yukawas, $`(\stackrel{~}{Y}_D\stackrel{~}{Y}_D^{}\stackrel{~}{Y}_U)`$ and hence sufficiently small. Notice that this situation is not new for these more general MSSM models and it was already present even in the CMSSM. We can conclude, then, that the structure of the sfermion mass matrices at $`M_W`$ is not modified from the familiar structure already present in the CMSSM, irrespective of the presence of an arbitrary number of new SUSY phases.
### 3.1 Indirect CP violation
Next, we will consider indirect CP violation both in the $`K`$ and $`B`$ systems. In the SM, neutral meson mixing arises at one loop through the well–known $`W`$–box. However, in the MSSM, there are new contributions to $`\mathrm{\Delta }F=2`$ processes coming from boxes mediated by supersymmetric particles. These are: charged Higgs boxes ($`H^\pm `$), chargino boxes ($`\chi ^\pm `$) and gluino–neutralino boxes ($`\stackrel{~}{g}`$, $`\chi ^0`$). $``$$`\overline{}`$ mixing is correctly described by the $`\mathrm{\Delta }F=2`$ effective Hamiltonian, $`_{eff}^{\mathrm{\Delta }F=2}`$, which can be decomposed as,
$`_{eff}^{\mathrm{\Delta }F=2}={\displaystyle \frac{G_F^2M_W^2}{(2\pi )^2}}(K_{td}^{}K_{tq})^2(C_1(\mu )Q_1(\mu )+C_2(\mu )Q_2(\mu )+C_3(\mu )Q_3(\mu )).`$
With the relevant four–fermion operators given by
$`Q_1=\overline{d}_L^\alpha \gamma ^\mu q_L^\alpha \overline{d}_L^\beta \gamma _\mu q_L^\beta ,Q_2=\overline{d}_L^\alpha q_R^\alpha \overline{d}_L^\beta q_R^\beta ,Q_3=\overline{d}_L^\alpha q_R^\beta \overline{d}_L^\beta q_R^\alpha ,`$ (10)
where $`q=s,b`$ for the $`K`$ and $`B`$–systems respectively, $`\alpha ,\beta `$ are color indices and $`K_{ij}`$ is the CKM mixing matrix. In the CMSSM, these are the only three operators present in the limit of vanishing $`m_d`$. The Wilson coefficients, $`C_1(\mu )`$, $`C_2(\mu )`$ and $`C_3(\mu )`$, receive contributions from the different supersymmetric boxes,
$`C_1(M_W)`$ $`=`$ $`C_1^W(M_W)+C_1^H(M_W)+C_1^{\stackrel{~}{g},\chi ^0}(M_W)+C_1^\chi (M_W)`$
$`C_2(M_W)`$ $`=`$ $`C_2^H(M_W)+C_2^{\stackrel{~}{g}}(M_W)`$
$`C_3(M_W)`$ $`=`$ $`C_3^{\stackrel{~}{g},\chi ^0}(M_W)+C_3^\chi (M_W)`$ (11)
Both, the usual SM $`W`$–box and the charged Higgs box contribute to these operators. However, with $`\delta _{CKM}=0`$, these contributions do not contain any complex phase and hence cannot generate an imaginary part for these Wilson coefficients.
Then, Gluino and neutralino contributions are specifically supersymmetric. They involve the superpartners of quarks and gauge bosons. Here, the source of flavor mixing is not directly the usual CKM matrix. It is the presence of off–diagonal elements in the sfermion mass matrices, as discussed in section 3. To analyze these contributions, it is convenient to use the so–called Mass Insertion (MI) approximation . To define the MI we go to the SCKM basis. In this basis, off–diagonal flavor–changing effects can be estimated by insertion of flavor–off–diagonal components of the mass–squared matrices. By normalizing by an average squark mass–squared $`m_{\stackrel{~}{q}}^2`$, we define,
$`(\delta _{LL}^d)_{ij}={\displaystyle \frac{(m_{}^{2}{}_{LL}{}^{(d)})_{ij}}{m_{\stackrel{~}{q}}^2}}(\delta _{RR}^d)_{ij}={\displaystyle \frac{(m_{}^{2}{}_{RR}{}^{(d)})_{ij}}{m_{\stackrel{~}{q}}^2}}(\delta _{LR}^d)_{ij}={\displaystyle \frac{(m_{}^{2}{}_{LR}{}^{(d)})_{ij}}{m_{\stackrel{~}{q}}^2}}`$ (12)
with $`m_{}^{2}{}_{AB}{}^{(d)}`$ the squark mass matrices in the SCKM basis.
From the point of view of CP violation, we will always need a complex Wilson coefficient. In the SCKM basis all gluino vertices are flavor diagonal and real. Then, a complex MI in one of the sfermion lines is always required. Only $`L`$$`L`$ mass insertions enter at first order in the Wilson coefficient $`C_{1}^{}{}_{}{}^{\stackrel{~}{g},\chi ^0}(M_W)`$. From equation (3), the imaginary parts of these MI are at most $`𝒪(10^6)`$ for the $`b`$$`s`$ transitions and smaller otherwise . Comparing these values with the phenomenological bounds required to saturate the measured values of these processes we can easily see that we are always several orders of magnitude below.
In the case of the Wilson coefficients $`C_{2}^{}{}_{}{}^{\stackrel{~}{g}}(M_W)`$ and $`C_{3}^{}{}_{}{}^{\stackrel{~}{g}}(M_W)`$, the involved MI are $`L`$$`R`$. However, as explained in section 3, these MI are always suppressed by light masses of right handed squark plus two additional up Yukawas. Moreover, in the case of $`b`$$`s`$ transitions they are directly constrained by the $`bs\gamma `$ decay. Hence, gluino boxes, in the absence of new flavor structures, can never give sizeable contributions to indirect CP violation processes .
The chargino contributions to these Wilson coefficients were discussed in great detail in the CMSSM framework in reference . In this more general MSSM, we find very similar results due to the absence of new flavor structure.
Basically, in the chargino boxes, flavor mixing comes explicitly from the CKM mixing matrix, although off–diagonality in the sfermion mass matrix introduces a small additional source of flavor mixing.
$`C_1^\chi (M_W)={\displaystyle \underset{i,j=1}{\overset{2}{}}}{\displaystyle \underset{k,l=1}{\overset{6}{}}}{\displaystyle \underset{\alpha \gamma \alpha ^{}\gamma ^{}}{}}{\displaystyle \frac{K_{\alpha ^{}d}^{}K_{\alpha q}K_{\gamma ^{}d}^{}K_{\gamma q}}{(K_{td}^{}K_{tq})^2}}`$
$`[G^{(\alpha ,k)i}G^{(\alpha ^{},k)j}G^{(\gamma ^{},l)i}G^{(\gamma ,l)j}Y_1(z_k,z_l,s_i,s_j)]`$ (13)
where $`K_{\alpha q}G^{(\alpha ,k)i}`$ represent the coupling of chargino and squark $`k`$ to left–handed down quark $`q`$, $`z_k=M_{\stackrel{~}{u}_k}^2/M_W^2`$ and $`s_i=M_{\stackrel{~}{\chi }_i}^2/M_W^2`$. The explicit expressions for the loop functions can be found in reference . These couplings, in terms of the standard mixing matrices ,
$`G^{(\alpha ,k)i}`$ $`=`$ $`\left(\mathrm{\Gamma }_{UL}^{k\alpha }V_{i1}^{}{\displaystyle \frac{m_\alpha }{\sqrt{2}M_W\mathrm{sin}\beta }}\mathrm{\Gamma }_{UR}^{k\alpha }V_{i2}^{}\right).`$ (14)
$`G^{(\alpha ,k)i}`$ are in general complex, as both $`\phi _\mu `$ and $`\phi _{A_i}`$ are present in the different mixing matrices.
The main part of $`C_1^\chi `$ in equation (3.1) will be given by pure CKM flavor mixing, neglecting the additional flavor mixing in the squark mass matrix . This means, $`\alpha =\alpha ^{}`$ and $`\gamma =\gamma ^{}`$. In these conditions, using the symmetry of the loop function $`Y_1(a,b,c,d)`$ under the exchange of any two indices it is easy to prove that $`C_1^\chi `$ would be exactly real . This is not exactly true either in the CMSSM or in our more general MSSM, where there is additional flavor change in the sfermion mass matrices. Here, some imaginary parts appear in the $`C_1^\chi `$ in equation (3.1). In figure 1 we show in a scatter plot the size of imaginary and real parts of $`C_1^\chi `$ in the B system for a fixed value of $`\mathrm{tan}\beta =40`$. We see that this Wilson coefficient is always real up to a part in $`10^3`$. In any case, this is out of reach for the foreseen B–factories. For the K system, imaginary parts are still smaller due to smaller mixing angles with the stop.
Finally, chargino boxes contribute also to the quirality changing Wilson coefficient $`C_3^\chi (M_W)`$,
$`C_3^\chi (M_W)={\displaystyle \underset{i,j=1}{\overset{2}{}}}{\displaystyle \underset{k,l=1}{\overset{6}{}}}{\displaystyle \underset{\alpha \gamma \alpha ^{}\gamma ^{}}{}}{\displaystyle \frac{K_{\alpha ^{}d}^{}K_{\alpha q}K_{\gamma ^{}d}^{}K_{\gamma q}}{(K_{td}^{}K_{tq})^2}}{\displaystyle \frac{m_q^2}{2M_W^2\mathrm{cos}^2\beta }}`$
$`H^{(\alpha ,k)i}G^{(\alpha ^{},k)j}G^{(\gamma ^{},l)i}H^{(\gamma ,l)j}Y_2(z_k,z_l,s_i,s_j)`$ (15)
where $`m_q/(\sqrt{2}M_W\mathrm{cos}\beta )K_{\alpha q}H^{(\alpha ,k)i}`$ is the coupling of chargino and squark to the right–handed down quark $`q`$ ,
$$H^{(\alpha ,k)i}=U_{i2}\mathrm{\Gamma }_{UL}^{k\alpha }.$$
(16)
Unlike the $`C_1^\chi `$ Wilson coefficient, due to the differences between $`H`$ and $`G`$ couplings, $`C_3^\chi `$ is complex even in the absence of intergenerational mixing in the sfermion mass matrices . Then, the presence of these small flavor violating entries in the up–squark mass matrix hardly modifies the results obtained in their absence . In fact, in spite the presence of the Yukawa coupling squared, $`m_q^2/(2M_W^2\mathrm{cos}^2\beta )`$, this contribution could be relevant in the large $`\mathrm{tan}\beta `$ regime. For instance, in $`B^0`$$`\overline{B}^0`$ mixing we have $`m_b^2/(2M_W^2\mathrm{cos}^2\beta )`$ that for $`\mathrm{tan}\beta >25`$ is larger than 1 and so, it is not suppressed at all when compared with the $`C_1^\chi `$ Wilson Coefficient. This means that this contribution can be very important in the large $`\mathrm{tan}\beta `$ regime and could have observable effects in CP violation experiments in the new B–factories. However, we will show next that when we include the constraints coming from $`bs\gamma `$ these chargino contributions are also reduced to an unobservable level.
The chargino contributes to the $`bs\gamma `$ decay through the Wilson coefficients $`𝒞_7`$ and $`𝒞_8`$, corresponding to the photon and gluon dipole penguins respectively . In the large $`\mathrm{tan}\beta `$ regime, we can approximate these Wilson coefficients as ,
$`𝒞_7^{\chi ^\pm }(M_W)={\displaystyle \underset{k=1}{\overset{6}{}}}{\displaystyle \underset{i=1}{\overset{2}{}}}{\displaystyle \underset{\alpha ,\beta =u,c,t}{}}{\displaystyle \frac{K_{\alpha b}K_{\beta s}^{}}{K_{tb}K_{ts}^{}}}{\displaystyle \frac{m_b}{\sqrt{2}M_W\mathrm{cos}\beta }}`$
$`H^{(\alpha ,k)i}G_{}^{}{}_{}{}^{(\beta ,k)i}{\displaystyle \frac{M_{\chi ^i}}{m_b}}F_R^7(z_k,s_i)`$
$`𝒞_8^{\chi ^\pm }(M_W)={\displaystyle \underset{k=1}{\overset{6}{}}}{\displaystyle \underset{i=1}{\overset{2}{}}}{\displaystyle \underset{\alpha ,\beta =u,c,t}{}}{\displaystyle \frac{K_{\alpha b}K_{\beta s}^{}}{K_{tb}K_{ts}^{}}}{\displaystyle \frac{m_b}{\sqrt{2}M_W\mathrm{cos}\beta }}`$
$`H^{(\alpha ,k)i}G_{}^{}{}_{}{}^{(\beta ,k)i}{\displaystyle \frac{M_{\chi ^i}}{m_b}}F_R^8(z_k,s_i)`$ (17)
Now, if we compare the chargino contributions to these Wilson coefficients and to the coefficient $`C_3`$, equations (3.1) and (17), we can see that they are deeply related. In fact, in the approximation where the two different loop functions involved are of the same order, we have,
$`C_3(M_W)(𝒞_7(M_W))^2{\displaystyle \frac{m_q^2}{M_W^2}}`$ (18)
In figure 2, we show a scatter plot of the allowed values of $`Re(𝒞_7)`$ versus $`Im(𝒞_7)`$ in the CMSSM for a fixed value of $`\mathrm{tan}\beta =40`$ with the constraints from the decay $`BX_s\gamma `$ taken from the reference . Notice that a relatively large value of $`\mathrm{tan}\beta `$, for example, $`\mathrm{tan}\beta >10`$, is needed to compensate the $`W`$ and charged Higgs contributions and cover the whole allowed area with positive and negative values. However, the shape of the plot is clearly independent of $`\mathrm{tan}\beta `$, only the number of allowed points and its location in the allowed area depend on the value considered. Then, figure 3 shows the allowed values for a re–scaled Wilson coefficient $`\overline{C}_3(M_W)=M_W^2/m_q^2C_3(M_W)`$ corresponding to the same allowed points of the SUSY parameter space in figure 2. As we anticipated previously, the allowed values for $`\overline{C}_3`$ are close to the square of the values of $`𝒞_7`$ in figure 2 slightly scaled by different values of the loop functions.
We can immediately translate this result to a constraint on the size of the chargino contributions to $`\epsilon _{}`$.
$`\epsilon _{}={\displaystyle \frac{G_F^2M_W^2}{4\pi ^2\sqrt{2}\mathrm{\Delta }M_{}}}{\displaystyle \frac{(K_{td}K_{tq})^2}{24}}F_{}^2M_{}`$
$`{\displaystyle \frac{M_{}^2}{m_q^2(\mu )+m_d^2(\mu )}}\eta _3(\mu )B_3(\mu )Im[C_3]`$ (19)
In this expression $`M_{}`$, $`\mathrm{\Delta }M_{}`$ and $`F_{}`$ denote the mass, mass difference and decay constant of the neutral meson $`^0`$. The coefficient $`\eta _3(\mu )=2.93`$ includes the RGE effects from $`M_W`$ to the meson mass scale, $`\mu `$, and $`B_3(\mu )`$ is the B–parameter associated with the matrix element of the $`Q_3`$ operator .
For the $`K`$ system, using the experimentally measured value of $`\mathrm{\Delta }M_K`$ we obtain,
$`\epsilon _K^\chi =1.7\times 10^2{\displaystyle \frac{m_s^2}{M_W^2}}Im[\overline{C}_3]0.4\times 10^7Im[\overline{C}_3]`$ (20)
Given the allowed values of $`\overline{C}_3`$ in figure 3, this means that in the MSSM, even with large SUSY phases, chargino cannot produce a sizeable contribution to $`\epsilon _K`$.
The case of $`B^0`$$`\overline{B}^0`$ mixing has a particular interest due to the arrival of new data from the B–factories. In fact, in the large $`\mathrm{tan}\beta `$ regime chargino contributions to indirect CP violation can be very important. However, for any value of $`\mathrm{tan}\beta `$, we must satisfy the bounds from the $`bs\gamma `$ decay. Then, if we apply these constraints to the $`B^0`$$`\overline{B}^0`$ mixing,
$`\epsilon _B^\chi =0.17{\displaystyle \frac{m_b^2}{M_W^2}}Im[\overline{C}_3]0.5\times 10^3Im[\overline{C}_3]`$ (21)
where once again, with the allowed values of figure 3 we get a very small contribution to CP violation in the mixing. We must take into account that the mixing–induced CP phase, $`\theta _M`$, measurable in $`B^0`$ CP asymmetries, is related to $`\epsilon _B`$ by $`\theta _M=\mathrm{arcsin}\{2\sqrt{2}\epsilon _B\}`$. The expected sensitivities on the CP phases at the B factories are around $`\pm 0.1`$ radians, so this supersymmetric chargino contribution will be absolutely out of reach.
### 3.2 Direct CP violation
To complete our analysis, we consider now direct CP violation. In this case, the different decay processes are described by a $`\mathrm{\Delta }F=1`$ effective Hamiltonian. A complete operator basis for these transitions in a general MSSM involves 14 different operators . The main difference with the case of indirect CP violation is that these operators receive contributions both from box and penguin diagrams. Nevertheless, the discussion of the presence of imaginary parts is completely analogous to the case of indirect CP violation.
Once more, in the gluino case, $`L`$$`L`$ transitions are real to a very good approximation, and several orders of magnitude below the phenomenological bounds . On the other hand, $`L`$$`R`$ transitions are suppressed by two up Yukawas and a down quark mass and $`bs\gamma `$ decay. This is always true for the squark mass matrices obtained in section 3, and valid both for boxes and penguins.
Finally, we are left with chargino contributions. The analysis of chargino boxes is exactly the same as in the previous section. In fact, even the Wilson coefficients are identical except some CKM elements that can always be factored out. Then, for the penguins, $`L`$$`L`$ transitions are exactly real if we neglect inter–generational mixing in the squark mass matrices. Taking into account this small mixing we find, for the very same reasons as in the indirect CP violation case, that imaginary parts are far too small. The relation of the $`bs\gamma `$ decay with the $`L`$$`R`$ chargino penguins is in this case even more transparent than for the boxes. So, our conclusion is again that no new supersymmetric CP violation effects are possible in $`\epsilon ^{}/\epsilon `$ or hadronic $`B`$ CP asymmetries.
However, there is still one possibility to observe the effects of the new supersymmetric phases even in the absence of new flavor structure. We have seen that the reason for the smallness of the contributions of chargino $`L`$$`R`$ transitions is the experimental bound from the $`BX_s\gamma `$ branching ratio. This bound makes the chirality changing transitions, although complex, too small to compete with $`L`$$`L`$ transitions. Hence, in these conditions, just the processes where only chirality changing operators contribute (EDMs or $`bs\gamma `$), or observables where chirality flip operators are relevant ($`bsl^+l^{}`$) can show the effects of new supersymmetric phases .
## 4 CP Violation in the presence of new Flavor Structures
In section 3, we have shown that CP violation effects are always small in models with flavor blind soft–breaking terms. However, as soon as one introduces some new flavor structure in the soft breaking sector, it is indeed possible to get sizeable CP contribution for large Susy phases and $`\delta _{CKM}=0`$ . To show this, we will mainly concentrate in new supersymmetric contributions to $`\epsilon ^{}/\epsilon `$.
In the CMSSM, the SUSY contribution to $`\epsilon ^{}/\epsilon `$ is small . However in a MSSM with a more general framework of flavor structure it is relatively easy to obtain larger SUSY effects to $`\epsilon ^{}/\epsilon `$. In ref. it was shown that such large SUSY contributions arise once one assumes that: i) hierarchical quark Yukawa matrices are protected by flavor symmetry, ii) a generic dependence of Yukawa matrices on Polonyi/moduli fields is present (as expected in many supergravity/superstring theories), iii) the Cabibbo rotation originates from the down–sector and iv) the phases are of order unity. In fact, in , it was illustrated how the observed $`\epsilon ^{}/\epsilon `$ could be mostly or entirely due to the SUSY contribution.
The universality of the breaking is a strong assumption and is known not to be true in many supergravity and string inspired models . In these models, we expect at least some non–universality in the squark mass matrices or tri–linear terms at the supersymmetry breaking scale. Hence, sizeable flavor–off-diagonal entries will appear in the squark mass matrices. In this regard, gluino contributions to $`\epsilon ^{}/\epsilon `$ are especially sensitive to $`(\delta _{12}^d)_{LR}`$; even $`|\mathrm{Im}(\delta _{12}^d)_{LR}^2|10^5`$ gives a significant contribution to $`\epsilon ^{}/\epsilon `$ while keeping the contributions from this MI to $`\mathrm{\Delta }m_K`$ and $`\epsilon _K`$ well bellow the phenomenological bounds. The situation is the opposite for $`L`$$`L`$ and $`R`$$`R`$ mass insertions; the stringent bounds on $`(\delta _{12}^d)_{LL}`$ and $`(\delta _{12}^d)_{RR}`$ from $`\mathrm{\Delta }m_K`$ and $`\epsilon _K`$ prevent them to contribute significantly to $`\epsilon ^{}/\epsilon `$.
The LR squark mass matrix has the same flavor structure as the fermion Yukawa matrix and both, in fact, originate from the superpotential couplings. It may be appealing to invoke the presence of an underlying flavor symmetry restricting the form of the Yukawa matrices to explain their hierarchical forms. Then, the LR mass matrix is expected to have a very similar form as the Yukawa matrix. Indeed, we expect the components of the LR mass matrix to be roughly the SUSY breaking scale (e.g., the gravitino mass) times the corresponding component of the quark mass matrix. However, there is no reason for them to be simultaneously diagonalizable based on this general argument. To make an order of magnitude estimate, we take the down quark mass matrix for the first and second generations to be (following our assumption iii)),
$$Y^dv_1\left(\begin{array}{cc}m_d& m_sV_{us}\\ & m_s\end{array}\right),$$
(22)
where the (2,1) element is unknown due to our lack of knowledge on the mixings among right–handed quarks (if we neglect small terms $`m_dV_{cd}`$). Based on the general considerations on the LR mass matrix above, we expect
$$m_{}^{2}{}_{LR}{}^{(d)}m_{3/2}\left(\begin{array}{cc}am_d& bm_sV_{us}\\ & cm_s\end{array}\right),$$
(23)
where $`a`$, $`b`$, $`c`$ are constants of order unity. Unless $`a=b=c`$ exactly, $`M_d`$ and $`m_{LR}^{2,d}`$ are not simultaneously diagonalizable and we find
$`(\delta _{12}^d)_{LR}{\displaystyle \frac{m_{3/2}m_sV_{us}}{m_{\stackrel{~}{q}}^2}}=2\times 10^5\left({\displaystyle \frac{m_s(M_{Pl})}{50\mathrm{MeV}}}\right)\left({\displaystyle \frac{m_{3/2}}{m_{\stackrel{~}{q}}}}\right)\left({\displaystyle \frac{500\mathrm{GeV}}{m_{\stackrel{~}{q}}}}\right).`$ (24)
It turns out that, following the simplest implementation along the lines of the above described idea, the amount of flavor changing LR mass insertion in the s and d–squark propagator results to roughly saturate the bound from $`\epsilon ^{}/\epsilon `$ if a SUSY phase of order unity is present .
This line of work has received a great deal of attention in recent times, after the last experimental measurements of $`\epsilon ^{}/\epsilon `$ in KTeV and NA31 . The effects of non–universal $`A`$ terms in CP violation experiments were previously analyzed by Abel and Frere and after this new measurement discussed in many different works . In the following we show a complete realization of the above Masiero–Murayama (MM) mechanism from a Type I string–derived model recently presented by one of the authors .
### 4.1 Type I string model and $`\epsilon ^{}/\epsilon `$
In first place we explain our starting model, which is based on type I string models. Our purpose is to study explicitly CP violation effects in models with non–universal gaugino masses and $`A`$–terms. Type I models can realize such initial conditions. These models contain nine–branes and three types of five–branes ($`5_a`$, $`a=1,2,3`$). Here we assume that the gauge group $`SU(3)\times U(1)_Y`$ is on a 9–brane and the gauge group $`SU(2)`$ on the $`5_1`$–brane like in Ref. , in order to get non–universal gaugino masses between $`SU(3)`$ and $`SU(2)`$. We call these branes the $`SU(3)`$–brane and the $`SU(2)`$–brane, respectively.
Chiral matter fields correspond to open strings spanning between branes. Thus, they must be assigned accordingly to their quantum numbers. For example, the chiral field corresponding to the open string between the $`SU(3)`$ and $`SU(2)`$ branes has non–trivial representations under both $`SU(3)`$ and $`SU(2)`$, while the chiral field corresponding to the open string, which starts and ends on the $`SU(3)`$–brane, should be an $`SU(2)`$–singlet. There is only one type of the open string that spans between the 9 and 5–branes, that we denote as the $`C^{95_1}`$. However, there are three types of open strings which start and end on the 9–brane, that is, the $`C_i^9`$ sectors (i=1,2,3), corresponding to the $`i`$–th complex compact dimension among the three complex dimensions. If we assign the three families to the different $`C_i^9`$ sectors we obtain non–universality in the right–handed sector. Notice that, in this model, we can not derive non–universality for the squark doublets, i.e. the left–handed sector. In particular, we assign the $`C_1^9`$ sector to the third family and the $`C_3^9`$ and $`C_2^9`$, to the first and second families, respectively.
Under the above assignment of the gauge multiplets and the matter fields, soft SUSY breaking terms are obtained, following the formulae in Ref. . The gaugino masses are obtained
$`M_3`$ $`=`$ $`M_1=\sqrt{3}m_{3/2}\mathrm{sin}\theta e^{i\alpha _S},`$ (25)
$`M_2`$ $`=`$ $`\sqrt{3}m_{3/2}\mathrm{cos}\theta \mathrm{\Theta }_1e^{i\alpha _1}.`$ (26)
While the $`A`$–terms are obtained as
$$A_{C_1^9}=\sqrt{3}m_{3/2}\mathrm{sin}\theta e^{i\alpha _S}=M_3,$$
(27)
for the coupling including $`C_1^9`$, i.e. the third family,
$$A_{C_2^9}=\sqrt{3}m_{3/2}(\mathrm{sin}\theta e^{i\alpha _S}+\mathrm{cos}\theta (\mathrm{\Theta }_1e^{i\alpha _1}\mathrm{\Theta }_2e^{i\alpha _2})),$$
(28)
for the coupling including $`C_2^9`$, i.e. the second family and
$$A_{C_3^9}=\sqrt{3}m_{3/2}(\mathrm{sin}\theta e^{i\alpha _S}+\mathrm{cos}\theta (\mathrm{\Theta }_1e^{i\alpha _1}\mathrm{\Theta }_3e^{i\alpha _3})),$$
(29)
for the coupling including $`C_3^9`$, i.e. the first family. Here $`m_{3/2}`$ is the gravitino mass, $`\alpha _S`$ and $`\alpha _i`$ are the CP phases of the F–terms of the dilaton field $`S`$ and the three moduli fields $`T_i`$, and $`\theta `$ and $`\mathrm{\Theta }_i`$ are goldstino angles, and we have the constraint, $`\mathrm{\Theta }_i^2=1`$. Thus, if quark fields correspond to different $`C_i^9`$ sectors, we have non–universal A–terms. We obtain the following A–matrix for both of the up and down sectors,
$`A=\left(\begin{array}{ccc}A_{C_3^9}& A_{C_2^9}& A_{C_1^9}\\ A_{C_3^9}& A_{C_2^9}& A_{C_1^9}\\ A_{C_3^9}& A_{C_2^9}& A_{C_1^9}\end{array}\right).`$ (33)
The trilinear SUSY breaking matrix, $`(Y^A)_{ij}=(Y)_{ij}(A)_{ij}`$, itself is obtained
$$Y^A=\left(\begin{array}{ccc}& & \\ & Y_{ij}& \end{array}\right)\left(\begin{array}{ccc}A_{C_3^9}& 0& 0\\ 0& A_{C_2^9}& 0\\ 0& 0& A_{C_1^9}\end{array}\right),$$
(34)
in matrix notation.
In addition, soft scalar masses for quark doublets and the Higgs fields are obtained,
$$m_{C^{95_1}}^2=m_{3/2}^2(1\frac{3}{2}\mathrm{cos}^2\theta (1\mathrm{\Theta }_1^2)).$$
(35)
The soft scalar masses for quark singlets are obtained as
$$m_{C_i^9}^2=m_{3/2}^2(13\mathrm{cos}^2\theta \mathrm{\Theta }_i^2),$$
(36)
if it corresponds to the $`C_i^9`$ sector.
Now, below the string or SUSY breaking scale, this model is simply a MSSM with non–trivial soft–breaking terms from the point of view of flavor. Scalar mass matrices and tri–linear terms have completely new flavor structures, as opposed to the super–gravity inspired CMSSM or the SM, where the only connection between different generations is provided by the Yukawa matrices.
This model includes, in the quark sector, 7 different structures of flavor, $`M_Q^2`$, $`M_U^2`$, $`M_D^2`$, $`Y_d`$, $`Y_u`$, $`Y_d^A`$ and $`Y_u^A`$. From these matrices, $`M_Q^2`$, the squark doublet mass matrix, is proportional to the identity matrix, and hence trivial, then we are left with 6 non–trivial flavor matrices. Notice that we have always the freedom to diagonalize the hermitian squark mass matrices (as we have done in the previous section, Eqs.(35,36)) and fix some general form for the Yukawa and tri–linear matrices. In this case, these four matrices are completely observable, unlike in the SM or CMSSM case.
At this point, to specify completely the model, we need not only the soft–breaking terms but also the complete Yukawa textures. The only available experimental information is the Cabbibo–Kobayashi–Maskawa (CKM) mixing matrix and the quark masses. Here, we choose our Yukawa texture following two simple assumptions : i) the CKM mixing matrix originates from the down Yukawa couplings (as done in the MM case) and ii) our Yukawa matrices are hermitian . With these two assumptions we fix completely the Yukawa matrices,
$`\begin{array}{cc}Y_u=\frac{1}{v_2}\left(\begin{array}{ccc}m_u& 0& 0\\ 0& m_c& 0\\ 0& 0& m_t\end{array}\right)\hfill & \hfill Y_d=\frac{1}{v_1}K^{}.\left(\begin{array}{ccc}m_d& 0& 0\\ 0& m_s& 0\\ 0& 0& m_b\end{array}\right).K\end{array}`$ (44)
with $`v=v_1/(\mathrm{cos}\beta )=v_2/(\mathrm{sin}\beta )=\sqrt{2}M_W/g`$, and $`K`$ the CKM matrix. We take $`\mathrm{tan}\beta =v_2/v_1=2`$ in the following in all numerical examples. In this basis we can analyze the down tri–linear matrix,
$$Y_d^A(M_{St})=\frac{1}{v_1}K^{}.M_d.K.\left(\begin{array}{ccc}A_{C_3^9}& 0& 0\\ 0& A_{C_2^9}& 0\\ 0& 0& A_{C_1^9}\end{array}\right)$$
(45)
with $`M_d=diag.(m_d,m_s,m_b)`$.
Hence, together with the up tri–linear matrix we have our MSSM completely defined. The next step is simply to use the MSSM Renormalization Group Equations to obtain the whole spectrum and couplings at the low scale, $`M_W`$. The dominant effect in the tri–linear terms renormalization is due to the gluino mass which produces the well–known alignment among A–terms and gaugino phases. However, this renormalization is always proportional to the Yukawa couplings and not to the tri–linear terms, Eq.(3). This implies that, in the SCKM basis, the gluino effects will be diagonalized in excellent approximation, while due to the different flavor structure of the tri–linear terms large off–diagonal elements will remain with phases $`𝒪(1)`$ . To see this more explicitly, we can roughly approximate the RGE effects as,
$$Y_d^A(M_W)=c_{\stackrel{~}{g}}m_{\stackrel{~}{g}}Y_d+c_AY_d.\left(\begin{array}{ccc}A_{C_3^9}& 0& 0\\ 0& A_{C_2^9}& 0\\ 0& 0& A_{C_1^9}\end{array}\right)$$
(46)
with $`m_{\stackrel{~}{g}}`$ the gluino mass and $`c_{\stackrel{~}{g}}`$, $`c_A`$ coefficients order 1 (typically $`c_{\stackrel{~}{g}}5`$ and $`c_A1`$).
We go to the SCKM basis after diagonalizing all the Yukawa matrices (that is, $`K.Y_d.K^{}=M_d/v_1`$). In this basis, we obtain the tri–linear couplings as,
$$v_1Y_d^A(M_W)=(c_{\stackrel{~}{g}}m_{\stackrel{~}{g}}M_d+c_AM_d.K.\left(\begin{array}{ccc}A_{C_3^9}& 0& 0\\ 0& A_{C_2^9}& 0\\ 0& 0& A_{C_1^9}\end{array}\right).K^{})$$
(47)
From this equation we can get the $`L`$$`R`$ down squark mass matrix
$$m_{LR}^{2}{}_{}{}^{(d)}=v_1Y_{d}^{A}{}_{}{}^{}\mu e^{i\phi _\mu }\mathrm{tan}\beta M_d$$
(48)
And finally using unitarity of $`K`$ we obtain for the $`L`$$`R`$ Mass Insertions,
$`(\delta _{LR}^{(d)})_{ij}={\displaystyle \frac{1}{m_{\stackrel{~}{q}}^2}}m_i(\delta _{ij}(c_AA_{C_3^9}^{}+c_{\stackrel{~}{g}}m_{\stackrel{~}{g}}^{}\mu e^{i\phi _\mu }\mathrm{tan}\beta )+`$
$`K_{i2}K_{j2}^{}c_A(A_{C_2^9}^{}A_{C_3^9}^{})+K_{i3}K_{j3}^{}c_A(A_{C_1^9}^{}A_{C_3^9}^{}))`$ (49)
where $`m_{\stackrel{~}{q}}^2`$ is an average squark mass and $`m_i`$ the quark mass. The same rotation must be applied to the $`L`$$`L`$ and $`R`$$`R`$ squark mass matrices,
$`M_{LL}^{(d)}{}_{}{}^{2}(M_W)=K.M_Q^2(M_W).K^{}`$
$`M_{RR}^{(d)}{}_{}{}^{2}(M_W)=K.M_D^2(M_W).K^{}`$ (50)
However, the off–diagonal MI in these matrices are sufficiently small in this case thanks to the universal and dominant contribution from gluino to the squark mass matrices in the RGE.
At this point, with the explicit expressions for $`(\delta _{LR}^{(d)})_{ij}`$, we can study the gluino mediated contributions to EDMs and $`\epsilon ^{}/\epsilon `$. In this non–universal scenario, it is relatively easy to maintain the SUSY contributions to the EDM of the electron and the neutron below the experimental bounds while having large SUSY phases that contribute to $`\epsilon ^{}/\epsilon `$. This is due to the fact the EDM are mainly controled by flavor–diagonal MI, while gluino contributions to $`\epsilon ^{}/\epsilon `$ are controled by $`(\delta _{LR}^{(d)})_{12}`$ and $`(\delta _{LR}^{(d)})_{21}`$. Here, we can have a very small phase for $`(\delta _{LR}^{(d)})_{11}`$ and $`(\delta _{LR}^{(u)})_{11}`$ and phases $`𝒪(1)`$ for the off–diagonal elements without any fine–tuning . It is important to remember that the observable phase is always the relative phase between these mass insertions and the relevant gaugino mass involved. In Eq.(4.1) we can see that the diagonal elements tend to align with the gluino phase, hence to have a small EDM, it is enough to have the phases of the gauginos and the $`\mu `$ term approximately equal, $`\alpha _S=\alpha _1=\phi _\mu `$. However $`\alpha _2`$ and $`\alpha _3`$ can still contribute to off–diagonal elements. In figure 4 we show the allowed values for $`\alpha _S`$, $`\alpha _2`$ and $`\alpha _3`$ assuming $`\alpha _1=\phi _\mu =0`$. We impose the EDM, $`\epsilon _K`$ and $`bs\gamma `$ bounds separately for gluino and chargino contributions together with the usual bounds on SUSY masses. We can see that, similarly to the CMSSM situation, $`\phi _\mu `$ is constrained to be very close to the gluino and chargino phases (in the plot $`\alpha _S0,\pi `$), but $`\alpha _2`$ and $`\alpha _3`$ are completely unconstrained.
Finally, in figure 5, we show the effects of these phases in the $`(\delta _{LR}^{(d)})_{21}`$ MI as a function of the gravitino mass. All the points in this plot satisfy all CP–conserving constraints besides EDM and $`\epsilon _K`$ constraints. We must remember that a value of $`|\mathrm{Im}(\delta _{12}^d)_{LR}^2|10^5`$ gives a significant contribution to $`\epsilon ^{}/\epsilon `$. In this plot, we can see a large percentage of points above or close to $`1\times 10^5`$. Hence, we can conclude that, in the presence of new flavor structures in the SUSY soft–breaking terms, it is not difficult to obtain sizeable SUSY contributions to CP violation observables and specially to $`\epsilon ^{}/\epsilon `$ .<sup>2</sup><sup>2</sup>2With these $`L`$$`R`$ mass insertions alone, it is in general difficult to saturate $`\epsilon _K`$ . However, in some special situations, it is still possible to have large contributions
## 5 Conclusions and Outlook
Here we summarize the main points of these lectures:
* There exist strong theoretical and “observational” reasons to go beyond the SM.
* The gauge hierarchy and coupling unification problems favor the presence of low–energy SUSY (either in its minimal version, CMSSM, or more naturally, in some less constrained realization).
* Flavor and CP problems constrain low–energy SUSY, but, at the same time, provide new tools to search for SUSY indirectly.
* In all generality, we expect new CP violating phases in the SUSY sector. However, these new phases are not going to produce sizeable effects as long as the SUSY model we consider does not exhibit a new flavor structure in addition to the SM Yukawa matrices.
* In the presence of a new flavor structure in SUSY, we showed that large contributions to CP violating observables are indeed possible.
In summary, given the fact that LEP searches for SUSY particles are close to their conclusion and that for Tevatron it may be rather challenging to find a SUSY evidence, we consider CP violation a potentially precious ground for SUSY searches before the advent of the “SUSY machine”, LHC.
## Acknowledgments
We thank D. Demir, T. Kobayashi, S. Khalil and H. Murayama as co–authors of some recent works reported in these lectures. We are grateful to S. Bertolini, L. Silvestrini and F.J. Botella for enlightening conversations. A.M. thanks the organizers for the stimulating settling in which the school took place. The work of A.M. was partly supported by the TMR project “Beyond the Standard Model” contract number ERBFMRX CT96 0090; O.V. acknowledges financial support from a Marie Curie EC grant (TMR-ERBFMBI CT98 3087).
|
warning/0003/astro-ph0003293.html
|
ar5iv
|
text
|
# Constraints on the Stellar/Sub-stellar Mass Function in the Inner Orion Nebula Cluster
## 1 Introduction
The Orion Nebula is one of the most famous objects in the sky, and has been the target of innumerable astronomical observations at virtually all wavelengths over the past 100 years. Yet it is only within the past few years that we have begun to discover the extent of the young stellar population just emerging from the ambient molecular cloud, and to characterize its nature. Work by Herbig & Terndrup (1986), Prosser et al. (1994) and Hillenbrand (1997) has established that the mean age of stars projected within $``$2 pc of the massive Trapezium stars is $`<`$1 Myr. The mass distribution derived for $``$1000 ONC stars located on a theoretical HR diagram by Hillenbrand (1997) rises to $`0.2M_{}`$ and shows some evidence for flattening or turning over towards lower masses (see, however, the reinterpretation of these data in Figure 10 of the current paper using updated tracks/isochrones and updated transformations from observational to theoretical quantities). Our former study was complete to just above the hydrogren burning limit and did not constrain the mass function across the stellar/sub-stellar boundary into the brown dwarf regime. Existence of brown dwarfs in the ONC has been discussed previously by McCaughrean et al. (1995).
Star forming regions like the ONC provide one of the best environments for investigating the shape of the stellar mass function into the brown dwarf regime. Unlike the case in older clusters and associations, star-forming regions are essentially unperturbed by dynamical evolution that selectively remove the lowest mass objects. Further, contracting low-mass pre-main sequence stars and brown dwarfs are 2-3.5 orders of magnitude more luminous that their counterparts on the main sequence and hence can be readily detected, especially in the near-infrared. Star forming regions are also less affected by field star contamination compared to older clusters due to their small angular extent and their association with obscuring molecular material. The ONC cluster in particular affords several distinct advantages compared to any other young stellar cluster for measuring the initial mass function. First, since it is located at high galactic latitude toward the outer Galaxy, contamination from field stars is minimized. Further, the winds and ionization from the central OB stars have dispersed much of the surrounding gas and dust, drastically reducing the extinction to the cluster members. A high-column density of obscuring molecular material does remain intact behind the stellar cluster. Foremost, however, as the nearest massive star-forming region to the Sun and the most populous young cluster within at least 2 kpc, the ONC is the one region where one can assemble a statistically robust assessment of the mass distribution well into the brown dwarf regime.
In this contribution we investigate whether the distribution of stars in the K–(H-K) color-magnitude diagram for the ONC is consistent with a mass function that rises across the stellar/sub-stellar boundary and into the brown dwarf regime, or if the data demand that the mass function turns over. After describing the observations, image analysis, construction of the point source list, and extraction of photometry, we present a new approach for constraining the stellar/sub-stellar mass function. We consider that the location of a particular star in the K–(H-K) diagram depends on four parameters: stellar mass, stellar age, presence and properties of a circumstellar disk, and extinction. De-reddening the stars along a known reddening vector in the K-(H-K) diagram enables us to compute the probability that a star could be of a certain mass given the distributions in age and near-infrared excess that characteristize the ONC cluster. Summing of these individual mass probability distributions yields the mass function for the entire cluster. We believe that our technique produces the most rigorously derived constraint yet from photometry alone on the inner ONC initial mass function.
## 2 Observations
Images were obtained on 8, 9, 10 February, 1999 using NIRC (Matthews & Soifer 1994) mounted on the Keck I 10-m telescope. Data were taken in H-band (1.50-1.82$`\mu `$m) on the first night, K- (2.00-2.43$`\mu `$m) and H-bands on the second night, and Z-band (0.95-1.11$`\mu `$m) on the third night. The field of view of a NIRC frame is 38$`\mathrm{"}`$ x 38$`\mathrm{"}`$ at 0.152$`\mathrm{"}`$/pixel. For each filter, the observations consisted of a 15 x 15 grid of such frames aligned with the equatorial coordinate system to produce a 5.1 x 5.1 mosaic. Adjacent rows and columns in the grid were spaced by one-half of the array, so that any one pixel within the final mosaic was nominally observed on four different frames (modulo minor telescope drift). The integration times per frame were 0.5 sec with 50 co-adds (25 seconds total) at H- and K-bands, and 2 sec with 20 coadds (40 seconds total) at Z-band. Such short integrations per frame were necessitated by the large number density of relatively bright (K $`<`$ 12 mag) stars in the region whose saturation effects we wished to minimize.
The observing sequence for each declination row in the grid was to center the array on a previously chosen setup star, offset to the beginning of a row, scan in right ascension across the row, then offset to and five point dither on an off-field sky position, return to the setup star, and repeat. Local sky was measured after every row (at 10-12 minute spacings in time) from a location $``$15’ northwest of the ONC which we had determined via examination of the 2MASS Image Atlas to be free of nebulosity and relatively free of infrared point sources. We chose to include in our sky field a star of magnitude K$``$14 mag in order to monitor the atmospheric extinction with airmass locally (see, however, the results in Appendix A). In addition, we observed absolute photometric standards from Persson et al. (1998). Sky conditions on all three nights were photometric, with standard star solutions matching nominal NIRC zero points and nominal Mauna Kea extinction curves with airmass. Flat-field, bias, and linearity calibration data were also obtained.
## 3 Image Processing
The NIRC images were processed in IRAF first by determining the detector gain and readnoise from raw flat-field and bias frames, and establishing the linearity from a series of exposures taken with different integration times of the tertiary mirror cover. The latter tests indicated NIRC is linear between within 0.5% up to 90% of the full-well depth. Next, a median-filtered, normalized flat-field image was constructed for each filter from a series of 10 bias-subtracted dome flats. Bad pixel masks were made from the flat-fields by identifying all pixels more than 12% above or 15% below the mean value. Approximately 1.5% of the pixels were flagged as bad, with most of these located in a 1 pixel border around the edge of the array. Sky images were constructed for each of the 15 rows in our mosaic (in each filter) using the 5 dithered off-field sky frames, after dark-subtracting, median-filtering, and bad pixel exclusion.
Each of the on-field ONC data images was then sky- and dark-subtracted, and flat-divided. Next, on a frame-by-frame basis we identified and interpolated over intermittently appearing “warm” pixels. These pixels had values 20-50% above the local background and thus would affect our averaging of overlapping pixels in the final mosaic. The “warm” pixels numbered between $``$5-30 per frame with no discernible pattern and appeared as “warm” for a sequence of $``$3-10 frames before returning to normal values.
Next, we corrected for a feature of NIRC known as “bleeding” (Liu & Graham 1997). The readout electronics of NIRC are such that bright stars exhibit an exponentially decaying trail to the right which wraps around the right edge of the array and continues from the left, one row higher. A similar effect trailing downwards and wrapping around to the top of the array is associated with the brightest of stars. This “bleeding” behavior must be present at a visually imperceptible level for every star, and must be in the background counts as well. It can be thought of as part of the point-spread-function. But we need to model and correct for it since “bleeding” from the brighter stars can extend over other sources in the field and, furthermore, will adversely affect our frame-to-frame flux adjustments in the mosaicing process and our ability to model the point-spread-function using standard methods. We have used the empirical solution developed by Liu & Graham of subtracting from every pixel in the array, the exponentially decaying contribution of every pixel further back in the NIRC readout scheme. The coefficients of the correction are such that each pixel contributes 0.25% of its counts to the next pixel which is read out (physically located four pixels to the right on the array) with a pixel of, for example, 20,000 ADU contributing 51 ADU to a pixel 4 downstream and 1 ADU to a pixel 256 (1 row) downstream. Application of this “de-bleeding” method does not affect the photometry, or if it does, the level of any difference is well less than 1%.
Next, we corrected the frames for the effects of optical path distortion which amounts to $``$1 pixel from the center to the edge of the array. A. Ghez generously provided to us a subroutine for this step. Correction for image distortion improves the photometry by 0.02$`\pm `$0.01 mag. The final step in the raw frame processing was to correct each image by a multiplicative factor representing the flux adjustment from the observed airmass to zero airmass. At K-band this factor was $`10^{0.4\times 0.092\times AIRMASS}`$ while at H-band it was $`10^{0.4\times 0.065\times AIRMASS}`$.
Custom C programs were written to co-add the 225 images in each band into a mosaic in order to detect faint point sources and to improve the signal-to-noise of the photometry. The mosaics were constructed by determining the relative positional and sky offsets between the overlapping frames, and making the appropriate shifts in order to tile them together.
Positional offsets were established using stars identified in the overlap regions of neighboring frames. To reduce the random-walk errors in stitching the images together, the relative offsets of all 225 frames per band were solved simultaneously using a linear least-squares fit that minimized the position residuals for all stellar matches in all the overlap regions. Histograms of the resulting positional residuals at K band have a 1$`\sigma `$ standard deviation of 0.017<sup>′′</sup> in right ascension and 0.028<sup>′′</sup> in declination, and 0.022<sup>′′</sup> and 0.026<sup>′′</sup>, respectively, for H band.
As the images were placed in the mosaic, the sky background was adjusted by an additive constant to match the background in the surrounding frames. The relative intensity offsets were determined by fitting a gaussian to the difference in image intensity in the overlap region between neighboring frames. As with the positional offsets, the sky offsets for all frames were determined simultaneously using a linear least-squares fit to all the overlap regions. We were unable to obtain an acceptable solution over the entire mosaic for the sky offsets at Z-band. Thus our Z-band photometry is derived from the individual frames and is not as deep as it would be if derived from a co-added mosaic. The Z-band data and supplemental calibration information are presented in a separate paper. The H- and K-band mosaics are shown in Figure 1 along with an extinction map which is described in section 5.2. Figure 2 shows the spatial distribution of stars in our sample, whose identification and photometry we describe next. Table 1 contains the coordinates and HK photometry.
## 4 Mosaic Analysis
### 4.1 Identification of Point Sources
Point sources were identified on the K-band mosaic using DAOFIND in IRAF with a 7$`\sigma `$ threshold. The initial source list was hand-edited to remove nebular knots, multiple listings of bright stars, diffraction spikes, edge effects, etc. We then examined contour plots of each point source to look for extended or double-peaked structure, and added any newly found sources. The final source list consists of 778 stellar point sources over the 5.1x5.1 field.
### 4.2 Aperture and Point-Spread Function Fitting Photometry
Aperture photometry was derived using PHOT with a 6 pixel radius aperture and a sky annulus extending radially from 7-12 pixels; contribution from the sky was determined from the mode of these values. The small aperture and the close sky annulus were necessitated mainly by the spatially variable background from the Orion Nebula, and to a lesser extent by the high source density. Aperture corrections were needed in order to correct from the 6 pixel radius used to measure the data to the 20 pixel radius used to measure the standard stars. The size of the aperture correction is directly related to the size (e.g. full-width-half-maximum; FWHM) of the point source. In our image mosaics, however, the point-spread-function (PSF) undergoes large and non-systematic spatial variations due to random wandering in time of the seeing compared to the 0.15$`\mathrm{"}`$ platescale, and to a systematic gradient in airmass. We derived an empirical calibration between the aperture correction and the image FWHM, as follows.
First, recall that each stellar image in our final mosaic is synthesized from several (ideally four) separate observations of the star. Each of these observations may have a different PSF size. We verified that the PSFs of point sources in the final mosaic indeed have the mean value of the PSFs and the mean value of photometry through a fixed aperture size, characterizing the individual images from which they were created. In order to determine the appropriate aperture corrections for photometry from the mosaiced data, we therefore need only measure each of the PSFs in the final mosaic and apply a correlation between PSF and aperture corrections.
We measured the size of each stellar PSF using a variety of IRAF tasks – IMEXAMINE, RADPROF, and FITSPSF. From the 40 most isolated K $`<`$ 13 mag stars in the mosaic we found the tightest correlation (error in slope $`<`$ 0.01 mag) to be between the aperture correction and the “enclosed” gaussian fit FWHM of the IMEXAMINE task. In cases where this primary FWHM – aperture correction correlation could not be applied (e.g. failure of the gaussian fit to converge due to crowding and/or high nebulosity) we used secondary correlations. The FWHM values of $`>`$90% of the stars are between 3-5 pixels at K-band and between 4.5-6.5 pixels at H-band. The full range of the aperture corrections at K is from $``$ -0.10 mag to -0.60 mag with a mean value of -0.350 mag and at H is from $``$ -0.25 to -0.55 mag with a mean value of -0.325 mag. Variable aperture corrections are less of a problem at H compared to K since the overall size of the stellar images is larger which acts to decrease the percentage of PSF change as the seeing and airmass vary. Despite the complicated nature of our process for applying aperture corrections, we believe it is the proper one based on significant improvement in correlations between our NIRC photometry and previous photometry (described below). We did attempt to correlate aperture corrections with a measure of the difference in width between each stellar image and a PSF constructed from the data itself (the “sharpness” parameter of the PEAK task). But the correlation was too loose to be useful.
The point-spread-function was constructed using the PSF task and 30 stars distributed over the outer, less crowded, regions of our mosaics and having FWHM values distributed like the data as a whole. Accurate characterization of the noise in the mosaiced images was done so as to achieve the best possible fits of the point-spread-function to the stellar sources. In practice, this means adding a constant to the mosaics such that their standard deviation equals the sky counts plus the square of the effective read noise. A moffat function with $`\beta `$=1.5 gave the best residuals among the moffat, lorentzian, gaussian, and penny functions. As just discussed, the point-spread-function varies across our image due to seeing fluctuations and airmass changes. Because the variations are random and not smooth, we can not model them in any useful way and we are forced to fit the same point-spread-function to every star. The constructed point-spread-function was fit using the PEAK task with a fitting radius equivalent to twice the average FWHM. Both positional re-centering and sky re-calculation were permitted. Unfortunately, the resultant photometry was strongly influenced by which star was chosen as the first in constructing the PSF.
Comparisons between the aperture photometry and the point-spread-function fitting photometry are generally poor, due to the varying PSF. We have decided to use in our analysis the results from aperture photometry calibrated as described above. Our final photometry list was hand-edited to remove the measurements in cases of contamination from bright stars and/or overlapping apertures (44 stars) or nonlinearity (36 stars).
### 4.3 Integrity of Photometry
In this section we discuss the internal errors of our photometry and the comparison of our photometry to previous work.
Figure 3 shows the run of photometric error with magnitude and color. These internal errors are those produced by the PHOT routine in IRAF and simply reflect photon statistics of the source and sky determination in the fully processed images; they do not include other errors such as those in the zero point. At K-band, 75% of the stars in our source list have errors $`<`$0.02 mag, 90% have errors $`<`$0.05 mag, and 97% have errors $`<`$0.1 mag. At H-band, 69% of the stars in our source list have errors $`<`$0.02 mag, 85% have errors $`<`$0.05 mag, and 91% have errors $`<`$0.1 mag. For the PSF fitting photometry the percentages of stars with internal errors of the magnitudes given above were all down by 5 to 45 points, with the worst results at K-band. This is due to the generally poor fit of any single point-spread-function to all stellar images in the mosaic and part of our justification for rejecting the PSF photometry.
Figure 4 shows the comparison of our photometry to photometry from the 2MASS survey. Although the scatter is large ($``$0.2 mag) for the full sample of stars in common, it drops to $`<`$0.1 mag when only spatially well-isolated stars are considered (filled circles). Many of the largest deviations ($`>`$1 mag) are found in crowded regions where our photometry is always fainter than the 2MASS photometry, presumably due to our higher spatial resolution which permits better source separation and better sky determination. Nonetheless, there are also well-isolated stars with rather large differences in the photometry. Note, for example, the star at K<sub>NIRC</sub>=10.3, K<sub>NIRC</sub>-K<sub>2MASS</sub>=-0.8. This is a very bright, very well-isolated star whose photometry differs by almost 1 mag at K between our NIRC data, the 2MASS data, our previous photometry with SQIID/NICMASS (Hillenbrand et al. 1998), and that published by Ali & DePoy (1995) and Hyland, Allen, & Bailey (1993; as reported by Samuel 1993). Incidentally, this star (JW 737) is a 4.5 mag variable at I-band (W. Herbst, private communication). We have no choice but to interpret cases such as this as examples of real infrared photometric variability. Variability may well be the cause of much of the spread along the ordinate in Figure 4. Indeed, we have strong evidence from Figure 18 (discussed in the Appendix) that short-term variations of order 0.1 mag are present in at least some stars in the Orion A molecular cloud. Comparisons between our NIRC photometry and that presented in Hillenbrand et al. (1998), and between 2MASS and Hillenbrand et al. (1998), show somewhat larger scatter in the magnitudes but similar scatter in the colors.
In summary, we believe that our photometry for bright, isolated stars is within $`<`$0.1 mag of that determined by others. Much of this scatter may be attributed to photometric variability, although some role is probably played by the variable PSF which plagued our photometry extraction. We note that our careful attention to aperture corrections improved considerably the correlations between our NIRC results, our previously published data, and the 2MASS survey. Based on the comparisons in Figure 4, however, for the analysis presented below we conservatively assume minimum K magnitude errors of $`0.09/\sqrt{2}=0.06`$ mag and minimum H-K color errors of $`0.18/\sqrt{2}=0.13`$ for all stars in our sample. The $`\sqrt{2}`$ factor implies that we ascribe equal errors to our data and to 2MASS data even though the formal 2MASS errors are much larger than the formal NIRC errors for these relatively bright stars.
### 4.4 Artificial Star Experiments
We determined the completeness of our point source list and the completeness of our photometry using results from extensive experimentation with artificial stars. First, fake source lists were generated by randomly distributing 300 stars over the $``$2080 $`\times `$ 2080 pixels<sup>2</sup> in our mosaics, with the caveat that no star be placed within 3 pixels of a known star or within 30 pixels of the edge of the frame. Next, stars with the point-spread-function derived as described above were added to the image at these 300 locations using ADDSTAR. The enhanced image was then run through the DAOFIND, PHOT, and PEAK tasks in a manner identical to that used to extract photometry from the unaltered data. A separate artificial star test was conducted for every 0.25 magnitude interval in the range 14-18.5 mag. Our results for finding fake stellar point sources are somewhat different from our results for photometering fake stellar point sources.
The DAOFIND results are that for a detection threshold of 20$`\sigma `$, we can identify fake stellar point sources in our images at the 90% completeness level down to K = 16.8 mag and at the 20% completeness level down to K = 17 mag. For a detection threshold of 7$`\sigma `$, we can identify fake stellar point sources at the 90% completeness level down to K = 17.7 mag and at the 20% completeness level down to K = 18.2 mag.
The DAOPHOT results are presented in Tables 2 and 3. We assess both the internal errors, the formal uncertainties in the output magnitudes, and the external errors, the differences between the input and the recovered magnitudes, as a function of magnitude and also radial position in the cluster. The strong and variable nebular background in combination with extreme point source crowding makes these experiments somewhat more difficult to interpret than in the usual case. Nevertheless, we conclude based on internal error estimates (Table 2) that at the limit of our ability to detect 90% of the stellar point sources (K = 16.8 mag for the 20$`\sigma `$ threshold which produces 94% of the stars in our source list), we are able to do photometry accurate to 0.02 mag for 15% of them, 0.05 mag for 68% of them, and 0.1 mag for 85% of them. For brighter stars the internal error estimates are lower. For fainter stars, at the K = 17.7 mag limit 7$`\sigma `$ detection threshold, the photometry is accurate to 0.02 mag for 0% of the sources, to 0.05 mag for 19% of them, and 0.1 mag for 59% of them. However, based on the external errors (Table 3), we seem perfectly capable of recovering the input stellar magnitudes to within a few percent down to K$``$17.5 mag. The numbers listed in each column of this table are the median offsets between input and output magnitudes, and also the standard deviations. Of note is that any bias in our photometry, when it appears at a moderately important level for stars K $`<`$ 17.5 mag, is such that the MEDIAN offset is positive, meaning that we seem to measure the star as slightly fainter than it really is, probably through over-subtraction of background. By contrast, the MEAN offsets are positive (meaning that we measure the star as being too bright), but we note that the MEAN values are dominated by just a few data points with special problems such as proximity to very bright stars; hence we prefer to quote MEDIAN offsets.
In conclusion, based on artificial star experiments we adopt a conservative K = 17.5 mag as the completeness limit for source detection. The 75% completeness limits for photometry accurate to 10% are K = 17.3 mag and H = 17.4 mag while the 50% numbers for photometry accurate to 10% are K = 18.0 mag and H = 18.1 mag. These numbers are averaged over the 5.1 x 5.1 field and mask the systematic gradient with radius caused by variable source crowding and nebular strength. We estimate that we have determined the completeness to an accuracy of a few tenths of a magnitude only, due in part to the radial gradient and in part to the spatially variable PSF. Note that we found our ability to reproduce the magnitudes assigned to fake stars in the input stage varied significantly with the parameters given to the PHOT, PSF, and PEAK tasks in IRAF. The parameters producing the most accurate results in the artifical star experiments were those then used to extract the real source photometry as described previously.
### 4.5 Astrometry
Our astrometry for the NIRC mosaics is referenced to the 2MASS database, which in turn is referenced to the ACT catalog. The nominal 1$`\sigma `$ 2MASS position error for bright, isolated sources is $``$ 0.1$`\mathrm{"}`$ in each of right ascension and declination. An edited list of $``$230 stars in common between 2MASS and our Table 1 was used to derive the final astrometric solution, producing a platescale of 0.152$`\mathrm{"}`$/pixel and a total r.m.s. error in the positions of 0.10$`\mathrm{"}`$.
We note that the current astrometric system is shifted by about +1.5$`\mathrm{"}`$ in right ascension and -0.3$`\mathrm{"}`$ in declination compared to that presented by us previously. In an optical study of stars located within $``$20’ of the ONC core, Hillenbrand (1997) derived astrometry using the HST Guide Star Catalog (Lasker et al. 1988) which is known to suffer some inaccuracies in this region of the sky. We have found that our previous astrometric solution is offset to the west and to the south compared to other studies of the ONC (e.g. Jones & Walker 1988; McCaughrean & Stauffer 1994; Prosser et al. 1994; Ali & DePoy 1995; O’Dell & Wong 1996) but by various amounts as several of these studies have their own astrometric inaccuracies. The Jones & Walker positions and the McCaughrean & Stauffer positions are each internally consistent, although offset from one another. The McCaughrean & Stauffer astrometry matches our 2MASS-referenced astrometry. The HST positions of Prosser et al. and of O’Dell & Wong, however, are not internally consistent and suffer random excursions of about 1$`\mathrm{"}`$ which were propogated into the Hillenbrand (1997) database. Likewise, the coordinates of Ali & DePoy suffer large random errors, of order 1.5-2$`\mathrm{"}`$. Furthermore, approximately 1/2 of the sources supposedly located within our survey area as listed by Ali & DePoy simply do not exist in our higher resolution data; we do not discuss this catalog further. In summary, we believe that the positions quoted in Table 1 are both internally consistent and properly referenced to the ACT reference frame.
### 4.6 Properties of Final Source List
A fundamental result of this paper is a list of coordinates and available HK photometry for 778 stars in the inner 5.1 x 5.1 of the ONC. A total of 687 stars have measurements at both H and K, with 647 stars having errors $`<`$0.15 mag in both H and K. These data are presented in Table 1 along with cross-identifications to previously published optical and infrared source lists.
In the optical (Optical ID column), Jones & Walker (1988) numbers are listed with first priority, then Parenago (1954), Prosser et al. (1994), Hillenbrand et al (1997), and finally O’Dell & Wong (1996) if no other designation exists. Several of the previously cataloged optical stars appear not to be real sources based on their absence in our NIRC images. Another, although unlikely, possibility is that these are large-amplitude variables which have faded to K$`<`$17.5 mag. As listed in Hillenbrand (1997) these are 459 and 699 (Jones & Walker sources), 3071 and 3089 (Hillenbrand sources), and 9081 and 9326 (Prosser et al. sources). One other object, 3083, is the head of a teardrop-shaped “proplyd” identified from high-resolution HST images; we have left this spatially extended source in our photometry list along with several other “proplyds” which may not be point sources (e.g. OW-114-426, also extended in our images and clearly seen in silhouette in the Z-band data). In the infrared (Alternate Infrared ID column) McCaughrean & Stauffer numbers are given; we recover all stars from that survey except for a few close pairs – MS-86 is a close companion to P-1889, MS-65 is a close companion to P-1891 ($`\mathrm{\Theta }^1`$C Ori), and MS-75 and MS-77 appear as a single source in our images. Finally, $``$250 of the sources with K$`<`$14.5 mag in Table 1 are also listed as point sources by 2MASS.
In summary, of the 778 stars in Table 1, $`350`$ are previously known from optical studies conducted to varying survey depths while $`430`$ are more heavily “embedded” in molecular cloud and/or circumstellar material. Of the embedded sources, approximately $`125`$ were previously catalogued (McCaughrean & Stauffer over the inner 1.4’ x 1.4’; Downes et al. 1981 and Rieke, Low, and Kleinmann 1973 in the BN/KL region) and all those K $`<`$ 13.5 mag over the full area of the current NIRC survey were found in previous studies at lower spatial resolution and lower sensitivity (e.g. Hillenbrand et al. 1998). There are $``$175 sources newly catalogued here. Nearly all of these do appear in images available from McCaughrean or the NAOJ / Subaru Telescope first light press release.
In deriving the ONC mass function, we have edited down the list of 778 in Table 1 to remove those with one or more of the following features: 1) no photometry at either K or H (21 sources); 2) photometry at K or H only, but not both (32 sources); 3) photometry given only as lower or upper limits at either K or H (38 sources); 4) photometry with internal errors $`>`$0.5 mag at either K or H (14 sources); and 5) brighter counterparts in close pairs where we can not derive photometry for the fainter counterpart because of contamination by the brighter one (15 sources). This last criteria was imposed so that the edited source listed would not be biased in any way against fainter, presumably lower mass, objects. The number of stars remaining for our derivation of the ONC stellar/sub-stellar mass function is 658.
## 5 Basic Results
### 5.1 The K Histogram and K–(K-H) Color-magnitude Diagram
In Figure 5 we present the histogram of K magnitudes for all objects with measurable photometry over our 5.1 x 5.1 field. Consistent with previous near-infrared studies of the ONC (McCaughrean et al 1995; Ali & DePoy 1995; Lada et al. 1996) we find that the K magnitude histogram rises to a peak around K = 12-12.5 mag and then declines. The minor peak at K$``$14.5 mag is also seen in Figure 2 of McCaughrean et al. (1995). The hatched portion of the Figure shows the sample remaining after removing stars without suitable photometry at both H and K according to the criteria listed above. This sample is not substantially different from the unhatched distribution representing the full photometric database.
The K–(H-K) diagram for all stars with both H and K photometry from this study is presented in Figure 6. The 100 Myr isochrone (equivalent to the zero-age main sequence for masses M$`>`$0.35 M) and the 1 Myr pre-main sequence isochrone from D’Antona & Mazzitelli (1997, 1998), translated into this color-magnitude plane as described in section 6.1.3, are shown. Reddening vectors originating from the 1 Myr isochrone at masses of 2.5 M, 0.08 M, and 0.02 M are indicated. Considering only stellar photospheres for the moment, our data are sensitive to all objects (stars and brown dwarfs) with ages $``$1 Myr and masses M$`>`$0.02 M seen through values of extinction A$`{}_{V}{}^{}<`$ 10 mag. The observed colors of most of the objects are substantially redder than the expectations from pre-main sequence isochrones, a fact which can be attributed to a combination of extinction and excess near-infrared emission due to a circumstellar disk as discussed in section 6.2. Nevertheless, several tens of reddened objects located below the hydrogen burning limit at 0.08 M are present. Most of these are probable young brown dwarfs, although some may be field stars.
### 5.2 Field Star Contamination
Our images and the resulting K-magnitude histogram and K–(H-K) color-magnitude diagram contain both ONC cluster members and unrelated field stars. Since H- and K-band photometry alone can not distinguish between cluster members and nonmembers, we assessed contamination to the star counts from field stars using a modified version the Galactic star count model of Wainscoat et al. (1992). While the nominal Wainscoat et al. model includes a smooth Galactic extinction distribution, the line of sight toward the ONC contains a substantial and spatially variable extinction component from the Orion molecular cloud, as was shown in Figure 1c. This extinction map was generated from the C<sup>18</sup>O column density data of Goldsmith, Bergin, & Lis (1997) by assuming a C<sup>18</sup>O/H<sub>2</sub> abundance of $`1.7\times 10^7`$ (Frerking, Langer, & Wilson 1982) and that an H<sub>2</sub> column density of 10<sup>21</sup> cm<sup>-2</sup> corresponds to 1 magnitude of visual extinction (Bohlin, Savage, & Drake 1978). The visual extinction peaks along the western part of the inner ONC with a maximum value A<sub>V</sub> = 75 mag, and falls off sharply to the east to a minimum value of A<sub>V</sub> = 3 mag at the edge of our NIRC map. Obviously the contribution from background field stars to the observed star counts will vary substantially across the ONC, in inverse relation to this extinction distribution. The number and near-infrared magnitudes and colors of expected field stars were obtained by convolving this extinction map with the nominal Wainscoat et al. star count model, although the additional extinction from the Orion molecular cloud was added only to the background field star population.
The distribution of K magnitudes for the field stars was shown in Figure 5 (dashed curves), both before and after convolution with the molecular extinction map. With addition of the proper amount of extinction at the distance of the Orion cloud, the numbers predicted for foreground/background contamination in our data are reduced to 0.08 stars arcmin<sup>-2</sup> at K$`<`$13 mag (from 0.17 stars arcmin<sup>-2</sup>), 0.27 stars arcmin<sup>-2</sup> at K$`<`$15 mag (from 0.76 stars arcmin<sup>-2</sup>), and 1.54 stars arcmin<sup>-2</sup> at K$`<`$18 mag (from 3.41 stars arcmin<sup>-2</sup>). The total number of stars predicted to contaminate our NIRC photometry down to the K completeness limit (17.5 mag) is 34 (with 43 down to K=18 mag), representing a small but non-negligible 5% of our survey sample.
In Figure 7 we compare the ONC data (panel a) to the model field star population (panel b). For the data we have used the edit source list of 658 stars discussed above and for the field stars we have convolved the model with the photometric errors as a function of magnitude characterizing the NIRC photometry. A Hess diagram format was adopted, where individual points have been smoothed with an elliptical gaussian corresponding to the photometric uncertainties. This Figure highlights the large concentration of observed stars with K $``$ 12 and H-K $``$ 0.5, and affirms that the the density of stars in the K–(H-K) diagram is dominated by ONC cluster members at all but the faintest mangitudes. To derive the K–(H-K) distribution of stars actually associated with the ONC we subtracted panel (b) from panel (a), as shown in Figure 11a.
## 6 Analysis: The ONC Mass Spectrum Across the Hydrogen-Burning Limit
The goal of this study is to translate the information contained in the K–(H-K) diagram shown in discrete format in Figure 6 and in Hess format in Figure 7a, into information on the stellar/sub-stellar mass function. This is not a trivial transformation since the location of a young star in the K–(H-K) diagram depends on four parameters: stellar mass, stellar age, presence and properties (e.g. accretion rate) of a circumstellar disk, and extinction. A moderately bright, red object, for example, while usually thought of as a massive star seen through large extinction, can also be a much lower mass star with a large near-infrared excess and significantly lower extinction. Indeed, bright red objects can be reproduced by a number of combinations of stellar mass, stellar age, near-infrared excess, and foreground extinction. Faint, blue stars on the other hand, can come only from the lower masses, older ages, smaller near-infrared excesses, and lower extinctions. The distribution of data points across the K–(H-K) diagram is dictated by the interaction occuring for each star in the cluster of these four primary physical parameters.
Photometry alone can not be used to deconvolve the age, near-infrared excess, and extinction distributions to obtain uniquely the mass of the object. Such an effort would require spectroscopic observations for the entire cluster population, which currently do not exist. Therefore, a variety of techniques making various assumptions about these parameters have been developed in order to constrain the initial mass function. The most common of these approaches is to determine if the peak and the width of a K-band histogram are consistent with both an assumed initial mass function and a “reasonable” stellar age distribution (Zinnecker & McCaughrean 1991). One drawback to this approach is that the information inherent to multi-band photometry often is not used to constrain other two parameters: the extinction and the near-infrared excess properties of the individual stars. Other approaches attempt to use multi-band photometry to de-redden individual stars and to consider the effects of near-infrared excess in estimating individual stellar masses from a single mean mass-luminosity relationship (e.g. Meyer 1996). A drawback of this approach is that it assumes that the near-infrared excess at the shortest wavelength is negligible, and that the mean age of the entire cluster is a good approximation for each individual member of the cluster.
Here, we describe a new method for deriving the stellar mass function that acknowledges the existence a distribution of stellar ages, near-infrared excesses, and extinction values in star-forming regions. We use these distributions to determine the probability a star could be a certain mass based on its location in the K–(H-K) diagram, as opposed to estimating unique masses for individual stars. The advantages to this approach are that we use all the photometric information available, we make no a priori assumptions about the shape of the mass function, and we incorporate the inherent photometric uncertainties. Before describing our method for inverting the observed K–(H-K) diagram to derive the mass function, we first establish the stellar age and near-infrared excess distribitions appropriate for the ONC needed for our analysis.
### 6.1 Assumptions
#### 6.1.1 Stellar Ages
Hillenbrand (1997) used the D’Antona & Mazzitelli (1994) tracks to find a mean age for low-mass optically visible ONC stars of 0.8 Myr with an age spread of up to 2 Myr. We show in Figure 8 the age distribution derived using the more recent D’Antona & Mazzitelli (1997, 1998) calculations and updated spectral-type – temperature – color – bolometric correction relations, described below. This Figure includes only stars within the area of our NIRC survey. Hillenbrand (1997) discussed the presence of a radial gradient in the stellar ages where the mean age for stars in the inner ONC is slightly younger (by 0.25 dex or so) than the mean age for the ensemble ONC. This trend is also present using the updated theory and observational-to-theoretical transformations. Based on Figure 8, we adopt in what follows an age distribution which is uniform in log between 10<sup>5</sup> and 10<sup>6</sup> yr; we also consider a distribution which is uniform in log between 3x10<sup>4</sup> and 3x10<sup>6</sup> yr with little difference in the results.
#### 6.1.2 Near-Infrared Excess
We quantify the near-infrared excess using the H-K color excess, defined as $`\mathrm{\Delta }(HK)=(HK)_{observed}(HK)_{reddening}(HK)_{photosphere}`$. Spectroscopic and photometric data presented in Hillenbrand (1997) and Hillenbrand et al. (1998) allow us to compute this quantify for those optically visible stars within the area of our NIRC mosaic. $`(HK)_{observed}`$ is the tabulated color. $`(HK)_{reddening}=0.065A_V`$ where $`A_V`$ is derived from the spectral type and observed V-I color in comparison to expected V-I color and $`A_V=2.56E(VI)`$. $`(HK)_{photosphere}`$ comes from the relation between temperature and intrinsic H-K color described below. A histogram of the derived H-K excesses is shown in the top panel of Figure 9. We find that the observed near-infrared excesses can be well represented by a half-gaussian with a dispersion $`\sigma `$=0.4 mag, as shown by the solid line. In practice, we truncate the gaussian at $`\mathrm{\Delta }`$(H-K) = 1 mag, which is the maximum H-K excess observed in the inner ONC. Hillenbrand et al. (1998) discussed the presence of a radial gradient in ONC near-infrared excess values with the mean near-infrared excess (measured as $`\mathrm{\Delta }(IK)`$ instead of the $`\mathrm{\Delta }(HK)`$ used here) larger for stars in the inner ONC than for the ensemble ONC. We emphasize, therefore, that the H-K excess distribution presented in Figure 9 is known to be accurate for the inner ONC only, although we note that the distributions similarly calculated for young stellar populations in Taurus-Auriga, IC348, L1641, NGC2264, NGC2024, MonR2, and Chamaeleon using literature data (see description of samples and procedure in Hillenbrand & Meyer, 2000) are generally similar in form although slightly narrower in width.
In addition to the H-K excess, we must estimate the excess at K band alone in order to properly model the K–(H-K) diagram. Since the K magnitude excess is more difficult to compute accurately than the H-K color excess, we have used the K excesses tabulated for pre-main sequence stars in the Taurus molecular cloud by Strom et al. (1989) and the H-K excesses calculated as above using data from Kenyon & Hartmann (1995) to establish an empirical relation between these quantities. The bottom panel in Figure 9 shows the correlation between the K and H-K excess derived for stars in Taurus, which can be represented by a linear fit of $`\mathrm{\Delta }`$K = 1.785 $`\times `$ $`\mathrm{\Delta }`$(H-K) + 0.134 with a scatter of $`\pm `$ 0.25 mag. We assume that this relationship also holds for stars in the ONC.
#### 6.1.3 Translations from Theoretical to Observational Quantities
The final step before we can create models of the K–(H-K) diagram is conversion of theoretical pre-main sequence evolution into the observational plane. We use the theoretical description of luminosity and effective temperature evolution with mass according to D’Antona & Mazzitelli (1997,1998). These tracks are the only set available which cover the full range of masses sampled by our data, at the numerical resolution needed. We note, however, that the most recently circulated calculations of pre-main sequence evolution by various groups (D’Antona & Mazzitelli; Burrows et al.; Barraffe et al.;) do seem to be converging in the ranges where they overlap. Nevertheless, we must note that the details of our results likely are sensitive to the set of tracks/isochrones we have adopted.
We have transformed the D’Antona & Mazzitelli (1997,1998) calculations of L/L and T<sub>eff</sub>/K into K magnitude and H-K color using Chebyshev fits (Press et al. 1989) to bolometric correction, V-I color, I-K color, and H-K color vs effective temperature. For the mass range of interest in this paper we have taken the empirical data on bolometric corrections from Bessell (1991), Bessell & Brett (1988), and Tinney, Mould, & Reid (1993); on colors from Bessell & Brett (1988), Bessell (1991), Bessell (1995), Kirkpatrick & McCarthy (1994), and Leggett, Allard, and Hauschildt(1998); and on effective temperatures from Cohen & Kuhi (1979) – effectively Bessell (1991) – Wilking, Greene, & Meyer (1999), and Reid et al. (1999).
Note that these relationships are somewhat different than those used in Hillenbrand (1997, 1998). We have now shifted the temperature scale cooler and the bolometric corrections slightly smaller at the latest spectral types, in keeping with current consensus that was not well-established at the time of our earlier work. The combination of updated tracks and updated transformations between observations and theory have caused a shift in our interpretation of the optical data presented by Hillenbrand (1997). As we show in Figure 10, instead of a mass function that rises to a peak, flattens, and shows evidence for a turnover (bottom panel), the same data now appear to suggest a mass function for the greater ONC which continue to rise to the mass limit of our previous survey (top panel).
Finally, in the current analysis, we use the Cohen et al. (1981) reddening vector, A<sub>K</sub>=0.090 A<sub>V</sub> and A<sub>H</sub>=0.155 A<sub>V</sub>, and assume the Genzel et al. (1981) distance of 480$`\pm `$80 pc to the ONC (distance modulus = 8.41 mag).
### 6.2 A Model for the Distribution of Stars in the K–(H-K) Diagram
Given the above assumptions concerning the age and near-infrared excess distributions for the inner ONC, and the translation of theoretical tracks/isochrones into the K–(H-K) plane, we are now in a position to construct model K–(H-K) diagrams. In Figure 11 we illustrate the effects of various age and near-infrared excess distributions on the appearance of the K–(H-K) diagram. Panel (a) shows discrete isochrones from the calculations of D’Antona & Mazzitelli (1997, 1998) for ages of 10<sup>5</sup> and 10<sup>6</sup> year; panel (b) shows a sample of stars uniformly distributed in log-mass between 0.02-3.0 M and uniformly distributed in log age between 10<sup>5</sup> and 10<sup>6</sup> year; panel (c) shows the same mass and age distribution of (b) but now includes the near-infrared excess distribution parameterized in Figure 9. No extinction is included in these panels. Note that in the case of a uniform age distribution (panel b), the K–(H-K) diagram is not uniformly populated between the limiting isochrones. As originally shown by Zinnecker & McCaughrean (1991) in an analysis of K band histograms, the onset of deuterium burning occurs at different times for different masses, leading to distinctive peaks in the magnitude and color-magnitude distribution for pre-main sequence stars. These peaks become less distinctive when a near-infrared distribution is added (as shown in panel c) and even less distinctive when extinction is added (as shown next).
We incorporate elements of Figure 11 to show in Figure 12 (note the change in scale, now set to match the range of our data) two model K–(H-K) diagrams in comparison to our NIRC observations. Panel (a) shows the ONC data of Figure 7a with the field star model of Figure 7b subtracted. Panel (b) shows a stellar population distributed in mass according to the Miller-Scalo mass function and distributed in age between 10<sup>5</sup>-10<sup>6</sup> yr log-uniform, then also having the near-infrared excess distribution parameterized in Figure 9 and seen through extinction uniformly distributed between A<sub>V</sub>=0-5 mag. Panel (c) is the same as panel (b) except that the mass distribution is now a power-law function instead of a Miller-Scalo function. The salient difference between these two mass functions is that the Miller-Scalo function (N(log M) $``$ e$`^{C1(logMC2)^2}`$; C=1.14, C2=-0.88 as in Miller & Scalo, 1979) slowly declines across the hydrogen burning limit as N(log M) $``$ M<sup>0.37</sup> if forced to a power-law, while the straight power-law function (N(log M) $``$ M<sup>-0.35</sup>) slowly rises. In creating Figure 12 we have not attempted to reproduce the observations; we wish merely to illustrate the combined effects in the K–(H-K) diagram of different assumptions about the mass, age, near-infrared excess, and extinction distributions. Note in particular that there are many stars observed (panel a) through higher values of A<sub>V</sub> than we have considered in the models (panels b and c). Nevertheless, if we accept that our assumptions about the age and near-infrared excess distributions (as derived from optically visible stars in exactly this region) are approximately correct, then we must conclude that a declining mass distribution such as the Miller-Scalo function is a much better match to the data than a rising power-law (or even a flat) function. We quantify these impressions in the following section.
### 6.3 Implementation and Tests
#### 6.3.1 Calculating Mass Probability Distributions
How can the effects of stellar age, near-infrared excess, and extinction be disentangled to derive the mass function? As already discussed, the main difficulty is that more than one stellar mass can contribute power to any particular location in the K–(H-K) diagram through various combinations of these variables. Fortunately, however, the range of stellar masses that a given H-K,K data point could represent is constrained by the stellar age and near-infrared excess distributions, which for the inner ONC we are able to measure (Figures 8 and 9), and by the slope of the reddening vector.
In practice we calculate the stellar mass function as follows. We take as a starting point the observed K–(H-K) grid shown in Figure 12a. Recall that each star has been smeared out in this diagram by an elliptical gaussian corresponding to its photometric error; increasing the error even by a factor of 3 in each direction does not change the form of the distribution. We project every 0.01 mag wide pixel populated by data back along the reddening vector to establish which of the other pixels are crossed, and hence which combinations of unreddened K magnitudes and H-K colors the star or star-plus-disk system could have. Using a model K–(H-K) diagram, we keep track of the probability that a star of given mass can occupy that H-K,K combination given the assumed stellar age and near-infrared excess distributions. We sum the probabilities for all of the possible H-K,K combinations along the reddening vector, and then normalize to unity the integerated probability over all masses 0.02-3.0 M; i.e., the star must have have some mass within the considered range. By weighting the mass distribution derived for each pixel in Figure 12a by the relative density of observational data it represents, and summing the probability distribution for all pixels, we produce the cluster mass function. In a similar manner, we calculate probability distributions in A<sub>V</sub> for each pixel, which we also density weight and sum to produce the cluster extinction distribution.
Examples of individual stellar mass probability distributions obtained by de-reddening a star in the K–(H-K) diagram are shown in Figure 13 for a representative set of K magnitudes and H-K colors. A K=15 mag relatively blue (H-K=0.5 mag) star is permitted to have a mass anywhere in the range $``$0.02-0.04 M with a most likely value just above 0.02 M, while a K=15 mag much redder (H-K=3.0 mag) star has a broader range of permitted masses, $``$0.03-0.6 M with a most likely value $``$0.2 M. Note the tails upward at the lower and upper mass extrema in the panels for K=16, H-K=0.5 and K=9, H-K=3.0, respectively. These are caused by our imposition of integrated probability equal to unity over the mass range contained in the theoretical grid; in reality such stars have some probability of coming from smaller and higher masses (respectively) than the 0.02-3.0 M range considered here. The requirement of an integrated mass probability of unity means that the outer few bins of our resultant mass distribution may be unreliable.
Examples of individual extinction probability distributions are not shown since extinction is essentially the independent variable in our technique. Because we project each “star” back along the reddening vector and keep track of the stellar mass, stellar age, and near-infrared excess combinations that can conspire to produce that color-magnitude location, any given star can have any value of extinction ranging from a minimum of zero (in general, although it is not always true that there is a zero-extinction solution) to a maximum set by the case of de-reddening to the oldest considered age (i.e. bluest possible original location in the K–(H-K) diagram) and having no near-infrared excess. The result is that an extinction distribution which is uniform will be recovered using our methodology as an extinction distribution that has an extended tail induced by a combination of the age range and the near-infrared excess range considered in the de-reddening process.
#### 6.3.2 Tests of Methodology
Before applying our newly developed methodology for deriving the stellar/sub-stellar mass function, we wish to test how accurately this method can recover a known mass function. For these tests we generated cluster models with various stellar mass, stellar age, near-infrared excess, and extinction distributions, and then attempted to recover the underlying mass function using the procedure described above. Figures 14 and 15 illustrate a sampling of the results. In general, we are fairly confident in our ability to recover the general form of the input mass distribution for masses 0.03 $`<`$ M/M $`<`$ 1. Outside of these mass limits we suffer problems due to “edge effects” given the 0.02-3.0 M range of the theoretical models we employ, and also due to saturation in our data at K $`<`$ 9 mag. In particular we note that in all cases we easily distinguish between mass functions that slowly fall across the hydrogen burning limit into the brown dwarf regime, as N(log M) $``$ M<sup>0.37</sup>, and those that slowly rise, as N(log M) $``$ M<sup>-0.35</sup>.
In Figure 14, we present test results where the stellar age and near-infrared excess distribution assumed in extracting the mass function from the K–(H-K) diagram is the same as the input cluster model. Thus these tests probe the success of the method when the cluster properties are accurately known a priori. Each of these models contain a log-uniform age distribution between 10<sup>5</sup>-10<sup>6</sup> yr. The left panels are for models where the input mass function is Miller-Scalo while the right panels are for a power law of form N(log M) $``$ M<sup>-0.35</sup>. The top panels are models with no extinction and no near-infrared excess; the middle panels include uniform extinction between A<sub>V</sub>=0-5 mag; and the bottom panels include both uniform extinction and the near-infrared excess distribution parameterized in Figure 9. Note that the two bottom panels correspond to the cases shown in the K–(H-K) diagrams of Figure 12bc. Looking at the difference between the top and middle panels, addition of extinction to a model (surprisingly) helps our method to recover the input mass function. Looking at the difference between the middle and bottom panels, addition of near-infrared excess hurts slightly but only at the tails in mass. These test results illustrate that our method can never perfectly recover the input mass function as long as there is a spread of ages or of near-infrared excesses – even if these distributions are known. H- and K-band photometry alone can not uniquely determine the mass, age, near-infrared excess, and extinction which go into producing the observed color-magnitude location, and hence our method considers all possible combinations of these parameters. The result is imperfect; nevertheless, it seems clear from Figure 14 that we do reasonably well in recovering the general shape of the input mass function.
In Figure 15, we present test results where we deliberately choose an incorrect cluster age and/or near-infrared distribution to de-redden the K–(H-K) diagram. These tests probe how robust our procedure is in recovering the mass function when faced with uncertainties in characterizing the actual cluster properties. In each of these tests, the input cluster contains the same age distribution and near-infrared excess distribution that we assumed for the ONC. In addition, we added uniform extinction between A<sub>V</sub>=0-5 mag. The left panels are tests results for the Miller-Scalo mass function, and the right panels for a power law mass function. The top three panels (left and right) show the effects of incorrect assumptions about the cluster age in the de-reddening process. Single-age assumptions give the worst results with two effects occurring. The first is a general shift of the recovered mass function towards higher masses as the age assumption is moved to older ages, due simply to the decrease in luminosity with age for a given mass star. The second effect is a “kinking” in the recovered mass function which is caused by considerable flattening of single isochrones in the 0.3-0.1 M range compared to higher and lower masses in the K–(H-K) diagram. When a range of ages is assumed in the de-reddening, instead of just a single age, this effect is smeared out. Note that there is little difference between the panels which assume the correct age distribution, log-uniform between 10<sup>5</sup> and 10<sup>6</sup> yr, and the panels which assume a somewhat broader age distribution, log-uniform between 3x10<sup>4</sup> and 3x10<sup>6</sup> yr. The bottom panels (left and right) show the effects of an incorrect assumption about the near-infrared excess. When no infrared excess is allowed for in the de-reddening process too much power is given in the recovered mass function to higher masses relative to lower masses.
To summarize our test results, we find that we can recover the input mass function with some reasonableness in all cases where we know the correct stellar age and near-infrared excess distributions, and in most cases where we assume somewhat (but not grossly) incorrect representations for these distributions. The worst results are obtained when the cluster consists of a uniform age distribution, but a single age is assumed to derive the mass function. Since the spectroscopic data for the ONC indicate an age spread, this is in fact the least applicable case for this study. Based upon these test results, we expect that our procedure to recover the input mass function performs well over the mass range 0.03$`<`$M/M$`<`$1, and that it can distinguish between mass functions that slowly fall across the hydrogen burning limit, as N(log M) $``$ M<sup>0.37</sup>, and those that slowly rise, as N(log M) $``$ M<sup>-0.35</sup>.
### 6.4 Results on the ONC Stellar/Sub-stellar Mass Function
Using the procedure we have described and tested above, we present in Figure 16 the ONC mass function resulting from our best determinations of the appropriate stellar age (Figure 8) and near-infrared excess (Figure 9) distributions. Of the 658 stars with suitably good H and K photometry going in to this analysis, we recover 598 when we integrate over this mass function. The loss of $``$9% is due to color-magnitude diagram locations (spread by photometric errors; see Figure 7a) with no solution inside the bounds of the mass grid considered in this analysis given the assumed age and near-infrared excess distributions. Since bright massive stars can be detected through larger values of extinction than faint brown dwarfs, we also plot the mass function for only those objects meeting certain extinction criteria: first, only those with A$`{}_{V}{}^{}<10`$ mag, the highest extinction level to which 0.02 M objects can be detected given the sensitivity limits of our survey, and second, only those with A$`{}_{V}{}^{}<2.5`$ mag, the extinction limit to which 0.1 M objects could be detected in the optical spectroscopic survey by Hillenbrand (1997), to which we compare our infrared photometric results below. Of the total of 598 sources in the mass function of Figure 16 (open histogram), 67% have A$`{}_{V}{}^{}<10`$ mag while 28% have A$`{}_{V}{}^{}<2.5`$ mag.
As shown in Figure 16, the stellar/sub-stellar mass function in the ONC peaks near $``$ 0.15 M and is clearly falling across the hydrogen burning limit into the brown dwarf regime – regardless of the adopted extinction limit, which affects the shape of the mass function only at the higher masses. We have investigated the robustness of Figure 16 for different plausible age ranges (e.g. log-uniform between 3x10<sup>4</sup> and 3x10<sup>6</sup> yr instead of between 1x10<sup>5</sup> and 1x10<sup>6</sup> yr), with and without a near-infrared excess distribution, and also with and without subtraction of field stars. The same basic conclusion is found. A power law fit to the declining inner ONC mass function for A$`{}_{V}{}^{}<`$ 10 mag between 0.03 M and 0.2 M has a slope of 0.57 $`\pm `$ 0.05 (in logarithmic units), where the uncertainties reflect only the residuals of the least squares fit to the data. Our best determination of the inner ONC mass function is inconsistent at the $`>10\sigma `$ level with a mass function that is flat or rising across the hydrogen burning limit.
According to the tests of our methodology (section 6.3.2), there are two ways to add power at low masses relative to higher masses and thus produce a less steeply declining or even flat slope across the hydrogen-burning limit: by making the cluster age much younger than we have assumed, and/or by making the near-infrared excesses much larger than we have assumed. We find neither of these options probable given the characteristics of the optically visible stars in the region, and hence we conclude that the inner ONC mass function is indeed declining. We have shown the accuracy to which our methodology recovers a known input mass function in Figure 14. Based on fits over the same 0.03-0.2 M mass range we consider for the data, we conclude that our method recovers the correct slope of the input mass function to within $`<`$0.05. Combining this methodology error with the r.m.s. fitting error of $`\pm `$0.05 discussed above, we estimate the total error on the slope derived here for the ONC mass function across the hydrogen burning limit at $`<`$0.1. We offer the following two additional cautions to any interpreters of our results.
First, we emphasize that the detailed shape of the mass function derived from data is still subject to dependence on theoretical tracks and isochrones (D’Antona & Mazzitelli; 1997, 1998 in this case), and on the calibrations used in converting between effective temperature / luminosity and K–(H-K) color/magnitude (discussed in section 6.1.3).
Second, we emphasize that our derived mass function is valid only for the inner 0.71 pc $`\times `$ 0.71 pc of the ONC cluster, which extends at least 8-10 pc in length and 3-5 pc in width. Our conclusions may not apply to the ONC as a whole where some evidence for general mass segregation has been found by Hillenbrand (1997) and Hillenbrand & Hartmann (1998). In Figure 17 we compare the mass function derived here for the inner cluster using near-infrared photometry to that derived previously by us using optical photometry and spectroscopy. The histogram is the mass function of Figure 16 with an extinction limit of A$`{}_{V}{}^{}<`$ 2.5 mag, for consistency with the effective extinction limit of the optical data. Solid symbols represent the full dataset from Hillenbrand (1997) while open symbols represent only that portion of the data which are spatially coincident with the near-infrared photometry used to derive the histogram (i.e. the inner 0.35 pc or so). The same A$`{}_{V}{}^{}<`$ 2.5 mag imposed on the infrared data has also been imposed on the optical data. As noted in reference to Figure 10, the updated pre-main sequence tracks and the updated transformations between observational and theoretical quantities adopted in this paper have caused a shift in our interpretation of the data presented by Hillenbrand (1997). The large-scale ONC mass function (solid symbols) now appears to be rising to the limit of that survey. The inner ONC mass function (open symbols), however, appears to flatten below $``$0.3 M. This flattening is confirmed by the near-infrared photometric analysis presented here, and in fact is the beginning of a turnover in the mass function above the hydrogen burning limit and extending down to at least 30 M<sub>Jupiter</sub>.
## 7 Discussion
Our analysis of the mass distribution in the inner ONC agrees with that of McCaughrean et al. (1995) in that there is “a substantial but not dominant population of young hot brown dwarfs” in the inner ONC. Although we do find $``$80 objects with masses in the range 0.02-0.08 M, the overall distribution of masses is inconsistent with a mass function that rises across the stellar/sub-stellar boundary. Instead, we find that the most likely form of the mass function in the inner ONC is one that peaks around 0.15 M and then declines across the hydrogen-burning limit to the mass limit of our survey, 0.02 M. The best-fit power-law for the decline, N(log M) $``$ M<sup>0.57</sup>, is steeper than that predicted by the log-normal representation of the Miller-Scalo initial mass function, N(log M) $``$ M<sup>0.37</sup> if forced to a power-law (see Figure 16).
How do our results compare to other determinations of the sub-stellar mass function? Thusfar there have been few actual measurements of the sub-stellar mass function which are not either lower limits or dominated by incompleteness corrections or small-number statistics. We can compare our results for the inner ONC only to those in the Pleiades (Bouvier et al. 1998; Festin 1998) and the solar neighborhood (Reid et al. 1999), and we find some differences. Converting the logarithmic units used thusfar in this paper (N(log M) $``$ M<sup>Γ</sup>) to the linear units adopted by others (N(M) $``$ M<sup>α</sup>) we find a mass function slope across the hydrogen-burning limit of $`\alpha =\mathrm{\Gamma }1`$ = -0.43. In the Pleiades, Bouvier et al. find $`\alpha `$ = -0.6 while Festin finds $`\alpha `$ in the range 0 to -1.0. In the solar neighborhood, Reid et al. find $`\alpha `$ in the range -1.0 to -2.0 with some preference for the former value. The methods used by these different authors for arriving at the slope of the mass function are very different, thus rendering somewhat difficult any interpretation of the comparison. Furthermore, it is not clear that the mass function in the center of a dense and violent star-forming environment should bear any similarity to the mass function in a lower-density, quiescent older cluster, or that either of these cluster mass functions should look anything like the well-mixed, much older local field star population. Nevertheless, if comparisons can be made, the inner ONC seems to have a shallower slope than that found in any other region where measurements have been made; recall as well that the inner ONC mass function appears shallower than the overall ONC mass function (see Figure 17).
## 8 Conclusions
We have introduced a new method for constraining the stellar/sub-stellar mass distribution for optically invisible stars in a star-forming region. A comparative review of the various techniques already in use for measuring mass functions in star-forming regions is presented by Meyer et al. (2000). These techniques range from studies of observed K-magnitude histograms (e.g. Muench et al. 2000), to discrete de-reddening of infrared color-magnitude diagrams (e.g. Comeron, Rieke, & Rieke 1996), to the assembly of photometric and spectroscopic data from which HR diagrams are created (e.g. Luhman & Rieke 1998). Our method is a variation on and an improvement to the discrete de-reddening of color-magnitude diagrams since we fully account for distributions in the relevant parameters instead of assuming a mean value for them. However, our method is not as good as a complete photometric - plus - spectroscopic survey since we produce only a mass probability distribution for each star, not a uniquely determined mass. Nonetheless, we believe that the statistical nature of our method does provide the most rigorously established constraint to date from photometry alone on the stellar mass function in a star-forming region.
We have used information from previous studies of optically visible stars in the ONC to derive plausible functional forms for the stellar age and the circumstellar near-infrared excess distributions in the innermost regions studied here. We assume that these distributions apply equally well to the optically invisible population. We find a mass function for the inner 0.71 pc x 0.71 pc of the ONC which rises to a peak around 0.15 M and then declines across the stellar/sub-stellar boundary as N(log M) $``$ M<sup>Γ</sup> with slope $`\mathrm{\Gamma }=0.57`$. This measurement is of the primary star/sub-star mass function only, and should be adjusted by the (currently unknown) companion mass function in order to derive the “single star mass function,” if desired.
We find strong evidence that the shape of the mass function for this inner ONC region is different from that characterizing the ONC as a whole, in the sense that the flattening and turning over of the mass function occurs at higher mass in the inner region than in the overall ONC. In fact, the shape of mass function for the overall ONC is currently unconstrained across the stellar/sub-stellar boundary, and appears now based on the most recent theoretical tracks and conversions between the theory and observables used in this paper, to continue to rise to at least 0.12 M.
###### Acknowledgements.
We thank Mike Liu and James Graham for sharing their method and source code for “de-bleeding” of NIRC images. We thank Andrea Ghez for providing her image distortion coefficients. We thank Keith Matthews for consultation regarding these and other NIRC features. Shri Kulkarni and Ben Oppenheimer suggested a mutually advantageous exchange of telescope time which enabled us to obtain the Z-band observations. Ted Bergin kindly provided his C<sup>18</sup>O data to us and Richard Wainscoat gave us a base code for his star count model. This publication makes use of data products from the Two Micron All Sky Survey, which is a joint project of the University of Massachusetts and the Infrared Processing and Analysis Center, funded by the National Aeronautics and Space Administration and the National Science Foundation. LAH acknowledges support from NASA Origins of Solar Systems grant #NAG5-7501. JMC acknowledges support from NASA Long Term Space Astrophysics grant #NAG5-8217, and from the Owens Valley Radio Observatory which is operated by the California Institute of Technology through NSF grant #96-13717.
## Appendix A The Infrared Variable 2MASSJ053448-050900 = AD95-1961
In this appendix we report on short timescale variability at infrared wavelengths of a star located approximately 15’ north-west of the ONC. Our observing procedure in constructing our 5.1’ x 5.1’ mosaic with NIRC was to scan across a row at constant declination, then move off to a sky position and obtain five measurements of the sky which were then averaged and subtracted from each frame in our ONC mosaic. We intentionally chose a sky field which included a relatively bright (K$``$ 14.0 mag) star in order to monitor the atmospheric extinction as part of our normal data acquisition. However, while our set of absolute standards from Persson et al. (1998) matched nominal NIRC zero points and nominal Mauna Kea extinction curves with airmass, our local standard exhibited significant flux variations (Figure 18). The amplitude of the variations is about 0.1 mag, and the timescale is less than the separation of our observations, about 10-12 minutes.
2MASSJ053448-050900 is also catalogued as AD95-1961 (Ali & DePoy 1995). This star has infrared fluxes of K=14.03 mag, H=14.43 mag, J=15.46 mag from the 2MASS survey and optical fluxes of I$``$17.5 mag, V$``$21.1 mag from our own unpublished CCD observations. These colors are consistent with those of low-mass ONC proper motion members.
The short-term photometric behavior of this relatively isolated and otherwise nondescript star located in the outer regions of the ONC may in fact be a general feature of all young stellar objects. Infrared monitoring studies of young clusters are needed in order to quantify the nature and constrain the causes of this variability.
\[THIS TABLE CAN BE FOUND IN ascii FORMAT AT http://astro.caltech.edu/$``$lah/papers.html \]
|
warning/0003/astro-ph0003200.html
|
ar5iv
|
text
|
# Using HI to probe large scale structures at 𝑧∼3 .
## 1 Introduction
The problem of determining the distribution of matter on large scales in the universe and understanding the large scale structure (LSS) formation is of prime importance in modern cosmology. Observing the angular positions and redshifts of galaxies has been the most straightforward method of mapping the LSS in the present day universe (Peebles 1993, Peebles 1980), but the interpretation of these observations is complicated by the fact that the relation between the distribution of galaxies and the distribution of underlying matter is not fully understood (Bardeen et al. 1986). Other direct methods use galaxy clusters or super-clusters as tracers to map the large scale structures. A large variety of techniques have been developed and applied to quantify the distribution of galaxies, and amongst these the correlation functions (two-, three-point correlation functions, etc) and their Fourier counterparts (power spectrum, bi-spectrum, etc.) have been most popular (Peebles 1980). Much of the work comparing the observations with different theories has been based on these.
An alternative approach is to use observations of fluctuations in the diffuse extragalactic background radiation at different wavelengths to probe the large scale structure. Here, the observations of anisotropies in the cosmic microwave background radiation (CMBR) have been most important. These observations probe the clustering of matter at the last scattering surface (e.g. Bond 1996), and combined with the information from the galaxy surveys, they have been successful in narrowing down the allowed class of theoretical models of LSS formation (e.g., Peacock 1999 and reference therein). The study of diffuse background at other wavelengths has been used to infer the clustering properties of matter at more recent epochs. Shectman (1974) observed fluctuations in the optical background; the results of this observation are consistent with predictions from galaxy counts (Peebles 1980). This method has also been applied to far-infrared background (Kashlinsky et al. 1997) and recently the first detection of fluctuations in this background has been reported (Kashlinsky et al. 1999). Gorjian, Wright and Chary (1999) have recently reported a tentative detection of a near-infrared background. There are similar predictions for fluctuations in the X-ray background (Barcons et al. 2000, Barcons et al. 1998).
In this paper we investigate the possibility of using the extragalactic background radiation at low radio frequencies (meter wave) to probe the LSS. This is largely motivated by the fact that the Giant Meter-wave Radio Telescope (GMRT; Swarup et al. 1991) which is designed to observe in several frequency bands in the interval $`150\mathrm{MHz}`$ to $`1420\mathrm{MHz}`$ has recently started functioning. This frequency range corresponds to the $`1420\mathrm{M}\mathrm{H}\mathrm{z}`$ emission from HI in the redshift range $`0z8.5`$. Absorption studies along sight lines to quasars indicate that most of the HI in the redshift range $`0z3`$ is in damped Lyman-$`\alpha `$ (DLA) clouds and the density of HI in DLAs has been determined reasonably well from absorption studies (Lanzetta, Wolfe, & Turnshek 1995). Although the HI emission from individual DLAs at high $`z`$ is too faint to detect using presently available telescopes, the redshifted HI emission from unresolved DLAs will contribute to the background radiation at low frequencies. In this paper we investigate how the brightness temperature of this radiation is related to the density and peculiar velocity of the HI. We consider the possibility of detecting correlations in the fluctuations in this component of the background radiation and using this to probe LSS at high $`z`$.
The possibility of observing the redshifted 21 cm emission from HI at high $`z`$ has been discussed earlier by many authors in a large variety of contexts. One of the first papers (Sunyaev & Zel’dovich, 1975) considers the possibility of meter-wave observations to detect protogalaxies and protoclusters at $`z10`$. There have been several attempts to detect the HI in proto-clusters and proto-super clusters (see Subrahmanyan & Anantharamaiah 1990 and reference therein).
More recently Subramanian and Padmanabhan (1993) have calculated the abundance of protoclusters as a function of the redshifted HI flux density for various redshifts for both the CDM and HDM models. In a later paper Kumar, Padmanabhan and Subramanian (1995) have calculated the line profile of the HI emission from a spherically symmetric protocluster. Bagla, Nath and Padmanabhan (1997) and Bagla (1998) have used a combination of N-body simulations and a model for the behaviour of the baryons to calculate the abundance and the expected flux from the HI in structures like protoclusters at high redshifts. A major uncertainty in all of these works is in the assumptions about the HI content of the universe at high redshifts. The main focus of all of these works has been on individual peaks of the density fluctuations (protoclusters) which will manifest themselves as detectable features in low-frequency radio maps. Subramanian and Padmanabhan (1993) have also studied the possibility of detecting the excess variance in radio images due to the density fluctuations in the HI at high $`z`$.
Katz, Weinberg and Hernquist (1998) have used smoothed particle hydrodynamic simulations to study the distribution of HI gas at high $`z`$ and they consider the possibility of detecting HI emission from galaxies at $`z>2`$.
Tozzi et al. (1999) and Shaver et al. (1999) have studied the possibility of detecting the HI in the IGM at $`z>5`$. The state of the HI at these redshifts is unknown and these works are based on different scenarios for the reionization of the universe.
The work presented here differs significantly from the previous papers in that:
* It is restricted to $`z<3.5`$ where the HI content of the universe is well determined from DLA absorption studies. We use the results of these observations as inputs to our calculations.
* Instead of looking at the possibility of detecting individual features (as has been the focus of a large number of previous papers) we have studied the statistical properties of the fluctuations in the brightness temperature in low frequency radio maps. The statistical quantity we have studied is the cross-correlation between the temperature fluctuations along different lines of sight in radio map made at different frequencies. Individual features corresponding to protoclusters are rare events and protoclusters with flux in the range $`1.53\mathrm{m}\mathrm{Jy}`$ are predicted to occur with abundances in the range $`10^810^7\mathrm{Mpc}^3`$ in the CDM model (Subramanian and Padmanabhan, 1993). Even small density fluctuations will contribute to the cross-correlation signal and our proposal has the advantage of simultaneously using the signal in all the pixels in all the frequency channels across the bandwidth of observation.
The structure of the paper is as follows: in §2 we discuss the basic formalism of calculating the brightness temperature and fluctuations in the redshifted 21 cm radiation, in §3 we discuss the numerical results for two currently favoured cosmological models. In §4, the observational difficulties in the presence galactic and extragalactic foregrounds are presented, and we discuss a possible strategy for overcoming these. §5 gives a summary of our main results.
## 2 Formalism.
We treat the HI in damped Lyman-$`\alpha `$ (DLA) systems as a continuous distribution with $`n_{\mathrm{HI}}(𝒙,t)`$ denoting the comoving number density of HI atoms in the excited state of the hyperfine transition. Such a treatment is justified in a situation where the resolution of the observations is not sufficient to detect individual DLAs. In addition, the fact that the HI actually does not have a continuous distribution but is distributed in discrete objects can, if required. be taken into account when calculating statistical measures of the fluctuation in the HI distribution.
The HI emission which is at a frequency $`\nu _e=1420\mathrm{M}\mathrm{H}\mathrm{z}`$ in the rest frame of the gas it is emitted from is redshifted to a frequency $`\nu `$ for an observer located at the origin of the coordinate system. Taking into account the effects of both the expansion of the universe as well as the peculiar velocity $`𝒗(𝒙,t)`$ of the HI at the time of emission and $`𝒗(0,t_0)`$ of the observer at the present time, the comoving coordinate $`𝒙`$ of the HI and $`\nu `$ are related by
$$𝒙=𝒏c_{\frac{\nu }{\nu _e(1W)}}^1\frac{da}{a^2H(a)}$$
(1)
in a spatially flat universe. Here $`𝒏=𝒙/x`$ is a unit vector along the line of sight to the HI, $`W(𝒙)=𝒏[𝒗(𝒙,t)𝒗(0,t_0)]/c`$ accounts for the effects of the peculiar velocities and
$$H(a)=\frac{\dot{a}(t)}{a(t)}=H_0\sqrt{\mathrm{\Omega }_{m0}a^3+\lambda _0}$$
(2)
is the Hubble parameter at the epoch when the scale factor has value $`a`$.
When discussing observations of the redshifted HI emission, it is convenient to use $`𝝂=𝒏\nu `$ to simultaneously denote the frequency and the direction of the observation. The vector $`𝝂`$ fixes both the comoving position $`𝒙`$ of the HI from where the radiation originates and the time $`t`$ at which the radiation originates, and we shall use $`𝝂`$ and $`(𝒙,t)`$ interchangeably. Here we calculate how $`T(𝝂)`$ the brightness temperature of the radiation is related to the density $`n_{\mathrm{HI}}(𝒙,t)`$ and the the peculiar velocity $`𝒗(𝒙,t)`$ of the HI.
The energy flux in the frequency interval $`d^3\nu `$ can be calculated if we know the number of excited HI atoms in the comoving volume $`d^3x`$ from where the radiation originates and it is given by
$$\mathrm{Flux}=\frac{h_P\nu _eA_{21}n_{\mathrm{HI}}(𝒙,t)d^3x}{4\pi r_L^2(𝒙)}$$
(3)
Here $`r_L`$ is the luminosity distance which is given by
$$r_L(𝒙)=x[1W]\frac{\nu _e}{\nu }$$
(4)
The flux is also related to the specific intensity as follows
$$\mathrm{Flux}=\frac{I(𝝂)}{\nu ^2}d^3\nu $$
(5)
which allows us to calculate the specific intensity to be
$$\frac{I(𝝂)}{\nu ^2}=\frac{h_P\nu _eA_{21}n_{\mathrm{HI}}(𝝂)}{4\pi \nu _e^2}\left\{\frac{1}{[1W(𝝂)]}\frac{\nu }{x}\right\}^2\left|\frac{𝒙}{𝝂}\right|.$$
(6)
where $`|\frac{𝒙}{𝝂}|`$ is the Jacobian of the transformation from $`𝝂`$ to $`𝒙`$ given in equation (1).
It should be noted here that the effect of the peculiar velocity $`(i.e.W)`$ can be neglected if we restrict the analysis to scales which are much smaller than the horizon. This is not true for the spatial derivatives of the peculiar velocities which appears in the Jacobian and we retain such terms in our analysis. Calculating the Jacobian gives us
$$\left|\frac{𝒙}{𝝂}\right|=\left(\frac{x}{\nu }\right)^2\frac{c}{H(\frac{\nu }{\nu _e})}\frac{\nu _e}{\nu ^2}\left[1\frac{\nu _e}{\nu }\frac{(𝒏\mathbf{})(𝒏𝒗)}{H(\frac{\nu }{\nu _e})}\right].$$
(7)
Equations (6) and (7) allows us to calculate the brightness temperature of the radiation
$$T(𝝂)=\left(\frac{c^2}{2k_B}\right)\frac{I(𝝂)}{\nu ^2}$$
(8)
which gives us
$$T(𝝂)=\frac{T_{21}N_{21}(𝝂)}{8\pi }\frac{A_{21}}{H(\frac{\nu }{\nu _e})}\left(\frac{\nu _e}{\nu }\right)^2\left[1\frac{\nu _e}{\nu }\frac{(𝒏\mathbf{})(𝒏𝒗)}{H(\frac{\nu }{\nu _e})}\right].$$
(9)
where $`T_{21}=h_p\nu _e/k_B`$ and $`N_{21}(𝝂)=(c/\nu _e)^3n_{HI}(𝝂)`$ is the number of HI atoms in the excited state in a comoving volume $`(21\mathrm{c}\mathrm{m})^3`$.
The number density $`n_{\mathrm{HI}}(𝝂)`$ can be written as a sum of two parts, namely the mean $`\overline{n}_{HI}(\nu )`$ and the fluctuation $`\mathrm{\Delta }n_{HI}(𝝂)`$. We assume that the fluctuation in the number density of HI can be related to the perturbations in the underlying dark matter distribution $`\delta (𝝂)`$ through a time dependent linear bias parameter $`b(\nu )`$ which gives us
$$N_{21}(𝝂)=\overline{N}_{21}(\nu )[1+b(\nu )\delta (𝝂)].$$
(10)
We use this to calculate the isotropic part of the temperature
$$\overline{T}(\nu )=\frac{2.38\mathrm{K}h^1\overline{N}_{21}(\nu )}{\sqrt{\mathrm{\Omega }_{m0}(\frac{\nu _e}{\nu })^3+\lambda _0}}\left(\frac{\nu _e}{\nu }\right)^2$$
(11)
and the fluctuation
$$\mathrm{\Delta }T(𝝂)=\overline{T}(\nu )\left[b(\nu )\delta (𝝂)\frac{\nu _e}{\nu }\frac{(𝒏\mathbf{})(𝒏𝒗)}{H(\frac{\nu }{\nu _e})}\right]$$
(12)
Now, in the linear theory of density perturbations, both the perturbation and the peculiar velocity can be expressed in terms of a potential i.e.
$$\delta (𝒙,t)=D(t)^2\psi (𝒙)$$
(13)
and
$$𝒗(𝒙,t)=f(\mathrm{\Omega }_m)H(t)a(t)D(t)\psi (𝐱).$$
(14)
where $`D(t)`$ is the growing mode of linear density perturbations (Peebles 1980). In a spatially flat universe the function $`f`$ can be well approximated by the form (Lahav et. al. 1991)
$$f(\mathrm{\Omega }_m)=\mathrm{\Omega }_m^{0.6}+\frac{1}{70}[1\frac{1}{2}\mathrm{\Omega }_m(1+\mathrm{\Omega }_m)]$$
(15)
where the time dependence of $`\mathrm{\Omega }_m`$ can be expressed as
$$\mathrm{\Omega }_m(\nu )=\left(\frac{H_0}{H(\frac{\nu }{\nu _e})}\right)^2\left(\frac{\nu _e}{\nu }\right)^3\mathrm{\Omega }_{m0}.$$
(16)
Using this and defining a time dependent linear redshift space distortion parameter $`\beta (\nu )=f(\mathrm{\Omega }_m)/b(\nu )`$ we can express the fluctuation in the temperature as
$$\mathrm{\Delta }T(𝝂)=\overline{T}_A(\nu )b(\nu )D(\nu )\left[^2+\beta (\nu )(𝒏)^\mathrm{𝟐}\right]\psi (𝒙)$$
(17)
We next calculate the cross-correlation between the temperature fluctuations along two different lines of sight $`𝒏_1`$ and $`𝒏_2`$ at two different frequencies $`\nu _1`$ and $`\nu _2`$. The quantity we consider is the correlation
$$w(𝝂_1,𝝂_2)=\mathrm{\Delta }T(𝝂_1)\mathrm{\Delta }T(𝝂_2)$$
(18)
which is a function of the two frequencies $`\nu _1`$, $`\nu _2`$ and $`\theta `$ the angle between the two lines of sight. Using equation (17) and defining $`\varphi (x_{12})=\psi (𝒙_1)\psi (𝒙_2)`$ to be the two point correlation of the potential $`\psi (𝒙)`$, we obtain
$$w(𝝂_1,𝝂_2)=\overline{T}_1\overline{T}_2D_1D_2b_1b_2\left[^2+\beta _1(𝒏_1\mathbf{})^2\right]\left[^2+\beta _2(𝒏_2\mathbf{})^2\right]\varphi (x_{12})$$
(19)
where $`x_{12}=𝒙_1𝒙_2`$ and we have used the notation $`\overline{T}_1=\overline{T}(\nu _1)`$, etc.
The function $`\varphi (x)`$ is related to the correlation of the perturbations in the underlying dark matter distribution, and $`\xi (𝒙_1,t_1,𝒙_2,t_2)=\delta (𝒙_1,t_1)\delta (𝒙_2,t_2)`$ the two point correlation of the perturbation in the dark matter density at the point $`𝒙_1`$ at the epoch $`t_1`$ and $`𝒙_2`$ at the epoch $`t_2`$ can be written in terms of the potential $`\varphi `$ as
$$\xi (𝒙_1,t_1,𝒙_2,t_2)=D(t_1)D(t_2)\stackrel{~}{\xi }(x_{12})=D(t_1)D(t_2)^4\varphi (x_{12}).$$
(20)
We have introduced the function $`\stackrel{~}{\xi }(x_{12})`$ so that we can write $`\xi (𝒙_1,t_1,𝒙_2,t_2)`$ as a product of two parts, one which has the temporal variation and another the spatial variation. Equation (20) can be inverted to express the different derivatives of $`\varphi (x_{12})`$ which appear in equation (19) in terms of moments of $`\stackrel{~}{\xi }(x_{12})`$ which are defined as
$$\stackrel{~}{\xi }_n(x_{12})=\frac{n+1}{x_{12}^{n+1}}_0^{x_{12}}\stackrel{~}{\xi }(y)y^n𝑑y.$$
(21)
The form of the angular correlation is further simplified if we restrict $`\theta `$ to be very small. Under this assumption $`𝒏_1𝒏_2`$, and we use $`𝒏`$ to denote the common line of sight. We also use $`\mu =𝒏(𝒙_1𝒙_2)/x_{12}`$ for the cosine of the angle between the line of sight $`𝒏`$ and the vector $`𝒙_1𝒙_2`$ joining the two points between which we are measuring the correlation. The relation between the different distances and angles is shown in figure 1, and we have $`x_{12}=\sqrt{x_1^2+x_2^22x_1x_2\mathrm{cos}(\theta )}`$ and $`\mu =(x_1x_2\mathrm{cos}(\theta ))/x_{12}`$. Using these we finally obtain the following expression for the two point correlation of the temperature
$`w(\nu _1,\nu _2,\theta )`$ $`=`$ $`\overline{T}_1\overline{T}_2D_1D_2b_2b_2\{\left[(1+\beta _1\mu ^2)(1+\beta _2\mu ^2)\right]\stackrel{~}{\xi }(x_{12})`$ (22)
$`+`$ $`\left[({\displaystyle \frac{1}{3}}\mu ^2)(\beta _1+\beta _2)+({\displaystyle \frac{1}{2}}3\mu ^2+{\displaystyle \frac{5}{2}}\mu ^4)\beta _1\beta _2\right]\stackrel{~}{\xi _2}(x_{12})`$
$``$ $`[{\displaystyle \frac{3}{10}}+3\mu ^2{\displaystyle \frac{7}{2}}\mu ^4]\stackrel{~}{\xi _4}(x_{12})\}`$
## 3 Predictions for different models.
The density of HI ($`\mathrm{\Omega }_{\mathrm{HI}}`$) in DLAs has been determined reasonably well for $`0z3`$ from absorption studies (Lanzetta, Wolfe, & Turnshek 1995) and they find that the observed evolution of $`\mathrm{\Omega }_{\mathrm{HI}}`$ is well approximated by $`\mathrm{\Omega }_{\mathrm{HI}}(z)=\mathrm{\Omega }_{\mathrm{HI0}}\mathrm{exp}(\alpha z)`$ with $`\mathrm{\Omega }_{\mathrm{HI0}}=0.18\pm 0.04\times 10^3h^1`$ and $`\alpha =0.60\pm 0.15`$ for $`q_0=0`$, and $`\mathrm{\Omega }_{\mathrm{HI0}}=0.19\pm 0.04\times 10^3h^1`$ and $`\alpha =0.83\pm 0.15`$ for $`q_0=0.5`$, We have used this to calculate the isotropic part of the background brightness temperature due to the redshifted HI emission (equation 17). We have considered two spatially flat FRW cosmological models with parameters (I) $`\mathrm{\Omega }_{m0}=1.0`$ and $`\lambda _0=0`$, and (II) $`\mathrm{\Omega }_{m0}=0.3`$ and $`\lambda _0=0.7`$, with $`h=0.5`$ in both cases. It should be pointed out that the fits given Lanzetta et al. (1995) are not valid for the model with a cosmological constant, and we use the $`q_0=0`$ fit in this case.
Figure 2 shows the brightness temperature as a function of frequency for both the cases.
We find that the temperature increases rapidly as we go to lower frequencies (higher $`z`$) and it is around $`1\mathrm{mK}`$ at $`\nu 330\mathrm{M}\mathrm{H}\mathrm{z}`$ which corresponds to $`z3`$. The increase in the temperature is a direct consequence of the increase of HI density with increasing redshift. The HI density is not very well determined at $`z>3`$ and there is evidence that $`\mathrm{\Omega }_{\mathrm{HI}}`$ falls off at higher redshifts (Storrie–Lombardi et al. 1996).
We next consider $`w(\nu _1,\nu _2,\theta )`$ the cross-correlation in the temperature fluctuation at different frequencies which we have calculated for the two cases discussed above. We have calculated the dark-matter two point correlation function (equation 20) using the analytic fitting form for the CDM power spectrum given by Efstathiou, Bond and White (1992). For case (I) we use a value of the shape parameter $`\mathrm{\Gamma }=0.25`$ and for (II) we use $`\mathrm{\Gamma }=0.14`$. The power spectra are normalized using the results of Bunn and White (1996) based on the 4-year COBE data. It should be noted that he redshift evolution of the matter correlation function and $`f(\mathrm{\Omega }_m)`$ are quite different in the two cases that we have considered. The predictions for $`w(\nu _1,\nu _2,\theta )`$ are shown in figures 3 and 4. We have assumed that the HI faithfully traces the matter distribution and set $`b(\nu )=1`$ throughout.
In our analysis we have kept $`\nu _1`$ fixed at $`320\mathrm{M}\mathrm{H}\mathrm{z}`$ and let $`\nu _2`$ vary over a band of $`16\mathrm{M}\mathrm{H}\mathrm{z}`$ centered around $`\nu _1`$, while $`\theta `$ takes values upto $`2^{}`$. In order to measure $`w(\nu _1,\nu _2,\theta )`$ over this range we would need radio images of the background temperature fluctuation in a $`2^{}\times 2^{}`$ field at different frequencies in a $`16\mathrm{M}\mathrm{H}\mathrm{z}`$ band centered around $`320\mathrm{M}\mathrm{H}\mathrm{z}`$ and $`w(\nu _1,\nu _2,\theta )`$ could be estimated using
$$w(\nu _1,\nu _2,\theta )=\mathrm{\Delta }T_{\nu _1}(𝒏_1)\mathrm{\Delta }T_{\nu _2}(𝒏_2)$$
(23)
where $`\mathrm{\Delta }T_{\nu _1}(𝒏_1)`$ refers to the temperature fluctuation along the direction $`𝒏_1`$ in the image made at frequency $`\nu _1`$, and the angular brackets denote average over all pairs of directions $`𝒏_1`$ and $`𝒏_2`$ which are separated by an angle $`\theta `$. The central frequency, bandwidth and angular range have been chosen keeping in mind the Giant Meter wave Radio Telescope (GMRT, Swarup et al. 1991) which has recently become operational. The frequency intervals and angular separation can be converted to a corresponding comoving length-scale and in case (I) $`1\mathrm{M}\mathrm{H}\mathrm{z}=8.9\mathrm{Mpc}`$ and $`1^{^{}}=1.8\mathrm{Mpc}`$, while in case (II) $`1\mathrm{M}\mathrm{H}\mathrm{z}=16.0\mathrm{Mpc}`$ and $`1^{^{}}=2.7\mathrm{Mpc}`$. The linear theory of density perturbations which we have used here can be applied at scales $`10\mathrm{M}\mathrm{p}\mathrm{c}`$, which covers most of the region shown in the figures.
We find that $`w(\nu _1,\nu _2,\theta )`$ is between $`10^7\mathrm{K}^2`$ to $`10^8\mathrm{K}^2`$ when the comoving-distance corresponding to the separation in frequencies and direction is less between $`10\mathrm{M}\mathrm{p}\mathrm{c}`$ and $`40\mathrm{M}\mathrm{p}\mathrm{c}`$ beyond which it falls off. The cross-correlation between the temperature fluctuations is expected to be larger at small scales where the linear theory of density perturbations will not be valid and we have not considered these scales here. The temperature fluctuations are anti-correlated when the separation in frequencies exceeds the angular separation. This occurs because of the effect of the peculiar velocity which produces a “distortion” very similar to the effect it has on the two-point correlation function in redshift surveys.
## 4 Observational prospects
In this section we discuss the prospects of actually observing the background radiation from the HI in damped Lyman-$`\alpha `$ clouds. Our discussion is restricted to observations at around $`320\mathrm{MHz}`$ largely because we have reliable estimates of the HI content at $`z3`$ and there are indications that $`z3`$ might be the redshift where the HI content of the universe is maximum (Storrie–Lombardi et al. 1996). Another reason for our choice of this frequency range is that the GMRT (Swarup et al. 1991) which is already functioning at these frequencies is expected to have an angular resolution of $`10^{^{\prime \prime }}`$ and reach noise levels of $`100\mu \mathrm{K}`$ in around 10 hrs of integration. The angular resolution and sensitivity of the GMRT will be sufficient for detecting both the isotropic component as well as the correlations in the component of the background radiation arising from the HI at $`z3`$ (Figures 23 and 4) provided we can distinguish this component from the contribution due to other sources.
### 4.1 Galactic and extra-galactic foregrounds
Any observations at low frequencies will have a very large contribution from the synchrotron radiation from our own galaxy. Observations at $`408\mathrm{MHz}`$ (Haslam et al. 1982) with resolution $`1^{}\text{x}1^{}`$ indicate a minimum temperature of $`10\mathrm{K}`$ at $`408\mathrm{MH}z`$. Using $`T\nu ^{2.7}`$ as indicated by the spectral index of galactic synchrotron radiation, the temperature at $`\nu 320\mathrm{MHz}`$ is $`20\mathrm{K}`$. Comparing this number with the expected background from redshifted HI radiation (Figure 2) the galactic foreground is seen to be several orders above the expected signal. Another quantity of interest to us is the fluctuations in the galactic radiation at angles $`\stackrel{<}{}1^{}`$. However as the resolution of Haslam maps is $`1^{}`$ it cannot be used to make any predictions about the fluctuations at the angular scales of interest.
Another source of contamination is the continuum radiation from unresolved extragalactic sources (the resolved ones will be removed from the image before analysis) and we use results from the recent FIRST survey (for details see White et al. 1997 and references therein) to estimate this. Tegmark & Estathiou (1995) provide an analytic fit for the number of sources per unit flux $`\varphi [\mathrm{mJy}]`$ per steradian in this survey, and at flux levels $`\varphi 100\mathrm{m}\mathrm{J}\mathrm{y}`$ this can be approximated by
$$\frac{dn}{d\varphi }=\frac{5.24\times 10^5}{\mathrm{mJy}\mathrm{sr}}\left(\frac{\varphi }{0.75\mathrm{mJy}}\right)^{1.65}.$$
(24)
We use this to estimate the total contribution from sources fainter than $`100\mu \mathrm{Jy}`$, and converting this to brightness temperature we find that the continuum emission from unresolved radio sources is expected to produce a background radiation with brightness temperature $`.1\mathrm{K}`$ at $`1.5\mathrm{Ghz}`$. For the purposes of estimating an order-of-magnitude we assume this value to be representative of what we expect at $`320\mathrm{MHz}`$.
The two-point correlation function of the sources detected by FIRST has been estimated (Cress et al. 1996): $`w(\theta )0.18\theta ^{1.1}`$, with $`\theta `$ in $`\mathrm{arcminutes}`$. Little is known about the correlations expected in the fluctuations in the contribution to the background radiation from unresolved radio sources.
It is clear from the foregoing discussion that both galactic and extragalactic continuum signals are likely to be so high that they would totally swamp the signal which we want to detect and therefore it is not possible to directly detect the radiation from HI unless we find some method for removing the foregrounds.
At present the most promising strategy is to use the proposed observations themselves to determine the foreground and remove this from the data, and we next discuss a possible method for doing this.
### 4.2 Removing foregrounds
We consider the GMRT band centered at $`325\mathrm{MHz}`$ for our discussion in this section. The total bandwidth of $`16\mathrm{MHz}`$ is divided in 128 channels each with $`\mathrm{\Delta }\nu =1.25\mathrm{kHz}`$, and the signal in any frequency channel is expected to be dominated by the foregrounds.
Here we briefly review the standard method of continuum subtraction which is generally used in spectral line detections and then discuss how this method can be used for the analysis proposed here. The reader is referred to Subrahmanyam and Anantharamaiah (1990) and references therein for examples of how this method is applied in searches for HI emission from protoclusters.
We represent the observed signal $`y_i(𝒏)`$ in the $`i`$ th frequency channel in the direction $`𝒏`$ as:
$$y_i(𝒏)=x_i(𝒏)+f_i(𝒏)+N_i$$
(25)
where $`x_i(𝒏)`$ and $`f_i(𝒏)`$ are the contribution to the output signal from the HI and the foregrounds respectively, and $`N_i`$ is the receiver noise.
We should at this stage remind ourselves that foregrounds are sources which are emitting continuum radiation—the primary such source being the synchrotron radiation from our own galaxy. We also have the contribution from unresolved extragalactic radio sources which, as we have seen above, is a smaller contribution. The sum total of these along any line of sight is expected to be a smooth function of the frequency. On the contrary the contribution from the HI will come from individual damped Lyman $`\alpha `$ clouds along the line of sight with a velocity width $`200\mathrm{km}\mathrm{sec}^1`$ (Prochaska & Wolfe 1998) and this will correspond to lines of width:
$$\mathrm{\Delta }\nu \frac{\mathrm{\Delta }V}{c}\nu 0.2\mathrm{MHz}.$$
(26)
Although the total number of damped Lyman-$`\alpha `$ clouds in a $`1^{}\times 1^{}\times \mathrm{\hspace{0.17em}16}\mathrm{MHz}`$ field is expected to be quite large $`(3\times 10^4)`$ we expect to have only one damped Lyman-$`\alpha `$ cloud along a single $`20^{^{\prime \prime }}\times 20^{^{\prime \prime }}\times 16\mathrm{MHz}`$ synthesized GMRT beam. The probability of two damped Lyman-$`\alpha `$ clouds occurring right next to each other in the same synthesized beam and causing confusion is quite small, and hence the signal we are looking for will come in the form of lines of width $`0.2\mathrm{MHz}`$ which we will not be able to detect as the individual lines will be swamped by both the foreground as well as the noise. What we can do is to fit the smooth component of the signal $`y_i(𝒏)`$ by a function $`F_i(𝒏)`$ and subtract this from the output signal to remove the effect of the foregrounds
$$S_i(𝒏)=y_i(𝒏)F_i(𝒏)=x_i(𝒏)+N_i(𝒏)+E_i(𝒏)$$
(27)
where we call $`S_i(𝒏)`$ the reduced signal and $`E_i(𝒏)`$ is the possible error in the fitting procedure. If our foreground subtraction works correctly, the reduced signal should have contributions from only the HI signal and the noise.
We first determine the mean reduced signal
$$\overline{S}_i=S_i(𝒏)=x_i(𝒏)+N_i(𝒏)+E_i(𝒏)$$
(28)
where the angular brackets denote average over all lines of sight. If the errors in the foreground subtraction can be made smaller than the signal i.e. $`E_i(𝒏)<x_i(𝒏)`$ then this will give an estimate of the background brightness temperature due to the HI emission i.e. $`\overline{T}(\nu )=\overline{S}_i`$ where the index $`i`$ refers to the channel with central frequency $`\nu `$.
We next consider the quantity
$$\mathrm{\Delta }S_i(𝒏)=S_i(𝒏)\overline{S}_i$$
(29)
which gives the fluctuation in the brightness temperature. We can use this to estimate the cross-correlation in the fluctuations at different frequencies
$$w(\nu _1,\nu _2,\theta )=\mathrm{\Delta }S_i(𝒏)\mathrm{\Delta }S_j(𝐦)$$
(30)
where $`i`$ and $`j`$ are the channel with frequency $`\nu _1`$ and $`\nu _2`$ respectively, and $`𝒏`$ and $`𝐦`$ are two lines of sight separated by an angle $`\theta `$ and the angular brackets denote average over all such pairs of lines of sight.
## 5 Summary and Conclusions
We have investigated the contribution from the HI in unresolved damped Lyman-$`\alpha `$ clouds at high redshifts to the background radiation at low frequency radio waves (meter waves). The isotropic part of this radiation depends on the density of HI and the background cosmological model, while the fluctuations in this component of the background radiation have an added dependence on the fluctuations in the distribution of the damped Lyman-$`\alpha `$ clouds and their peculiar velocities.
We have used estimates of the HI density available from absorption studies to calculate the brightness temperature of this radiation. We find that this has a value $`1\mathrm{mK}`$ at $`320\mathrm{MHz}`$ which corresponds to $`z3`$. The distribution of damped Lyman-$`\alpha `$ clouds is assumed to trace the underlying dark matter distribution which also determines the peculiar velocities. Using this and the linear theory of density perturbations, we have calculated the relation between the fluctuation in this component of the background radiation and the density perturbations at high $`z`$. Observations of the cross-correlations of the fluctuations at different sight lines across images produced at different frequencies holds the possibility of allowing us to probe the two point correlation function (or power spectrum) at high redshifts. We have calculated the expected cross-correlations for two currently acceptable CDM models and find it to be in the range $`10^7\mathrm{K}^2`$ to $`10^8\mathrm{K}^2`$ at $`\nu 320\mathrm{M}\mathrm{H}\mathrm{z}`$ for separations in sight lines and frequencies such that the corresponding spatial separation is in the range $`10\mathrm{Mpc}`$ and $`40\mathrm{Mpc}`$. The cross-correlations are expected to be larger at smaller scales where the linear theory cannot be applied. Our results show that both the isotropic background (Figure 2) and its fluctuations (Figure 3 and 4) can be detected by GMRT which is the largest telescope operating at meter waves at present, provided this signal can be distinguished from other sources which contribute to the low frequency background radiation.
The biggest obstacles in detecting the HI contribution are the galactic and extra-galactic foregrounds, both of which are many orders larger than the signal we want to detect. The fact that both those sources of contamination emit continuum radiation while the HI contribution is from individual damped Lyman-$`\alpha `$ clouds each of which emits a spectral line with a relatively small velocity width keeps alive the possibility of being able to distinguish this signal from the contamination. We have, in this paper, considered one possible approach which might allow us to model and subtract the foreground along any line of sight. More work is needed in this direction and work is currently underway in investigating other viable possibilities for foreground removal.
All the authors would like to thank Jayaram Chengalur for many useful discussions on issues related to the GMRT and foregrounds.
|
warning/0003/cond-mat0003453.html
|
ar5iv
|
text
|
# References
Napoli DSF-T-44/99
A Conformal Field Theory description of the Paired and parafermionic states in the Quantum Hall Effect <sup>1</sup><sup>1</sup>1Work supported in part by EC n. FMRX-CT96-0045.
Gerardo Cristofano Giuseppe Maiella
and
Vincenzo Marotta
Dipartimento di Scienze Fisiche
Universitá di Napoli “Federico II”
and
INFN, Sezione di Napoli
Mostra d’Oltremare Pad.19-I-80125 Napoli, Italy<sup>2</sup><sup>2</sup>2E:mail: gerardo.cristofano(giuseppe.maiella;vincenzo.marotta)@na.infn.it
> Abstract
>
> We extend the construction of the effective conformal field theory for the Jain hierarchical fillings proposed in to the description of a quantum Hall fluid at non standard fillings $`\nu =\frac{m}{pm+2}`$. The chiral primary fields are found by using a procedure which induces twisted boundary conditions on the $`m`$ scalar fields; they appear as composite operators of a charged and neutral component. The neutral modes describe parafermions and contribute to the ground state wave function with a generalized Pfaffian term. Correlators of $`N_e`$ electrons in the presence of quasi-hole excitations are explicitly given for $`m=2`$.
>
> Keyword: Vertex operator, Kac-Moody algebra, Quantum Hall Effect
The experimental evidence of a Hall plateau at filling $`\nu =\frac{5}{2}`$ has recently spurred a renewed interest in a deeper understanding of the underlying physics at plateaus which do not fall into the hierarchical scheme . To such an extent a pairing picture, in which pairs of spinless or spin-polarized fermions condense, has been presented for the non-standard fillings $`\nu =\frac{1}{q}`$, $`q>0`$ and even. As a result the ground state gets described in terms of the Pfaffian (the so called Pfaffian state) and the non-Abelian statistics of the fractional charged excitations evidenced . More recently it has been argued that the non-Abelian statistics might come out from constraining an Abelian theory, by employing the Meissner effect in the neutral sector; that would also account for the pairing phenomenon.
Lately a Conformal Field Theory (CFT) description in terms of composite fermions for the Jain filling fractions $`\nu =\frac{m}{2pm+1}`$ of a Quantum Hall Fluid (QHF) has been proposed has been proposed . The composite fermions are described by composite vertex operators, which are (the image of) primary fields of a CFT with central charge $`c=m`$ in the Lowest Landau Level (LLL), and factorize into a charged and a neutral component. The neutral component assures the locality properties of the composite electron field and the single-valuedness of the ground state wave function.
Further due to the $`Z_m`$ symmetry in the neutral sector the general $`N_e`$ electron correlation function shows clustering properties which have been proposed in the context of paired Hall states .
In this letter we present a natural extension of the $`m`$-reduction procedure (previously applied in ref. to the description of the Jain fillings) to the case of the non-standard fillings $`\nu =\frac{m}{pm+2}`$. In particular for $`m=2`$ the neutral degrees of freedom of the composite electrons describe Majorana fermions and contribute to the ground state wave function with a Pfaffian term, in agreement with an early proposal . Further, for generic $`m`$, the neutral degrees of freedom are parafermions and give rise to the clustering phenomenon of $`m`$ objects, which simply reproduces the picture presented in ref. by employing a hamiltonian formalism with $`m+1`$ body interactions.
The letter is organized as follows: after a brief review of the $`m`$-reduction procedure we construct the parafermion vertex operators which together with the charged $`U(1)`$ component give rise to the primary fields of the CFT. It should be noticed an interesting phenomenon related to the existence of two classes of primary fields: the $`Z_m`$ twist invariant and the $`Z_m`$ non-invariant ones. In fact on physical grounds one can argue that only the invariant fields are relevant to the description of the pairing phenomenon. As a result the central charge of neutral sector of our CFT is given by $`c=\frac{2(m1)}{m+2}`$. An interpretation of such a phenomenon has been proposed in ref. in terms of a strong coupling between the symmetric modes.
We then give the ground state wave function as a correlator of the primary fields with integer electric charge and the quasi-holes wave function for the elementary excitations, with fractional charge, up to the correlator of the twist fields.
Our approach is meant to describe all the plateaus with even denominator starting from the bosonic Laughlin filling $`\nu =1/2`$, which is described by a CFT with $`c=1`$, in terms of a scalar chiral field compactified on a circle with radius $`R^2=1/\nu =2`$ (or the dual $`R^2=1/2`$). Then the $`U(1)`$ current is given by $`J(z|1,2)=i_zQ(z)`$, where $`Q(z)`$ is the compactified Fubini field with the standard mode expansion:
$$Q(z)=qiplnz+\underset{n0}{}\frac{a_n}{n}z^n$$
(1)
with $`a_n`$, $`q`$ and $`p`$ satisfying the commutation relations $`[a_n,a_n^{}]=n\delta _{n,n^{}}`$ and $`[q,p]=i`$.
The representations are realized by the vertex operators $`U^\alpha (z)=:e^{i\alpha Q(z)}:`$ with $`\alpha ^2=2`$ and conformal dimension $`h=1`$, with a consequent extension of the $`U(1)`$ symmetry to the $`SU(2)_1`$ affine one. Furthermore the theory contains the Virasoro algebra generated by the stress-energy tensor $`T(z|1,2)=\frac{1}{2}:\left(_zQ(z)\right)^2:`$.
In order to construct the $`\nu =m/2`$ filling we start with the set of fields in the above CFT (mother). Using the $`m`$-reduction procedure, which consists in considering the subalgebra generated only by the modes which are divided by an integer $`m`$, we get the image of an orbifold of a $`c=m`$ CFT (see ref. and references therein).
Also for the $`SU(2)`$ case the fields in the mother CFT can be factorized into irreducible orbits of the discrete group $`Z_m`$ which is a symmetry of the daughter theory and can be organized into components which have well defined transformation properties under this group.
In order to compare the image so obtained to the $`c=m`$ CFT, we map $`zz^{1/m}`$ and we will indicate the components in the base $`\widehat{z}=z^m`$ with an hatted symbol (for instance, $`\varphi (z)\widehat{\varphi }(z)`$). In particular any component in the subalgebra is a function only of the variable $`z^m`$.
In ref. it was also defined an isomorphism between fields on the $`z`$ complex plane and fields on the $`z^m`$ plane by means of the following identifications:
$$a_{nm+l}\sqrt{m}a_{n+l/m}q\frac{1}{\sqrt{m}}q$$
(2)
Let us first introduce the invariant scalar field
$$X(z|m,2)=\frac{1}{m}\underset{j=1}{\overset{m}{}}Q(\epsilon ^jz)$$
(3)
where $`\epsilon ^j=e^{i\frac{2\pi j}{m}}`$, corresponding to a compactified boson on a circle with radius now equal to $`R_X^2=2/m`$. This field depends only on powers of $`z^m`$ and satisfies trivial boundary conditions. It is the basic field of the $`U(1)`$ electrically charged sector of the theory where the charge is measured by the zero mode.
The non-invariant components, as a resulting image of $`m`$ constrained bosons, are expressed by
$$\varphi ^j(z|m,2)=Q(\epsilon ^jz)X(z|m,2)$$
(4)
with the condition $`_{j=1}^m\varphi ^j(z|m,2)=0`$.
These fields satisfy non-trivial twisted boundary conditions
$$\alpha \varphi ^j(\epsilon z|m,2)=\alpha \varphi ^{j+1}(z|m,2)+2\pi n\alpha pnZ$$
(5)
where the shift is due to the definition of index $`j`$ mod $`m`$.
The $`J(z|1,2)`$ current of the mother theory decomposes into a charged current given by $`J(z|m,2)=i_zX(z|m,2)`$ and $`m1`$ neutral ones $`_z\varphi ^j(z|m,2)`$.
In the same way every vertex operator in the mother theory can be factorized in a vertex that depends only on the invariant field:
$$𝒰^\alpha (z|m,2)=z^{\frac{\alpha ^2}{2}\frac{(m1)}{m}}:e^{i\alpha X(z|m,2)}:\alpha ^2=2$$
(6)
and in vertex operators depending on the $`\varphi ^j(z|m,2)`$ fields.
We also introduce the neutral components:
$$\psi _1^\alpha (z|m,2)=\frac{z^{\frac{\alpha ^2}{2}\frac{(1m)}{m}}}{m}\underset{j=1}{\overset{m}{}}\epsilon ^{\frac{\alpha ^2j}{2}}:e^{i\alpha \varphi ^j(z|m,2)}:$$
(7)
which satisfy the fundamental product:
$$\psi _1^\alpha (z|m,2)\psi _1^\beta (\xi |m,2)=\frac{z^{\frac{\alpha ^2}{2}\frac{1m}{m}}\xi ^{\frac{\beta ^2}{2}\frac{1m}{m}}}{m^2}\underset{j,j^{}=1}{\overset{m}{}}\epsilon ^{\frac{\alpha ^2j+\beta ^2j^{}}{2}}:e^{i\alpha \varphi ^j^{}(z|m,2)}e^{i\beta \varphi ^j(\xi |m,2)}:\frac{(\epsilon ^j^{}z\epsilon ^j\xi )^{\alpha \beta }}{(z^m\xi ^m)^{\frac{\alpha \beta }{m}}}$$
The set of primary fields generated by this product can be expressed in terms of the fundamental representations $`\mathrm{\Lambda }^i`$ of $`SU(m)`$ Lie algebra. In fact, defining
$$\varphi ^{\mathrm{\Lambda }^i}(z|m,2)=\underset{j=1}{\overset{i}{}}\varphi ^j(z|m,2)$$
(8)
and $`\varphi ^\mathrm{\Lambda }(z|m,2)`$, where $`\mathrm{\Lambda }=_{i=1}^{m1}l_i\mathrm{\Lambda }_i`$, and introducing the $`m`$-ality parameter $`a=_{i=1}^{m1}il_i`$ (mod $`m`$), which is invariant under the addition of any vector in the root lattice, the exact form of these fields can be deduced by the analysis of the OPE of eq.(S0.Ex1) for $`\alpha =\beta `$ to get the $`a=2`$ field. By repeated application of this analysis we can obtain the full set of fields:
$$\widehat{\psi }_a^\alpha (z|m,2)=\underset{j_1>j_2>\mathrm{}>j_a=1}{\overset{m}{}}f(\epsilon ^{j_1},\mathrm{},\epsilon ^{j_a},z^{1/m}):e^{i\alpha \widehat{\varphi }^{j_1}(z|m,2)}\mathrm{}e^{i\alpha \widehat{\varphi }^{j_k}(z|m,2)}:$$
(9)
where the functions $`f(\epsilon ^{j_1},\mathrm{},\epsilon ^{j_a},z)`$ can be extracted from the OPE relations.
The sum takes into account the fact that any field can be associated to the $`a`$-th fundamental representation of $`SU(m)`$ (namely, the antisymmetric tensor representation). In it was shown that these are a realization of parafermions and satisfy the operator product algebra :
$$\widehat{\psi }_a^\alpha (z|m,2)\widehat{\psi }_a^{}^\alpha (\xi |m,2)=\frac{C_{a,a^{}}}{(z\xi )^{\frac{2aa^{}}{m}}}\left[\widehat{\psi }_{a+a^{}}^\alpha (z|m,2)+O(z\xi )\right]a+a^{}<m$$
(10)
and for $`a+a^{}=m`$ the OPE contains the Virasoro algebra generators:
$$\widehat{\psi }_a^\alpha (z|m,2)\widehat{\psi }_{ma}^\alpha (\xi |m,2)=\frac{C_{a,ma}}{(z\xi )^{\frac{2a(ma)}{m}}}\left[1+\frac{2h_a}{c_\psi }\widehat{T}_\psi (z|m,2)(z\xi )^2+O(z\xi )^3\right]$$
(11)
where $`C_{a,a^{}}`$ are the structure constants and $`h_a`$, $`c_\psi `$ are the conformal dimensions and central charge of parafermions .
Moreover the $`SU(m)`$ representations that can appear are the fundamental ones $`\mathrm{\Lambda }_a`$, because in the OPE algebra of $`\psi _1^\alpha (z|m,2)`$ (corresponding to the $`\mathrm{\Lambda }_1`$ representation) only the fields $`\psi _a^\alpha (z|m,2)`$ with $`a\{1,\mathrm{},m1\}`$ appear while $`\psi _m^\alpha (z|m,2)`$ is the identity operator. That is an effect of the $`Z_m`$ invariance of the parafermions algebra.
Notice that no neutral currents are present in the above OPE, as one would expect from symmetry considerations.
It is well known that the $`c_X=1`$ rationalCFT with $`R_X^2=2/m`$ has $`2m`$ primary fields which can be parametrized by $`\alpha =\sqrt{\frac{2}{m}}a`$ and $`\alpha =\sqrt{\frac{1}{2m}}a`$, $`a=1,\mathrm{},m`$. In our formalism these fields appear together with the neutral ones into the composite operators:
$$\widehat{V}^{\sqrt{\frac{2}{m}}a}(z|m,2)=\widehat{𝒰}^{\sqrt{\frac{2}{m}}a}(z|m,2)\widehat{\psi }_a(z|m,2)$$
(12)
for quasi-particles and
$$\widehat{V}_{qh}^{\frac{a}{\sqrt{2m}}}(z|m,2)=\widehat{𝒰}^{\frac{a}{\sqrt{2m}}}(z|m,2)\widehat{\sigma }_a(z|m,2)$$
(13)
for quasi-holes; $`\widehat{\sigma }_a(z|m,2)`$ are the parafermionic twist fields. The electric charges and magnetic flux contents of such fields is given after eqs.(20, 21).
In eq.(13) we have not included the contribution coming from the $`\overline{\sigma }`$ fields of the $`Z_m`$ non-invariant theory. In fact they give rise to a term which does not survive after the projection to the LLL<sup>3</sup><sup>3</sup>3We notice that we do not have an explicit realization of the $`\overline{\sigma }`$ fields because we have projected the mother $`c=1`$ CFT onto the $`m`$-covering of the plane, so that those fields appear as branch cuts ..
The currents $`\widehat{V}^{\pm \sqrt{\frac{2}{m}}}(z|m,2)`$ and $`J(z|m,2)`$ generate the $`SU(2)_m`$ affine algebra while $`\widehat{V}_{qh}^{\frac{a}{\sqrt{2m}}}`$ are the primary fields of this algebra (see ref. for details).
The generator of the Virasoro algebra $`\widehat{T}(z|m,2)`$ was given in as the sum of two independent operators, one depending on the charged sector:
$$\widehat{T}_X(z|m,2)=\frac{1}{2}:\left(\widehat{_zX}(z|m,2)\right)^2:$$
(14)
and the other given in terms of the $`Z_m`$ twisted bosons $`\widehat{\varphi }^j(z|m,2)`$:
$$\widehat{T}_\psi (z|m,2)=\frac{2}{m+2}\left(\underset{j=1}{\overset{m}{}}\frac{:(\widehat{_z\varphi }^j(z))^2:}{2m^2}+\underset{j^{}j=1}{\overset{m}{}}\frac{\epsilon ^{j^{}+j}:e^{i\alpha \widehat{\varphi }^j^{}(z)}e^{i\alpha \widehat{\varphi }^j(z)}:}{2m^2z^2(\epsilon ^j^{}\epsilon ^j)^2}+\frac{m^21}{24mz^2}\right)$$
(15)
while the higher integer spin operators generating the full parafermionic $`𝒲_m`$ algebra was given in terms of the $`\widehat{\varphi }^j(z|m,2)`$ fields in . Notice that the vacuum expectation value of $`\widehat{T}_\psi `$ is zero due to the cancellation between the second and the third term in eq.(15).
It is not very hard to verify that the conformal dimensions of the fields of eq.(12) and eq.(13) are given by:
$$h_a=\frac{a^2}{m}+a\left(\frac{ma}{m}\right)=aa\{1,\mathrm{},m\}$$
(16)
for quasi-particles and
$$h_a^{qh}=\frac{a^2}{4m}+\frac{a(ma)}{2m(m+2)}=\frac{a(a+2)}{4(m+2)}a\{1,\mathrm{},m\}$$
(17)
for quasi-holes.
The contribution to the central charge $`c`$ is given by
$$c=1+\frac{2(m1)}{m+2}=\frac{3m}{m+2}$$
(18)
where 1 and $`{\displaystyle \frac{2(m1)}{m+2}}`$ come from the charged and neutral degrees of a freedom respectively.
We must notice that the contribution of the $`Z_m`$ invariant fields to the central charge is not equal to $`m`$. In fact the missing part is associated with the non-invariant fields defined as:
$$\psi ^{\alpha (i,j)}(z|m,2)=\frac{z^{\frac{\alpha ^2}{2}\frac{(1m)}{m}}}{m}(\epsilon ^{\frac{\alpha ^2i}{2}}:e^{i\alpha \varphi ^i(z|m,2)}:\epsilon ^{\frac{\alpha ^2j}{2}}:e^{i\alpha \varphi ^j(z|m,2)}:)$$
(19)
with $`ij=\{1,\mathrm{},m\}`$, which are primary fields of the complementary $`\overline{c}=\frac{m(m1)}{m+2}`$ theory. One can see also that these extra degrees of freedom can be gauged away by means of a coset reduction.
This appears to be a generalization of the Low-Barrier limit of a double-layer sample described in ref. where a strong tunnel effect was introduced. In our case the strong coupling between electrons of the different layers is simply a consequence of the induced $`Z_m`$ symmetry of the composite electrons states. This phenomenon was suggested to be relevant for CFT’s on algebraic curves in . It would be interesting to have a clear picture of the symmetries of the vacuum and excited states, so to understand the mechanism which decouples the $`Z_m`$ non-invariant theory from the invariant one.
We can now describe the generic filling $`\nu ={\displaystyle \frac{m}{pm+2}}`$ by flux attachment starting from $`\nu =m/2`$. In order to do so, we factorize the fields into two parts, the first describing the $`c_X=1`$ charged sector with radius $`R_X^2=\frac{pm+2}{m}`$, the second describing the neutral excitations, which are parafermions, with central charge $`c_\psi =\frac{2(m1)}{m+2}`$, for any $`pN`$.
The $`U(1)`$ sector is now described by the compactified boson $`X(z|m,pm+2)`$ and its related vertex operators $`𝒰^{\pm \alpha _l}(z|m,pm+2)`$, with $`\alpha _l=l/\sqrt{m(pm+2)}`$, $`l=1,\mathrm{},m(pm+2)`$, which produce excitations with anyonic statistics $`\theta =\pi \alpha _l^2`$. While the $`m1`$ neutral bosons $`\varphi ^j(z|m,2)`$ are independent from the flux number $`p`$.
To obtain a pure holomorphic function we will consider the correlator of the composite operators $`\widehat{V}^{\alpha _l}(z|m,pm+2)`$ with conformal dimensions:
$$h_l=\frac{l^2}{2m(pm+2)}+a\left(\frac{ma}{m}\right)l=(pm+2)aa=1,2,\mathrm{},m$$
(20)
They describe dressed $`a`$-electrons with electric charge $`q_a^e=a`$ and “magnetic charge” $`q_a^m=\frac{pm+2}{m}a`$ interacting through the neutral “cloud” associated with them (see eq.(12)).
Also we will consider correlators in which are present quasi-hole operators given by $`\widehat{V}_{qh}^{\alpha _l}(z|m,pm+2)`$ having conformal dimensions:
$$h_l^{qh}=\frac{l^2}{2m(pm+2)}+\frac{a(ma)}{2m(m+2)}$$
(21)
where $`(pm+2)(a1)<l<(pm+2)a`$ and $`a=1,2,\mathrm{},m`$.
Their electric and magnetic charges are $`q_l^e=\frac{l}{pm+2}`$ and $`q_l^m=\frac{l}{m}`$.
We should point out that $`m`$-ality in the neutral sector is coupled to the charged one in analogy to the case of Jain hierarchical fillings in order to assure the locality of electrons with respect to all the edge excitations . This follows from the fact that our projection when applied to a local field automatically couples the discrete $`Z_m`$ charge of $`U(1)`$ with the neutral sector, in order to give a totally single-valued composite field.
Also notice that the $`m`$-electron vertex operator does not contain any neutral field. Therefore, the $`m`$-electron wave function is realized only by means of the $`c_X=1`$ charged sector as proposed in ref..
We are now ready to give the holomorphic part of the ground state wave function for the generic filling $`\nu =\frac{m}{pm+2}`$. To such an extent we consider the $`N_e`$ single($`a=1`$)-electrons correlator which factorizes into a Laughlin-Jastrow type term coming from the charged sector:
$$<N_e\alpha |\underset{i=1}{\overset{N_e}{}}\widehat{𝒰}^\alpha (z_i|m,2pm+2)|0>=\underset{i<i^{}=1}{\overset{N_e}{}}(z_iz_i^{})^{p+\frac{2}{m}}$$
(22)
and a contribution coming from the neutral excitations:
$$<0|\underset{i=1}{\overset{N_e}{}}\widehat{\psi }_1^\alpha (z_i|m,2)|0>=\frac{_{\{j_i\}=1}^m\epsilon ^{\frac{\alpha ^2}{2}(2i1)j_i+j_i}_{\{j_i,j_i^{}\}=1}^m(\epsilon ^{j_i}z_i^{1/m}\epsilon ^{j_i^{}}z_i^{}^{1/m})^{\alpha ^2}}{_{i<i^{}=1}^{N_e}(z_iz_i^{})^{\frac{2}{m}}}$$
(23)
For $`N_e`$ multiple of $`m`$ we observe that the non analytic part of the neutral fields $`\widehat{\psi }^\alpha (z_i|m,2)`$ is necessary to eliminate the non integer part of the exponent in the correlator of the charged fields. We also point out that in our formalism the correlators are given, in principle, for any $`m`$ and $`N_e`$ and that follows from the projection procedure only.
By considering the case $`m=2`$, $`p`$ odd that is for $`\nu =1/q`$, $`q=p+1`$ we get for the correlator of $`N_e`$ single($`a=1`$)-electrons:
$$<N_e\alpha |\underset{i=1}{\overset{N_e}{}}\widehat{V}^{\sqrt{2q}}(z_i|2,2q)|0>=\underset{i<i^{}=1}{\overset{N_e}{}}(z_iz_i^{})^qPf\left(\frac{1}{z_iz_i^{}}\right)$$
(24)
(where $`Pf\left(\frac{1}{z_iz_i^{}}\right)=𝒜\left(\frac{1}{z_1z_2}\frac{1}{z_3z_4}\mathrm{}\right)`$ is the antisymmetrized product over pairs of electrons) which is in agreement with previous results . For the generic $`m`$ case, even though it is hard to work out explicitly the sum over the phases in eq.(23), for the $`N_e`$ point functions we found the clustering properties previously given in ref. in another framework by grouping particles into clusters of $`m`$.
In this last case the neutral modes describe parafermions and contribute to the ground state wave function with a generalized Pfaffian term. It would be very interesting to give an interpretation of the Pfaffian term (but also of its generalization presented in eq.(23)) in the context of a plasma description a la Laughlin in order to better understand the physics of the paired states .
In a similar way we also are able to evaluate correlators of $`N_e`$ single($`a=1`$)-electrons in the presence of quasi-hole excitations. In particular for $`m=2`$ and for two quasi-holes we get:
$`{\displaystyle \frac{<N_e\sqrt{2q}+2/\sqrt{2q}|_{i=1}^{N_e}\widehat{V}^{\sqrt{2q}}(z_i|2,2q)\widehat{V}^{1/\sqrt{2q}}(w_1|2,2q)\widehat{V}^{1/\sqrt{2q}}(w_2|2,2q)|0>}{<2/\sqrt{2q}|\widehat{V}^{1/\sqrt{2q}}(w_1|2,2q)\widehat{V}^{1/\sqrt{2q}}(w_2|2,2q)|0>}}`$
$`={\displaystyle \underset{i<i^{}=1}{\overset{N_e}{}}}(z_iz_i^{})^qPf\left({\displaystyle \frac{(z_iw_1)(z_i^{}w_2)+(z_i^{}w_1)(z_iw_2)}{z_iz_i^{}}}\right)`$ (25)
in agreement with the wave functions proposed in Ref..
The explicit evaluation of correlation functions for the twist fields lies outside the scope of the present paper, but it does not affect the main results given in eq.(25). On the other hand, it is fundamental to understand the non-Abelian statistics of the quasi-holes. We will analyze this aspect in a forthcoming paper.
We also point out that we are not considering the full set of primary fields in the theory. Indeed, also for $`p=0`$ there are neutral fields which correspond to the “termal” fields of parafermionic theory .
Acknowledgments \- We thank M. Huerta, G. Zemba and A. Sciarrino for useful comments and for reading the manuscript.
|
warning/0003/hep-th0003041.html
|
ar5iv
|
text
|
# 1 Introduction
## 1 <br>Introduction
Recently, there has been renewed activities in the study of gauged supergravity in various dimensions as well as in their solutions. This, to a large extent, is motivated by the conjectured equivalence between string theory on anti-de Sitter (AdS) spaces (times some compact manifold) and certain superconformal gauge theories living on the boundary of AdS . The theories of extended supergravity in various dimensions possess rigid symmetries. A subgroup of these symmetries can be gauged by the vector fields present in the ungauged theory. Gauged supergravity theories exist in space-time dimensions where supersymmetry allows the existence of a cosmological constant. In $`D=11`$, $`D=10`$ and $`D=9,`$ a cosmological constant is not possible.
In five dimensions, $`N=2`$ supergravity theories can be obtained by gauging the $`U(1)`$ subgroup of the $`SU(2)`$ automorphism group of the $`N=2`$ supersymmetry algebra, thus breaking $`SU(2)`$ down to the $`U(1)`$ group. The $`U(1)`$ gauge field introduced to gauge the theory can be taken as a linear combination of the abelian vector fields of the ungauged theory with a coupling constant $`g`$. The additional couplings of the fermi-fields of the ungauged theory to the $`U(1)`$ gauge vector field breaks supersymmetry which necessitates the addition of $`g`$-dependent gauge invariant terms in order to restore $`N=2`$ supersymmetry. The purely bosonic terms added produce the scalar potential of the theory .
Gauged supergravity can be obtained by compactifying higher dimensional supergravity on a group manifold. It is usually difficult to find a consistent ansatz for the compactification<sup>1</sup><sup>1</sup>1see and references therein.. In particular it was shown recently that the Freedman-Schwarz gauged $`N=4`$ supergravity can be obtained by compactifying ten dimensional supergravity on the $`SU(2)\times SU(2)`$ group manifold. The difficulty arises because one has to identify the vector fields coming from the compactification of the metric with the vector fields coming from the antisymmetric tensor. The vector fields coming from compactifying the metric on a group manifold behave properly as $`SU(2)\times SU(2)`$ gauge fields. The components of the antisymmetric tensor do not usually behave like $`SU(2)\times SU(2)`$ gauge fields, and for this to happen a very precise form for the ansatz of the antisymmetric tensor must be taken. This prescription also works in obtaining $`D=7`$ gauged supergravity by compactifying ten dimensional supergravity on $`SU(2)`$ group manifold as was recently shown in . In a related matter, there are now many known black hole solutions for gauged $`D=5`$ supergravity . These solutions could be promoted to solutions of ten and eleven dimensional supergravity if one can embed the five dimensional supergravity in the higher dimensional ones.
It is our purpose in this paper to show that by compactifying and truncating $`D=10`$ supergravity to $`D=5`$ on an $`SU(2)\times `$ $`U(1)^2`$ group manifold one obtains a gauged $`D=5`$ supegravity theory with one vector multiplet. The solutions for this model can then be lifted to seven, ten and eleven dimensions. This work is organized as follow. In section two, it is shown how to reduce $`D=10`$ supergravity to a particular gauged N=2 five dimensional supergravity theory coupled to one vector multiplet. The gauged theory obtained is the $`U(1)`$ gauged version of the model introduced by Strominger and Vafa. The five dimensional model obtained is then reformulated in the framework of very special geometry in section three. Some particular solutions of the five dimensional theory, as examples, are lifted to seven, ten and eleven dimensions in section four. Finally our results are summarized and discussed.
## 2 A $`D=5`$ Gauged Supergravity From $`D=10`$ Supergravity
In this section we consider the dimensional reduction of $`D=10`$ supergravity down to $`N=2`$, $`U(1)`$ gauged $`D=5`$ supergravity coupled to one vector multiplet. First, the bosonic part of $`N=1`$ supergravity action in ten dimensions is
$`S_{10}`$ $`=`$ $`{\displaystyle \widehat{e}\left(\frac{1}{4}\widehat{R}+\frac{1}{2}_M\widehat{\varphi }^M\widehat{\varphi }+\frac{1}{12}e^{2\widehat{\varphi }}\widehat{H}_{MNP}\widehat{H}^{MNP}\right)d^4xd^6z}`$ (1)
$``$ $`S_{\widehat{G}}+S_{\widehat{\varphi }}+S_{\widehat{H}}.`$
The notation used in this paper is as follows. We denote ten-dimensional quantities by hatted symbols. Base and tangent space indices are denoted by late and early capital Latin letters, respectively. For the four-dimensional space-time, we use late and early Greek letters, respectively, to denote base space and tangent space indices. Similarly, the internal base space and tangent space indices are denoted by late and early Latin letters, respectively.
$$\{M\}=\{\mu =0,\mathrm{},3;m=1,\mathrm{},6\},\{A\}=\{\alpha =0,\mathrm{},3;a=1,\mathrm{},6\}.$$
(2)
The general coordinates $`\widehat{x}^M`$ consist of spacetime coordinates $`x^\mu `$ and internal coordinates $`z^m`$. The flat Lorentz metric of the tangent space is chosen to be $`(+,,\mathrm{},)`$ with the internal dimensions all spacelike. Thus the metric is related to the vielbein by
$$\widehat{𝐠}_{MN}=\widehat{\eta }_{AB}\widehat{e}_M^A\widehat{e}_N^B=\eta _{\alpha \beta }\widehat{e}_M^\alpha \widehat{e}_N^\beta \delta _{ab}\widehat{e}_M^a\widehat{e}_N^b,$$
(3)
and the antisymmetric tensor field strength is
$$\widehat{H}_{MNP}=_M\widehat{B}_{NP}+_N\widehat{B}_{PM}+_P\widehat{B}_{MN}.$$
(4)
The coordinates $`z^m`$ span the internal compact group space. Thus we introduce the functions $`U_m^a(z)`$ which satisfy the condition
$$\left(U^1\right)_b^m\left(U^1\right)_c^n\left(_mU_n^a_nU_m^a\right)=\frac{f_{abc}}{\sqrt{2}},$$
(5)
Here $`f_{abc}`$ are the group structure constants and the internal space volume is $`\mathrm{\Omega }=|U_m^a|d^6𝐳`$. In the maximal case, i. e., $`SU(2)\times SU(2)`$, each $`S^3`$ factor admits invariant 1-form $`\theta ^a=\theta _i^adz^i`$, which satisfies
$$d\theta ^a+\frac{1}{2}ϵ_{abc}\theta ^b\theta ^c=0.$$
(6)
If one chooses
$$U_m^aU_i^a=\frac{\sqrt{2}}{g}\theta _i^a,$$
(7)
where $`g`$ is a coupling constant, then the structure constants will be given in terms of the coupling constant by $`f_{abc}=gϵ_{abc}`$. For the case where the coupling constant of one of the $`SU(2)`$ factors vanishes, the internal space becomes the group manifold $`SU(2)\times [U(1)]^3`$.
Our ansatz for the reduction to five dimensions is given by the following parameterization of the vielbein
$$\widehat{e}_M^A=\left(\begin{array}{cc}e^{\frac{5}{6}\widehat{\varphi }}e_\mu ^\alpha & \sqrt{2}A_\mu ^ae^{\frac{1}{2}\widehat{\varphi }}\\ 0& e^{\frac{1}{2}\widehat{\varphi }}U_m^a\end{array}\right).$$
(8)
Here the function $`U`$ depends only on the internal coordinates $`(5,\mathrm{},9)`$. As an internal space we take the group manifold $`SU(2)\times [U(1)]^2`$, this means that the following choice for the structure constants is taken,
$`f_{mnp}`$ $`=`$ $`gϵ_{mnp},m,n,p=5,6,7,`$
$`f_{mnp}`$ $`=`$ $`0,m,n,p=8,9`$ (9)
The $`S_{\widehat{G}}`$ term of $`D=10`$ supergravity gives upon reduction, the following $`D=5`$ Lagrangian
$$_{5G}=e\left(\frac{1}{4}R\frac{1}{8}e^{\frac{8}{3}\widehat{\varphi }}F_{\mu \nu }^aF^{\mu \nu a}+\frac{5}{6}g^{\mu \nu }_\mu \widehat{\varphi }_\nu \widehat{\varphi }+\frac{3g^2}{16}e^{\frac{8}{3}\widehat{\varphi }}\right)$$
(10)
$$F_{\mu \nu }^a=_\mu A_\nu ^a_\nu A_\mu ^a+f_{abc}A_\mu ^bA_\nu ^c$$
(11)
This is obtained by using the general formula of Scherk and Schwarz for reducing gravity from higher dimensions. For the antisymmetric tensor, we take the ansatz
$$\widehat{B}_{\mu \nu }=B_{\mu \nu },\widehat{B}_{\mu m}=\frac{1}{\sqrt{2}}A_\mu ^aU_m^a(z),\widehat{B}_{mn}=\stackrel{~}{B}_{mn}(z)$$
(12)
such that $`\widehat{H}_{mnp}=\frac{1}{2\sqrt{2}}f_{mnp}.`$ To evaluate $`\widehat{H}_{MNP}\widehat{H}^{MNP}`$ we first define
$$\widehat{H}_{ABC}=\widehat{e}_A^M\widehat{e}_B^N\widehat{e}_C^P\widehat{H}_{MNP}.$$
A direct substitution of the ansatz $`\left(\text{8}\right)`$ and $`\left(\text{12}\right)`$ gives
$`\widehat{H}_{\alpha \beta \gamma }`$ $`=`$ $`e^{\frac{5}{2}\widehat{\varphi }}e_\alpha ^\mu e_\beta ^\nu e_\gamma ^\rho H_{\mu \nu \rho }^{},`$
$`\widehat{H}_{\alpha \beta c}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}e^{\frac{7}{6}\widehat{\varphi }}e_\alpha ^\mu e_\beta ^\nu F_{\mu \nu }^a,`$
$`\widehat{H}_{\alpha bc}`$ $`=`$ $`0,`$
$`\widehat{H}_{abc}`$ $`=`$ $`{\displaystyle \frac{1}{2\sqrt{2}}}f_{abc}`$
where
$`H_{\mu \nu \rho }^{}`$ $`=`$ $`H_{\mu \nu \rho }\omega _{\mu \nu \rho },`$
$`\omega _{\mu \nu \rho }`$ $`=`$ $`6(A_{[\mu }^a_\nu A_{\rho ]}^a+{\displaystyle \frac{1}{3}}f_{abc}A_\mu ^aA_\nu ^bA_\rho ^c).`$ (13)
These results are a very strong consistency checks on the ansatz, especially in the form of $`H_{\mu \nu \rho }^{}`$ and $`\widehat{H}_{\alpha \beta c}.`$ We now have
$$_{5B}=e\left(\frac{1}{12}e^{\frac{16}{3}\widehat{\varphi }}H_{\mu \nu \rho }^{}H^{\mu \nu \rho }\frac{1}{8}e^{\frac{8\widehat{\varphi }}{3}}F_{\mu \nu }^aF^{\mu \nu a}\frac{g^2}{16}e^{\frac{8}{3}\widehat{\varphi }}\right).$$
(14)
The scalar part of $`D=10`$ supergravity, gives the following contribution to the five dimensional theory
$$_{5S}=\frac{e}{2}_\mu \widehat{\varphi }^\mu \widehat{\varphi }.$$
(15)
Therefore, combining all terms, the five dimensional theory is described by the Lagrangian
$$_5=e\left(\frac{1}{4}R\frac{1}{4}e^{\frac{8\widehat{\varphi }}{3}}F_{\mu \nu }^aF^{\mu \nu a}+\frac{4}{3}g^{\mu \nu }_\mu \widehat{\varphi }_\nu \widehat{\varphi }+\frac{1}{12}e^{\frac{16\widehat{\varphi }}{3}}H_{\mu \nu \rho }^{}H^{\mu \nu \rho }+\frac{g^2}{8}e^{\frac{8\widehat{\varphi }}{3}}\right).$$
(16)
Since the potential of the resulting theory depends only on one scalar field, a consistent truncation can be achieved by setting all gauge fields but one to zero. Therefore the index $`a`$ in the above Lagrangian will take one value. In order to compare with the standard Lagrangian, we multiply our Lagrangian by a factor of $`2`$. Therefore the resulting $`N=2`$ five dimensional theory is described by the Lagrangian
$$_5=e\left(\frac{1}{2}R\frac{1}{2}e^{\frac{8\widehat{\varphi }}{3}}F_{\mu \nu }^1F^{\mu \nu 1}+\frac{8}{3}g^{\mu \nu }_\mu \widehat{\varphi }_\nu \widehat{\varphi }+\frac{1}{6}e^{\frac{16\widehat{\varphi }}{3}}H_{\mu \nu \rho }^{}H^{\mu \nu \rho }+\frac{g^2}{4}e^{\frac{8\widehat{\varphi }}{3}}\right),$$
(17)
where $`\omega _{\mu \nu \rho }=6A_{[\mu }^1_\nu A_{\rho ]}^1`$ in $`H_{\mu \nu \rho }^{}.`$ We now apply a duality transformation by adding to the above Lagrangian the term
$$\frac{1}{3}ϵ^{\mu \nu \rho \sigma \kappa }_\mu A_\nu ^0H_{\rho \sigma \kappa }.$$
Integrating $`A_\nu ^0`$ forces $`H_{\mu \nu \rho }`$ to be of the form $`3_{[\mu }B_{\nu \rho ]}`$ giving the five dimensional theory with the $`B_{\mu \nu }`$ field. On the other hand by integrating the independent field $`H_{\mu \nu \rho }`$ in the path integral as it appears linearly and quadratically is equivalent to the substitution
$$H_{\mu \nu \rho }=\omega _{\mu \nu \rho }+e^{\frac{16}{3}\widehat{\varphi }}ϵ_{\mu \nu \rho }^{\sigma \kappa }_\sigma A_\kappa ^0.$$
This gives the dual Lagrangian
$`_5`$ $`=`$ $`e({\displaystyle \frac{1}{2}}R{\displaystyle \frac{1}{2}}e^{\frac{8\widehat{\varphi }}{3}}F_{\mu \nu }^1F^{\mu \nu 1}{\displaystyle \frac{1}{2}}e^{\frac{16\widehat{\varphi }}{3}}F_{\mu \nu }^0F^{\mu \nu 0}+{\displaystyle \frac{8}{3}}g^{\mu \nu }_\mu \widehat{\varphi }_\nu \widehat{\varphi }`$ (18)
$`+{\displaystyle \frac{e^1}{2}}ϵ^{\mu \nu \rho \sigma \lambda }F_{\mu \nu }^1F_{\rho \sigma }^1A_\lambda ^0+{\displaystyle \frac{g^2}{4}}e^{\frac{8\widehat{\varphi }}{3}}).`$
where $`0`$ and $`1`$ are to label the graviphoton and the additional vector multiplet gauge fields respectively. This is the gauged $`U(1)`$ five dimensional $`N=2`$ supergravity with one vector multiplet. It is the gauged version of the five dimensional theory initially introduced by Strominger and Vafa .
We close this section by noting that it is possible to obtain $`D=5`$ gauged supergravity from $`D=7`$ gauged supergravity. In a recent work we have derived $`D=7`$ gauged supergravity by compactifying $`D=10`$ on an $`SU(2)`$ group manifold . The Lagrangian in seven dimensions is given by
$$_7=e\left(\frac{1}{4}R\frac{1}{4}e^{\frac{8}{5}\widehat{\varphi }}F_{MN}^aF^{MNa}+\frac{4}{5}g^{MN}_M\widehat{\varphi }_N\widehat{\varphi }+\frac{1}{12}e^{\frac{16}{5}\widehat{\varphi }}H_{MNP}^{}H^{{}_{}{}^{}MNP}+\frac{g^2}{8}e^{\frac{8}{5}\widehat{\varphi }}\right).$$
This is reduced and truncated to $`D=5`$ gauged supergravity by taking the following ansatz
$$e_M^A=\left(\begin{array}{cc}e^{\frac{8}{15}\widehat{\varphi }}e_\mu ^\alpha & 0\\ 0& e^{\frac{4}{5}\widehat{\varphi }}\delta _m^i\end{array}\right),m=5,6.$$
as well as $`B_{\mu m}=0`$ and $`A_m^a=0`$. One can easily show that the reduced Lagrangian is given by $`\left(\text{16}\right)`$
## 3 <br>Embedding Into Very Special Geometry
The solutions of five dimensional $`N=2`$ supergravity theory with vector multiplets theory have been discussed within the framework of very special geometry . Therefore, before we discuss solutions of the five dimensional theory and their embedding into ten dimensional supergravity and M-theory, it is essential to consider our compactified Lagrangian in this framework.
A class of five-dimensional $`N=2`$ supergravity coupled to abelian vector supermultiplets can be obtained by compactifying eleven-dimensional supergravity, the low-energy theory of M-theory, on a Calabi-Yau three-folds . The massless spectrum of the theory contains $`(h_{(1,1)}1)`$ vector multiplets with real scalar components, and thus $`h_{(1,1)}`$ vector bosons (the additional vector boson is the graviphoton). The theory also contains $`h_{(2,1)}+1`$ hypermultiplets, where $`h_{(1,1)}`$ and $`h_{(2,1)},`$ are the Calabi-Yau Hodge numbers.
The bosonic part of the effective gauged supersymmetric $`N=2`$ Lagrangian which describes the coupling of vector multiplets to supergravity is given by
$$=e\left(\frac{1}{2}R+g^2V\frac{1}{4}G_{IJ}F_{\mu \nu }{}_{}{}^{I}F_{}^{\mu \nu J}\frac{1}{2}𝒢_{ij}_\mu \varphi ^i^\mu \varphi ^j+\frac{e^1}{48}ϵ^{\mu \nu \rho \sigma \lambda }C_{IJK}F_{\mu \nu }^IF_{\rho \sigma }^JA_\lambda ^K\right),$$
(19)
$`R`$ is the scalar curvature, $`F_{\mu \nu }^I`$ are the Abelian field-strength tensor, $`V`$ is the potential given by
$$V(X)=V_IV_J\left(6X^IX^J\frac{9}{2}𝒢^{ij}_iX^I_jX^J\right),$$
(20)
where $`X^I`$ represent the real scalar fields which have to satisfy the constraint
$$𝒱=\frac{1}{6}C_{IJK}X^IX^JX^K=1.$$
(21)
Also:
$$G_{IJ}=\frac{1}{2}_I_J\mathrm{log}𝒱|_{𝒱=1},𝒢_{ij}=_iX^I_jX^JG_{IJ}|_{𝒱=1},$$
(22)
where $`_i`$ refers to a partial derivative with respect to the scalar field $`\varphi ^i`$. The physical quantities in (19) can all be expressed in terms of the homogeneous cubic polynomial $`𝒱`$.
Further useful relations are
$$_iX_I=\frac{2}{3}G_{IJ}_iX^J,X_I=\frac{2}{3}G_{IJ}X^J.$$
(23)
It is worth pointing out that for Calabi-Yau compactification, $`𝒱`$ is the intersection form, $`X^I`$ and $`X_I=\frac{1}{6}C_{IJK}X^JX^K`$ correspond to the size of the two- and four-cycles and $`C_{IJK}`$ are the intersection numbers of the Calabi-Yau threefold.
A very useful relation of very special geometry is
$$𝒢^{ij}_jX^I_jX^J=G^{IJ}\frac{2}{3}X^IX^J.$$
(24)
The potential can also be written as
$$V(X)=9V_IV_J\left(X^IX^J\frac{1}{2}G^{IJ}\right).$$
(25)
The Lagrangian (18) correspond to the following identifications<sup>2</sup><sup>2</sup>2note that the sign difference of the kinetic terms is due to using a different metric signature
$$G_{00}=2e^{\frac{16\widehat{\varphi }}{3}},G_{11}=2e^{8\frac{\widehat{\varphi }}{3}},𝒢_{11}=\frac{16}{3},C_{011}=8.$$
In order to determine the $`X^I,`$ we use the relations (23) together with $`X_IX^I=1,`$ this gives $`X^0=c_1e^{8\frac{\widehat{\varphi }}{3}}`$, $`X^1=c_2e^{4\frac{\widehat{\varphi }}{3}}.`$ Upon using (24) and the above identification we get $`c_1=\frac{1}{2},`$ $`c_2=\frac{1}{\sqrt{2}}.`$ Finally, from the expression of the potential one obtains that $`V_0=0,`$ $`V_1=\frac{1}{3}.`$ Therefore we have
$$X^0=\frac{1}{2}e^{8\frac{\widehat{\varphi }}{3}},X^1=\frac{1}{\sqrt{2}}e^{4\frac{\widehat{\varphi }}{3}},X_0=\frac{2}{3}e^{8\frac{\widehat{\varphi }}{3}},X_1=\frac{2\sqrt{2}}{3}e^{4\frac{\widehat{\varphi }}{3}}.$$
## 4 Lifting Solutions to ten and eleven dimensions
Our previous results suggest that any solution of gauged supergravity in seven and five dimension given in terms of the metric, gauge field and scalar fields can be lifted to ten dimensions as a solution of $`N=1`$ ten dimensional supergravity. By also noting the relation between the ten and eleven dimensional theory, one can then lift all solutions to eleven dimensions.
To lift the solutions to ten dimensions we have to express the ten dimensional fields in terms of five dimensional ones. To start with we write the ten dimensional metric as
$`\widehat{g}_{\mu \nu }`$ $`=`$ $`e^{\frac{5}{3}\widehat{\varphi }}g_{\mu \nu }2e^{\widehat{\varphi }}A_\mu ^aA_\nu ^a,`$
$`\widehat{g}_{\mu m}`$ $`=`$ $`\sqrt{2}e^{\widehat{\varphi }}A_\mu ^aU_m^a\left(z\right),`$
$`\widehat{g}_{mn}`$ $`=`$ $`e^{\widehat{\varphi }}U_m^a\left(z\right)U_n^a\left(z\right).`$
The five dimensional model is given in terms of $`X^0,X^1,A_\mu ^0`$ and $`A_\mu ^1.`$ Since $`A_\mu ^0`$ is the dual of $`B_{\mu \nu }`$ we first evaluate
$$H_{\mu \nu \rho }^{}=H_{\mu \nu \rho }\omega _{\mu \nu \rho }=e^{\frac{16}{3}\widehat{\varphi }}ϵ_{\mu \nu \rho }^{\sigma \kappa }_\sigma A_\kappa ^0.$$
Or in terms of the ten dimensional fields
$`\widehat{H}_{\alpha \beta \gamma }`$ $`=`$ $`e^{\frac{5}{2}\widehat{\varphi }}e_\alpha ^\mu e_\beta ^\nu e_\gamma ^\rho H_{\mu \nu \rho }^{}=e^{\frac{17}{6}\widehat{\varphi }}ϵ_{\alpha \beta \gamma }^{\delta \eta }e_\delta ^\sigma e_\eta ^\kappa _\sigma A_\kappa ^0,`$
$`\widehat{H}_{\alpha \beta c}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\delta _c^1e^{\frac{7}{6}\widehat{\varphi }}e_\alpha ^\mu e_\beta ^\nu F_{\mu \nu }^1.`$
The dilaton field is determined from the solution to $`X^0`$ and $`X^1.`$
### 4.1 <br>Electrically charged solutions
The spherically symmetric BPS electric solutions as well as magnetic string solutions were obtained in by solving for the vanishing of the gravitino and gaugino supersymmetry variation for a particular choice of the supersymmetry parameter. These are given by
$`ds^2`$ $`=`$ $`𝒱^{4/3}(1+g^2r^2𝒱^2)dt^2+𝒱^{2/3}\left[{\displaystyle \frac{dr^2}{1+g^2r^2𝒱^2}}+r^2(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2+\mathrm{cos}^2\theta d\psi ^2)\right]`$
$`F_{tm}^I`$ $`=`$ $`_m(𝒱^1Y^I),`$
$`𝒱`$ $`=`$ $`{\displaystyle \frac{1}{6}}C_{IJK}Y^IY^JY^K,{\displaystyle \frac{1}{2}}C_{IJK}Y^JY^K=H_I=3V_I+{\displaystyle \frac{q_I}{r^2}}`$ (26)
where
$$Y^I=𝒱^{\frac{2}{3}}X^I,Y_I=𝒱^{\frac{1}{3}}X_I,𝒱=e^{3U}$$
Let us go back to the solutions of our five dimensional gauged theory with one vector multiplet. As for the electric solutions, we have the following equations
$`𝒱`$ $`=`$ $`{\displaystyle \frac{1}{2}}C_{011}Y^0Y^1Y^1=4Y^0(Y^1)^2,`$ (27)
$`{\displaystyle \frac{1}{2}}C_{011}(Y^1)^2=4(Y^1)^2=H_0={\displaystyle \frac{q_0}{r^2}},`$
$`C_{110}Y^1Y^0=8(Y^1)Y^0=H_1=1+{\displaystyle \frac{q_1}{r^2}}.`$
This gives
$`Y^0`$ $`=`$ $`{\displaystyle \frac{r}{4q_0^{\frac{1}{2}}}}\left(1+{\displaystyle \frac{q_1}{r^2}}\right),`$
$`Y^1`$ $`=`$ $`{\displaystyle \frac{q_0^{\frac{1}{2}}}{2r}},`$
$`𝒱`$ $`=`$ $`{\displaystyle \frac{q_0^{\frac{1}{2}}}{4r}}\left(1+{\displaystyle \frac{q_1}{r^2}}\right).`$
The metric depends only on $`𝒱`$. The gauge field strengths are given by
$`F_{tr}^0`$ $`=`$ $`_r(𝒱^1Y^0)={\displaystyle \frac{2r}{q_0}},`$
$`F_{tr}^1`$ $`=`$ $`_r(𝒱^1Y^1)={\displaystyle \frac{4q_1}{r^3\left(1+\frac{q_1}{r^2}\right)^2}}.`$ (28)
To lift our solution to ten dimensions first we write
$$Y^0=𝒱^{\frac{1}{3}}X^0,Y^1=𝒱^{\frac{1}{3}}X^1,$$
from which we deduce that
$$\frac{X^0}{X^1}=\frac{Y^0}{Y^1}=\frac{1}{\sqrt{2}}e^{4\widehat{\varphi }}=\frac{r^2}{2q^0}\left(1+\frac{q_1}{r^2}\right).$$
The gauge fields are
$$A_t^0=\frac{r^2}{q_0},A_t^1=\frac{2}{1+\frac{q_1}{r^2}}.$$
The ten dimensional metric is then
$`d\widehat{s}^2`$ $`=`$ $`e^{\frac{5}{3}\widehat{\varphi }}\left(𝒱^{4/3}(1+g^2r^2𝒱^2)dt^2+𝒱^{2/3}\left[{\displaystyle \frac{dr^2}{1+g^2r^2𝒱^2}}+r^2(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2+\mathrm{cos}^2\theta d\psi ^2)\right]\right)`$
$`+e^{\widehat{\varphi }}\left({\displaystyle \frac{q_0}{2r^2𝒱^2}}dt^2\right)+e^{\widehat{\varphi }}\left({\displaystyle \frac{\sqrt{2q_0}}{r𝒱}}U_m^1\left(z\right)dtdz^m\right)+e^{\widehat{\varphi }}U_m^a\left(z\right)U_n^a\left(z\right)dz^mdz^n,`$
where $`𝒱=\frac{q_0^{\frac{1}{2}}}{4r}\left(1+\frac{q_1}{r^2}\right).`$ The non vanishing components of the field strength of the antisymmetric tensor $`B_{MN}`$ are
$`\widehat{H}_{234}`$ $`=`$ $`e^{\frac{17}{6}\widehat{\varphi }}e_0^te_1^r_rA_t^0={\displaystyle \frac{q_0^{\frac{7}{8}}}{2^{\frac{5}{16}}r^{\frac{7}{4}}\left(1+\frac{q_1}{r^2}\right)^{\frac{3}{8}}}},`$
$`\widehat{H}_{015}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}e^{\frac{7}{6}\widehat{\varphi }}e_0^te_1^r_rA_t^1={\displaystyle \frac{2^{\frac{11}{16}}q_1}{q_0^{\frac{1}{8}}r^{\frac{11}{4}}\left(1+\frac{q_1}{r^2}\right)^{\frac{11}{8}}}},`$
$`\widehat{H}_{567}`$ $`=`$ $`{\displaystyle \frac{g}{2\sqrt{2}}}.`$
### 4.2 <br>Magnetic solutions
Here we will discuss the magnetic string solution found in . This is given by the metric
$`ds^2`$ $`=`$ $`(gr)^{\frac{1}{2}}e^{\frac{3U}{2}}(dt^2+dz^2)+e^{2U}dr^2+r^2\left(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2\right)`$
$`e^U`$ $`=`$ $`{\displaystyle \frac{1}{3gr}}+gr.`$ (29)
The gauge fields are given by $`A_\varphi ^I=q^I\mathrm{cos}\theta `$ and the scalar fields and the magnetic charges satisfy
$$X^IV_I=1,3gq^IV_I=1.$$
(30)
As for the magnetic solution of our model one finds that
$$X^0=\frac{1}{36},X_0=12,X^1=3,X_1=\frac{2}{9},e^{\frac{4\widehat{\varphi }}{3}}=3\sqrt{2}.$$
(31)
The above solution can be lifted to ten dimensions and we get
$`d\widehat{s}^2`$ $`=`$ $`\left(18\right)^{\frac{5}{8}}\left((gr)^{\frac{1}{2}}e^{\frac{3U}{2}}(dt^2+dz^2)+e^{2U}dr^2+r^2\left(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2\right)\right)`$
$`+2\left(18\right)^{\frac{1}{8}}q_1^2\mathrm{cos}^2\theta d\varphi ^22^{\frac{3}{4}}\left(18\right)^{\frac{1}{8}}q_1\mathrm{cos}\theta U_m^1\left(z\right)d\varphi dz^m`$
$`+\left(18\right)^{\frac{1}{8}}U_m^a\left(z\right)U_n^a\left(z\right)dz^mdz^n.`$
$$\widehat{H}_{012}=e^{\frac{17}{6}\widehat{\varphi }}e_3^\theta e_4^\varphi _\theta A_\varphi ^0=18^{\frac{17}{48}}\frac{q_0}{r^2}.$$
We note that the numerical factors for the solution can be absorbed by rescaling the coordinates and the charge $`q_1.`$ More solutions found in can also be lifted.
To lift these solutions further to eleven dimensions we have to first write the dimensional reduction of eleven dimensional supergravity to $`N=1`$ ten dimensional supergravity. These are
$$E_M^A=e^{\frac{1}{6}\widehat{\varphi }}\widehat{e}_M^A,E_{\stackrel{.}{11}}^{11}=e^{\frac{4}{3}\widehat{\varphi }},A_{MN\stackrel{.}{11}}=\widehat{B}_{MN}$$
and the eleven dimensional metric is related to the ten dimensional one by
$$ds^{2\left(eleven\right)}=e^{\frac{1}{3}\widehat{\varphi }}d\widehat{s}^2+e^{\frac{8}{3}\widehat{\varphi }}\left(dx^{11}\right)^2.$$
The non vanishing components of the antisymmetric tensor field-strengths are
$$F_{MNP\stackrel{.}{11}}=\widehat{H}_{MNP}.$$
The lifting of solutions from five to seven dimensions is very simple. We write
$$ds_7^2=e^{\frac{16}{15}\widehat{\varphi }}ds_5^2+e^{\frac{8}{5}\widehat{\varphi }}\left(\left(dx^5\right)^2+\left(dx^6\right)^2\right).$$
## 5 Conclusions
In this work we have shown that it is possible to obtain $`U(1)`$ gauged $`N=2`$ five dimensional supergravity interacting with one vector multiplet by compactifying and truncating ten dimensional supergravity on the group manifold $`SU(2)\times `$ $`U(1)^2`$. The model obtained is the gauged version of the supergravity model introduced by Strominger and Vafa. Using the relation between the higher dimensional fields and the lower ones, it becomes possible to lift known solutions such as black holes, string solutions and domain walls of the five dimensional theory to seven, ten and eleven dimensional supergravity theories to ten and eleven dimensions. Some known electric and magnetic solutions for our gauged $`D=5`$ supergravity compactified model, formulated in the framework of special geometry, were lifted to higher dimensions. Such solutions are not easy to find directly by studying the seven, ten and eleven dimensional supergravity theories. At this stage it would be useful to study some of the properties of these solutions and give their interpretation in terms of D-brane and M-theory dynamics.
|
warning/0003/gr-qc0003027.html
|
ar5iv
|
text
|
# Gravitational waves from black hole collisions via an eclectic approach
|
warning/0003/astro-ph0003315.html
|
ar5iv
|
text
|
# I Introduction
## I Introduction
The confirmed detection of cosmic rays above the Greisen-Zatsepin-Kuz’min cutoff gives confidence in the existence of neutrinos of energies reaching the EeV scale and above. Such neutrinos are expected both in models in which the protons are accelerated to the highest energies , such as in Active Galactic Nuclei or Gamma Ray Bursts and in “top bottom” scenarios in which cosmic rays are basically produced through quark fragmentation in events such as the decay of long lived heavy relic particles or the annihilation of topological defects . If the highest energy component of the cosmic rays are protons, as suggested by increasing experimental evidence , they are expected to produce neutrinos in their interactions with the cosmic microwave background . Neutrino detection would provide extremely valuable information on fundamental questions, both in astrophysics, such as the origin of the highest energy cosmic rays and in particle physics.
Detecting high energy neutrinos may be a reality in the immediate future as many efforts are being made to develop large scale Čerenkov detectors under water or ice , designed to challenge the low neutrino cross section exploiting the long range of the high energy muons produced in charged current muon neutrino interactions. For EeV neutrinos these detectors are also capable of detecting light from high energy showers produced by neutrinos of any flavor in both neutral and charged current interactions, but the effective acceptance of the detector is reduced because the shower must be produced very close or within the instrumented volume.
It has been known for long that the development of showers in dense media produces an excess charge which generates a coherent Čerenkov pulse in the radiowave frequency when it propagates through the medium . The detection of these pulses provides a possible alternative to neutrino detection particularly appropriate for very high energies because the signal scales with the square of the primary energy . The method is attractive because of the good transmission properties of large natural volumes of ice and sand and because much information about the charge distribution in the shower is preserved in the frequency and angular distribution of the pulses. This last property can be used to extract information about shower energy and neutrino flavor . The technique faces a number of technical difficulties however and several attempts are currently being made to test the theoretical predictions and to study the feasibility of the technique in Antarctic ice .
Theoretical calculations are also difficult because a complete interference calculation calls for simulations capable of following electrons and positrons to the Čerenkov threshold ($`100`$keV). For high energy showers this is unfortunately out of question because of the large number of particles involved and approximations have been specifically deviced to study the radio emission of high energy showers in ice. The calculation of radio pulses from EeV showers has been possible in the one dimensional (1-D) approximation which consists on neglecting both the lateral distribution and the subluminal velocity of shower particles . All the calculations of radio pulses have been made so far in the Fraunhofer limit. In this limit the dependence of the electric field on distance to shower is trivial and the characterization of the angular distribution of the radio pulse at a given frequency is effectively only dependent on one variable, namely the angle between the shower axis and the observation direction what simplifies the simulations . Clearly Fresnel type interference will take place if the showers are close enough to the detectors, but the calculation of these effects becomes even more time consuming.
In this paper we firstly give a brief introduction to coherent radio emission in Section II (fuller details can be found in Refs. ) accounting for the approximations made. In Section III we make extensive tests and explore the validity of the 1-D approximation in the Fraunhofer limit by direct comparison with complete simulations, and we discuss the approximation pointing out the connections between the radio emission and shower fluctuations, what gives new and useful insight into the radioemission processes. In Section IV we use the 1-D approximation without taking the Fraunhofer limit to study the radiopulse as a function of the distance to observation point. In Section V we summarize and conclude, commenting on the implications of our results for neutrino detection.
## II Čerenkov Radio Pulses
When a charged particle travels through a dielectric medium of refraction index $`n`$ with speed $`\beta c`$ greater than the phase velocity of light in that medium ($`c/n`$), then Čerenkov radiation is emitted in a frequency band over which the $`\beta n>1`$ condition is satisfied without large absorbtion. The calculation of the Čerenkov electric field associated to the particle is a problem of classical electromagnetism that has been addressed elsewhere . Solving the inhomogeneous Maxwell’s equations in the transverse gauge, it is easy to obtain the Fourier components of the electric field produced by a current density $`\stackrel{}{J}(\stackrel{}{x}^{},t^{})`$:
$$\stackrel{}{E}(\stackrel{}{x},\omega )=\frac{e\mu _\mathrm{r}}{2\pi ϵ_0\mathrm{c}^2}i\omega 𝑑t^{}d^3\stackrel{}{x}^{}\frac{e^{i\omega t^{}+ik|\stackrel{}{x}\stackrel{}{x}^{}|}}{|\stackrel{}{x}\stackrel{}{x}^{}|}\stackrel{}{J}_{}(\stackrel{}{x}^{},t^{})$$
(1)
where $`\stackrel{}{J}_{}(\stackrel{}{x}^{},t^{})`$ is the component of the current transverse to the direction of observation $`\stackrel{}{x}`$. Also $`\nu `$ ($`\omega `$) is the frequency (angular frequency), $`k`$ is the modulus of the wave vector $`\stackrel{}{k}`$, $`\mu _r`$ is the relative permeability of the medium and $`ϵ_0`$ and $`c`$ is the permittivity and velocity of light in the vacuum.
A powerful approach to the simulation problem can be obtained neglecting the lateral distributions in shower particles and assuming all particles move at constant speed $`c`$ in one dimension. We obtain a useful compact expression relating the charge distribution of the shower and its associated electric field. Crude as it may look, this approximation (1-D approximation in brief) will be shown to give very good results particularly around the Čerenkov angle and it has allowed the possibility of establishing the radioemission from EeV showers . The method naturally relates different features of shower development to the spectrum and angular distribution of the radio emission in an interesting way, giving insight into the complexity of the calculated angular pulses. For simplicity we are going to take $`\stackrel{}{x}^{}=\stackrel{}{z}{}_{}{}^{}=z^{}\widehat{n}_z^{}`$ where $`\widehat{n}_z^{}`$ is a unitary vector along the shower axis. The current associated to the shower development in this approximation is then given by:
$$\stackrel{}{J}_{}(\stackrel{}{z}{}_{}{}^{},t^{})=Q(z^{})\stackrel{}{c}_{}\delta ^3(\stackrel{}{z}{}_{}{}^{}\stackrel{}{c}t^{})$$
(2)
where $`Q(z^{})`$ is the longitudinal development of the excess charge in the shower. The substitution of this current into Eq. 1 leads to:
$$\stackrel{}{E}(\stackrel{}{x},\omega )=\frac{e\mu _\mathrm{r}}{2\pi ϵ_0\mathrm{c}^2}i\omega \mathrm{sin}\theta \widehat{n}_{}𝑑z^{}Q(z^{})\frac{e^{i\frac{\omega }{c}z^{}+ik|\stackrel{}{x}z^{}\widehat{n}_z^{}|}}{|\stackrel{}{x}z^{}\widehat{n}_z^{}|}$$
(3)
where $`\theta `$ is the angle between the shower axis and the direction of observation $`\stackrel{}{x}`$ and $`\widehat{n}_{}`$ is a unitary vector perpendicular to $`\stackrel{}{x}`$.
We can use this expression to obtain the Čerenkov electric field emitted by a particle shower propagating along a medium. Eq. 3 accounts for the correct phase factors and distances for showers that are close to the observer (Fresnel region). In the Fraunhofer limit the phase factor in Eq. 3 can be approximated by $`ik|\stackrel{}{x}\stackrel{}{z}{}_{}{}^{}|ikRi\stackrel{}{k}\stackrel{}{z}^{}`$, where $`R=|\stackrel{}{x}|`$ is the distance from the center of the shower to the observation point. It corresponds to the condition that observation distance $`R`$ exceeds the Fresnel distance $`R_F=\pi n\nu (L_s\mathrm{sin}\theta /2)^2/c`$, where $`L_s`$ is the typical length of the shower. In this limit it is straightforward to show that the electric field emitted by a shower in the 1-D approximation can be related to the Fourier transform of the longitudinal charge distribution:
$$\stackrel{}{E}(\omega ,\stackrel{}{\mathrm{x}})=\frac{e\mu _\mathrm{r}}{2\pi ϵ_0\mathrm{c}^2}i\omega \mathrm{sin}\theta \frac{\mathrm{e}^{ikR}}{R}\widehat{n}_{}𝑑z^{}Q(z^{})\mathrm{e}^{ipz^{}}$$
(4)
where we have introduced for convenience the parameter $`p(\theta ,\omega )=(1n\mathrm{cos}\theta )\omega /c`$ in Eq. 4 to stress the connection between the radio emission spectrum and the Fourier transform of the (excess) charge distribution. This allows a simple analogy to the classical diffraction pattern of an aperture function and helps understanding many of the complex features of the results obtained by simulation.
For the case of a single particle moving between two fixed points this expression (replacing $`c`$ by an arbitrary particle velocity $`v`$) reproduces the formula obtained in :
$$\stackrel{}{E}(\omega ,\stackrel{}{\mathrm{x}})=\frac{e\mu _\mathrm{r}i\omega }{2\pi ϵ_0\mathrm{c}^2}\frac{\mathrm{e}^{ikR}}{R}\stackrel{}{v}_{}\left[\frac{\mathrm{e}^{i(\omega \stackrel{}{k}\stackrel{}{v})\mathrm{t}_2}\mathrm{e}^{i(\omega \stackrel{}{k}\stackrel{}{v})\mathrm{t}_1}}{i(\omega \stackrel{}{k}\stackrel{}{v})}\right]$$
(5)
where $`\stackrel{}{v}_{}`$ refers to the particle’s velocity projected in a plane perpendicular to the observing direction and $`\mathrm{t}_2`$ $`(\mathrm{t}_1)`$ is the time corresponding to the final (initial) point of the track. This is the basic expression used for the numerical simulation of radio pulses from individual tracks (see appendix A).
## III The one-dimensional approach
We will firstly explore the validity of the 1-D approximation by direct comparison with simulation results in three dimensions. The program we use for the full simulation of electromagnetic showers in homogeneous ice, is described in Ref. . The results of the simulation will be compared to those obtained using Eq. 4 with different curves for the excess charge development function $`Q(z)`$ what will turn out to be quite illustrative.
Fig. 1 compares the angular distributions of the pulses for showers initiated by different energy electrons using the full simulation and using Eq. 4 with $`Q(z)`$ directly from the excess charge depth distribution as obtained in the same simulations. Fig. 2 displays the frequency spectra at different observation angles for a 10 TeV shower again for both approaches. Several conclusions can be drawn from these graphs with respect to the validity of the 1-D approximation. Clearly the electric field amplitude around the Čerenkov cone is well reproduced in shape by the 1-D approximation except in the Čerenkov direction where the approximation overestimates the amplitude by a factor that increases with frequency. Below 100 MHz the effect is negligible becoming of order $`20\%`$ (a factor of 2) for 300 MHz (1 GHz). The angular interval over which the approximation is valid slowly increases with shower energy and scales with the inverse of the frequency. Well outside the Čerenkov cone no agreement can be claimed but the order of magnitude of the approximation agrees with the simulation.
For completitude we give a new parameterization for the frequency spectrum in the Čerenkov direction using a finer subdivision of individual electron tracks (approximation $`a`$ see appendix A), which represents slight increase at frequencies above 500 MHz from that given in Ref. :
$$R|\stackrel{}{E}(\omega ,R,\theta _C)|2.53\times 10^7\left[\frac{E_{\mathrm{em}}}{1\mathrm{TeV}}\right]\left[\frac{\nu }{\nu _0}\right]\left[\frac{1}{1+(\nu /\nu _0)^{1.44}}\right]\mathrm{V}\mathrm{MHz}^1$$
(6)
where $`\nu _0=1.15`$ GHz. This parameterization is valid to frequencies below $`5`$GHz.
It is worth discussing the interpretation of the behavior of this approximation before we attempt to understand its validity in more complicated showers such as those having strong LPM effects . In the Čerenkov direction, corresponding to $`p=0`$, the agreement between the approximation and the full simulation is excellent for frequencies below about 100 MHz. This corresponds to complete constructive interference characterized by a spectrum that increases linearly with frequency as shown in Fig. 2. Above 100 MHz the simulated frequency spectrum deviates from linear behavior because the wavelength becomes comparable to the transverse deviation of shower particles and to a lesser extent because of time delays <sup>*</sup><sup>*</sup>*It has been checked by direct simulation that the time delays only become important for frequencies in the 10 GHz range at the Čerenkov direction.. Both these effects are ignored in the 1-D approximation that keeps on rising linearly. Away from the Čerenkov cone the approximation becomes valid even to higher frequencies. This is because destructive interference is in this case due to the longitudinal excess charge distribution which is correctly taken into account by the approximation.
In spite of the approximation overestimating the amplitude of the electric field in the Čerenkov direction for frequencies above $`100`$MHz, an ad-hoc correction can be implemented based on the shape of the frequency spectrum as obtained in the simulations. Since this effect is due to the lateral distribution of the electromagnetic component of the showers, it can be corrected with a unique function for each frequency. We have explicitly checked that the lateral distribution of electromagnetic showers is similar for showers with and without the LPM effect .
We have calculated the difference between the 1-D approximation and the full simulation in the Čerenkov direction as a function of frequency what is shown in Fig. 3 for two different shower energies. Note that the difference is (up to a factor that scales with shower energy) the same for showers of different energies. For this calculation we have actually improved the simulation by splitting the individual tracks in small subintervals (approximation $`c`$, see appendix A). Also shown is the calculation without track subdivisions (approximation $`a`$) for comparison. The angular behavior of the correction at a particular frequency can be also shown to be fairly independent of energy.
The needed correction basically consists of rescaling the pulse just in the region around the Čerenkov direction. It can be achieved for instance dividing the result of Eq. 4 by a gaussian correction factor:
$$1+\left[\frac{1DFS}{FS}\right]\mathrm{e}^{\frac{1}{2}\left[\frac{\theta \theta _C}{\sigma _\theta }\right]^2}$$
(7)
The expression in brackets symbolically represents the relative difference between the frequency spectra as given by the 1-D approximation ($`1D`$) and the full simulation ($`FS`$) calculated in the Čerenkov direction. It simply sets the scale of the correction. The numerator is shown in Fig. 3 for two test cases, showing that it also scales with energy at least in the energy interval checked. The half width of the gaussian term is approximately given by:
$$\sigma _\theta =2.2^\mathrm{o}\left[\frac{1\mathrm{GHz}}{\nu }\right]$$
(8)
For frequencies above the 100 MHz scale and high energies when the full simulation is not viable, one would implement the correction taking Eq. 6 instead of the full simulated result.
The 1-D approximation also works for complicated showers such as those initiated by electrons and photons of EeV energies with strong LPM effects . This has been explicitly checked by artificially lowering $`E_{\mathrm{LPM}}`$, the onset energy for LPM effects, so that showers with energies that allow full three dimensional simulations display the characteristic LPM elongations . The agreement between the full simulation and the 1-D approximations is illustrated in Fig. 4 and it is clear that it is not limited to the central peak but also applies to the secondary peaks that appear in the angular distribution of the radiated pulse. The above correction prescription also works for these fictious elongated showers with a mild reduction in precision.
Lastly, the simulation of the excess charge in an EeV shower can also be extremely time consuming because particles have to be followed at least to MeV energies when the interactions responsible for the excess charge become dominant over pair production and bremsstrahlung . According to simulations the pulse scales with the excess tracklength and this is practically only due to an excess of MeV electrons. The excess number of electrons can be approximately obtained by rescaling the total number of electrons and positrons in a shower by the fraction of excess and total tracklengths. This factor is very stable and has a value of $`25\%`$ in ice This value corrects the previous conservative estimates used in that quoted instead the ratio of excess projected tracklength to total tracklength as the relevant number (21$`\%`$).. As convenient parameterizations of the number of electrons and photons in showers are readily available it is possible to calculate shower size distributions for very large showers using them . In spite of the small gradual rise in the excess charge as the shower develops shown by simulations , the effects of this approximation are mild, a slight narrowing of the pulse which is negligible compared to the other approximations made (see Fig. 5).
Finally it is remarkably fortunate that neglecting lateral distributions and time delays is a very good way of approaching the problem if some considerations are cautiously taken into account, namely:
* Take the Fourier transform of the longitudinal distribution of the excess charge $`Q(z)`$ (or one fourth of the total number of electrons and positrons if $`Q(z)`$ is not available) as given by Eq. 4.
* For frequencies above 100 MHz divide the 1-D approximation by a correction factor as indicated by Eq. 7 taking Eq. 6 instead of the full simulation ($`FS`$) value.
### A Discussion: The relation between radio pulses and shower fluctuations
In the 1-D approximation the Fourier transform for $`p=0`$ becomes the integral of $`Q(z)`$, i.e. the excess tracklength. Simulations have shown that the excess tracklength scales extremely well with electromagnetic energy in the shower ($`E_{\mathrm{em}}`$) for both electromagnetic and hadronic showers up to energies exceeding 100 EeV with small fluctuations:
$$t=6400\left[\frac{E_{\mathrm{em}}}{1\mathrm{TeV}}\right]\mathrm{m}$$
(9)
Incidentally this nice property of the excess charge together with the fact that the radioemission in the Čerenkov direction is proportional to the excess tracklength, make such measurements excellent candidates for electromagnetic energy estimators. The breaking of the approximation at high frequencies is telling us that the lateral distribution is playing a significant role.
A simple limit of the 1-D approximation is obtained by taking an analytical expression for $`Q(z)`$ such as Greisen’s parameterization for the average development of an electromagnetic shower . The result just gives the radiation in the Čerenkov cone but no radiation outside, just like the Fourier transform of a gaussian.
Invoking superposition we can subtract from a given shower development curve a smooth Greisen-like curve having the same tracklength. The result displays the “roughness” of the depth development curve and we shall refer to it as the difference function. The electromagnetic pulse is the sum of an isolated Čerenkov peak due to the Greisen-like curve and an extra contribution from the Fourier spectrum of the difference function which precisely vanishes at the Čerenkov direction because it does not contribute to the total tracklength. Moreoever, for ordinary showers the amplitude of the difference function becomes smaller relative to shower size as the shower energy increases. This is just an statistical effect of having a larger number of particles and it indicates that the “spatial correlations” The name stresses the fact that they are different from standard fluctuations in shower theory because they refer to variations in shower size for the same shower at different positions rather than comparing shower size at the same spatial position for different showers. contained in the difference function must be related to fluctuations in shower size.
The fact that the magnitude of the difference function becomes smaller relative to shower size as the shower energy $`E_0`$ increases has the effect of ”illuminating” the Čerenkov cone much more sharply with respect to directions well outside the Čerenkov direction. This effect can also be understood in terms of coherence. In the Čerenkov direction of greatest coherence the electric field amplitude scales with the shower energy $`E_0`$, but when the radiation is incoherent, i.e. well outside the Čerenkov direction, the electric field should add incoherently and hence scale with $`\sqrt{E_0}`$. This roughly agrees with simulations and nicely connects the properties of the radio emission to spatial correlations in shower development.
For LPM showers the structure of the pulse outside the central narrow peak is still dominated by the longitudinal development of these showers because the amplitude of the difference function is much larger than for a conventional shower. This is because LPM showers fluctuate a great deal. (One can picture a characteristic LPM shower as a superposition of smaller subshowers with typical smooth profiles with random starting points along the shower length.)
In other words, the Fourier modes of the excess charge distribution are probed by the electric field at a given value of $`p`$ and hence at a given value of $`\theta `$ for a fixed frequency. The scale of the correlations in the distribution (the “wavelength” of the corresponding mode) is inversely proportional to $`p`$. As long as the scale of these correlations is larger than the characteristic lateral structure of the shower, the 1-D approximation is expected to work. This is precisely what happens for the LPM fluctuations.
In summary there are two angular regions for the electromagnetic pulse with a not very well defined boundary. One angular region corresponds to the surroundings of the Čerenkov cone where the 1-D approximation has powerful predictive power when one accounts for the correction described above. There is another region well outside the Čerenkov cone in which the calculated electric field amplitude drops considerably and behaves erratically, as some kind of “white noise” corresponding to the incoherent regime. In this region the short scale correlations of the excess track distribution are being probed and here the predictive power is lost with the approximations discussed. To calculate the radiopulse in such regions one needs three dimensional simulation programs which must sample tracks in small subintervals and which must follow all particles to the 100 keV region (approximation $`c`$ described in appendix A). All these requirements make it impossible with current computing power to simulate beyond 100 TeV. However the region outside the Čerenkov cone, having much reduced radioemission, is not very relevant for shower detection.
## IV The validity of the Fraunhofer approximation
All the calculations made of radio pulses have been made in the Fraunhofer approximation which corresponds to the limit:
$$R>R_F=3\mathrm{m}\left[\frac{L_s}{1\mathrm{m}}\right]^2\left[\frac{\nu }{1\mathrm{GHz}}\right]$$
(10)
Taking $`L_s3.8`$m, corresponding to the nominal ten radiation lengths in ice of a electromagnetic (hadronic) shower below about 10 PeV (10 EeV), a frequency of 1 GHz and $`\theta `$ equal to the Čerenkov angle then $`R_F45`$m. This distance is to be compared with the km scale set by the small absorption coefficient of radio waves in cold ice which tells us that the Fraunhofer condition is clearly satisfied. For very long showers (such as those that display a very strong LPM effect) and high frequencies, $`R_F`$ exceeds the typical attenuation scale. As the distance $`R`$ is reduced to values below $`R_F`$, the diffraction pattern gradually turns into a Fresnel pattern in which the angular features become blurred.
It is possible to calculate diffraction patterns for such showers with the typical restrictions that apply to these simulations. A full calculation is again not viable for the shower energies at which this effect becomes important at km scale distances. We have calculated the radio pulses as observed at distances in which the Fraunhoffer approximation breaks down, using simulated electron showers of different energies. We apply Eq. 3 for calculating electric field amplitudes at distances of order the Fresnel distance $`R_F`$, (a one dimensional transform that does not take the Fraunhofer limit). We calculate the effects for a range of energies and observation distances to specify the conditions under which the properties of the emission in the Fraunhofer limit are still valid.
In Fig. 6 we display the Čerenkov peak structure at 100 MHz for a range of distances around the Fraunhofer limit for a 1 EeV electromagnetic shower spanning 135 radiation lengths. We define the distance in relation to the center of charge of the shower. The calculated pattern has a reduced amplitude at the peak and becomes broader as expected. The Fraunhofer approximation is good to better than $`10\%`$ in absolute value for distances above $`400`$m. For a 100 EeV (10 PeV) shower the distance increases to 5 km (decreases to 20 m) for a roughly similar accuracy. The angular width of the pulse in the near field case increases with respect to the Fraunhofer case roughly by $`20\%`$ when the amplitude reduces by $`10\%`$.
In Fig. 7 we plot the ratio of the calculated and Fraunhofer amplitudes at the Čerenkov peak as a function of distance to the shower for different frequencies and shower energies. Also indicated are the absorption lengths at three different temperatures for reference. This graph sumarizes the results, for 1 km distance and energies above a few hundred PeV, Fresnel effects will become a serious concern for GHz frequencies. Provided that the distance to the shower and its direction can be determined, Fresnel effects could be corrected for, but this would clearly complicate and limit the analysis. This suggests that lower frequencies in the 100 MHz or even below may be appropriate for EeV showers. For hadronic type showers however no effects are foreseen for energies up to the 10 EeV range except for few abnormally long showers that are unlikely to happen .
## V Summary and Conclusions
We have shown that the calculation of coherent Čerenkov radio pulses from high energy showers in ice in the Fraunhoffer limit can be well approximated by neglecting the lateral distributions of the particles assuming that they travel at constant speed ($`c`$). The electric field amplitude simply becomes the one dimensional Fourier transform of the excess charge depth distribution. For the most relevant region around the Čerenkov direction, the approximation is correct for frequencies below 100 MHz. At higher frequencies the approximation is still relatively good but systematically overestimates the pulse in the Čerenkov direction. We have shown that the model can be made to agree at least up to 1 GHz by subtracting a simple ad-hoc gaussian correction that is proportional to the shower energy and otherwise only dependent on frequency. We have reported the relevant parameters for the correction and have presented an improved parameterization for the electric field amplitude in the Čerenkov direction.
We have also shown that instead of the actual charge excess distribution one can use the shower size longitudinal development curve which is more conventional than the excess charge, scaling the amplitude of the central peak by the excess tracklength fraction $`0.25`$.
We have developped a similar approximation for the region in which the Fraunhofer limit ceases to be valid. We have finally studied the behavior of the radiopulses of long electromagnetic showers in this region. Our results are again suggesting to use low frequencies for EeV showers as concluded in Ref. These frequencies have a number of advantages because they are less attenuated, they allow observation of the angular structure with less detectors, and they have milder Fresnel effects at a given distance. Because of Fresnel corrections, the possibility of extracting the mixed character of electron neutrino interactions suggested in requires frequencies below 100 MHz if the electron initiated subshower exceeds about 10 EeV.
Lowering the frequency implies a higher threshold for detection because the Čerenkov spectrum increases with frequency but for EeV showers this should not be a problem. It has been estimated that the threshold for detecting showers at 1 km distance with 1 GHz broadband antennas is in the 10 PeV range . Since the signal to noise roughly scales with the square root of the bandwith which directly relates to the central frequency, a factor of 100 reduction in frequency will only call for about a factor of 10 enhancement of the threshold still giving a very large signal to noise ratio for EeV showers.
Although our tests of the 1-D approximation rely heavily on a specific simulation program , our claim on the validity of the 1-D approximation is model independent. For testing purposes we used the charge excess distribution and the emitted radiopulses as obtained by the same routine. Numerically our results only apply for ice but it is only natural to expect that the same procedures can be applied to calculate the radiation in other materials.
Acknowledgements: We thank P. Gorham for many early discussions about Fresnel corrections and D. Besson, D.W. McKay, J.P. Ralston, S. Razzaque, D. Seckel and S. Seunarine for constructive criticism of the Montecarlo and many discussions. This work was supported in part by CICYT (AEN99-0589-C02-02) and by Xunta de Galicia (XUGA-20602B98). J. A. thanks the Department of Physics, University of Wisconsin, Madison and the Fundación Caixa Galicia for financial support. E. Z. thanks the Department of Physics, University of Wisconsin, Madison, where this work was finished for its hospitality, and the Xunta de Galicia for partially supporting this trip.
APPENDIX A: The ZHS Montecarlo
The simulation program used described in is a specifically deviced program for calculating radio-pulses from electromagnetic showers that follows particles to $`100`$ keV, taking into account low energy processes and timing. The depth development results have been compared to analytical parameterizations given in the Particle Data Book , with which they agree to a few percent.
The calculation of the radio emission uses Eq. 5 for electron and positron tracks. Several approximations can be made according to different choices in the subdivision of the individual charged particle tracks. In Ref. three different choices, named approximations $`a`$, $`b`$, and $`c`$ have been compared, testing for convergence as the subtracks become smaller.
Approximation $`a`$ is the standard that has been used in Refs. . It corresponds to taking the end points of all the tracks, and it just uses the average velocity for the corresponding effective track in Eq. 5. This is the standard reference calculation used throughout in this article except for Fig. 3. Note that this approximation gives the correct result provided the particle velocity is constant along the track.
Approximation $`b`$ subdivides the electron tracks according the different interaction points found along the track, (multiscattering is not considered as an interaction here). This approximation subdivides the track in finer subintervals as the energy becomes smaller, because the low energy electron scattering cross sections exceed bremsstrahlung and pair production. For each subtrack the average velocity is calculated between the corresponding end points of the track. Finally approximation $`c`$ subdivides each interaction according to a convenient algorithm for spliting the propagation of particles designed to better calculate the multiple scattering at low energies.
The three approximations are compared in Fig. 8 illustrating the convergence of the method and how the approximation $`a`$ is valid in the Čerenkov cone to a precision better than about $`10\%`$ for frequencies below 1 GHz. Full simulations in approximation $`c`$ are much more time consuming and have to be done for shower energies below $`100`$TeV. At low energies fluctuations from shower to shower are more important so that these tests are inevitably subject to larger uncertainties because of such fluctuations.
APPENDIX B: The gaussian approximation
For electromagnetic (hadronic) showers below 10 PeV (10 EeV), that is having no important deviations from Greisen behavior, the electric fied around the Čerenkov cone can be accurately determined with a gaussian approximation. The precise width of the cone inversely relates to the width (in $`z`$) of the excess charge depth distribution, $`Q(z)`$. As $`p`$ is directly related to the observation angle $`\theta `$ with an expression that involves the frequency as an overall factor, the width of the angular distribution of the ”central peak” becomes inversely proportional to $`\omega `$.
For small deviations from the Čerenkov angle ($`\mathrm{\Delta }\theta `$) the expression for $`p`$ to first order is :
$$p=\frac{\omega }{\mathrm{c}}\sqrt{n^21}\mathrm{\Delta }\theta +O(\mathrm{\Delta }\theta ^2)30.8\left[\frac{\nu }{1\mathrm{GHz}}\right]\mathrm{\Delta }\theta (\mathrm{m}^1)$$
(11)
The numerical value given in this expression corresponds to showers in ice with $`n=1.78`$. Defining the gaussian width by the points in which the amplitude drops by a factor $`\sqrt{e}`$ a gaussian of half-width $`\sigma _z`$ transforms to another gaussian of half-width $`\sigma _p=(\sigma _z)^1`$. We can fit a gaussian to the excess charge depth development curve identifying the shower length by the width $`l=2\sigma _z`$ and the angular full width of the radiopulse is then:
$$\sigma _\theta 3.72^{}\left[\frac{1\mathrm{GHz}}{\nu }\right]\left[\frac{1\mathrm{m}}{l}\right]$$
(12)
using approximation given by Eq. 11. For a typical shower length of 8 radiation lengths ($`3.1`$m in ice) the angular width of the pulse is about $`1^{}`$ at 1 GHz, in agreement with Ref. .
|
warning/0003/cs0003046.html
|
ar5iv
|
text
|
# Linear Tabulated Resolution Based on Prolog Control Strategy
## 1 Introduction
While Prolog has many distinct advantages, it suffers from some serious problems, among the best-known of which are infinite loops and redundant computations. Infinite loops cause users (especially less skilled users) to lose confidence in writing terminating Prolog programs, whereas redundant computations greatly reduce the efficiency of Prolog. The existing approaches to resolving these problems can be classified into two categories: loop checking and tabling.
Loop checking is a direct way to cut infinite loops. It locates nodes at which SLD-derivations step into a loop and prunes them from SLD-trees. Informally, an SLD-derivation $`G_0_{C_1,\theta _1}G_1\mathrm{}`$ $`_{C_i,\theta _i}G_i\mathrm{}`$ $`_{C_k,\theta _k}G_k\mathrm{}`$ is said to step into a loop at a node $`N_k`$ labeled with a goal $`G_k`$ if there is a node $`N_i`$ ($`0i<k`$) labeled with a goal $`G_i`$ in the derivation such that $`G_i`$ and $`G_k`$ are sufficiently similar. Many loop checking mechanisms have been presented in the literature (e.g. ). However, no loop checking mechanism can be both (weakly) sound and complete because the loop checking problem itself is undecidable in general even for function-free logic programs .
The main idea of tabling is that during top-down query evaluation, we store intermediate results of some subgoals and look them up to solve variants of the subgoals that occur later. Since no variant subgoals will be recomputed by applying the same set of program clauses, infinite loops can be avoided. As a result, termination can be guaranteed for bounded-term-size programs and redundant computations substantially reduced .
There are many ways to formulate tabling, each leading to a tabulated resolution (e.g. OLDT-resolution , SLG-resolution , Tabulated SLS-resolution , etc.). However, although existing tabulated resolutions differ in one aspect or another, all of them rely on the so called solution-lookup mode. That is, all nodes in a search tree/forest are partitioned into two subsets, solution nodes and lookup nodes; solution nodes produce child nodes using program clauses, whereas lookup nodes produce child nodes using answers in tables.
Our investigation shows that the principal disadvantage of the solution-lookup mode is that it makes tabulated resolutions non-linear. Let $`G_0_{C_1,\theta _1}G_1`$ $`\mathrm{}_{C_i,\theta _i}G_i`$ be the current derivation with $`G_i`$ being the latest generated goal. A tabulated resolution is said to be linear<sup>1</sup><sup>1</sup>1 The concept of “linear” here is different from the one used for SL-resolution . if it makes the next derivation step either by expanding $`G_i`$ by resolving a subgoal in $`G_i`$ against a program clause or a tabled answer, which yields $`G_i_{C_{i+1},\theta _{i+1}}G_{i+1}`$, or by expanding $`G_{i1}`$ via backtracking. It is due to such non-linearity that the underlying tabulated resolutions cannot be implemented in the same way as SLD-resolution (Prolog) using a simple stack-based memory structure. Moreover, some strictly sequential operators such as cuts ($`!`$) may not be handled as easily as in Prolog. For instance, in the well-known tabulated resolution system XSB, clauses like
$`p(.)\mathrm{},t(.),!,\mathrm{}`$
where $`t(.)`$ is a tabled subgoal, are not allowed because the tabled predicate $`t`$ occurs in the scope of a cut .
The objective of our research is to establish a hybrid approach to resolving infinite loops and redundant computations and develop a linear tabulated Prolog system. In this paper, we establish a theoretical framework for such a system, focusing on a linear tabulated resolution $``$ TP-resolution for positive logic programs (TP for Tabulated Prolog).
###### Remark 1.1
In this paper we will use the prefix TP to name some key concepts such as TP-strategy, TP-tree, TP-derivation and TP-resolution, in contrast to the standard Prolog control strategy, Prolog-tree (i.e. SLD-tree generated under Prolog-strategy), Prolog-derivation and Prolog-resolution (i.e. SLD-resolution controlled by Prolog-strategy), respectively.
In TP-resolution, each node in a search tree can act not only as a solution node but also as a lookup node, regardless of when and where it is generated. In fact, we do not distinguish between solution and lookup nodes in TP-resolution. This shows an essential difference from existing tabulated resolutions using the solution-lookup mode. The main idea is as follows: for any selected tabled subgoal $`A`$ at a node $`N_i`$ labeled with a goal $`G_i`$, it always first uses an answer $`I`$ in a table to generate a child node $`N_{i+1}`$ ($`N_i`$ acts as a lookup node), which is labeled by the resolvent of $`G_i`$ and $`I`$; if no new answers are available in the table, it resolves against program clauses to produce child nodes ($`N_i`$ then acts as a solution node). The order in which answers in a table are used is based on first-generated-first-use and the order in which program clauses are applied is from top to bottom except for the case where the derivation steps into a loop at $`N_i`$. In such a case, the subgoal $`A`$ skips the clause that is being used by its closest ancestor subgoal that is a variant of $`A`$. Like OLDT-resolution, TP-resolution is sound and complete for positive logic programs with the bounded-term-size property.
The plan of this paper is as follows. In Section 2 we present a typical example to illustrate the main idea of TP-resolution and its key differences from existing tabulated resolutions. In Section 3, we formally define TP-resolution. In Section 3.1 we discuss how to represent tables and how to operate on tables. In Section 3.2 we first introduce the so called PMF mode for resolving tabled subgoals with program clauses, which lays the basis for a linear tabulated resolution. We then define a tabulated control strategy called TP-strategy, which enhances Prolog-strategy with proper policies for the selection of answers in tables. Next we present a constructive definition (an algorithm) of a TP-tree based on TP-strategy. Finally, based on TP-trees we define TP-derivations and TP-resolution.
Section 4 is devoted to showing some major characteristics of TP-resolution, including its termination property and soundness and completeness. We also discuss in detail how TP-resolution deals with the cut operator.
We assume familiarity with the basic concepts of logic programming, as presented in . Here and throughout, variables begin with a capital letter, and predicates, functions and constants with a lower case letter. By $`\stackrel{}{E}`$ we denote a list/tuple ($`E_1,\mathrm{},E_m`$) of elements. Let $`\stackrel{}{X}=(X_1,\mathrm{},X_m)`$ be a list of variables and $`\stackrel{}{I}=(I_1,\mathrm{},I_m)`$ a list of terms. By $`\stackrel{}{X}/\stackrel{}{I}`$ we denote a substitution $`\{X_1/I_1,\mathrm{},X_m/I_m\}`$. By $`p(.)`$ we refer to any atom with the predicate $`p`$ and by $`p(\stackrel{}{X})`$ to an atom $`p(.)`$ that contains the list $`\stackrel{}{X}`$ of distinct variables. For instance, if $`p(\stackrel{}{X})=p(W,a,f(Y),W)`$, then $`\stackrel{}{X}=(W,Y)`$. Let $`G=A_1,\mathrm{},A_m`$ be a goal and $`B`$ a subgoal. By $`G+B`$ we denote the goal $`A_1,\mathrm{},A_m,B`$. By a variant of an atom (resp. a subgoal or a term) $`A`$ we mean an atom (resp. a subgoal or a term) $`A^{}`$ that is the same as $`A`$ up to variable renaming.<sup>2</sup><sup>2</sup>2By this definition, $`A`$ is a variant of itself. Let $`V`$ be a set of atoms (resp. subgoals or terms) that are variants of each other; then they are called variant atoms (resp. variant subgoals or variant terms). Moreover, clauses with the same head predicate $`p`$ are numbered sequentially, with $`C_{p_i}`$ referring to its $`i`$-th clause $`(i>0)`$. Finally, unless otherwise stated, by a (logic) program we refer to a positive logic program with a finite set of clauses.
## 2 An Illustrative Example
We use the following simple program to illustrate the basic idea of the TP approach. For convenience of presentation, we choose OLDT-resolution for a side-by-side comparison (other typical tabulated resolutions, such as SLG-resolution and Tabulated SLS-resolution , have similar effects).
| $`P_1`$: | $`reach(X,Y)reach(X,Z),edge(Z,Y).`$ $`C_{r_1}`$ |
| --- | --- |
| | $`reach(X,X).`$ $`C_{r_2}`$ |
| | $`reach(X,d).`$ $`C_{r_3}`$ |
| | $`edge(a,b).`$ $`C_{e_1}`$ |
| | $`edge(d,e).`$ $`C_{e_2}`$ |
Let $`G_0=reach(a,X)`$ be the query (top goal). Then Prolog will step into an infinite loop right after the application of the first clause $`C_{r_1}`$. We now show how it works using OLDT-resolution (under the depth-first control strategy). Starting from the root node $`N_0`$ labeled with the goal $`reach(a,X)`$, the application of the clause $`C_{r_1}`$ gives a child node $`N_1`$ labeled with the goal $`reach(a,Z),edge(Z,X)`$ (see Figure 1). Since the subgoal $`reach(a,Z)`$ is a variant of $`reach(a,X)`$ that occurred earlier, it is suspended to wait for $`reach(a,X)`$ to produce answers. $`N_0`$ and $`N_1`$ (resp. $`reach(a,X)`$ and $`reach(a,Z)`$) are then called solution and lookup nodes (resp. subgoals), respectively. So the derivation goes back to $`N_0`$ and resolves $`reach(a,X)`$ with the second clause $`C_{r_2}`$, which gives a sibling node $`N_2`$ labeled with an empty clause $`\mathrm{}`$. Since $`reach(a,a)`$ is an answer to the subgoal $`reach(a,X)`$, it is memorized in a table, say $`TB(reach(a,X))`$. The derivation then jumps back to $`N_1`$ and uses the answer $`reach(a,a)`$ in the table to resolve with the lookup subgoal $`reach(a,Z)`$, which gives a new node $`N_3`$ labeled with $`edge(a,X)`$. Next, the node $`N_4`$ labeled with $`\mathrm{}`$ is derived from $`N_3`$ by resolving the subgoal $`edge(a,X)`$ with the clause $`C_{e_1}`$. So the answer $`reach(a,b)`$ is added to the table $`TB(reach(a,X))`$. After these steps, the OLDT-derivation evolves into a tree as depicted in Figure 1, which is clearly not linear.
We now explain how TP-resolution works. Starting from the root node $`N_0`$ labeled with the goal $`reach(a,X)`$ we apply the clause $`C_{r_1}`$ to derive a child node $`N_1`$ labeled with the goal $`reach(a,Z),edge(Z,X)`$ (see Figure 2). As the subgoal $`reach(a,Z)`$ is a variant of $`reach(a,X)`$ and the latter is an ancestor of the former (i.e., the derivation steps into a loop at $`N_1`$ ), we choose $`C_{r_2}`$, the clause from the backtracking point of the subgoal $`reach(a,X)`$, to resolve with $`reach(a,Z)`$, which gives a child node $`N_2`$ labeled with $`edge(a,X)`$. Since $`reach(a,a)`$ is an answer to the subgoal $`reach(a,Z)`$, it is memorized in a table $`TB(reach(a,X))`$. We then resolve the subgoal $`edge(a,X)`$ against the clause $`C_{e_1}`$, which gives the leaf $`N_3`$ labeled with $`\mathrm{}`$. So the answer $`reach(a,b)`$ to the subgoal $`reach(a,X)`$ is added to the table $`TB(reach(a,X))`$. After these steps, we get a path as shown in Figure 2, which is clearly linear.
Now consider backtracking. Remember that after the above derivation steps, the table $`TB(`$ $`reach(a,X))`$ consists of two answers, $`reach(a,a)`$ and $`reach(a,b)`$. For the OLDT approach, it first backtracks to $`N_3`$ and then to $`N_1`$ (Figure 1). Since the subgoal $`reach(a,Z)`$ has used the first answer in the table before, it resolves with the second, $`reach(a,b)`$, which gives a new node labeled with the goal $`edge(b,X)`$. Obviously, this goal will fail, so it backtracks to $`N_1`$ again. This time no new answers in the table are available to the subgoal $`reach(a,Z)`$, so it is suspended and the derivation goes to the solution node $`N_0`$. The third clause $`C_{r_3}`$ is then selected to resolve with the subgoal $`reach(a,X)`$, yielding a new answer $`reach(a,d)`$, which is added to the table. The derivation then goes back to $`N_1`$ where the new answer is used in the same way as described before.
The TP approach does backtracking in the same way as the OLDT approach except for the following key differences: (1) Because we do not distinguish between solution and lookup nodes/subgoals, when no new answers in the table are available to the subgoal $`reach(a,Z)`$ at $`N_1`$, we backtrack the subgoal by resolving it against the next clause $`C_{r_3}`$. This guarantees that TP-derivations are always linear. (2) Since there is a loop between $`N_0`$ and $`N_1`$, before failing the subgoal $`reach(a,X)`$ at $`N_0`$ via backtracking we need to be sure that the subgoal has got its complete set of answers. This is achieved by performing answer iteration via the loop. That is, we regenerate the loop to see if any new answers can be derived until we reach a fixpoint. Figure 3 shows the first part of TP-resolution, where the following answers to $`G_0`$ are derived: $`X=a`$, $`X=b`$, $`X=d`$ and $`X=e`$. Figure 4 shows the part of answer iteration. Since no new answer is derived during the iteration (i.e. no answer is added to any tables), we fail the subgoal $`reach(a,X)`$ at $`N_0`$.
###### Remark 2.1
From the above illustration, we see that in OLDT-resolution, solution nodes are those at which the left-most subgoals are generated earliest among all their variant subgoals. In SLG-resolution, however, solution nodes are roots of trees in a search forest, each labeled by a special clause of the form $`AA`$ . In Tabulated SLS-resolution, any root of a tree in a forest is itself labeled by an instance, say $`AB_1,\mathrm{},B_n`$ $`(n0)`$, of a program clause and no nodes in the tree will produce child nodes using program clauses . However, for any atom $`A`$ we can assume a virtual super-root labeled with $`AA`$, which takes all the roots in the forest labeled by $`A\mathrm{}`$ as its child nodes. In this sense, the search forest in Tabulated SLS-resolution is the same as that in SLG-resolution for positive logic programs. Therefore, we can consider all virtual super-roots as solution nodes.
## 3 TP-Resolution
This section formally defines the TP approach to tabulated resolution, mainly including the representation of tables, the strategy for controlling tabulated derivations (TP-strategy), and the algorithm for making tabulated derivations based on the control strategy (TP-trees).
### 3.1 Tabled Predicates and Tables
Predicates in a program $`P`$ are classified as tabled predicates and non-tabled predicates. The classification is made based on a dependency graph . Informally, for any predicates $`p`$ and $`q`$, there is an edge $`pq`$ in a dependency graph $`G_P`$ if there is a clause in $`P`$ of the form $`p(.)\mathrm{},q(.),\mathrm{}`$ Then a predicate $`p`$ is to be tabled if $`G_P`$ contains a cycle with a node $`p`$.
Any atom/subgoal with a tabled predicate is called a tabled atom/subgoal. During tabulated resolution, we will create a table for each tabled subgoal, $`A`$. Apparently, the table must contain $`A`$ (as an index) and have space to store intermediate answers of $`A`$. Note that in our tabling approach, any tabled subgoal can act both as a solution subgoal and as a lookup subgoal, so a table can be viewed as a blackboard on which a set of variant subgoals will read and write answers. In order to guarantee not losing answers for any tabled subgoals (i.e. the table should contain all answers that $`A`$ is supposed to have by applying its related clauses), while avoiding redundant computations (i.e. after a clause has been used by $`A`$, it should not be re-used by any other variant subgoal $`A^{}`$), a third component is needed in the table that keeps the status of the clauses related to $`A`$. Therefore, after a clause $`C_i`$ has been used by $`A`$, we change its status. Then when evaluating a new subgoal $`A^{}`$ that is a variant of $`A`$, $`C_i`$ will be ignored because all answers of $`A`$ derived via $`C_i`$ have already been stored in the table. For any clause whose head is a tabled atom, its status can be “no longer available” or “still available.” We say that $`C_i`$ is “no longer available” to $`A`$ if all answers of $`A`$ through the application of $`C_i`$ have already been stored in the table of $`A`$. Otherwise, we say $`C_i`$ is “still available” to $`A`$. Finally, we need a flag variable $`COMP`$ in the table to indicate if all answers through the application of all clauses related to $`A`$ have been completely stored in the table. This leads to the following.
###### Definition 3.1
Let $`P`$ be a logic program and $`p(\stackrel{}{X})`$ a tabled subgoal. Let $`P`$ contain exactly $`M`$ clauses, $`C_{p_1},\mathrm{},C_{p_M}`$, with a head $`p(.)`$. A table for $`p(\stackrel{}{X})`$, denoted $`TB(p(\stackrel{}{X}))`$, is a four-tuple $`(p(\stackrel{}{X}),T,C,COMP)`$, where
1. $`T`$ consists of tuples that are instances of $`\stackrel{}{X}`$, each $`\stackrel{}{I}`$ of which represents an answer, $`p(\stackrel{}{X})\stackrel{}{X}/\stackrel{}{I}`$, to the subgoal.
2. $`C`$ is a vector of $`M`$ elements, with $`C[i]=0`$ (resp. $`=1`$) representing that the status of $`C_{p_i}`$ w.r.t. $`p(\stackrel{}{X})`$ is “no longer available” (resp. “still available”).
3. $`COMP\{0,1\}`$, with $`COMP=1`$ indicating that the answers of $`p(\stackrel{}{X})`$ have been completed.
For convenience, we use $`TB(p(\stackrel{}{X}))answer\mathrm{\_}tuple[i]`$ to refer to the $`i`$-th answer tuple in $`T`$, $`TB(p(\stackrel{}{X}))clause\mathrm{\_}status[i]`$ to the status of $`C_{p_i}`$ w.r.t. $`p(\stackrel{}{X})`$, and $`TB(p(\stackrel{}{X}))COMP`$ to the flag $`COMP`$.
###### Example 3.1
Let $`P`$ be a logic program that contains exactly three clauses, $`C_{p_1},C_{p_2}`$ and $`C_{p_3}`$, with a head $`p(.)`$. The table
$`TB(p(X,Y)):`$ $`(p(X,Y),`$ $`\{(a,b),(b,a),(b,c)\},(1,0,0),0)`$
represents that there are three answers to $`p(X,Y)`$, namely $`p(a,b)`$, $`p(b,a)`$ and $`p(b,c)`$, and that $`C_{p_2}`$ and $`C_{p_3}`$ have already been used by $`p(X,Y)`$ (or its variant subgoals) and $`C_{p_1}`$ is still available to $`p(X,Y)`$. Obviously, the answers of $`p(X,Y)`$ have not yet been completed. The table
$`TB(p(a,b)):`$ $`(p(a,b),`$ $`\{()\},(0,1,1),1)`$
shows that $`p(a,b)`$ has been proved true after applying $`C_{p_1}`$. Note that since $`p(a,b)`$ contains no variables, its answer is a 0-ary tuple. Finally, the table
$`TB(p(a,X)):`$ $`(p(a,X),`$ $`\{\},(0,0,0),1)`$
represents that $`p(a,X)`$ has no answer at all.
Before introducing operations on tables, we define the structure of nodes used in TP-resolution.
###### Definition 3.2
Let $`P`$ be a logic program and $`G_i`$ a goal $`p(\stackrel{}{X}),A_2,\mathrm{},A_m`$. By “register a node $`N_i`$ with $`G_i`$” we do the following: (1) label $`N_i`$ with $`G_i`$, i.e. $`N_i:p(\stackrel{}{X}),A_2,\mathrm{},A_m`$; and (2) create the following structure for $`N_i`$:
$``$ $`answer\mathrm{\_}ptr`$, a pointer that points to an answer tuple in $`TB(p(\stackrel{}{X}))`$.
$``$ $`clause\mathrm{\_}ptr`$, a pointer that points to a clause in $`P`$ with a head $`p(.)`$.
$``$ $`clause\mathrm{\_}SUSP`$ (initially =0), a flag used for the update of clause status.
$``$ $`node\mathrm{\_}LOOP`$ (initially =0), a flag showing if $`N_i`$ is a loop node.
$``$ $`node\mathrm{\_}ITER`$ (initially =0), a flag showing if $`N_i`$ is an iteration node.
$``$ $`node\mathrm{\_}ANC`$ (initially =$`1`$), a flag showing if $`N_i`$ has any ancestor variant subgoals.
For any field $`F`$ in the structure of $`N_i`$, we refer to it by $`N_iF`$. The meaning of $`N_ianswer\mathrm{\_}ptr`$ and $`N_iclause\mathrm{\_}ptr`$ is obvious. The remaining fields will be defined by Definition 3.8 followed by the procedure $`nodetype\mathrm{\_}update(.)`$. We are now ready to define operations on tables.
###### Definition 3.3
Let $`P`$ be a logic program with $`M`$ clauses with a head $`p(.)`$ and $`N_i`$ a node labeled by a goal $`p(\stackrel{}{X}),\mathrm{},A_m.`$ Let $`NEW`$ be a global flag variable used for answer iteration (see Algorithm 2 for details). We have the following basic operations on a table.
1. $`create(p(\stackrel{}{X}))`$. Create a table $`TB(p(\stackrel{}{X})):`$ $`(p(\stackrel{}{X}),T,C,COMP)`$, with $`T=\{\}`$, $`COMP=0`$, and $`C[j]=1`$ for all $`1jM`$.
2. $`memo(p(\stackrel{}{X}),\stackrel{}{I})`$, where $`\stackrel{}{I}`$ is an instance of $`\stackrel{}{X}`$. When $`\stackrel{}{I}`$ is not in $`TB(p(\stackrel{}{X}))`$, add it to the end of the table, set $`NEW=1`$, and if $`\stackrel{}{I}`$ is a variant of $`\stackrel{}{X}`$, set $`TB(p(\stackrel{}{X}))COMP=1`$.
3. $`lookup(N_i,\stackrel{}{I_i})`$. Fetch the next answer tuple in $`TB(p(\stackrel{}{X}))`$, which is pointed by $`N_ianswer\mathrm{\_}ptr`$, into $`\stackrel{}{I_i}`$. If there is no next tuple, $`\stackrel{}{I_i}=null`$.
4. $`memo\mathrm{\_}look(N_i,p(\stackrel{}{X}),\stackrel{}{I},\theta _i)`$. It is a compact operator, which combines $`memo(.)`$ and $`lookup(.)`$. That is, it first performs $`memo(p(\stackrel{}{X}),\stackrel{}{I})`$ and then gets the next answer tuple $`\stackrel{}{F}`$ from $`TB(p(\stackrel{}{X}))`$, which together with $`\stackrel{}{X}`$ forms a substitution $`\theta _i=\stackrel{}{X}/\stackrel{}{F}`$. If there is no next tuple, $`\theta _i=null`$.
First, the procedure $`create(p(\stackrel{}{X}))`$ is called only when the subgoal $`p(\stackrel{}{X})`$ occurs the first time and no variant subgoals occurred before. Therefore, up to the time when we call $`create(p(\stackrel{}{X}))`$, no clauses with a head $`p(.)`$ in $`P`$ have been selected by any variant subgoals of $`p(\stackrel{}{X})`$, so their status should be set to $`1`$. Second, whenever an answer $`p(\stackrel{}{I})`$ of $`p(\stackrel{}{X})`$ is derived, we call the procedure $`memo(p(\stackrel{}{X}),\stackrel{}{I})`$. If the answer is new, it is appended to the end of the table. The flag $`NEW`$ is then set to $`1`$, showing that a new answer has been derived. If the new tuple $`\stackrel{}{I}`$ is a variant of $`\stackrel{}{X}`$, which means that $`p(\stackrel{}{X})`$ is true for any instances of $`\stackrel{}{X}`$, the answers of $`p(\stackrel{}{X})`$ are completed so $`TB(p(\stackrel{}{X}))COMP`$ is set to 1. Finally, $`lookup(N_i,\stackrel{}{I_i})`$ is used to fetch an answer tuple from the table for the subgoal $`p(\stackrel{}{X})`$ at $`N_i`$.
$`memo(.)`$ and $`lookup(.)`$ can be used independently. They can also be used in pairs, i.e. $`memo(.)`$ immediately followed by $`lookup(.)`$. In the latter case, it would be more convenient to use $`memo\mathrm{\_}look(.)`$.
### 3.2 TP-Strategy and TP-Trees
In this subsection, we introduce the tabulated control strategy and the way to make tabulated derivations based on this strategy. We begin by discussing how to resolve subgoals with program clauses and answers in tables.
Let $`N_i`$ be a node labeled by a goal $`G_i=A_1,\mathrm{},A_m`$ with $`A_1=p(\stackrel{}{X})`$ a tabled subgoal. Consider evaluating $`A_1`$ using a program clause $`C_p=AB_1,\mathrm{},B_n`$ $`(n0)`$, where $`A_1\theta =A\theta `$.<sup>3</sup><sup>3</sup>3Here and throughout, we assume that $`C_p`$ has been standardized apart to share no variables with $`G_i`$. If we use SLD-resolution, we would obtain a new node labeled with the goal $`G_{i+1}=(B_1,\mathrm{},B_n,A_2,\mathrm{},A_m)\theta `$, where the mgu $`\theta `$ is consumed by all $`A_j`$s $`(j>1)`$, although the proof of $`A_1\theta `$ has not yet been completed (produced). In order to avoid such kind of pre-consumption, we propose the so called PMF (for Prove-Memorize-Fetch) mode for resolving tabled subgoals with clauses. That is, we first prove $`(B_1,\mathrm{},B_n)\theta `$. If it is true with an mgu $`\theta _1`$, which means $`A_1\theta \theta _1`$ is true, we memorize the answer $`A_1\theta \theta _1`$ in the table $`TB(A_1)`$ if it is new. We then fetch an answer from $`TB(A_1)`$ to apply to the remaining subgoals of $`G_i`$. Obviously modifying SLD-resolution by the PMF mode preserves the original answers to $`G_i`$. Moreover, since only new answers are added to $`TB(A_1)`$, all repeated answers of $`A_1`$ will be precluded to apply to the remaining subgoals of $`G_i`$, so that redundant computations are avoided.
The PMF mode can readily be realized by using the two table procedures, $`memo(.)`$ and $`lookup(.)`$, or using the compact operator $`memo\mathrm{\_}look(.)`$. That is, after resolving the subgoal $`A_1`$ with the clause $`C_p`$, $`N_i`$ gets a child node $`N_{i+1}`$ labeled with the goal
$`G_{i+1}=(B_1,\mathrm{},B_n)\theta ,memo\mathrm{\_}look(N_i,p(\stackrel{}{X}),\stackrel{}{X}\theta ,\theta _i),A_2,\mathrm{},A_m`$.
Note that the application of $`\theta `$ is blocked by the subgoal $`memo\mathrm{\_}look(.)`$ because the consumption (fetch) must follow the production (prove and memorize). We now explain how it works.
Assume that after some resolution steps from $`N_{i+1}`$ we reach a node $`N_k`$ that is labeled by the goal $`G_k=memo\mathrm{\_}look(N_i,p(\stackrel{}{X}),\stackrel{}{X}\theta \theta _1,\theta _i),A_2,\mathrm{},A_m`$. This means that $`(B_1,\mathrm{},B_n)\theta `$ has been proved true with the mgu $`\theta _1`$. That is, $`A_1\theta \theta _1`$ is an answer of $`A_1`$. By the left-most computation rule, $`memo\mathrm{\_}look(N_i,p(\stackrel{}{X}),\stackrel{}{X}\theta \theta _1,\theta _i)`$ is executed, which adds to the table $`TB(A_1)`$ the answer tuple $`\stackrel{}{X}\theta \theta _1`$ if it is new, gets from $`TB(A_1)`$ the next tuple $`\stackrel{}{I}`$, and then sets $`\theta _i=\stackrel{}{X}/\stackrel{}{I}`$. Since $`A_1\theta _i`$ is an answer to the subgoal $`A_1`$ of $`G_i`$, the mgu $`\theta _i`$ needs to be applied to the remaining $`A_j`$s of $`G_i`$. We distinguish between two cases.
1. From $`A_2`$ to $`A_m,`$ $`A_j=memo\mathrm{\_}look(N_f,B,\mathrm{\_},\theta _f)`$ is the first subgoal of the form $`memo\mathrm{\_}look(.)`$. According to the PMF mode, there must be a node $`N_f`$, which occurred earlier than $`N_i`$, labeled with a goal $`G_f=B,A_{j+1},\mathrm{},A_m`$ such that $`B`$ is a tabled subgoal and $`A_j=memo\mathrm{\_}look(N_f,B,\mathrm{\_},\theta _f)`$ resulted from resolving $`B`$ with a program clause. This means that the proof of $`B`$ is now reduced to the proof of $`(A_2,\mathrm{},A_{j1})\theta _i`$. Therefore, by the PMF mode $`\theta _i`$ should be applied to the subgoals $`A_2`$ until $`A_j`$. That is, $`N_k`$ has a child node $`N_{k+1}`$ labeled with a goal $`G_{k+1}=(A_2,\mathrm{},A_j)\theta _i,A_{j+1},\mathrm{},A_m`$.
2. For no $`j2`$ $`A_j`$ is of the form $`memo\mathrm{\_}look(.)`$. This means that no $`A_j`$ is a descendant of any tabled subgoal, so the mgu $`\theta _i`$ should be applied to all the $`A_j`$s. That is, $`G_{k+1}=(A_2,\mathrm{},A_m)\theta _i`$.
Note that by Definition 3.3 the atom $`p(\stackrel{}{X})`$ in $`memo(p(\stackrel{}{X}),\mathrm{\_})`$ and $`memo\mathrm{\_}look(\mathrm{\_},p(\stackrel{}{X}),\mathrm{\_},\mathrm{\_})`$ is merely used to index the table $`TB(p(\stackrel{}{X}))`$, so it cannot be instantiated during the resolution. That is, for any mgu $`\theta `$, $`memo(p(\stackrel{}{X}),\stackrel{}{I})\theta =memo(p(\stackrel{}{X}),\stackrel{}{I}\theta )`$ and $`memo\mathrm{\_}look(N_i,p(\stackrel{}{X}),\stackrel{}{I},\theta _i)\theta =memo\mathrm{\_}look(N_i,p(\stackrel{}{X}),\stackrel{}{I}\theta ,\theta _i)`$
The above discussion shows how to resolve the tabled subgoal $`A_1`$ at $`N_i`$ against a program clause using the PMF mode. The same principle can be applied to resolve $`A_1`$ with an answer tuple $`\stackrel{}{I}`$ in $`TB(A_1)`$ and to resolve $`A_1`$ with a program clause when $`A_1`$ is a non-tabled subgoal. Therefore, we have the following definition of resolvents for TP-resolution.
###### Definition 3.4
Let $`N_i`$ be a node labeled by a goal $`G_i=A_1,\mathrm{},A_m`$ $`(m1)`$.
1. If $`A_1`$ is $`memo\mathrm{\_}look(N_h,p(\stackrel{}{X}),\stackrel{}{I},\theta _h)`$, then the resolvent of $`G_i`$ and $`\theta _h`$ ($`\theta _hnull`$) is the goal $`G_{i+1}=(A_2,\mathrm{},A_k)\theta _h,A_{k+1},\mathrm{},A_m`$, where $`A_k`$ $`(k>1)`$ is the left-most subgoal of the form $`memo\mathrm{\_}look(.)`$.
Otherwise, let $`A_1=p(\stackrel{}{X})`$ and $`C_p`$ be a program clause $`AB_1,\mathrm{},B_n`$ with $`A\theta =A_1\theta `$.
2. If $`A_1`$ is a non-tabled subgoal, the resolvent of $`G_i`$ and $`C_p`$ is the goal $`G_{i+1}=(B_1,\mathrm{},B_n,`$ $`A_2,\mathrm{},A_k)\theta ,`$ $`A_{k+1},\mathrm{},A_m`$, where $`A_k`$ is the left-most subgoal of the form $`memo\mathrm{\_}look(.)`$.
3. If $`A_1`$ is a tabled subgoal, the resolvent of $`G_i`$ and $`C_p`$ is the goal $`G_{i+1}=(B_1,\mathrm{},B_n)\theta ,`$ $`memo\mathrm{\_}look(N_i,p(\stackrel{}{X}),\stackrel{}{X}\theta ,\theta _i),A_2,\mathrm{},A_m`$.
4. If $`A_1`$ is a tabled subgoal, let $`\stackrel{}{I}`$ ($`\stackrel{}{I}null`$) be an answer tuple in $`TB(A_1)`$, then the resolvent of $`G_i`$ and $`\stackrel{}{I}`$ is the goal $`G_{i+1}=(A_2,\mathrm{},A_k)\stackrel{}{X}/\stackrel{}{I},A_{k+1},\mathrm{},A_m`$, where $`A_k`$ is the left-most subgoal of the form $`memo\mathrm{\_}look(.)`$.
We now discuss tabulated control strategies. Recall that Prolog implements SLD-resolution by sequentially searching an SLD-tree using the Prolog control strategy (Prolog-strategy, for short): Depth-first (for goal selection) + Left-most (for subgoal selection) \+ Top-down (for clause selection) + Last-first (for backtracking). Let “register a node $`N_i`$ with $`G_i`$ be as defined by Definition 3.2 except that the structure of $`N_i`$ only contains the pointer $`clause\mathrm{\_}ptr`$. Let $`return(\stackrel{}{Z})`$ be a procedure that returns $`\stackrel{}{Z}`$ when $`\stackrel{}{Z}()`$ and YES otherwise. Then the way that Prolog makes SLD-derivations based on Prolog-strategy can be formulated as follows.
###### Definition 3.5 (Algorithm 1)
Let $`P`$ be a logic program and $`G_0`$ a top goal with the list $`\stackrel{}{Y}`$ of variables. The Prolog-tree $`T_{G_0}`$ of $`P\{G_0\}`$ is constructed by recursively performing the following steps until the answer $`NO`$ is returned.
1. (Root node) Register the root $`N_0`$ with $`G_0+return(\stackrel{}{Y})`$ and goto 2.
2. (Node expansion) Let $`N_i`$ be the latest registered node labeled by $`G_i=A_1,\mathrm{},A_m`$ $`(i0,m>0)`$. Register $`N_{i+1}`$ as a child of $`N_i`$ with $`G_{i+1}`$ if $`G_{i+1}`$ can be obtained as follows.
* Case 1: $`A_1`$ is $`return(.)`$. Execute the procedure $`return(.)`$, set $`G_{i+1}=\mathrm{}`$ (an empty clause), and goto $`3`$ with $`N=N_i`$.
* Case 2: $`A_1`$ is an atom. Get a program clause $`AB_1,\mathrm{},B_n`$ (top-down via the pointer $`N_iclause\mathrm{\_}ptr`$) such that $`A_1\theta =A\theta `$. If no such a clause exists, then goto $`3`$ with $`N=N_i`$; else set $`G_{i+1}=(B_1,\mathrm{}B_n,A_2,\mathrm{},A_m)\theta `$ and goto 2.
3. (Backtracking) If $`N`$ is the root, then return $`NO`$; else goto $`2`$ with its parent node as the latest registered node.
Let $`ST_{G_0}`$ be the SLD-tree of $`P\{G_0\}`$ via the left-most computation rule.<sup>4</sup><sup>4</sup>4In , it is called an OLD-tree. It is easy to prove that when $`P`$ has the bounded-term-size property and $`ST_{G_0}`$ contains no infinite loops, Algorithm 1 is sound and complete in that $`T_{G_0}=ST_{G_0}`$. Moreover, Algorithm 1 has the following distinct advantages: (1) since SLD-resolution is linear, Algorithm 1 can be efficiently implemented using a simple stack-based memory structure; (2) due to its linearity and regular sequentiality, some useful control mechanisms, such as the well-known cut operator !, can be used to heuristically reduce search space. Unfortunately, Algorithm 1 suffers from two serious problems. One is that it is easy to get into infinite loops even for very simple programs such as $`P=\{p(X)p(X)\}`$, which makes it incomplete in many cases. The second problem is that it unnecessarily re-applies the same set of clauses to variant subgoals such as in the query $`p(X),p(Y)`$, which leads to unacceptable performance.
As tabling has a distinct advantage of resolving infinite loops and redundant derivations, one interesting question then arises: Can we enhance Algorithm 1 with tabling, making it free from infinite loops and redundant computations while preserving the above two advantages? In the rest of this subsection, we give a constructive answer to this question. We first discuss how to enhance Prolog-strategy with tabling.
Observe that in a tabling system, we will have both program clauses and tables. For convenience, we refer to answer tuples in tables as tabled facts. Therefore, in addition to the existing policies in Prolog-strategy, we need to have the following two additional policies: (1) when both program clauses and tabled facts are available, first use tabled facts (i.e. Table-first for program and table selection); (2) when there are more than one tabled fact available, first use the one that is earliest memorized. Since we always add new answers to the end of tables (see Definition 3.3 for $`memo(.)`$), policy (2) amounts to saying Top-down selection for tabled facts. This leads to the following control strategy for tabulated derivations.
###### Definition 3.6
By TP-strategy we mean: Depth-first (for goal selection) + Left-most (for subgoal selection) \+ Table-first (for program and table selection) \+ Top-down (for the selection of tabled facts and program clauses) + Last-first (for backtracking).
Our goal is to extend Algorithm 1 to make linear tabulated derivations based on TP-strategy. To this end, we need to review a few concepts concerning loop checking.
###### Definition 3.7 ( with slight modification)
An ancestor list $`AL_A`$ of pairs $`(N,B)`$ is associated with each tabled subgoal $`A`$ at a node $`N_i`$ in a tree (see the TP-tree below), which is defined recursively as follows.
1. If $`A`$ is at the root, then $`AL_A=\{\}`$.
2. If $`A`$ inherits a subgoal $`A^{}`$ (by copying or instantiation) from its parent node, then $`AL_A=AL_A^{}`$.
3. Let $`A`$ be in the resolvent of a subgoal $`B`$ at $`N_f`$ against a clause $`B^{}A_1,\mathrm{},A_n`$ with $`B\theta =B^{}\theta `$ (i.e. $`A=A_i\theta `$ for some $`1in`$). If $`B`$ is a tabled subgoal, $`AL_A=AL_B\{(N_f,B)\}`$; otherwise $`AL_A=\{\}`$.
We see that for any tabled subgoals $`A`$ and $`A^{}`$, if $`A`$ is in the ancestor list of $`A^{}`$, i.e. $`(\mathrm{\_},A)AL_A^{}`$, the proof of $`A`$ needs the proof of $`A^{}`$. Particularly, if $`(\mathrm{\_},A)AL_A^{}`$ and $`A^{}`$ is a variant of $`A`$, the derivation goes into a loop. This leads to the following.
###### Definition 3.8
Let $`G_i`$ at $`N_i`$ and $`G_k`$ at $`N_k`$ be two goals in a derivation and $`A_i`$ and $`A_k`$ be the left-most subgoals of $`G_i`$ and $`G_k`$, respectively. We say $`A_i`$ (resp. $`N_i`$) is an ancestor subgoal of $`A_k`$ (resp. an ancestor node of $`N_k`$) if $`(N_i,A_i)AL_{A_k}`$. If $`A_i`$ is both an ancestor subgoal and a variant, i.e. an ancestor variant subgoal, of $`A_k`$, we say the derivation goes into a loop, denoted $`L(N_i,N_k)`$. Then, $`N_k`$ and all its ancestor nodes involved in the loop are called loop nodes. $`N_i`$ is also called the top loop node of the loop. Finally, a loop node is called an iteration node if by the time the node is about to fail through backtracking, it is the top loop node of all loops containing the node that were generated before.
###### Example 3.2
Figure 5 shows four loops, $`L_1`$, …, $`L_4`$, with $`N_1`$, …, $`N_4`$ their respective top loop nodes. We see that only $`N_1`$ and $`N_4`$ are iteration nodes.
Information about the types and ancestors of nodes is the basis on which we make tabulated resolution. Such information is kept in the structure of each node $`N_i`$ (see Definition 3.2). The flag $`N_inode\mathrm{\_}LOOP=1`$ shows that $`N_i`$ is a loop node. The flag $`N_inode\mathrm{\_}ITER=1`$ shows that $`N_i`$ is an (candidate) iteration node. Let $`A_1=p(\stackrel{}{X})`$ be the left-most subgoal at $`N_i`$. The flag $`N_inode\mathrm{\_}ANC=1`$ represents that it is unknown whether $`A_1`$ has any ancestor variant subgoal; $`N_inode\mathrm{\_}ANC=0`$ shows that $`A_1`$ has no ancestor variant subgoal; and $`N_inode\mathrm{\_}ANC=j`$ $`(j>0)`$ indicates that $`A_1`$ has ancestor variant subgoals and that $`C_{p_j}`$ is the clause that is being used by its closest ancestor variant subgoal (i.e., let $`A_h`$ at $`N_h`$ be the closest ancestor variant subgoal of $`A_1`$, then $`N_inode\mathrm{\_}ANC=j`$ represents that $`N_i`$ is derived from $`N_h`$ via $`C_{p_j}`$).
Once a loop, say $`L(N_1,N_m)`$, of the form
$`(N_1:A_1,\mathrm{})_{C_{p_j},\theta _1}(N_2:A_2,\mathrm{})\mathrm{}(N_m:A_m,\mathrm{})`$
occurs, where all $`N_i`$s ($`i<m`$) are ancestor nodes of $`N_m`$ and $`A_1=p(\stackrel{}{X})`$ is the closest ancestor variant subgoal of $`A_m`$, we update the flags of all nodes, $`N_1,\mathrm{},N_m`$, involved in the loop by calling the following procedure.
| Procedure $`nodetype\mathrm{\_}update(L(N_1,N_m))`$ |
| --- |
| | (1) For all $`i>1`$ set $`N_inode\mathrm{\_}LOOP=1`$ and $`N_inode\mathrm{\_}ITER=0`$. |
| | (2) If $`N_1node\mathrm{\_}LOOP=0`$, set $`N_1node\mathrm{\_}LOOP=1`$ and $`N_1node\mathrm{\_}ITER=1`$. |
| | (3) Set $`N_mnode\mathrm{\_}ANC=j`$. |
| | (4) For all $`i<m`$ set $`N_iclause\mathrm{\_}SUSP=1`$. |
Point (1) is straightforward, where since $`N_1`$ is the top loop node of $`L(N_1,N_m)`$, all the remaining nodes in the loop cannot be an iteration node (see Definition 3.8).
If $`N_1node\mathrm{\_}LOOP=0`$, meaning that $`N_1`$ is not involved in any loop that occurred before, $`N_1`$ is considered as a candidate iteration node (point (2)). A candidate iteration node becomes an iteration node if the node keeps its candidacy by the time it is about to fail through backtracking (by that time it must be the top loop node of all previously generated loops containing it).
Since $`A_1`$ is the closest ancestor variant subgoal of $`A_m`$ and $`C_{p_j}`$ is the clause that is being used by $`A_1`$, we set the flag $`N_mnode\mathrm{\_}ANC=j`$ (point (3)).
As mentioned in Section 2, during TP-resolution when a loop $`L(N_1,N_m)`$ occurs, where the left-most subgoal $`A_1=p(\stackrel{}{X})`$ at $`N_1`$ is the closest ancestor variant subgoal of the left-most subgoal $`A_m`$ at $`N_m`$, $`A_m`$ will skip the clause $`C_{p_j}`$ that is being used by $`A_1`$. In order to ensure that such a skip will not lead to loss of answers to $`A_1`$, we will do answer iteration before failing $`N_1`$ via backtracking until we reach a fixpoint of answers. Answer iteration is done by regenerating $`L(N_1,N_m)`$. This requires keeping the status of all clauses being used by the loop nodes to “still available” during backtracking. Point (4) is used for such a purpose. After the flag $`N_iclause\mathrm{\_}SUSP`$ is set to 1, which indicates that $`N_i`$ is currently involved in a loop, the status of the clause being currently used by $`N_i`$ will not be set to “no longer available” when backtracking on $`N_i`$ (see Case B3 of Algorithm 2).
###### Remark 3.1
We do answer iteration only at iteration nodes because they are the top nodes of all loops involving them. If we did answer iteration at a non-iteration loop node $`N`$, we would have to do it again at some top loop node $`N_{top}`$ of $`N`$, in order to reach a fixpoint at $`N_{top}`$ (see Figure 5). This would certainly lead to more redundant computations.
We are now in a position to define the TP-tree, which is constructed based on the TP-strategy using the following algorithm.
###### Definition 3.9 (Algorithm 2)
Let $`P`$ be a logic program and $`G_0`$ a top goal with the list $`\stackrel{}{Y}`$ of variables. The TP-tree $`TP_{G_0}`$ of $`P\{G_0\}`$ is constructed by recursively performing the following steps until the answer $`NO`$ is returned.
1. (Root node) Register the root $`N_0`$ with $`G_0+return(\stackrel{}{Y})`$, set $`NEW=0`$, and goto 2.
2. (Node expansion) Let $`N_i`$ be the latest registered node labeled by $`G_i=A_1,\mathrm{},A_m`$ $`(m>0)`$. Register $`N_{i+1}`$ as a child of $`N_i`$ with $`G_{i+1}`$ if $`G_{i+1}`$ can be obtained as follows.
* Case 1: $`A_1`$ is $`return(.)`$. Execute the procedure $`return(.)`$, set $`G_{i+1}=\mathrm{}`$ (an empty clause), and goto $`3`$ with $`N=N_i`$.
* Case 2: $`A_1`$ is $`memo\mathrm{\_}look(N_h,p(\stackrel{}{X}),\stackrel{}{I},\theta _h)`$. Execute the procedure.<sup>5</sup><sup>5</sup>5See Definition 3.3, where the flags $`NEW`$ and $`TB(p(\stackrel{}{X}))COMP`$ will be updated. If $`\theta _h=null`$ then goto $`3`$ with $`N=N_i`$; else set $`G_{i+1}`$ to the resolvent of $`G_i`$ and $`\theta _h`$ and goto 2.
* Case 3: $`A_1`$ is a non-tabled subgoal. Get a clause $`C`$ whose head is unifiable with $`A_1`$.<sup>6</sup><sup>6</sup>6Here and throughout, clauses and answers in tables are selected top-down via the pointers $`N_iclause\mathrm{\_}ptr`$ and $`N_ianswer\mathrm{\_}ptr`$, respectively. If no such a clause exists then goto $`3`$ with $`N=N_i`$; else set $`G_{i+1}`$ to the resolvent of $`G_i`$ and $`C`$ and goto 2.
* Case 4: $`A_1=p(\stackrel{}{X})`$ is a tabled subgoal. Get an instance $`\stackrel{}{I}`$ of $`\stackrel{}{X}`$ from the table $`TB(A_1)`$. If $`\stackrel{}{I}null`$ then set $`G_{i+1}`$ to the resolvent of $`G_i`$ and $`\stackrel{}{I}`$ and goto 2. Otherwise, if $`TB(A_1)COMP=1`$ then goto 3 with $`N=N_i`$; else
+ Case 4.1: $`N_inode\mathrm{\_}ANC=1`$. If $`A_1`$ has no ancestor variant subgoal, set $`N_inode\mathrm{\_}ANC=0`$ and goto Case 4.2. Otherwise, let $`N_h`$ be the closest ancestor node of $`N_i`$ such that $`L(N_h,N_i)`$ is a loop. Call $`nodetype\mathrm{\_}update(L(N_h,N_i))`$ and goto Case 4.3.
+ Case 4.2: $`N_inode\mathrm{\_}ANC=0`$. Get a clause $`C_{p_j}`$ whose head is unifiable with $`A_1`$ such that $`TB(A_1)clause\mathrm{\_}status[j]=1`$. If such a clause exists, set $`G_{i+1}`$ to the resolvent of $`G_i`$ and $`C_{p_j}`$ and goto 2. Otherwise, if $`N_inode\mathrm{\_}ITER=0`$ then goto 3 with $`N=N_i`$; else
- Case 4.2.1: $`NEW=0`$. Set $`TB(A_1)COMP=1`$ and goto $`3`$ with $`N=N_i`$.
- Case 4.2.2: $`NEW=1`$. Set $`NEW=0`$, reset $`N_iclause\mathrm{\_}ptr`$ to pointing to the first clause $`C_{p_j}`$ whose status is “still available”, and goto Case 4.2.
+ Case 4.3: $`N_inode\mathrm{\_}ANC=j`$ $`(j>0)`$. Get a clause $`C_{p_k}`$ ($`k>j`$) whose head is unifiable with $`A_1`$ such that $`TB(A_1)clause\mathrm{\_}status[k]=1`$. If such a clause exists then set $`G_{i+1}`$ to the resolvent of $`G_i`$ and $`C_{p_k}`$ and goto 2; else goto $`3`$ with $`N=N_i`$.
3. (Backtracking) If $`N`$ is the root, return $`NO`$. Otherwise let $`N_f`$ be the parent node of $`N`$ with the left-most subgoal $`A_f`$.
* Case B1: $`A_f`$ is $`memo\mathrm{\_}look(.)`$. Goto $`3`$ with $`N=N_f`$.
* Case B2: $`A_f`$ is a non-tabled subgoal. Goto $`2`$ with $`N_f`$ as the latest registered node.
* Case B3: $`A_f=q(\stackrel{}{Z})`$ is a tabled subgoal. Let $`N`$ be generated from $`N_f`$ by resolving $`A_f`$ with a clause $`C_{q_j}`$. If $`N_fnode\mathrm{\_}SUSP=0`$ then set $`TB(A_f)clause\mathrm{\_}status[j]=0`$; else set $`N_fnode\mathrm{\_}SUSP=0`$. Goto $`2`$ with $`N_f`$ as the latest registered node.
Obviously, Algorithm 2 reduces to Algorithm 1 when $`P`$ contains no tabled predicates. We now explain Algorithm 2 briefly. First we set up the root $`N_0`$ via registration (see Definition 3.2). The global variable $`NEW`$ is initialized to $`0`$, meaning that up to now no new answer has been derived for any subgoal. Then by the Depth-first policy we select the latest registered node, say $`N_i`$ labeled with the goal $`G_i`$, for expansion (point 2). If the left-most subgoal $`A_1`$ of $`G_i`$ is $`return(\stackrel{}{I})`$ (Case 1), which means the top goal $`G_0`$ has been proved true with the answer substitution $`\stackrel{}{Y}/\stackrel{}{I}`$, we reach a success leaf $`N_{i+1}`$ labeled with $`\mathrm{}`$. We then do backtracking (point 3) to derive alternative answers to $`G_0`$.
If $`A_1`$ is $`memo\mathrm{\_}look(N_h,p(\stackrel{}{X}),\stackrel{}{I},\theta _h)`$ (Case 2), which means that the left-most subgoal $`p(\stackrel{}{X})`$ at node $`N_h`$ is proved true with the answer substitution $`\stackrel{}{X}/\stackrel{}{I}`$, we memorize $`\stackrel{}{I}`$ in the table $`TB(p(\stackrel{}{X}))`$ and set $`NEW=1`$ if the answer is new. Meanwhile, if the answer $`p(\stackrel{}{I})`$ is a variant of the subgoal $`p(\stackrel{}{X})`$, we set the flag $`TB(p(\stackrel{}{X}))COMP=1`$, indicating that the answers of $`p(\stackrel{}{X})`$ have been completed. After memorization, we fetch the next answer from the table and then prove the resolvent $`G_{i+1}`$ of $`G_i`$ and the new answer.
Case 3 is straightforward, so we move to Case 4. By the Table-first policy, we first look up answers for $`A_1`$ from the table $`TB(p(\stackrel{}{X}))`$. When available, we fetch the next unused answer for $`A_1`$ and create the resolvent $`G_{i+1}`$. Otherwise, we check the flag $`TB(p(\stackrel{}{X}))COMP`$ to see if the answers of $`p(\stackrel{}{X})`$ have been completed. If yes, which means that the subgoal $`A_1`$ at $`N_i`$ has used all its answers, we backtrack to its parent node. Otherwise, we continue to derive new answers by resolving $`A_1`$ with the remaining clauses. Based on whether $`A_1`$ has any ancestor variant subgoal, we distinguish three cases (Cases 4.1, 4.2 and 4.3).
At the time that $`N_i`$ is registered (see Definition 3.2), we do not know if $`A_1`$ at $`N_i`$ has any ancestor variant subgoal (i.e. $`N_inode\mathrm{\_}ANC=1`$ initially). So we check it via the ancestor list $`AL_{A_1}`$ (see Definition 3.7) and update the flag $`N_inode\mathrm{\_}ANC`$ accordingly. If $`N_h`$ is the closest ancestor node of $`N_i`$ such that $`L(N_h,N_i)`$ is a loop, the other flags of $`N_i`$, namely $`node\mathrm{\_}LOOP`$, $`node\mathrm{\_}ITER`$ and $`node\mathrm{\_}SUSP`$, will also be updated by the procedure $`nodetype\mathrm{\_}update(.)`$ (see Case 4.1).
For Case 4.2, $`A_1`$ has no ancestor variant subgoal, which implies that the derivation does not get into a loop at $`N_i`$. So we seek the next clause whose head is unifiable with $`A_1`$ and whose status is “still available,” and use it to build the resolvent $`G_{i+1}`$. Now consider the case that no such a clause exists, which means that the subgoal $`A_1`$ at $`N_i`$ has used all its answers and clauses available. In Prolog, we would fail the subgoal immediately and backtrack to its parent node. In TP-resolution, however, we cannot do this unless $`N_i`$ is a non-iteration node. Suppose $`N_i`$ is an iteration node (i.e. $`N_inode\mathrm{\_}ITER=1`$). Before failing $`A_1`$ via backtracking, we do answer iteration to complete its answers. The process is quite simple. We start an iteration simply by initializing $`NEW`$ to 0 and resetting the pointer $`N_iclause\mathrm{\_}ptr`$ to pointing to the first clause $`C_{p_j}`$ whose status remains to “still available” (Case 4.2.2). Since the status of all clauses involved in loops are kept to “still available” during backtracking, all the loops can be regenerated by the iteration. By the end of an iteration, i.e. when we come back to $`N_i`$ again and try to fail $`A_1`$ via backtracking, we check the flag $`NEW`$ to see if the termination condition is satisfied. If $`NEW=0`$, meaning that a fixpoint has been reached so that the answers of $`A_1`$ have been completed, we stop answer iteration by failing $`A_1`$ via backtracking (see Cases 4.2.1). Otherwise, we start a new iteration to seek more answers (Case 4.2.2).
For Case 4.3, $`A_1`$ has an ancestor variant subgoal, so the derivation has gone into a loop, say $`L(N_h,N_i)`$. In order to break the loop, we skip the clause $`C_{p_j}`$ that is being used by the top loop node $`N_h`$. The skip of clauses may lead to loss of answers, which is the only reason why answer iteration is required. (Remark: Algorithm 2 uses loop checking to cut loops and adopts answer iteration to iteratively regenerate loops that are pruned by loop checking. Such a complementary use of loop checking and answer iteration is an effective way of cutting infinite loops while guaranteeing the completeness of answers.)
Backtracking (point 3) is done as usual except that the status of the clauses that have been used should be set to “no longer available” (Case B3). Let $`C_{q_j}`$ be the clause that is being used by $`N_f`$. If no loop occurred that went through $`N_f`$ via $`C_{q_j}`$, the flag $`N_fnode\mathrm{\_}SUSP`$ must remain to 0. In this case, we set the status of $`C_{q_j}`$ in $`TB(A_f)`$ to “no longer available” because all answers of $`A_f`$ by the application of $`C_{q_j}`$ have been exhausted. Otherwise, when a loop occurred before that went through $`N_f`$ via $`C_{q_j}`$, $`N_fnode\mathrm{\_}SUSP`$ must be 1 (see the procedure $`nodetype\mathrm{\_}update(.))`$. So we keep the status of $`C_{q_j}`$ to “still available” while setting $`N_fnode\mathrm{\_}SUSP`$ to 0 again.
Based on TP-trees, we have the following standard definitions.
###### Definition 3.10
Let $`TP_{G_0}`$ be a TP-tree of $`P\{G_0\}`$. All leaves of $`TP_{G_0}`$ labeled by $`\mathrm{}`$ are success leaves and all other leaves are failure leaves. A TP-derivation, denoted by $`G_0_{C_1,\theta _1}G_1\mathrm{}`$ $`_{C_i,\theta _i}G_i\mathrm{}`$ $`_{C_n,\theta _n}G_n`$, is a partial branch in $`TP_{G_0}`$ starting at the root, where each $`G_i`$ is a goal labeling a node $`N_i`$ and for each $`0i<n`$, $`G_{i+1}`$ is the resolvent of $`G_i`$ and $`C_{i+1}`$ with the mgu $`\theta _{i+1}`$, where $`C_{i+1}`$ may be a program clause or a tabled fact or blank (when the left-most subgoal of $`G_i`$ is a procedure). A TP-derivation is successful if it ends with a success leaf and failed, otherwise. The process of constructing TP-derivations is called TP-resolution.
###### Example 3.3
Consider the example program $`P_1`$ again (see Section 2). Based on the dependency graph of $`P_1`$, we choose $`reach`$ as a tabled predicate and $`edge`$ as a non-tabled one. Now consider applying Algorithm 2 to the top goal $`G_0=reach(a,X)`$.
We first set up the root $`N_0`$ labeled with $`reach(a,X),return((X))`$ and set $`NEW=0`$ (point 1). Then we expand $`N_0`$ using the clause $`C_{r_1}`$ (point 2, Cases 4, 4.1 and 4.2), which creates a table
$`TB(reach(a,X)):(reach(a,X),\{\},(1,1,1),0)`$
and a child node $`N_1`$ (see Figure 6). Obviously there is a loop $`L(N_0,N_1)`$, so we call the procedure $`notetype\mathrm{\_}update(L(N_0,N_1))`$, which marks $`N_0`$ as a candidate iteration node, sets $`N_0clause\mathrm{\_}SUSP=1`$ and $`N_1node\mathrm{\_}ANC=1`$. Then by Case 4.3 the clause $`C_{r_2}`$ (instead of $`C_{r_1}`$) is applied to $`reach(a,Z)`$ at $`N_1`$, which gives a node $`N_2`$. Next, by Case 2 the answer $`reach(a,a)`$ is memorized in the table (so $`NEW=1`$), yielding
$`TB(reach(a,X)):(reach(a,X),\{(a)\},(1,1,1),0)`$
and the node $`N_3`$ is derived using the first tabled fact. By successively performing Cases 3, 2 and 1, we reach a success leaf $`N_6`$ with the first answer $`X=a`$ to the top goal. After these steps, the table looks like
$`TB(reach(a,X)):(reach(a,X),\{(a),(b)\},(1,1,1),0)`$.
Now we do backtracking. By Cases B1 and B2 we go back until $`N_3`$. Since $`C_{e_2}`$ is not unifiable with the subgoal $`edge(a,X)`$, we go back to $`N_2`$ and then to $`N_1`$. From $`N_1`$ we consecutively derive a failure leaf $`N_7`$ (Figure 7), a success leaf $`N_{12}`$ (Figure 8) and another failure leaf $`N_{13}`$ (Figure 9). After these steps, the table becomes
$`TB(reach(a,X)):(reach(a,X),\{(a),(b),(d),(e)\},(1,0,0),0)`$.
Now $`reach(a,Z)`$ at $`N_1`$ has used all answers in the table and has no more clause available. So we return to the root $`N_0`$. Note that since the flag $`N_0clause\mathrm{\_}SUSP=1`$, which shows the clause $`C_{r_1}`$ that is being used by $`N_0`$ is involved in a loop, the status of $`C_{r_1}`$ in $`TB(reach(a,X))`$ remains to “still available” when backtracking from $`N_1`$ to $`N_0`$ (see Case B3).
From Figures 6$``$9, we see that $`N_0`$ has used only the first two answers in $`TB(reach(a,X))`$, namely $`reach(a,a)`$ and $`reach(a,b)`$. So it continues to use the other two. By repeating Case 4, Case 1 and point 3 twice, we get another two successful derivations as depicted in Figures 10 and 11.
Now $`reach(a,X)`$ at $`N_0`$ has used all tabled facts in $`TB(reach(a,X))`$ and has no more clause available (note that $`C_{r_2}`$ and $`C_{r_3}`$ have already been used by $`N_1`$). Before failing it via backtracking, we check if $`N_0`$ is an iteration node (i.e. we see if $`N_0node\mathrm{\_}ITER`$ remains to the value 1). Since $`N_0`$ is an iteration node and the flag $`NEW=1`$, by Case 4.2.2 we do answer iteration. It is easy to check that no new answer will be derived (see Figure 4), so by the end of the first iteration $`NEW`$ remains to the value $`0`$. Thus by Case 4.2.1, the flag $`COMP`$ of $`TB(reach(a,X))`$ is changed to 1, showing that the answers of $`reach(a,X)`$ have been completed.
Finally, by point 3 the answer $`NO`$ is returned, which terminates the algorithm. Therefore by putting together Figures 6$``$11 and the figures for answer iteration (which are omitted here) we obtain the TP-tree $`TP_{G_0}`$ of $`P_1\{G_0\}`$.
The following example is also useful in illustrating TP-resolution.<sup>7</sup><sup>7</sup>7This program is suggested by an anonymous referee. To simplify the presentation, in the sequel, in depicting derivations we omit subgoals like $`memo\mathrm{\_}look(.)`$ and $`return(.)`$ unless they are required to be explicitly present.
###### Example 3.4
Consider the logic program
| $`P_2`$: | $`p(a,b,c).`$ $`C_{p_1}`$ |
| --- | --- |
| | $`p(X,Y,Z)p(Z,X,Y).`$ $`C_{p_2}`$ |
Choose $`p`$ as a tabled predicate. Let $`G_0=p(X,Y,Z)`$ be the top goal. The TP-tree of $`P_2\{G_0\}`$ consists of Figures 12 and 13, which yields three answers, $`p(a,b,c)`$, $`p(b,c,a)`$ and $`p(c,a,b)`$.
Note that in the above examples, no new answers are derived during answer iteration (i.e. Algorithm 2 stops by the end of the first iteration). We now give another example, which shows that answer iteration is indispensable.
###### Example 3.5
Consider the following logic program
| $`P_3`$: | $`p(X,Y)q(X,Y).`$ $`C_{p_1}`$ |
| --- | --- |
| | $`q(X,Y)p(X,Z),t(Z,Y).`$ $`C_{q_1}`$ |
| | $`q(a,b).`$ $`C_{q_2}`$ |
| | $`t(b,c).`$ $`C_{t_1}`$ |
Choose $`p`$ and $`q`$ as tabled predicates and apply Algorithm 2 to the top goal $`G_0=p(X,Y).`$ After applying the clauses $`C_{p_1}`$ and $`C_{q_1}`$, we generate the derivation shown in Figure 14. We see that a loop $`L(N_0,N2)`$ occurs. So we do not use $`C_{p_1}`$ to expand $`N_2`$ because that would repeat the loop. Instead, we try alternative clauses. Since there is no other clause in $`P_3`$ that is unifiable with $`p(X,Z)`$, we fail $`N_2`$ and backtrack to its parent node $`N_1`$, which leads to the derivation of Figure 15. Now, since there is no more clause available for $`q(X,Y)`$, we fail $`N_1`$ and go back to $`N_0`$. Note that the flag $`NEW`$ has been set to $`1`$ because new answers, $`q(a,b)`$ and $`p(a,b)`$, have been derived. Moreover, $`C_{q_2}`$ is no longer available to $`q(X,Y)`$, whereas both $`C_{p_1}`$ and $`C_{q_1}`$ are still available because they are involved in a loop.
At $`N_0`$ answer iteration is performed. The first iteration is shown in Figure 16, where two new answers, $`q(a,c)`$ and $`p(a,c)`$, are derived. The second iteration will derive no new answers, so the algorithm stops with the flag $`COMP`$ of $`TB(p(X,Y))`$ set to 1.
## 4 Characteristics of TP-Resolution
In this section, we prove the termination of Algorithm 2 and the soundness and completeness of TP-resolution. We also discuss the way to deal with the cut operator in TP-resolution.
### 4.1 Soundness and Completeness
In order to guarantee termination of Algorithm 2, we restrict ourselves to logic programs with the bounded-term-size property. The following definition is adapted from .
###### Definition 4.1
A logic program $`P`$ has the bounded-term-size property if there is a function $`f(n)`$ such that whenever a top goal $`G_0`$ has no argument whose term size exceeds $`n`$, then no subgoal in the TP-tree $`TP_{G_0}`$ and no answer tuple in any table have an argument whose term size exceeds $`f(n)`$.
Obviously, all function-free logic programs have the bounded-term-size property.
###### Theorem 4.1 (Termination)
Let $`P`$ be a logic program with the bounded-term-size property and $`G_0`$ a top goal. Algorithm 2 terminates with a finite TP-tree $`TP_{G_0}`$.
The following lemma is required to prove this theorem.
###### Lemma 4.2
Let $`G_i`$ and $`G_k`$ be two goals in a TP-derivation of $`P\{G_0\}`$ and $`A_i`$ and $`A_k`$ be the left-most subgoals of $`G_i`$ and $`G_k`$, respectively. If $`A_i`$ is an ancestor variant subgoal of $`A_k`$ then $`A_i`$ is a tabled subgoal.
Proof. Let $`A_i=p(.)`$. By Definitions 3.7 and 3.8, $`A_i`$ being an ancestor variant subgoal of $`A_k`$ implies that there is a cycle of the form $`p\mathrm{}p`$ in the dependency graph $`G_P`$. So $`p`$ is a tabled predicate and thus $`A_i`$ is a tabled subgoal. $`\mathrm{}`$
Proof of Theorem 4.1. Assume, on the contrary, that Algorithm 2 does not terminate. Then it generates an infinite TP-tree. This can occur only in two cases: (1) it memorizes infinitely many (new) answers in tables, so we do backtracking at some nodes infinite times; and (2) it traps into an infinite derivation. We first show that the first case is not possible. Since $`P`$ has the bounded-term-size property, all tabled facts have finite term size. Then, in view of the fact that any logic program has only a finite number of predicate, function and constant symbols, all tabled facts having finite term size implies that any table has only a finite number of tabled facts.
We now assume the second case. Since $`P`$ has the bounded-term-size property and contains only a finite number of clauses, any infinite derivation must contain an infinite loop, i.e. an infinite set of subgoals, $`A_0,A_1,\mathrm{},A_k,\mathrm{},`$ such that for any $`i0`$, $`A_i`$ is both an ancestor subgoal and a variant of $`A_{i+1}`$. This means that all the $`A_i`$s are tabled subgoals (Lemma 4.2). However, from Cases 4 and 4.3 of Algorithm 2 we see that such a set of subgoals will never be generated unless $`P`$ contains an infinite set of clauses whose heads are unifiable with the $`A_i`$s, a contradiction. $`\mathrm{}`$
To simplify the proof of soundness and completeness, we assume, in the sequel, that all predicates are tabled predicates.
###### Theorem 4.3 (Soundness and Completeness)
Let $`P`$ be a logic program with the bounded-term-size property and $`G_0=A_1,\mathrm{},A_m`$ a top goal with the list $`\stackrel{}{Y}`$ of variables. Let $`TP_{G_0}`$ be the TP-tree of $`P\{G_0\}`$ and $`ST_{G_0}`$ the SLD-tree of $`P\{G_0\}`$ via the left-most computation rule. Then $`TP_{G_0}`$ and $`ST_{G_0}`$ have the same set of answers to $`G_0`$.
Proof. (Soundness) By the PMF mode, each tabled fact is an intermediate answer of some subgoal (called sub-refutation in ) in an SLD-derivation in $`ST_{G_0}`$. Since the answer $`\stackrel{}{I}`$ returned at any success leaf in $`TP_{G_0}`$ is an instance of $`\stackrel{}{Y}`$ such that each $`A_i\stackrel{}{Y}/\stackrel{}{I}`$ is an instance of a tabled fact, $`\stackrel{}{Y}/\stackrel{}{I}`$ must be the answer substitution of some successful SLD-derivation in $`ST_{G_0}`$.
(Completeness) Algorithm 2 works in the same way as Algorithm 1 (i.e. it expands and backtracks on nodes in the same way as Algorithm 1) except (1) it is based on the PMF mode, (2) after finishing backtracking for answers of a subgoal $`A_f=q(.)`$ through the application of a clause $`C_{q_j}`$, the status of $`C_{q_j}`$ w.r.t. $`A_f`$ will be set to “no longer available” (see Case B3), and (3) loops are handled by skipping repeated clauses and doing answer iteration. Since the PMF mode preserves the answers of SLD-resolution and point (2) is only for the purpose of avoiding redundant computations (i.e. when variant subgoals of $`A_f`$ later occur, they will directly use the tabled answers instead of recomputing them by applying $`C_{q_j}`$), it suffices to prove that point (3) does not lose any answers to $`G_0`$.
Let $`SD`$ be an arbitrary successful SLD-derivation in $`ST_{G_0}`$ with loops as shown in Figure 17, where $`m>0`$, $`N_{l_0}`$ is an iteration node and for any $`0i<m`$ $`p(\stackrel{}{X_i})`$ is an ancestor variant subgoal of $`p(\stackrel{}{X}_{i+1})`$. Note that the SLD-derivation starts looping at $`N_{l_1}`$ by applying $`C_{p_j}`$ to $`p(\stackrel{}{X_1})`$. However, Algorithm 2 will handle such loops by skipping $`C_{p_j}`$ at $`N_{l_1}`$ and doing answer iteration at $`N_{l_0}`$. Before showing that no answers to $`p(\stackrel{}{X_0})`$ will be lost using the skipping-iterating technique, we further explain the structure of the loops in $`SD`$ as follows.
1. For $`0i<m`$ from $`N_{l_i}`$ to $`N_{l_{i+1}}`$ the proof of $`p(\stackrel{}{X_i})`$ reduces to the proof of $`(p(\stackrel{}{X}_{i+1}),B_{i+1})`$ with a substitution $`\theta _i`$ for $`p(\stackrel{}{X_i})`$, where each $`B_k`$ $`(0km)`$ is a set of subgoals.
2. The sub-refutation between $`N_{l_m}`$ and $`N_{x_m}`$ contains no loops and yields an answer $`p(\stackrel{}{X}_m)\gamma _m`$ to $`p(\stackrel{}{X}_m)`$. The answer substitution $`\gamma _m`$ for $`p(\stackrel{}{X}_m)`$ is then applied to the remaining subgoals of $`N_{l_m}`$ (see node $`N_{x_m}`$), which leads to an answer $`p(\stackrel{}{X}_{m1})\gamma _m\gamma _{m1}\theta _{m1}`$ to $`p(\stackrel{}{X}_{m1})`$. Such process continues recursively until an answer $`p(\stackrel{}{X_0})\gamma _m\mathrm{}\gamma _0\theta _{m1}\mathrm{}\theta _0`$ to $`p(\stackrel{}{X_0})`$ is produced at $`N_{x_0}`$.
We now prove that a variant of the answer $`p(\stackrel{}{X_0})\gamma _m\mathrm{}\gamma _0\theta _{m1}\mathrm{}\theta _0`$ to $`p(\stackrel{}{X_0})`$ will be produced by Algorithm 2 by means of answer iteration.
Since $`p(\stackrel{}{X}_0)`$ and $`p(\stackrel{}{X}_m)`$ are variants, via backtracking from $`N_{l_1}`$ up to $`N_{l_0}`$ a variant of the sub-refutation between $`N_{l_m}`$ and $`N_{x_m}`$ can be generated, which starts from $`N_{l_0}`$ via $`C_{p_j}`$. This means that a variant of the answer $`p(\stackrel{}{X}_m)\gamma _m`$ to $`p(\stackrel{}{X}_m)`$ can be derived via backtracking from $`N_{l_1}`$ up to $`N_{l_0}`$, independently of the sub-derivation below $`N_{l_1}`$.
Let us do backtracking from $`N_{l_1}`$ up to $`N_{l_0}`$ and store all intermediate answers in their tables. So $`p(\stackrel{}{X}_m)\gamma _m`$ is in $`TB(p(\stackrel{}{X}_0))`$. Now we regenerate the loop $`L(N_{l_0},N_{l_1})`$ (the first iteration).
Since $`p(\stackrel{}{X}_0)`$ and $`p(\stackrel{}{X}_{m1})`$ are variants, a variant of the sub-refutation between $`N_{l_{m1}}`$ and $`N_{x_{m1}}`$, where the sub-refutation between $`N_{l_m}`$ and $`N_{x_m}`$ is replaced by directly using the answer $`p(\stackrel{}{X}_m)\gamma _m`$, can be generated via backtracking from $`N_{l_1}`$ up to $`N_{l_0}`$. That is, a variant of the answer $`p(\stackrel{}{X}_{m1})\gamma _m\gamma _{m1}\theta _{m1}`$ to $`p(\stackrel{}{X}_{m1})`$ can also be derived via backtracking from $`N_{l_1}`$ up to $`N_{l_0}`$ when the tabled answer $`p(\stackrel{}{X}_m)\gamma _m`$ is used. So we do the backtracking, store $`p(\stackrel{}{X}_{m1})\gamma _m\gamma _{m1}\theta _{m1}`$ in $`TB(p(\stackrel{}{X}_0))`$, and then regenerate the loop $`L(N_{l_0},N_{l_1})`$ (the second iteration).
Continue the above process recursively. After (at most) $`m`$ iterations, a variant of the answer $`p(\stackrel{}{X_0})\gamma _m\mathrm{}\gamma _0\theta _{m1}\mathrm{}\theta _0`$ to $`p(\stackrel{}{X_0})`$ will be derived and stored in $`TB(p(\stackrel{}{X}_0))`$.
The above arguments show that although the branch below $`N_{l_1}`$ via $`C_{p_j}`$ is skipped by Algorithm 2, by means of answer iteration along with tabling no answers will be lost to $`p(\stackrel{}{X_0})`$. Therefore, when a fixpoint is reached at $`N_{l_0}`$, which means no new answers to $`p(\stackrel{}{X}_0)`$ can be derived via iterations, all answers of $`p(\stackrel{}{X}_0)`$ must be exhausted and stored in $`TB(p(\stackrel{}{X}_0))`$ (in such a case, the flag $`TB(p(\stackrel{}{X}_0))COMP`$ is set to 1). We now prove that the fixpoint can be reached in finite time even if $`m\mathrm{}`$.
Let $`m\mathrm{}`$. Then $`SD`$ contains infinite loops. Since $`P`$ has the bounded-term-size property and only a finite number of clauses, we have only a finite number of subgoals and any subgoal has only a finite number of answers (up to variable renaming). Let $`N`$ be the number of all answers of all subgoals. Since before the fixpoint is reached, in each iteration at $`N_{l_0}`$ at least one new answer to some subgoal will be derived, the fixpoint will be reached after at most $`N`$ iterations.
To sum up, Algorithm 2 traverses $`ST_{G_0}`$ as follows: For any SLD-derivation $`SD`$ in $`ST_{G_0}`$, if it has no loops Algorithm 2 will generate it based on the PMF mode while removing redundant application of clauses; otherwise, Algorithm 2 will derive the answers of subgoals involved in the loops by means of answer iteration. In either case, Algorithm 2 terminates and preserves the answers of SLD-resolution. As a result, if $`SD`$ is successful with an answer to $`G_0`$, there must be a successful TP-derivation in $`TP_{G_0}`$ with the same answer (up to variable renaming). $`\mathrm{}`$
### 4.2 Dealing with Cuts
The cut operator, !, is very popular in Prolog programming. It basically serves two purposes. One is to simulate the if-then-else statement, which is one of the key flow control statements in procedural languages. For example, in order to realize the statement if-$`A`$-then-$`B`$-else-$`C`$, we define the following:
| | $`HA,!,B`$. |
| --- | --- |
| | $`HC`$. |
The other, perhaps more important, purpose of using cuts is to prune the search space by aborting further exploration of some remaining branches, which may lead to significant computational savings. For instance, the following clauses
| | $`p(\stackrel{}{X})A_1,\mathrm{},A_m,!`$. |
| --- | --- |
| | $`C_{p(.)}`$: the remaining clauses defining $`p(.)`$. |
achieve the effect that for any $`\stackrel{}{X}`$ whenever $`A_1,\mathrm{},A_m`$ is true with an mgu $`\theta `$, we return $`p(\stackrel{}{X})\theta `$ and stop searching the remaining space (via backtracking on the $`A_i`$s and using the remaining clauses $`C_{p(.)}`$) for any additional answers of $`p(\stackrel{}{X})`$.
The cut operator requires a strictly sequential strategy $``$ Prolog-strategy for the selection of goals, subgoals and program clauses. TP-strategy is an enhancement of Prolog-strategy with the following two policies for dealing with tabled facts (see Definition 3.6): Table-first when both tabled facts and program clauses are available and Top-down for the selection of tabled facts. Since new answers are always appended to the end of tables, by the PMF mode, such an enhancement does not affect the original sequentiality of Prolog-strategy. That is, TP-strategy supports the cut operator as well.
Before enhancing Algorithm 2 with mechanisms for handling cuts, we recall the operational semantics of cuts.
###### Definition 4.2
Let $`P`$ be a logic program that contains the following clauses with a head $`p(.)`$:
| $`p(.)`$ | $`\mathrm{}.`$ $`C_{p_1}`$ |
| --- | --- |
| | |
| $`p(\stackrel{}{Y})B_1,\mathrm{},B_m,!,B_{m+2},\mathrm{}.,B_{m+k}`$ $`C_{p_i}`$ |
| | |
| $`p(.)\mathrm{}.`$ $`C_{p_n}`$ |
Let $`p(\stackrel{}{X})`$ be a subgoal such that $`p(\stackrel{}{X})\theta =p(\stackrel{}{Y})\theta `$. The semantics of ! in $`C_{p_i}`$ is defined as follows: During top-down evaluation of $`p(\stackrel{}{X})`$, by the left-most computation rule whenever $`(B_1,\mathrm{},B_m)\theta `$ succeeds with an mgu $`\theta _1`$, all the remaining answers to the subgoal $`p(\stackrel{}{X})`$ are obtained by computing $`(B_{m+2},\mathrm{}.,B_{m+k})\theta \theta _1`$, with the backtracking on the $`B_j`$s ($`1jm`$) and the remaining clauses $`C_{p_j}`$s ($`i<jn`$) ignored. In other words, we force two skips when backtracking on the cut: the skip of all $`B_j`$s ($`1jm`$) and the skip of all $`C_{p_j}`$s ($`i<jn`$).
It is quite easy to realize cuts in TP-resolution. Let $`N_h`$ be a node labeled by a goal
$`G_h=p(\stackrel{}{X}),\mathrm{}`$
and the clauses for $`p(.)`$ be as in Definition 4.2. Let
$`G_{h+1}=(B_1,\mathrm{},B_m)\theta ,!,(B_{m+2},\mathrm{},B_{m+k})\theta ,\mathrm{}`$
be the resolvent of $`G_h`$ and $`C_{p_i}`$. When evaluated as a subgoal for forward node expansion, $`!`$ is unconditionally true. However, during backtracking, by Definition 4.2 it will skip all $`B_j`$s by directly jumping back to the node $`N_h`$. In order to formalize such a jump, we attach to the subgoal $`!`$ a node name $`N_h`$ as a directive for backtracking. That is, we create a subgoal $`!(N_h)`$, instead of $`!`$, in the resolvent $`G_{h+1}`$.
Then cuts are realized in TP-resolution simply by adding to Algorithm 2, before Case 1 in point 2, the case
* Case 0: $`A_1`$ is $`!(N_h)`$. Set $`G_{i+1}=A_2,\mathrm{},A_m`$ and $`N_hnode\mathrm{\_}SUSP=0`$, and goto 2.
and, before Case B1 in point 3, the case
* Case B0: $`A_f`$ is $`!(N_h)`$. Let $`A_h=p(\stackrel{}{X})`$ be the left-most subgoal at $`N_h`$ and $`C_{p_i}`$ be the clause that is being used by $`A_h`$. If $`A_h`$ is a non-tabled subgoal then goto 3 with $`N=N_h`$. Otherwise, if $`N_hnode\mathrm{\_}SUSP=0`$ then set $`TB(A_h)clause\mathrm{\_}status[j]=0`$ for all $`ji`$; else set $`N_hnode\mathrm{\_}SUSP=0`$ and $`N_hclause\mathrm{\_}ptr=null`$. Goto $`2`$ with $`N_h`$ as the latest registered node.
For Case 0, since $`!`$ is unconditionally true, $`G_{i+1}=A_2,\mathrm{},A_m`$. For Case B0, we do backtracking on the subgoal $`!(N_h)`$ at node $`N_f`$. By Definition 4.2, we will skip all nodes used for evaluating $`(B_1,\mathrm{},B_m)\theta `$ and then skip all clauses $`C_{p_j}`$s with $`j>i`$. The first skip is done by jumping from $`N_f`$ back to $`N_h`$. If $`p(\stackrel{}{X})`$ at $`N_h`$ is a non-tabled subgoal, the second skip is done by failing the subgoal via backtracking. Otherwise, we consider two cases.
1. Assume $`N_hnode\mathrm{\_}SUSP=0`$. This means the evaluation of $`(B_{m+2},\mathrm{}.,B_{m+k})\theta \theta _1`$ did not encounter any loop that goes through $`N_h`$ via $`C_{p_i}`$, so that all answers of $`(B_{m+2},\mathrm{}.,B_{m+k})\theta \theta _1`$ must have been exhausted via backtracking. Thus there will be no new answers of $`p(\stackrel{}{X})`$ that can be derived by applying the clauses $`C_{p_j}`$s $`(ji)`$. Therefore, in this case the second skip is achieved by changing the status of the $`C_{p_j}`$s in $`TB(p(\stackrel{}{X}))`$ to “no longer available”.
2. Assume $`N_hnode\mathrm{\_}SUSP=1`$. Since the flag $`N_hnode\mathrm{\_}SUSP`$ is initialized to 0 after the evaluation of $`(B_1,\mathrm{},B_m)\theta `$ (see Case 0), $`N_hnode\mathrm{\_}SUSP=1`$ means that the evaluation of $`(B_{m+2},\mathrm{}.,B_{m+k})\theta \theta _1`$ encountered loops that go through $`N_h`$ via $`C_{p_i}`$. So answer iteration is required to exhaust the answers of $`(B_{m+2},\mathrm{}.,B_{m+k})\theta \theta _1`$. Hence, in this case the second skip is done simply by clearing the pointer $`N_hclause\mathrm{\_}ptr`$, so that no more clauses will be available to $`p(\stackrel{}{X})`$ at $`N_h`$.
###### Example 4.1
Consider the following logic program:
| $`P_4`$: | $`p(X,Y)p(X,Z),t(Z,Y).`$ $`C_{p_1}`$ |
| --- | --- |
| | $`p(X,Y)p(X,Y),!.`$ $`C_{p_2}`$ |
| | $`p(a,b).`$ $`C_{p_3}`$ |
| | $`p(f,g).`$ $`C_{p_4}`$ |
| | $`t(b,c).`$ $`C_{t_1}`$ |
Choose $`p`$ as a tabled predicate. Let $`G_0=p(X,Y)`$ be the top goal. By applying $`C_{p_1}`$ to the root $`N_0`$ we generate $`N_1`$, where the first loop $`L(N_0,N_1)`$ occurs (see Figure 18). Then $`C_{p_2}`$ is applied, which yields the second loop $`L(N_1,N_2)`$. Since $`C_{p_2}`$ is being used by $`N_1`$, $`C_{p_3}`$ is used to expand $`N_2`$, which gives the first tabled fact $`p(a,b)`$. At $`N_3`$, the cut succeeds unconditionally, which leads to $`N_4`$. Then $`C_{t_1}`$ is applied, giving the first success leaf $`N_5`$ with the second tabled fact $`p(a,c)`$ added to $`TB(p(X,Y))`$.
We backtrack to $`N_4`$ and then to $`N_3`$. Due to the subgoal $`!(N_1)`$, we directly backtrack to $`N_1`$ (the first skip). The status of $`C_{p_2}`$, $`C_{p_3}`$ and $`C_{p_4}`$ in $`TB(p(X,Y))`$ is then changed to “no longer available” (the second skip). At $`N_1`$, the second tabled fact $`p(a,c)`$ is used, which yields a failure leaf $`N_6`$. Next we go back to $`N_0`$, where the second tabled fact $`p(a,c)`$ is used, which gives the second success leaf $`N_7`$.
Similar extension can be made to Algorithm 1 to deal with cuts in Prolog. By comparison of the two, we see that without loops, cuts in TP-resolution achieve the same effect as in Prolog. When there are loops, however, TP-resolution still reaches conclusions, whereas Prolog will never stop. The following representative example illustrates such a difference.
###### Example 4.2
The following two clauses
| $`not\mathrm{\_}p(X)p(X),!,fail.`$$`C_{np_1}`$ |
| --- |
| $`not\mathrm{\_}p(X).`$ $`C_{np_2}`$ |
define the predicate $`not\mathrm{\_}p`$ which says that for any object $`X`$, $`not\mathrm{\_}p(X)`$ succeeds if and only if $`p(X)`$ fails. Let $`G_0=not\mathrm{\_}p(a)`$ be the top goal and the programs $`P_{5_i}`$ be defined as follows.
1. $`P_{5_1}=\{C_{np_1},C_{np_2}\}`$. As $`p(a)`$ fails, $`C_{np_2}`$ is applied, so that both Prolog and Algorithm 2 give an answer $`YES`$ to $`G_0`$.
2. $`P_{5_2}=\{C_{np_1},C_{np_2},p(a)\}`$. As $`p(a)`$ succeeds, the cut $`!`$ in $`C_{np_1}`$ is executed. Since the subgoal $`fail`$ always fails, the backtracking on $`!`$ skips $`C_{np_2}`$, so that both Prolog and Algorithm 2 give an answer $`NO`$ to $`G_0`$.
3. $`P_{5_3}=\{C_{np_1},C_{np_2},p(X)p(X)\}`$. Note that $`p`$ is a tabled predicate. As Prolog goes into an infinite loop in proving the subgoal $`p(a)`$, no answer to $`G_0`$ can be obtained. However, Algorithm 2 breaks the loop by deriving a negative answer to $`p(a)`$, so that $`C_{np_2}`$ is applied, which leads to an answer $`YES`$ to $`G_0`$.
As we mentioned earlier, cuts are used for two main purposes: (1) simulate the if-$`A`$-then-$`B`$-else-$`C`$ statement, i.e. treat $`B`$ and $`C`$ to be two exclusive objects; (2) prune the search space, i.e. force the two skips when backtracking on cuts (see Definition 4.2). Since the second purpose exactly corresponds to the operational semantics of cuts, it is achieved by both Prolog and TP-resolution in any situations. It turns out, however, that the first purpose cannot be achieved in arbitrary situations. The following example illustrates this.
###### Example 4.3
Consider the following logic program:
| $`P_6`$: | $`p(X)q(X),p(b),!,B`$. $`C_{p_1}`$ |
| --- | --- |
| | $`p(X)C`$. $`C_{p_2}`$ |
| | $`q(a).`$ $`C_{q_1}`$ |
| | $`B.`$ $`C_{B_1}`$ |
| | $`C.`$ $`C_{C_1}`$ |
It is easy to check that this program will generate no loops. However, the two clauses $`C_{p_1}`$ and $`C_{p_2}`$ do not represent
if $`q(X)`$ and $`p(b)`$ then $`B`$ else $`C`$
because evaluating $`p(X)`$ by Prolog/TP-resolution will lead to both $`C`$ and $`B`$ being executed, which violates the intension that they are exclusive objects.
###### Definition 4.3
Let $`P`$ be a program. We say that the effect of if-A-then-B-else-C is achieved using clauses of the form
| | $`HA,!,B`$. |
| --- | --- |
| | $`HC`$. |
if when evaluating $`H`$ against $`P`$, either $`B`$ (i.e. when $`A`$ is true) or $`C`$ (i.e. when $`A`$ is false) but not both will be executed.
Based on this criterion, we give the following characterizations of the classes of programs for which cuts are effectively handled by Prolog/TP-resolution to achieve the effect of if-A-then-B-else-C.
###### Theorem 4.4
Let $`P`$ be a program with the bounded-term-size property. Let $`A=A_1,\mathrm{},A_m`$, $`B=B_1,\mathrm{},B_n`$ and $`C=C_1,\mathrm{},C_q`$. TP-resolution achieves the effect of if-$`A`$-then-$`B`$-else-$`C`$ using the following clauses in $`P`$
| | $`HA,!,B`$. |
| --- | --- |
| | $`HC`$. |
if and only if (1) if $`A`$ is true with the first answer substitution $`\theta `$ then the evaluation of $`A`$ for the first answer and the evaluation of $`B\theta `$ will not invoke $`C`$; (2) if $`A`$ is false then the evaluation of $`A`$ and the evaluation of $`C`$ will not invoke $`B`$.
Proof. $`()`$ Straightforward.
$`()`$ Since TP-resolution always terminates, the truth value ($`true`$ or $`false`$) of $`A`$ can be definitely determined. So, for point (1), $`B\theta `$ will be executed with $`C`$ excluded; and for point (2), $`C`$ will be executed with $`B`$ excluded. Therefore, the effect of if-$`A`$-then-$`B`$-else-$`C`$ is achieved. $`\mathrm{}`$
###### Theorem 4.5
The conditions of Prolog achieving the effect of if-$`A`$-then-$`B`$-else-$`C`$ using the following clauses
| | $`HA,!,B`$. |
| --- | --- |
| | $`HC`$. |
are the two conditions for TP-resolution plus a third one: (3) the evaluation of $`A`$ for its first answer will not go into a loop.
Proof. Without loops in evaluating $`A`$ for its first answer, the truth value ($`true`$ or $`false`$) of $`A`$ can be definitely determined. Otherwise, neither $`B`$ nor $`C`$ will be executed, which violates the criterion of Definition 4.3. $`\mathrm{}`$
By Theorem 4.4, for programs $`P_{5_1}`$, $`P_{5_2}`$ and $`P_{5_3}`$ (see Example 4.2) the two clauses $`C_{np_1}`$ and $`C_{np_2}`$ can be used by TP-resolution to represent if-$`p(X)`$-then-$`fail`$-else-$`true`$. By Theorem 4.5, however, Prolog cannot achieve such effect for $`P_{5_3}`$ because the evaluation of $`p(X)`$ will go into a loop. Moreover, neither TP-resolution nor Prolog can use $`C_{p_1}`$ and $`C_{p_2}`$ in $`P_6`$ (see Example 4.3) to represent if-($`q(X)`$ and $`p(b))`$-then-$`B`$-else-$`C`$ because the evaluation of $`p(b)`$ will invoke $`C`$, which violates point (1) of Theorem 4.4.
Summarizing the above discussion leads to the following conclusion.
###### Corollary 4.6
Let $`P`$ be a program with the bounded-term-size property. If Prolog effectively handles cuts for $`P`$ w.r.t. the two intended purposes, so does TP-resolution; but the converse is not true w.r.t. the first purpose.
Proof. The second purpose of using cuts is achieved by both Prolog and TP-resolution for any programs. For the first purpose, this corollary follows immediately from Theorems 4.4 and 4.5. $`\mathrm{}`$
## 5 Conclusions and Further Work
Existing tabulated resolutions, such as OLDT-resolution, SLG-resolution and Tabulated SLS-resolution, rely on the solution-lookup mode in formulating tabling. Because lookup nodes are not allowed to resolve tabled subgoals against program clauses, the underlying tabulated resolutions cannot be linear, so that it is impossible to implement such resolutions using a simple stack-based memory structure like that in Prolog. This may make their implementation much more complicated (SLG-WAM for XSB is a typical example , in contrast to WAM/ATOAM for Prolog ). Moreover, because lookup nodes totally depend on solution nodes, without any autonomy, it may be difficult to handle some strictly sequential operators such as cuts as effectively as in Prolog ().
In contrast, TP-resolution presented in this paper has the following novel properties.
1. It does not distinguish between solution and lookup nodes. Any nodes can resolve tabled subgoals against program clauses as well as answers in tables provided that they abide by the Table-first policy, regardless of when and where they are generated.
2. It makes linear tabulated derivations based on TP-strategy in the same way as Prolog except that infinite loops are broken and redundant computations are reduced. The resolution algorithm (Algorithm 2) is sound and complete for positive logic programs with the bounded-term-size property and can be implemented by an extension to any existing Prolog abstract machines such as WAM or ATOAM .
3. Due to its linearity, cuts can be easily realized. It handles cuts as effectively as Prolog in the case that cuts are used for pruning the search space, and better than Prolog in the case for simulating the if-then-else statement.
However, TP-resolution has some disadvantages. In particular, an efficient implementation requires further investigation of the following issues.
1. Because it is a mixture of loop checking and tabling, ancestor checking is required to see if a TP-derivation has gone into a loop. That could be costly. Therefore, fast ancestor checking algorithms remain to be explored in further investigation.
2. Answer iteration introduces redundant computations for those programs and goals where the iteration is totally redundant (see, for example, the programs $`P_1`$ and $`P_2`$ in Examples 3.3 and 3.4 where no new answers can be derived through the iteration). Methods of determining in what cases answer iteration can be ignored remain an interesting open problem.
We have recently extended TP-resolution to compute the well-founded semantics of general logic programs. A preliminary report on the extension appears in . We are also working on the implementation of TP-resolution to realize a linear tabulated Prolog system.
Acknowledgements
We are grateful to the three anonymous referees for their insightful comments, which have greatly improved the presentation. The first author is supported in part by Chinese National Natural Science Foundation and Trans-Century Training Program Foundation for the Talents by the Chinese Ministry of Education.
|
warning/0003/hep-th0003111.html
|
ar5iv
|
text
|
# 1 Introduction
## 1 Introduction
Let $`G`$ be a simple compact Lie group, of dimension $`r`$ and rank $`l`$, of Lie algebra $`𝔤`$ defined by
$$[X_a,X_b]=iC_{abc}X_c,$$
(1)
where $`a,b,c=1,2,\mathrm{},r`$, and the generators $`X_a`$ are hermitian. We wish to examine here the realisation $`X_aS_a`$ of $`𝔤`$ given by
$$S_a=\frac{1}{4}iC_{abc}\gamma _b\gamma _c,$$
(2)
where the $`\gamma _a`$ are the hermitian Dirac matrices
$$\{\gamma _a,\gamma _b\}=2\delta _{ab}$$
(3)
of a euclidean space of dimension $`r\text{dim}𝔤`$. We wish to describe exactly which representation $`𝒟`$ of $`𝔤`$ is provided by (2), and how the vector space $`𝒱_𝒟`$ which carries it may be constructed explicitly. We shall address first (and more fully) the even case in which both $`r`$ (and hence the rank $`l`$) are even, contrasting it afterwards with the somewhat different, but also allowed and interesting, odd case in which $`r`$ (and $`l`$) are odd.
We have a general reason for undertaking this study and one specific class of applications immediately in mind. The reason is the following: the Fock space of any quantum mechanical system in which there exists $`\text{dim}𝔤`$ fermions transforming according to $`G`$ is, necessarily, the carrier space of the representation $`𝒟`$ under investigation here. The specific application that provides much of our motivation is the study of the hidden supercharges $`Q_s`$ in $`G`$-invariant supersymmetric quantum mechanical systems involving a set of exactly $`\text{dim}𝔤`$ fermions. There is one such $`Q_s`$ for each Lie algebra cohomology cocycle of $`𝔤`$ and its construction involves the associated completely antisymmetric invariant tensors fully contracted with the fermionic variables. These matters will be treated in a forthcoming publication.
The first result of this paper (Sec. 2) is a rather intriguing one: in the even case, $`𝒟`$ is the direct sum of $`2^{{\scriptscriptstyle \frac{l}{2}}}`$ copies of the irreducible representation (irrep) of $`𝔤`$ whose highest weight $`\mathrm{\Lambda }`$ is equal to the Weyl principal vector $`\delta =(1,\stackrel{l}{},1)`$. Here the Weyl principal vector of $`𝔤`$, equal to half the sum of the $`(rl)/2`$ positive roots of $`𝔤`$, has been referred to a basis of the fundamental dominant weights of $`𝔤`$. Below, in Sec. 2, we establish this result, giving a systematic description of $`𝒱_𝒟`$ as a fermionic Fock space.
Sec. 3 gives full treatment of the $`𝔤=A_2=su(3)`$ case, which clarifies many of the issues typical of the even case. In particular, we exhibit the role of both the chirality operator and of the Kostant fermionic $`SU(3)`$-invariant cubic operator $`K_3`$. This operator is defined by
$$K_3=\frac{1}{12}iC_{abc}\gamma _a\gamma _b\gamma _c=\frac{1}{3}\gamma _aS_a.$$
(4)
For any $`𝔤`$ of dimension larger than 3, however, $`K_3`$ is not the only relevant fermionic $`G`$-invariant operator. Indeed, for any $`𝔤`$ of rank $`l`$ we may introduce $`l`$ fermionic operators by using its $`l`$ primitive cocycles. These cocycles are given by skewsymmetric tensors of odd dimension $`2m_i1`$ ($`i=1,\mathrm{},l`$) where, for each $`i`$, $`m_i`$ is the order of the symmetric $`G`$-invariant polynomial giving the corresponding primitive Casimir-Racah operator. The structure constants used in (4) simply define the lowest order operator $`K_3`$, for which the ($`m_1`$=2)-order tensor is the Killing metric giving the quadratic Casimir. The higher order Casimirs and their associated ($`2m_i1`$)-cocycles for a simple compact $`𝔤`$ are all known, and known to be relevant in many areas of physics, as in the mathematical description of anomalies (see ), current algebras and Schwinger terms (see e.g. , and references therein), Wess-Zumino terms and effective actions (), $`W`$-algebras (), principal chiral models (see and earlier references therein), and others. For $`𝔤=su(n)`$, for example, the Casimir-Racah tensors of order $`m_i=2,\mathrm{},n1`$ give cocycles of orders $`3,5,\mathrm{},(2n1)`$. Writing $`C_{abc}=f_{abc}`$ for $`su(n)`$, the five-cocycle is determined by the third-order invariant tensor of coordinates $`d_{abc}`$ through
$$\mathrm{\Omega }_{abcde}=f_{xa[b}f_{cd]y}d_{xye},$$
(5)
and is tabulated in for $`n=3`$ and $`n=4`$. Such a quantity can be seen e.g. as a Schwinger term in a two-dimensional chiral $`SU(n)\times SU(n)`$ model: see eq. (28) of and . Using (5) we may then construct, for some $`k`$
$$K_5=\frac{k}{5!}\mathrm{\Omega }_{abcde}\gamma _a\gamma _b\gamma _c\gamma _d\gamma _e.$$
(6)
As we shall see, $`K_5`$ plays a non-trivial role in the understanding of the representation (1) already for $`su(3)`$, even though, in this case, $`K_5`$ is related to $`K_3`$ by Hodge duality as is proved in Sec. 3 and discussed in Sec. 6. For a given $`𝔤`$, higher order $`K_{(2m_i1)}`$ operators may be constructed similarly using the corresponding $`\mathrm{\Omega }_{(2m_i1)}`$ cocycles provided that the rank $`l`$ is high enough. They all have odd character since they are given by skewsymmetric tensors in the basis of the $`r`$ different anticommuting $`\gamma `$’s.
To develop a deeper view, Sec. 4 is devoted to analyse the next smallest even case, that of $`𝔤=su(5)`$, for which there are four fermionic scalars $`K_3,K_5,K_7,K_9`$, one for each of the four primitive cocycles of $`su(5)`$. The higher cocycles give rise to operators which feature in any study of the system essentially on the same footing as the cubic Kostant operator (4). With the aid of some MAPLE programs, we are able to analyse various aspects of the role played by all the fermionic scalar operators $`K_{(2m_i1)}`$.
Our treatment will make this fermionic character explicit by replacing $`\gamma `$ matrices, two at a time, as e.g. in (20) below, by Dirac fermions $`A`$ such that $`\{A,A^{}\}=1`$. This is easy to do for $`r`$ even. A similar approach to the $`\text{dim}𝔤=r`$ odd case is evidently complicated by the fact that it leaves over, unpaired, the last matrix $`\gamma _r`$ or, put otherwise, a single Majorana fermion type entity. This does not mean that the odd case cannot be treated in a satisfactory way. It is however essential in treating the odd case to do so in a fashion that respects fully the fermionic nature of the unpartnered Majorana fermion, and also of the $`K_3`$ operator and the above generalisations. With this in mind, we treat in Sec.5 the $`\mathrm{dim}𝔤=r=3`$ case of $`su(2)`$ in detail to indicate how to handle the general odd case. In fact, when $`\mathrm{dim}𝔤`$ is odd, $`𝒟`$ involves the irrep of $`𝔤`$ with highest weight $`\delta `$ repeated in direct sum $`2^{{\scriptscriptstyle \frac{l1}{2}}}`$ times.
We have referred to $`K_3`$ above as the Kostant operator. Strictly speaking, however, the operator $`K`$ recently introduced by Kostant (see also ) contains, besides the cubic term $`K_3`$, an additional representation dependent piece that we ignore here by restricting our attention to the purely geometrical part $`K_3`$ of $`K`$. Operators such as $`K_3`$, however, are already familiar in non-relativistic supersymmetric quantum mechanics of particles with spin-$`\frac{1}{2}`$ (see e.g. ) or colour degrees of freedom (see e.g. ). In such theories there are Majorana fermion variables $`\psi _i`$ with anticommutation relations
$$\{\psi _i,\psi _j\}=\delta _{ij},$$
(7)
for which there is a representation $`\psi _i\gamma _i/\sqrt{2}`$ in terms of hermitian Dirac matrices. The supercharges of such theories contain or consist of a term proportional to
$$iC_{ijk}\psi _i\psi _j\psi _k.$$
(8)
In the case of $`su(2)`$, where $`C_{ijk}=ϵ_{ijk}`$, then $`S_i\sigma _i/2`$ describes spin one-half . In this context, it may be remarked, because it underlies our interest, that more general models of particles with colour can possess supercharges involving the higher fermionic scalars. Demonstrations of how this arises, and its relationship to the mentioned hidden supersymmetries that do not close on the Hamiltonian of the model, will be addressed elsewhere.
## 2 The representation $`𝒟`$
It is easy to use (2) and (3) to show that
$`[S_a,\gamma _b]`$ $`=`$ $`iC_{abc}\gamma _c,`$
$`[S_a,S_b]`$ $`=`$ $`iC_{abc}S_c,`$ (9)
so that $`X_aS_a`$ is indeed a representation of $`𝔤`$. Next we define the quadratic Casimir operator
$$C_2=X_aX_a.$$
(10)
Then using only (3) and the Jacobi identity for the structure constants of $`𝔤`$, it follows that for $`𝒟`$
$$C_2(𝒟)=S_aS_a$$
(11)
has the form
$$C_2(𝒟)=\frac{1}{8}c_2(ad).\text{dim}𝔤\mathbf{.1}_{\mathrm{dim}\gamma },$$
(12)
where $`c_2(ad)`$, the eigenvalue of $`C_2`$ for the adjoint representation $`X_aad(X_a)`$ given by $`(adX_a)_{bc}=iC_{abc}`$, enters via
$$\mathrm{Tr}(adX_aadX_b)=c_2(ad)\delta _{ab}.$$
(13)
Thus, for $`su(n)`$,
$$C_2(𝒟)=\frac{1}{8}n(n^21)\mathbf{.1}_{\mathrm{dim}\gamma }.$$
(14)
To identify the representation $`𝒟`$ in terms of the irreps of $`𝔤`$ is convenient to discuss first the $`r`$ even case, even though the answer for the odd-dimensional $`𝔤`$ case is very similar. Since the irreducible $`\gamma _a`$ matrices of an $`(r=2s)`$-dimensional space have dimension $`2^s`$, $`𝒟`$ must be at least of dimension $`2^s`$ (larger if we used non-irreducible $`\gamma `$’s). So we ask, what irreps of $`𝔤`$ are contained in $`𝒟`$? The Weyl formula for the dimension of the irrep of $`𝔤`$ with highest weight $`\mathrm{\Lambda }`$ is given by the Weyl formula (see, e.g., , eq. (5.5)):
$$N(\mathrm{\Lambda })=\underset{positiveroots}{}(1+\frac{(\mathrm{\Lambda },\alpha )}{(\delta ,\alpha )}),$$
(15)
where the $`\alpha `$’s are the positive roots of $`𝔤`$. For the irrep corresponding to the Weyl principal vector $`\mathrm{\Lambda }=\delta =(1,\stackrel{l}{},1)`$ we find, since there are $`(rl)/2`$ positive roots,
$$N(\delta )=2^{(rl)/2}.$$
(16)
Further this irrep does have the correct eigenvalue of the Casimir operator (11) to agree with (13). To see this we use the general result (see , eq. (5.10))
$$c_2(\mathrm{\Lambda })=(\mathrm{\Lambda },\mathrm{\Lambda }+2\delta ),C_2(\mathrm{\Lambda })=c_2(\mathrm{\Lambda })\mathbf{.1}_{\mathrm{dim}_\delta },$$
(17)
to find
$$c_2(\delta )=3(\delta ,\delta )=\frac{1}{8}C_2(ad)r$$
(18)
upon use of the ‘strange formula’ of Freudenthal and de Vries (see , eqs. (7.20) and (14.30)). For $`su(n)`$, (18) reproduces the scalar in (14) since $`c_2(ad)=n`$. Moreover, we see from eq.(12) that, for any $`𝔤`$, $`c_2(\delta )=c_2(𝒟)`$. Since, for $`r`$ even, the irreducible gammas that determine the dimension of $`𝒟`$ are $`p\times p`$ matrices with $`p=2^{\frac{r}{2}}`$, and the irrep $`(1,\stackrel{l}{},1)`$ with the correct value for its Casimir operator involves $`q\times q`$ matrices with $`q=2^{{\scriptscriptstyle \frac{rl}{2}}}`$, the natural conclusion is that $`𝒟`$ is the direct sum of exactly $`2^{{\scriptscriptstyle \frac{l}{2}}}`$ copies of $`(1,\stackrel{l}{},1)`$. When $`𝔤`$ is odd-dimensional, $`r=2s+1`$, $`𝒟`$ is of dimension $`2^s=2^{{\scriptscriptstyle \frac{r1}{2}}}`$ and the same reasoning shows that the irrep of highest weight $`\mathrm{\Lambda }`$=$`\delta `$, of dimension $`2^{{\scriptscriptstyle \frac{rl}{2}}}`$ is contained $`2^{{\scriptscriptstyle \frac{l1}{2}}}`$ times in $`𝒟`$.
We may see this explicitly by constructing the required number of copies of the irrep $`(1,\stackrel{l}{},1)`$ in a fermionic Fock space description of $`𝒟`$. We first consider the even case $`𝔤=su(3)`$ fully; this not only shows how our realisation works in the simplest non-trivial case but also indicates exactly how the general even case is handled. The odd $`r`$ case will be discussed in Sec. 5.
## 3 The fermionic Fock space $`𝒱_𝒟`$ for $`su(3)`$
We build the fermionic Fock space of $`𝒟`$ for $`su(3)`$ by constructing four Dirac fermions out of the eight available gamma matrices $`\gamma _a`$, where the index $`a=1,\mathrm{},8`$ also labels the $`su(3)`$ generators (say, the Gell-Mann $`\lambda _a`$ matrices, $`[\lambda _a,\lambda _b]=2if_{abc}\lambda _c`$). However not all four Dirac fermions enter on the same footing, nor do they have the same role to play. So, on the one hand, we set
$$2B=\gamma _3+i\gamma _8,2B^{}=\gamma _3i\gamma _8,$$
(19)
corresponding to the $`su(3)`$ Cartan subalgebra generators $`\lambda _3`$ and $`\lambda _8`$, and, on the other, corresponding to the positive roots, we define
$$2A_5=\gamma _4i\gamma _5,2A_2=\gamma _6+i\gamma _7,2A_1=\gamma _1+i\gamma _2$$
(20)
(the somewhat strange labelling is adapted to the $`su(5)`$ case in next section). In this way, (3) translates into
$`\{A_i,A_{j}^{}{}_{}{}^{}\}`$ $`=`$ $`\delta _{ij},\{A_i,A_j\}=0,`$
$`\{B,B^{}\}`$ $`=`$ $`1,B^2=0,\{B,A_j\}=0,`$ (21)
which exhibit the fermionic nature of our realisation. Defining the fermion number operators
$$N_i=A_{i}^{}{}_{}{}^{}A_i(i=5,2,1),N_B=B^{}B,$$
(22)
$$\gamma _4\gamma _5=i(2N_51),\gamma _6\gamma _7=i(2N_21),\gamma _1\gamma _2=i(2N_11),\gamma _3\gamma _8=i(2N_B1)$$
(23)
we find
$$I_z=S_3=\frac{1}{2}(N_2+N_52N_1),$$
(24)
$$Y=\frac{2}{\sqrt{3}}S_8=N_5N_2,$$
(25)
independently of $`N_B`$.
It is now easy to observe the correspondence between the weights of an $`su(3)`$ octet (irrep $`(1,1)=\{8\}`$) and fermionic states labelled by occupation numbers
$`|n_5,n_2,n_1)`$ $`(I_z,Y)`$
$`|0,0,0)`$ $`(0,0)`$
$`|1,0,0)`$ $`(\frac{1}{2},1)`$
$`|0,1,0)`$ $`(\frac{1}{2},1)`$
$`|0,0,1)`$ $`(1,0)`$
$`|0,1,1)`$ $`(\frac{1}{2},1)`$
$`|1,0,1)`$ $`(\frac{1}{2},1)`$
$`|1,1,0)`$ $`(1,0)`$
$`|1,1,1)`$ $`(0,0).`$ (26)
Thus, apart from the eigenvalue $`I`$ of the $`su(2)`$ Casimir, which is used to distinguish the octet central states, the $`|n_5,n_2,n_1)`$ part of the full Fermi Fock space determines the bulk of the expected $`su(3)`$ characteristics. To complete the characterization we may use the $`su(3)`$ standard $`I,U,V`$-spin raising and lowering operators. These are
$`I_{}=S_1iS_2`$ $`=`$ $`(B+B^{})A_{1}^{}{}_{}{}^{}+A_2A_5,`$
$`V_{}=S_4iS_5`$ $`=`$ $`(\omega B+\omega ^2B^{})A_5+A_{2}^{}{}_{}{}^{}A_{1}^{}{}_{}{}^{},`$
$`U_+=S_6+iS_7`$ $`=`$ $`(\omega ^2B+\omega B^{})A_2+A_{1}^{}{}_{}{}^{}A_{5}^{}{}_{}{}^{},`$ (27)
where $`\omega =\frac{1}{2}(1i\sqrt{3})`$, $`\omega ^3=1`$, and $`I_+=I_{}^{},V_+=V_{}^{},U_{}=U_+^{}`$. The way $`B`$ enters here points the way directly towards the construction of two sets of orthogonal octet states in the fermionic Fock space of $`𝒟`$. We write $`|n_B,n_5,n_2,n_1`$ to denote the $`𝒱_𝒟`$ Fock space states that are simultaneous eigenstates of the number operators $`N_B`$ and the $`N_i`$; the standard $`su(3)`$ octet states of chirality $`\pm `$ are labelled $`|8,\pm ,I,I_z,Y`$. We define chirality by
$$\gamma _9\underset{a=1}{\overset{8}{}}\gamma _a=(2N_B1)\underset{k=5,2,1}{}(2N_k1)=(1)^{N_B+N_5+N_2+N_1}=(1)^{tot.Fermi\mathrm{\#}},$$
(28)
which commutes with the generators $`S_a`$ of the representation $`𝒟`$ of $`su(3)`$. To complete the construction we label the states of the two octets of different chiralities $`|8,\pm ,II_zY`$. To relate these to the Fock $`|n_B,n_5,n_2,n_1`$ states we refer to (3), and identify the highest weight $`\pm `$ octet states with the Fock space states of the correct chirality using
$$|8,\pm ,I=1,I_z=1Y=0|n_B=\genfrac{}{}{0pt}{}{0}{1},n_5=1,n_2=1,n_1=0.$$
(29)
The other octet states are constructed using $`V_{}`$ and $`U_+`$ each once and then $`I_{}`$ as needed to get all other states except the central state with $`I=I_z=Y=0`$. One gets this last state by orthogonality with the state with $`I=1,I_z=Y=0`$. The results that give the states $`|8,\pm ,I,I_z,Y`$ in terms of $`|n_B,n_5,n_2,n_1`$ with the correct phases thus are
$`|8,\pm ,\frac{1}{2},\frac{1}{2},1`$ $`=`$ $`U_+|\genfrac{}{}{0pt}{}{0}{1},1,1,0=\genfrac{}{}{0pt}{}{\omega }{\omega ^2}|\genfrac{}{}{0pt}{}{1}{0},1,0,0,`$
$`|8,\pm ,\frac{1}{2},\frac{1}{2},1`$ $`=`$ $`V_{}|\genfrac{}{}{0pt}{}{0}{1},1,1,0=\genfrac{}{}{0pt}{}{\omega ^2}{\omega }|\genfrac{}{}{0pt}{}{1}{0},0,1,0,`$
$`|8,\pm ,\frac{1}{2},\frac{1}{2},1`$ $`=`$ $`I_{}|8,\pm ,\frac{1}{2},\frac{1}{2}\mathrm{\hspace{0.17em}1}=\genfrac{}{}{0pt}{}{\omega }{\omega ^2}|\genfrac{}{}{0pt}{}{0}{1},1,0,1,`$
$`|8,\pm ,\frac{1}{2},\frac{1}{2},1`$ $`=`$ $`I_{}|8,\pm ,\frac{1}{2},\frac{1}{2},1=\genfrac{}{}{0pt}{}{\omega ^2}{\omega }|\genfrac{}{}{0pt}{}{0}{1},0,1,1,`$
$`|8,\pm ,1,\mathrm{0\hspace{0.17em}0}`$ $`=`$ $`\frac{1}{\sqrt{2}}I_{}|\genfrac{}{}{0pt}{}{0}{1},1,1,0=\frac{1}{\sqrt{2}}\left(|\genfrac{}{}{0pt}{}{1}{0},1,1,1+|\genfrac{}{}{0pt}{}{0}{1},0,0,0\right),`$
$`|8,\pm ,1,1,0`$ $`=`$ $`\frac{1}{\sqrt{2}}I_{}|8,\pm ,1,0,0=\pm |\genfrac{}{}{0pt}{}{1}{0},0,0,1`$ (30)
together with
$$|8,\pm ,0,0,0=\frac{1}{\sqrt{2}}|\genfrac{}{}{0pt}{}{1}{0},1,1,1\pm |\genfrac{}{}{0pt}{}{0}{1},0,0,0.$$
(31)
Since the $`I,U`$ and $`V`$ spin operators are all even, as any generator should be, all states within a given octet have the same chirality by (28). Thus, chirality distinguishes the two octets $`\{8,\pm \}`$ of even/odd total fermion number that span $`𝒟`$ in this case.
Let us now look for the role played by the two fermionic scalars $`K_3`$, $`K_5`$. For $`K_3`$ we may use the results given so far to find its complete operator expression in the full fermionic Fock space $`𝒱_𝒟`$. This gives
$$K_3=L_{125}B(N_1+\omega N_5+\omega ^2N_2)B^{}(N_1+\omega ^2N_5+\omega N_2),$$
(32)
where $`L_{125}=A_1A_2A_5+A_{5}^{}{}_{}{}^{}A_{2}^{}{}_{}{}^{}A_{1}^{}{}_{}{}^{}`$. $`L_{125}`$ in eq. (32) affects only the two central octet states, while the other terms affect only the states of the hexagonal rim of the octet (i.e. the states of the orbit of the state (29) under the Weyl group of $`su(3)`$). Eq. (3) leads easily to the result
$$K_3|8,\pm ,II_zY=|8,,II_zY.$$
(33)
Thus, $`K_3`$ changes the chirality of the state while respecting the $`su(3)`$ labels. This does not depend on $`𝔤`$ being $`su(3)`$: it is a consequence of $`K_3`$ being realized by a three-cocycle of a $`𝔤`$, since a cocycle is both odd (so that $`\{K_3,\gamma _9\}=0`$) and $`G`$-invariant. Thus, similar remarks apply also for the higher order fermionic operators $`K_{(2m_i1)}`$ associated with the various primitive cocycles of any $`𝔤`$, since they are all odd and $`G`$-invariant and, in particular, to $`K_5`$ for $`su(3)`$.
Since for larger $`𝔤`$ and higher cocycles the task of finding complete operator expressions for quantities like $`K_3,K_5,\mathrm{}`$ becomes rapidly more time-consuming, it is better to seek expressions only for their parts that are non-trivial in their action on highest weight states. We illustrate this for $`K_3`$ first. We have
$$K_3=\frac{1}{12}if_{abc}\gamma _a\gamma _b\gamma _c=\frac{1}{2}i\underset{triples}{}f_{abc}\gamma _a\gamma _b\gamma _c,$$
(34)
over non-trivial triples such that $`a<b<c`$, so that
$`K_3`$ $`=`$ $`{\displaystyle \frac{1}{2}}i\{\gamma _{123}+{\displaystyle \frac{1}{2}}(\gamma _{147}+\gamma _{165}+\gamma _{257}+\gamma _{246})+{\displaystyle \frac{1}{2}}(\gamma _{345}+\gamma _{376})+{\displaystyle \frac{\sqrt{3}}{2}}(\gamma _{458}+\gamma _{678})\}`$ (35)
$`=`$ $`{\displaystyle \frac{1}{2}}(N_2+N_52N_1)\gamma _3+{\displaystyle \frac{\sqrt{3}}{2}}(N_5N_2)\gamma _8+\mathrm{}`$
$`=`$ $`\gamma _3S_3+\gamma _8S_8+\mathrm{},`$
by (24), (25). Here the dots indicate terms (coming from the second bracket of (35)) that give vanishing contribution to the action of $`K_3`$ on the Weyl orbit of (29). Hence on (29)
$$K_3=\gamma _3=B+B^{}\mathrm{and}K_3^2=1,$$
(36)
in agreement with (32).
Similarly, for actions on (29) we have
$$K_5=\frac{k}{5!}\mathrm{\Omega }_{abcde}\gamma _a\gamma _b\gamma _c\gamma _d\gamma _e=k\underset{pentuples}{}\mathrm{\Omega }_{abcde}\gamma _a\gamma _b\gamma _c\gamma _d\gamma _e,$$
(37)
over pentuples such that $`a\{3,8\},c=b+1,c<d,e=d+1`$. This is justified: the total antisymmetry of $`\mathrm{\Omega }_{sabcd}`$ allows us to order the indices in any convenient way and sum over pentuples that both respect this order and give terms that fail to annihilate (29). We find
$$K_5=k\gamma _3(\mathrm{\Omega }_{31245}\mathrm{\Omega }_{31267})+k\gamma _8(\mathrm{\Omega }_{81245}\mathrm{\Omega }_{81267}\mathrm{\Omega }_{84567})=k\frac{\sqrt{3}}{12}\gamma _8(1+1+2),$$
(38)
using data from table 3 of and the results
$$i\gamma _1\gamma _2=2N_111,i\gamma _4\gamma _5=12N_51,i\gamma _6\gamma _7=2N_211,$$
(39)
for actions on (29). Thus
$$K_5=\frac{k}{\sqrt{3}}\gamma _8\mathrm{and}K_5^2=\frac{1}{3}k^2$$
(40)
on (29). It was also checked explicitly that non-zero coefficients of the type $`\mathrm{\Omega }_{38cde}`$ do not give rise to any non-vanishing contributions to (40).
Thus, setting $`k=\sqrt{3}`$, and acting on the vector space spanned by the two highest weight vectors of (29),
$$|n_b=0,HW,|n_b=1,HW,$$
(41)
we find $`K_3\sigma _1`$, and $`K_5\sigma _2`$, where $`2B=\gamma _3+i\gamma _8`$. Hence $`\chi =iK_3K_5i\sigma _1\sigma _2=\sigma _3`$. We may also use Hodge duality to give a direct proof
$`K_5`$ $`=`$ $`({\displaystyle \frac{\sqrt{3}}{5!}})\mathrm{\Omega }_{pqrst}\gamma _{pqrst}`$ (42)
$`=`$ $`\left({\displaystyle \frac{1}{12.5!}}\right)\gamma _{[pqrst]}ϵ_{pqrstxyz}f_{xyz}`$
$`=`$ $`{\displaystyle \frac{1}{12}}f_{xyz}\gamma _{[xyz]}\gamma _9=iK_3\gamma _9.`$
Here, we have used Eq.(8.14) of , the definition (28) of $`\gamma _9`$ and the result
$$ϵ_{pqrstxyz}\gamma _{[pqrst]}=5!\gamma _{[xyz]}\gamma _9.$$
(43)
Eq, (42) implies
$$iK_3K_5=\gamma _9$$
(44)
known to be represented by $`\sigma _3`$. As a final remark on the $`su(3)`$ case, we notice that insertion of (24) to (3) into (11) gives rise to the result
$$C_2(𝒟)=3\mathbf{\hspace{0.17em}1}_{16},$$
(45)
as it should by eq. (14), upon cancellation of all number operator terms.
## 4 The general $`𝔤`$ even case
### 4.1 General remarks
The method of Sec. 3 extends directly to a general even-dimensional $`𝔤`$ and indeed, without much modification, to the odd case (Sec. 5). For $`r`$ ($`l`$) even, we define $`\frac{l}{2}`$ operators $`B_1,\mathrm{},B_{l/2}`$, and $`(rl)/2`$ operators $`A_1,\mathrm{},A_{(rl)/2}`$, their adjoints and the corresponding number operators $`N_{B_\mu }`$ and $`N_{A_i}`$. To label the states within an irrep one needs the $`l`$ labels provided by the Cartan subalgebra generators plus $`(r3l)/2`$ additional ones to sort out the possible degeneracy. Since $`𝒟`$ contains $`(1,\stackrel{l}{},1)`$ $`2^{l/2}`$ times (Sec. 2), we still need $`l/2`$ labels taking two possible values to distinguish the states of the different copies of $`(1,\stackrel{l}{},1)`$ in $`𝒟`$. This means, in all, $`r/2`$ labels, provided by the $`l/2`$ operators $`N_{B_\mu }`$ and the $`(rl)/2`$ operators $`N_{A_i}`$. There is thus one $`A`$ for each of the positive roots and we can establish of a one to one correspondence like (3) between the weights of the irrep $`(1,\stackrel{l}{},1)`$ of $`𝔤`$ and the simultaneous eigenstates of the commuting number operators $`N_{A_i}`$. As seen in already for the $`su(3)`$ case (eq.(3)), however, both the $`B_\mu `$’s and the $`A_i`$’s appear in the definition of the various ladder operators that generate the states of the $`(1,\stackrel{l}{},1)`$ representation. The eigenvalues $`n_1,\mathrm{},n_{l/2}`$ of the $`N_{B_\mu }`$ can be used as in (18) to define the Fock space states equal to the highest weight states ($`2^{{\scriptscriptstyle \frac{l}{2}}}`$ of them) of the different copies of $`(1,\stackrel{l}{},1)`$ in $`𝒟`$. If one builds the rest of the states of each copy by application of lowering operators, one will find that any one state of any copy is orthogonal to all states of any other copy, as well, of course, to all of the states of its own copy.
The details of the next simplest even $`𝔤`$ of higher rank, $`su(5)`$, offers further insight into the situation surrounding the fermionic $`SU(5)`$-invariant operators $`K_3,K_5,K_7,K_9`$ built with the aid of its four non-trivial cocycles. Our discussion for $`su(5)`$ proceeds along lines similar to those followed for $`su(3)`$.
### 4.2 Basic definitions for $`su(5)`$
We use an explicit and essentially standard (cf. ) set of Gell-Mann lambda matrices for $`su(5)`$. For the diagonal matrices $`\lambda _{(p^21)},(p=2,3,4,5),`$ we have
$`\lambda _3=\mathrm{diag}(1,1,0,0,0)`$ , $`\sqrt{3}\lambda _8=\mathrm{diag}(1,1,2,0,0),`$
$`\sqrt{6}\lambda _{15}=\mathrm{diag}(1,1,1,3,0,)`$ , $`\sqrt{10}\lambda _{24}=\mathrm{diag}(1,1,1,1,4).`$ (46)
The index pairs
$$(1,2),(4,5),(6,7),(9,10),(11,12),(13,14),(16,17),(18,19),(20,21),(22,23),$$
(47)
are associated with the remaining $`\lambda `$’s in a way that can be inferred from the following array:
$$\left(\begin{array}{ccccc}& (1,2)& (4,5)& (9,10)& (16,17)\\ & & (6,7)& (11,12)& (18,19)\\ & & & (13,14)& (20,21)\\ & & & & (22,23)\\ & & & & \end{array}\right).$$
(48)
For example, $`\lambda _6`$ is a symmetric matrix whose only non-zero element above the main diagonal is a $`1`$ at the place marked $`(6,7)`$ in (48), while $`\lambda _7`$ is antisymmetric with a single entry $`i`$ in the same place. Next, follow the succesive ‘diagonals’ of the array to define the $`su(5)`$ Dirac fermions for the various pairs as follows.
$`2A_1=\gamma _1+i\gamma _2`$ , $`2A_2=\gamma _6+i\gamma _7,`$
$`2A_3=\gamma _{13}i\gamma _{14}`$ , $`2A_4=\gamma _{22}+i\gamma _{23},`$
$`2A_5=\gamma _4i\gamma _5`$ , $`2A_6=\gamma _{11}+i\gamma _{12},`$
$`2A_7=\gamma _{20}+i\gamma _{21}`$ , $`2A_8=\gamma _9+i\gamma _{10},`$
$`2A_9=\gamma _{18}i\gamma _{19}`$ , $`2A_{10}=\gamma _{16}i\gamma _{17},`$ (49)
plus
$$2B_1=\gamma _3+i\gamma _8,2B_2=\gamma _{15}i\gamma _{24}.$$
(50)
This leads to relations with the number operators $`N_\alpha =A_{\alpha }^{}{}_{}{}^{}A_\alpha `$, like
$`i\gamma _1\gamma _2=2N_11\mathrm{and}\mathrm{similarly}\mathrm{for}`$ $`N_\alpha `$ $`\mathrm{when}\alpha \{1,2,4,6,7,8\}`$
$`i\gamma _4\gamma _5=2N_51\mathrm{and}\mathrm{similarly}\mathrm{for}`$ $`N_\alpha `$ $`\mathrm{when}\alpha \{3,5,9,10\}.`$ (51)
Then $`S_a=\frac{1}{4}if_{abc}\gamma _b\gamma _c`$, with the aid of MAPLE output for structure constants, allows us to derive
$`2S_3=2I_3`$ $`=`$ $`N_2+N_52N_1N_8+N_{10}N_9+N_6,`$
$`S_8={\displaystyle \frac{\sqrt{3}}{2}}Y,3Y`$ $`=`$ $`(3N_53N_2N_8N_6+N_9+N_{10}2N_3+2N_7),`$
$`S_{15}={\displaystyle \frac{\sqrt{6}}{3}}Z_3,4Z_3`$ $`=`$ $`(4N_34N_64N_8+N_9+N_{10}N_7+3N_4),`$
$`S_{24}={\displaystyle \frac{\sqrt{10}}{4}}Z_4,Z_4`$ $`=`$ $`N_9+N_{10}N_4N_7,`$ (52)
which reproduce the $`su(3)`$ expressions when only the (1,2 and 5)-labelled quantities are retained. Our choice of signs in (49) depends upon the signs of the structure constants for $`su(5)`$ – our choice of lambda-matrices was made above to yield agreement with the tables given in , upon our wish of avoiding constant terms in the definitions (52) and upon our desire of having all entries in (57) below equal to +1.
We define the highest weight state of any irrep of $`su(5)`$ by taking first the highest $`Z_4`$ eigenvalue, then the highest $`Z_3`$ eigenvalue that can arise for that $`Z_4`$ eigenvalue. Next the highest $`Y`$ and finally the highest $`I_3`$. Thus we get $`Z_4=2`$ for $`N_4=N_7=0,N_9=N_{10}=1`$. Hence $`Z_3=N_8N_6+N_3+\frac{1}{2}=\frac{3}{2}`$ for $`N_8=N_6=0,N_3=1`$, and $`Y=N_5N_2=1`$ for $`N_2=0,N_5=1`$. Finally $`I_3=\frac{1}{2}`$ for $`N_1=0`$. Hence our highest weight state for any of the four possible irreps $`(1,1,1,1)`$ of $`su(5)`$ in $`𝒟`$ includes the $`(rl)/2=10`$ labels
$$|0,0,1,0,1,0,0,0,1,1),$$
(53)
where we have written the eigenvalues $`n_\alpha `$ of the $`N_\alpha `$ with the $`\alpha `$ in the standard order $`1,\mathrm{},10`$. We remark that our method of definition corresponds to the use of the subgroup chain in which $`(1,1,1,1)`$ of $`su(5)`$ reduces to $`(1,1,1)_2`$ of $`su(4)\times u(1)`$, plus irreps of lower $`Z_4`$, and then $`(1,1,1)`$ of $`su(4)`$ reduces to $`(1,1)_{\frac{3}{2}}`$ of $`su(3)\times u(1)`$, plus irreps of $`su(3)`$ of lower $`Z_3`$, and so on.
For $`𝔤=su(5)`$, there are two operators $`N_{B_1}`$ and $`N_{B_2}`$, providing the $`l/2=2`$ additional labels $`|n_{B_1},n_{B_2})`$ that give rise to the possibilities
$$|0,0),|1,1),|0,1),|1,0),$$
(54)
available to construct the highest weight states of chiralities $`+,+,,`$ of the four different copies of the $`(1,1,1,1)`$ irrep contained in $`𝒟`$. The complete irreps can be built from these by lowerings.
The $`K_3`$ operator commutes with all $`su(5)`$ actions and changes chiralities. This does not furnish a complete picture; there are three other fermionic scalar operators that deserve to be treated on the same footing as the simplest one $`K_3`$. We turn to them next.
### 4.3 $`su(5)`$ fermionic operators
We define the 5th order fermionic operator for $`su(5)`$ via the five-cocycle $`\mathrm{\Omega }_{sabcd}`$ by
$`K_5`$ $`=`$ $`{\displaystyle \frac{k_5}{5!}}\mathrm{\Omega }_{sabcd}\gamma _s\gamma _a\gamma _b\gamma _c\gamma _d,k_5R,`$ (55)
$`=`$ $`k_5{\displaystyle \underset{pentuples}{}}\mathrm{\Omega }_{sabcd}(i\gamma _a\gamma _b)(i\gamma _c\gamma _d)\gamma _s,`$
over pentuples such that $`s\{3,8,15,24\}`$, and $`a\{1,4,6,9,11,13,16,18,20\}`$, and $`b=a+1,c>b,d=b+1`$. These are the only ones that can fail to annihilate the highest weight state. In fact, for the action of $`K_5`$ on (53), it follows that
$$(i\gamma _x\gamma _{x+1})=1\text{for all}x\{1,4,6,9,11,13,16,18,20,22\}.$$
(56)
Thus if define an array $`N_{ab}`$ whose only non-vanishing entries on (53) are
$$N_{12}=N_{45}=N_{67}=N_{9,10}=N_{11,12}=N_{13,14}=N_{16,17}=N_{18,19}=N_{20,21}=N_{22,23}=+1,$$
(57)
we may do a computation of the coefficients in $`K_5`$ of the $`\gamma _s`$, for each of $`s\{3,8,15,24\}`$, by a separate MAPLE run. This yields a result valid for the action of $`K_5`$ on highest weight states:
$$2K_5=k_5(3\gamma _3+\frac{5}{3}\sqrt{3}\gamma _8+\frac{1}{3}\sqrt{6}\gamma _{15}\sqrt{10}\gamma _{24}).$$
(58)
The same procedure works for the 7th order fermionic scalar $`K_7`$ of $`su(5)`$, given by
$`K_7`$ $`=`$ $`i{\displaystyle \frac{k_7}{7!}}\mathrm{\Omega }_{sabcdef}(i\gamma _{ab})(i\gamma _{cd})(i\gamma _{ef})\gamma _s`$ (59)
$`=`$ $`k_7{\displaystyle \underset{heptuples}{}}\mathrm{\Omega }_{sabcdefg}(i\gamma _{ab})(i\gamma _{cd})(i\gamma _{ef})\gamma _s,`$
over heptuples with $`b=a+1,d=c+1,f=e+1,c>b,e>d`$ and
$$a\{1,4,6,9,11,13,16,18\}.$$
This leads, by four runs of the corresponding MAPLE program, to
$$K_7=k_7(\frac{2}{5}\gamma _3\frac{2}{45}\sqrt{3}\gamma _8\frac{7}{45}\sqrt{6}\gamma _{15}+\frac{1}{15}\sqrt{10}\gamma _{24})$$
(60)
on the HW states (53).
It is of course possible to derive (by MAPLE program) that
$$2K_3=\gamma _3+\sqrt{3}\gamma _8+\sqrt{6}\gamma _{15}+\sqrt{10}\gamma _{24}$$
(61)
on (53). But it is easier to derive this directly from
$$K_3=S_3\gamma _3+S_8\gamma _8+S_{15}\gamma _{15}+S_{24}\gamma _{24},$$
(62)
and insertion of the known eigenvalues of the $`S_s`$ for (53) gives back (61). To see that (62) for $`K_3`$ is correct, think in terms of triples $`s,x,x+1`$ with $`s\{3,8,15,24\}`$ and
$$x\{1,4,6,9,11,13,16,18,20,22\}.$$
(63)
The relevant $`f_{abc}`$ are in the tables of and we used data from MAPLE, which agreed with these tables.
### 4.4 Discussion of $`su(5)`$ results
By giving a direct evaluation that uses Jacobi identites, it may be shown that the different fermionic scalars $`K_3`$, etc., anticommute. It is wise to verify that the results (58), (60), (63) are in agreement with this. This requires only that the numbers that come from terms like $`\gamma _{3}^{}{}_{}{}^{2}=1`$, etc. sum to zero, which they indeed do. Our results also give the squares of the fermionic scalars
$$K_{3}^{}{}_{}{}^{2}=5,K_{5}^{}{}_{}{}^{2}=7k_{5}^{}{}_{}{}^{2},K_{7}^{}{}_{}{}^{2}=\frac{16}{45}k_{7}^{}{}_{}{}^{2}.$$
(64)
Fortunately we know sufficient cocycle identities valid for all $`su(n)`$ to complete a worthwhile check upon our MAPLE computation methods.
We have $`K_3=\frac{1}{12}if_{abc}\gamma _a\gamma _b\gamma _c`$, and also $`K_{3}^{}{}_{}{}^{2}=xI`$ for some $`x`$, since $`K_{3}^{}{}_{}{}^{2}`$ is a scalar. Hence
$$(\mathrm{dim}\gamma )x=\mathrm{tr}K_{3}^{}{}_{}{}^{2}=\frac{1}{24}(\mathrm{dim}\gamma )f_{abc}f_{abc},$$
(65)
using Dirac trace methods with $`\mathrm{tr}I=\mathrm{dim}\gamma `$. Since $`f_{abc}f_{abc}=n(n^21)=5.24`$, we find $`x=5`$, as required. To confirm the answers (64) for the cases of $`K_5`$ and $`K_7`$, with $`k_5`$ and $`k_7`$ set equal to $`1`$, we need the identities which are the $`n=5`$ special cases of
$$\mathrm{\Omega }_{abcde}\mathrm{\Omega }_{abcde}=\frac{1}{3}n(n^21)(n^24),$$
(66)
$$\mathrm{\Omega }_{abcdefg}\mathrm{\Omega }_{abcdefg}=\frac{2}{45}n(n^21)(n^24)(n^29).$$
(67)
Thus, setting $`K_{5}^{}{}_{}{}^{2}=yI`$ for some $`y`$, we obtain
$$(\mathrm{dim}\gamma )y=\mathrm{tr}K_{5}^{}{}_{}{}^{2}=\frac{1}{5!}\mathrm{\Omega }_{abcde}\mathrm{\Omega }_{abcde}(\mathrm{dim}\gamma )$$
(68)
upon evaluating the trace in a fashion that takes full advantage of antisymmetries. Hence, using (66) at $`n=5`$, we find $`y=7`$, as expected.
We note the right hand sides of (66) and (67) are zero for low enough $`n`$ as consistency requires. We note also that (66) evaluated for $`n=3`$ allows us to see agreement with the result (40), when one puts $`k=1`$ in (40). Also, we check from (61),(58), and (60) that $`K_3,K_5`$ and $`K_7`$ anticommute with each other. Next, by asking for the unique linear combination that anticommutes with $`K_3,K_5`$ and $`K_7`$, we can show that on highest weight states (53)
$$K_9=k_9(10\gamma _310\sqrt{3}\gamma _8+5\sqrt{6}\gamma _{15}\sqrt{10}\gamma _{24}),$$
(69)
so that
$$K_{9}^{}{}_{}{}^{2}=560k_{9}^{}{}_{}{}^{2}.$$
(70)
Setting $`k_5,k_7,k_9=1`$, we can compute the effect on (53) of the product $`K_3K_5K_7K_9`$. It is found that, as expected, only terms containing permutations of $`\gamma _3\gamma _8\gamma _{15}\gamma _{24}`$ survive, so that
$$K_3K_5K_7K_9=\frac{16}{3}.7.\sqrt{5}.\gamma _3\gamma _8\gamma _{15}\gamma _{24}.$$
(71)
If we define $`L_3`$ so that $`L_3=c_3K_3`$ and $`L_3^2=1`$, and so on, then above results allow the choice
$$c_3=\sqrt{5},c_5=\sqrt{7},c_7=\frac{3}{4}\sqrt{5},c_9=\frac{1}{4\sqrt{35}},$$
(72)
so that
$$\chi =L_3L_5L_7L_9=\gamma _3\gamma _8\gamma _{15}\gamma _{24},$$
(73)
and $`\chi ^2=1`$.
In view of the direct computation (42) for $`su(3)`$ and of (71), we might have expected to find
$$\chi =\gamma _{25}=\underset{\alpha =1}{\overset{24}{}}\gamma _\alpha $$
(74)
rather than (73) for (factors apart) the product of the $`K`$’s in (71). However, recalling from (56) that for action upon (53) $`\gamma _1\gamma _2=i`$ etc., we see that all the $`\gamma `$’s absent from (73) can be smuggled back (73) at the sole cost of a factor $`i^{10}=1`$, so that the two expressions for $`\chi `$ (73), (74) coincide. We shall come back to this point in Sec. 6.
Since the $`L`$ operators commute with the $`su(5)`$ action, we can consider the effect on the vector of $`su(5)`$ highest weight states
$$|00,HW,|11,HW,|01,HW,|10,HW,$$
(75)
where the first two labels provide the eigenvalues of the number operators $`N_{B1}`$ and $`N_{B2}`$, associated with the Dirac fermions $`2B_1`$ and $`2B_2`$ in (50), and $`HW`$ indicates the specification given previously (53) of the highest weight state of the $`su(5)`$ irrep $`(1,1,1,1)`$.
Thus the four $`L`$-operators may be represented by a set of four Dirac matrices $`\mathrm{\Gamma }_\mu `$, where $`\mu \{1,2,3,4\}`$ with $`\chi `$ represented by $`\mathrm{\Gamma }_5=\mathrm{\Gamma }_1\mathrm{\Gamma }_2\mathrm{\Gamma }_3\mathrm{\Gamma }_4`$. While the $`\mathrm{\Gamma }_\mu `$ have complicated looking actions, although with the correct anti-commutation relations
$$\{\mathrm{\Gamma }_\mu ,\mathrm{\Gamma }_\nu \}=2\delta _{\mu \nu },$$
(76)
the bosonic scalar $`\chi `$ is given by
$$\chi =(2N_{B1}1)(2N_{B2}1)I\sigma _3$$
(77)
To compare with the action of $`K_3`$, $`K_5`$ and chirality on the highest weight states (29), we recall that from Sec. 3 that these actions are given simply by Pauli matrices. Also the eigenvalues of chirality distinguish between the two octets that comprise the Brink-Ramond representation of $`su(3)`$. While we have found a fairly natural analogue of this picture for $`su(5)`$, no nice representation of the $`\mathrm{\Gamma }`$ matrices just discussed emerges.
## 5 The odd case
Again doing the simplest odd case of $`su(2)`$ with $`r=3`$ and $`l=1`$ with sufficient care indicates the pattern that governs the general odd case quite clearly.
Given $`S_c=\frac{1}{4}iϵ_{abc}\gamma _a\gamma _b`$, we have
$$2A=\gamma _1i\gamma _2,\mathrm{\hspace{0.33em}2}A^{}=\gamma _1+i\gamma _2,N_A=A^{}A,i\gamma _1\gamma _2=2(N_A1),$$
(78)
$$S_3=\frac{1}{2}i\gamma _1\gamma _2=N_A\frac{1}{2},S_+=\gamma _3A^{},S_{}=A\gamma _3.$$
(79)
It then follows algebraically that $`𝐒^2=\frac{3}{4}I`$, without raising the question as to whether or not we use Pauli matrices $`\sigma _1`$ and $`\sigma _2`$ for $`\gamma _1`$ and $`\gamma _2`$. The $`K_3`$ operator is now given by
$$K_3=\frac{i}{12}ϵ_{abc}\gamma _a\gamma _b\gamma _c$$
(80)
One way to get a complete Fock space description introduces $`\varphi `$ such that $`\varphi ^2=1,\varphi ^{}=\varphi ,\{\varphi ,\gamma _a\}=0`$, for $`a\{1,2,3\}`$. Defining $`B`$ by
$$2B=\gamma _3i\varphi ,$$
(81)
so that $`N_B=B^{}B`$, $`i\gamma _3\varphi =2(N_B1)`$, we find the representation
$$S_+=(B+B^{})A^{},S_{}=(S_+)^{}.$$
(82)
Now in the $`^2^2`$ Fock space of the $`B`$ and $`A`$ fermions, we build two equivalent $`j=\frac{1}{2}`$ irreps of $`su(2)`$ of chiralities $`\pm `$, using the definition of chirality
$$\gamma _5\gamma _1\gamma _2\gamma _3\varphi =(1)^{N_B+N_A}=(1)^{tot.Fermi\mathrm{\#}},$$
(83)
Labelling the Fock states $`|n_B,n_A`$, we set
$$|\pm ,j=\frac{1}{2}m=\frac{1}{2}=|\genfrac{}{}{0pt}{}{0}{1},0.$$
(84)
Then,
$$|\pm ,j=m=\frac{1}{2}=S_+|\pm ,j=\frac{1}{2}m=\frac{1}{2}=\pm |\genfrac{}{}{0pt}{}{1}{0},1$$
(85)
with a typical and essential fermionic minus sign making its appearance. It is clear that we have two orthogonal $`j=\frac{1}{2}`$ irreps. It is easy too to compute the matrices of our Fock space representation ($`\sigma _0=\mathrm{𝟏}_2`$)
$$2S_3\sigma _0\sigma _3,\mathrm{\hspace{0.33em}2}S_1\sigma _3\sigma _1,\mathrm{\hspace{0.33em}2}S_2\sigma _3\sigma _2,$$
(86)
$$\gamma _1\sigma _2\sigma _2,\gamma _2\sigma _2\sigma _1,\gamma _3\sigma _1\sigma _0,,$$
(87)
$$\varphi =\gamma _4\sigma _2\sigma _3,\gamma _5\sigma _3\sigma _0,\mathrm{\hspace{0.33em}2}K_3\sigma _1\sigma _3,$$
(88)
with all the expected properties. The action of $`K_3`$ (eq. (80)) on any state of our Fock space basis is to produce a state of opposite chirality. We note that all quantities that should anticommute with other fermionic quantities actually do so.
To achieve all the features noted above, a certain price has had to be paid. In case it may not be thought worthwhile to pay it, especially in simple contexts, we offer the following argument. To reach a Fock space description of $`𝒟`$ for $`su(2)`$, we chose to introduce a dynamical variable $`\varphi `$ of Majorana type not present in the original formalism. This led us to use $`4\times 4`$ rather than $`2\times 2`$ gammas. Had we used Pauli matrices, $`\gamma _a=\sigma _a`$, the Kostant operator (80) would have read
$$K_3=\frac{1}{12}iϵ_{abc}\sigma _a\sigma _b\sigma _c=\frac{1}{2}\mathbf{\hspace{0.17em}1}_2,$$
(89)
losing sight of the dynamical role of $`K_3`$, and of its fermionic nature as well. Using the irreducible two-dimensional picture is much the same as using only the chirality plus piece of the system, and forgetting about the negative chirality piece. This does represent $`𝒟`$ itself more or less satisfactorily, but fails, as our formalism does not, in the treatment of a composite system of two independent systems of type $`𝒟`$ .
Passing to the general odd case, we see that it is only the last gamma $`\gamma _r`$, viewed as a lone Majorana fermion, that makes us do more than we have already done for the even case. And it is to be treated just as we described for ($`\gamma _3,\varphi `$) in the $`su(2)`$ example. Thus we have $`(rl)/2`$ fermions of $`A`$-type as well as $`(l+1)/2`$ of $`B`$-type, allowing for $`(l+1)/2`$ copies of the irrep $`(1,\stackrel{l}{},1)`$. The latter irrep has dimension $`2^{(rl)/2}`$ and uses $`r=2s+1`$ gammas that may represented minimally by $`2^s\times 2^s`$ matrices. This itself would imply that $`𝒟`$ contains $`2^p`$ copies of $`(1,\stackrel{l}{},1)`$ with $`p=s(rl)/2=(l1)/2`$. But we have in our Fock space realisation of $`𝒟`$ twice as many copies in virtue of the fact that we chose to adjoin to the original dynamical system one additional gamma, thereby doubling the size of all the gammas. We have the choice for $`𝒟`$ itself of accepting this, or else, of using only the positive chirality half of the full fermionic Fock space that carries $`2^{(l1)/2}`$ copies of $`(1,\stackrel{l}{},1)`$ for $`𝔤`$. As above, this fails to bring into clear focus the full fermionic nature of all the relevant fermionic quantities for the system $`𝒟`$.
## 6 The chirality operator $`\gamma _{r+1}`$ and final remarks
The previous discussion has exhibited the existence of $`l`$ fermionic anticommuting scalars $`K_{(2m_i1)}`$ which are constructed from the Lie algebra cohomology cocycles of a Lie algebra $`𝔤`$ of even dimension $`r`$, and the associated set of Dirac matrices. In the even case, to which the rest of this section mainly refers, all these fermionic scalar operators change the chirality of the states, since they anticommute with $`\gamma _{r+1}`$. In the simplest case of $`su(3)`$, we have seen in eq. (44) that $`iK_3K_5=\gamma _9`$, while for $`su(5)`$ we see, from (71)-(74), that for actions on highest weight states
$$K_3K_5K_7K_9\chi =\gamma _{25}=\underset{\alpha =1}{\overset{24}{}}\gamma _\alpha .$$
(90)
But, since all $`su(5)`$ generators and in particular the raising and lowering operators commute with the ($`𝔤`$-invariant) $`K`$’s, the same applies to actions on all states. In general, we expect, for each even $`𝔤`$, that $`\gamma _{r+1}`$ is given to within normalisation by the product of all the available $`K`$’s, $`_{i=1}^lK_{(2m_i1)}`$; notice that $`_{i=1}^l(2m_i1)=r`$.
The above properties of the various $`K`$’s reflect the underlying group geometry. The $`l`$ cocycles of the Lie algebra cohomology of $`𝔤`$ may be looked at as invariant forms on the compact group manifold $`G`$ associated with $`𝔤`$ and, although not all the form properties are transported to the $`K`$’s (for instance, unlike forms, $`K^21`$), some of them are. The product of all the $`(2m_i1)`$-forms associated with the cocycles is the volume $`r`$-form on $`G`$, an even form for each even $`𝔤`$. This accounts for duality relations such as (44), which can be read as $`K_3(K_3)\gamma _9`$, expressing the fact that $`K_3`$ and $`K_5`$ are dual to each other. It also shows that, when defined, the even chirality operator takes over the role of the volume form on the group manifold in the present context.
The original motivation of was to study the invariant cubic Kostant operator on Lie algebra (symmetric) cosets, to understand the physical degrees of freedom of certain supersymmetric theories; these appeared as solutions to the Kostant-Dirac equation associated with specific cosets. The generalisations introduced in this paper retain many of the properties of the representation independent part of the cubic Kostant cubic operator, in particular that their square is given by Casimir invariants. At the same time, however, they have rich geometrical properties that reflect their Lie algebra cohomology origin (as e.g., that the product of the $`l`$ operators $`K_i`$ in the even case is represented by the chirality/volume form). It seems worthwhile to extend the geometrical methods of this paper to the coset case, not discussed here, and to the possible full higher order Kostant operators. This, and the analysis of the hidden supersymmetries mentioned in the introduction, will be discussed elsewhere.
Acknowledgements. This work was partly supported by the DGICYT, Spain ($`\mathrm{\#}`$PB 96-0756) and PPARC, UK.
|
warning/0003/physics0003074.html
|
ar5iv
|
text
|
# Untitled Document
HIERARCHIC THEORY OF CONDENSED MATTER :
New state equation &Interrelation between mesoscopic and macroscopic properties
Alex Kaivarainen
JBL, University of Turku, FIN-20520, Turku, Finland
URL: http://www.karelia.ru/~alexk
H2o@karelia.ru
> Materials, presented in this original article are based on:
>
> . Book by A. Kaivarainen: Hierarchic Concept of Matter and Field. Water, biosystems and elementary particles. New York, 1995 and two articles:
>
> . New Hierarchic Theory of Matter General for Liquids and Solids:
>
> dynamics, thermodynamics and mesoscopic structure of water and ice
>
> (see: http://arXiv.org/abs/physics/0003044 and http://www.karelia.ru/~alexk \[New articles\]);
>
> . Hierarchic Concept of Condensed Matter and its Interaction with Light: New Theories of Light Refraction, Brillouin Scattering and Mössbauer effect (see URL: http://www.karelia.ru/~alexk \[New articles\]).
>
> Computerized verification of described here new theories are presented on examples of WATER and ICE, using special computer program (copyright, 1997, A. Kaivarainen).
>
> CONTENTS
>
> Summary of Hierarchic Theory of Matter and Field
>
> 1 The state equation for real gas
>
> 2 New state equation for condensed matter
>
> 3 Vapor pressure
>
> 4 Surface tension
>
> 5 Mesoscopic theory of thermal conductivity
>
> 6 Mesoscopic theory of viscosity for liquids and solids
>
> 7 Brownian diffusion
>
> 8 Self-diffusion in liquids and solids
>
> 9 Mesoscopic approach to proton conductivity in water, ice and other systems containing hydrogen bonds
>
> 10 Regulation of pH and shining of water by electromagnetic and acoustic fields
Summary of
Hierarchic Theory of Condensed Matter (http://arXiv.org/abs/physics/0003044)
A basically new hierarchic quantitative theory, general for solids and liquids, has been developed.
It is assumed, that unharmonic oscillations of particles in any condensed matter lead to emergence of three-dimensional (3D) superposition of standing de Broglie waves of molecules, electromagnetic and acoustic waves. Consequently, any condensed matter could be considered as a gas of 3D standing waves of corresponding nature. Our approach unifies and develops strongly the Einstein’s and Debye’s models.
Collective excitations, like 3D standing de Broglie waves of molecules, representing at certain conditions the mesoscopic molecular Bose condensate, were analyzed, as a background of hierarchic model of condensed matter.
The most probable de Broglie wave (wave B) length is determined by the ratio of Plank constant to the most probable impulse of molecules, or by ratio of its most probable phase velocity to frequency. The waves B are related to molecular translations (tr) and librations (lb).
As far the quantum dynamics of condensed matter does not follow in general case the classical Maxwell-Boltzmann distribution, the real most probable de Broglie wave length can exceed the classical thermal de Broglie wave length and the distance between centers of molecules many times.
This makes possible the atomic and molecular Bose condensation in solids and liquids at temperatures, below boiling point. It is one of the most important results of new theory, confirmed by computer simulations on examples of water and ice.
Four strongly interrelated new types of quasiparticles (collective excitations) were introduced in our hierarchic model:
1. Effectons (tr and lb), existing in ”acoustic” (a) and ”optic” (b) states represent the coherent clusters in general case;
2. Convertons, corresponding to interconversions between tr and lb types of the effectons (flickering clusters);
3. Transitons are the intermediate $`\left[ab\right]`$ transition states of the tr and lb effectons;
4. Deformons are the 3D superposition of IR electromagnetic or acoustic waves, activated by transitons and convertons.
Primary effectons (tr and lb) are formed by 3D superposition of the most probable standing de Broglie waves of the oscillating ions, atoms or molecules. The volume of effectons (tr and lb) may contain from less than one, to tens and even thousands of molecules. The first condition means validity of classical approximation in description of the subsystems of the effectons. The second one points to quantum properties of coherent clusters due to molecular Bose condensation.
The liquids are semiclassical systems because their primary (tr) effectons contain less than one molecule and primary (lb) effectons - more than one molecule. The solids are quantum systems totally because both kind of their primary effectons (tr and lb) are molecular Bose condensates. These consequences of our theory are confirmed by computer calculations.
The 1st order $`\left[gasliquid\right]`$ transition is accompanied by strong decreasing of rotational (librational) degrees of freedom due to emergence of primary (lb) effectons and $`\left[liquidsolid\right]`$ transition - by decreasing of translational degrees of freedom due to mesoscopic Bose-condensation in form of primary (tr) effectons.
In the general case the effecton can be approximated by parallelepiped with edges corresponding to de Broglie waves length in three selected directions (1, 2, 3), related to the symmetry of the molecular dynamics. In the case of isotropic molecular motion the effectons’ shape may be approximated by cube. The edge-length of primary effectons (tr and lb) can be considered as the ”parameter of order”.
The in-phase oscillations of molecules in the effectons correspond to the effecton’s (a) - acoustic state and the counterphase oscillations correspond to their (b) - optic state. States (a) and (b) of the effectons differ in potential energy only, however, their kinetic energies, impulses and spatial dimensions - are the same. The b-state of the effectons has a common feature with Frölich’s polar mode.
The $`(ab)`$ or $`(ba)`$ transition states of the primary effectons (tr and lb), defined as primary transitons, are accompanied by a change in molecule polarizability and dipole moment without density fluctuations. At this case they lead to absorption or radiation of IR photons, respectively. Superposition (interception) of three internal standing IR photons of different directions (1,2,3) - forms primary electromagnetic deformons (tr and lb).
On the other hand, the \[lb$``$tr\] convertons and secondary transitons are accompanied by the density fluctuations, leading to absorption or radiation of phonons.
Superposition, resulting from interception of standing phonons in three directions (1,2,3) is termed: secondary acoustic deformons (tr and lb).
Correlated collective excitations of primary and secondary effectons and deformons (tr and lb), localized in the volume of primary tr and lb electromagnetic deformons, lead to origination of macroeffectons, macrotransitons and macrodeformons (tr and lb respectively).
Correlated simultaneous excitations of tr and lb macroeffectons in the volume of superimposed tr and lb electromagnetic deformons lead to origination of supereffectons.
In turn, the coherent excitation of both: tr and lb macrodeformons and macroconvertons in the same volume means creation of superdeformons. Superdeformons are the biggest (cavitational) fluctuations, leading to microbubbles in liquids and to local defects in solids.
Total number of quasiparticles of condensed matter equal to 4!=24, reflects all of possible combinations of the four basic ones \[1-4\], introduced above. This set of collective excitations in the form of ”gas” of 3D standing waves of three types: de Broglie, acoustic and electromagnetic - is shown to be able to explain virtually all the properties of all condensed matter.
The important positive feature of our hierarchic model of matter is that it does not need the semi-empiric intermolecular potentials for calculations, which are unavoidable in existing theories of many body systems. The potential energy of intermolecular interaction is involved indirectly in dimensions and stability of quasiparticles, introduced in our model.
The main formulae of theory are the same for liquids and solids and include following experimental parameters, which take into account their different properties:
$`\left[1\right]`$\- Positions of (tr) and (lb) bands in oscillatory spectra;
$`\left[2\right]`$\- Sound velocity;
$`\left[3\right]`$\- Density;
$`\left[4\right]`$\- Refraction index (extrapolated to the infinitive wave length of photon$`)`$.
The knowledge of these four basic parameters at the same temperature and pressure makes it possible using our computer program, to evaluate more than 300 important characteristics of any condensed matter. Among them are such as: total internal energy, kinetic and potential energies, heat capacity and thermal conductivity, surface tension, vapor pressure, viscosity, coefficient of self-diffusion, osmotic pressure, solvent activity, etc. Most of calculated parameters are hidden, i.e. inaccessible to direct experimental measurement.
The new interpretation and evaluation of Brillouin light scattering and Mössbauer effect parameters may also be done on the basis of hierarchic theory. Mesoscopic scenarios of turbulence, superconductivity and superfluity are elaborated.
Some original aspects of water in organization and large-scale dynamics of biosystems - such as proteins, DNA, microtubules, membranes and regulative role of water in cytoplasm, cancer development, quantum neurodynamics, etc. have been analyzed in the framework of Hierarchic theory.
Computerized verification of our Hierarchic theory of matter on examples of water and ice is performed, using special computer program: Comprehensive Analyzer of Matter Properties (CAMP, copyright, 1997, Kaivarainen). The new optoacoustic device, based on this program, with possibilities much wider, than that of IR, Raman and Brillouin spectrometers, has been proposed (see URL: http://www.karelia.ru/~alexk \[CAMP\]).
This is the first theory able to predict all known experimental temperature anomalies for water and ice. The conformity between theory and experiment is very good even without any adjustable parameters.
The hierarchic concept creates a bridge between micro- and macro- phenomena, dynamics and thermodynamics, liquids and solids in terms of quantum physics.
\********************************************************************
1. The state equation for real gas
The Clapeyrone-Mendeleyev equation sets the relationship between pressure $`(P)`$, volume $`(V)`$ and temperature ($`T`$) values for the ideal gas containing $`N_0`$ molecules (one mole):
$$PV=N_0kT=RT$$
(1)
In the real gases interactions between the molecules and their sizes should be taken into account. It can be achieved by entering the corresponding amendments into the left part, to the right or to the both parts of eq. (1).
It was Van der Waals who choused the first way more than a hundred years ago and derived the equation:
$$\left(P+\frac{a}{V^2}\right)\left(\begin{array}{c}Vb\end{array}\right)=RT$$
(2)
where the attraction forces are accounted for by the amending term $`(a/V^2)`$, while the repulsion forces and the effects of the excluded volume accounted for the term (b).
Equation (2) correctly describes changes in P,V and T related to liquid-gas transitions on the qualitative level. However, the quantitative analysis by means of (2) is approximate and needs the fitting parameters. The parameters (a) and (b) are not constant for the given substance and depend on temperature. Hence, the Van der Waals equation is only some approximation describing the state of a real gas.
We propose a way to modify the right part of eq.(1), substituting it for the part of the kinetic energy (T) of 1 mole of the substance (eq.4.31 in ) in real gas phase formed only by secondary effectons and deformons with nonzero impulse, affecting the pressure:
$$PV=\frac{2}{3}\overline{T}_{\text{kin}}=\frac{2}{3}V_0\frac{1}{Z}\underset{tr,lb}{}[\overline{n}_{ef}\frac{_1^3\left(\overline{E}_{1,2,3}^a\right)^2}{2m\left(\overline{v}_{ph}^a\right)^2}(\overline{P}_{ef}^a+\overline{P}_{ef}^b)+$$
$$+\overline{n}_d\frac{_1^3\left(\overline{E}_d^{1,2,3}\right)^2}{2m\left(v_s\right)^2}\overline{P}_d]_{tr,lb}$$
(3)
The contribution to pressure caused by primary quasiparticles as Bose-condensate with the zero resulting impulse is equal to zero also.
It is assumed when using such approach that for real gases the model of a system of weakly interacted oscillator pairs is valid. The validity of such an approach for water is confirmed by available experimental data indicating the presence of dimers, trimers and larger $`H_2O`$ clusters in the water vapor (Eisenberg and Kauzmann, 1975).
Water vapor has an intensive band in oscillatory spectra at $`\stackrel{~}{\nu }=200cm^1`$. Possibly, it is this band that characterizes the frequencies of quantum beats between ”acoustic” (a) and ”optic” (b) translational oscillations in pairs of molecules and small clusters. The frequencies of librational collective modes in vapor are absent.
The energies of primary gas quasiparticles $`(h\nu _a`$ and $`h\nu _b)`$ can be calculated on the basis of the formulae used for a liquid (Chapter 4 of or ).
However, to calculate the energies of secondary quasiparticles in (ā) and (b̄) states the Bose-Einstein distribution must be used for the case when the temperature is higher than the Bose-condensation temperature $`(T>T_0)`$ and the chemical potential is not equal to zero $`(\mu <0)`$. According to this distribution:
$$\begin{array}{c}\left\{\overline{E}^a=h\overline{\nu }^a=\frac{h\nu ^a}{\mathrm{exp}\left(\frac{h\nu ^a\mu }{kT}\right)1}\right\}_{tr,lb}\hfill \\ \left\{\overline{E}^b=h\overline{\nu }^b=\frac{h\nu ^b}{\mathrm{exp}\left(\frac{h\nu ^b\mu }{kT}\right)1}\right\}_{tr,lb}\hfill \end{array}$$
(4)
The kinetic energies of effectons $`(\overline{a})_{tr,lb}`$ and $`(\overline{b})_{tr,lb}`$ states are equal, only the potential energies differ as in the case of condensed matter.
All other parameters in basic equation (3) can be calculated as previously described .
2. New state equation for condensed matter
Using our eq.(4.3 from ) for the total internal energy of condensed matter $`(U_{\text{tot}})`$, we can present state equation in a more general form than (3).
For this end we introduce the notions of internal pressure $`(P_{\text{in}})`$, including all type of interactions between particles of matter and excluded molar volume $`(V_{\text{exc}})`$:
$$V_{\text{exc}}=\frac{4}{3}\pi \alpha ^{}N_0=V_0\left(\frac{n^21}{n^2}\right)$$
(5)
where $`\alpha ^{}`$ is the acting polarizability of molecules in condensed matter $`(`$see Part 1 of $`);N_0`$ is Avogadro number, and $`V_0`$ is molar volume.
The general state equation can be expressed in the following form:
$$P_{\text{tot}}V_{fr}=(P_{\text{ext}}+P_{\text{in}})(V_0V_{\text{exc}})=U_{ef}$$
(6)
where: $`U_{ef}=U_{\text{tot}}(1+V/T_{\text{kin}}^t)=U_{\text{tot}}^2/T_{\text{kin}}`$ is the effective internal energy and:
$$(1+V/T_{\text{kin}})=U_{\text{tot}}/T_{\text{kin}}=S^1$$
is the reciprocal value of the total structural factor $`(eq\mathrm{.2.46}aof[1]);P_{\text{tot}}=P_{\text{ext}}+P_{\text{in}}`$ is total pressure, $`P_{\text{ext}}`$ and $`P_{\text{in}}`$ are external and internal pressures; $`V_{fr}=V_0V_{\text{exc}}=V_0/n^2(`$see eq.5) is a free molar volume; $`U_{\text{tot}}=V+T_{\text{kin}}`$ is the total internal energy, V and $`T_{\text{kin}}`$ are total potential and kinetic energies of one mole of matter.
For the limit case of ideal gas, when $`P_{\text{in}}=0;V_{\text{exc}}=0`$; and the potential energy $`V=0`$, we get from (6) the Clapeyrone - Mendeleyev equation (see 1):
$$P_{\text{ext}}V_0=T_{\text{kin}}=RT$$
One can use equation of state (6) for estimation of sum of all types of internal matter interactions, which determines the internal pressure $`P_{\text{in}}`$:
$$P_{\text{in}}=\frac{U_{ef}}{V_{fr}}P_{\text{ext}}=\frac{n^2U_{\text{tot}}^2}{V_0T_{\text{kin}}}P_{\text{ext}}$$
(7)
where: the molar free volume: $`V_{fr}=V_0V_{\text{exc}}=V_0/n^2`$;
and the effective total energy: $`U_{ef}=U_{\text{tot}}^2/T_{\text{kin}}=U_{\text{tot}}/`$S.
For solids and most of liquids with a good approximation: $`P_{\text{in}}[P_{\text{ext}}1`$ atm. $`=10^5Pa]`$. Then from (7) we have:
$$P_{\text{in}}\frac{n^2U_{\text{tot}}}{V_0S}=\frac{n^2}{V_0}U_{\text{tot}}\left(1+\frac{V}{T_{\text{kin}}}\right)$$
(8)
where $`S=T_{\text{kin}}/U_{\text{tot}}`$ is a total structural factor; $`T_{\text{kin}}`$ and V are total kinetic and potential energies, respectively.
For example for 1 mole of water under standard conditions we obtain:
$`V_{\text{exc}}=8.4cm^3;V_{fr}=9.6cm^3;V_0=V_{\text{exc}}+V_{fr}=18cm^3;`$
$`P_{\text{in}}380000`$ atm. $`=3.810^{10}Pa`$ (1 atm. =10$`{}_{}{}^{5}Pa).`$
The parameters such as sound velocity, molar volume, and the positions of translational and librational bands in oscillatory spectra that determine $`U_{ef}(4.3)`$ depend on external pressure and temperature.
The results of computer calculations of $`P_{\text{in}}(eq.7)`$ for ice and water are presented on Fig. 1 a,b.
Polarizability and, consequently, free volume $`(V_{fr})`$ and $`P_{\text{in}}`$ in (6) depend on energy of external electromagnetic fields (see Part 1 of .
> Fig. 1. (a) Theoretical temperature dependence of internal pressure $`(P_{\text{in}})`$ in ice including the point of \[ice $``$ water\] phase transition; (b) Theoretical temperature dependence of internal pressure $`(P_{\text{in}})`$ in water. Computer calculations were performed using eq. (7).
The minima of $`P_{\text{in}}(T)`$ for ice at $`140^0`$ and $`50^0C`$ in accordance with eq.(9) correspond to the most stable structure of this matter, related to temperature transition. In water some kind of transition appears at $`35^0C`$, near physiological temperature.
There may exist conditions when the derivatives of internal pressure P$`_{\text{in}}`$ are equal to zero:
$$(a):\left(\frac{P_{\text{in}}}{P_{\text{ext}}}\right)_T=0\text{ and }(b):\left(\frac{P_{\text{in}}}{T}\right)_{P_{\text{ext}}}=0$$
(9)
This condition corresponds to the minima of potential energy, i.e. to the most stable structure of given matter. In a general case there may be a few metastable states when conditions (9) are fulfilled.
Equation of state (7) may be useful for the study of mechanical properties of condensed matter and their change under different influences.
Differentiation of (6) by external pressure gives us at T = const:
$$V_{fr}+\frac{P_{fr}}{P_{\text{ext}}}(P_{\text{ex}}+P_{\text{in}})+V_{fr}\frac{P_{\text{in}}}{P_{\text{ext}}}=\frac{P_{ef}}{P_{\text{ext}}}$$
(10)
Dividing the left and right part of (10) by free volume $`V_{fr}`$ we obtain:
$$\left(\frac{P_{\text{in}}}{P_{\text{ext}}}\right)_T=\left(\frac{P_{ef}}{P_{\text{ext}}}\right)_T\left[\begin{array}{c}1+\beta _T(P_{\text{ext}}+P_{\text{in}})\end{array}\right]_T$$
(11)
where: $`\beta _T=(V_{fr}/P_{\text{ext}})/V_{fr}`$ is isothermal compressibility. From (9) and (11) we derive condition for the maximum stability of matter structure:
$$\left(\frac{P_{ef}}{P_{\text{ext}}}\right)_T=1+\beta _T^0P_{\text{tot}}^{\text{opt}}$$
(12)
where: $`P_{\text{tot}}^{\text{opt}}=P_{\text{ext}}+P_{\text{in}}^{\text{opt}}`$ is the ”optimum” total pressure.
The derivative of (6) by temperature gives us at $`P_{\text{ext}}=`$ const:
$$P_{\text{tot}}\left(\frac{V_{fr}}{T}\right)_{P_{\text{ext}}}+V_{fr}\left(\frac{P_{\text{in}}}{T}\right)_{P_{\text{ext}}}=\left(\frac{U_{ef}}{T}\right)_{P_{\text{ext}}}=C_V$$
(13)
where
$`\left({\displaystyle \frac{V_{fr}}{T}}\right)_{P_{\text{ext}}}`$ $`=\left({\displaystyle \frac{V_0}{T}}\right)_{P_{\text{ext}}}{\displaystyle \frac{4}{3}}\pi N_0\left({\displaystyle \frac{\alpha ^{}}{T}}\right)_{P_{\text{ext}}}`$ (14)
$`\text{and }\left({\displaystyle \frac{V_{\text{tot}}}{T}}\right)_{P_{\text{ext}}}`$ $`={\displaystyle \frac{P_{\text{in}}}{T}}`$ (14a)
From our mesoscopic theory of refraction index (Part 1 of ) the acting polarizability $`\alpha ^{}`$ is:
$$\alpha ^{}=\frac{\left(\frac{n^21}{n^2}\right)}{\frac{4}{3}\pi \frac{N_0}{V_0}}$$
(15)
When condition (9b) is fulfilled, we obtain for optimum internal pressure $`(P_{\text{in}}^{\text{opt}})`$ from (13):
$$P_{\text{in}}^{\text{opt}}=C_V/\left(\frac{V_{fr}}{T}\right)_{P_{\text{ext}}}P_{\text{ext}}$$
(16)
or
$$P_{\text{in}}^{\text{opt}}=\frac{C}{V_{fr}\gamma }P_{\text{ext}},$$
(17)
where
$$\gamma =(V_{fr}/T)/V_{fr}$$
(18)
is the thermal expansion coefficient;
V<sub>fr</sub> is the total free volume in 1 mole of condensed matter:
$$V_{fr}=V_0V_{\text{exc}}=V_0/n^2$$
(19)
It is taken into account in (13) and (19) that
$$(V_{\text{exc}}/T)0$$
(20)
because, as has been shown earlier (Fig.25a of and Part 1 of ),
$$\alpha ^{}/T0$$
Dividing the left and right parts of (13) by $`P_{\text{tot}}V_{fr}=U_{ef}`$, we obtain for the heat expansion coefficient:
$$\gamma =\frac{C_V}{U_{ef}}\frac{1}{P_{\text{tot}}}\left(\frac{P_{\text{in}}}{T}\right)_{P_{\text{ext}}}$$
(21)
Under metastable states, when condition (9 b) is fulfilled,
$$\gamma ^0=C_V/U_{ef}$$
(22)
Putting (8) into (12), we obtain for isothermal compressibility of metastable states corresponding to (9a) following formula:
$$\beta _T^0=\frac{V_0T_{\text{kin}}}{n^2U_{\text{tot}}^2}\left(\frac{U_{ef}}{P_{\text{ext}}}1\right)$$
(23)
It seems that our equation of state (7) may be used to study different types of external influences (pressure, temperature, electromagnetic radiation, deformation, etc.) on the thermodynamic and mechanic properties of solids and liquids.
3. Vapor pressure
When a liquid is incubated long enough in a closed vessel at constant temperature, then an equilibrium between the liquid and vapor is attained.
At this moment, the number of molecules evaporated and condensed back to liquid is equal. The same is true of the process of sublimation.
There is still no satisfactory quantitative theory for vapor pressure calculation.
We can suggest such a theory using our notion of superdeformons, representing the biggest thermal fluctuations (see Table 1 and Introduction). The basic idea is that the external equilibrium vapor pressure is related to internal one $`(P_{\text{in}}^S)`$ with coefficient determined by the probability of cavitational fluctuations (superdeformons) in the surface layer of liquids or solids.
In other words due to excitation of superdeformons with probability $`(P_D^S)`$, the internal pressure $`(P_{\text{in}}^S)`$ in surface layers, determined by the total contributions of all intramolecular interactions turns to external one - vapor pressure $`(P_V)`$. It is something like a compressed spring energy realization due to trigger switching off.
For taking into account the difference between the surface and bulk internal pressure $`(P_{\text{in}})`$ we introduce the semiempirical surface pressure factor $`(q^S)`$ as:
$$P_{\text{in}}^S=q^SP_{\text{in}}P_{\text{ext}}=q^S\frac{n^2U_{\text{tot}}}{V_0S}P_{\text{ext}}$$
(24)
where: P$`_{\text{in}}`$ corresponds to $`eq.(7);S=T_{\text{kin}}/U_{\text{tot}}`$ is a total structure factor.
The value of surface factor $`(q^S)`$ for liquid and solid states is not the same:
$$q_{\text{liq}}^S<q_{\text{sol}}^S$$
(25)
> Fig. 2. a) Theoretical $`()`$ and experimental ($``$) temperature dependences of vapor pressure $`(P_{\text{vap}})`$ for ice (a) and water (b) including phase transition region. Computer calculations were performed using eq. (26).
Multiplying (24) to probability of superdeformons excitation we obtain for vapor pressure, resulting from evaporation or sublimation, the following formulae:
$$P_{vap}=P_{\text{in}}^SP_D^S=\left(q^S\frac{n^2U_{\text{tot}}^2}{V_0T_{\text{kin}}}P_{\text{ext}}\right)\mathrm{exp}\left(\frac{E_D^S}{kT}\right)$$
(26)
where:
$$P_D^S=\mathrm{exp}\left(\frac{E_D^S}{kT}\right)$$
(27)
is a probability of superdeformons excitation (see eqs. 3.37, 3.32 and 3.33).
We can assume, that the difference in the surface and bulk internal pressure is determined mainly by difference in total internal energy $`(U_{\text{tot}})`$ but not in kinetic one $`(T_k)`$. Then a pressure surface factor could be presented as:
$$q^S=\gamma ^2=(U_{\text{in}}/U_{\text{tot}})^2$$
where: $`\gamma =U_{\text{tot}}^S/U_{\text{tot}}`$ is the surface energy factor, reflecting the ratio of surface and bulk total energy.
Theoretical calculated temperature dependences of vapor pressure, described by (26) coincide very well with experimental ones for water at $`q_{\text{liq}}^S=3.1(\gamma _l=1.76)`$ and for ice at $`q_{\text{sol}}^S=18(\gamma _s=4.24)(`$Fig. 2).
The almost five-times difference between q$`{}_{}{}^{S}{}_{\text{sol}}{}^{}`$ and $`q_{\text{liq}}^S`$ means that the surface properties of ice differ from bulk ones much more than for liquid water.
The surface factors $`q_{\text{liq}}^S`$ and $`q_{\text{sol}}^S`$ should be considered as a fit parameters. The $`q^S=`$ $`\gamma ^2`$ is the only one fit parameter that was used in our hierarchic mesoscopic theory. Its calculation from the known vapor pressure or surface tension can give an important information itself.
4. Surface tension
The resulting surface tension is introduced in our mesoscopic model as a sum:
$$\sigma =(\sigma _{tr}+\sigma _{lb})$$
(28)
where: $`\sigma _{tr}`$ and $`\sigma _{lb}`$ are translational and librational contributions to surface tension. Each of these components can be expressed using our mesoscopic state equation (6, 7), taking into account the difference between surface and bulk total energies $`(q^S)`$, introduced in previous section:
$$\sigma _{tr,lb}=\frac{1}{\frac{1}{\pi }(V_{ef})_{tr,lb}^{2/3}}\left[\frac{q^SP_{\text{tot}}(P_{ef}V_{ef})_{tr,lb}P_{\text{tot}}(P_{ef}V_{ef})_{tr,lb}}{(P_{ef}+P_t)_{tr}+(P_{ef}+P_t)_{lb}+(P_{\text{con}}+P_{\text{cMt}})}\right]$$
(29)
where $`(V_{ef})_{tr,lb}`$ are volumes of primary tr and lib effectons, related to their concentration $`(n_{ef})_{tr,lb}`$ as:
$$(V_{ef})_{tr,lb}=(1/n_{ef})_{tr,lb};$$
$$r_{tr,lb}=\frac{1}{\pi }(V_{ef})_{tr,lb}^{2/3}$$
is an effective radius of the primary translational and librational effectons, localized on the surface of condensed matter; $`q^S`$ is the surface factor, equal to that used in eq.(24 -26)$`;[P_{\text{tot}}=P_{\text{in}}+P_{\text{ext}}]`$ is a total pressure, corresponding to eq.(6);$`(P_{ef})_{tr,lb}`$ is a total probability of primary effecton excitations in the (a) and (b) states:
$$(P_{ef})_{tr}=(P_{ef}^a+P_{ef}^b)_{tr}$$
$$(P_{ef})_{lb}=(P_{ef}^a+P_{ef}^b)_{lb}$$
$`(P_t)_{tr}`$ and $`(P_t)_{lb}`$ in (29) are the probabilities of corresponding transiton excitation;
$`P_{\text{con}}=P_{ac}+P_{bc}`$ is the sum of probabilities of $`\left[a\right]`$ and $`\left[b\right]`$ convertons; $`P_{\text{cMt}}=P_{ac}P_{bc}`$ is the probability of Macroconverton (see Introduction and Chapter 4).
The eq. (29) contains the ratio:
$$(V_{ef}/V_{ef}^{2/3})_{tr,lb}=l_{tr,lb}$$
(30)
where: $`l_{tr}=(1/n_{ef})_{tr}^{1/3}`$ and $`l_{lb}=(1/n_{ef})_{\text{lib}}^{1/3}`$ are the length of the ribs of the primary translational and librational effectons, approximated by cube.
Using (30) and (29) the resulting surface tension (28) can be presented as:
$$\sigma =\sigma _{tr}+\sigma _{lb}=\pi \frac{\begin{array}{c}P_{\text{tot}}(q^S1)\left[\begin{array}{c}(P_{ef})_{tr}l_{tr}+(P_{ef})l_{lb}\end{array}\right]\end{array}}{(P_{ef}+P_t)_{tr}+(P_{ef}+P_t)_{lb}+(P_{\text{con}}+P_{\text{cMt}})}$$
(31)
where translational component of surface tension is:
$$\sigma _{tr}=\pi \frac{P_{\text{tot}}(q^s1)(P_{ef})_{tr}l_{tr}}{(P_{ef}+P_t)_{tr}+(P_{ef}+P_t)_{lb}+(P_{\text{con}}+P_{\text{cMt}})}$$
(32)
and librational component of $`\sigma `$ is:
$$\sigma _{lb}=\pi \frac{P_{\text{tot}}(q^S1)(P_{ef})_{lb}l_{lb}}{(P_{ef}+P_t)_{lb}+(P_{ef}+P_t)_{lb}+(P_{\text{con}}+P_{\text{cMt}})}$$
(33)
Under the boiling condition when q$`{}_{}{}^{S}`$ 1 as a result of $`(U_{\text{tot}}^SU_{\text{tot}})`$, then $`\sigma _{tr},\sigma _{lb}`$ and $`\sigma `$ tends to zero. The maximum depth of the surface layer, which determines the $`\sigma _{lb}`$is equal to the length of edge of cube $`\left(l_{lb}\right)`$, that approximates the shape of primary librational effectons. It decreases from about 20 Å at 0$`{}_{}{}^{0}C`$ till about 2.5 Å at 100$`{}_{}{}^{0}C(`$see Fig. 7b of or Fig. 4b of ). Monotonic decrease of $`\left(l_{lb}\right)`$with temperature could be accompanied by nonmonotonic change of probabilities of \[lb/tr\] convertons and macroconvertons excitations (see comments to Fig. 7a of or to Fig 4a of ). Consequently, the temperature dependence of surface tension on temperature can display anomalies at definite temperatures. This consequence of our theory is confirmed experimentally (Adamson, 1982; Drost-Hansen and Lin Singleton, 1992).
The thickness of layer $`\left(l_{tr}\right)`$, responsible for contribution of translational effectons in surface tension $`\left(\sigma _{tr}\right)`$has the dimension of one molecule in all temperature interval for liquid water.
The results of computer calculations of $`\sigma `$(eq.31) for water and experimental data are presented at Fig.3.
> Fig. 3. Experimental $`(^{\mathrm{\_}\mathrm{\_}\mathrm{\_}})`$ and theoretical (- - -) temperature dependences of the surface tension for water, calculated from eq.(31).
It is obvious, that the correspondence between theory and experiment is very good, confirming in such a way the correctness of our model and Hierarchic concept in general.
5. Mesoscopic theory of thermal conductivity
Thermal conductivity may be related to phonons, photons, free electrons, holes and \[electron-hole\] pairs movement.
We will discuss here only the main type of thermal conductivity in condensed matter, related to phonons.
The analogy with the known formula for thermal conductivity ($`\kappa `$) in the framework of the kinetic theory for gas is used :
$$\kappa =\frac{1}{3}C_vv_s\mathrm{\Lambda }$$
(34)
where C<sub>v</sub> is the heat capacity of condensed matter, $`v_s`$ is sound velocity, characterizing the speed of phonon propagation in matter, and $`\mathrm{\Lambda }`$ is the average length of free run of phonons.
The value of $`\mathrm{\Lambda }`$ depends on the scattering and dissipation of phonons at other phonons and different types of defects. Usually decreasing temperature increases $`\mathrm{\Lambda }`$.
Different factors influencing a thermal equilibrium in the system of phonons are discussed. Among them are the so called U- and N- processes describing the types of phonon-phonon interaction. However, the traditional theories are unable to calculate $`\mathrm{\Lambda }`$ directly.
Mesoscopic theory introduce two contributions to thermal conductivity: related to phonons, irradiated by secondary effectons and forming secondary translational and librational deformons ($`\kappa _{sd}`$)<sub>tr,lb</sub> and to phonons, irradiated by $`a`$ and $`b`$ convertons $`[tr/lb]`$, forming the convertons-induced deformons $`(\kappa _{cd})_{ac.bc}`$:
$$\kappa =(\kappa _{sd})_{tr,lb}+(\kappa _{cd})_{ac.bc}=\frac{1}{3}C_vv_s[(\mathrm{\Lambda }_{sd})_{tr,lb}+(\mathrm{\Lambda }_{cd})_{ac,bc}]$$
(35)
where: free runs of secondary phonons (tr and lb) are represented as:
$$1/(\mathrm{\Lambda }_{sd})_{tr,lb}=1/(\mathrm{\Lambda }_{tr})+1/(\mathrm{\Lambda }_{lb})=(\overline{\nu }_d)_{tr}/v_s+(\overline{\nu }_d)_{lb}/v_s$$
consequently:
$$1/(\mathrm{\Lambda }_{sd})_{tr,lb}=\frac{v_s}{(\overline{\nu }_d)_{tr}+(\overline{\nu }_d)_{lb}}$$
(36)
and free runs of convertons-induced phonons:
$$1/(\mathrm{\Lambda }_{cd})_{ac,bc}=1/(\mathrm{\Lambda }_{ac})+1/(\mathrm{\Lambda }_{bc})=(\nu _{ac})/v_s+(\nu _{bc})/v_s$$
$$\text{consequently: }(\mathrm{\Lambda }_{sd})_{tr,lb}=\frac{v_s}{(\nu _d)_{tr}+(\nu _d)_{lb}}$$
(37)
The heat capacity: $`C_V=U_{\text{tot}}/T`$ can be calculated also from our theory (see Chapter 4 and 5).
> Fig. 4. Temperature dependences of total thermal conductivity for water and contributions, related to acoustic deformons and $`[lb/tr]`$convertons. The dependences were calculated, using eq. (37).
Quantitative calculations show that formula (35), based on our mesoscopic model, works well for water (Fig. 4). It could be used for any other condensed matter also if positions of translational and librational bands, sound velocity and molar volume for this matter at the same temperature interval are known.
The small difference between experimental and theoretical data can reflect the contributions of non-phonon process in thermal conductivity, related to macrodeformons, superdeformons and macroconvertons, i.e. big fluctuations.
6. Mesoscopic theory of viscosity for liquids and solids
The viscosity is determined by the energy dissipation as a result of medium (liquid or solid) structure deformation. Viscosity corresponding to the shift deformation is named shear viscosity. So- called bulk viscosity is related to deformation of volume parameters and corresponding dissipation. These types of viscosity have not the same values and nature.
The statistical theory of irreversible process leads to the following expression for shear viscosity (Prokhorov, 1988):
$$\eta =\text{ }nkT\tau _p+(\mu _{\mathrm{}}\text{ }\text{ }nkT)\tau _q$$
(38)
where n is the concentration of particles, $`\mu _{\mathrm{}}`$ is the modulus of instant shift characterizing the instant elastic reaction of medium, $`\tau _p`$ and $`\tau _q`$ are the relaxation times of impulses and coordinates, respectively.
However, eq.(38) is inconvenient for practical purposes due to difficulties in determination of $`\tau _p,\tau _q`$ and $`\mu _{\mathrm{}}`$.
Sometimes in a narrow temperature interval the empiric Ondrade equation is working:
$$\eta =A(T)\mathrm{exp}(\beta /T)$$
(39)
A(T) is a function poorly dependent on temperature.
A good results in study the microviscosity problem were obtained by combining the model of molecular rotational relaxation and the Kramers equation (Åkesson et al., 1991). However, the using of the fit parameters was necessarily in this case also.
We present here our mesoscopic theory of viscosity. To this end the dissipation processes, related to ($`AB)_{tr.lb}`$ cycles of translational and librational macroeffectons and (a,b)-convertons excitations were used. The same approach was employed for elaboration of mesoscopic theory of diffusion in condensed matter (see next section).
In contrast to liquid state, the viscosity of solids is determined by the biggest fluctuations: supereffectons and superdeformons, resulting from simultaneous excitations of translational and librational macroeffectons and macrodeformons in the same volume.
The dissipation phenomena and ability of particles or molecules to diffusion are related to the local fluctuations of the free volume $`(\mathrm{\Delta }v_f)_{tr,lb}`$. According to mesoscopic theory, the fluctuations of free volume and that of density occur in the almost macroscopic volumes of translational and librational macrodeformons and in mesoscopic volumes of macroconvertons, equal to volume of primary librational effecton at the given conditions. Translational and librational types of macroeffectons determine two types of viscosity, i.e. translational $`(\eta _{tr})`$ and librational $`(\eta _{lb})`$ ones. They can be attributed to the bulk viscosity. The contribution to viscosity, determined by (a and b)- convertons is much more local and may be responsible for microviscosity and mesoviscosity.
Let us start from calculation of the additional free volumes $`(\mathrm{\Delta }v_f)`$ originating from fluctuations of density, accompanied the translational and librational macrodeformons (macrotransitons).
For 1 mole of condensed matter the following ratio between free volume and concentration fluctuations is true:
$$\left(\frac{\mathrm{\Delta }v_f}{v_f}\right)_{tr,lb}=\left(\frac{\mathrm{\Delta }N_0}{N_0}\right)_{tr,lb}$$
(40)
where $`N_0`$ is the average number of molecules in 1 mole of matter
$$\text{and }(\mathrm{\Delta }N_0)_{tr,lb}=N_0\left(\frac{P_D^M}{Z}\right)_{tr,lb}$$
(41)
is the number of molecules changing their concentration as a result of translational and librational macrodeformons excitation.
The probability of translational and librational macroeffectons excitation (see eqs. 3.23; 3.24):
$$\left(\frac{P_D^M}{Z}\right)_{tr,lb}=\frac{1}{Z}\mathrm{exp}(\frac{ϵ_D^M}{kT})_{tr,lb}$$
(42)
where Z is the total partition function of the system (Chapter 4 of ).
Putting (41) to (40) and dividing to Avogadro number $`(N_0)`$, we obtain the fluctuating free volume, reduced to 1 molecule of matter:
$$\mathrm{\Delta }v_f^0=\frac{\mathrm{\Delta }v_f}{N_0}=\left[\frac{v_f}{N_0}\left(\frac{P_D^M}{Z}\right)\right]_{tr,lb}$$
(43)
It has been shown above (eq.19) that the average value of free volume in 1 mole of matter is:
$$v_f=V_0/n^2$$
Consequently, for reduced fluctuating (additional) volume we have:
$$(\mathrm{\Delta }v_f^0)_{tr,lb}=\frac{V_0}{N_0n^2}\frac{1}{Z}\mathrm{exp}(\frac{ϵ_D^M}{kT})_{tr,lb}$$
(44)
Taking into account the dimensions of viscosity and its physical sense, it should be proportional to the work (activation energy) of fluctuation-dissipation, necessary for creating the unit of additional free volume: $`(E_D^M/\mathrm{\Delta }v_f^0)`$, and the period of ($`AB)_{tr.lb}`$ cycles of translational and librational macroeffectons $`\tau _{AB},`$determined by the life-times of all intermediate states (eq.46).
In turn, the energy of dissipation should be strongly dependent on the structural factor (S): the ratio of kinetic energy of matter to its total internal energy. We assume here that this dependence for viscosity calculation is cubical: $`(T_k/U_{\text{tot}})^3=S^3`$.
Consequently, the contributions of translational and librational macrodeformons to resulting viscosity we present in the following way:
$$\eta _{tr,lb}^M=\left[\frac{E_D^M}{\mathrm{\Delta }v_f^0}\tau ^M\left(\frac{T_k}{U_{\text{tot}}}\right)^3\right]_{tr,lb}$$
(45)
where: reduced fluctuating volume $`(\mathrm{\Delta }v_f^0)`$ corresponds to (44); the energy of macrodeformons: $`[E_D^M=kT(\mathrm{ln}P_D^M)]_{tr,lb}`$.
The cycle-periods of the tr and lib macroeffectons has been introduced as:
$$\left[\tau ^M=\tau _A+\tau _B+\tau _D\right]_{tr,lb}$$
(46)
where: characteristic life-times of macroeffectons in A, B-states and that of transition state in the volume of primary electromagnetic deformons can be presented, correspondingly, as follows:
$$\left[\tau _A=\left(\tau _a\tau _{\overline{a}}\right)^{1/2}\right]_{tr,lb}\text{ and }\left[\tau _A=\left(\tau _a\tau _{\overline{a}}\right)^{1/2}\right]_{tr,lb}$$
(47)
$$\left[\tau _D=\left|(1/\tau _A)(1/\tau _B)\right|^1\right]_{tr,lb}$$
Using (47, 46 and 44) it is possible to calculate the contributions of $`\left(AB\right)`$ cycles of translational and librational macroeffectons to viscosity separately, using (45).
The averaged contribution of macroexcitations (tr and lb)in viscosity is:
$$\eta ^M=\left[\begin{array}{c}(\eta )_{tr}^M(\eta )_{lb}^M\end{array}\right]^{1/2}$$
(48)
The contribution of a and b convertons to viscosity of liquids could be presented in a similar to (44-48) manner after substituting the parameters of tr and lb macroeffectons with parameters of a and b convertons:
$$\eta _{ac,bc}=\left[\frac{E_c}{\mathrm{\Delta }v_f^0}\tau _c\left(\frac{T_k}{U_{\text{tot}}}\right)^3\right]_{ac,bc}$$
(49)
where: reduced fluctuating volume of (a and b) convertons $`(\mathrm{\Delta }v_f^0)_{ac,bc}`$ corresponds to:
$$(\mathrm{\Delta }v_f^0)_{ac,bc}=\frac{V_0}{N_0n^2}\frac{1}{Z}P_{ac,bc}$$
(50)
where: $`P_{ac}`$ and $`P_{bc}`$ are the relative probabilities of tr/lib interconversions between a and b states of translational and librational primary effectons (see Introduction and Chapter $`4);E_{ac}`$ and $`E_{bc}`$ are the excitation energies of (a and b) convertons correspondingly (see Chapter 4 of and );
Characteristic life-times for ac-convertons and bc-convertons $`[tr/lb]`$ in the volume of primary librational effectons (”flickering clusters”) could be presented as:
$$\begin{array}{c}\tau _{ac}=(\tau _a)_{tr}+(\tau _a)_{lb}=(1/\nu _a)_{tr}+(1/\nu _a)_{lb}\hfill \\ \tau _{bc}=(\tau _b)_{tr}+(\tau _b)_{lb}=(1/\nu _b)_{tr}+(1/\nu _b)_{lb}\hfill \end{array}$$
(51)
The averaged contribution of the both types of convertons in viscosity is:
$$\eta _c=(\eta _{ac}\eta _{bc})^{1/2}$$
(52)
This contribution could be responsible for microviscosity or better term: mesoviscosity, related to volumes, equal to that of primary librational effectons.
The resulting viscosity (Fig.5) is a sum of the averaged contributions of macrodeformons and convertons:
$$\eta =\eta ^M+\eta _c$$
(53)
> Fig. 5. Theoretical and experimental temperature dependences of viscosities for water. Computer calculations were performed using eqs. (44 - 53) and (4.3; 4.36).
The best correlation between theoretical and experimental data was achieved after assuming that only ($`\pi /2=2\pi /4)`$ part of the period of above described fluctuation cycles is important for dissipation and viscosity. Introducing this factor to equations for viscosity calculations gives up very good correspondence between theory and experiment in all temperature interval (0-100<sup>0</sup>C) for water (Fig.5).
As will be shown below the same factor, introducing the effective time of fluctuations \[$`\frac{\tau }{\pi /2}`$\], leads to best results for self-diffusion coefficient calculation.
In the classical hydrodynamic theory the sound absorption coefficient ($`\alpha `$) obtained by Stokes includes share $`(\eta )`$ and bulk $`(\eta _b)`$ averaged macroviscosities:
$$\alpha =\frac{\mathrm{\Omega }}{2\rho v_s^3}\left(\frac{4}{3}\eta +\eta _b\right),$$
(54)
where $`\mathrm{\Omega }`$ is the angular frequency of sound waves; $`\rho `$ is the density of liquid.
Bulk viscosity ($`\eta _b`$) is usually calculated from the experimental $`\eta `$ and $`\alpha `$. It is known that for water:
$$(\eta _b/\eta )3.$$
The viscosity of solids
In accordance with our model, the biggest fluctuations: supereffectons and superdeformons (see Introduction) are responsible for viscosity and diffusion phenomena in solid state. Superdeformons are accompanied by the emergency of cavitational fluctuations in liquids and the defects in solids. The presentation of viscosity formula in solids $`(\eta _s)`$ is similar to that for liquids:
$$\eta _S=\frac{E_S}{(\mathrm{\Delta }v_f^0)_S}\tau _S\left[\frac{T_k}{U_{\text{tot}}}\right]^3$$
(55)
where: reduced fluctuating volume, related to superdeformons excitation $`(\mathrm{\Delta }v_f^0)_s`$ is:
$$(\mathrm{\Delta }v_f^0)_S=\frac{V_0}{N_0n^2}\frac{1}{Z}P_S$$
(56)
where: $`P_s=(P_D^M)_{tr}(P_D^M)_{lb}`$ is the relative probability of superdeformons, equal to product of probabilities of tr and lb macrodeformons excitation (see $`42);E_s=kT\mathrm{ln}P_s`$ is the energy of superdeformons (see Chapter 4);
Characteristic cycle-period of $`(A^{}B^{})`$ transition of supereffectons is related to its life-times in A$`^{},`$Band transition D$`^{}`$states (see eq.46) as was shown in section 4.3:
$$\tau _S=\tau _A^{}+\tau _B^{}+\tau _D^{}$$
(56a)
The viscosity of ice, calculated from eq.(55) is bigger than that of water (eq.53) to about $`10^5`$ times. This result is in accordance with available experimental data.
7. Brownian diffusion
The important formula obtained by Einstein in his theory of Brownian motion is for translational motion of particle:
$$r^2=6Dt=\frac{kT}{\pi \eta a}t$$
(57)
and that for rotational Brownian motion:
$$\phi ^2=\frac{kT}{4\pi \eta a^3}t$$
(58)
where: a \- radius of spherical particle, much larger than dimension of molecules of liquid. The coefficient of diffusion D for Brownian motion is equal to:
$$D=\frac{kT}{6\pi \eta a}$$
(59)
If we take the angle $`\overline{\phi }^2=1/3`$ in (59), then the corresponding rotational correlation time comes to the form of the known Stokes- Einstein equation:
$$\tau =\frac{4}{3}\pi a^3\frac{1}{k}\left(\frac{\eta }{T}\right)$$
(60)
All these formulas (57 - 60) include macroscopic share viscosity $`(\eta )`$ corresponding to our (53).
8. Self-diffusion in liquids and solids
Molecular theory of self-diffusion, as well as general concept of transfer phenomena in condensed matter is extremely important, but still unresolved problem.
Simple semiempirical approach developed by Frenkel leads to following expression for diffusion coefficient in liquid and solid:
$$D=\frac{a^2}{\tau _0}\mathrm{exp}(W/kT)$$
(61)
where \[a\] is the distance of fluctuation jump; $`\tau _0(10^{12}÷10^{13})s`$ is the average period of molecule oscillations between jumps; W - activation energy of jump.
The parameters: a, $`\tau _0`$ and W should be considered as a fit parameters.
In accordance with mesoscopic theory, the process of self-diffusion in liquids, like that of viscosity, described above, is determined by two contributions:
a) the collective, nonlocal contribution, related to translational and librational macrodeformons $`(D_{tr,lb})`$;
b) the local contribution, related to coherent clusters flickering: \[dissociation/association\] of primary librational effectons (a and b)- convertons $`(D_{ac,bc})`$.
Each component of the resulting coefficient of self-diffusion (D) in liquid could be presented as the ratio of fluctuation volume cross-section surface: $`[\mathrm{\Delta }v_f^0]^{2/3}`$ to the period of macrofluctuation $`(\tau )`$. The first contribution to coefficient D, produced by translational and librational macrodeformons is:
$$D_{tr,lb}=\left[\left(\mathrm{\Delta }v_f^0\right)^{2/3}\frac{1}{\tau ^M}\right]_{tr,lb}$$
(62)
where: the surface cross-sections of reduced fluctuating free volumes (see eq.43) fluctuations in composition of macrodeformons (tr and lb) are:
$$(\mathrm{\Delta }v_f^0)_{tr,lb}^{2/3}=\left[\frac{V_0}{N_0n^2}\frac{1}{Z}\mathrm{exp}(\frac{ϵ_D^M}{kT})_{tr,lb}\right]^{2/3}$$
(63)
($`\tau ^M`$)<sub>tr,lb</sub> are the characteristic $`(AB)`$ cycle-periods of translational and librational macroeffectons (see eqs. 46 and 47).
The averaged component of self-diffusion coefficient, which takes into account both types of nonlocal fluctuations, related to translational and librational macroeffectons and macrodeformons, can be find as:
$$D^M=[(D)_{tr}^M(D)_{lb}^M]^{1/2}$$
(64)
The formulae for the second, local contribution to self-diffusion in liquids, related to (a and b) convertons $`(D_{ac,bc})`$ are symmetrical by form to that, presented above for nonlocal processes:
$$D_{ac,bc}=\left[(\mathrm{\Delta }v_f^0)^{2/3}\frac{1}{\tau _S}\right]_{ac,bc}$$
(65)
where: reduced fluctuating free volume of (a and b) convertons $`(\mathrm{\Delta }v_f^0)_{ac,bc}`$ is the same as was used above in mesoscopic theory of viscosity (eq.50):
$$(\mathrm{\Delta }v_f^0)_{ac,bc}=\frac{V_0}{N_0n^2}\frac{1}{Z}P_{ac,bc}$$
(66)
where:$`P_{ac}`$ and $`P_{bc}`$ are the relative probabilities of tr/lib interconversions between a and b states of translational and librational primary effectons (see Introduction and Chapter 4)
The averaged local component of self-diffusion coefficient, which takes into account both types of convertons (ac and bc) is:
$$D_C=[(D)_{ac}(D)_{bc}]^{1/2}$$
(67)
In similar way we should take into account the contribution of macroconvertons $`(D_{Mc})`$:
$$D_{Mc}=\left(\frac{V_0}{N_0n^2}\frac{1}{Z}P_{Mc}\right)^{2/3}\frac{1}{\tau _{Mc}}$$
(67a)
where: $`P_{Mc}=P_{ac}P_{bc}`$ is a probability of macroconverton excitation;
the life-time of macroconverton is:
$$\tau _{Mc}=(\tau _{ac}\tau _{bc})^{1/2}$$
(67b)
The cycle-period of $`(ac)`$ and $`(bc)`$ convertons are determined by the sum of life-times of intermediate states of primary translational and librational effectons:
$$\tau _{ac}=(\tau _a)_{tr}+(\tau _a)_{lb};\text{ and }\tau _{bc}=(\tau _b)_{tr}+(\tau _b)_{lb}$$
(67c)
The life-times of primary and secondary effectons (lb and tr) in a\- and b-states are the reciprocal values of corresponding state frequencies:
$$\text{[}\tau _a=1/\nu _a;\text{ }\tau _{\overline{a}}=1/\nu _{\overline{a}};\text{ and }\tau _b=1/\nu _b;\text{ }\tau \overline{_b}=1/\nu _{\overline{b}}\text{]}_{tr,lb}$$
(67d)
\[$`\nu _a`$ and $`\nu _b`$\]<sub>tr,lb</sub> correspond to eqs. 4.8 and 4.9; \[$`\nu _{\overline{a}}`$ and $`\nu _{\overline{b}}`$\]<sub>tr,lb</sub> could be calculated using eqs.2.54 and 2.55.
The resulting coefficient of self-diffusion in liquids (D) is a sum of nonlocal $`(D^M)`$ and local $`(D_c,D_{Mc})`$ effects contributions (see eqs.64 and 67):
$$D=D^M+D_c+D_{Mc}$$
(68)
The effective fluctuation-times were taken the same as in previous section for viscosity calculation, using the correction factor \[($`\pi /2)\tau ]`$.
> Fig. 6. Theoretical and experimental temperature dependences of self-diffusion coefficients in water. Theoretical coefficient was calculated using eq. 68.
Like in the cases of thermal conductivity, viscosity and vapor pressure, the results of theoretical calculations of self-diffusion coefficient coincide well with experimental data for water (Fig. 6) in temperature interval $`(0100^0C)`$.
The self-diffusion in solids
In solid state only the biggest fluctuations: superdeformons, representing simultaneous excitation of translational and librational macrodeformons in the same volumes of matter are responsible for diffusion and the viscosity phenomena. They are related to origination and migration of the defects in solids. The formal presentation of superdeformons contribution to self-diffusion in solids $`(D_s)`$ is similar to that of macrodeformons for liquids:
$$D_S=(\mathrm{\Delta }v_f^0)_S^{2/3}\frac{1}{\tau _S}$$
(69)
where: reduced fluctuating free volume in composition of superdeformons $`(\mathrm{\Delta }v_f^0)_S`$ is the same as was used above in mesoscopic theory of viscosity (eq.56):
$$(\mathrm{\Delta }v_f^0)_S=\frac{V_0}{N_0n^2}\frac{1}{Z}P_S$$
(70)
where: $`P_S=(P_D^M)_{tr}(P_D^M)_{lb}`$ is the relative probability of superdeformons, equal to product of probabilities of tr and lb macrodeformons excitation (see 42).
Characteristic cycle-period of supereffectons is related to that of tr and lb macroeffectons like it was presented in eq.(56a):
$$\tau _s=\tau _A^{}+\tau _B^{}+\tau _D^{}$$
(71)
The self-diffusion coefficient for ice, calculated from eq.69 is less than that of water (eq.53) to about $`10^5`$ times. This result is in accordance with available experimental data.
Strong decreasing of D in a course of phase transition: \[water $``$ ice\] predicted by our mesoscopic theory also is in accordance with experiment (Fig. 7).
Fig. 7. Theoretical temperature dependences of self-diffusion coefficients in ice.
All these results allow to consider our mesoscopic theory of transfer phenomena as a quantitatively confirmed one. They point that the ”mesoscopic bridge” between Micro- and Macro Worlds is wide and reliable indeed. It gives a new possibilities for understanding and detailed description of very different phenomena in solids and liquids.
One of the important consequences of our theory of viscosity and diffusion is the possibility of explaining numerous nonmonotonic temperature changes, registered by a number of physicochemical methods in various aqueous systems during the study of temperature dependences (, , , ; ; , $`[12]`$).
Most of them are related to diffusion or viscosity processes and may be explained by nonmonotonic changes of the refraction index, included in our equations: 44, 45, 50 for viscosity and eqs. 69, 70 for self-diffusion. For water these temperature anomalies of refraction index were revealed experimentally, using few wave lengths in the temperature interval $`395^0`$. They are close to Drost-Hansen temperatures. The explanation of these effects, related to periodic variation of primary librational effectons stability with monotonic temperature change was presented as comments to Fig. 7a of or Fig.4a of .
Another consequence of our theory is the elucidation of a big difference between librational $`\eta _{lb}(48)`$, translational $`\eta _{tr}(45)`$ viscosities and mesoviscosity, determined by $`[lb/tr]`$ convertons (49 and 52).
The effect of mesoviscosity can be checked as long as the volume of a Brownian particle does not exceed much the volume of primary librational effectons (eq. 15). If we take a Brownian particle, much bigger than the librational primary effecton, then its motion will reflect only averaged share viscosity (eq.53).
The third consequence of the mesoscopic theory of viscosity is the prediction of nonmonotonic temperature behavior of the sound absorption coefficient $`\alpha `$ (51). Its temperature dependence must have anomalies in the same regions, where the refraction index has.
The experimentally revealed temperature anomalies of (n) also follow from our theory as a result of nonmonotonic $`(ab)_{lb}`$ equilibrium behavior, stability of primary lb effectons and probability of \[lb/tr\] convertons excitation (see Discussion to Fig.7a of or to Fig.4a of ).
Our model predicts also that in the course of transition from the laminar type of flow to the turbulent one the share viscosity ($`\eta `$) will increases due to increasing of structural factor $`(T_k/U_{\text{tot}})`$ in eq. 45.
The superfluidity $`(\eta 0)`$ in the liquid helium could be a result of inability of this liquid at the very low temperature for translational and librational macroeffectons excitations, i.e. $`\tau ^M0`$.
In turn, it is a consequence of tending to zero the life-times of secondary effectons and deformons in eqs.(45), responsible for dissipation processes, due to their Bose-condensation and transformation to primary ones (see Chapter 12 of ). The polyeffectons, stabilized by Josephson’s junctions between primary effectons form the superfluid component of liquid helium.
9. Mesoscopic approach to proton conductivity in water, ice and other systems, containing hydrogen bonds
The numerous models of proton transitions in water and ice are usually related to migration of two types of defects in the ideal Bernal-Fouler structure :
1. Ionic defects originated as a result of $`2H_2O`$ dissociation to hydroxonium and hydroxyl ions:
$$2H_2OH_3O^++\text{OH}^{}$$
2. Orientational Bjerrum defects are subdivided to D (dopplet) and L (leer) ones.
D-defect (positive) corresponds to situation, when 2 protons are placed between two oxygen atoms, instead of the normal structure of hydrogen bond: $`O\mathrm{}HO`$ containing 1 proton.
L-defect (negative) corresponds to opposite anomaly, when even 1 proton between two oxygens is absent. Reorientation of dipole moment of $`H_2O`$ in the case of D- and L-defects leads to origination of charges:
$$q_B=q_D^+=q_L^{}=0.38e$$
(72)
The interrelation between the charge of electron (e), Bjerrum charge $`(q_B)`$ and ionic charge $`(e_I)(`$Onsager, Dupius, 1962) is:
$$e=e_I+q_B$$
(73)
The general approach to problem of proton transition takes into account both types of defects: ionic and orientational. It was assumed that orientational defects originate and annihilate in the process of continuous migration of ions H<sup>+</sup> and OH<sup>-</sup> through the water medium. Krjachko (1987) considers DL-pairs as a cooperative water cluster with linear dimensions of about $`15\AA `$ and with ”kink”. The Bjerrum’s DL-pair is a limit case of such model.
The protons conductivity in water must decrease with temperature increasing due to decreasing and disordering of water clusters and chains.
The kink-soliton model of orientational defects migration along the $`H_2O`$ chain was developed by Sergienko (1986). Mobility of ionic defects exceeds the orientational ones about 10 times.
But it is important to point out that the strong experimental evidence confirming the existence of just Bjerrum type orientational defects are still absent.
Our mesoscopic model of proton diffusion in ice, water and other hydrogen bonds containing systems includes following stages:
1. Ionization of water molecules in composition of superdeformons and ionic defects origination;
2. Bordering by $`H_3^+O`$ and $`HO^{}`$ the opposite surface-sides of primary librational effectons;
3. Tunneling of proton through the volume of primary effecton as a coherent water cluster (Bose-particle);
4. Diffusion of ions $`H_3^+O`$ and $`HO^{}`$in the less ordered medium between primary effectons can be realized in accordance with fluctuation mechanism described above in Section 8. The velocity of this stage is less than tunneling.
Transitions of protons and hydroxyl groups can occur also due to exchange processes (Antonchenko, 1991) like:
$$H_3^+O+H_2OH_2O+H_3^+O$$
(74)
$$H_2O+HO^{}\text{ }HO^{}+H_2O$$
(75)
The rate of ions transferring due to exchange is about 10 times more, than diffusion velocity, but slower than that, determined by tunneling jumps.
5. The orientational defects can originate as a result of $`H_2O`$ molecules rearrangements and conversions between translational and librational effectons in composition of superdeformons. Activation energy of superdeformons and macroconvertons in water is 10.2 kcal/M and about $`\mathrm{\hspace{0.17em}12}`$ kcal/M in ice (see 6.12; 6.13). The additional activation energy about 2-3 kcal/M is necessary for subsequent reorientation of surrounding molecules (Bjerrum, 1951).
Like the ionic defects, positive (D) and negative (L) defects can form a separated pairs on the opposite sides of primary effectons, approximated by parallelepiped. Such pairs means the effectons polarization.
Probability of $`H^+orHO^{}`$ tunneling through the coherent cluster - primary effecton in the (a)-state is higher than that in the (b)-state as far (see 1.30-1.32 of ):
$$[E_a=T_{\text{kin}}^a+V_a]<[E_b=T_{\text{kin}}^b+V_b]$$
(76)
where: $`T_{\text{kin}}^a=T_{\text{kin}}^b`$ are the kinetic energies of (a) and (b) states; $`E_bE_a=V_bV_a`$. is the difference between total and potential energies of these states.
In accordance with known theory of tunneling, the probability of passing the particle with mass $`(m)`$ through the barrier with wideness (a) and height $`(ϵ)`$ has a following dependence on these parameters:
$$\psi _a\mathrm{exp}\left(\frac{a}{b}\right)=\mathrm{exp}\left(\frac{a(2mϵ)^{1/2}}{\mathrm{}}\right)$$
(77)
where:
$$b=\mathrm{}/(2mϵ)^{1/2}$$
(78)
is the effective wave function fading length.
Parameter ($`b`$) is similar to wave B most probable amplitude $`(A_B)`$ with total energy $`E_B=ϵ`$ (see eq. 2.22 of ):
$$b=A_B=\mathrm{}/(2mE_B)^{1/2}$$
(79)
With temperature decreasing the $`(ab)_{tr,lb}`$ equilibrium of primary effectons shifts to the left:
$$K_{a}^{}{}_{}{}^{}{}_{b}{}^{}=(P_a/P_b)\mathrm{}$$
(80)
where $`P_a1`$and $`P_b0`$ are the thermoaccessibilities of (a) and (b) states of primary effectons (see eqs. 4.10-4.12). The linear dimensions of primary effectons of ice also tend to infinity at T$`0`$.
In water the tunneling stage of proton conductivity can be related to primary librational effectons only and their role increase with temperature decreasing. Dimensions of translational effectons in water does not exceed that of one molecule as it leads from our computer calculations.
Increasing of protons conductivity in ice with respect to water, in accordance with our model, is a consequence of participation of translational primary effectons in tunneling of $`\left[H^+\right]`$besides librational ones, as well as significant elevation of primary librational effectons dimensions. Increasing of the total contribution of tunneling process in protons migration in ice rise up their resulting transferring velocity comparing to water.
The external electric field induce:
a) redistribution of positive and negative charges on the surface of primary effectons determined by ionic defects and corresponding orientational defects;
b) orientation of polarized primary effectons in field, making quasi-continuous polyeffectons chains and that of the effectons orchestrated superclusters.
These effects create the conditions for relay mechanism of $`\left[H^+\right]`$ and $`H_3^+O`$ transmitting in the direction of electric field and $`\left[HO^{}\right]`$ in the opposite one. In accordance with our hierarchic model, the $`\left[H^+\right]`$ transition mechanism includes the alternation of tunneling, exchange and usual diffusion processes.
10. Regulation of pH and shining of water by electromagnetic and acoustic fields
In accordance with our model, water dissociation reaction:
$$H_2OH^++HO^{}$$
leading to increase of protons concentration is dependent on probability of $`\left[A_S^{}B_S^{}\right]`$transitions in supereffectons. This means that stimulation of $`\left[A_S^{}B_S^{}\right]`$ transitions (superdeformons) by ultrasound with resonant frequencies, corresponding to frequency of these transitions, should lead to decreasing of pH, i.e. to increasing the concentration of protons $`[H^+]`$.
The $`\left[A_SB_S\right]`$ transitions of supereffectons can be accompanied by origination of cavitational fluctuations (cavitational microbubbles). The opposite $`\left[B_SA_S\right]`$ transitions are related to the collapse of these microbubbles. As a result of this adiabatic process, water vapor in the bubbles is heated up to $`40006000^0`$K. The usual energy of superdeformons in water (Section 6.3):
$$ϵ_D^S=10.2\text{ kcal}/MRT^{}$$
(81)
correspond to local temperature $`T^{}5000^0`$K. For the other hand it is known, that even $`2000^0K`$ is enough already for partial dissociation of water molecules inside bubbles (about 0.01% of total amount of bubble water).
The variable pressure (P), generated by ultrasound in liquid is dependent on its intensity $`(I,wt/cm^2)`$ like:
$$P=(\rho v_sI)^{1/2}4.610^3(\text{atm)}$$
(82)
where $`\rho `$ is density of liquid; $`v_s`$ \- sound velocity $`(m/s)`$.
$`\left[A_S^{}B_S^{}\right]`$ transitions and cavitational bubbles origination can be stimulated also by IR radiation with frequency, corresponding to the activation energy of corresponding big fluctuations, described in mesoscopic theory by superdeformons and macroconvertons.
In such a way, using IR radiation and ultrasound it is possible to regulate a lot of different processes in aqueous systems, depending on pH and water activity.
The increasing of ultrasound intensity leads to increased cavitational bubble concentration. The dependence of the resonance cavity radius $`(R_{\text{res}})`$ on ultrasound frequency (f) can be approximately expressed as:
$$R_{\text{res}}=3000/f$$
(83)
At certain conditions the water placed in the ultrasound field, begins to shine in the region: $`300600nm`$ . This shining (sonoluminescense) is a consequence of electronic excitation of water ions and molecules in the volume of cavitational bubbles.
When the conditions of ultrasound standing wave exist, the number of bubbles and intensity of sonoluminescense is maximal.
The intensity of shining is nonmonotonicly dependent on temperature with maxima around $`15,\mathrm{\hspace{0.17em}30},\mathrm{\hspace{0.17em}45}`$ and $`65^0`$ . This temperature corresponds to extremes of stability of primary librational effectons, related to the number of $`H_2O`$ per effecton’s edge $`(\kappa )`$ (see comments to Fig. 7a of or to Fig 4a of ). An increase of inorganic ion concentration, destabilizing (a)-state of these effectons, elevate the probability of superdeformons and consequently, shining intensity.
The most probable reason of photon radiation is recombination of water molecules, turning it into exited state:
$${}_{}{}^{}OH\text{ }+H^+H_2O^{}H_2O+h\nu _p$$
(84)
Very different chemical reactions can be stimulated in the volume of cavitational fluctuation by the external fields. The optimal resonant parameters of these fields could be calculated using hierarchic theory.
We propose here that the reaction of water molecules recombination (84) could be responsible for coherent ”biophotons” radiation by cell’s and microbes cultures and living organisms in visible and ultraviolet (UV) range. The advances in biophoton research are described by Popp et al., 1992 .
In accordance to our model, the cell’s body filaments - microtubules (MTs) ”catastrophe” (cooperative reversible disassembly of MTs) is a result of the internal water cavitational fluctuations due to superdeformons excitation. Such collective process should be accompanied by dissociation and recombination (84) of part of water molecules, localized in the hollow core of microtubules, leading to high-frequency electromagnetic radiation (see: http://arXiv.org/abs/physics/0003045). The coherent biophotons in the infrared (IR) range are a consequence of $`(ab)_{tr,lb}`$transitions of the water primary effectons in microtubules.
We can see that lot of well working new theoretical models for different physical phenomena, based on our Hierarchic theory of condensed matter, can be elaborated. It means that this theory may serve as new convenient scientific language.
=======================================================================
> REFERENCES
>
> . Kaivarainen A. Hierarchic Concept of Matter and Field. Water, biosystems and elementary particles. New York, NY,1995, pp. 485.
>
> . Kaivarainen A. New Hierarchic Theory of Matter General for Liquids and Solids: dynamics, thermodynamics and mesoscopic structure of water and ice
>
> (see URL: http://www.karelia.ru/~alexk) and:
>
> . Kaivarainen A. Hierarchic Concept of Condensed Matter and its Interaction with Light: New Theories of Light Refraction, Brillouin Scattering and Mössbauer effect
>
> (see URL: http://www.karelia.ru/~alexk).
>
> . Blakemore J.S. Solid state physics. Cambridge University Press, Cambridge, N.Y. e.a, 1985.
>
> . Dote J.L., Kivelson D., Schwartz H. J.Phys.Chem. 1981, 85, 2169.
>
> . Drost-Hansen W. In: Colloid and Interface Science. Ed. Kerker M. Academic Press, New York, 1976, p.267.
>
> . Drost-Hansen W., Singleton J. Lin. Our aqueous heritage: evidence for vicinal water in cells. In: Fundamentals of Medical Cell Biology, v.3A, Chemistry of the living cell, JAI Press Inc.,1992, p.157-180.
>
> . Johri G.K., Roberts J.A. Study of the dielectric response of water using a resonant microwave cavity as a probe. J.Phys.Chem. $`\mathrm{\hspace{0.17em}1990},\mathrm{\hspace{0.17em}94},7386`$.
>
> . Aksnes G., Asaad A.N. Influence of the water structure on chemical reactions in water. A study of proton-catalyzed acetal hydrolysis. Acta Chem. Scand. $`1989,43,726734`$.
>
> . Aksnes G., Libnau O. Temperature dependence of esther hydrolysis in water. Acta Chem.Scand. $`1991,45,463467`$.
>
> . Käiväräinen A.I. Solvent-dependent flexibility of proteins and principles of their function. D.Reidel Publ.Co., Dordrecht, Boston, Lancaster, 1985, pp.290.
>
> . Käiväräinen A., Fradkova L., Korpela T. Separate contributions of large- and small-scale dynamics to the heat capacity of proteins. A new viscosity approach. Acta Chem.Scand. $`1993,47,456460`$.
>
> . Frontas’ev V.P., Schreiber L.S. J. Struct. Chem. (USSR$`)`$6(1966)512.
>
> . Antonchenko V.Ya. Physics of water. Naukova dumka, Kiev, 1986.
>
> . Guravlev A.I. and Akopjan V.B. Ultrasound shining. Nauka, Moscow, 1977.
>
> . Popp F.A., Li K.H. and Gu Q. Recent advances in biophoton research. Singapore: World Scientific, 1992.
|
warning/0003/astro-ph0003097.html
|
ar5iv
|
text
|
# Giant Cyclones in Gaseous Discs of Spiral Galaxies
## 1 Introduction
In previous papers , authors have shown that besides the spiral arms, well known for more than one and a half century in spiral galaxies as structures in the brightness distribution, there also exist structures, revealing themselves in the velocity field. The latter structures have the appearance of giant anticyclonic vortices: their rotation is opposite to the rotation of the galactic disc. The centers of these vortices are located near the zero points of the force field (Lagrange points $`L_4`$ and $`L_5`$, see ). In spiral galaxies these points are fixed in the close vicinity of the corotation circle, where the velocity of the rigid-body rotating spiral pattern coincides with the velocity of the differentially rotating disc. Being in such a position, the centers of the anticyclones turn out to be fixed with respect to the spiral arms and stationary if the spiral structure is stationary. If the spiral structure develops due to some instability in the disc, then the arms and vortices arise simultaneously and grow with the same growth rate . As a result a unified spiral–vortex structure forms.
The growth of the perturbation amplitude leads rather often to the saturation of the instability, after which a spiral–vortex structure becomes stationary. From the very beginning this scenario implies that spiral arms in discs of galaxies have a wave nature, that satisfies modern theoretical conceptions and agrees well with the results of the velocity field analysis in spiral galaxies (, ).
A typical rotation velocity field of a galactic disc is represented schematically in Fig. 1a. There is no way to close the trajectories of particles in the laboratory (inertial) reference frame by any radial component of the velocity (arrows radially directed in Fig. 1a). The situation is opposite in the rotating reference frame, where anticyclonic motion takes place under condition of appropriate azimuthal variation of the radial component (see Fig. 1b). Note that in the latter case it is not necessary to have large radial velocities to form a vortex near the circle corotating with the reference frame due to small values of the azimuthal velocities in this region.
It is not accidental that the radial velocity was chosen as a periodic function of the azimuth: in the spiral density wave such periodicity should take place both for the radial $`\stackrel{~}{V}_r`$ and for the azimuthal $`\stackrel{~}{V}_\varphi `$ perturbed velocities. Arrows for $`\stackrel{~}{V_\varphi }`$ are not shown in Fig.1. Taking into account the perturbed azimuthal velocity will not change the picture qualitatively, if
$$|\stackrel{~}{V_\varphi }/r|<|dV_{circ}/dr|,$$
(1)
where $`V_{circ}`$ $``$ $`V_{rot}\mathrm{\Omega }_{rf}r`$ is the circular velocity in the reference frame rotating with the angular velocity $`\mathrm{\Omega }_{rf}`$. Note that even a weak inequality is quite enough to form only anticyclonic vortices like in Fig. 1b (for more details see ). Under the opposite condition:
$$|\stackrel{~}{V_\varphi }/r|>|dV_{circ}/dr|$$
(2)
not only anticyclonic, but also cyclonic vortices can appear . In the next section we will consider both these cases in a schematic model of a galactic disc with a well-defined two-armed spiral structure (Grand-Design galaxy).
## 2 Qualitative description of anticyclone and cyclone formation in a disc of Grand-Design galaxy
Let us assume that two-armed Grand-Design spiral structure exists in a differentially rotating gaseous disc. Then the functions describing the perturbed surface density and velocity components, related to the density wave, may be approximated in the form:
$$\stackrel{~}{\sigma }(r,\phi )=C_\sigma (r)\mathrm{cos}[2\phi F_\sigma (r)],$$
(3)
$$\stackrel{~}{V}_r(r,\phi )=C_r(r)\mathrm{cos}[2\phi F_r(r)],$$
(4)
$$\stackrel{~}{V}_\phi (r,\phi )=C_\phi (r)\mathrm{cos}[2\phi F_\phi (r)].$$
(5)
As we consider the momentary picture of the spiral structure, a time dependence is absent in the equations above.
It is natural to ask, whether it is possible to use expressions (3)–(5) for the description of the dynamics of spiral density waves with a finite amplitude. We have checked it in , by the example of the galaxy NGC 3631, where cyclones have been found just now (see Section III below). The operating sequence of the check was the following. First, we analysed the azimuthal expansion in the harmonic series of the observed surface brightness of this galaxy at different galactocentric radii. A domination of the second Fourier harmonic over others (see Fig. 1 in ) proved the correctness of the approximation expressed in Eq.(3).
Then the observed line-of-sight velocity field $`V^{obs}(r,\phi )`$ was analysed. This velocity can be expressed as (see , ):
$$V^{obs}(r,\phi )=V_s+V_\phi (r,\phi )\mathrm{cos}\phi +V_r(r,\phi )\mathrm{sin}\phi \mathrm{sin}i+V_z(r,\phi )\mathrm{cos}i.$$
(6)
Here $`i`$ is an inclination angle (angle between the rotation axis of galaxy and line of sight) and $`V_s`$ is systemic velocity of a galaxy (velocity of the center of mass of the galaxy).
Within the frame of the present model
$$V_r(r,\phi )=\stackrel{~}{V}_r(r,\phi ),$$
$$V_\phi (r,\phi )=V_{rot}(r)+\stackrel{~}{V}_\phi (r,\phi ),$$
(7)
$$V_z(r,\phi )=\stackrel{~}{V}_z(r,\phi )=C_z(r)\mathrm{cos}[2\phi F_z(r)],$$
where $`\stackrel{~}{V}_r`$, $`\stackrel{~}{V}_\phi `$ are determined by (4) and (5). Substituting (7) in (6) one obtains the model representation of the Fourier expansion for $`V^{obs}(r,\phi )`$.
Equating the coefficients of the model expansion to the coefficients calculated from observational data,
$$V^{obs}(r,\phi )=V_s+\underset{n}{}\left(a_n^{obs}\mathrm{cos}n\phi +b_n^{obs}\mathrm{sin}n\phi \right),$$
(8)
we obtain the following set of relations between unknown characteristics of the vector velocity field of the galaxy and the Fourier coefficients of the observed line-of-sight velocity field:
$$V_{rot}+\frac{1}{2}(C_r\mathrm{sin}F_r+C_\phi \mathrm{cos}F_\phi )=a_1^{obs},$$
(9)
$$C_r\mathrm{cos}F_r+C_\phi \mathrm{sin}F_\phi =2b_1^{obs},$$
(10)
$$C_z\mathrm{cos}F_z\mathrm{cos}i=a_2^{obs},$$
(11)
$$C_z\mathrm{sin}F_z\mathrm{cos}i=b_2^{obs},$$
(12)
$$C_r\mathrm{sin}F_r+C_\phi \mathrm{cos}F_\phi =2a_3^{obs},$$
(13)
$$C_r\mathrm{cos}F_r+C_\phi \mathrm{sin}F_\phi =2b_3^{obs}.$$
(14)
From (9)-(14) we see that the perturbed velocity components $`\stackrel{~}{V_r,}\stackrel{~}{V_\phi }`$ and $`\stackrel{~}{V_z}`$, written in the form of (4), (5) and (7) give contributions only in the first three harmonics of the Fourier expansion. This is really observed in the case of NGC 3631 (see Fig. 3 in ), where the first three Fourier harmonics of the line-of-sight velocity field dominate over the higher-order harmonics. This gives evidence for the correctness of the representation of the perturbed velocity components in the form (4), (5) and (7) for the galaxy NGC 3631.
To draw schematically the spiral–vortex structure in different reference frames we need relations between the phases $`F_r(r),F_\phi (r)`$ and $`F_\sigma (r)`$. These relations were derived in from the set of linearized equations. In spite of this simplification, an excellent qualitative agreement was found between the phase of the ”modified” third harmonic (see below) of the observed line-of-sight velocity field and $`F_\sigma (r)`$ (see Fig. 5 in ), which enables to use the phase relations obtained in for a qualitative description of the residual velocities. The modified third harmonic is a two-armed spiral with phase $`F_3\pi /2,`$ i.e. $`\mathrm{cos}(2\phi F_3+\pi /2),`$ where $`F_3`$ is the phase of the original radial-velocity-field third harmonic.
Below a list of these relations is presented.
a) Relations between the phases of the radial velocity $`\stackrel{~}{V}_r`$ and the surface density perturbation $`\stackrel{~}{\sigma }`$.
* On the corotation circle $`r=r_c`$ :
$$F_r=F_\sigma \frac{\pi }{2},$$
(15)
whence
$$\stackrel{~}{V}_r(r_c,\phi )=C_r(r_c)\mathrm{sin}[2\phi F_\sigma (r_c)],$$
(16)
i.e. on the corotation circle the extrema of $`\stackrel{~}{V}_r`$ coincide with the zeroes of $`\stackrel{~}{\sigma }`$ and vice versa.
* Inside the corotation circle, $`r<r_c`$,
$$F_r=F_\sigma +\pi ,$$
(17)
whence
$$\stackrel{~}{V}_r(r,\phi )=C_r(r)\mathrm{cos}[2\phi F_\sigma (r)],r<r_c,$$
(18)
i.e. inside the corotation circle the maxima of $`\stackrel{~}{V}_r`$ coincide with the minima of $`\stackrel{~}{\sigma }`$ and vice versa.
* Outside the corotation circle, $`r>r_c`$,
$$F_r=F_\sigma ,$$
(19)
whence
$$\stackrel{~}{V}_r(r,\phi )=C_r(r)\mathrm{cos}[2\phi F_\sigma (r)],r>r_c,$$
(20)
i.e. outside the corotation circle the maxima and minima of $`\stackrel{~}{V}_r`$ coincide, respectively, with the maxima and minima of $`\stackrel{~}{\sigma }`$.
b) Relation between the phases of the azimuthal velocity $`\stackrel{~}{V}_\phi `$ and the surface density $`\stackrel{~}{\sigma }`$.
* In all cases, both on the corotation circle, $`r=r_c`$, and out of it, $`rr_c`$, we have
$$F_\phi =F_\sigma +\frac{\pi }{2},$$
(21)
therefore
$$\stackrel{~}{V}_\phi (r,\phi )=C_\phi (r)\mathrm{sin}[2\phi F_\sigma (r)],$$
(22)
i.e. the extrema of $`\stackrel{~}{V}_\phi `$ coincide with the zeroes of $`\stackrel{~}{\sigma }`$ and vice versa throughout the disc.
Using the expressions (16), (18), (20), and (22) we can draw a scheme for the velocity field with vortices overlayed on the scheme of the surface density distribution (Fig. 2). In Fig. 2b we can see the presence of cyclones and anticyclones in the field of residual velocity (that corresponds a reference frame locally corotating with the disc at each radius, $`V_{circ}(r)=0`$). This fact by itself does not depend on the amplitude of the spiral density wave, although the sizes of the vortices depend on this amplitude.
Lindblad and Langebartel were the first to calculate the field of displacements of the stars in the gravitational potential of a bar. The form of the star displacement field resembles the system of two cyclones and two anticyclones, which is similar our schematic Fig. 2b. Refering on Lynden- Bell has noted the existence of cyclonic and anticyclonic trajectories in Fig. 9 of. A similar picture of four vortices in the velocity field of a gaseous disc – two cyclones on the bar and two anticyclones between the spiral arms – was obtained in a recent paper .
In the reference frame rotating with the angular velocity $`\mathrm{\Omega }_{ph}`$ of the two-armed spiral pattern the picture of vortices differs qualitatively from the previous one. In the case (1) when the gradient of the perturbed azimuthal velocity is less than the circular velocity gradient, i.e. the anticyclonic shear dominates over the cyclonic one, the appearance of cyclones is impossible – we can see only two anticyclones with centers near the corotation circle between two spiral arms (Fig. 2c). In the opposite case (2) when the gradient of the perturbed azimuthal velocity dominates over the circular velocity gradient, the location of the cyclones is determined by the place of this domination. Cyclones can survive almost in the same places as in Fig. 2b, i.e. on the corotation circle, but their sizes will be smaller. There are also two other variants when the cyclone center moves either inside or outside the corotation circle along the zero lines of the radial velocity. Finally, both last these variants can be realized simultaneously. In this case four cyclones exist in the velocity field as is shown in Fig. 2d.
## 3 Discovery of Giant Cyclones in the Gaseous Disc of the Spiral Galaxy NGC 3631 with the 6m telescope in SAO
Finding giant vortices in a spiral galaxy requires the reconstruction of the vector velocity field in the reference frame rotating with the angular velocity of the spiral pattern. Thus, first of all the corotation radius should be determined. In this paper for NGC 3631 we use the position of the corotation derived in on the base of the Fourier analysis of the observed line-of-sight velocity field. The observations of NGC 3631 were carried out with the scanning Fabry-Perot interferometer with Russian 6m telescope. The details of observations were described in .
It is necessary to determine $`\stackrel{~}{V}_r`$ and $`V_\phi =V_{rot}+\stackrel{~}{V}_\phi `$, i.e. to find five unknown functions: $`V_{rot}(r)`$, $`C_r(r)`$, $`C_\phi (r)`$, $`F_r(r)`$, $`F_\phi (r)`$ (see Eqs. (4), (5), and (7)). These five functions are connected with the characteristics of the observed velocity field by the four relations (9), (10), (13), (14). An additional condition required to close the system should have a theoretical origin. Unfortunately, up to now a reliable condition valid for any density wave amplitude is not available. Several possibilities discussed in the literature , have limited applicability. To overcome this difficulty, we propose the following approach.
Among the functions listed above, $`V_{rot}(r)`$ could be the most reliably estimated on the base of independent observational data. For this purpose the equilibrium condition of the gaseous disc rotating in a gravitational potential $`\mathrm{\Psi }`$ can be used
$$\frac{V_{rot}^2}{r}=\frac{\mathrm{\Psi }}{r}.$$
(23)
The latter is determined from the mass distribution in a galaxy or its surface brightness maps assuming a costant mass-to-light ratio to be known. To calculate the right hand side of the equation (23) we use a three-component dynamical model of a spiral galaxy similar to . In spite of the roughness of the model, the $`V_{rot}`$ curve found by this method is in the region between $`(a_1)_{\mathrm{min}}`$ and $`(a_1)_{\mathrm{max}}`$ (see Fig. 3; a peculiarity of the thick curve is described below).
The same result is obtained by another way. From equation (9) it follows that the difference $`|a_1V_{rot}|`$ cannot exceed the amplitudes $`C_r`$ and $`C_\phi `$, which, in turn, are connected by equations (10), (13) and (14) with the Fourier coefficients $`b_1^{obs}`$, $`a_3^{obs}`$, and $`b_3^{obs}`$, determined from our observations. In Fig. 4 the radial behaviour of $`b_3^{obs}b_1^{obs}`$ $`=`$ $`C_r\mathrm{cos}F_r`$ and $`b_3^{obs}+b_1^{obs}`$ $`=`$ $`C_\phi \mathrm{sin}F_\phi `$ is demonstrated. The extrema of these functions allow to estimate the amplitudes $`C_r`$ and $`C_\phi `$. According to Fig. 4, one can conclude that in NGC 3631 a maximum value of the amplitude of the residual velocities can reach at $`60`$ km/s, i.e.
$$|a_1V_{rot}|_{\mathrm{max}}60\mathrm{km}/\mathrm{s}.$$
(24)
The conditions (23) and (24) do not allow to calculate the function $`V_{rot}(r)`$ exactly. Nevertheless, they set limits on the variations of both the amplitude and the form of $`V_{rot}(r)`$. Within these limits, we choose a set of trial curves (Fig. 3) and analyse the velocity field obtained from Eqs. (9)–(14) for a given $`V_{rot}(r)`$.
The general views of the restored velocity fields are qualitatively similar for the full set of the rotation curves. In particular, in all cases, except the most extreme ones, distinct areas of anticyclonic as well as cyclonic shear flows exist. However, only rare examples demonstrate the presence of areas with trapped fluid particles moving along stream lines bounded by closed separatrices. Such behavior could be a natural consequence of the quasi-stationary character of the density wave. It is an argument in favour of choosing just these examples of $`V_{rot}(r)`$ as the best-corresponding to the real dynamical model of the galaxy. The best example of $`U_{rot}`$ from the point of view of the stationarity of the stream lines is shown by the thick curve in Fig.3. In other examples the vortices have not close separatrices and their stream lines are spirals - strongly inwards or outwards.
In Fig. 5 we can see the amplitudes of radial ($`\stackrel{~}{V}_r`$) and azimuthal ($`\stackrel{~}{V}_\phi `$) components of the residual velocity. These amplitudes, $`C_r`$ and $`C_\phi `$ were obtained from four equations (9), (10), (13), and (14) after substituting in them $`V_{rot}(r)`$ in the form drawn in Fig. 3 by the thickest line.
Using the functions $`C_r`$ and $`C_\phi `$ in Fig.5 let us show, that the necessary condition (2) of the cyclone formation is satisfied in the galaxy NGC 3631 outside and inside the corotation circle and is not satisfied on the corotation circle. In the reference frame, rotating with the angular velocity of the spiral pattern ($`\mathrm{\Omega }`$ $`=`$ $`\mathrm{\Omega }_{ph}`$), the inequality (2) can be rewritten in the form
$$\left|\frac{\stackrel{~}{V}_\phi }{r}\right|>\left|\frac{dV_{circ}}{dr}\right|\left|\frac{d}{dr}(V_{rot}\mathrm{\Omega }_{ph}r)\right|.$$
(25)
Substituting here the expression (5) for $`\stackrel{~}{V}_\phi `$, we obtain:
$$\left|C_\phi ^{}\mathrm{cos}(2\phi F_\phi )+F_\phi ^{}C_\phi \mathrm{sin}(2\phi F_\phi )\right|>\left|V_{circ}^{}\right|,$$
(26)
where ”prime” denotes the derivative with respect to $`r`$.
In order that cyclones should be possible the condition (26) should be fulfilled at least near the center of the vortex. In the centers
$$V_r=\stackrel{~}{V}_r=C_r\mathrm{cos}(2\phi F_r)=0.$$
(27)
This equation gives a possibility to reduce the left hand side of inequality (26) to a form more suitable for an estimation.
In the vicinity of the center of the cyclone in residual velocities (Fig. 2 b) on the corotation circle from (15) and (21) we have $`F_\phi `$ $`=`$ $`F_r+\pi `$. Thus with the help of (27), the condition (26) is reduced to
$$\left|F_\phi ^{}C_\phi \right|>\left|V_{circ}^{}\right|.$$
(28)
As follows from Fig. 5 b at the corotation radius ($`r`$ $``$ $`42`$ arcsec) $`C_\phi `$ $``$ $`0`$. Hence, as the value of $`F_\phi ^{}`$ is very small the inequality (28) cannot be fulfilled.
In the case of ”external” cyclones (see Fig. 2 d) from (19) and (21) we have $`F_\phi `$ $`=`$ $`F_r+\pi /2`$. Then using (27) from (26) we obtain
$$\left|C_\phi ^{}\right|>\left|V_{circ}^{}\right|.$$
(29)
This relation explains the coincidence of the location of the external cyclones in NGC 3631 (52 $`<`$ $`r`$ $`<`$ 59 arcsec in Fig. 7) with the region of rapid growth of $`C_\phi `$ (50 $`<`$ $`r`$ $`<`$ 60 arcsec in Fig. 5b).
In the case of the ”internal” cyclones from (17) and (21) it follows $`F_\phi `$ $`=`$ $`F_r\pi /2`$. Then in this case the condition (26) also has the form (29). According to Fig. 5b it is fulfilled in the region 22 $`<`$ $`r`$ $`<`$ 31 arcsec, which really contains the internal cyclone region 25 $`<`$ $`r`$ $`<`$ 30 arcsec (Fig. 7).
The vector velocity field of the galaxy NGC 3631, when $`V_{rot}(r)`$ is taken in the form presented in Fig. 3 by the thickest line, is shown in Figs. 6 and 7. In the figures some streamlines of the velocity fields are presented. If the velocity field is stationary the trajectories of fluid particles coincide with the correspondent streamlines. One can see that the residual velocity field of NGC 3631 (Fig. 6) corresponds to that of Fig. 2b and the full velocity field in the reference frame of the spirals (Fig. 7) corresponds to either of two cases presented in Fig. 2d when the centers of two cyclones lie outside the corotation circle. Internal cyclones in the case shown in Fig. 7 are not enveloped by close streamlines.
## 4 Conclusions
I. An analysis of velocity fields in Grand-Design galaxies shows that:
1) the field of residual velocities contains two cyclones and two anticyclones with centers on the corotation circle;
2) in the reference frame rotating with the spiral pattern the velocity field belongs to one of two types:
a) under the condition (1) it contains only two anticyclones with centers near the corotation circle;
b) under the condition (2) besides the anticyclones the field contains also either two or four cyclones. In the former case the cyclone centers lay either on, or outside, or inside the corotation circle. In the latter case, two pairs of cyclones appear with centers inside and outside the corotation circle.
II. Our analysis of the velocity field data obtained by our team at the 6m telescope in the Special Astrophysical Observatory of the Russian Academy of Sciences shows that the Grand-Design galaxy NGC 3631 belongs to the type b).
## 5 Acknowledgement
We thank V. I. Arnold, B. V. Chirikov, V. L. Polyachenko, M. I. Rabinovich, Ya. G. Sinai, and specially G. Contopoulos for fruitful discussions. We thank also J.Bolesteix for his kind placing at our disposal a collection of interferometric filters.
This work was performed under partial financial support of RFBR grant 99–02–18432, grant ”Leading Scientific Schools” 96–15–96648, and the grant ”Fundamental Space Researches. Astronomy” for the 1999 1.2.3.1 and 1.7.4.3.
Figure Captions
Figure 1. A scheme for the anticyclonic vortices formation in a reference frame rotating with an arbitrary angular velocity. Unperturbed rotation velocities and radial perturbed velocities are marked by solid arrows and dashed arrows correspondingly. a) In the laboratory system of coordinates the rotation velocity of the disc varies with radius without change of sign. In this case the radial perturbed velocity can not provide a vortex formation despite the regular change of its sign along azimuth. b) In a rotating reference frame, where the dotted circle is at rest, the same field of radial perturbed velocities as in Fig.1a participates in the anticyclone formation.
Figure 2. A scheme for the velocity field with vortices in different reference frames. Vectors of unperturbed rotation velocity field are shown by solid arrows, radial and azimuthal components of the residual velocity field — by dashed arrows. Solid curves of different thickness – the thickest, the thinnest, and intermediate – tracing the azimuthal locations of the maxima, minima and zero values of the perturbed surface density defined at every radius are denoted by $`\stackrel{~}{\sigma }_{max}`$, $`\stackrel{~}{\sigma }_{min}`$, and $`0(\stackrel{~}{\sigma })`$ respectively. A and C denote anticyclones and cyclones correspondingly. a) The angular velocities of a disc ($`\mathrm{\Omega }(r)`$) and of the spiral pattern ($`\mathrm{\Omega }_{ph}`$) in the laboratory reference frame. The dashed line represents the corotation circle. b) The residual velocity field, which is the result of the subtraction of the rotation velocity from the full velocity field. Dotted lines show the boundaries of vortices. c) The velocity field in the reference frame rotating with the angular velocity $`\mathrm{\Omega }_{ph}`$ of the two-armed spiral pattern. Only two anticyclones can be seen in the vicinity of the corotation circle, when the gradient of perturbed azimuthal velocity is lower than the rotation velocity gradient. d) The velocity field in the same reference frame as in the previous figure, but in the case, when the gradient of the perturbed azimuthal velocity exceeds the rotation velocity gradient. In this case one can see two cyclones and two anticyclones. Depending on the residual velocity field geometry, the domination of the cyclonic shear over the anticyclonic one can take place either near corotation, or inside (or outside) the corotation circle. Figure d) demonstrates one more possibility when four cyclones exist simultaneously - two inside and two outside the corotation circle together with two anticyclones situated at corotation.
Figure 3. Examples of trial curves used to represent the rotation curve ($`V_{rot}(r)`$) in NGC 3631 are shown by solid lines together with the observed behaviour of $`a_1^{obs}(r)`$ (triangles). The thickest line marks the rotation curve corresponding to a quasi-stationary gaseous disk.
Figure 4. The radial dependence of $`b_3^{obs}b_1^{obs}`$ $`=`$ $`C_r\mathrm{cos}F_r`$ and $`b_3^{obs}+b_1^{obs}`$ $`=`$ $`C_\phi \mathrm{sin}F_\phi `$ observed in the spiral galaxy NGC 3631. An estimation of the amplitudes of the velocity components from the extrema of the presented functions gives $`max(C_r)`$ $``$ $`max(C_\phi )`$ $``$ 60 km/s.
Figure 5. Some parameters of the velocity field which correspond to a quasi-stationary regime of the vortex structure.
(a) The calculated amplitude of the radial velocity as a function of $`r`$.
(b) The same for the azimuthal velocity (squares) with overlaid profile of $`\left|V_{circ}(r)\right|`$.
Figure 6. The residual velocity field in the plane of the gaseous disc of the galaxy NGC 3631. One can see two cyclones and two anticyclones. Solid lines show some streamlines of the residual velocity field.
Figure 7. The full velocity field of the same galaxy in the reference frame rotating with the angular velocity of the two-armed spiral pattern. Overlayed squares show the position of maxima of the second Fourier harmonics of brightness map of NGC 3631 in H<sub>α</sub> line. Thin circle shows the position of the corotation. Presented streamlines (solid curves) were calculated using reverse ”time” direction that allows to reveal separatrices evidently. Two cyclones and two anticyclones are enveloped by close streamlines. Other two cyclones are not enveloped by close streamlines. The anticyclone centers lie almost at corotation; the cyclone centers lie outside corotation Cyclones lie in the vicinity of spiral arms, anticyclones are situated between them. The agreement of the vortices position with theoretical predictions (scheme in Fig. 2d) is very good.
|
warning/0003/nucl-th0003066.html
|
ar5iv
|
text
|
# Quark-Hadron Phase Transitions in Young and Old Neutron Stars
## Figure Captions
Figure 1: Left panels: The volume fraction of hadrons as a function of density in units of $`n_0`$ within the quark-hadron mixed phase for cold, catalyzed matter ($`s=0,Y_{\nu _e}=0`$) without hyperons (npQ). Three choices for the parameters $`\zeta `$ and $`\xi `$ in the Müller-Serot (MS) hadronic Lagrangian are illustrated, and the upper panel refers to the Nambu Jones-Lasinio (NJL) model and the lower panel to the MIT bag model with $`B=200`$ MeV fm<sup>-3</sup>. Right panels: The volume fraction of hadrons in the star’s center as a function of stellar mass for the same configurations and quark models.
Figure 2: The same as Figure 1, except that results are compared for three choices of the bag constant $`B`$ (in units of MeV fm<sup>-3</sup>) in the MIT bag model. Hyperons are ignored in the top panels (npQ) and included in the bottom panels (npHQ). The parameters $`\zeta =\xi =0`$ in the Müller-Serot (MS) hadronic Lagrangian are chosen.
Figure 3: Pressure versus density in units of $`n_0`$ for three representative snapshots during the evolution of a proto-neutron star. The top (bottom) panels display results without (with) hyperons, and the left (right) panels utilize the NJL (MIT bag) quark EOS. The parameters $`\zeta =\xi =0`$ in the Müller-Serot (MS) hadronic Lagrangian are chosen. Bold curves indicate the mixed phase region.
Figure 4: Temperature versus density in units of $`n_0`$ for two PNS evolutionary snapshots. The upper (lower) panel displays results for the NJL (MIT bag) Lagrangian. The parameters $`\zeta =\xi =0`$ in the Müller-Serot (MS) hadronic Lagrangian are chosen. Results are compared for matter containing only nucleons (np), nucleons plus hyperons (npH), nucleons plus quarks (npQ) and nucleons, hyperons and quarks (npHQ). Bold curves indicate the mixed phase region.
Figure 5: The concentrations of hadrons, quarks, and leptons as functions of density in units of $`n_0`$. Three representative snapshots during the evolution of a proto-neutron star are displayed. Matter is assumed to contain nucleons and quarks (npQ). The parameters $`\zeta =\xi =0`$ in the Müller-Serot (MS) hadronic Lagrangian are chosen. Bold curves indicate the mixed phase region.
Figure 6: The same as Figure 5, except that hyperons are included (npHQ).
Figure 7: The gravitational mass versus radius, for three representative snapshots during the PNS evolution. The left (right) panels are for the NJL (MIT bag) quark EOS, and hyperons are (are not) included in the bottom (top) panels. The parameters $`\zeta =\xi =0`$ in the Müller-Serot (MS) hadronic Lagrangian are chosen. Bold lines indicate configurations with a mixed phase at the star’s center.
Figure 8: The phase diagram of the quark-hadron transition in the baryon number density - neutrino concentration plane for three representative snapshots during the evolution of a proto-neutron star. The left (right) panels are for the NJL (MIT bag) quark EOS, and hyperons are (are not) included in the bottom (top) panels. The parameters $`\zeta =\xi =0`$ in the Müller-Serot (MS) hadronic Lagrangian are chosen. The lower- and upper-density boundaries of the mixed phase are indicated by bold curves. The central densities of maximum mass configurations are shown by thin curves.
|
warning/0003/cond-mat0003471.html
|
ar5iv
|
text
|
# Fully quantum mechanical moment of inertia of a mesoscopic ideal Bose gas
## I Introduction
Despite recent achievements in preparing Bose condensed atomic gases the question of superfluidity still escapes direct experimental observation. This is partly due to the difficulty to define and/or access the appropriate experimental observables.
Our understanding of superfluidity is formed by the physics of macroscopic systems such as liquid helium. In those systems standard theoretical methods could be applied that required a thermodynamic limit procedure.
However, for the finite size systems prepared with trapped atomic Bose gases the standard answers had to be made more precise. E.g., one can speak of a phase transition, but a discontinuous change of system observables does not occur in mesoscopic systems. The phase transition was linked to the ground state population, but the most striking effect of superfluidity cannot be observed in a direct way.
In a spatially homogeneous situation the phenomenon of superfluidity is defined as the suppression of friction for a linear motion slower than the velocity of sound . This phenomenon is constrained to a fraction of the fluid only, the so-called superfluid fraction. In linear response theory the latter is calculated via the quantum mechanical dispersion of the momentum distribution.
The suppression of friction itself can be traced back to the appearance of excitations with a linear dispersion relation for an interacting Bose gas in the presence of a condensate. It is thus inseparably connected with interparticle interactions. Nevertheless the superfluid fraction mentioned above does not vanish for an ideal Bose gas due to Bose-Einstein statistics at least for mesoscopic samples.
In the case of a atomic cloud trapped e.g. in a harmonic potential a similar qualitative change of the spectrum of low lying excitations can not be observed, and the understanding of superfluidity in rotational motion of trapped atoms is more subtle. Instead of regarding the response to a Galilei shift one has to consider the response of the gas to rotations. The response coefficient then is the moment of inertia of the trapped gas.
The approach of Brosens et. al. to the moment of inertia is based on the classical expectation value $`x^2+y^2`$ . Their analysis focuses on the difference in the moment of inertia of a totally classical Boltzmann gas in a trap and the expectation value of $`x^2+y^2`$ for a Bose gas (cf. Eq. (16)). Therefore they miss the true superfluid effects that may only be analyzed by calculating the moment of inertia from quantum mechanical response to rotations.
By contrast, Stringari’s work is based on linear response theory. He obtained the different contributions from the condensate and the thermal cloud to the response coefficient both for an ideal and an interacting Bose gas. His use of the grand canonical ensemble is certainly justified for the study of $`10^6`$ particles, a characteristic number for present BEC experiments.
Instead of considering the relation between rotations and superfluidity one might investigate the relation between dissipation and superfluidity like in where the onset of dissipation is analyzed depending on the velocity of an external perturbation. The numerical work in favors vortex creation as the main dissipation mechanism but it is still under debate at which critical velocity dissipation sets in.
The rotational properties have recently been analyzed in various works either focusing on vortices or on the so-called scissors mode . The analysis of this mode led to a connection between the quadrupole excitations and the moment of inertia of the normal fluid fraction which might open a way to measure the moment of inertia.
In this paper we present a calculation of the fully quantum mechanical moment of inertia for a mesoscopic cloud of non-interacting atoms in a cylindrically symmetrical trap. Finite size effects are allowed for by calculating the canonical ensemble averages appropriate for this regime. In this respect, our calculations are complementary to and show markedly different results for the superfluid fraction. It is of particular interest that all relevant averages are expressed using permutation cycles which have already played a crucial rule in previous Path-Integral-Monte-Carlo (PIMC) studies . Our analytical results are compared to and corroborated by numerically exact results computed by the PIMC method.
In Section II we first present the method of permutation cycles and apply it to the evaluation of the moment of inertia in Section III. In Sec. IV, we finally compare the results obtained by the different methods.
## II Canonical averages
We want to perform our calculations using the permutation cycle analysis introduced by Feynman and the canonical ensemble. Let us consider $`N`$ particles living in a exterior potential $`V(𝐫)`$ ($`𝐫R^D`$), so the single particle Hamiltonian is given as usual by
$$H=\frac{𝐩^2}{2m}+V(𝐫).$$
(1)
With the help of the eigenvalues $`E_i`$ and eigenfunctions $`|\varphi _i`$ $`H`$ can also be written as
$$H=\underset{i}{}E_i|\varphi _i\varphi _i|.$$
(2)
The total Hamiltonian for $`N`$ particles is then given by the sum $`H_N=_{j=1}^NH^{(j)}`$ over all the particles.
The central physical quantity in statistical mechanics is the partition function $`Z_N(\beta )`$ at inverse temperature $`\beta =1/(kT)`$. For a gas of $`N`$ bosons $`Z_N(\beta )`$ is given by
$$Z_N(\beta )=\frac{1}{N!}\underset{PS_n}{}\mathrm{dR}\rho (\mathrm{R},\mathrm{PR}),\mathrm{R}=(𝐫_1,\mathrm{},𝐫_\mathrm{N})$$
(3)
where $`\rho (R,PR)=R|\mathrm{e}^{\beta H_N}|PR`$ is the density matrix between the point $`R`$ and the permuted point $`PR=(𝐫_{P1},\mathrm{},𝐫_{PN})`$ and $`P`$ is a permutation of the first $`N`$ integer numbers. As the total Hamiltonian $`H_N`$ is a sum of independent single-particle Hamiltonians the integral factorizes
$$Z_N(\beta )=\frac{1}{N!}\underset{PS_n}{}\underset{j=1}{\overset{N}{}}\mathrm{d}^\mathrm{D}𝐫_\mathrm{j}\rho _1(𝐫_\mathrm{j},𝐫_{\mathrm{Pj}}).$$
(4)
Here, $`\rho _1(𝐫_,𝐫_{Pj})=𝐫_j|\mathrm{e}^{\beta H}|𝐫_{Pj}`$ is the single-particle density matrix. Now, we break up the permutations into so-called ”cycles”, that is subsets of the number from 1 to $`N`$ that are invariant under the action of a permutation $`P`$. If we break up $`P`$ in this way, we may get $`C_q`$ cycles of length $`q`$; as we are working in the canonical ensemble these numbers are restricted by $`_{q=1}^NqC_q=N`$. Rearranging the integrand of (4) one arrives at
$$Z_N(\beta )=\underset{\genfrac{}{}{0pt}{}{C_1,\mathrm{},C_N;}{_{q=1}^NqC_q=N}}{}\underset{q}{}\frac{Z_1(q\beta )^{C_q}}{C_q!q^{C_q}}$$
(5)
for the partition function (see ). Here, the sum over all combinations of ”cycle populations” $`C_1,\mathrm{},C_N`$ is restricted by $`_{q=1}^NqC_q=N`$.
By calculating the derivative with respect to $`\beta E_i`$ one gets the formula for $`N_i`$ for later evaluations (see also )
$$N_i=\frac{1}{Z_N(\beta )}\frac{Z_N(\beta )}{\beta E_i}.$$
(6)
We now apply this expression to (5) and use the fact that $`Z_1(q\beta )=_i\mathrm{e}^{q\beta E_i}`$ to finally obtain
$$N_i=\underset{q=1}{\overset{N}{}}\frac{\mathrm{e}^{q\beta E_i}}{Z_1(q\beta )}qC_q.$$
(7)
So we need to know the mean number of q-cycles $`C_q`$ to compute $`N_i`$.
Evidently, $`C_q`$ is defined by
$$C_q=\frac{1}{Z_N(\beta )}\underset{\genfrac{}{}{0pt}{}{C_1,\mathrm{},C_N;}{_{r=1}^NrC_r=N}}{}\underset{r=1}{\overset{N}{}}\frac{Z_1(r\beta )^{C_r}}{C_r!r^{C_r}}C_q.$$
(8)
To calculate this expression, we split the product into the factors with $`rq`$ and the factor with $`r=q`$. Note also that terms with $`C_q=0`$ do not contribute. So one gets
$`C_q`$ $`=`$ $`{\displaystyle \frac{1}{Z_N(\beta )}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{C_1,\mathrm{},C_N;}{_{r=1}^NrC_r=N}}{}}`$ (10)
$`{\displaystyle \underset{r=1,rq}{\overset{N}{}}}{\displaystyle \frac{Z_1(r\beta )^{C_r}}{C_r!r^{C_r}}}{\displaystyle \frac{Z_1(q\beta )^{C_q1}}{(C_q1)!q^{C_q1}}}{\displaystyle \frac{Z_1(q\beta )}{q}}.`$
As $`C_q>0`$, we can substitute $`C_q1`$ by $`C_q`$ and do the sum again from $`C_q=0`$ to $`\mathrm{}`$. But this means that we consider one $`q`$-cycle less, so $`_rrC_r=Nq`$. We end up with the final formula for the cycle occupation number
$$C_q=\frac{Z_{Nq}(\beta )}{Z_N(\beta )}\frac{Z_1(q\beta )}{q}$$
(11)
where $`Z_{Nq}(\beta )`$ originates from the sum over the products in (10). This equation together with the constraint on the $`C_q`$’s constitute the well-known recursion relations for $`Z_N(\beta )`$ .
## III Suprafluidity in a harmonic trap
In this section we want to compute the superfluid fraction $`\rho _\mathrm{s}/\rho `$ of a gas of noninteracting bosons in a harmonic trap. This will be done by using the permutation cycles introduced in the last section.
The superfluid fraction can be defined via the response of the system to infinitesimal rotations just like in the usual case for translations . The superfluid part shows no response to rotations at all while its density distribution contributes to the classical moment of inertia. Therefore, one has
$$\frac{\rho _\mathrm{s}}{\rho }=1\frac{\rho _\mathrm{n}}{\rho },$$
(12)
with the normal fluid fraction defined by the quotient of the quantum mechanical and the classical moment of inertia for rotations around the symmetry axis ($`z`$-axis)
$$\frac{\rho _\mathrm{n}}{\rho }=\frac{I_{\mathrm{qm}}}{I_{\mathrm{class}}}.$$
(13)
One can calculate $`I_{\mathrm{qm}}`$ via the response to rotations , this yields
$$I_{\mathrm{qm}}=\beta \left(L_z^2L_z^2\right)$$
(14)
or
$$I_{\mathrm{qm}}=\beta L_z^2,$$
(15)
because we only consider non-rotating situations with $`L_z=0`$. The classical moment of inertia is defined as usual by
$$I_{\mathrm{class}}=m\underset{j=1}{\overset{N}{}}(x_j^2+y_j^2).$$
(16)
We now want to compute both (15) and (16) by using the permutation cycles of section II.
We consider the single-particle Hamiltonian for a deformed, harmonic potential in three dimensions
$$H=\frac{𝐩^2}{2m}+\frac{1}{2}m\left(\omega _{}^2(x^2+y^2)+\omega _{}^2z^2\right).$$
(17)
Its eigenfunctions can be classified by three quantum numbers $`n_r=0,1,2,\mathrm{},m=0,\pm 1,\mathrm{},n_z=0,1,\mathrm{}`$ with
$`H|n_r,m,n_z`$ $`=`$ $`\{\mathrm{}\omega _{}(2n_r+|m|+1)`$ (19)
$`+\mathrm{}\omega _{}(n_z+1/2)\}|n_r,m,n_z;`$
they are also eigenfunctions of the angular momentum operator around the $`z`$-direction
$$l_z|n_r,m,n_z=m\mathrm{}|n_r,m,n_z.$$
(20)
The total angular momentum is given by the sum over the angular momentum operators for the $`N`$ particles in the trap
$$L_z=\underset{j=1}{\overset{N}{}}l_z^{(j)}.$$
(21)
We first turn to $`L_z^2`$. Instead of expressing the sum over all states in the thermodynamic averaging as integrals over the particle positions like in (3) we here use the basis of the single-particle states in (19) to calculate the expectation value of $`L_z^2`$ (we use $`i_j=(n_{r,j},m_j,n_{z,j})`$ to denote the states of particle $`j`$)
$`L_z^2`$ $`=`$ $`{\displaystyle \frac{1}{Z_NN!}}{\displaystyle \underset{PS_N}{}}{\displaystyle \underset{i_1,\mathrm{},i_N}{}}`$ (23)
$`i_1,\mathrm{},i_N|{\displaystyle \underset{j,k=1}{\overset{N}{}}}l_z^{(j)}l_z^{(k)}\mathrm{e}^{\beta H_N}|i_{P1},\mathrm{},i_{PN}.`$
Again, one can factorize the matrix element to a product of matrix elements for only one particle. The factors are either of the form $`i_j|\mathrm{e}^{\beta H^{(j)}}|i_{Pj}=\mathrm{e}^{\beta E_{i_j}}\delta _{i_j,i_{P_j}}`$ or $`i_j|l_z^{(j)}\mathrm{e}^{\beta H^{(j)}}|i_{Pj}=\mathrm{}m_j\mathrm{e}^{\beta E_{i_j}}\delta _{i_j,i_{P_j}}`$ or $`i_j|l_{z}^{(j)}{}_{}{}^{2}\mathrm{e}^{\beta H^{(j)}}|i_{Pj}=(\mathrm{}m_j)^2\mathrm{e}^{\beta E_{i_j}}\delta _{i_j,i_{P_j}}`$.
If now the sum of the permutations is expressed as a sum over all cycle occupations one can see, that due to the Kronecker-$`\delta `$s in the factors, all particles on the same $`q`$-cycle have the same state. For one $`q`$-cycle, there are again three different possibilities: if there is no $`l_z^{(j)}`$ associated to one of the particles on the cycle, then (23) gets the contribution
$$\underset{i_j}{}\mathrm{e}^{q\beta E_{i_j}}=Z_1(q\beta ).$$
(24)
Or there may be only one such $`l_z^{(j)}`$. Then the contribution is
$$\underset{i_j}{}\mathrm{}m_j\mathrm{e}^{q\beta E_{i_j}}=0,$$
(25)
which vanishes due to the symmetry $`E_{n_r,m,n_z}=E_{n_r,m,n_z}`$. The third possibility is to have two angular momentum operators acting on two particles on the same q-cycle. This contributes a factor
$$\underset{i_j}{}(\mathrm{}m_j)^2\mathrm{e}^{q\beta E_{i_j}}=\frac{2\mathrm{}^2\mathrm{e}^{q\beta \mathrm{}\omega _{}}}{\left(1\mathrm{e}^{q\beta \mathrm{}\omega _{}}\right)^2}Z_1(q\beta ).$$
(26)
So all $`q`$-cycles contribute a factor of $`Z_1(q\beta )`$, those involving two angular momentum operators additionally contribute a factor of $`\frac{2\mathrm{}^2\mathrm{e}^{q\beta \mathrm{}\omega _{}}}{\left(1\mathrm{e}^{q\beta \mathrm{}\omega _{}}\right)^2}`$.
Before we can write down the formula for $`L_z^2`$ we must count the number of ways how the two $`l_z`$-operators may be distributed among the cycles: for a given number $`C_q`$ of $`q`$-cycles there are $`qC_q`$ particles sitting on these cycles. They can be paired with $`q`$ other particles on their own cycle (including themselves), so there are $`q^2C_q`$ ways to pair the two angular momentum operators to get (26). This leads to
$`L_z^2`$ $`=`$ $`{\displaystyle \frac{1}{Z_NN!}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{C_1,\mathrm{},C_N;}{_{q=1}^NqC_q=N}}{}}M(C_r){\displaystyle \underset{q=1}{\overset{N}{}}}`$ (28)
$`{\displaystyle \frac{2\mathrm{}^2\mathrm{e}^{q\beta \mathrm{}\omega _{}}}{\left(1\mathrm{e}^{q\beta \mathrm{}\omega _{}}\right)^2}}q^2C_q{\displaystyle \underset{r=1}{\overset{N}{}}}Z_1(r\beta )^{C_r},`$
where
$$M(C_1,\mathrm{},C_N)=\frac{N!}{_qC_q!q^{C_q}}$$
(29)
is the number of permutations with $`C_1`$ 1-cycles, $`C_2`$ 2-cycles etc (see \[18, above Eq. (2.154)\]). By using this and Eqs. (8) and (15)
$$I_{\mathrm{qm}}=\beta L_z^2=2\mathrm{}^2\underset{q=1}{\overset{N}{}}\frac{q\beta \mathrm{e}^{q\beta \mathrm{}\omega _{}}}{\left(1\mathrm{e}^{q\beta \mathrm{}\omega _{}}\right)^2}qC_q.$$
(30)
So we can calculate the quantum mechanical value of the momentum of inertia by using the cycle occupations in (11).
Now we turn to the classical moment of inertia. One can do the analogous analysis as before for $`I_{\mathrm{qm}}`$. Here, one meets terms like
$$\underset{i_j}{}i_j|(x_j^2+y_j^2)\mathrm{e}^{q\beta H}|i_j=\frac{\mathrm{}}{m\omega _{}}\frac{1+\mathrm{e}^{q\beta \mathrm{}\omega _{}}}{1\mathrm{e}^{q\beta \mathrm{}\omega _{}}}Z_1(q\beta )$$
(31)
in analogy to (25) or (26). Eq. (31) is most easily computed using the eigenstates $`|n_x,n_y,n_z`$ of the harmonic trap Hamiltonian.
Finally, $`I_{\mathrm{class}}`$ can be written as
$$I_{\mathrm{class}}=\frac{\mathrm{}}{\omega _{}}\underset{q=1}{\overset{N}{}}\frac{1+\mathrm{e}^{q\beta \mathrm{}\omega _{}}}{1\mathrm{e}^{q\beta \mathrm{}\omega _{}}}qC_q,$$
(32)
which again depends on the cycle occupation numbers. Eqs. (30) and (32) constitute the main result of the present work, they allow the computation of the superfluid fraction from Eqs. (12), (13) totally based on the cycle occupation numbers in (11).
## IV Numerical results and comparison to path integral Monte-Carlo results
We now turn to the comparison of the cycle approach with results from other methods for the description of the trapped Bose gas.
In Fig. 1 we show the superfluid fraction as a function of temperature for $`N=25`$ particles in a spherical trap. The full line shows the result originating from our cycle analysis. It is obtained by using (12, 13, 30, 32). $`qC_q`$ has been calculated from the recursion relations for $`Z_N(\beta )`$.
The first comparison is with the model of Stringari who has given a formula for $`\rho _\mathrm{s}/\rho `$ based on grand canonical considerations. This result is plotted with short dashes in Fig. 1.
The difference between the canonical and grand canonical values is clearly visible both for $`N=25`$ particles and for $`N=100`$. The two chosen examples already illustrate that the difference between the two ensembles will vanish in the limit $`N\mathrm{}`$.
It is interesting that by using the canonical expectation values we can also reconcile the cycle analysis given above with a simple two-fluid model of the inhomogeneous Bose gas . In that model the superfluid fraction is totally made up of the condensed part of the system and the normal fluid is identical to the non-condensed part. We here modify the two-fluid model by inserting the number of condensed particles $`N_0`$ (see (7)) from the canonical averages. For the case of the harmonic trap, $`N_0`$ can be written as
$$N_0=\underset{q=1}{\overset{N}{}}\left(1\mathrm{e}^{q\beta \mathrm{}\omega _{}}\right)^2\left(1\mathrm{e}^{q\beta \mathrm{}\omega _{}}\right)qC_q.$$
(33)
To compute $`\rho _\mathrm{s}/\rho `$ we need the moments of inertia $`I_0`$ of the condensate and $`I_{\mathrm{nc}}`$ of the non-condensed part. The condensate particles all reside in the ground state of the trap whose moment of inertia for rotations around the $`z`$-axis is $`\mathrm{}/\omega _{}`$, so
$$I_0=N_0\frac{\mathrm{}}{\omega _{}}.$$
(34)
$`I_{\mathrm{nc}}`$ is estimated by assuming that the non-condensed particles behave like a Boltzmann gas in the harmonic trap. For the inverse temperature $`\beta `$ the moment of inertia then equals $`2/(\beta \omega _{}^2)`$ so
$$I_{\mathrm{nc}}=N_{\mathrm{nc}}\frac{2}{\beta \omega _{}^2}=N_{\mathrm{nc}}\frac{2kT}{\omega _{}^2},$$
(35)
where $`N_{\mathrm{nc}}=NN_0`$.
Finally, by noting that the condensate only contributes to the classical moment of inertia $`I_{\mathrm{class}}`$ but does not take part in the rotation, we can estimate $`\rho _\mathrm{s}/\rho `$ by
$`{\displaystyle \frac{\rho _\mathrm{s}}{\rho }}`$ $``$ $`1{\displaystyle \frac{I_{\mathrm{nc}}}{I_0+I_{\mathrm{nc}}}}={\displaystyle \frac{1}{1+\frac{NN_0}{N_0}\frac{2kT}{\mathrm{}\omega _{}}}}.`$ (36)
The long dashed lines in Fig. 1 give a plot of this formula. It fits the exact result for the canonical ensemble surprisingly well. The two-fluid model differs from the exact result only due to the small difference between the true non-condensed part and its quasi-classical approximation.
A third way to calculate the superfluid fraction via moments of inertia is the path integral representation of the density matrix . This approach works for both non-interacting and interacting systems. It represents the most important point of comparison as it is a potentially exact method. We have implemented a Path Integral Monte-Carlo (PIMC) code that relies on the bisectioning ideas of Ceperley and on the factorization of the complete density matrix into a non-interacting part and the interaction correction .
Here, we will only show that the results discussed in the previous section agree well — as they should — with the data obtained from PIMC calculations without interactions.
The data points (crosses) in Fig. 1 show our results for the superfluid fraction. They have been computed by using the so-called “area formula” which is the most appropriate method for our investigations. The superfluid fraction is obtained from
$$\frac{\rho _s}{\rho }=\frac{4m^2}{\mathrm{}^2\beta }\frac{A_z^2}{I_z}.$$
(37)
Here
$$I_z=\frac{m}{M}\underset{t=0,i=1}{\overset{M1,N}{}}\left(x_i(t)x_i(t+1)+y_i(t)y_i(t+1)\right)$$
(38)
denotes the PIMC approximation of classical moment of inertia. $`m`$ is again the mass of the particles. $`x_i(t)`$ and $`y_i(t)`$ are the coordinates of the $`i`$-th particle on time slice $`t`$ of the PIMC simulation and there are $`M`$ such time slices. $`I_z`$ clearly converges to $`I_{\mathrm{class}}`$ as $`M\mathrm{}`$.
The expression
$$A_z=\frac{1}{2}\underset{t=0,i=1}{\overset{M1,N}{}}\left(x_i(t)y_i(t+1)y_i(t)x_i(t+1)\right)$$
(39)
is the projected area perpendicular to the rotation axis $`z`$. $`A_z^2`$ is the portion of the moment of inertia that can be traced back to the superfluid fraction .
As Fig. 1 shows, all three methods are in good agreement with each other. We have furthermore calculated the density distribution of the particles in the trap both in the two-fluid model and with PIMC. They also exhibit a nice agreement thus indicating the validity of the empirical two-fluid model.
## V Summary and Conclusion
The main result of our paper is the calculation of the superfluid fraction from the permutation cycles. We have compared this approach to the grand-canonical prediction by Stringari and to PIMC calculations. For small particle numbers our results are in good agreement with the exact (canonical) PIMC results and we were able to reproduce them to a very good accuracy with a two-fluid model which divides the gas into a condensed and a non-condensed part where the latter is treated as a classical Boltzmann gas. For small particle numbers we find a distinct difference between our results and Stringari’s grand canonical approach.
The techniques and results presented in this paper have established a solid starting point of PIMC investigations including interactions. As the calculation of the condensate fraction in PIMC calculations of inhomogeneous Bose gases is still under debate, the role of the permutation cycles deserves further investigations also in the interacting case (see e.g. for a related discussion).
###### Acknowledgements.
J.S. thanks M. Holzmann for a stimulating discussion on the subject. We gratefully acknowledge financial support by DFG under Grant Nr. SCHE 128/7-1.
|
warning/0003/hep-th0003089.html
|
ar5iv
|
text
|
# The Central Term in 3D Simple Superalgebra
## 1 Self–Interacting Scalar Model and Vortex <br>Configurations
The component-field expansion for a scalar superfield reads
$$\mathrm{\Phi }(x,\theta )=A(x)+\theta ^a\psi _a(x)\theta ^2F(x);$$
(6)
where $`\theta `$ is a (real) Grassmann-valued Majorana spinor and $`A(x)`$ is a physical scalar, $`\psi _a(x)`$ is a physical fermion and $`F(x)`$ is an auxiliary field. The supersymmetry covariant derivative is represented as
$`D_a=_a+i\theta ^b_{ab},`$ (7)
$`\{D_a,D_b\}=\mathrm{\hspace{0.17em}2}P_{ab}.`$ (8)
The most general $`N`$=$`1`$-supersymmetric action with renormalizable matter self-interactions is given by
$$S_{scalar}=d^3xd^2\theta \left\{\frac{1}{2}(D_a\mathrm{\Phi })^2+\frac{1}{2}m\mathrm{\Phi }^2+\frac{\lambda }{8}\mathrm{\Phi }^4\right\},$$
(9)
and the supersymmetry transformations on the components $`A`$, $`\psi _a`$ and $`F`$ read as below:
$`\delta A`$ $`=`$ $`ϵ^a\psi _a,`$
$`\delta \psi _a`$ $`=`$ $`ϵ^b(C_{ab}F+i_{ab}A),`$
$`\delta F`$ $`=`$ $`ϵ^bi_{b}^{}{}_{}{}^{a}\psi _a.`$ (10)
The action for the physical fields,
$`S_{scalar}`$ $`=`$ $`{\displaystyle d^3x\{\begin{array}{c}\frac{1}{2}\left[\frac{1}{2}(_{ab}A)(^{ab}A)+\psi ^ai_{a}^{}{}_{}{}^{b}\psi _b\right]+m\psi ^2+\frac{3}{2}\lambda \psi ^2A^2\hfill \end{array}}+`$ (14)
$`\begin{array}{c}\hfill \frac{1}{2}m^2A^2\frac{1}{2}m\lambda A^4\frac{1}{8}\lambda ^2A^6\end{array}\},`$
can be shown to be invariant under the non-linear “on-shell” transformations,
$`\delta A`$ $`=`$ $`ϵ^a\psi _a,`$
$`\delta \psi _a`$ $`=`$ $`ϵ^b\left[C_{ab}\left({\displaystyle \frac{1}{2}}mA{\displaystyle \frac{1}{2}}\lambda A^3\right)+i_{ab}A\right].`$ (15)
Now , taking into account that supersymmetry is a symmetry of the action (the Lagrangian density transforms as a total derivative), it can be shown that the Noether current associated to $`N`$=$`1`$-supersymmetry turns out to be:
$`J_{}^{\mu }{}_{c}{}^{}`$ $`=`$ $`i\psi ^a(\gamma ^\mu )_{ac}\left(mA+{\displaystyle \frac{\lambda }{2}}A^3\right){\displaystyle \frac{i}{2}}\epsilon ^{\mu \nu \rho }\psi ^b(\gamma _\rho )_{bc}_\nu A+`$ (16)
$`+{\displaystyle \frac{1}{2}}\psi _c^\mu A+{\displaystyle \frac{1}{2}}A^\mu \psi _c{\displaystyle \frac{i}{2}}\epsilon ^{\mu \nu \rho }A_\nu \psi ^a(\gamma _\rho )_{ac}.`$
The supercharge is defined as
$`Q_c`$ $`=`$ $`{\displaystyle d^2\stackrel{}{x}J_{}^{0}{}_{c}{}^{}}`$ (17)
$`=`$ $`{\displaystyle }d^2\stackrel{}{x}\{i\psi ^a(\gamma ^0)_{ac}(mA+{\displaystyle \frac{\lambda }{2}}A^3)+{\displaystyle \frac{1}{2}}\psi _c^0A+{\displaystyle \frac{1}{2}}A^0\psi _c+`$
$`{\displaystyle \frac{i}{2}}\epsilon ^{0\nu \rho }\psi ^a(\gamma _\rho )_{ac}_\nu A+{\displaystyle \frac{i}{2}}\epsilon ^{0\nu \rho }A_\nu \psi ^a(\gamma _\rho )_{ac}\}.`$
With the help of the canonical commutation (and anticommutation) relations for the physical fields, a tedious calculation yields the following expression for the algebra of supersymmetry charges:
$`\{Q_a,Q_b\}={\displaystyle d^2\stackrel{}{x}\times }`$
$`2i\{{\displaystyle \frac{1}{4}}[\begin{array}{c}2i\psi ^a(\gamma ^0)_{a}^{}{}_{}{}^{b}^0\psi _b+(^0A)(^0A)+\mathrm{\hspace{0.17em}2}i\psi ^a(\gamma ^i)_{a}^{}{}_{}{}^{b}^i\psi _b+(^iA)(^iA)+\hfill \end{array}`$ (19)
$`\begin{array}{c}\begin{array}{c}\psi ^2\left(m+\frac{3}{2}\lambda A^2\right)+\left(\frac{1}{2}m^2A^2+\frac{1}{2}m\lambda A^4+\frac{1}{4}\lambda ^2A^6\right)\hfill \end{array}]\hfill \end{array}\}(\gamma ^0)_{ab}+`$ (22)
$`\mathrm{\hspace{0.17em}2}i\left\{{\displaystyle \frac{1}{4}}\begin{array}{c}2i\psi ^a(\gamma ^{(0})_{a}^{}{}_{}{}^{b}^{i)}\psi _b+(^{(0}A)(^{i)}A)\hfill \end{array}\right\}(\gamma _i)_{ab}.`$ (24)
If we compare this expression with the $`0\mu `$ component of the “improved” energy-momentum tensor
$$T_{\mu \nu }\frac{1}{e}\frac{\delta S}{\delta e_a^{}^{(\mu }}e_{\nu )a^{}}=\frac{2}{\sqrt{g}}\frac{\delta S}{\delta g^{\mu \nu }}|_{g^{\mu \nu }=\eta ^{\mu \nu }},$$
(25)
where $`S`$ in this expression indicates the action (14), we may rewrite (24) as
$$\{Q_a,Q_b\}=\mathrm{\hspace{0.17em}2}iP^\mu (\gamma _\mu )_{ab}+\mathrm{\hspace{0.17em}2}iϵ^{ij}d^2\stackrel{}{x}(_iA)(\gamma _j)_{ab}\left(mA+\frac{\lambda }{2}A^3\right).$$
(26)
In terms of the chiral components of the supersymmetry charge, the algebra takes over the form:
$`\{Q^+,Q^+\}`$ $`=`$ $`2i(P^0+P^1)\mathrm{\hspace{0.17em}2}{\displaystyle d^2\stackrel{}{x}\left(mA+\frac{\lambda }{2}A^3\right)_2A},`$ (27)
$`\{Q^{},Q^{}\}`$ $`=`$ $`2i(P^0P^1)\mathrm{\hspace{0.17em}2}{\displaystyle d^2\stackrel{}{x}\left(mA+\frac{\lambda }{2}A^3\right)_2A},`$
$`\{Q^+,Q^{}\}`$ $`=`$ $`2iP^2\mathrm{\hspace{0.17em}2}{\displaystyle d^2x\left(mA+\frac{\lambda }{2}A^3\right)_1A}.`$ (28)
Expressions (26) and (28) sign the presence of a central charge that is sensitive to a topologically non-trivial behavior of the scalar sector at infinity:
$`T_2`$ $`=`$ $`{\displaystyle 𝑑x_1𝑑x_2\frac{}{_{x^2}}\left(mA^2+\frac{\lambda }{4}A^4\right)},`$
$`T_1`$ $`=`$ $`{\displaystyle 𝑑x_2𝑑x_1\frac{}{_{x^1}}\left(mA^2+\frac{\lambda }{4}A^4\right)},`$ (29)
where we observe the topological character of the central charge terms, which has its origin in the mass and self-interacting terms of the scalar field in Lagrangian. Bearing in mind this result, we shall now consider the introduction of an Abelian gauge field with $`N`$=$`1`$–supersymmetry. We also know that such a coupling is fundamental to stabilize the soliton-like solutions in the form of magnetic vortices with finite energy.
## 2 On the N=1 Super–QED<sub>3</sub>
The $`N`$=$`1`$–supersymmetric version of QED<sub>3</sub> is achieved upon the complexification of the scalar superfield in eq. (9) and the gauge–covariantization of the spinor derivative:
$$_aD_ai\mathrm{\Gamma }_a,$$
(30)
where $`\mathrm{\Gamma }_a`$ is a gauge superconnection with super–helicity $`h=\frac{1}{2}`$, and the signs $``$ and $`+`$ indicate that the derivative is acting in the superfields $`\mathrm{\Phi }`$ and $`\overline{\mathrm{\Phi }}`$ respectively ($`\overline{\mathrm{\Phi }}\mathrm{\Phi }^{}`$). $`\mathrm{\Gamma }_a`$ admits the following $`\theta `$–expansion:
$$\mathrm{\Gamma }_a=\chi _a+\theta ^b(C_{ab}B+iV_{ab})+\theta ^2(2\lambda _ai_{a}^{}{}_{}{}^{b}\chi _b),$$
(31)
where $`\lambda _a`$ is the gaugino field, $`V_{ab}`$ is the usual gauge field; $`B`$ and $`\chi _a`$ are compensating component–fields. In the so–called Wess–Zumino gauge , $`\mathrm{\Gamma }_a`$ is reduced to:
$$\mathrm{\Gamma }_a=i\theta ^bV_{ab}2\theta ^2\lambda _a,$$
(32)
where the supersymmetry transformations read:
$`\delta V_{ab}=iϵ_{(b}\lambda _{b)},`$
$`\delta \lambda _a={\displaystyle \frac{1}{2}}ϵ^c_{c(a}V_{c)b}.`$ (33)
The covariantized vectorial derivative is written as
$$_{ab}=D_{ab}i\mathrm{\Gamma }_{ab},$$
(34)
where $`\mathrm{\Gamma }_{ab}`$ is the vector superconnection. As we know, in order to have irreducible representations of symmetry, we need constraints in the model. In the supersymmetric case, we have the so–called conventional constraint, that acts in such away that the supersymmetric algebra of the spinor derivatives, $`\{_a,_b\}=\mathrm{\hspace{0.17em}2}i_{ab}+F_{ab}`$, will have $`F_{ab}=\mathrm{\hspace{0.17em}0}`$. Then, we easily compute that
$$\mathrm{\Gamma }_{ab}=\frac{i}{2}D_{(a}\mathrm{\Gamma }_{b)},$$
(35)
implying that in the Wess–Zumino gauge we have
$$\mathrm{\Gamma }_{ab}=V_{ab}+i\theta _{(a}\lambda _{b)}\frac{i}{2}\theta ^2_{c(a}V_{b)}^{}{}_{}{}^{c}.$$
(36)
By the graded Bianchi identity, we redefine the gauge field as
$$W_a=\frac{1}{2}D^bD_a\mathrm{\Gamma }_b,$$
(37)
with the constraint $`D^aW_a=\mathrm{\hspace{0.17em}0}`$ ($`D^aD_bD_a=\mathrm{\hspace{0.17em}0}`$), implying as in the usual Lorentz gauge that exists only one independent component of the field $`W_a`$. Using the projector method, we write
$$W_a|=\lambda _a,D_aW_b|=\frac{1}{2}(_{ca}V_{b}^{}{}_{}{}^{c}+_{cb}V_{a}^{}{}_{}{}^{c})f_{ab},$$
(38)
with $`f_{ab}`$ the usual gauge field strength. Another relations that will be very important and that may straightforwardly be obtained are (conf. ()):
$$_a^2=i_{a}^{}{}_{}{}^{b}_b\pm iW_a\text{e}(^2)^2=\mathrm{}iW^a_a,$$
(39)
where $`\mathrm{}`$ means the gauge–covariant d’Alembertian. Now, we are ready to discuss the supersymmetry algebra in the framework of $`N`$=$`1`$ Super-QED<sub>3</sub>.
### 2.1 Scalar Superaction with a Background Gauge Field
The U(1)–invariant superfield action without kinetic term for the gauge sector is given as below:
$$S_{scalar}=\frac{1}{2}d^3xd^2\theta \left\{(^a\overline{\mathrm{\Phi }})(_a\mathrm{\Phi })\right\},$$
(40)
Redefining the component field, (we can always do this) by the projections
$`\mathrm{\Phi }|=A`$ , $`\overline{\mathrm{\Phi }}|=\overline{A},`$
$`_a\mathrm{\Phi }|=\psi _a`$ , $`_a\overline{\mathrm{\Phi }}|=\overline{\psi }_a,`$
$`^2\mathrm{\Phi }|=F`$ , $`\overline{\mathrm{\Phi }}|=\overline{F},`$ (41)
the gauge–field component action takes the form:
$$S=\frac{1}{2}d^3x\left\{\overline{\psi }^aiD_{a}^{}{}_{}{}^{b}\psi _b+\psi ^aiD_{a}^{}{}_{}{}^{b}\overline{\psi }_b+A\mathrm{}\overline{A}+\overline{A}\mathrm{}A\right\},$$
(42)
with the “on shell” supersymmetric transformations:
$`\delta \overline{\psi }^a=ϵ_biD^{ab}\overline{A}`$ , $`\delta \psi _b=ϵ^ciD_{bc}A,`$
$`\delta A=ϵ^a\psi _a`$ , $`\delta \overline{A}=ϵ^a\overline{\psi }_a,`$ (43)
The Noether current associated to $`N`$=$`1`$–supersymmetry is now:
$`J_{}^{\mu }{}_{c}{}^{}`$ $`=`$ $`{\displaystyle \frac{i}{2}}\epsilon ^{\mu \nu \rho }[\overline{\psi }^a(\gamma _\rho )_{ac}_\nu A+\psi ^a(\gamma _\rho )_{ac}_\nu \overline{A}]{\displaystyle \frac{1}{2}}(\psi _c^\mu \overline{A}+\overline{\psi }_c^\mu A)+`$ (44)
$``$ $`{\displaystyle \frac{1}{2}}(\overline{A}^\mu \psi _c+A^\mu \overline{\psi }_c)+{\displaystyle \frac{i}{2}}\epsilon ^{\mu \nu \rho }(\overline{A}_\nu \psi ^a+A_\nu \overline{\psi }^a)(\gamma _\rho )_{ac},`$
yielding the supercharge
$`Q_c`$ $`=`$ $`{\displaystyle d^2\stackrel{}{x}J_{}^{0}{}_{c}{}^{}}=`$ (45)
$`=`$ $`{\displaystyle }d^2\stackrel{}{x}{\displaystyle \frac{1}{2}}\{(\psi _c^0\overline{A}+\overline{\psi }_c^0A)+i\epsilon ^{0ij}[\overline{\psi }^a(\gamma _j)_{ac}_iA+\psi ^a(\gamma _j)_{ac}_i\overline{A}]+`$
$`(\overline{A}^0\psi _c+A^0\overline{\psi }_c)+i\epsilon ^{0ij}(\overline{A}_i\psi ^a+A_i\overline{\psi }^a)(\gamma _j)_{ac}\}`$
The canonical conjugate momenta that will be necessary for the supercharge algebra are
$`\mathrm{\Pi }_{\psi _d}=i\overline{\psi }^a(\gamma ^0)_{a}^{}{}_{}{}^{d}`$ , $`\mathrm{\Pi }_{\overline{\psi }_d}=i\psi ^a(\gamma ^0)_{a}^{}{}_{}{}^{d},`$
$`\mathrm{\Pi }_A=(D^0\overline{A})`$ , $`\mathrm{\Pi }_{\overline{A}}=(D^0A),`$ (46)
giving the canonical commutation and anticommutation relations
$`\{\psi ^d,\overline{\psi }^a\}=i(\gamma ^0)^{ad}\delta ^2(\stackrel{}{x}\stackrel{}{y})`$ , $`\{\overline{\psi }^d,\psi ^a\}=i(\gamma ^0)^{ad}\delta ^2(\stackrel{}{x}\stackrel{}{y}),`$
$`[A,D^0\overline{A}]=\delta ^2(\stackrel{}{x}\stackrel{}{y})`$ , $`[\overline{A},D^0A]=\delta ^2(\stackrel{}{x}\stackrel{}{y}).`$ (47)
After a lengthy computation, using the $`\gamma `$–matrices Clifford algebra, we reach the result
$$\{Q_a,Q_b\}=2iP^\mu (\gamma _\mu )_{ab},$$
(48)
where the momentum operator $`P^\mu `$ appearing in the RHS includes now contributions from the gauge field minimally coupled to matter through (42). Nevertheless, no term in the form of a central charge arises from the action (42); this means that the central charge operator of eq. (26) is not modified by the introduction of the $`U(1)`$ gauge superfield. The role of the latter is to stabilize the topological configurations associated to the action (9) in the form of vortex-like solitons, as already known from the works quoted in ref. .
## 3 The Supersymmetric Chern–Simons Term
Now, we shall add a supersymmetric Chern–Simons (CS) term to the action eq. (42)and we will verify how it modifies the supercharge algebra. For this purpose, we begin with the (gauge–invariant) CS term in superspace
$$S_{CS}=\frac{M}{g^2}d^3xd^2\theta \mathrm{\Gamma }^aW_a,$$
(49)
where $`M`$ is a mass parameter and $`g`$ is the gauge coupling constant. In components, using the Wess–Zumino gauge, the action eq. (49) leads to the expression
$$S_{CS}=\frac{M}{g^2}d^3x\left[iV^{ab}(_{ac}V_{}^{c}{}_{b}{}^{})+\mathrm{\hspace{0.17em}4}\lambda ^2\right],$$
(50)
where the first term in the r.h.s. is the well–known CS term. Now, including the term (49) in the action (42), and then taking into account the equations of motion for the $`F`$ , $`\overline{F}`$ (which are not affected by the CS term) and of the $`\lambda _a`$–field, the complete Lagrangian reads:
$`L_{din}`$ $`=`$ $`{\displaystyle \frac{i}{2}}\overline{\psi }^a(\gamma ^\mu )_{a}^{}{}_{}{}^{b}_\mu \psi _b+{\displaystyle \frac{i}{2}}\psi ^a(\gamma ^\mu )_{a}^{}{}_{}{}^{b}_\mu \overline{\psi }_b{\displaystyle \frac{1}{2}}(^\mu \overline{A})(_\mu A)+{\displaystyle \frac{i}{2}}(^\mu \overline{A})V_\mu A+`$ (51)
$`{\displaystyle \frac{i}{2}}V^\mu \overline{A}(_\mu A)+{\displaystyle \frac{iM}{g^2}}ϵ^{\mu \nu \rho }V_\mu _\nu V_\rho ,`$
with the “on shell” supersymmetric transformations:
$`\delta \overline{\psi }^a=iϵ_bD^{ab}\overline{A}`$ , $`\delta \psi _b=iϵ^cD_{bc}A,`$
$`\delta A=ϵ^a\psi _a`$ , $`\delta \overline{A}=ϵ^a\overline{\psi }_a,`$
$`\delta V_{ab}`$ $`=`$ $`iϵ_{(a}\lambda _{b)}.`$ (52)
Following again the procedure to read off the Noether current associated to the transformations (52), it can be found out that the contribution of the CS term yields:
$$(J_{SCS})_c^\mu =\frac{i}{2}V^\mu \left(\psi _c\overline{A}\overline{\psi }_cA\right),$$
(53)
with the supercharge
$$(Q_{SCS})_c=d^2\stackrel{}{x}(J_{SCS})_c^0=d^2\stackrel{}{x}\left\{\frac{i}{2}V^0\left(\psi _c\overline{A}\overline{\psi }_cA\right)\right\},$$
(54)
whence
$$\{Q_a,Q_b\}_{SCS}=\mathrm{\hspace{0.17em}2}iP^\mu [V^0](\gamma _\mu )_{ab}.$$
(55)
where $`P^\mu [V^0]`$ means a functional that depends only on the time–component of the potential vector.
What we observe is that the $`V^0`$ potential field is completely eliminated from the algebra, implying that the “corrected” Chern-Simons $`T^{0\mu }`$ component of energy-momentum tensor, defined as “new” $`P^\mu `$ becomes independent on the potential gauge field. It is possible to say that the the conjugate momenta of the $`A`$ and $`\overline{A}`$ fields are in fact “corrected” by the CS term to become $`\mathrm{\Pi }_A\alpha ^0\overline{A}`$ and $`\mathrm{\Pi }_{\overline{A}}\alpha ^0A`$ . This indicates that the CS term plays a role similar to a partial gauge fixing, eliminating one degree of freedom of the gauge field, referring to the algebra.
## 4 Conclusions
The basic motivation of this paper was to analyze the 3-dimensional counterpart of a well-known result by Olive and Witten , namely, the appearance of a central charge in the algebra of simple supersymmetry as originated from non-trivial topological field configurations. To this aim we have analyzed the full supersymmetric model in 3D. We have computed the Noether supersymmetric charges, and the “improved” energy-momentum tensor using the gravitation approach. From the canonical commutations (and anti-commutations) relations of the superfields we obtained the supercharge algebra. Here, with and without gauge fields, we could conclude that vortex-like field configurations are responsible for a central charge in the supersymmetry algebra, even in the case of a $`N=1`$-supersymmetry. It is worthwhile to mention the results obtained by Lee, Lee and Weinberg , where a central charge comes out in context of an $`N=2`$ extended supersymmetric model. We would like to point out that the calculations of Section 3 recall that the Chern-Simons term for the gauge field does not give contribution to the central charge appearing in the algebra. In fact it could represent a sample of gauge fixing eliminating the time direction of the vector potencial in the algebra. So the central charge of this model arises exclusively from the matter sector and its existence to the vortex configurations of the scalar fields. Clearly, the role of the gauge fields is to render finite the vortex energy . The analysis of the BPS bounds and its consequences will appear in a forthcoming paper.
Next, it might be of relevance to analyze the presence of central charges in the models proposed by Dorey and Mavromatos to study $`P,T`$ conserving superconducting gauge models whenever the latter are supersymmetrised. One could perhaps understand whether or not central charges may be related to some physical aspects of superconductivity.
Acknowledgements: The author is grateful to Dr. J. A. Helayël-Neto for all the discussions and the critical comments on the original manuscript; thanks are also due to Dr. O. Piguet for very pertinent and helpful remarks, and to the colleagues at DCP-CBPF for the general discussions. CAPES-Brasil is also acknowledged for the invaluable financial help.
|
warning/0003/math-ph0003021.html
|
ar5iv
|
text
|
# 1 Introduction
## 1 Introduction
The study of lattice Bose gas models was proposed in order to understand the effect of interaction on Bose-Einstein Condensation (BEC). It was argumented that, if the interaction has a hard core, i.e. each site can be occupied by at most one particle, then the model can be mapped onto a quantum spin-$`\frac{1}{2}`$ XY-model on a $`Z^d`$-lattice.
The exact solution, exhibiting condensation, for the hard core Bose gas on a complete graph has been obtained . In spin language this work corresponds to the mean field spin-$`\frac{1}{2}`$ XY-model. Other exact results on the full model are an upper bound on the condensate and an existence proof of condensation for the system in a half filled $`Z^d`$-lattice for dimensions $`d3`$ . In this context we mention also the relevant result about the occurrence of long-range order in the XY-model with spins and dimensions greater than one.
In this paper we turn our attention to a system of two types of interacting hard core Bose particles on a complete graph and interacting through a repulsive interaction if the particles are of different type. (Compare this with the Hubbard model for electron systems.) The model is exactly soluble. We study in detail the equilibrium states in the case of a half filled lattice and equal density for the two particle types. If the system is weakly interacting we obtain an explicit and complete analysis of the equilibrium states. We analyse the BEC phenomenon and in particular we are interested in the appearance of spontaneous symmetry breaking (SSB) linked with it. The interesting aspect of this model is that the broken symmetry is two-dimensional. While, as is well known for a single type Bose system, the broken symmetry is the one-dimensional gauge symmetry. We compute rigorously the plasmon frequencies related to these modes, and find the quantum normal modes of these plasmon motions. We find out that there are two different frequencies, but that both frequencies are two-fold degenerate. The four quantum degrees of freedom have normal modes which we build up as the quantum fluctuation operators of the components of the infinitesimal generators of the spontaneously broken gauge symmetries, with adjoints the quantum fluctuation operators of corresponding order parameter operators. The degeneracy of the plasmon spectrum is the new phenomenon in this model study of Bose-Einstein condensation. Another revealing fact is that the adjoint fluctuation operators are the total and the relative current fluctuations between the two types of bosons.
## 2 The Model and its Equilibrium States
### 2.1 The Model
The model is defined on the lattice $`Z`$, by the local Hamiltonians $`H_\mathrm{\Lambda }`$, with $`\mathrm{\Lambda }`$ a finite interval in $`Z`$:
$$H_\mathrm{\Lambda }=\frac{t}{|\mathrm{\Lambda }|}\underset{\alpha \{1,2\}}{}\underset{ij\mathrm{\Lambda }}{}\sigma _{i\alpha }^+\sigma _{j\alpha }^{}\mu \underset{\alpha \{1,2\}}{}\underset{i\mathrm{\Lambda }}{}\sigma _{i\alpha }^+\sigma _{i\alpha }^{}+U\underset{i\mathrm{\Lambda }}{}\sigma _{i1}^+\sigma _{i1}^{}\sigma _{i2}^+\sigma _{i2}^{}.$$
(1)
The first term is the hopping contribution, the particles jump freely from one site to another with a hopping rate $`t/|\mathrm{\Lambda }|`$, $`|\mathrm{\Lambda }|`$ is the number of points in $`\mathrm{\Lambda }`$; the index $`\alpha `$ specifies the type of particles; $`\mu `$ is the chemical potential; $`U`$ is taken to be positive and is the coupling constant for the interaction between the types of particles. If two particles of different type are on the same lattice point the potential energy is raised with an amount $`U`$. Clearly the $`\sigma ^\pm ,\sigma ^z`$ are the usual Pauli matrices and are here the creation an annihilation operators for the Bose particles with hard core, satisfying:
$`[\sigma _{i\alpha }^+,\sigma _{j\beta }^{}]=\sigma _{i\alpha }^z\delta _{ij}\delta _{\alpha \beta };`$
$`\left(\sigma _{i\alpha }^+\right)^2=0;\left(\sigma _{i\alpha }^+\right)^{}=\sigma _{i\alpha }^{};`$
$`\sigma _{i\alpha }^+\sigma _{i\alpha }^{}={\displaystyle \frac{1}{2}}\left(\sigma _{i\alpha }^z+\mathrm{𝟏}\right).`$
We look first for the equilibrium states of this model in the thermodynamic limit. The Hamiltonians (1) are permutation invariant with respect to the lattice index. This means that the equilibrium states at inverse temperatures $`\beta `$ are convex combinations of product states on the infinite tensor product algebra of observables $`=_{iZ}\left(M_2M_2\right)_i`$ of the system; $`\left(M_2M_2\right)_i`$ represents the tensor product of the $`2\times 2`$ complex matrices $`M_2`$ with itself. Each set $`M_2`$ is referring to a different type of hard core bosons, both at site $`iZ`$. In particular, if one looks only at the product states which are invariant for the interchange of the two types of particles, then the equilibrium states are given by the states $`\omega _{\beta \lambda }(.)`$ of the form :
$$\omega _{\beta \lambda }(A_1A_2\mathrm{}A_n)=\underset{i=1}{\overset{n}{}}\mathrm{tr}\rho _{\beta \lambda }A_i;\text{for any }A_i\left(M_2M_2\right)_i,i=1,\mathrm{}n$$
where $`\rho _{\beta \lambda }`$ is a density matrix in $`M_2M_2`$, satisfying the self-consistency equation:
$$\rho _{\beta \lambda }=\frac{\mathrm{exp}[\beta h_\lambda ]}{\mathrm{tr}\mathrm{exp}[\beta h_\lambda ]},$$
(2)
where
$$h_\lambda =t\lambda \left(\sigma _1^++\sigma _2^+\right)+t\overline{\lambda }\left(\sigma _1^{}+\sigma _2^{}\right)\mu \left(\sigma _1^+\sigma _1^{}+\sigma _2^+\sigma _2^{}\right)+U\sigma _1^+\sigma _1^{}\sigma _2^+\sigma _2^{}$$
(3)
and
$$\lambda =\mathrm{tr}\rho _{\beta \lambda }\sigma _1^{}=\mathrm{tr}\rho _{\beta \lambda }\sigma _2^{}.$$
(4)
In fact equation (2) for $`\rho _{\beta \lambda }`$ is equivalent to the equation (4) for $`\lambda `$. Remark that (4) also results from the assumption of symmetry in the particle types. Looking for equilibrium states breaking this symmetry would violate (4).
Moreover, we will limit here our study to the case of a half filled lattice. Mathematically this means that we require:
$$\omega _{\beta \lambda }(\sigma _{i\alpha }^z)=0\text{for all }iZ,\alpha =1,2.$$
This fixes the chemical potential $`\mu =U/2`$. In order to simplify the notations, we adapt our Hamiltonian by adding a constant term and by scaling the constant $`U`$ to $`U/4`$. In this case one gets for (3) a more simple expression:
$$h_\lambda =t\lambda \left(\sigma _1^++\sigma _2^+\right)+t\overline{\lambda }\left(\sigma _1^{}+\sigma _2^{}\right)+U\sigma _1^z\sigma _2^z.$$
(5)
On a finite subset $`\mathrm{\Lambda }Z`$ we define Fourier-transformed operators:
$$a_{k\alpha }=\frac{1}{\sqrt{|\mathrm{\Lambda }|}}\underset{j\mathrm{\Lambda }}{}\sigma _j^{}\mathrm{exp}\left[\mathrm{i}\frac{2\pi }{|\mathrm{\Lambda }|}kj\right];k\mathrm{\Lambda }^{}\text{the dual of }\mathrm{\Lambda },$$
the annihilation operator of the $`\alpha `$-particle with quasi-momentum $`k`$.
Then,
$$n_k=\underset{\alpha =1}{\overset{2}{}}a_{k\alpha }^{}a_{k\alpha }$$
is the number operator of particles with quasi-momentum $`k`$.
Compute in a state $`\omega _{\beta \lambda }`$ (2):
$$\rho _0=\underset{|\mathrm{\Lambda }|\mathrm{}}{lim}\omega _{\beta \lambda }\left(\frac{n_0}{|\mathrm{\Lambda }|}\right)=2\omega _{\beta \lambda }(\sigma ^+)\omega _{\beta \lambda }(\sigma ^{})=2|\lambda |^2;$$
(6)
i.e. the $`0`$-mode has a density $`\rho _0`$ equal to $`2|\lambda |^2`$. A solution $`\lambda 0`$ of (4) yields a macroscopic occupation of the $`0`$-mode and the density of the condensate $`\rho _0`$ equals $`2|\lambda |^2`$. Therefore, if one finds solutions of (2) or (4) such that $`\lambda 0`$, one proved Bose-Einstein condensation.
### 2.2 Equilibrium States
Now we proceed to look for the solution of (4) and its properties. It is clear that based on general results we reduced our problem to a one-site problem, namely looking for the density matrix state $`\mathrm{tr}\rho _{\beta \lambda }.`$ with $`\rho _{\beta \lambda }M_2M_2`$.
In general $`M_2M_2=M_4`$ and we use the following isomorphism:
$`\left(\begin{array}{cc}a_{11}& a_{12}\\ a_{21}& a_{22}\end{array}\right)\left(\begin{array}{cc}b_{11}& b_{12}\\ b_{21}& b_{22}\end{array}\right)=\left(\begin{array}{cccc}a_{11}b_{11}& a_{12}b_{11}& a_{11}b_{12}& a_{12}b_{12}\\ a_{21}b_{11}& a_{22}b_{11}& a_{21}b_{12}& a_{22}b_{12}\\ a_{11}b_{21}& a_{12}b_{21}& a_{11}b_{22}& a_{12}b_{22}\\ a_{21}b_{21}& a_{22}b_{21}& a_{21}b_{22}& a_{22}b_{22}\end{array}\right).`$
This yields the following matrix representation for $`h_\lambda `$:
$$h_\lambda =\left(\begin{array}{cccc}U& t\lambda & t\lambda & 0\\ t\overline{\lambda }& U& 0& t\lambda \\ t\overline{\lambda }& 0& U& t\lambda \\ 0& t\overline{\lambda }& t\overline{\lambda }& U\end{array}\right).$$
(8)
Then a straightforward diagonalisation of $`h_\lambda `$ yields the following spectrum:
$$ϵ_0=U,ϵ_1=U,ϵ_2=\eta ,ϵ_3=\eta ,$$
(9)
where $`\eta =\sqrt{U^2+4|t\lambda |^2}`$.
The corresponding eigenvectors are:
$$\begin{array}{cc}\hfill |\varphi _0& =\frac{1}{\sqrt{2}}\left(\begin{array}{c}0\\ 1\\ 1\\ 0\end{array}\right),\hfill \\ \hfill |\varphi _2& =\frac{1}{\sqrt{\eta (\eta U)}}\left(\begin{array}{c}t\lambda \\ \frac{\eta U}{2}\\ \frac{\eta U}{2}\\ t\overline{\lambda }\end{array}\right),\hfill \end{array}\begin{array}{cc}\hfill |\varphi _1& =\frac{1}{\sqrt{2}|t\lambda |}\left(\begin{array}{c}t\lambda \\ 0\\ 0\\ t\overline{\lambda }\end{array}\right),\hfill \\ \hfill |\varphi _3& =\frac{1}{\sqrt{\eta (\eta +U)}}\left(\begin{array}{c}t\lambda \\ \frac{\eta +U}{2}\\ \frac{\eta +U}{2}\\ t\overline{\lambda }\end{array}\right).\hfill \end{array}$$
(10)
Using (9) and (10) one gets for equation (4) the following explicit and nonlinear equation for $`\lambda `$:
$$\lambda \left[1+\frac{t}{\eta }\frac{\mathrm{sinh}(\beta \eta )}{\mathrm{cosh}(\beta U)+\mathrm{cosh}(\beta \eta )}\right]=0.$$
(11)
Remark that $`\lambda =0`$ is always a solution. It corresponds to a state without condensation (see (6)).
Looking for solutions $`\lambda 0`$ amounts to look for a solution $`\eta `$ satisfying:
$$\eta =f_\beta (\eta );$$
(12)
where
$$f_\beta (\eta )=t\frac{\mathrm{sinh}(\beta \eta )}{\mathrm{cosh}(\beta U)+\mathrm{cosh}(\beta \eta )}.$$
Since the condensate density $`\rho _0=2|\lambda |^2`$ (cfr. (6)) cannot exceed the total particle density $`\rho `$ (which is equal to one in the half filled lattice), $`\sqrt{U^2+2t^2}`$ is an upper bound on $`\eta =\sqrt{U^2+4|t\lambda |^2}`$. The minimum value of $`\eta `$ is $`U`$. This defines $`I`$, the range of accessible values for $`\eta `$:
$$\eta I=[U,\sqrt{U^2+2t^2}].$$
Remark that $`f_\beta (\eta )`$ is unaffected by the sign of $`U`$. On the other hand, equation (12) has no solutions if $`t0`$. Therefore we take $`t<0`$ and $`U>0`$.
Discussion of the solutions of (11)
We are looking for solutions yielding a second order phase transition.
Remark first that for $`\eta I`$:
$`{\displaystyle \frac{f_\beta (\eta )}{\eta }}`$ $`>0`$ (13)
$`{\displaystyle \frac{^2f_\beta (\eta )}{\eta ^2}}`$ $`<0`$ (14)
$`{\displaystyle \frac{f_\beta (\eta )}{\beta }}`$ $`>0`$ (15)
i.e. $`\eta f_\beta (\eta )`$ is a monotonically increasing concave function on $`I`$; also $`\beta f_\beta (\eta )`$ is monotonically increasing for $`T`$ decreasing.
The problem here is that, due to the presence of the potential $`U`$, these properties are not valid outside the interval $`I`$, such that the situation is more complicated than for the usual mean field case.
First we consider temperature $`T>0`$, and we distinguish two cases:
case a):
$$f_\beta (U)>U$$
(16)
yielding
$$\frac{1}{2U}\mathrm{ln}\left(\frac{t+2U}{t2U}\right)<\beta ;$$
(17)
and implying the condition on the interparticle interaction:
$$2U<t$$
(18)
i.e. the interparticle interaction must be relatively small compared with the hopping term (1); this condition is a meaningful situation for this kind of systems, if the interaction becomes to strong, the typical quantum effects like BEC disappear.
One derives also
$$f_\beta \left(\sqrt{U^2+2t^2}\right)\sqrt{U^2+2t^2}.$$
(19)
Together, condition (16) and (19) yield a unique solution of $`f_\beta (\eta )=\eta I`$, and because of (15) a second order phase transition with $`\lambda 0`$ as an order parameter. Equation (17) learns that there exists a critical $`\beta _c`$ determinated by:
$$\frac{1}{k_BT_c}=\beta _c=\frac{1}{2U}\mathrm{ln}\left(\frac{t+2U}{t2U}\right).$$
(20)
This defines the critical temperature $`T_c`$.
case b):
$$f_\beta (U)<U$$
(21)
This case is somewhat more complicated, a careful analysis learns that apart from a second order phase transition, a first order one can occur, even if (18) is satisfied. In order to avoid non-natural first order transitions, one has to impose:
$$\frac{f_{\beta _c}(\eta )}{\eta }|_U1;f_{\beta _c}(U)=U.$$
The analysis of these yields a stronger condition on $`U`$ than (18):
$$U<(t)\kappa ,$$
where $`\kappa `$ can be computed numerically as the solution of the equation:
$$\begin{array}{cc}& \mathrm{ln}\left(\frac{1+2\kappa }{12\kappa }\right)=\frac{4\kappa }{12\kappa ^2};\hfill \\ & \kappa 0.461292.\hfill \end{array}$$
Finally we consider the ground state $`(T=0)`$ situation. The self-consistency equation (11) becomes:
$$\eta =\underset{\beta \mathrm{}}{lim}f_\beta (\eta )=f_{\mathrm{}}(\eta );$$
with:
$$f_{\mathrm{}}(\eta )=\{\begin{array}{cc}0\hfill & \text{if}\eta <U\hfill \\ t\hfill & \text{if}\eta >U.\hfill \end{array}$$
If $`U>t`$, then there are no non-trivial solutions $`\lambda 0`$. On the other hand, if $`U<t`$, then there exists always a unique solution $`\lambda 0`$. It is obtained as the solution of:
$$t=\sqrt{U^2+4t^2|\lambda |^2}$$
given by:
$$|\lambda |=\frac{1}{2t}\left(t^2U^2\right)^{1/2}<\frac{1}{2}.$$
Since $`f_\beta (\eta )`$ is increasing with $`T`$ decreasing (15), this also implies an expression for a supremum on the condensate density $`\rho _0`$ (6):
$$\rho _0\frac{\left(t^2U^2\right)}{2t^2}<\rho .$$
(22)
The condensate density is always strictly smaller than the total particle density, which is compatible with former results on this topic . This finishes the study of the equilibrium properties of the model (1).
## 3 Spontaneous Symmetry Breaking
The model $`H_\mathrm{\Lambda }`$ (1) and its equilibrium states have a discrete symmetry consisting of the interchange of the two particle types into each other. It is implemented by the following unitary matrix: $`u=\sigma _1^+\sigma _1^{}\sigma _2^+\sigma _2^{}+\sigma _1^{}\sigma _1^+\sigma _2^{}\sigma _2^++\sigma _1^+\sigma _2^{}+\sigma _1^{}\sigma _2^+`$, then for any local observable $`A_{12}`$ depending on the types $`1`$ and $`2`$, one gets:
$$\gamma _{12}(A_{12})=\left(\underset{i\mathrm{\Lambda }}{}(u_{12})_i\right)A_{12}\left(\underset{i\mathrm{\Lambda }}{}(u_{12})_i\right)=A_{21}.$$
The eigenvectors (10) of $`h_\lambda `$ satisfy
$$u_{12}|\varphi _0=|\varphi _0;u_{12}|\varphi _1=|\varphi _1;u_{12}|\varphi _2=|\varphi _2;u_{12}|\varphi _3=|\varphi _3.$$
A second discrete symmetry is the particle-hole ($`ph`$) symmetry. This symmetry interchanges the creation and annihilation operators ($`\sigma _i^+`$ and $`\sigma _i^{}`$) and is implemented by the unitaries $`\sigma _{12}^x=\sigma _1^x\sigma _2^x`$ ($`\sigma ^x`$ is the x-component Pauli matrix). The invariance of $`H_\mathrm{\Lambda }`$ (1) is expressed by:
$$\begin{array}{cc}\hfill \gamma _{ph}(H_\mathrm{\Lambda })& =\left(\underset{i\mathrm{\Lambda }}{}(\sigma _{12}^x)_i\right)H_\mathrm{\Lambda }\left(\underset{i\mathrm{\Lambda }}{}(\sigma _{12}^x)_i\right)\hfill \\ & =H_\mathrm{\Lambda }.\hfill \end{array}$$
The eigenfunctions (10) of $`h_\lambda `$ satisfy
$$\sigma _{12}^x|\varphi _0=|\varphi _0;\sigma _{12}^x|\varphi _1=|\varphi _1;\sigma _{12}^x|\varphi _2=|\varphi _2;\sigma _{12}^x|\varphi _3=|\varphi _3.$$
The Hamiltonian $`H_\mathrm{\Lambda }`$ (1) has also a two-dimensional continuous gauge symmetry group given by the actions:
$$\gamma _{\varphi _1\varphi _2}(\sigma _k^+)=\mathrm{e}^{\mathrm{i}\varphi _k}\sigma _k^+,\varphi _k[0,2\pi ],k=1,2.$$
(23)
For a finite volume $`\mathrm{\Lambda }`$, it is implemented by the unitaries:
$$U_{\mathrm{\Lambda }(\varphi _1,\varphi _2)}=\mathrm{exp}\frac{\mathrm{i}}{2}\left(\varphi _1Q_{1,\mathrm{\Lambda }}+\varphi _2Q_{2,\mathrm{\Lambda }}\right);$$
where
$$Q_{k,\mathrm{\Lambda }}=\underset{j\mathrm{\Lambda }}{}(\sigma _k^z)_j;k=1,2.$$
Clearly, $`\gamma _{\varphi _1\varphi _2}(H_\mathrm{\Lambda })=U_{\mathrm{\Lambda }(\varphi _1,\varphi _2)}H_\mathrm{\Lambda }U_{\mathrm{\Lambda }(\varphi _1,\varphi _2)}^{}=H_\mathrm{\Lambda }`$. Hence, per lattice site one has a commutative two-dimensional Lie-algebra of this symmetry group, generated by $`\sigma _1^z`$ and $`\sigma _2^z`$. However, the equilibrium states $`\omega _{\beta \lambda }`$ with $`\lambda 0`$ (see section 2) break this symmetry down. Indeed we have here for $`k=1,2`$:
$$\omega _{\beta \lambda }\left(\gamma _{\varphi _1\varphi _2}\left(\sigma _k^+\right)\right)=\mathrm{e}^{\mathrm{i}\varphi _k}\omega _{\beta \lambda }(\sigma _k^+)\omega _{\beta \lambda }(\sigma _k^+)$$
In contrast to all other models for which we have been studying the Goldstone phenomenon , here we have the spontaneous symmetry breakdown of a two-dimensional group. In the following, we construct again the canonical normal modes of the Goldstone bosons together with their corresponding plasmon frequencies. We take a symmetry breaking state $`\omega _{\beta \lambda }`$ with $`\lambda 0`$ and we consider fluctuation mode operators. Before constructing the relevant canonical normal coordinates of the Goldstone particles in the next section, we introduce here the notion of fluctuation operators.
For the mathematical details we refer to . Denote by $`M_4^{sa}=(M_2M_2)^{sa}`$ the vectorspace of the Hermitian $`4\times 4`$ matrices, i.e. these are one-site observables.
For any $`AM_4^{sa}`$, its fluctuation operator $`F_\mathrm{\Lambda }(A)`$ for the volume $`\mathrm{\Lambda }`$ in the state $`\omega _{\beta \lambda }`$ is as usual defined by:
$$F_{\mathrm{\Lambda },\lambda }(A)=\frac{1}{\sqrt{|\mathrm{\Lambda }|}}\underset{x\mathrm{\Lambda }}{}\left(A_x\omega _{\beta \lambda }(A)\right),$$
where $`A_x`$ is a copy of $`A`$ on site $`xZ`$. In it is proved that the limit:
$$\underset{\mathrm{\Lambda }Z}{lim}F_{\mathrm{\Lambda },\lambda }(A)F_\lambda (A)$$
exists in the sense of a non-commutative central limit as an operator on a Hilbert space $`_\lambda `$, generated by a normalised vector $`\stackrel{~}{\mathrm{\Omega }}_\lambda `$, and vectors $`F_\lambda (A_1)\mathrm{}F_\lambda (A_n)\stackrel{~}{\mathrm{\Omega }}_\lambda `$, with $`A_iM_4^{sa}`$, $`n`$ is arbitrary and a scalar product:
for $`A_i,B_jM_4^{sa}`$,
$$(F_\lambda (A_1)\mathrm{}F_\lambda (A_n)\stackrel{~}{\mathrm{\Omega }}_\lambda ,F_\lambda (B_1)\mathrm{}F_\lambda (B_m)\stackrel{~}{\mathrm{\Omega }}_\lambda )=\delta _{nm}\mathrm{Perm}.\left((\stackrel{~}{\mathrm{\Omega }}_\lambda ,F_\lambda (A_i)F_\lambda (B_j)\stackrel{~}{\mathrm{\Omega }}_\lambda )\right)_{i,j=1,\mathrm{}n}$$
(24)
where the two-point function is given by:
$$(\stackrel{~}{\mathrm{\Omega }}_\lambda ,F_\lambda (A_i)F_\lambda (B_j)\stackrel{~}{\mathrm{\Omega }}_\lambda )=\omega _{\beta \lambda }(AB)\omega _{\beta \lambda }(A)\omega _{\beta \lambda }(B).$$
(25)
i.e. by the truncated two-point functions of the given state $`\omega _{\beta \lambda }`$. Furthermore, as already could be understood from (24), the fluctuation operators $`F_\lambda (A),F_\lambda (B),\mathrm{}`$ satisfy the canonical commutation relations:
$$[F_\lambda (A),F_\lambda (B)]=\omega _{\beta \lambda }([A,B])\mathrm{𝟏}$$
(26)
i.e. the fluctuation operators are quantum canonical variables (compare $`[q,p]=\mathrm{i}\mathrm{}`$) with quantisation parameter $`\omega _{\beta \lambda }([A,B])`$ (instead of $`\mathrm{i}\mathrm{}`$) and the state $`(\stackrel{~}{\mathrm{\Omega }}_\lambda ,.\stackrel{~}{\mathrm{\Omega }}_\lambda )`$ is a generalised free or quasi-free state of a boson field algebra of observables. In section 4, we look for the relevant operators $`A,B,\mathrm{}`$ in order to obtain the normal canonical variables for the Goldstone particles, i.e. the particles generated by the spontaneous symmetry breaking. We derive the dynamical equations for these modes and we solve them.
## 4 Spontaneous Symmetry Breaking and Plasma Oscillations
We will consider fluctuations of the generators
$$Q_\pm =\sigma _1^z\pm \sigma _2^z$$
of the broken symmetry.
Let $`P_i,(i=0,1,2,3)`$ be the spectral projections on the eigenfunctions $`|\varphi _i`$ corresponding to the eigenvalues $`ϵ_i`$ (9) of $`h_\lambda `$ (8), i.e. $`h_\lambda =_{i=0}^3ϵ_iP_i`$. Consider for any local observable $`AM_4^{sa}`$ the Hermitian operators:
$$E_{ij}(A)=P_iAP_j+P_jAP_i;i<j;i,j=0,1,2,3;$$
(27)
and the adjoints:
$$JE_{ij}(A)=\mathrm{i}\left(P_iAP_jP_jAP_i\right).$$
(28)
An immediate calculation leads to:
$$\begin{array}{cc}& [h_\lambda ,E_{kl}(A)]=\mathrm{i}(ϵ_kϵ_l)JE_{kl}(A);\hfill \\ & [h_\lambda ,JE_{kl}(A)]=\mathrm{i}(ϵ_kϵ_l)E_{kl}(A).\hfill \end{array}$$
(29)
Now we take $`A=Q_+`$, respectively $`A=Q_{}`$. Due to the particle-hole and type-exchange symmetries (section 3), the only non-zero operators of the type (27-28) are:
$$E_{0i}(Q_{}),JE_{0i}(Q_{}),E_{1i}(Q_+),JE_{1i}(Q_+)i=2,3.$$
(30)
It is interesting to write out these operators in terms of the Pauli matrices $`\sigma _i^\pm ,\sigma _i^z,i=1,2`$. As an example we give the expression for the first pair of these operators $`(E_{02}(Q_{}),JE_{02}(Q_{}))`$, the others are left as an exercise for the reader:
$`E_{02}(Q_{})`$ $`={\displaystyle \frac{\xi _{}}{2\eta }}Q_{}+{\displaystyle \frac{t\lambda }{\eta }}(\sigma _1^z\sigma _2^+\sigma _2^z\sigma _1^+)+{\displaystyle \frac{t\overline{\lambda }}{\eta }}(\sigma _1^z\sigma _2^{}\sigma _2^z\sigma _1^{});`$
$`JE_{02}(Q_{})`$ $`=\mathrm{i}{\displaystyle \frac{\xi _{}}{2\eta }}(\sigma _1^+\sigma _2^{}\sigma _2^+\sigma _1^{})+\mathrm{i}{\displaystyle \frac{t\lambda }{\eta }}(\sigma _1^+\sigma _2^+)\mathrm{i}{\displaystyle \frac{t\overline{\lambda }}{\eta }}(\sigma _1^{}\sigma _2^{}).`$
$`E_{02}(Q_{})`$ is a component of the infinitesimal generator $`Q_{}`$ and $`JE_{02}(Q_{})`$ is the adjoint order parameter operator to $`E_{02}(Q_{})`$.
It is clear from (9) that the dynamical equations (29) for the operators (30) contain only the relative energies $`ϵ_kϵ_l`$ given by:
$$\xi _\pm =\eta \pm U.$$
In particular, the energy gap $`\xi _+`$ is entering in the equations:
$$[h_\lambda ,E_{02}(Q_{})]=\mathrm{i}\xi _+JE_{02}(Q_{}),[h_\lambda ,JE_{02}(Q_{})]=\mathrm{i}\xi _+E_{02}(Q_{}),$$
(31)
$$[h_\lambda ,E_{13}(Q_+)]=\mathrm{i}\xi _+JE_{13}(Q_+),[h_\lambda ,JE_{13}(Q_+)]=\mathrm{i}\xi _+E_{13}(Q_+);$$
(32)
and the energy gap $`\xi _{}`$ in:
$$[h_\lambda ,E_{03}(Q_{})]=\mathrm{i}\xi _{}JE_{03}(Q_{}),[h_\lambda ,JE_{03}(Q_{})]=\mathrm{i}\xi _{}E_{03}(Q_{}),$$
(33)
$$[h_\lambda ,E_{12}(Q_+)]=\mathrm{i}\xi _{}JE_{12}(Q_+),[h_\lambda ,JE_{12}(Q_+)]=\mathrm{i}\xi _{}E_{12}(Q_+).$$
(34)
All these results are on the level of the microscopic observables. In particular one discovers in (31-34) that the dynamics leaves invariant the pairs $`(E_{02}(Q_{}),JE_{02}(Q_{}))`$, $`(E_{13}(Q_+),JE_{13}(Q_+))`$, $`(E_{03}(Q_{}),JE_{03}(Q_{}))`$ and $`(E_{12}(Q_+),JE_{12}(Q_+))`$. On the other hand the different pairs do not form canonical pairs of variables. Therefore, we have to consider the macroscopic variables, namely the fluctuations of these quantities. We know from section 3 that the fluctuation operators are canonical variables, satisfying the canonical commutation relations.
We treat explicitly the first pair appearing in the equation (31), and give the results for the other pairs without technical details.
Plasmon frequency $`\xi _+`$
With the notations of section 3, denote the fluctuations:
$`X_{02}`$ $`{\displaystyle \frac{1}{n_{02}}}F_\lambda \left(E_{02}(Q_{})\right);`$ $`P_{02}`$ $`{\displaystyle \frac{1}{n_{02}}}F_\lambda \left(JE_{02}(Q_{})\right),`$ (35)
where
$$n_{02}=\left(\frac{\mathrm{e}^{\beta U}+\mathrm{e}^{\beta \eta }}{Z}\right)^{1/2};Z=2\mathrm{cosh}(\beta U)+2\mathrm{cosh}(\beta \eta ).$$
These satisfy (following (26)) the canonical commutation relation:
$$[X_{02},P_{02}]=\frac{1}{n_{02}^2}\omega _{\beta \lambda }\left([E_{02}(Q_{}),JE_{02}(Q_{})]\right)\mathrm{𝟏}=\mathrm{i}\mathrm{}_+$$
(36)
with quantisation parameter:
$$\mathrm{}_+=\frac{4\xi _{}}{\eta }\mathrm{tanh}\left(\beta \xi _+/2\right).$$
(37)
Following (25), the calculations for the variances are straightforward:
$$(\stackrel{~}{\mathrm{\Omega }}_\lambda ,X_{02}^2\stackrel{~}{\mathrm{\Omega }}_\lambda )=(\stackrel{~}{\mathrm{\Omega }}_\lambda ,P_{02}^2\stackrel{~}{\mathrm{\Omega }}_\lambda )=2\frac{\xi _{}}{\eta }.$$
(38)
Now compute the dynamics of these fluctuation operators induced by the microdynamics determined by the Hamiltonian $`h_\lambda `$ (5). In differential form the dynamics follows immediately from the equation (31):
$$\frac{d}{idt}\left(F_\lambda \left(E_{02}(Q_{})\right)\right)=F_\lambda \left([h_\lambda ,E_{02}(Q_{})]\right)=\mathrm{i}\xi _+F_\lambda \left(JE_{02}(Q_{})\right);$$
$$\frac{d}{idt}\left(F_\lambda \left(JE_{02}(Q_{})\right)\right)=F_\lambda \left([h_\lambda ,JE_{02}(Q_{})]\right)=\mathrm{i}\xi _+F_\lambda \left(E_{02}(Q_{})\right).$$
Using the notations (35) and in integrated form one gets:
$$\begin{array}{cc}\hfill X_{02}(t)=& X_{02}(0)\mathrm{cos}(\xi _+t)+P_{02}(0)\mathrm{sin}(\xi _+t);\hfill \\ \hfill P_{02}(t)=& X_{02}(0)\mathrm{sin}(\xi _+t)+P_{02}(0)\mathrm{cos}(\xi _+t).\hfill \end{array}$$
(39)
i.e. the pair of macroscopic variables $`(X_{02},P_{02})`$ is a canonical pair (36), evolving in time, independently from all other degrees of freedom of the system, as a pair of quantum harmonic oscillator variables oscillating with a frequency $`\xi _+`$. Remark that if the temperature $`T`$ tends to the critical temperature $`T_c`$, i.e. when $`|\lambda |0`$, then $`\xi _{}0`$. It follows from the equations (38) that the pair disappears completely i.e. $`X_{02}=P_{02}=0`$ for $`TT_c`$.
This finishes the study of the first pair of fluctuation variables.
In a completely analogous way one defines the second pair of canonical fluctuation operators:
$`X_{13}`$ $`={\displaystyle \frac{1}{n_{13}}}F_\lambda \left(E_{13}(Q_+)\right);`$ $`P_{13}`$ $`={\displaystyle \frac{1}{n_{13}}}F_\lambda \left(JE_{13}(Q_+)\right);`$ (40)
with
$$n_{13}=\left(\frac{\mathrm{e}^{\beta U}+\mathrm{e}^{\beta \eta }}{Z}\right)^{1/2}.$$
One gets, comparable to (36), (38) and (39):
$$[X_{13},P_{13}]=\mathrm{i}\mathrm{}_+;$$
(41)
$$(\stackrel{~}{\mathrm{\Omega }}_\lambda ,X_{13}^2\stackrel{~}{\mathrm{\Omega }}_\lambda )=(\stackrel{~}{\mathrm{\Omega }}_\lambda ,P_{13}^2\stackrel{~}{\mathrm{\Omega }}_\lambda )=2\frac{\xi _{}}{\eta };$$
(42)
$$\begin{array}{cc}\hfill X_{13}(t)=& X_{13}(0)\mathrm{cos}(\xi _+t)+P_{13}(0)\mathrm{sin}(\xi _+t);\hfill \\ \hfill P_{13}(t)=& X_{13}(0)\mathrm{sin}(\xi _+t)+P_{13}(0)\mathrm{cos}(\xi _+t).\hfill \end{array}$$
(43)
Plasmon frequency $`\xi _{}`$
Following (33), consider the third pair of canonical fluctuation operators
$`X_{03}`$ $`{\displaystyle \frac{1}{n_{03}}}F_\lambda \left(E_{03}(Q_{})\right);`$ $`P_{03}`$ $`{\displaystyle \frac{1}{n_{03}}}F_\lambda \left(JE_{03}(Q_{})\right),`$ (44)
where
$$n_{03}=\left(\frac{\mathrm{e}^{\beta U}+\mathrm{e}^{\beta \eta }}{Z}\right)^{1/2}.$$
One gets
$$[X_{03},P_{03}]=\mathrm{i}\mathrm{}_{}$$
(45)
with a new quantisation parameter,
$$\mathrm{}_{}=\frac{4\xi _+}{\eta }\mathrm{tanh}\left(\beta \xi _{}/2\right);$$
(46)
variances,
$$(\stackrel{~}{\mathrm{\Omega }}_\lambda ,X_{03}^2\stackrel{~}{\mathrm{\Omega }}_\lambda )=(\stackrel{~}{\mathrm{\Omega }}_\lambda ,P_{03}^2\stackrel{~}{\mathrm{\Omega }}_\lambda )=2\frac{\xi _+}{\eta };$$
(47)
and dynamics
$$\begin{array}{cc}\hfill X_{03}(t)=& X_{03}(0)\mathrm{cos}(\xi _{}t)+P_{03}(0)\mathrm{sin}(\xi _{}t);\hfill \\ \hfill P_{03}(t)=& X_{03}(0)\mathrm{sin}(\xi _{}t)+P_{03}(0)\mathrm{cos}(\xi _{}t).\hfill \end{array}$$
(48)
Finally, following (34), consider the fourth pair of canonical fluctuation operators:
$`X_{12}`$ $`{\displaystyle \frac{1}{n_{12}}}F_\lambda \left(E_{12}(Q_+)\right);`$ $`P_{12}`$ $`{\displaystyle \frac{1}{n_{12}}}F_\lambda \left(JE_{12}(Q_+)\right),`$ (49)
where
$$n_{12}=\left(\frac{\mathrm{e}^{\beta U}+\mathrm{e}^{\beta \eta }}{Z}\right)^{1/2}.$$
One computes the commutator:
$$[X_{12},P_{12}]=\mathrm{i}\mathrm{}_{};$$
(50)
the variances,
$$(\stackrel{~}{\mathrm{\Omega }}_\lambda ,X_{12}^2\stackrel{~}{\mathrm{\Omega }}_\lambda )=(\stackrel{~}{\mathrm{\Omega }}_\lambda ,P_{12}^2\stackrel{~}{\mathrm{\Omega }}_\lambda )=2\frac{\xi _+}{\eta };$$
(51)
and the dynamics
$$\begin{array}{cc}\hfill X_{12}(t)=& X_{12}(0)\mathrm{cos}(\xi _{}t)+P_{12}(0)\mathrm{sin}(\xi _{}t);\hfill \\ \hfill P_{12}(t)=& X_{12}(0)\mathrm{sin}(\xi _{}t)+P_{12}(0)\mathrm{cos}(\xi _{}t).\hfill \end{array}$$
(52)
This completes the construction of the plasmon spectrum together with all its quantum normal modes. The next question is: are these four pairs (35,40,44,49) all different ? The question is relevant, because the mode variables are all fluctuations and coarse graining might make some of them equivalent, i.e. although $`A,BM_4^{sa}`$ and $`AB`$, $`F_\lambda (A)`$ might be equal to $`F_\lambda (B)`$. The problem of equivalence of fluctuation operators is considered before (see ). It turns out that they are equivalent, i.e. $`F_\lambda (A)=F_\lambda (B)`$, if and only if
$$(\stackrel{~}{\mathrm{\Omega }}_\lambda ,F_\lambda (AB)^2\stackrel{~}{\mathrm{\Omega }}_\lambda )=\omega _{\beta \lambda }\left(\left(AB\omega _{\beta \lambda }\left(AB\right)\right)^2\right)=0.$$
(53)
If we take for $`A,B`$ the operators $`E_{ij}(Q_\pm )`$ and $`JE_{ij}(Q_\pm )`$ with $`i<j`$, then clearly
$$\omega _{\beta \lambda }(E_{ij}(Q_\pm ))=\omega _{\beta \lambda }(JE_{ij}(Q_\pm ))=0$$
and we have to check
$$\omega _{\beta \lambda }\left(\left(AB\right)^2\right)=\omega _{\beta \lambda }\left(A^2+B^2ABBA\right)=0;$$
(54)
for $`A`$ and $`B`$ taken out of different pairs $`(E_{02}(Q_{}),JE_{02}(Q_{}))`$, $`(E_{13}(Q_+),JE_{13}(Q_+))`$, $`(E_{03}(Q_{}),JE_{03}(Q_{}))`$, $`(E_{12}(Q_+),JE_{12}(Q_+))`$. It is immediately checked that in this case always $`\omega _{\beta \lambda }\left(AB\right)=\omega _{\beta \lambda }\left(BA\right)=0`$. Hence, (54) reduces to
$$\omega _{\beta \lambda }\left(\left(AB\right)^2\right)=\omega _{\beta \lambda }\left(A^2+B^2\right)0;$$
which follows from (38),(42),(47) and (51). Hence, the four pairs of canonical fluctuation operators are all independent, proving the two-fold degeneracy of the plasmon frequency $`\xi _+`$ and $`\xi _{}`$.
Finally, remark that if the temperature $`T`$ tends to the critical temperature, i.e. when $`\lambda 0`$, then the pairs $`(X_{02},P_{02})`$ (35) and $`(X_{13},P_{13})`$ (40) disappear completely. On the other hand the pairs $`(X_{03},P_{03})`$ (44) and $`(X_{12},P_{12})`$ (49) become classical variables (see (45) and (50)), and become constants of the motions (see (48) and (52)). In fact for $`TT_c`$ the continuous symmetry group (23) is restored. The fluctuation operators $`X_{03}`$ and $`X_{12}`$ are the fluctuations of the generators of this symmetry. The fluctuation operators $`P_{03}`$ and $`P_{12}`$ are the total particle current and relative particle current fluctuations. For $`T<T_c`$ the physical interpretation of these operators remains the same with the basic difference that they show a definite collective or macroscopic quantum character, and a non-trivial oscillatory time behaviour. We remind that the results derived for this model are mathematically rigorous. In view of the recent interest in interferences between Bose-Einstein condensates with Josephson-like coupling an extension of our model including such a coupling raises over as the natural thing to do next. We leave it for a future contribution.
Acknowledgement: The authors thank Tom Michoel for helpful discussions.
|
warning/0003/math0003208.html
|
ar5iv
|
text
|
# Energy levels of steady states for thin film type equations
## 1. Introduction
We study the evolution equation
(1)
$$h_t=(f(h)h_{xxx})_x(g(h)h_x)_x.$$
This is the one dimensional version of $`h_t=(f(h)\mathrm{\Delta }h)(g(h)h)`$, which has been used to model the dynamics of a thin film of viscous liquid. The air/liquid interface is at height $`z=h(x,y,t)0`$ and the liquid/solid interface is at $`z=0`$. The one dimensional equation (1) applies if the liquid film is uniform in the $`y`$ direction.
The fourth order term in the equation reflects surface tension effects and the second order term can reflect gravity, van der Waals interactions, thermocapillary effects or the geometry of the solid substrate, for example. Typically $`f(h)=h^3+\beta h^p`$ where $`0<p<3,\beta 0`$, and $`g(h)\pm h^m`$ as $`h0`$, where $`m`$. In certain applications $`g(h)`$ changes sign at some positive $`h`$. We refer to for reviews of the physical and modeling literature.
The extensively studied Cahn–Hilliard equation also has the form (1), with $`f1`$ and $`g(h)=13h^2`$. See for further references on the Cahn–Hilliard equation.
Equations like (1) are of mathematical as well as physical interest: for example, Bertozzi and Pugh conjectured that blow-up ($`h(,t)_{\mathrm{}}\mathrm{}`$) is possible in some cases (e.g. if $`f(h)=h^n,g(h)=h^m`$ with $`m>n+2`$), and they have proved finite time blow-up for $`f(h)=h`$ and $`g(h)=h^m`$ when $`m3`$. In \[17, §8\] we related the steady states and some of their properties to this blow-up conjecture.
### Background and goals.
In we proved linear stability and instability results for the positive periodic steady states of (1). Periodicity should not be regarded as a constraint, since if $`f,g>0`$ then positive bounded steady states must be periodic or constant, by \[17, Theorem B.1\]. And periodic steady states do exist for many equations of type (1), by the methods of \[17, §2.2\] or , for example. Compactly supported ‘droplet’ steady states only exist, though, if $`g/f`$ satisfies additional constraints \[17, §2.2\], and can have relatively low regularity at the contact line.
In this paper we concentrate mostly on positive periodic steady states, and on droplet steady states with zero contact angles. Our main investigative tool is the energy
$$(h(,t))=_0^X\left[\frac{1}{2}h_x(x,t)^2H(h(x,t))\right]𝑑x;$$
here $`h(x,t)`$ is a smooth solution of (1) that is $`X`$-periodic in $`x`$, and $`H(y)`$ satisfies $`H^{\prime \prime }=g/f`$. This energy is strictly dissipated: $`(d/dt)(h(,t))0`$ with equality if and only if $`h`$ is a steady state (see §2.1). Thus the energy is a Liapunov function for the evolution.
We address two questions about the energy landscape of the evolution (1).
1. Which steady states are local minima of the energy? Which are saddle points?
2. Among steady states having the same period and the same area (i.e. fluid volume), which has the lowest energy?
Answering these questions will help clarify the phase portrait of the evolution. For example, if all small zero–mean perturbations of a steady state can be shown to raise the energy, then the steady state might be asymptotically stable: it might be that every smooth solution starting from near the steady state relaxes back to the steady state as $`t\mathrm{}`$. But if some zero–mean perturbation decreases the energy, then asymptotic stability definitely fails.
Our requirement that the perturbations have zero mean seems reasonable from a physical standpoint, because it corresponds to a disturbance of the fluid that alters the profile without adding additional fluid. Mathematically it is reasonable because the evolution equation (1) preserves volume for spatially periodic solutions: $`h(x,t)𝑑x=h(x,0)𝑑x`$ for all time $`t`$. Thus zero–mean perturbations allow the possibility of relaxation back to the original steady state, while nonzero–mean perturbations do not.
### A sketch of definitions and results.
Take $`X>0`$. If there is an $`X`$-periodic zero–mean perturbation $`v`$ such that $`(h_{\mathrm{ss}}+\epsilon v)<(h_{\mathrm{ss}})`$ for all small $`ϵ>0`$, then we call the steady state $`h_{\mathrm{ss}}`$ ‘energy unstable’ at period $`X`$. If instead $`(h_{\mathrm{ss}}+\epsilon v)>(h_{\mathrm{ss}})`$ for all sufficiently small $`\epsilon >0`$, for each $`X`$-periodic zero–mean perturbation $`v`$, then we call the steady state ‘energy stable’ at period $`X`$. (Some authors call this formal stability .) It is conceivable that a steady state might be energy stable and yet not be a local minimum of the energy.
Our main stability results, in Section 2, are roughly stated as follows.
* Theorem 1. For positive periodic steady states, linear instability implies energy instability. Hence our linear instability results in imply that every positive $`X^{}`$-periodic steady state is energy unstable at periods $`X=2X^{},3X^{},\mathrm{}`$, and is also energy unstable at period $`X^{}`$ if $`g/f`$ is a strictly convex function.
* Theorems 23. Further, for the ‘power law’ coefficients $`f(y)=y^n`$ and $`g(y)=y^m`$ with $`>0`$, we completely characterize energy stability at period $`X^{}`$ even when $`g/f`$ is not convex, that is, when $`mn(0,1)`$.
In Section 2.6 we explain how these results for periodic boundary conditions cover the case of Neumann (‘no flux’) boundary conditions as well.
Then in Section 3 we determine the relative energy levels of three different kinds of steady state: constant steady states, positive periodic steady states, and zero contact angle droplet steady states. Figure 1 illustrates these three steady states, as well as showing a nonzero contact angle droplet steady state (about which we say little in this paper).
We have found it too difficult to determine energy levels when working with arbitrary coefficients $`f`$ and $`g`$, but have obtained fairly complete answers for power law coefficients. This provides at least some insight into the general case. Further insight comes from the work of Grinfeld and Novick–Cohen on the energy levels of steady states for the Cahn–Hilliard equation, which has non–power–law coefficients. The earlier work of Mischaikow applies to a variety of gradient-like bistable equations.
Our main energy level results, in Section 3, for $`f(y)=y^n`$ and $`g(y)=y^m`$ are:
* Theorem 6. If $`m<n`$ or $`mn+1`$ then a positive periodic steady state always has higher energy than the constant steady state with the same mean value. For $`n<mn+0.75`$ our analytical and numerical work suggest the positive periodic steady state has lower energy than the constant steady state.
* Theorem 9. When $`mn+0.77`$ there can be two steady states with the same period and area, $`h_{\mathrm{ss1}}`$ and $`h_{\mathrm{ss2}}`$, with $`\mathrm{min}_xh_{\mathrm{ss1}}(x)<\mathrm{min}_xh_{\mathrm{ss2}}(x)`$; we essentially prove $`h_{\mathrm{ss2}}`$ is energy unstable and has higher energy than $`h_{\mathrm{ss1}}`$, which is energy stable.
* Theorem 7. If $`mn(2,0)[1,2)`$ then a positive periodic steady state always has higher energy than a zero contact angle droplet steady state with the same mean value.
* Theorem 11. The constant steady state can have higher energy than the zero contact angle droplet steady state with the same mean value. When $`mn[1,2)`$, for example, a mountain pass scenario can occur, in which the constant steady state is a local minimum of the energy, the positive periodic steady state is an energy unstable saddle, and there is a zero contact angle droplet having lower energy than either of them.
For example, Theorem 6 with $`m=1<n=3`$ covers the ‘van der Waals’ case
(2)
$$h_t=(h^3h_{xxx})_xA(h^1h_x)_x$$
with $`A>0`$. This equation has been studied by a number of other authors, e.g. , mostly with regard to similarity solutions and film rupture (where the solution goes to zero in finite time, at some point). Our numerical work on this equation pays particular attention to behavior near steady states, and how this evolves into rupture.
Our energy level results suggest possible basins of attraction around the stable steady states, and possible heteroclinic connections between steady states. In the companion article \[19, §4\], we investigate such possibilities with numerical simulations. For example, when $`mn+0.77`$ as in Theorem 9, we find robust heteroclinic connections between the unstable positive periodic steady state $`h_{\mathrm{ss2}}`$ and the stable one $`h_{\mathrm{ss1}}`$. For the mountain pass scenario in Theorem 11, we find that perturbing the saddle point (the periodic steady state) in one direction leads to relaxation to the constant steady state and perturbing in the opposite direction gives apparent relaxation to a droplet. A similar dichotomy was found for axisymmetric surface diffusion by Bernoff, Bertozzi and Witelski, \[4, p. 744\], with perturbed unduloids either relaxing to a cylinder or else pinching off in finite time.
We also present in \[19, §5\] simulations suggesting that small changes in the ‘mobility’ coefficient $`f`$ do not break heteroclinic orbits, but can affect whether or not the solution remains positive as it evolves.
We discuss some of these conclusions and future directions further in Section 6.
### Terminology.
We write $`𝕋_X`$ for a circle of circumference $`X>0`$. As usual, one identifies functions on $`𝕋_X`$ with functions on $``$ that are $`X`$-periodic and calls them even or odd according to whether they are even or odd on $``$.
A positive periodic steady state is assumed to satisfy the steady state equation classically. A droplet steady state $`h_{\mathrm{ss}}(x)`$ (see Figures 1c and d) is by definition positive on some interval $`(a,b)`$ and zero elsewhere, with $`h_{\mathrm{ss}}C^1[a,b]`$; we require $`h_{\mathrm{ss}}`$ to satisfy the steady state equation on the open interval $`(a,b)`$ only, and to have equal acute contact angles: $`0h_{\mathrm{ss}}^{}(a)=h_{\mathrm{ss}}^{}(b)<\mathrm{}`$. (Throughout the paper, if a function has only one independent variable then we use to denote differentiation with respect to that variable: $`h_{\mathrm{ss}}^{}=(h_{\mathrm{ss}})_x`$.)
We say a droplet steady state $`h_{\mathrm{ss}}`$ has ‘zero contact angle’ if $`0=h_{\mathrm{ss}}^{}(a)=h_{\mathrm{ss}}^{}(b)`$, and ‘nonzero contact angle’ otherwise. A ‘configuration’ of droplet steady states is a collection of steady droplets whose supports are disjoint. For more on the steady states and their properties, see .
## 2. Energy stability for periodic steady states
We assume throughout this section that $`f(y)`$ and $`g(y)`$ are $`C^2`$-smooth for $`y>0`$, and that $`f>0`$. Define
$$r=\frac{g}{f}.$$
Take $`X>0`$.
We investigate stability and dynamical questions by means of a Liapunov energy. A few of the theorems follow directly from our linear stability results in , but most are quite different and complementary.
### 2.1. Definition of the energy, and of energy instability
The energy function for the evolution equation (1) is defined for $`\mathrm{}H^1(𝕋_X)`$ to be
(3)
$$(\mathrm{})=_0^X\left[\frac{1}{2}(\mathrm{}^{})^2H(\mathrm{})\right]𝑑x,$$
where $`H(y)`$ is a function with $`H^{\prime \prime }=r=g/f`$.
To verify the energy $``$ is a Liapunov function for the evolution (1), suppose $`h(x,t)`$ is a positive smooth solution of (1) that is $`X`$-periodic in $`x`$. Bertozzi and Pugh \[5, §2\] observed (generalizing ) that $``$ is dissipated by the evolution:
$$\frac{d}{dt}(h(,t))=_0^X\frac{1}{f(h(x,t))}\left[f(h(x,t))h_{xxx}(x,t)+g(h(x,t))h_x(x,t)\right]^2𝑑x0.$$
The dissipation is strict at each time $`t`$ unless $`f(h)h_{xxx}+g(h)h_x=0`$ for all $`x`$. For smooth positive periodic solutions this occurs only when $`h(,t)`$ is a steady state.
Let $`h_{\mathrm{ss}}C^4(𝕋_X)`$ be a positive periodic steady state of (1). It is easy to see (cf. formula (4.1)) that $`h_{\mathrm{ss}}`$ is a critical point for the energy $``$, with respect to zero-mean perturbations.
###### Definition .
Call $`h_{\mathrm{ss}}`$ an energy unstable critical point (with respect to zero–mean perturbations at period $`X`$) if there exists a smooth $`X`$-periodic perturbation $`u(x)`$ with mean value zero such that
$$(h_{\mathrm{ss}}+\epsilon u)<(h_{\mathrm{ss}})\text{for all small }\epsilon >0\text{.}$$
That is, small perturbations in the direction $`u`$ decrease the energy. (Some authors call this formal instability .)
An energy unstable steady state is necessarily a saddle point in the energy landscape, since $``$ is increased by the perturbation $`u(x)=\epsilon \mathrm{cos}(2\pi kx/X)`$ for $`k1`$.
Energy unstable steady states are not asymptotically stable in $`H^1(𝕋_X)`$, in the following sense: suppose $`h(x,t)`$ is a positive smooth solution with initial data $`h_{\mathrm{ss}}+\epsilon u`$; then $`h(,t)\to ̸h_{\mathrm{ss}}()`$ in $`H^1(𝕋_X)`$ since for all $`t`$, $`(h(,t))(h(,0))=(h_{\mathrm{ss}}+\epsilon u)<(h_{\mathrm{ss}})`$ (convergence in $`H^1`$ would imply convergence in $`L^{\mathrm{}}`$ and hence convergence of the energy). In fact, $`h(,t)`$ cannot converge to any translate of $`h_{\mathrm{ss}}`$, for the same reason.
The last paragraph extends to nonnegative weak solutions if they also dissipate the energy. Weak solutions must sometimes be considered because solutions of (1) that are initially positive might not always remain so, and where they go to zero they can lose regularity. See Bertozzi and Pugh for existence of nonnegative weak solutions that dissipate the energy.
### 2.2. Energy instability results
In \[18, §2\] we linearized the evolution equation (1) around the positive periodic steady state $`h_{\mathrm{ss}}`$ and then reduced the linear stability question to determining the sign of the first eigenvalue of a certain self-adjoint fourth order linear operator. We will not repeat the linearization here, or re-state the linear stability results of . However we warn readers that when we say a steady state is ‘linearly stable’, we are including the neutrally stable case in which the first eigenvalue of the linearized operator is zero. This is unavoidable: the operator always has a zero eigenvalue in its spectrum, corresponding to an infinitesimal translation of the steady state in space (the evolution equation is translation invariant).
The next theorem states that if a steady state is linearly unstable then it is energy unstable. Also, we present some unstable directions, when $`g/f`$ is strongly convex.
###### Theorem 1.
Let $`f,gC^2(0,\mathrm{})`$ with $`f>0`$. Take $`X>0`$ and suppose $`h_{\mathrm{ss}}C^4()`$ is an $`X`$-periodic positive steady state of (1).
If $`h_{\mathrm{ss}}`$ is linearly unstable with respect to zero–mean perturbations at period $`X`$, then it is also energy unstable at period $`X`$.
In particular, $`h_{\mathrm{ss}}`$ is energy unstable at period $`X`$ if it is non-constant and either: the least period of $`h_{\mathrm{ss}}`$ is $`X/j`$ for some integer $`j2`$ or else $`r=g/f`$ is convex ($`r^{\prime \prime }0`$) and non-constant on the range of $`h_{\mathrm{ss}}`$. For example, if $`h_{\mathrm{ss}}`$ is non-constant and $`r^{\prime \prime }>0`$, then $`h_{\mathrm{ss}}`$ is energy unstable in the directions $`u=\pm h_{\mathrm{ss}}^{}`$ and $`\pm h_{\mathrm{ss}}^{\prime \prime }`$.
This is proved in Section 4.1. Note for example that the theorem covers the van der Waals evolution (2), since there $`r(y)=Ay^1/y^3=Ay^4`$ is strongly convex. Thus all positive periodic steady states of the van der Waals evolution are energy unstable.
These energy unstable steady states (which we observed above are not asymptotically stable) are presumably nonlinearly unstable, in general, but we cannot prove this.
In Theorem 1 we have assumed $`f,gC^2(0,\mathrm{})`$, which is generally the case for the thin film equations that are our main motivation. But our arguments are all local (involving only small perturbations of the steady state), and so the theorem still holds if the coefficient functions $`f(y)`$ and $`g(y)`$ are defined and $`C^2`$ merely for $`y`$-values in a neighborhood of $`[h_{\mathrm{ss}}^{}{}_{min}{}^{},h_{\mathrm{ss}}^{}{}_{max}{}^{}]`$.
### 2.3. Review of power law steady states and their rescalings
We now turn to power law coefficients: $`f(y)=y^n`$ and $`g(y)=y^m`$ for some exponents $`n,m`$ and some positive constant $`>0`$. Here
$$r(y)=y^{q1}$$
where
$$\overline{)q:=mn+1}$$
This exponent $`q`$ determines many properties of the steady states, including (usually) their linear stability.
The evolution equation (1) becomes
$$h_t=(h^nh_{xxx})_x(h^mh_x)_x.$$
To state our results on energy stability for this power law evolution, we first review some properties of the steady states and explain how to rescale them to solutions of a canonical nonlinear oscillator ODE, given as equation (89) below.
We start with a non-constant positive periodic steady state $`h_{\mathrm{ss}}C^4(𝕋_X)`$ of the general evolution (1). The steady state condition for (1) integrates to give $`f(h_{\mathrm{ss}})h_{\mathrm{ss}}^{\prime \prime \prime }+g(h_{\mathrm{ss}})h_{\mathrm{ss}}^{}=C`$ for some constant $`C`$. The least period of $`h_{\mathrm{ss}}`$ is $`X/j`$ for some integer $`j1`$.
One finds that the constant $`C`$ (the flux) equals zero, by dividing $`f(h_{\mathrm{ss}})h_{\mathrm{ss}}^{\prime \prime \prime }+g(h_{\mathrm{ss}})h_{\mathrm{ss}}^{}=C`$ by $`f(h_{\mathrm{ss}})>0`$ and integrating over a period (cf. ). Hence the steady state satisfies
(4)
$$h_{\mathrm{ss}}^{\prime \prime \prime }+r(h_{\mathrm{ss}})h_{\mathrm{ss}}^{}=0.$$
\[If $`h_{\mathrm{ss}}`$ were a droplet steady state then again $`C=0`$, by \[17, Theorem 2.1\], and equation (4) would hold wherever $`h_{\mathrm{ss}}`$ is positive.\]
Integrating, the steady states have a nonlinear oscillator formulation:
(5)
$$h_{\mathrm{ss}}^{\prime \prime }+H^{}(h_{\mathrm{ss}})=0$$
holds wherever $`h_{\mathrm{ss}}`$ is positive. Here $`H(y)`$ is a function with $`H^{\prime \prime }=r=g/f`$; if we regard $`x`$ as a ‘time’ variable then $`\frac{1}{2}h_{\mathrm{ss}}^{}(x)^2+H(h_{\mathrm{ss}}(x))`$ is a conserved quantity.
Returning to the power law evolution, remember $`r(y)=y^{q1}`$. Thus for $`q0`$ we can write the steady state equation (5) as
(6)
$$h_{\mathrm{ss}}^{\prime \prime }+\frac{h_{\mathrm{ss}}^qD}{q}=0$$
for some constant $`D`$. For $`q=0`$ the analogous equation is $`h_{\mathrm{ss}}^{\prime \prime }+\mathrm{log}h_{\mathrm{ss}}D=0`$. This oscillator equation involves three constants: $`q`$, $``$, and $`D`$. We remove $``$ and $`D`$ by rescaling: let
(7)
$$k(x)=\{\begin{array}{cc}\hfill \left(\frac{}{D}\right)^{1/q}h_{\mathrm{ss}}\left(\left(\frac{D}{}\right)^{1/2q}\frac{x}{D^{1/2}}\right),& q0,\hfill \\ \hfill e^{D/}h_{\mathrm{ss}}\left(e^{D/2}\frac{x}{^{1/2}}\right),& q=0.\hfill \end{array}$$
For $`q0`$ this rescaling uses that $`D>0`$, by \[17, §3.1\]. \[A different rescaling would be used to study droplet steady states \[17, §§3.2,4\].\] The steady state equation (6) rescales to
(8) $`k^{\prime \prime }+{\displaystyle \frac{k^q1}{q}}`$ $`=`$ $`0,q0,`$
(9) $`k^{\prime \prime }+\mathrm{log}k`$ $`=`$ $`0,q=0.`$
Differentiating, we find for all $`q`$ that $`k^{\prime \prime \prime }+k^{q1}k^{}=0`$, and so $`k`$ satisfies $`\left(k^nk^{\prime \prime \prime }+k^mk^{}\right)^{}=0`$, i.e. it is a steady state of $`k_t=(k^nk_{xxx})_x(k^mk_x)_x`$.
Since $`h_{\mathrm{ss}}`$ is non-constant, positive and periodic, we know $`k^{\prime \prime }(x_0)>0`$ for some point $`x_0`$. Evaluating (89) at $`x_0`$ shows the minimum value of $`k`$ is less than $`1`$. Also $`k^{}(0)=0`$ since (after a suitable translation) $`h_{\mathrm{ss}}`$ has its minimum at $`x=0`$. Introducing the notation $`k_\alpha `$ for the solution $`k`$ that has minimum value $`\alpha (0,1)`$, at $`x=0`$, we have
(10)
$$0<k_\alpha (0)=\alpha <1,k_\alpha ^{}(0)=0.$$
Thus every steady state $`h_{\mathrm{ss}}`$ can be rescaled to a $`k_\alpha `$, as above. Conversely, for each $`q`$ and $`\alpha (0,1)`$ there exists a unique smooth positive periodic $`k_\alpha `$ satisfying equations (89) and (10) (see \[17, Proposition 3.1\]). The same holds for $`\alpha =0`$ when $`q>1`$, except that $`k_0`$ may be only $`C^1`$-smooth at $`x=0`$ (see \[17, Theorem 3.2\]). To illustrate, Figure 2 plots the steady states $`k_\alpha `$ over two periods, for $`q=3`$ and eight $`\alpha `$-values between $`0`$ and $`1`$; see \[17, §6.1\] for details.
Note that the map $`(\alpha ,x)k_\alpha (x)`$ is $`C^{\mathrm{}}`$-smooth for $`(\alpha ,x)(0,1)\times `$, by an ODE theorem giving smooth dependence on the initial data \[15, Ch. V §4\]. We write
$$P=P(\alpha )\text{and}A=A(\alpha )$$
for the least period of $`k_\alpha `$ and for the area under its graph, $`A=_0^Pk_\alpha (x)𝑑x`$, respectively. Then $`P`$ and $`A`$ are smooth functions of $`\alpha `$ that approach $`2\pi `$ as $`\alpha 1`$, by \[18, Lemma 6\]. The function
$$E(\alpha ):=P(\alpha )^{3q}A(\alpha )^{q1}=P(\alpha )^2[A(\alpha )/P(\alpha )]^{q1}$$
determines whether the steady state is energy unstable or stable, in several results below.
The above rescaling ideas are a useful tool throughout the paper. We hope this tool does not obscure the fact that stability and energy level properties for equations of type (1) seem to be determined by the period map of a family of steady states $`h_{\mathrm{ss}}`$ with fixed area but varying amplitudes and periods, or alternatively the area map of a family of steady states with fixed period but varying amplitudes and areas. This is how one should think of the function $`E(\alpha )`$; see \[18, §6.3\] for more on this. The same underlying idea appears in the work of Grinfeld and Novick–Cohen on the Cahn–Hilliard equation, an equation which is not amenable to rescaling in the same way.
### 2.4. Energy in/stability for the power law evolution
###### Theorem 2.
Let $`h_{\mathrm{ss}}C^4()`$ be a non-constant positive $`X`$-periodic steady state of the power law evolution $`h_t=(h^nh_{xxx})_x(h^mh_x)_x`$. Translate $`h_{\mathrm{ss}}`$ to put its minimum at $`x=0`$ so that $`h_{\mathrm{ss}}`$ rescales to $`k_\alpha `$ for some $`\alpha (0,1)`$, as in §2.3.
If $`q<1`$ or $`q>2`$ then $`h_{\mathrm{ss}}`$ is energy unstable in the directions $`u=\pm h_{\mathrm{ss}}^{}`$ and $`u=\pm h_{\mathrm{ss}}^{\prime \prime }`$.
If $`q=2`$, or if $`q>1`$ and $`E^{}(\alpha )>0`$, then $`h_{\mathrm{ss}}`$ is energy unstable.
We prove the theorem in Section 4.2. Its first statement follows immediately from Theorem 1, since $`r(y)=y^{q1}`$ is strongly convex ($`r^{\prime \prime }>0`$) when $`q<1`$ or $`q>2`$. The final statement of the theorem certainly applies when $`q2`$, since then $`E^{}>0`$ by \[18, Theorem 11\]. Computational studies \[17, §6.1\] suggest $`E^{}(\alpha )>0`$ for all $`\alpha (0,1)`$ except when $`q[1,1.794]`$ (approximately); Figure 5 plots $`E(\alpha )`$ for certain of these $`q`$-values.
To help explain the appearance of the criterion $`E^{}(\alpha )>0`$ for energy instability, in the theorem, we refer the reader to \[18, §6.3\] and to \[19, §3\] for a ‘bifurcation diagram’ interpretation of the function $`\alpha E(\alpha )`$ in terms of a family of steady states $`h_{\mathrm{ss}}`$ with fixed area but varying amplitudes and periods, or fixed period but varying amplitudes and areas.
The case $`1<q<2`$ with $`E^{}(\alpha )0`$ is not covered by Theorem 2. By \[18, Theorem 9\] we know $`h_{\mathrm{ss}}`$ is linearly stable in this case. We have been unable to prove nonlinear stability, but we do prove in Theorem 3 below that if $`E^{}(\alpha )<0`$ then small perturbations of $`h_{\mathrm{ss}}`$ in every possible direction strictly increase the energy; this is consistent with nonlinear stability.
###### Definition .
Let $`h_{\mathrm{ss}}C^4(𝕋_X)`$ be a non-constant positive periodic steady state of (1). Call $`h_{\mathrm{ss}}`$ energy stable (with respect to zero–mean perturbations at period $`X`$) if for each $`uH^1(𝕋_X)\{0\}`$ with $`_0^Xu𝑑x=0`$ we have
$$(h_{\mathrm{ss}}+\epsilon u)>(h_{\mathrm{ss}})\text{for all small }\epsilon >0\text{.}$$
(Some authors would call this formal stability .)
A steady state might perhaps be energy stable without being a local minimum of the energy. (For example, the function $`f(x,y)=(yx^2)(y3x^2)`$ on $`^2`$ has the property that origin is a local minimum point on each straight line through the origin, though it is not a local minimum in $`^2`$.) Another cautionary note is that the energy is insensitive to translation: in particular, a steady state and its translates have the same energy, consistent with the translation invariance of the evolution equation itself. Numerical simulations in our article \[19, §4.4\] demonstrate that perturbations of an energy stable steady state can evolve towards a translate of that steady state. This suggests that any asymptotic or nonlinear stability result that can be proved will hold only up to translation.
In Section 4.3 we prove:
###### Theorem 3.
Let $`h_{\mathrm{ss}}C^4()`$ be a non-constant positive periodic steady state of the power law equation $`h_t=(h^nh_{xxx})_x(h^mh_x)_x`$, with $`h_{\mathrm{ss}}`$ having least period $`X`$. Translate $`h_{\mathrm{ss}}`$ to put its minimum at $`x=0`$ so that $`h_{\mathrm{ss}}`$ rescales to $`k_\alpha `$ for some $`\alpha (0,1)`$, as in §2.3.
If $`1<q<2`$ and $`E^{}(\alpha )<0`$ then $`h_{\mathrm{ss}}`$ is energy stable.
The hypothesis $`E^{}(\alpha )<0`$ seems numerically to hold for all $`\alpha (0,1)`$ when $`1<q1.75`$, as indicated by Figure 3 in .
We have no energy stability result when $`1<q<2`$ and $`E^{}(\alpha )=0`$; fortunately it seems $`E^{}(\alpha )=0`$ for at most one $`\alpha `$-value, for each $`q`$, as shown numerically by Figures 3–5 in . Theorem 9 of does imply linear stability when $`E^{}(\alpha )=0`$, but \[18, Theorem 10(b)\] shows the space of neutrally stable directions is two dimensional (rather than one dimensional as when $`E^{}(\alpha )<0`$) and this might perhaps lead to instability.
We next address the $`q=1`$ case.
###### Lemma 4.
Let $`q=1`$ (i.e. $`m=n`$) and suppose $`h_{\mathrm{ss}}C^4()`$ is a non-constant positive periodic steady state of $`h_t=(h^nh_{xxx})_x(h^nh_x)_x`$ with least period $`X`$, and translate so that $`h_{\mathrm{ss}}`$ has its minimum at $`x=0`$. Then $`h_{\mathrm{ss}}`$ is not asymptotically stable in $`H^1(𝕋_X)`$ with respect to even perturbations.
###### Proof.
The steady state equation (6) with $`q=1`$ has general solution $`h_{\mathrm{ss}}(x)=D/+c\mathrm{cos}(\sqrt{}x)`$, where we have used that $`h_{\mathrm{ss}}`$ has an extremum at $`x=0`$. Hence the period is $`X=2\pi /\sqrt{}`$, and for all small $`\epsilon `$ we see that the perturbed function $`h_{\mathrm{ss}}(x)+\epsilon \mathrm{cos}(\sqrt{}x)`$ is another positive periodic steady state solution. Thus $`h_{\mathrm{ss}}`$ is not asymptotically stable in $`H^1(𝕋_X)`$ with respect to even perturbations. ∎
The steady state is of course not asymptotically stable with respect to general perturbations, since one can always perturb by translating the steady state a small distance (this remark applies for all $`q`$). However we are really not interested in such simple translational motion. Also, translational perturbations are not even permitted under the Neumann boundary conditions that we consider below, in Section 2.6.
To summarize, when $`q=1`$ the positive periodic steady states are linearly neutrally stable with respect to zero–mean perturbations of the same period, by or \[18, Lemma 8\], and are not asymptotically stable with respect to ‘even’ perturbations, by the above lemma. Our numerical simulations in \[19, §4.3\] and those of suggest that a wide range of small perturbations yield solutions relaxing to nearby positive periodic steady states, suggesting they are nonlinearly stable.
In the companion paper \[19, §4\] we illustrate the above stability and instability theorems for the power law evolution with a variety of numerical simulations. There we find not only the short time behavior suggested by the energy (in)stability results, but also some longer time limits that are suggested by the energy level results in Section 3.
### 2.5. Odd perturbations
Returning momentarily to general coefficients $`f`$ and $`g`$, in Section 4.4 we prove the energy increases under odd perturbations, when $`r`$ is concave.
###### Theorem 5.
Let $`h_{\mathrm{ss}}C^4()`$ be a non-constant positive periodic steady state of (1) with coefficient functions $`f,gC^2(0,\mathrm{}),f>0`$. Suppose $`h_{\mathrm{ss}}`$ has least period $`X`$, and translate $`h_{\mathrm{ss}}`$ to put its minimum at $`x=0`$.
If $`r=g/f`$ is strongly concave ($`r^{\prime \prime }<0`$) then for every nontrivial $`uH^1(𝕋_X)`$ that is odd and is such that $`h_{\mathrm{ss}}+u>0`$, we have $`(h_{\mathrm{ss}}+u)>(h_{\mathrm{ss}})`$.
The theorem is global since the perturbations are not required to be small, and it is consistent with asymptotic stability (although convergence to a translate of $`h_{\mathrm{ss}}`$ seems more likely than convergence to $`h_{\mathrm{ss}}`$ itself).
Theorem 5 applies in the power law case with $`1<q<2,r(y)=y^{q1}`$. Another example with strongly concave $`r`$ is an equation \[20, eqn. (24)\] describing the dynamics of a population of aphids, for which $`f(y)=y,g(y)=yc`$ and $`r(y)=1c/y`$.
### 2.6. Relation between the periodic and Neumann stability problems
Suppose $`h_{\mathrm{ss}}`$ is an even $`X`$-periodic steady state of the evolution equation (1) with extrema at $`x=0,\pm X/2,\mathrm{}`$, so that $`h_{\mathrm{ss}}^{}=h_{\mathrm{ss}}^{\prime \prime \prime }=0`$ at these points. As we observed at the end of \[18, §2.5\], linear instability of $`h_{\mathrm{ss}}`$ with respect to periodic boundary conditions on $`(X/2,X/2)`$ is equivalent to linear instability with respect to Neumann (‘no flux’) conditions on the half-interval $`(0,X/2)`$; these Neumann conditions are: $`h_x=h_{xxx}=0`$ at $`x=0,X/2`$.
The energy of a positive smooth solution is still dissipated in the case of Neumann boundary conditions, and obviously energy instability of $`h_{\mathrm{ss}}`$ in an even direction $`u(x)`$ for the periodic problem on $`(X/2,X/2)`$ is equivalent to energy instability in the direction $`u(x)`$ for the Neumann problem on $`(0,X/2)`$. (If the perturbation $`u(x)`$ is even and has mean value zero on $`(X/2,X/2)`$ then it has mean value zero on $`(0,X/2)`$ as well.) Similarly, ‘periodic’ energy stability in all even directions on $`(X/2,X/2)`$ is equivalent to ‘Neumann’ energy stability in all directions on $`(0,X/2)`$.
In short, for the Neumann problem on $`(0,X/2)`$, the stability result in Theorem 3 still holds, and the instability claims involving $`\pm h_{\mathrm{ss}}^{\prime \prime }`$ in Theorems 1 and 2 also still hold, since these are even functions; the claims involving $`\pm h_{\mathrm{ss}}^{}`$ do not carry over, since those are odd.
## 3. Relative energy levels of periodic, constant and droplet steady states
In this section we investigate the phase space of the power law equation $`h_t=(h^nh_{xxx})_x(h^mh_x)_x`$ by comparing the value of the energy at the positive periodic, constant and zero-angle droplet steady states. Let $`X>0`$ and recall $`q=mn+1`$.
### 3.1. Positive periodic vs. constant steady states
The fluid volume $`_0^Xh(x,t)𝑑x`$ is conserved by the evolution, under periodic boundary conditions, and so the mean value
$$\overline{h}:=\frac{1}{X}_0^Xh𝑑x$$
is constant in time. Suppose the initial data $`h(,0)`$ arises from a small zero-mean perturbation of $`h_{\mathrm{ss}}`$. It is natural to ask whether $`h`$ can converge (while staying positive and smooth) towards the constant steady state $`\overline{h_{\mathrm{ss}}}`$, as $`t\mathrm{}`$. This cannot happen if $`(h_{\mathrm{ss}})<(\overline{h_{\mathrm{ss}}})`$ and also $`h(,0)`$ is close enough to $`h_{\mathrm{ss}}`$ so that $`(h(,0))<(\overline{h_{\mathrm{ss}}})`$, because the energy is dissipated by the evolution.
###### Theorem 6.
Let $`h_{\mathrm{ss}}C^4()`$ be a non-constant positive periodic steady state of the power law equation $`h_t=(h^nh_{xxx})_x(h^mh_x)_x`$, with least period $`X`$. Translate $`h_{\mathrm{ss}}`$ to put its minimum at $`x=0`$ so that $`h_{\mathrm{ss}}`$ rescales to $`k_\alpha `$ for some $`\alpha (0,1)`$, as in §2.3.
If $`q2`$ or $`q<1`$ then $`(h_{\mathrm{ss}})>(\overline{h_{\mathrm{ss}}})`$.
If $`q=1`$ then $`(h_{\mathrm{ss}})=(\overline{h_{\mathrm{ss}}})`$.
If $`1<q<2`$ and $`E^{}>0`$ on $`(\alpha ,1)`$ then $`(h_{\mathrm{ss}})>E(\overline{h_{\mathrm{ss}}})`$.
If $`1<q<2`$ and $`E^{}<0`$ on $`(\alpha ,1)`$ then $`(h_{\mathrm{ss}})<E(\overline{h_{\mathrm{ss}}})`$.
The theorem is proved in Section 5.1. For $`q=3`$ it was observed numerically in \[32, §3\] that $`(h_{\mathrm{ss}})>(\overline{h_{\mathrm{ss}}})`$.
When $`1<q1.75`$, numerical evidence in Figure 3 of suggests $`E^{}(\alpha )<0`$ for all $`\alpha `$. If this is true, then $`(h_{\mathrm{ss}})<(\overline{h_{\mathrm{ss}}})`$ by the theorem and so there can be no heteroclinic connection from $`h_{\mathrm{ss}}`$ to the constant steady state $`\overline{h_{\mathrm{ss}}}`$; indeed in \[19, §4.4\] we find numerically for $`q=1.5`$ that $`h_{\mathrm{ss}}`$ is asymptotically stable, up to translation.
For $`1.795q<2`$, Figure 5 of suggests $`E^{}(\alpha )>0`$ for all $`\alpha `$ and so $`(h_{\mathrm{ss}})>(\overline{h_{\mathrm{ss}}})`$ by the theorem. Thus when $`q1.795`$ or $`q<1`$, the instability result Theorem 2 and the energy level result Theorem 6 together lead us to suspect the existence of a heteroclinic connection from $`h_{\mathrm{ss}}`$ to $`\overline{h_{\mathrm{ss}}}`$.
In Figure 4 we present numerical simulations of such heteroclinic orbits, taken from the companion article \[19, §4.6\]. The top part of the figure presents an orbit connecting the positive periodic steady state to the constant steady state. The bottom part relates to §3.2 and presents a solution connecting a perturbation of the same positive periodic steady state in finite time to a droplet profile (which will not in general be a steady state).
In particular the van der Waals equation (2) has $`q=13+1=3`$, and so by the first claim of Theorem 6, the energy of any positive periodic steady state is greater than that of the constant steady state. This was observed numerically by Witelski and Bernoff \[32, §3\]. See and our companion paper \[19, §4.1\] for numerical simulations of this equation and a discussion of the literature.
For $`q(1.75,1.794]`$, see the remarks around Theorem 9.
For the Cahn–Hilliard equation, the analogue of Theorem 6 (comparing energy levels of non-constant and constant steady states) can be found in Grinfeld and Novick–Cohen’s work \[14, Theorem 4.1\]; further, \[14, §7\] discusses a number of results on existence of heteroclinic connections. See also \[21, §3.4\].
### 3.2. Positive periodic vs. droplet steady states
We do not yet have stability results for droplet steady states. But here we show under certain conditions that the energy of a zero-angle droplet steady state must be lower than that of a positive periodic steady state $`h_{\mathrm{ss}}`$ whose period exceeds the length of the droplet. In our article \[19, §4\] we show numerically that in these cases, the droplet seems to be strongly attracting, with heteroclinic orbits from $`h_{\mathrm{ss}}`$ towards the droplet steady state as shown in the bottom part of Figure 4.
###### Theorem 7.
Let $`h_{\mathrm{ss}}C^4()`$ be a non-constant positive periodic steady state of the power law equation $`h_t=(h^nh_{xxx})_x(h^mh_x)_x`$, with $`h_{\mathrm{ss}}`$ having least period $`X`$ and area $`A_{\mathrm{ss}}=_0^Xh_{\mathrm{ss}}𝑑x`$.
If $`1<q<1`$ or $`2q<3`$, or if $`1<q<2`$ and $`E^{}(\alpha )>0`$ for all $`\alpha (0,1)`$, then there exists a zero contact angle droplet steady state $`\widehat{h}_{\mathrm{ss}}`$ with length $`\widehat{X}<X`$ and area $`\widehat{A}=A_{\mathrm{ss}}`$. Furthermore $`(\widehat{h}_{\mathrm{ss}})<(h_{\mathrm{ss}})`$.
We prove the theorem in Section 5.2.
Remarks
1. Steady states with zero contact angle can occur only for $`q>1`$, by \[17, §2.2\]. Hence we do not consider $`q1`$ in the theorem; indeed variational methods show there can be droplet steady states that do not have acute contact angles .
2. The theorem presumably applies to $`1.795q<2`$, since it seems numerically that $`E^{}>0`$ for those $`q`$-values. For $`q3`$ we think a similar theorem might hold but with $`\widehat{h}_{\mathrm{ss}}`$ being a configuration of disjoint zero–angle droplets.
3. Since $`(h_{\mathrm{ss}})>(\widehat{h}_{\mathrm{ss}})`$, there might be an orbit from $`h_{\mathrm{ss}}`$ to $`\widehat{h}_{\mathrm{ss}}`$. This orbit might describe a positive solution that converges to the nonnegative droplet profile as $`t\mathrm{}`$, or it might describe a positive solution that loses positivity in finite time and then approaches the droplet as a nonnegative weak solution. For $`1<q<1`$ and $`2q<3`$, Theorems 6 and 7 are consistent with the idea that the unstable positive periodic steady state $`h_{\mathrm{ss}}`$ and its stable manifold form a separatrix between the basin of attraction of the constant steady state and the basin of attraction of the droplet steady state. In particular, after perturbing $`h_{\mathrm{ss}}`$ in one direction one seems to find a solution that converges to the constant steady state, while perturbing in the opposite direction often yields a solution that converges to a droplet profile. We present some numerical evidence for such behavior in Figure 4, and discuss this at length in the companion article \[19, §4\]. See also the ‘mountain pass’ remark after Theorem 11.
The preceding theorem and its Remarks are definitely not valid for $`q(1,1.75]`$, since for these $`q`$-values we show there does not even exist a zero contact angle steady state with length less than $`X`$:
###### Theorem 8.
Let $`1<q1.75`$ and suppose $`E^{}(\alpha )<0`$ for all $`\alpha (0,1)`$. Let $`h_{\mathrm{ss}}C^4(𝕋_X)`$ be a non-constant positive periodic steady state of the power law equation $`h_t=(h^nh_{xxx})_x(h^mh_x)_x`$, with $`h_{\mathrm{ss}}`$ having least period $`X`$ and area $`A_{\mathrm{ss}}=_0^Xh_{\mathrm{ss}}𝑑x`$. Assume $`\widehat{h}_{\mathrm{ss}}`$ is nonnegative and piecewise-$`C^1`$ on $`𝕋_X`$, has area $`A_{\mathrm{ss}}`$, and is smooth on the set where it is positive and satisfies there the ‘nonlinear oscillator’ steady state equation (6).
Then $`\widehat{h}_{\mathrm{ss}}`$ is either constant or is a translate of $`h_{\mathrm{ss}}`$, or is a configuration of nonzero contact angle droplet steady states. Specifically, $`\widehat{h}_{\mathrm{ss}}`$ cannot be a zero contact angle droplet steady state.
We prove the theorem in Section 5.3. Note that the hypothesis $`E^{}<0`$ seems to hold for $`1<q1.75`$, by the numerical evidence in Figure 3 of .
Finally, for $`q(1.75,1.794]`$ approximately, we know by the analytical and numerical work in \[17, §5.1\] that there can be two positive periodic steady states with the same period and area. The next theorem shows that the steady state with smaller minimum value (and larger amplitude) is energy stable, and has lower energy than the other one, which is energy unstable. This leads us to suspect there exists an orbit from the energy unstable steady state to the energy stable one, at least when the steady states have been chosen to have their minima at the same location. The bottom plot of Figure 4 presents a numerical simulation of such an orbit, taken from \[19, §4.5\]. The top plot of the figure presents an orbit connecting the unstable positive periodic steady state to the constant steady state, similar to the top plot of Figure 4.
###### Theorem 9.
Assume $`1<q<2`$ and there exists $`\alpha _{crit}(0,1)`$ with $`E^{}(\alpha )<0`$ on $`(0,\alpha _{crit})`$ and $`E^{}(\alpha )>0`$ on $`(\alpha _{crit},1)`$, and assume $`\alpha \alpha P(\alpha )^{2/(q1)}`$ is strictly increasing for $`\alpha (0,1)`$.
Suppose $`h_{\mathrm{ss1}}`$ and $`h_{\mathrm{ss2}}`$ are non-constant positive periodic steady states of the power law equation $`h_t=(h^nh_{xxx})_x(h^mh_x)_x`$, with $`h_{\mathrm{ss1}}`$ and $`h_{\mathrm{ss2}}`$ having the same least period $`X`$ and same area $`_0^Xh_{\mathrm{ss1}}𝑑x=_0^Xh_{\mathrm{ss2}}𝑑x`$.
If $`h_{\mathrm{ss1}}(x)`$ has lower minimum value than $`h_{\mathrm{ss2}}(x)`$, then $`h_{\mathrm{ss1}}`$ is energy stable, $`h_{\mathrm{ss2}}`$ is energy unstable, and $`(h_{\mathrm{ss1}})<(h_{\mathrm{ss2}})`$. Furthermore, $`(h_{\mathrm{ss2}})>(\overline{h_{\mathrm{ss2}}})`$.
We prove this in Section 5.4. The hypothesis about $`E(\alpha )`$ being first strictly decreasing and then strictly increasing is confirmed numerically for $`q`$ in the interval $`(1.750,1.794]`$ by \[17, §6.1\]; see Figure 5. Numerical work also confirms the hypothesis about $`\alpha P(\alpha )^{2/(q1)}`$ being a strictly increasing function of $`\alpha `$ for all $`q(1,2)`$.
Theorem 9 is analogous to \[14, Theorem 4.1(v)\] for the transitional and metastable cases of the Cahn–Hilliard equation, where Grinfeld and Novick–Cohen show the energy of a monotonic ‘interface’ solution is less than that of a monotonic ‘spike’ solution having the same length and area.
### 3.3. Constant steady states: stability results
For any number $`\overline{h}>0`$ the constant function $`h_{\mathrm{ss}}\overline{h}`$ is a steady state of the general evolution equation (1). We now develop analogues for this constant steady state of our earlier stability results. Then in the following subsection we can compare the energies of the constant and droplet steady states, for power law coefficients.
We start by recalling a characterization of linear instability for steady states. From \[18, Appendix A\], an $`X`$-periodic positive steady state $`h_{\mathrm{ss}}C^4()`$ of (1) is linearly unstable with respect to zero–mean perturbations at period $`X`$ if and only if $`\tau _1(h_{\mathrm{ss}})<0`$, where
(11)
$$\tau _1(h_{\mathrm{ss}}):=\underset{u}{\mathrm{min}}\frac{_0^X\left[(u^{})^2r(h_{\mathrm{ss}})u^2\right]𝑑x}{_0^Xu^2𝑑x};$$
the minimum here is taken over nonzero $`uH^1(𝕋_X)`$ with $`_0^Xu𝑑x=0`$. The trial function $`u`$ corresponds to a perturbation of $`h_{\mathrm{ss}}`$.
###### Theorem 10.
Let $`f,gC^2(0,\mathrm{})`$ with $`f>0`$, and write $`r=g/f`$. Let $`\overline{h},X>0`$. Then for the constant steady state $`h_{\mathrm{ss}}\overline{h}`$ of (1), the eigenvalue $`\tau _1`$ in (11) is $`\tau _1(\overline{h})=(2\pi /X)^2r(\overline{h})`$. The $`\tau _1(\overline{h})`$-eigenspace is spanned by $`\mathrm{sin}(2\pi x/X)`$ and $`\mathrm{cos}(2\pi x/X)`$.
Thus with respect to zero–mean perturbations at period $`X`$, the constant steady state is
linearly unstable if $`r(\overline{h})X^2>4\pi ^2,`$
linearly neutrally stable if $`r(\overline{h})X^2=4\pi ^2,`$
linearly asymptotically stable if $`r(\overline{h})X^2<4\pi ^2.`$
(a) If $`r(\overline{h})X^2>4\pi ^2`$, or if $`r(\overline{h})X^2=4\pi ^2`$ and $`r^{\prime \prime }(\overline{h})>0`$, then the constant steady state is energy unstable in the directions $`\pm \mathrm{sin}(2\pi x/X)`$ and $`\pm \mathrm{cos}(2\pi x/X)`$.
(b) If $`r(\overline{h})X^2<4\pi ^2`$, or if $`r(\overline{h})X^2=4\pi ^2`$ and $`r^{\prime \prime }(\overline{h})<0`$, then the constant steady state is energy stable with respect to zero–mean perturbations of period $`X`$. In fact, if $`r(\overline{h})X^2<4\pi ^2`$ then the constant steady state is a strict local minimum of the energy with respect to zero–mean perturbations in $`H^1(𝕋_X)`$, and $`\overline{h}`$ is nonlinearly stable under the evolution (1), in an $`H^1(𝕋_X)`$-sense made precise in the proof.
We prove the theorem in Section 5.5. Its linear stability assertions are well-known and are included for the sake of completeness.
Goldstein, Pesci and Shelley \[12, §IIIB\] used the energy to prove nonlinear instability of the constant steady state for the $`q=1`$ case ($`f(y)=y^n,g(y)=y^n,X=2\pi `$) with either $`2<4`$ or $`=j^2`$ for some integer $`j2`$.
### 3.4. Constant vs. droplet steady states
Consider power law coefficients, so that $`r(y)=y^{q1}`$; then the previous theorem shows that the stability of the constant steady state $`\overline{h}`$ is determined by whether the quantity $`\overline{h}^{q1}X^2`$ is $`>,=,<4\pi ^2`$.
Fix $`X>0`$. Does a zero-angle droplet steady state exist with length at most $`X`$ and with the same area $`\overline{h}X`$ as the constant steady state? If such a droplet steady state exists, can it have lower energy than the constant steady state?
In this direction, in Section 5.6 we prove:
###### Theorem 11.
Let $`\overline{h},X>0`$, and consider the constant steady state $`h_{\mathrm{ss}}\overline{h}`$ of the power law evolution equation $`h_t=(h^nh_{xxx})_x(h^mh_x)_x`$.
(a) Suppose $`1<q<3`$. Then there exists a zero contact angle droplet steady state $`\widehat{h}_{\mathrm{ss}}`$ of length $`\widehat{X}X`$ and area $`\overline{h}X`$ if and only if
(12)
$$\overline{h}^{q1}X^2E(0)=:E_0(q).$$
If such a droplet steady state exists, then $`(\overline{h})>(\widehat{h}_{\mathrm{ss}})`$ if and only if
(13)
$$\overline{h}^{q1}X^2>A(0)^2\left[\frac{3+q}{(3q)(q+1)}\right]^{(3q)/q}=:L(q)\text{(for }1<q<3,q0\text{)}$$
or $`\overline{h}^1X^2>4e^2\pi /3=:L(0)`$ (for $`q=0`$).
(b) For $`q=3`$, such a steady state $`\widehat{h}_{\mathrm{ss}}`$ exists if and only if $`\overline{h}^2X^2=E(0)`$. For $`q>3`$, $`\widehat{h}_{\mathrm{ss}}`$ exists if and only if $`\overline{h}^{q1}X^2E(0)`$. For all $`q3`$, if $`\widehat{h}_{\mathrm{ss}}`$ exists then $`(\overline{h})<(\widehat{h}_{\mathrm{ss}})`$.
To understand conditions (12) and (13) see the plots of $`E_0(q)`$ and $`L(q)`$ in Figures 7 and 7
(constructed using the formulas for $`A(0),P(0),E(0)`$ in \[17, §3.1.2\]). The graphs of $`E_0`$ and $`L`$ intersect at $`q=1,1,3`$ and at $`q^{}1.775`$. For $`1<q<q^{}`$ the figure suggests $`E_0>L`$, and so if the droplet steady state in Theorem 11(a) exists then it certainly has lower energy than the constant steady state. On the other hand, it appears that $`L>E_0`$ when $`1<q<1`$ and when $`q^{}<q<3`$; in these cases the energy condition (13) may or may not hold when the existence condition (12) holds, so that the energy of the droplet steady state (if it exists) might be higher or lower than that of the constant steady state.
The dashed line at height $`4\pi ^2`$ in Figures 7 and 7 intersects $`E_0(q)`$ at $`1`$ and $`q^{}1.768`$, and intersects $`L(q)`$ at $`1`$ and $`1.761`$ (approx.). This line matters because the constant steady state is a strict local minimum for the energy when $`\overline{h}^{q1}X^2<4\pi ^2`$, by Theorem 10(b). For example, suppose that $`2q<3`$ and $`\overline{h}`$ and $`X`$ are such that $`E_0(q)<L(q)<\overline{h}^{q1}X^2<4\pi ^2`$. Then the constant steady state $`\overline{h}`$ is a strict local minimum of the energy but is not a global minimum since $`(\overline{h})>(\widehat{h}_{\mathrm{ss}})`$ by Theorem 11(a). A mountain pass idea then suggests the energy might have a saddle point at which its value is greater than $`(\overline{h})`$. Such a saddle ought to be an energy unstable positive periodic steady state, and should have period $`X`$ and area $`\overline{h}X`$. In fact we already know such a positive periodic steady state exists, by \[18, Theorem 12\] and Theorem 6 of this paper; we illustrate it in Figure 4. Perturbing from the saddle in one direction leads to relaxation to the constant (see the top part of Figure 4), while perturbing in the opposite direction yields apparent relaxation to a droplet (the bottom part of Figure 4). The companion article contains more extensive numerical investigations.
## 4. Proofs of Theorems 15
### 4.1. Proof of Theorem 1
Recall that the positive periodic steady state $`h_{\mathrm{ss}}`$ satisfies equations (4) and (5):
(14)
$$h_{\mathrm{ss}}^{\prime \prime \prime }+r(h_{\mathrm{ss}})h_{\mathrm{ss}}^{}=0,h_{\mathrm{ss}}^{\prime \prime }+H^{}(h_{\mathrm{ss}})=\text{const.},$$
where $`H^{\prime \prime }(y)=r(y)`$ is defined for $`y>0`$.
We compute the first four variations of the energy (3). Let $`uH^1(𝕋_X)`$ have mean value zero. For all $`\epsilon `$ small enough that $`h_{\mathrm{ss}}+\epsilon u>0`$ (so that $`H(h_{\mathrm{ss}}+\epsilon u)`$ makes sense), we have
$`{\displaystyle \frac{d}{d\epsilon }}(h_{\mathrm{ss}}+\epsilon u)|_{\epsilon =0}`$ $`=`$ $`{\displaystyle _0^X}\left[h_{\mathrm{ss}}^{\prime \prime }+H^{}(h_{\mathrm{ss}})\right]u𝑑x=0`$
by (14), since $`u`$ has mean value zero,
$`{\displaystyle \frac{d^2}{d\epsilon ^2}}(h_{\mathrm{ss}}+\epsilon u)|_{\epsilon =0}`$ $`=`$ $`{\displaystyle _0^X}\left[(u^{})^2r(h_{\mathrm{ss}})u^2\right]𝑑x`$
$`=`$ numerator of Rayleigh quotient (11) for $`\tau _1(h_{\mathrm{ss}})`$,
(17) $`{\displaystyle \frac{d^3}{d\epsilon ^3}}(h_{\mathrm{ss}}+\epsilon u)|_{\epsilon =0}`$ $`=`$ $`{\displaystyle _0^X}r^{}(h_{\mathrm{ss}})u^3𝑑x,`$
(18) $`{\displaystyle \frac{d^4}{d\epsilon ^4}}(h_{\mathrm{ss}}+\epsilon u)|_{\epsilon =0}`$ $`=`$ $`{\displaystyle _0^X}r^{\prime \prime }(h_{\mathrm{ss}})u^4𝑑x.`$
If the steady state $`h_{\mathrm{ss}}`$ is linearly unstable with respect to zero–mean perturbations at period $`X`$, then the numerator of the Rayleigh quotient (11) is negative for some zero–mean trial function $`uH^1(𝕋_X)\{0\}`$. We can assume $`uC^{\mathrm{}}(𝕋_X)`$. From (4.1), the second variation of the energy in the direction $`u`$ is negative, so that $`h_{\mathrm{ss}}`$ is energy unstable in the direction $`u`$ as desired.
From Theorems 1 and 3 in , $`h_{\mathrm{ss}}`$ is linearly unstable with respect to zero–mean perturbations at period $`X`$ if it is non-constant and: either the least period of $`h_{\mathrm{ss}}`$ is $`X/j`$ for some integer $`j2`$, or $`r=g/f`$ is convex ($`r^{\prime \prime }0`$) and non-constant on the range of $`h_{\mathrm{ss}}`$. Hence $`h_{\mathrm{ss}}`$ is energy unstable in those situations.
Now assume $`h_{\mathrm{ss}}`$ is non-constant and $`r`$ is strongly convex ($`r^{\prime \prime }>0`$). We consider the second variation of $``$ in the direction $`u=\pm h_{\mathrm{ss}}^{\prime \prime }`$:
(19) $`{\displaystyle \frac{d^2}{d\epsilon ^2}}(h_{\mathrm{ss}}\pm \epsilon h_{\mathrm{ss}}^{\prime \prime })|_{\epsilon =0}`$ $`=`$ $`{\displaystyle _0^X}\left[h_{\mathrm{ss}}^{\prime \prime \prime }{}_{}{}^{\mathrm{\hspace{0.17em}2}}r(h_{\mathrm{ss}})h_{\mathrm{ss}}^{\prime \prime }{}_{}{}^{2}\right]𝑑x`$
$`=`$ $`{\displaystyle _0^X}r(h_{\mathrm{ss}})\left[h_{\mathrm{ss}}^{}h_{\mathrm{ss}}^{\prime \prime \prime }+h_{\mathrm{ss}}^{\prime \prime }{}_{}{}^{2}\right]𝑑x\text{by (}\text{14}\text{)}`$
$`=`$ $`{\displaystyle _0^X}r(h_{\mathrm{ss}})\left[h_{\mathrm{ss}}^{}h_{\mathrm{ss}}^{\prime \prime }\right]^{}𝑑x={\displaystyle _0^X}r^{}(h_{\mathrm{ss}})h_{\mathrm{ss}}^{}{}_{}{}^{2}h_{\mathrm{ss}}^{\prime \prime }𝑑x`$
$`=`$ $`{\displaystyle \frac{1}{3}}{\displaystyle _0^X}r^{}(h_{\mathrm{ss}})\left[h_{\mathrm{ss}}^{}{}_{}{}^{3}\right]^{}𝑑x={\displaystyle \frac{1}{3}}{\displaystyle _0^X}r^{\prime \prime }(h_{\mathrm{ss}})h_{\mathrm{ss}}^{}{}_{}{}^{4}𝑑x.`$
By assumption, $`r^{\prime \prime }(h_{\mathrm{ss}}(x))>0`$ for all $`x`$. Since $`h_{\mathrm{ss}}^{}`$ is not identically zero, the second variation in the direction $`u=\pm h_{\mathrm{ss}}^{\prime \prime }`$ is negative, and so $`\pm h_{\mathrm{ss}}^{\prime \prime }`$ is an energy unstable direction for $`h_{\mathrm{ss}}`$.
It remains to prove $`h_{\mathrm{ss}}`$ is energy unstable in the directions $`u=\pm h_{\mathrm{ss}}^{}`$. Here the second variation is zero, since (4.1) becomes
$$_0^X\left[(h_{\mathrm{ss}}^{\prime \prime })^2r(h_{\mathrm{ss}})(h_{\mathrm{ss}}^{})^2\right]𝑑x=_0^X\left[h_{\mathrm{ss}}^{\prime \prime \prime }+r(h_{\mathrm{ss}})h_{\mathrm{ss}}^{}\right]h_{\mathrm{ss}}^{}𝑑x=0\text{by parts and (}\text{14}\text{)}.$$
The third variation is zero for $`u=\pm h_{\mathrm{ss}}^{}`$ by (17), because $`h_{\mathrm{ss}}`$ is even about its minimum point while $`(\pm h_{\mathrm{ss}}^{})^3`$ is odd, by uniqueness for the nonlinear oscillator equation $`h_{\mathrm{ss}}^{\prime \prime }+H^{}(h_{\mathrm{ss}})=\text{const}`$ in (14). The fourth variation is negative by (18), because $`r^{\prime \prime }(h_{\mathrm{ss}}(x))>0`$ by assumption. Hence $`h_{\mathrm{ss}}`$ is energy unstable in the directions $`u=\pm h_{\mathrm{ss}}^{}`$, completing the proof. ∎
A number of other authors, working on closely related topics, have noted that translation invariance of the evolution implies the second variation of the energy in the direction $`h_{\mathrm{ss}}^{}`$ is zero . Those authors then proved energy instability in the direction $`u=h_{\mathrm{ss}}^{}+\eta `$ for some small function $`\eta `$. Their arguments relied either on $`h_{\mathrm{ss}}`$ having least period $`X/j`$ for some $`j2`$ or else they did not impose a zero–mean requirement on the perturbation. Above, we have used instead the assumption $`r^{\prime \prime }>0`$.
### 4.2. Proof of Theorem 2
When $`q<1`$ or $`q>2`$, Theorem 2 follows from the last statement of Theorem 1 since $`r(y)=y^{q1}`$ is strongly convex.
When $`q=2`$ and $`u=h_{\mathrm{ss}}^{\prime \prime }`$, formula (19) shows the second variation of the energy is zero, since $`r(y)=y`$ and $`r^{\prime \prime }0`$. The third variation is negative because
$`{\displaystyle \frac{d^3}{d\epsilon ^3}}(h_{\mathrm{ss}}\epsilon h_{\mathrm{ss}}^{\prime \prime })|_{\epsilon =0}`$ $`=`$ $`{\displaystyle _0^X}(h_{\mathrm{ss}}^{\prime \prime })^3𝑑x\text{by (}\text{17}\text{)}`$
$`=`$ $`2{\displaystyle _0^X}h_{\mathrm{ss}}^{}h_{\mathrm{ss}}^{\prime \prime }h_{\mathrm{ss}}^{\prime \prime \prime }𝑑x\text{by parts}`$
$`=`$ $`2^2{\displaystyle _0^X}h_{\mathrm{ss}}(h_{\mathrm{ss}}^{})^2h_{\mathrm{ss}}^{\prime \prime }𝑑x\text{since }h_{\mathrm{ss}}^{\prime \prime \prime }=h_{\mathrm{ss}}h_{\mathrm{ss}}^{}\text{ by (}\text{14}\text{)}`$
$`=`$ $`{\displaystyle \frac{2}{3}}^2{\displaystyle _0^X}(h_{\mathrm{ss}}^{})^4𝑑x<0.`$
Thus the steady state is energy unstable in the direction $`u=h_{\mathrm{ss}}^{\prime \prime }`$.
Now suppose $`q>1`$ and $`E^{}(\alpha )>0`$. To obtain the unstable direction $`u`$, start by rescaling $`k_\alpha `$ to give
(20)
$$K_\alpha (x):=\frac{P(\alpha )}{A(\alpha )}k_\alpha \left(P(\alpha )x\right).$$
By construction, $`K_\alpha `$ has period $`1`$ and mean value $`1`$. Then define
$$\kappa _\alpha (x):=\frac{}{\alpha }K_\alpha (x);$$
$`\kappa _\alpha `$ is well-defined and smooth because $`P`$ and $`A`$ depend smoothly on $`\alpha `$ while $`k_\alpha (x)`$ is jointly smooth in $`(\alpha ,x)`$. Notice $`\kappa _\alpha `$ is even in $`x`$, has period $`1`$, and has mean value zero:
$$_0^1\kappa _\alpha (x)𝑑x=\frac{}{\alpha }_0^1K_\alpha (x)𝑑x=\frac{}{\alpha }(1)=0.$$
(See \[18, §5.4\] for more properties of $`\kappa _\alpha `$.) Let $`u(x)=\pm \kappa _\alpha (x/X)`$.
We want to show $`u`$ is an unstable direction. First we show that composing the rescalings of $`h_{\mathrm{ss}}`$ to $`k_\alpha `$ and then $`k_\alpha `$ to $`K_\alpha `$ yields
(21)
$$X^2h_{\mathrm{ss}}(xX)^{q1}=EK_\alpha (x)^{q1}.$$
Indeed, from the definition (20) of $`K_\alpha `$, the righthand side of (21) reduces to $`P^2k_\alpha (Px)^{q1}`$, and then one can substitute for $`k_\alpha `$ in terms of $`h_{\mathrm{ss}}`$ using (7). Next one obtains the lefthand side of (21) by using the relation $`P=(/D)^{1/2q}D^{1/2}X`$ which relates the periods of $`k_\alpha `$ and $`h_{\mathrm{ss}}`$ ($`P`$ and $`X`$, respectively).
Now that we have (21), use (4.1) to compute the second variation of the energy in the direction $`u(x)=\pm \kappa _\alpha (x/X)`$ as
$`{\displaystyle _0^X}\left[\kappa _\alpha ^{}(x/X)^2/X^2h_{\mathrm{ss}}(x)^{q1}\kappa _\alpha (x/X)^2\right]𝑑x`$ $`=`$ $`{\displaystyle _0^1}\left[\kappa _\alpha ^{}(x)^2X^2h_{\mathrm{ss}}(xX)^{q1}\kappa _\alpha (x)^2\right]𝑑x/X`$
$`=`$ $`{\displaystyle _0^1}\left[(\kappa _\alpha ^{})^2EK_\alpha ^{q1}\kappa _\alpha ^2\right]𝑑x/X\text{by (}\text{21}\text{)}`$
$`<`$ $`0`$
by the proof of \[18, Prop. 14\], which uses both $`q>1`$ and $`E^{}(\alpha )>0`$. Thus the steady state is energy unstable in the direction $`u=\pm \kappa _\alpha (x/X)`$. ∎
### 4.3. Proof of Theorem 3
Since $`E^{}(\alpha )<0`$ by hypothesis, Proposition 15 of implies $`\mu _1(\alpha )0`$, where
$$\mu _1(\alpha ):=\mathrm{min}\left\{\frac{_0^1\left[(v^{})^2E(\alpha )K_\alpha ^{q1}v^2\right]𝑑x}{_0^1v^2𝑑x}:vH^1(𝕋_1)\{0\},_0^1v(x)𝑑x=0\right\}.$$
(In fact $`E^{}(\alpha )0`$ would suffice to get this.) Notice $`\mu _1(\alpha )=X^2\tau _1(h_{\mathrm{ss}})`$, by letting $`v(x)=u(xX)`$ and using the identity (21) and the definition (11) of $`\tau _1`$. Hence $`\tau _1(h_{\mathrm{ss}})0`$.
Consider $`uH^1(𝕋_X)\{0\}`$ with mean value zero. The first variation of $``$ in the direction $`u`$ is zero by (4.1), and the second variation of $``$ in (4.1) is nonnegative because it equals the numerator of the Rayleigh quotient for $`\tau _1(h_{\mathrm{ss}})`$.
If the second variation is positive then $`h_{\mathrm{ss}}`$ is energy stable in the direction $`h_{\mathrm{ss}}`$ and we are done. If the second variation is zero then the Rayleigh quotient of $`u`$ is zero, and so $`\tau _1(h_{\mathrm{ss}})=0`$ and $`u`$ minimizes the Rayleigh quotient for $`\tau _1`$ in (11). Hence $`u`$ satisfies the Euler–Lagrange condition $`u^{\prime \prime }+r(h_{\mathrm{ss}})u=\text{const}`$, and so it satisfies $`u=0`$ where $``$ is the linearized operator defined in \[18, eq. (4)\] (take $`a=0`$ there). Theorem 10(a) in and the hypothesis $`E^{}(\alpha )<0`$ now imply $`u`$ is a multiple of $`h_{\mathrm{ss}}^{}`$, and so $`u`$ is odd. Therefore the third variation (17) of $``$ in the direction $`u`$ is zero, because $`u`$ is odd and $`h_{\mathrm{ss}}`$ is even. The fourth variation is positive by (18) because $`r(y)=y^{q1}`$ and $`r^{\prime \prime }(y)<0`$ for $`1<q<2`$. Thus $`h_{\mathrm{ss}}`$ is energy stable in the direction $`u`$. ∎
### 4.4. Proof of Theorem 5
As usual, the first variation of $``$ at $`h_{\mathrm{ss}}`$ is zero by (4.1). We now prove non-negativity of the second variation, given by (4.1).
First note that $`h_{\mathrm{ss}}`$ is symmetric about every point at which $`h_{\mathrm{ss}}^{}=0`$, by uniqueness for the ODE $`h_{\mathrm{ss}}^{\prime \prime }+H^{}(h_{\mathrm{ss}})=\text{const}`$ (see (14)); the uniqueness theorem applies here since the coefficient function $`H^{}`$ is $`C^1`$ (even $`C^3`$) on the range of the positive bounded function $`h_{\mathrm{ss}}`$.
Since $`h_{\mathrm{ss}}`$ has a minimum at $`x=0`$ by hypothesis, we conclude $`h_{\mathrm{ss}}`$ is even and $`h_{\mathrm{ss}}^{}>0`$ on $`(0,X/2)`$ (otherwise $`h_{\mathrm{ss}}`$ would have period less than $`X`$).
Consider a minimizer of the Rayleigh quotient (11) with respect to odd functions $`u`$; it is a smooth odd function $`\stackrel{~}{u}`$ satisfying $`\stackrel{~}{u}^{\prime \prime }+r(h_{\mathrm{ss}})\stackrel{~}{u}+\stackrel{~}{\tau }\stackrel{~}{u}=0`$ for some constant $`\stackrel{~}{\tau }`$. Since $`\stackrel{~}{u}(0)=0`$, one must have $`\stackrel{~}{u}^{}(0)0`$ because otherwise $`\stackrel{~}{u}0`$ by the uniqueness theorem for linear ODEs. Also, $`\stackrel{~}{u}(X/2)=0`$ by the oddness and periodicity of $`\stackrel{~}{u}`$, and so there is a point $`b(0,X/2]`$ with $`\stackrel{~}{u}(b)=0`$ and $`\stackrel{~}{u}0`$ between $`0`$ and $`b`$. Assume $`\stackrel{~}{u}>0`$ between $`0`$ and $`b`$ (otherwise consider $`\stackrel{~}{u}`$). Then
$`\stackrel{~}{\tau }{\displaystyle _0^b}\stackrel{~}{u}h_{\mathrm{ss}}^{}𝑑x`$ $`=`$ $`{\displaystyle _0^b}\left[\stackrel{~}{u}^{\prime \prime }+r(h_{\mathrm{ss}})\stackrel{~}{u}\right]h_{\mathrm{ss}}^{}𝑑x\text{since }\stackrel{~}{u}^{\prime \prime }+r(h_{\mathrm{ss}})\stackrel{~}{u}+\stackrel{~}{\tau }\stackrel{~}{u}=0`$
$`=`$ $`\stackrel{~}{u}^{}(b)h_{\mathrm{ss}}^{}(b){\displaystyle _0^b}\left[h_{\mathrm{ss}}^{\prime \prime \prime }+r(h_{\mathrm{ss}})h_{\mathrm{ss}}^{}\right]\stackrel{~}{u}𝑑x\text{by parts, since }h_{\mathrm{ss}}^{}(0)=0`$
$`=`$ $`\stackrel{~}{u}^{}(b)h_{\mathrm{ss}}^{}(b)\text{by (}\text{14}\text{)}`$
$``$ $`0,`$
because $`\stackrel{~}{u}^{}(b)0`$ and $`h_{\mathrm{ss}}^{}(b)0`$. Since $`\stackrel{~}{u}`$ and $`h_{\mathrm{ss}}^{}`$ are positive on $`(0,b)`$ it follows that $`\stackrel{~}{\tau }0`$. Hence the second variation of $``$ in (4.1) is nonnegative, as desired.
The third variation at $`\epsilon =0`$ is zero by (17), since $`u`$ is odd and $`h_{\mathrm{ss}}`$ is even.
The fourth variation (18) is positive not just at $`\epsilon =0`$ but at all $`\epsilon (0,1)`$:
$$\frac{d^4}{d\epsilon ^4}(h_{\mathrm{ss}}+\epsilon u)=_0^Xr^{\prime \prime }(h_{\mathrm{ss}}+\epsilon u)u^4𝑑x>0$$
by the strong concavity of $`r`$, provided $`u0`$. Taylor’s theorem completes the proof, since for some $`\stackrel{~}{\epsilon }(0,1)`$,
$$(h_{\mathrm{ss}}+u)=(h_{\mathrm{ss}})+\frac{1}{2!}\frac{d^2}{d\epsilon ^2}(h_{\mathrm{ss}}+\epsilon u)|_{\epsilon =0}+\frac{1}{4!}\frac{d^4}{d\epsilon ^4}(h_{\mathrm{ss}}+\epsilon u)|_{\epsilon =\stackrel{~}{\epsilon }}>(h_{\mathrm{ss}}).$$
## 5. Proofs of Theorems 611
### 5.1. Proof of Theorem 6
We start by relating $`(h_{\mathrm{ss}})`$ to $`(k_\alpha )`$. Take
(22)
$$H(y):=\{\begin{array}{cc}\hfill \frac{1}{q}\left[\frac{y^{q+1}}{q+1}y\right],& q0,1,\hfill \\ \hfill y\mathrm{log}yy,& q=0,\hfill \\ \hfill y\mathrm{log}y,& q=1,\hfill \end{array}$$
so that $`H^{\prime \prime }(y)=y^{q1}`$, and recall from the definition (3) of the energy that
(23)
$$(h_{\mathrm{ss}})=_0^X\left[\frac{1}{2}(h_{\mathrm{ss}}^{})^2H(h_{\mathrm{ss}})\right]𝑑x\text{and}(k_\alpha )=_0^{P(\alpha )}\left[\frac{1}{2}(k_\alpha ^{})^2H(k_\alpha )\right]𝑑x.$$
Denote the period of $`h_{\mathrm{ss}}`$ by $`P_{\mathrm{ss}}=X`$, and the area by $`A_{\mathrm{ss}}=_0^Xh_{\mathrm{ss}}𝑑x`$. Writing $`P=P(\alpha )`$ and $`A=A(\alpha )`$, the rescaling (7) implies
(24)
$$P=\{\begin{array}{cc}\left(\frac{}{D}\right)^{1/2q}D^{1/2}P_{\mathrm{ss}},& q0,\hfill \\ e^{D/2}^{1/2}P_{\mathrm{ss}},& q=0,\hfill \end{array}\text{and}A=\{\begin{array}{cc}\left(\frac{}{D}\right)^{3/2q}D^{1/2}A_{\mathrm{ss}},& q0,\hfill \\ e^{3D/2}^{1/2}A_{\mathrm{ss}},& q=0.\hfill \end{array}$$
Notice that the rescaling (7) can be written as
(25)
$$h_{\mathrm{ss}}(x)=\frac{A_{\mathrm{ss}}}{A}\frac{P}{P_{\mathrm{ss}}}k_\alpha \left(\frac{P}{P_{\mathrm{ss}}}x\right).$$
From (24) we obtain the invariance relation
(26)
$$P_{\mathrm{ss}}^{3q}A_{\mathrm{ss}}^{q1}=P^{3q}A^{q1}=E(\alpha ),$$
and this implies
(27)
$$\left(\frac{A_{\mathrm{ss}}}{A}\right)^2\left(\frac{P}{P_{\mathrm{ss}}}\right)^3=\left(\frac{A_{\mathrm{ss}}}{A}\right)^{q+1}\left(\frac{P}{P_{\mathrm{ss}}}\right)^q.$$
Using (25), (26), (27) and the definitions (23), we at last deduce a relation between $`(h_{\mathrm{ss}})`$ and $`(k_\alpha )`$:
(28)
$$^1\overline{h_{\mathrm{ss}}}^{(q+1)}\frac{(h_{\mathrm{ss}})(\overline{h_{\mathrm{ss}}})}{P_{\mathrm{ss}}}=\overline{k_\alpha }^{(q+1)}\frac{(k_\alpha )(\overline{k_\alpha })}{P(\alpha )},$$
where the mean values are $`\overline{h_{\mathrm{ss}}}:=A_{\mathrm{ss}}/P_{\mathrm{ss}}`$ and $`\overline{k_\alpha }:=A(\alpha )/P(\alpha )`$. (When checking (28), one can omit the linear terms in $`H(y)`$ from the calculations, since $`h_{\mathrm{ss}}`$ and $`\overline{h_{\mathrm{ss}}}`$ have the same mean value, as do $`k_\alpha `$ and $`\overline{k_\alpha }`$.)
In view of (28), then, Theorem 6 follows from:
###### Proposition 12.
Fix $`\alpha _1(0,1)`$.
If $`q2`$ or $`q<1`$ then $`(k_\alpha )>(\overline{k_\alpha })`$ for all $`\alpha (0,1)`$.
If $`q=1`$ then $`(k_\alpha )=(\overline{k_\alpha })`$ for all $`\alpha (0,1)`$.
If $`1<q<2`$ and $`E^{}(\alpha )>0`$ $`\alpha (\alpha _1,1)`$ then $`(k_\alpha )>(\overline{k_\alpha })`$ $`\alpha [\alpha _1,1)`$.
If $`1<q<2`$ and $`E^{}(\alpha )<0`$ $`\alpha (\alpha _1,1)`$ then $`(k_\alpha )<(\overline{k_\alpha })`$ $`\alpha [\alpha _1,1)`$.
Note that if $`q2`$ or $`q<1`$ then $`E^{}(\alpha )>0`$ for all $`\alpha `$ by \[18, Theorem 11\].
###### Proof of Proposition 12.
The proof depends on a number of elementary differential equations and inequalities that we derived in §§5.1–5.2 of , and the reader may wish to skim those sections before proceeding.
If $`q=1`$ then $`(k_\alpha )=(\overline{k_\alpha })`$, as one sees directly from the formula in (23), using that $`k_\alpha (x)=1+(\alpha 1)\mathrm{cos}x,P(\alpha )=2\pi `$ and $`\overline{k_\alpha }=1`$. So we assume $`q1`$ from now on, and $`\alpha (0,1)`$.
By definition,
(29) $`{\displaystyle \frac{(k_\alpha )(\overline{k_\alpha })}{P}}`$ $`=`$ $`{\displaystyle \frac{1}{P}}{\displaystyle _0^P}\left[{\displaystyle \frac{1}{2}}k_{\alpha }^{}{}_{}{}^{2}H(k_\alpha )\right]𝑑x+H(A/P)`$
$`=`$ $`{\displaystyle \frac{1}{P}}{\displaystyle _0^P}k_{\alpha }^{}{}_{}{}^{2}𝑑xH(\alpha )+H(A/P)\text{by }\text{[18, eq. (21)]}\text{.}`$
First assume $`q1`$; then
(30)
$$\frac{d}{d\alpha }\frac{(k_\alpha )(\overline{k_\alpha })}{P}=\frac{P^{}(\alpha )}{P^2}_0^Pk_{\alpha }^{}{}_{}{}^{2}𝑑x+\left(1+(q+1)H(A/P)\left(\frac{A}{P}\right)^1\right)\left(\frac{A}{P}\right)^{}(\alpha ),$$
where we have used \[18, eq. (35)\] and the identity $`H^{}(y)=1+(q+1)H(y)y^1`$ (valid for $`q1`$). Differentiating the function
(31)
$$(\alpha ):=\left(\frac{A}{P}\right)^{(q+1)}\frac{(k_\alpha )(\overline{k_\alpha })}{P}$$
(which is inspired by (28)) with respect to $`\alpha `$, we find from (29) and (30) that
$$^{}(\alpha )=\left(\frac{A}{P}\right)^{(q+2)}\left\{\frac{1}{P}_0^Pk_{\alpha }^{}{}_{}{}^{2}𝑑x\left[(q+1)\left(\frac{A}{P}\right)^{}\frac{AP^{}}{P^2}\right]+\left[\frac{A}{P}+(q+1)H(\alpha )\right]\left(\frac{A}{P}\right)^{}\right\}.$$
Substituting
$$\frac{A}{P}+(q+1)H(\alpha )=\frac{q+3}{2}\frac{1}{P}_0^Pk_{\alpha }^{}{}_{}{}^{2}𝑑x$$
from \[18, eqs. (31–32)\] yields
(32)
$$^{}(\alpha )=\frac{1}{2}\left(\frac{A}{P}\right)^{(q+2)}\frac{1}{P}_0^Pk_{\alpha }^{}{}_{}{}^{2}𝑑x\left[(q1)\left(\frac{A}{P}\right)^{}+2AP^2P^{}\right]\text{for }q1.$$
For $`q=1`$ we obtain exactly the same formula (32) for $`^{}(\alpha )`$, as follows:
$`^{}(\alpha )`$ $`=`$ $`{\displaystyle \frac{d}{d\alpha }}{\displaystyle \frac{(k_\alpha )(\overline{k_\alpha })}{P}}`$
$`=`$ $`{\displaystyle \frac{P^{}}{P^2}}{\displaystyle _0^P}k_{\alpha }^{}{}_{}{}^{2}𝑑x+\left(1P/A\right)\left({\displaystyle \frac{A}{P}}\right)^{}\text{from (}\text{29}\text{) and }\text{[18, eq. (35)]}`$
$`=`$ $`{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{A}{P}}\right)^1{\displaystyle \frac{1}{P}}{\displaystyle _0^P}k_{\alpha }^{}{}_{}{}^{2}𝑑x\left[2\left({\displaystyle \frac{A}{P}}\right)^{}+2AP^2P^{}\right]`$
using that $`(1P/A)=_0^P(k_\alpha ^{})^2𝑑x/A`$ when $`q=1`$, by \[18, eq. (28)\]. The last equation is (32) for $`q=1`$.
Equation (32) simplifies to
(33)
$$^{}(\alpha )=\left[\frac{1}{2}P^3\left(\frac{A}{P}\right)^{2q}_0^Pk_{\alpha }^{}{}_{}{}^{2}𝑑x\right]E^{}(\alpha ),$$
using that $`E=P^2(A/P)^{q1}`$ by (26). Hence
(34)
$$E^{}>0^{}<0\text{and}E^{}<0^{}>0.$$
Also
(35)
$$(\alpha )0\text{as}\alpha 1,$$
by the formula (29) together with the facts that $`P,A2\pi `$ and $`P/A1`$ as $`\alpha 1`$ (see \[18, Lemma 6\]) and that $`k_\alpha ^{}0`$ uniformly as $`\alpha 1`$, by \[18, eq. (21)\]. Proposition 12 now follows from (34), (35) and \[18, Theorem 11\] (which shows $`E^{}>0`$ when $`q2`$ or $`q<1`$). For example, if $`E^{}>0`$ on $`(\alpha _1,1)`$ then $`^{}<0`$ on $`(\alpha _1,1)`$; since $`(1)=0`$ we deduce $`>0`$ on $`[\alpha _1,1)`$, and so $`(k_\alpha )(\overline{k_\alpha })>0`$ for $`\alpha [\alpha _1,1)`$. ∎
### 5.2. Proof of Theorem 7
The proof relies on formulas for $`A^{}(\alpha )`$ and $`E^{}(\alpha )`$ that were derived in Lemmas 16 and 18 of : for all $`q1`$,
(36)
$$A^{}=(q+1)H(\alpha )P^{}\frac{q1}{2}H^{}(\alpha )P,$$
(37)
$$E^{}=\left(\frac{A}{P}\right)^{q2}\left\{P^{}\left[(q3)A+(q1)(q+1)H(\alpha )P\right]+\frac{1}{2}(q1)^2H^{}(\alpha )P^2\right\}.$$
Note that $`E^{}>0`$ if $`q<1`$ or $`q2`$, by \[18, Theorem 11\]. Our assumptions therefore imply $`1<q<3`$ and $`E^{}(\alpha )>0`$ for all $`\alpha (0,1)`$. Also $`E(0)>0`$ because $`q>1`$ \[17, §3.1.2\]. If we define
$$\widehat{X}=\left[E(0)/(A_{\mathrm{ss}}^{q1})\right]^{1/(3q)}$$
then
$$\widehat{X}^{3q}A_{\mathrm{ss}}^{q1}=E(0)<E(\alpha )=X^{3q}A_{\mathrm{ss}}^{q1}$$
by (26); here the value $`\alpha (0,1)`$ is determined by translating and rescaling $`h_{\mathrm{ss}}`$ to a particular $`k_\alpha `$, as in §2.3. Hence $`0<\widehat{X}<X`$ (using that $`q<3`$).
By rescaling the zero contact angle function $`k_0`$ on the interval $`[0,P(0)]`$ (as in §2.3; see \[17, Claim 5.1.2\] for details) we obtain a zero contact angle droplet steady state $`\widehat{h}_{\mathrm{ss}}`$ of $`h_t=(h^nh_{xxx})_x(h^mh_x)_x`$, with length $`\widehat{X}`$ and area $`A_{\mathrm{ss}}`$ as desired.
It remains to prove that the energy of this droplet steady state is lower than the energy of the positive periodic steady state $`h_{\mathrm{ss}}`$. That is, we want to prove
$$_0^X\left[\frac{1}{2}(\widehat{h}_{\mathrm{ss}}^{})^2G(\widehat{h}_{\mathrm{ss}})\right]𝑑x<_0^X\left[\frac{1}{2}(h_{\mathrm{ss}}^{})^2G(h_{\mathrm{ss}})\right]𝑑x$$
where
$$G(y)=\{\begin{array}{cc}\hfill \frac{y^{q+1}}{q(q+1)},& q0,1,\hfill \\ \hfill y\mathrm{log}yy,& q=0,\hfill \\ \hfill \mathrm{log}y,& q=1;\hfill \end{array}$$
note that we can use $`G`$ instead of $`H`$ in the energy because they differ only by a linear function (cf. the definition (22) of $`H`$) and $`\widehat{h}_{\mathrm{ss}}`$ and $`h_{\mathrm{ss}}`$ have the same area, $`A_{\mathrm{ss}}`$.
Since $`\widehat{h}_{\mathrm{ss}}`$ is supported on $`(0,\widehat{X})`$ and because $`G(0)=0`$, the desired inequality is
(38)
$$_0^{\widehat{X}}\left[\frac{1}{2}(\widehat{h}_{\mathrm{ss}}^{})^2G(\widehat{h}_{\mathrm{ss}})\right]𝑑x<_0^X\left[\frac{1}{2}(h_{\mathrm{ss}}^{})^2G(h_{\mathrm{ss}})\right]𝑑x.$$
Next rescale $`\widehat{h}_{\mathrm{ss}}`$ to $`k_0`$ and $`h_{\mathrm{ss}}`$ to $`k_\alpha `$: from (25) and (27) (with $`P_{\mathrm{ss}}`$ replaced by $`\widehat{X}`$ or $`X`$ as appropriate) we deduce (38) is equivalent to
$$\left[\frac{A_{\mathrm{ss}}}{A(0)}\right]^2\left[\frac{P(0)}{\widehat{X}}\right]^3_0^{P(0)}\left[\frac{1}{2}(k_0^{})^2G(k_0)\right]𝑑x<\left[\frac{A_{\mathrm{ss}}}{A(\alpha )}\right]^2\left[\frac{P(\alpha )}{X}\right]^3_0^{P(\alpha )}\left[\frac{1}{2}(k_\alpha ^{})^2G(k_\alpha )\right]𝑑x,$$
except that when $`q=0`$ we have to subtract
$$\left[\mathrm{log}\frac{A_{\mathrm{ss}}}{A(0)}\frac{P(0)}{\widehat{X}}\right]k_0\text{and}\left[\mathrm{log}\frac{A_{\mathrm{ss}}}{A(\alpha )}\frac{P(\alpha )}{X}\right]k_\alpha $$
from the integrands on the left and right sides, respectively. By substituting the relations
$$\widehat{X}=\left[E(0)/(A_{\mathrm{ss}}^{q1})\right]^{1/(3q)}\text{and}X=\left[E(\alpha )/(A_{\mathrm{ss}}^{q1})\right]^{1/(3q)}$$
into the last inequality and using the definition $`E(\alpha )=P(\alpha )^{3q}A(\alpha )^{q1}`$, we see the desired inequality reduces to $`𝒢(0)<𝒢(\alpha )`$ where (for $`q3`$)
(39)
$$𝒢(\alpha )=A(\alpha )^{(q+3)/(q3)}_0^{P(\alpha )}\left[\frac{1}{2}(k_\alpha ^{})^2G(k_\alpha )\right]𝑑x\left(+\frac{2}{3}\mathrm{log}A(\alpha )\text{, when }q=0\right).$$
Thus to show the energy of $`\widehat{h}_{\mathrm{ss}}`$ is lower than that of $`h_{\mathrm{ss}}`$, it suffices to show $`𝒢^{}(\alpha )>0`$ for all $`\alpha (0,1)`$, assuming $`1<q<3`$ and $`E^{}(\alpha )>0`$ for all $`\alpha (0,1)`$.
To show $`𝒢^{}>0`$, we substitute \[18, eq. (32)\] and \[18, eq. (30)\] into the definition (39) of $`𝒢`$, obtaining that
$$𝒢=\frac{q3}{q(q+3)}A^{2q/(q3)}+\frac{q1}{q+3}H(\alpha )PA^{(q+3)/(q3)}\text{when }q3,0,1,3\text{.}$$
After differentiating this formula with respect to $`\alpha `$ and then substituting for $`A^{}(\alpha )`$ from (36), we simplify with the help of (37) to obtain
(40)
$$𝒢^{}=\frac{1}{q3}H(\alpha )A^{6/(q3)}\left(\frac{A}{P}\right)^{q+2}E^{}\text{when }q3,0,1,3\text{.}$$
When $`q=0`$ we obtain $`𝒢=\frac{2}{3}\frac{1}{3}HPA^1+\frac{2}{3}\mathrm{log}A`$, by putting the $`q=0`$ versions of \[18, eq. (32)\] and \[18, eq. (30)\] into (39). Hence $`𝒢^{}=\frac{1}{3}HP^2E^{}`$, by differentiating and using (36) and (37). That is, (40) holds when $`q=0`$, also.
We conclude from (40) that $`𝒢^{}>0`$ as desired (provided $`1<q<3`$ and $`E^{}(\alpha )>0`$ for all $`\alpha (0,1)`$), since for $`q>1`$ we have $`H(\alpha )<0`$ by the definition (22). ∎
### 5.3. Proof of Theorem 8
The proof involves rescaling arguments rather than the energy. Write $`P_{\mathrm{ss}}=X`$. Assume $`\widehat{h}_{\mathrm{ss}}`$ is non-constant.
Suppose that in fact $`\widehat{h}_{\mathrm{ss}}`$ is positive (and so smooth by hypothesis). Then the least period of $`\widehat{h}_{\mathrm{ss}}`$ equals $`P_{\mathrm{ss}}/j`$ for some positive integer $`j`$, with the area per period equaling $`A_{\mathrm{ss}}/j`$. If $`j=1`$ then $`\widehat{h}_{\mathrm{ss}}`$ must be a translate of $`h_{\mathrm{ss}}`$, by modifying slightly the uniqueness remarks in \[18, §6.2\] (using the assumption that $`E^{}<0`$ to get strict monotonicity of $`E`$).
Thus we can assume $`j2`$. By rescaling $`h_{\mathrm{ss}}`$ and $`\widehat{h}_{\mathrm{ss}}`$ to $`k_\alpha `$ and $`k_{\widehat{\alpha }}`$ for some $`\alpha ,\widehat{\alpha }(0,1)`$, as in (7), we get from (26) that
$$E(\alpha )=P_{\mathrm{ss}}^{3q}A_{\mathrm{ss}}^{q1}\text{and}E(\widehat{\alpha })=(P_{\mathrm{ss}}/j)^{3q}(A_{\mathrm{ss}}/j)^{q1}.$$
Hence
$`4j^2={\displaystyle \frac{E(\alpha )}{E(\widehat{\alpha })}}`$ $`<`$ $`{\displaystyle \frac{E(0)}{E(1)}}\text{since }E^{}<0`$
$`=`$ $`{\displaystyle \frac{1}{4\pi ^2}}{\displaystyle \frac{2}{q}}(1+q)B({\displaystyle \frac{1}{2q}},{\displaystyle \frac{1}{2}})^{3q}B({\displaystyle \frac{3}{2q}},{\displaystyle \frac{1}{2}})^{q1}=:J(q)\text{say,}`$
by the formula for $`E(0)`$ in \[17, eq. (3.13)\] and since $`E(1)=P(1)^{3q}A(1)^{q1}=4\pi ^2`$ by \[18, Lemma 6\]. We will obtain a contradiction by showing $`J(q)<4`$, when $`1<q1.75`$; this will show $`\widehat{h}_{\mathrm{ss}}`$ is not positive.
For $`1<q1.5`$ we have
$$J(q)\frac{1}{4\pi ^2}\frac{2}{1}(1+1.5)B(\frac{1}{21.5},\frac{1}{2})^{31}B(\frac{3}{21.5},\frac{1}{2})^{1.51}3.17<4,$$
where we have used that the Beta function $`B(a,b)=_0^1t^{a1}(1t)^{b1}𝑑t`$ is decreasing in its arguments, and is bigger than $`1`$ when those arguments are less than $`1`$. For $`1.5<q1.75`$ we similarly have
$$J(q)\frac{1}{4\pi ^2}\frac{2}{1.5}(1+1.75)B(\frac{1}{21.75},\frac{1}{2})^{31.5}B(\frac{3}{21.75},\frac{1}{2})^{1.751}1.73<4,$$
completing the contradiction. (Actually the best upper bound for $`J(q)`$ seems numerically to be about $`1.04`$, when $`1<q1.75`$; see \[18, Figure 6\].)
The above contradiction implies $`\widehat{h}_{\mathrm{ss}}`$ is not positive everywhere. Consider therefore one component of the set $`\{x:\widehat{h}_{\mathrm{ss}}(x)>0\}`$, say an interval with length $`\widehat{P}P_{\mathrm{ss}}`$. Write $`\widehat{A}A_{\mathrm{ss}}`$ for the area under $`\widehat{h}_{\mathrm{ss}}`$ on this interval. Note the contact angles of $`\widehat{h}_{\mathrm{ss}}`$ must be the same at the two endpoints of the interval, as a consequence of the nonlinear oscillator equation (6) (see for example \[17, §2.2\]). Suppose these contact angles are zero, so that $`\widehat{h}_{\mathrm{ss}}`$ rescales to $`k_0`$ on the interval, using (7). Then
$$E(0)=\widehat{P}^{3q}\widehat{A}^{q1}P_{\mathrm{ss}}^{3q}A_{\mathrm{ss}}^{q1}=E(\alpha )$$
by applying (26) twice, but this contradicts our assumption that $`E^{}<0`$. Thus the contact angles of $`\widehat{h}_{\mathrm{ss}}`$ must all be nonzero, as desired. ∎
### 5.4. Proof of Theorem 9
Translate $`h_{\mathrm{ss1}}`$ and $`h_{\mathrm{ss2}}`$ so that they attain their minimum values at $`x=0`$, and then rescale as in §2.3 to obtain $`k_{\alpha _1}`$ and $`k_{\alpha _2}`$ respectively. Since $`h_{\mathrm{ss1}}`$ and $`h_{\mathrm{ss2}}`$ have the same period and area, for which we write $`P_{\mathrm{ss}}=X`$ and $`A_{\mathrm{ss}}`$ respectively, it follows from (26) that $`E(\alpha _1)=E(\alpha _2)`$. Notice $`h_{\mathrm{ss1}}h_{\mathrm{ss2}}\alpha _1\alpha _2`$, in view of the expression (25) for $`h_{\mathrm{ss}}`$ in terms of $`k_\alpha `$ and $`P_{\mathrm{ss}},A_{\mathrm{ss}},P(\alpha ),A(\alpha )`$. Since $`E(\alpha _1)=E(\alpha _2)`$ while $`E`$ is strictly decreasing on $`(0,\alpha _{crit})`$ and strictly increasing on $`(\alpha _{crit},1)`$, we conclude $`\alpha _{crit}`$ must lie between $`\alpha _1`$ and $`\alpha _2`$.
We show $`\alpha _1<\alpha _2`$. The hypothesis $`\mathrm{min}h_{\mathrm{ss1}}<\mathrm{min}h_{\mathrm{ss2}}`$ gives
$$\alpha _1\left(\frac{D_1}{}\right)^{1/q}<\alpha _2\left(\frac{D_2}{}\right)^{1/q}$$
by the rescaling (7). Next apply the first equation in (24) to solve for $`D_1`$ in terms of $`P(\alpha _1),P_{\mathrm{ss}}`$ and $``$, and similarly solve for $`D_2`$ in terms of $`P(\alpha _2),P_{\mathrm{ss}}`$ and $``$. Substituting into the above inequality gives $`\alpha _1P(\alpha _1)^{2/(q1)}<\alpha _2P(\alpha _2)^{2/(q1)}`$. The strict increase of $`\alpha \alpha P(\alpha )^{2/(q1)}`$ implies $`\alpha _1<\alpha _2`$.
Since $`\alpha _1<\alpha _{crit}<\alpha _2`$, our hypothesis on $`E^{}`$ implies $`E^{}(\alpha _1)<0`$ and $`E^{}(\alpha _2)>0`$. Theorem 3 then implies that $`h_{\mathrm{ss1}}`$ is energy stable, and Theorem 2 implies $`h_{\mathrm{ss2}}`$ is energy unstable.
Next we show $`(h_{\mathrm{ss1}})<(h_{\mathrm{ss2}})`$, or $`(h_{\mathrm{ss1}})(\overline{h_{\mathrm{ss1}}})<(h_{\mathrm{ss2}})(\overline{h_{\mathrm{ss2}}})`$. In view of the rescaling relation (28) for the energy, it suffices to prove $`(\alpha _1)<(\alpha _2)`$, where $``$ was defined in (31).
To prove $`(\alpha _1)<(\alpha _2)`$, we write (33) as $`^{}(\alpha )=\delta (\alpha )(1/E)^{}(\alpha )`$, where
$$\delta (\alpha )=\frac{1}{2}P^3A^2_0^Pk_{\alpha }^{}{}_{}{}^{2}𝑑x.$$
The point of this transformation is that $`\delta (\alpha )`$ is strictly decreasing: $`P^{}<0`$ and $`A^{}>0`$, by \[17, Props. 7.3 and 7.4\], while $`\alpha _0^Pk_{\alpha }^{}{}_{}{}^{2}𝑑x`$ is decreasing by \[18, eq. (35)\]. Also $`(1/E)^{}>0`$ on $`(0,\alpha _{crit})`$ and $`(1/E)^{}<0`$ on $`(\alpha _{crit},1)`$, by assumption. Hence
$$^{}(\alpha )>\delta (\alpha _{crit})(1/E)^{}(\alpha )\text{for }\alpha (0,1),\alpha \alpha _{crit}\text{.}$$
Integrating this inequality from $`\alpha _1`$ to $`\alpha _2`$ and using that $`E(\alpha _1)=E(\alpha _2)`$ gives $`(\alpha _2)>(\alpha _1)`$, as desired.
Finally, (33) shows $`^{}(\alpha )<0`$ on $`(\alpha _{crit},1)`$, and so $`(\alpha _2)>(1)=0`$ by (35). Thus (28) yields $`(h_{\mathrm{ss2}})>(\overline{h_{\mathrm{ss2}}})`$. ∎
### 5.5. Proof of Theorem 10
From the definition (11) we see
$$\tau _1(\overline{h})=\underset{u}{\mathrm{min}}\frac{_0^X(u^{})^2𝑑x}{_0^Xu^2𝑑x}r(\overline{h})=\left(\frac{2\pi }{X}\right)^2r(\overline{h}),$$
with the minimum being attained precisely for linear combinations of $`\mathrm{sin}(2\pi x/X)`$ and $`\mathrm{cos}(2\pi x/X)`$. The first two paragraphs of the theorem follow directly.
(a) Now suppose $`r(\overline{h})X^2>4\pi ^2`$. The variational formulas (4.1) and (4.1) in the proof of Theorem 1 show the constant steady state $`\overline{h}`$ is energy unstable in the sine and cosine directions, since these are $`\tau _1(\overline{h})`$-eigenfunctions and $`\tau _1(\overline{h})<0`$. Suppose next $`r(\overline{h})X^2=4\pi ^2`$ and $`r^{\prime \prime }(\overline{h})>0`$. Then the first two variations of the energy in the $`\pm \mathrm{sin}`$ directions are zero, by (4.1) and (4.1). The third variation equals $`r^{}(\overline{h})`$ times the integral of $`\mathrm{sin}^3`$, by (17); thus the third variation is also zero. The fourth variation is $`r^{\prime \prime }(\overline{h})_0^X\mathrm{sin}^4(2\pi x/X)𝑑x`$, which is negative because we assumed $`r^{\prime \prime }(\overline{h})>0`$. Thus the constant steady state is energy unstable in the $`\pm \mathrm{sin}`$ directions. Argue similarly for the $`\pm \mathrm{cos}`$ directions.
(b) If $`r(\overline{h})X^2<4\pi ^2`$, or if $`r(\overline{h})X^2=4\pi ^2`$ and $`r^{\prime \prime }(\overline{h})<0`$, then $`\tau _1(\overline{h})0`$ and so we get energy stability by modifying the argument of part (a) as follows. The first variation of the energy in a direction $`u`$ is always zero. If the second variation is positive then we are done. Otherwise the second variation must be zero, so that $`\tau _1(\overline{h})=0`$ and $`r^{\prime \prime }(\overline{h})<0`$. Then $`u`$ must be a linear combination of sines and cosines, and so the third variation is also zero. Then the fourth variation is positive.
Furthermore, for $`r(\overline{h})X^2<4\pi ^2`$ we will prove $`\overline{h}`$ is a strict local minimum of the energy, and is nonlinearly stable in $`H^1`$. In doing this, we will use below a certain sufficiently small number $`\delta (0,1)`$. Then considering $`uH^1(𝕋_X)`$ with mean value zero and $`u_{H^1(𝕋_X)}=1`$, we find for all $`\epsilon [0,\delta ]`$ that
$`{\displaystyle \frac{d^2}{d\epsilon ^2}}(\overline{h}+\epsilon u)`$ $`=`$ $`{\displaystyle _0^X}\left[(u^{})^2r(\overline{h}+\epsilon u)u^2\right]𝑑x`$
$`=`$ $`\delta {\displaystyle _0^X}\left[(u^{})^2+u^2\right]𝑑x+(1\delta ){\displaystyle _0^X}\left[(u^{})^2{\displaystyle \frac{r(\overline{h}+\epsilon u)+\delta }{1\delta }}u^2\right]𝑑x`$
$`>`$ $`\delta u_{H^1(𝕋_X)}^2+(1\delta ){\displaystyle _0^X}\left[(u^{})^2{\displaystyle \frac{4\pi ^2}{X^2}}u^2\right]𝑑x\delta u_{H^1(𝕋_X)}^2.`$
In the second-to-last step above, we used $`r(\overline{h})X^2<4\pi ^2`$ and $`u_{\mathrm{}}Cu_{H^1}`$ and we chose $`\delta `$ sufficiently small (independent of $`u`$ and $`\epsilon `$).
On the other hand, $`|\epsilon u|\delta u_L^{\mathrm{}}\overline{h}/2`$ provided $`\delta `$ is chosen small enough (independent of $`u`$), and hence
$`{\displaystyle \frac{d^2}{d\epsilon ^2}}(\overline{h}+\epsilon u)={\displaystyle _0^X}\left[(u^{})^2r(\overline{h}+\epsilon u)u^2\right]𝑑x`$ $``$ $`C(\overline{h}){\displaystyle _0^X}\left[(u^{})^2+u^2\right]𝑑x`$
for some constant $`C(\overline{h})1`$
$`=`$ $`C(\overline{h})u_{H^1(𝕋_X)}^2.`$
We deduce from the preceding estimates and the vanishing of the first variation in (4.1) that if $`uH^1(𝕋_X)`$ has mean value zero and $`u_{H^1(𝕋_X)}\delta `$, then $`\overline{h}+u>0`$ and the energy varies quadratically away from $`\overline{h}`$:
(41)
$$\frac{1}{2}\delta u_{H^1(𝕋_X)}^2(\overline{h}+u)(\overline{h})\frac{1}{2}C(\overline{h})u_{H^1(𝕋_X)}^2.$$
The lefthand estimate implies $`\overline{h}`$ is a strict local minimum of the energy, with respect to $`X`$-periodic zero-mean perturbations.
We now prove $`\overline{h}`$ is nonlinearly stable, in the sense that if $`h(x,t)`$ is a smooth positive solution of (1) for $`x𝕋_X`$ and $`t[0,T]`$, for some $`T>0`$, and if $`h(,0)`$ has mean value $`\overline{h}`$ and $`h(,0)\overline{h}_{H^1(𝕋_X)}<\sqrt{\delta ^3/4C(\overline{h})}`$, then $`h(,t)\overline{h}_{H^1(𝕋_X)}<\delta /2`$ for all $`t[0,T]`$. Indeed, the quadratic bounds (41) and the dissipation of the energy together imply
$`{\displaystyle \frac{1}{2}}\delta h(,t)\overline{h}_{H^1(𝕋_X)}^2`$ $``$ $`(h(,t))(\overline{h})`$
$``$ $`(h(,0))(\overline{h}){\displaystyle \frac{1}{2}}C(\overline{h})h(,0)\overline{h}_{H^1(𝕋_X)}^2<{\displaystyle \frac{1}{2}}{\displaystyle \frac{\delta ^3}{4}},`$
so that $`h(,t)\overline{h}_{H^1(𝕋_X)}<\delta /2`$ for all $`t[0,T]`$. This stability result holds for all sufficiently small $`\delta `$. ∎
This proof of nonlinear stability for a linearly stable constant steady state does not carry over to linearly stable positive periodic steady states $`h_{\mathrm{ss}}`$, because there $`\tau _1(h_{\mathrm{ss}})=0`$ (due to translational null directions). This zero eigenvalue is absent for the constant steady state, since translation of $`\overline{h}`$ gives $`\overline{h}`$ again — a trivial perturbation. Imposing Neumann boundary conditions eliminates the translational perturbations and their associated zero eigenvalue. Hence, the nonlinear stability proof would hold for a positive steady state under the Neumann boundary conditions $`h_x=h_{xxx}=0`$ at $`x=0,X/2`$ (see §2.6), provided the steady state is strictly linearly stable, i.e. the first eigenvalue is positive.
### 5.6. Proof of Theorem 11
Suppose first that $`1<q<3`$. By (26) with $`\alpha =0`$, a steady state $`\widehat{h}_{\mathrm{ss}}`$ supported on a single interval of length $`\widehat{X}X`$ and with area $`\widehat{A}_{\mathrm{ss}}=\overline{h}X`$ and zero contact angles exists if and only if for some length $`\widehat{X}X`$ we have
(42)
$$\widehat{X}^{3q}\widehat{A}_{\mathrm{ss}}^{q1}=E(0)=P(0)^{3q}A(0)^{q1}.$$
That is, if and only if $`X^{3q}(\overline{h}X)^{q1}E(0)`$, which is (12).
Suppose such a steady state $`\widehat{h}_{\mathrm{ss}}`$ exists, supported say on the interval $`(0,\widehat{X})`$. By above,
$$\widehat{X}=^{1/(q3)}P(0)[A(0)/\overline{h}X]^{(q1)/(3q)}.$$
We want to show that $`(\widehat{h}_{\mathrm{ss}})<(\overline{h})`$ if and only if (13) holds. For this, compute using $`G`$ like in §5.2 to find
(43) $`(\widehat{h}_{\mathrm{ss}})(\overline{h})`$ $`=`$ $`{\displaystyle _0^X}\left[{\displaystyle \frac{1}{2}}\left(\widehat{h}_{\mathrm{ss}}^{}\right)^2G(\widehat{h}_{\mathrm{ss}})+G(\overline{h})\right]𝑑x`$
$`=`$ $`{\displaystyle _0^{\widehat{X}}}\left[{\displaystyle \frac{1}{2}}\left(\widehat{h}_{\mathrm{ss}}^{}\right)^2G(\widehat{h}_{\mathrm{ss}})\right]𝑑x+G(\overline{h})X`$
$`=`$ $`{\displaystyle \frac{\widehat{A}_{\mathrm{ss}}^2P(0)^3}{A(0)^2\widehat{X}^3}}{\displaystyle _0^{P(0)}}\left[{\displaystyle \frac{1}{2}}(k_0^{})^2G(k_0)\right]𝑑x+G(\overline{h})X`$
when $`q>1,q0`$, by rescaling from $`\widehat{h}_{\mathrm{ss}}`$ to $`k_0`$ and using (25) and (27) (with $`P_{\mathrm{ss}}`$ replaced by $`\widehat{X}`$, and $`A_{\mathrm{ss}}`$ replaced by $`\widehat{A}_{\mathrm{ss}}`$). Putting $`\alpha =0`$ into \[18, eq. (32)\] and \[18, eq. (30)\] and using $`H(0)=0`$ enables us to evaluate $`(k_0^{})^2𝑑x`$ and $`G(k_0)𝑑x`$, and hence we deduce
(44)
$$(\widehat{h}_{\mathrm{ss}})(\overline{h})=\frac{\widehat{A}_{\mathrm{ss}}^2P(0)^3}{A(0)^2\widehat{X}^3}\frac{q3}{q(q+3)}A(0)+G(\overline{h})X\text{when }q>1,q0\text{.}$$
After substituting the definitions of $`\widehat{X}`$ and $`\widehat{A}_{\mathrm{ss}}=\overline{h}X`$ from above, we find
$$(\widehat{h}_{\mathrm{ss}})(\overline{h})=\frac{1}{q}^{3/(3q)}(\overline{h}X)^{(3+q)/(3q)}\left[\frac{q3}{q+3}A(0)^{2q/(q3)}+\frac{1}{q+1}\left(\overline{h}^{q1}X^2\right)^{q/(q3)}\right]$$
when $`1<q<3,q0`$. Plainly now $`(\widehat{h}_{\mathrm{ss}})<(\overline{h})`$ if and only if (13) holds.
When $`q=0`$ we find (43) has an extra term $`\widehat{A}_{\mathrm{ss}}\mathrm{log}\left[\widehat{A}_{\mathrm{ss}}P(0)/A(0)\widehat{X}\right]`$, so that after again using \[18, eq. (32)\] and \[18, eq. (30)\] and substituting for $`\widehat{X}`$ and $`\widehat{A}_{\mathrm{ss}}`$, we obtain
$$(\widehat{h}_{\mathrm{ss}})(\overline{h})=\frac{1}{3}\overline{h}X\mathrm{log}\frac{A(0)^2/e}{\overline{h}^1X^2}\text{when }q=0\text{.}$$
Remembering that $`A(0)=2e^{3/2}\sqrt{\pi /3}`$ when $`q=0`$, from \[17, §3.1.2\], we conclude $`(\widehat{h}_{\mathrm{ss}})<(\overline{h})`$ if and only if $`\overline{h}^1X^2>4e^2\pi /3`$. Incidentally, one can check $`L(q)`$ is continuous at $`q=0`$.
For $`q3`$, simply modify the above proof from the case $`1<q<3`$. Notice when $`q=3`$ that (42) becomes $`(\overline{h}X)^2=E(0)`$, which yields no formula for $`\widehat{X}`$. And when $`q3`$ we get $`(\widehat{h}_{\mathrm{ss}})>(\overline{h})`$ because the first term in (44) is nonnegative and the second is positive. ∎
## 6. Conclusions and Future Directions
If you perturb a positive periodic (or constant) steady state $`h_{\mathrm{ss}}`$ of the evolution equation (1), without changing its area, then towards which steady states might the solution subsequently evolve?
That is the broad question addressed by this paper. To answer it, we focused on three specific questions:
* Existence: Do there exist other steady states having the same area and same period as $`h_{\mathrm{ss}}`$, or having period a fraction of the period of $`h_{\mathrm{ss}}`$? If so, then these other steady states are plausible contenders for the long time limit. The constant steady state $`\overline{h_{\mathrm{ss}}}`$ obviously satisfies these requirements, but there might be another positive periodic steady state (different from $`h_{\mathrm{ss}}`$) that does also, or perhaps a array of droplet steady states having zero or nonzero contact angles.
* Stability: Are $`h_{\mathrm{ss}}`$ or these other steady states linearly stable? energy stable? If a steady state is the long time limit of some generic solution, then surely it must be stable under perturbations.
* Relative energy levels: Do any of these other steady states have lower energy than $`h_{\mathrm{ss}}`$? Only steady states with lower energy are accessible, when starting from a small perturbation of $`h_{\mathrm{ss}}`$.
The existence question was substantially answered for power law coefficients by \[18, Theorem 12\], Theorems 7 and 8 and \[18, Figures 3–5\], also using Theorem 11 when $`h_{\mathrm{ss}}`$ is the constant steady state. But the existence question remains open for droplet steady states with nonzero contact angles, if we wish to specify the area and the length of the support. Some information on nonzero angle droplet steady states is in our earlier paper \[17, §5.2\].
The stability question was resolved for positive periodic steady states in the power law case by \[18, Theorems 1,3,7,9\] and Theorems 15 here. In particular, Theorem 1 related linear instability to energy instability. Theorem 10 handled the case of constant steady states. Our numerical simulations in the companion article suggest that linearly unstable steady states are indeed nonlinearly unstable, with the linear behavior dominating near the steady state, but this observation is so far unsupported by a ‘linearization theorem’ for the power law evolution $`h_t=(h^nh_{xxx})_x(h^mh_x)_x`$. (A linearization theorem is known in the Cahn–Hilliard case $`f1`$, by using semilinear operator theory; see for example \[23, §6\].)
The energy level question has been largely settled in the power law case by Theorems 6, 7 and 9 when $`h_{\mathrm{ss}}`$ is positive and periodic, and by Theorems 6 and 11 when $`h_{\mathrm{ss}}`$ is constant. When $`h_{\mathrm{ss}}`$ has compact support with zero contact angle, use Theorems 7 and 11. For example, Theorems 9 suggests that when $`m=n+0.77`$, there can exist two positive periodic steady states with the same period and area, and that the unstable steady state has higher energy than the stable one. Our numerical simulations have found heteroclinic connections from the high energy steady state to the low energy one.
### Future directions.
The stability question for droplet steady states (with zero and nonzero contact angles) is open. So is the problem of computing relative energy levels of non-zero angle droplet steady states with respect to zero-angle droplets and constant and periodic steady states.
Another open problem is to answer the existence, stability and energy level questions for general coefficient functions $`f`$ and $`g`$. We have treated power law coefficients, and Grinfeld and Novick–Cohen cover the Cahn–Hilliard equation, with $`f1`$ and $`g(y)=13y^2`$. But for general coefficients, about all we know is that every positive periodic steady state is linearly and energy unstable when $`g/f`$ is a convex function, by Theorem 1.
Finally, many of our existence, stability and relative energy level theorems for the power law evolution would be improved if we knew $`E^{}(\alpha )<0`$ for all $`\alpha `$ when $`1<q1.75`$. We have not been able to prove this conjecture, though numerically it is clear from \[18, Figure 3\].
### Acknowledgments
The authors are grateful to Andrew Bernoff for stimulating comments on the energy landscape of phase space.
Laugesen was partially supported by NSF grant number DMS-9970228, and a grant from the University of Illinois Research Board. He is grateful for the hospitality of the Department of Mathematics at Washington University in St. Louis.
Pugh was partially supported by NSF grant number DMS-9971392, by the MRSEC Program of the NSF under Award Number DMR-9808595, by the ASCI Flash Center at the University of Chicago under DOE contract B341495, and by an Alfred P. Sloan fellowship. Some of the computations were done using a network of workstations paid for by an NSF SCREMS grant, DMS-9872029. Part of the research was conducted while enjoying the hospitality of the Mathematics Department and the James Franck Institute of the University of Chicago.
email contact: laugesen@math.uiuc.edu, mpugh@math.upenn.edu
|
warning/0003/nucl-th0003039.html
|
ar5iv
|
text
|
# References
Dispersive effects in neutron matter superfluidity
M. Baldo$``$ and A. Grasso$``$
$``$ INFN, Sez. Catania, Corso Italia 57 - 95129 Catania, Italy
$``$ Dipartimento di Fisica, Universita’ di Catania, Corso Italia 57, 95129 Catania, Italy
Abstract. The explicit energy dependence of the single particle self-energy (dispersive effects), due to short range correlations, is included in the treatment of neutron matter superfluidity. The method can be applied in general to strong interacting fermion systems, and it is expected to be valid whenever the pairing gap is substantially smaller than the Fermi kinetic energy. The results for neutron matter show that dispersive effects are strong in the density region near the gap closure.
PACS: 21.65+f , 26.60+c , 97.60Jd
1. Introduction
Neutron and nuclear matter superfluidity is one of the main issue in the physics of neutron stars. Superfluidity is expected to play a major role in some of the most striking phenomena occurring in neutron stars, like glitches and post-glitches transients , vortex pinning , neutron star cooling, and maybe strong magnetic field penetration . However, since the observational data are only indirectly related to superfluidity, or need explicit models for their interpretation, a firm theoretical prediction of the superfluidity strength, based on microscopic ab initio calculations, appears highly required. Unfortunately, neutron and nuclear matter are strongly correlated systems, where short range correlations dominate the overall interaction energy, even at densities well below the saturation value. The superfluidity problem turns out, therefore, to be a complex many-body problem, where the delicate balance between short range interactions and the long range pairing correlations needs an accurate treatment. This problem was firstly considered in the works of ref. , within the variational Jastrow method. The medium effect on the effective pairing interaction was investigated in refs. . The correlation effects were treated in ref. within a generalization of the Babu-Brown approach to the effective nucleon-nucleon (NN) interaction and in the weak coupling limit. In general, these microscopic approaches seem to indicate a reduction of the pairing gap due to the medium, with respect to BCS approximation with the bare interaction. The use in the BCS scheme of realistic bare nucleon-nucleon interactions, which reproduce the experimental phase shifts, can be a good starting point for a more sophisticated many-body treatments, and the connection between the pairing gap value and the phase shifts has been elucidated, in general, in ref. .
In all these microscopic approaches the single particle spectrum is usually considered within the effective mass approximation, or taken from normal state calculations. Dispersive effects, due to the energy dependence of the single particle self-energy, are usually neglected, or considered only in the weak coupling limit, and no first principle scheme to include them in the general gap equation has been proposed. Only recently a self-consistent scheme has been developed, where the short range correlations and the pairing problem are treated on the same footing. The method is numerically complex, and it has been solved only for a schematic interaction . It would be desirable to understand in simple terms the effect of the self-energy dispersion on the pairing strength, also in the case of strong coupling as in neutron matter. In this letter we present a general scheme for including dispersive effects in the gap equation, which is simple and accurate, provided the gap is substantially smaller than the Fermi kinetic energy. We then apply the method to neutron matter superfluidity with realistic interaction, and we show numerically that the size of the effect can be large in the vicinity of the gap closure.
2. Including self-energy in the gap equation The general many-body theory of pairing in fermion systems has been formulated just after the BCS solution has been introduced, and it can be found in standard textbooks . The main framework is the Green’ s function (GF) formalism, which generalizes the Gorkov’ s method beyond the BCS approximation. The single particle Green’s function $`𝒢`$ has a 2x2 matrix structure, with normal diagonal components $`F_1`$ and abnormal off-diagonal components $`F_2`$,
$$𝒢(k,\omega )=\left(\begin{array}{cc}F_1(k,\omega )F_2(k,\omega )& \\ F_2(k,\omega )F_1(k,\omega )& \end{array}\right),𝒢^1(k,\omega )=\left(\begin{array}{cc}\stackrel{~}{ϵ}_k\omega +M(k,\omega )\mathrm{\Delta }(k,\omega )& \\ \mathrm{\Delta }(k,\omega )(\stackrel{~}{ϵ}_k+\omega +M(k,\omega ))& \end{array}\right)$$
(1)
In the expression for the inverse Green’s function $`𝒢^1`$, we have introduced the quantity $`\stackrel{~}{ϵ}_k=\mathrm{}^2k^2/2m\mu `$ as the single particle kinetic energy, with respect to the chemical potential $`\mu `$, the diagonal single particle self-energy $`M(k,\omega )`$ and the momentum and energy dependent gap function $`\mathrm{\Delta }(k,\omega )`$. Here we are assuming s-wave singlet pairing, and therefore we omit spin indices. They simply express the coupling between the time-reversal states $`(k,)`$ and $`(k,)`$. Both $`M(k,\omega )`$ and $`\mathrm{\Delta }(k,\omega )`$ can be expanded in terms of the NN interaction and the full GF itself, which in general entails a self-consistent procedure. The gap function $`\mathrm{\Delta }(k,\omega )`$, however, is solution of the homogeneous generalized Bethe-Salpeter equation , and as such it satisfies the generalized gap equation
$$\mathrm{\Delta }(k,\omega )=\underset{k^{}}{}𝑑\omega ^{}\frac{I(k\omega ,k^{}\omega ^{})\mathrm{\Delta }(k^{},\omega ^{})}{(\stackrel{~}{ϵ}_k^{}\omega ^{}+M(k^{},\omega ^{}))(\stackrel{~}{ϵ}_k^{}+\omega ^{}+M(k,\omega ^{}))+\mathrm{\Delta }(k^{},\omega ^{})^2}$$
(2)
where $`I(k\omega ,k^{}\omega ^{})`$ is the irreducible NN interaction at zero total energy and momentum. If one takes the bare NN interaction for the interaction $`I`$, and the Hartree-Fock approximation for the diagonal self-energy $`M(k,\omega )`$, the standard BCS approximation is recovered. It has to be noticed that the energy dependence of the gap function $`\mathrm{\Delta }(k,\omega )`$ originates only from the energy dependence of the irreducible interaction $`I`$. In fact, if the interaction is taken as energy-independent, the gap function is also energy-independent, despite the possible energy dependence of the self-energy $`M(k,\omega )`$. Since we are looking for an estimate of the dispersive effects, we indeed assume the irreducible interaction as energy independent, while we keep the full energy dependence of the self-energies. Then, the energy integration appearing in the gap equation (2) can be performed with a good accuracy in the limit of small self-energy imaginary part, since than the main contribution is expected to come from the poles close to the real axis. The denominator is an even function of the energy $`\omega `$ (we remind again that single particle energies are measured with respect to $`\mu `$), and, therefore, the kernel presents two poles, symmetrical with respect to the origin in the complex $`\omega `$-plane. This is a feature typical of the superconducting phase. Formally, the pole energies $`\pm E_k`$ are the solutions of the implicit equation
$$\pm E_k=\frac{1}{2}(M(k,\pm E_k)M(k,E_k))\pm \sqrt{\left[\stackrel{~}{ϵ}_k+\frac{1}{2}(M(k,E_k)+M(k,E_k))\right]^2+\mathrm{\Delta }(k)^2}$$
(3)
If the energy dependence of $`M(k,\omega )`$ is neglected, than Eq. (3) reduces to the usual square root expression for the quasi-particle excitation energy of the BCS approximation. On the other hand, in the non-superconducting limit $`\mathrm{\Delta }0`$, and neglecting the imaginary part of $`M(k,\omega )`$, one can verify that Eq. (3) reduces to the usual self-consistent equation, e.g. Brueckner , for the single particle energy $`e_k`$
$$e_k=\stackrel{~}{ϵ}_k+M(k,e_k)$$
(4)
Equation (4) is valid whenever $`\mathrm{\Delta }`$ is negligible, in particular for momenta far away from the Fermi surface, since then $`|\mathrm{\Delta }_k||\stackrel{~}{ϵ}_k|`$. We now make use of the assumption of a gap $`\mathrm{\Delta }`$ smaller than the normal self-energy, which is mainly determined by short range correlations. The size of the normal self-energy is indeed of the same order of the Fermi kinetic energy $`E_F`$. In this case, Eq. (4) will be valid to order $`\mathrm{\Delta }(k_F)/E_F`$, and therefore on the right hand side of Eq. (3) we can replace $`E_k`$ with $`e_k`$, solution of Eq. (4), to get
$$E_k\frac{1}{2}(M(k,|e_k|)M(k,|e_k|))+\sqrt{\left[\stackrel{~}{ϵ}_k+\frac{1}{2}(M(k,e_k)+M(k,e_k))\right]^2+\mathrm{\Delta }(k)^2}$$
(5)
where the self-energy is now calculated in the normal phase. The procedure is justified, provided $`M(k,\omega )`$ is a smooth function of $`\omega `$. It will be further discussed below.
The residue $`R`$ of the kernel at each one of the pole can be easily calculated
$`R=\left[1{\displaystyle \frac{1}{2}}((1\mathrm{\Theta }_k)a_k+(1+\mathrm{\Theta }_k)b_k)\right]^1`$ (6)
$`\mathrm{\Theta }_k={\displaystyle \frac{\stackrel{~}{ϵ}_k+\frac{1}{2}(M(k,E_k)+M(k,E_k))}{E_k}}`$ (7)
$`a_k=\left({\displaystyle \frac{M}{\omega }}\right)_{\omega =E_k};b_k=\left({\displaystyle \frac{M}{\omega }}\right)_{\omega =E_k}`$ (8)
In the limit $`\mathrm{\Delta }0`$, Eq. (6) is the usual expression for the quasi-particle strength $`Z_k`$, provided the momentum $`k`$ is close enough to $`k_F`$. The corrections to the normal phase value of $`1R^1`$ are of the order $`\mathrm{\Delta }(k_F)/E_F`$, and therefore the residue $`R`$ can be identified with $`Z_k`$, at least in the vicinity of the Fermi momentum. Far away from the Fermi momentum, the residue $`R`$ has still the expression of Eq. (6) in the limit of small imaginary part, despite the quasi-particle concept becomes less meaningful, since its width can be large (but it can be still much smaller than the real part of the energy). In this case the procedure is just an accurate method of calculating the energy integral (pole approximation). For simplicity, the residue will be denoted by $`Z_k`$ in all cases. Within these approximations, the generalized gap equation (2) reads
$$\mathrm{\Delta }(k)=\underset{k^{}}{}I(k,k^{})Z_k^{}\frac{\mathrm{\Delta }_k^{}}{2\sqrt{\left[\stackrel{~}{ϵ}_k+\frac{1}{2}(M(k,e_k)+M(k,e_k))\right]^2+\mathrm{\Delta }(k)^2}}$$
(9)
Eqs. (5), (6), (9) contain the main result of the paper. It has to noticed that in the generalized gap equation (9) the square root in the denominator does not coincide with the quasi-particle energy of Eq. (5), in contrast with the usual BCS approximation. This feature is general and it is not bound to the approximation of Eq. (5). The approximation of Eq. (5) to the expression inside the square root is valid up to corrections of order $`\mathrm{\Delta }/E_F`$. They can be absorbed in large part by a small shift of the chemical potential $`\mu `$.
Before going to the application of the formalism to neutron matter, let us discuss Eq. (9) in the extreme weak coupling limit, where one assumes that the main contribution to the momentum integral is concentrated around the Fermi surface. In this limit, following the standard procedure of expanding the integrand of Eq. (9) around $`k_F`$, one gets
$$\mathrm{\Delta }_F=\mathrm{\hspace{0.17em}8}\frac{E_F}{m_F}\mathrm{exp}(\frac{1}{m_FZ_F\pi ^2n_0I(k_F)})$$
(10)
where $`n_0`$ is the density of state for the free Fermi gas and $`m_F`$ the so called $`kmass`$ (in units of the bare mass) . The interaction $`I(k_F)`$ is the diagonal matrix element of the NN potential in the considered channel ( $`{}_{}{}^{1}S_{0}^{}`$ for neutron matter), in the plane wave representation. The self-energy effects are, therefore, contained mainly in the factor $`m_FZ_F`$, which can be written also as $`m^{}Z_F^2`$, since the full effective mass $`m^{}=m_F/Z_F`$ . This is the standard result for the weak coupling limit . Eq. (9) generalizes the treatment to the case where the contribution from momenta far from the Fermi momentum is relevant. The appearance of the $`kmass`$ is a peculiar feature of the pairing phenomenon, and it is a direct consequence of the coupling between time-reversal states. In Eq. (9) the combination $`M(k,\omega )+M(k,\omega )`$ gives rise to the combined density of state of the pair $`\{(k,\omega );(k,\omega )\}`$, which is mainly determined by the $`kmass`$.
The weak coupling limit is not valid in general for neutron or nuclear matter , if one starts from the bare NN interaction. This can be seen directly from the observation that often the gap equation has a well defined solution even when the interaction matrix element $`I(k_F)`$ is positive. This is due to the dominant role of the off-diagonal matrix elements $`I(k,k^{})`$. Therefore, one must solve the more general equation (9) in this case. The above considerations are, anyhow, still valid.
3. Application to neutron matter superfluidity In order to estimate the dispersive effects on the superfluid gap of neutron matter, we have solved Eq. (9), with the bare Argonne v<sub>14</sub> potential as the pairing interaction $`I(k,k^{})`$ and with the self-energy calculated in the Brueckner approximation at the lowest order, see Fig. 1a, with the same interaction. The higher order contribution of Fig. 1b turns out to be indeed negligible in the relevant density range. In the superfluid phase, in principle, the diagonal self-energy $`M(k,\omega )`$ differs from the self-energy in the normal phase. The main contribution not present in the normal phase originates from the coupling of the single particle motion with the superfluid collective modes. The latter correspond mainly to the center of mass motion of the Cooper pairs and their possible “vibrations” . The branch starting at zero energy, in the long wave-length limit, is the branch of the Goldstone boson , corresponding to the gauge invariance symmetry breaking at the superfluid phase transition. This contribution to the diagonal single particle self-energy is expected to be at most of the order of the superfluid condensation energy per particle, and therefore negligible with respect to the typical short range correlation energy, as calculated e.g. in Brueckner theory, at least to the extent that $`\mathrm{\Delta }/E_F1`$. For the same reason, the deviation of the occupation number from the free gas value and the presence of a forbidden energy region, of order $`\mathrm{\Delta }`$, around the Fermi energy, typical of the pairing phenomenon, seem to play no relevant role in determining the size of the self-energy. It appears, therefore, justified to adopt for $`M(k,\omega )`$ its normal phase value.
The choice of the bare interaction for $`I(k,k^{})`$ is suggested by the observation that no ladder summation should be included in the irreducible interaction kernel $`I(k,k^{})`$ . Of course, other terms, like polarization diagrams, should be included , but here we want simply to single out the dispersive effects, and therefore it appears meaningful to compare the results obtained with and without self-energy, within the same scheme of approximation.
In Fig. 2 is reported, for $`k_F=0.9fm^1`$, the imaginary part of the neutron self-energy $`M_I(k,e(k))`$ at the quasi-particle pole, as a function of the momentum $`k`$, together with the real part of the quasi-particle energy $`e(k)`$ (calculated with respect to $`\mu `$). One can see that indeed the imaginary part is small with respect to the real part. The situation is completely similar for the other densities. The residue $`Z_k`$ of Eq. (6), which appears in Eq. (9) for the gap function, is reported in Fig. 3 for three densities. According to Migdal-Luttinger theorem , the value $`Z_F`$ of $`Z_k`$ at $`k=k_F`$ is the discontinuity of the momentum distribution at the Fermi momentum (in the normal phase). One must have, therefore, $`0<Z_F<1`$. One can see, however, that $`Z_k`$ exceeds 1 slightly in some interval well above $`k_F`$. This is not surprising, since for large $`\omega `$ values, at fixed $`k`$, the real part of the self-energy is an increasing function of $`\omega `$, and it approaches asymptotically an energy-independent value (and therefore $`Z_k1`$). As already mentioned, for large momentum the pole approximation is just an accurate method of calculating the relevant energy integral. The position of the pole is, of course, not exactly on the real axis, but numerical estimate of the second derivative of the self-energy shows that to calculate $`Z_k`$ on the real axis is an extremely good approximation. The factor $`Z_k`$ is also related to the single particle occupation number $`n(k)`$. In fact, the contour integral, closed in the lower complex plane, of the single particle GF equals $`1n(k)`$, $`k>k_F`$. Besides the pole, this integral receives contribution from the regular (non polar) part of the single particle GF . If $`Z(k)>1`$, this means simply that the regular part leads to a slightly negative contribution. We have checked, indeed, that the single particle spectral function, calculated with the same self-energy (including of course the imaginary part), satisfies the sum rules and gives well defined occupation numbers $`n(k)`$. A full account of the calculations will be reported elsewhere. Anyhow, this small deviation from 1 does not affect at all the final results, and one can take $`Z_k=1`$ in this momentum region.
Finally, in Fig. 4 is reported the pairing gap at the Fermi energy as a function of density for three different cases : i) without self-energy , ii) with the $`Z_k`$ factor in the numerator of the gap equation (9), iii) with both the self-energy in the denominator and the $`Z_k`$ factor. The reduction of the pairing gap is substantial at the highest densities, near the gap closure.
3. Discussion and conclusion We have developed a method to include the single particle self-energy in the gap equation, and in particular dispersive effects, beyond the usual BCS approximation. The method rely on the assumption of a small gap with respect to the Fermi kinetic energy and strong short range correlations, typical of the neutron matter in the inner crust of neutron stars. The results indicate that dispersive effects can strongly reduce the gap value near its closure. The effect is due both to the quasi-particle strength $`Z_k`$ and to the $`kmass`$, which enter in the generalized gap equation (9). This result appears in line with the work of ref. , where the self-consistent treatment of pairing and short range correlations seems indeed to reduce strongly the gap value mainly because of these two factors . Of course, before drawing any conclusion on the pairing strength in neutron matter, one should include, along a consistent scheme, the correlation effects on the irreducible interaction $`I(k,k^{})`$. Work in this direction is in progress.
The extension of the method to symmetric nuclear matter, possibly relevant for pairing in nuclei, appears problematic. In that case the approximation of a small imaginary part looks less justified away from the Fermi momentum . Furthermore, the neutron-proton pairing in the $`{}_{}{}^{3}S_{1}^{}^3D_1`$ channel is too strong, i.e. $`\mathrm{\Delta }E_F`$, to be treatable in the proposed approximate method. In both cases a self-consistent procedure, involving both the self-energy and the effective interaction, seems to be the only viable method.
Figure captions
Fig. 1.- One (a) and two (b) hole-line diagrams contributing to the nucleon self-energy in the Bethe-Brueckner-Goldstone expansion. The wavy lines indicate Brueckner G-matrices.
Fig. 2.- Real part e(k) and imaginary part $`M_I`$ of the quasi-particle pole, as a function of the momentum $`k`$ at the Fermi momentum $`k_F=0.9fm^1`$.
Fig. 3.- The residue $`Z_k`$ at the quasi-particle pole as a function of the momentum $`k`$ at three values of the Fermi momentum $`k_F`$.
Fig. 4.- The superfluid gap value, at the Fermi momentum, as a function of density, in the case of free single particle spectrum (diamonds), with the inclusion of the factor $`Z_k`$ (crosses) and with the inclusion of both $`Z_k`$ factor and the self-energy in the single particle spectrum (squares).
|
warning/0003/nucl-th0003063.html
|
ar5iv
|
text
|
# I Introduction and Conclusions
## I Introduction and Conclusions
The presence of attractive interactions in a degenerate fermionic system destabilizes the Fermi surface and leads to the formation of Cooper pairs — the system becomes a superconductor . When increasing the quark density in cold quark matter, asymptotic freedom implies that single-gluon exchange becomes the dominant interaction between quarks. Single-gluon exchange is attractive in the color-antitriplet channel, and therefore leads to color superconductivity in cold, dense quark matter .
Considerable activity has been recently generated by the observation that the zero-temperature color-superconducting gap $`\varphi _0`$ can be as large as 100 MeV . Gaps of this magnitude may have important consequences for the physics of nuclear collisions. The critical temperature for the onset of color superconductivity, $`T_c`$, is (to leading order in the strong coupling constant $`g`$) related to $`\varphi _0`$ in the same way as in BCS theory, $`T_c0.57\varphi _0`$ . Thus, for $`\varphi _0100`$ MeV, it cannot be excluded that color-superconducting quark matter could be created in nuclear collisions in the GSI–SIS or BNL–AGS energy range.
In order to compute the color-superconducting gap, one commonly solves a gap equation . For theories with point-like four-fermion interactions, $`\varphi _0\mu \mathrm{exp}(c_{4\mathrm{F}}/G^2)`$, where $`\mu `$ is the quark chemical potential, $`c_{4\mathrm{F}}`$ is a constant, and $`G^2`$ is the four-fermion coupling strength . On the other hand, in QCD, $`\varphi _0\mu \mathrm{exp}(c_{\mathrm{QCD}}/g)`$, where $`c_{\mathrm{QCD}}`$ is another constant . In weak coupling, there are then three different energy scales, $`\varphi _0m_g\mu `$, where $`m_g`$ is the gluon mass. At $`T=0`$ and for $`N_f`$ flavors of massless quarks ,
$$m_g^2\frac{N_f}{3}\frac{g^2\mu ^2}{2\pi ^2}.$$
(1)
The value of $`c_{\mathrm{QCD}}`$ depends on the form of the gluon propagator in the cold, dense quark medium. If one takes the gluon propagator in the standard “hard dense loop“ (HDL) approximation , one obtains $`c_{\mathrm{QCD}}=3\pi ^2/\sqrt{2}`$ . In this approximation, the quarks inside the HDL’s are assumed to be in the normal, and not the color-superconducting phase. Consequently, an important question that has to be addressed is how color superconductivity influences the propagation of gluons and whether this could change $`c_{\mathrm{QCD}}`$. In a recent work , I have derived a general expression for the gluon self-energy in a two-flavor color superconductor, and explicitly computed the self-energy in the static limit, $`p_0=0`$, for gluon momenta $`p|𝐩|0`$, and for $`p\varphi _0`$.
In a two-flavor color superconductor, condensation of Cooper pairs in a channel of total spin $`J=0`$ breaks $`SU(3)_c`$ to $`SU(2)_c`$. Then, the three gluons corresponding to the generators of the unbroken $`SU(2)_c`$ subgroup are expected to remain massless, while the other five should attain a mass through the Anderson–Higgs mechanism. An explicit computation of the gluon self-energy to one-loop order in perturbation theory confirms this qualitative expectation, but quantitatively reveals some surprising details . At $`T=0`$, the three gluons of the unbroken $`SU(2)_c`$ attain no Meissner mass, but also no Debye mass. This means that static, homogeneous color-electric fields of the unbroken $`SU(2)_c`$ subgroup are not screened. Furthermore, the Debye and Meissner masses of the remaining five gluons are not identical: four gluons have a Debye mass $`\sqrt{3/2}m_g`$ and a Meissner mass $`m_g/\sqrt{2}`$, while the last one has a Debye mass $`\sqrt{3}m_g`$, like in the non-superconducting phase, but a Meissner mass $`m_g/\sqrt{3}`$.
On the other hand, in a three-flavor color superconductor the color and flavor $`SU(3)_c\times U(3)_V\times U(3)_A`$ symmetry is broken to the diagonal subgroup $`SU(3)_{c+V}`$. This locks color and flavor rotations . From the 18 Goldstone bosons resulting from symmetry breaking, eight get “eaten” by the gluons, which consequently become massive. The purpose of this paper is to complement the results of for the two-flavor case with the computation of the Debye and Meissner masses of these gluons in the three-flavor case.
It is worthwhile mentioning that here, as well as in , the terms “Debye mass” and “Meissner mass” refer exclusively to the screening of color-electric and color-magnetic fields. The “ordinary” (electro-)magnetic Meissner effect was studied in . Similar to the mixing of weak and electromagnetic gauge bosons in the standard model, the electromagnetic field mixes with the eighth gluon to form a modified photon, under which the color-superconducting condensate is electrically neutral. The mass of the modified eighth gluon becomes slightly larger than that of the other seven. However, the mixing angle as well as this increase in mass is determined by the ratio of electromagnetic and strong coupling constants and consequently quite small. Therefore, effects from electromagnetism will be neglected throughout the following.
There is another reason why it is important to know the values for the Debye and Meissner mass. The relevant degrees of freedom in the color-flavor locked phase at energy scales much smaller than the gap, $`\varphi _0`$, are the remaining 10 Goldstone bosons resulting from the breaking of color and flavor symmetries. Apart from an additional Goldstone boson arising from breaking $`U(1)_V`$, these bosons are analogous to the pseudoscalar mesons which result from chiral symmetry breaking in the QCD vacuum . Consequently, the dynamics of these 10 Goldstone bosons is described by an effective theory which resembles the Lagrangian of the nonlinear sigma model, describing the dynamics of the chiral fields in the QCD vacuum . To lowest order, this Lagrangian contains only kinetic terms,
$$_{\mathrm{nl}\mathrm{\Sigma }}^{\mathrm{kin}}=\frac{f_\pi ^2}{4}\mathrm{Tr}\left(_0\mathrm{\Sigma }^{}_0\mathrm{\Sigma }v_\pi ^2\mathbf{}\mathrm{\Sigma }^{}\mathbf{}\mathrm{\Sigma }\right)+12f_\eta ^{}^2\left[\left(_0\theta \right)^2v_\eta ^{}^2\left(\mathbf{}\theta \right)^2\right]+12f_H^2\left[\left(_0\phi \right)^2v_H^2\left(\mathbf{}\phi \right)^2\right].$$
(2)
Here, $`\mathrm{\Sigma }\mathrm{exp}(i\lambda ^a/f_\pi )`$, where $`\pi ^a,a=1,\mathrm{},8,`$ are the fields corresponding to the meson octet in the QCD vacuum. These are the pions, the kaons, and the $`\eta `$ meson, which are the Goldstone bosons resulting from spontaneous breaking of the $`SU(3)_A`$ symmetry. $`\lambda ^a`$ are the Gell-Mann matrices. $`\theta `$ is the field corresponding to the meson singlet in the QCD vacuum, i.e., the $`\eta ^{}`$ meson, which is the Goldstone boson resulting from breaking $`U(1)_A`$ spontaneously. (In the QCD vacuum, $`U(1)_A`$ is also broken explicitly by instantons. At the quark densities relevant for the effective theory in the color-flavor locked phase, however, instantons play no longer any significant role .) Finally, $`\phi `$ is the Goldstone mode resulting from breaking $`U(1)_V`$. This field has no analogon in the QCD vacuum.
The coefficients of the kinetic terms are determined by the “decay constants” $`f_\pi `$, $`f_\eta ^{}`$, and $`f_H`$. The presence of a medium breaks Lorentz invariance, and the coefficients of the time-like and space-like terms may differ, i.e., the velocities $`v_\pi ^2`$, $`v_\eta ^{}^2`$, and $`v_H^2`$ are not necessarily equal to one. These velocities enter the dispersion relation of the Goldstone bosons as $`ϵ_i^2(k)=v_i^2k^2+m_i^2`$, $`i=\pi ,\eta ^{},H`$.
The decay constants as well as the velocities can be computed by matching the effective theory to the underlying microscopic theory. A convenient way to do this was proposed by Son and Stephanov . First, they observed that, by minimally gauging the model (2), for instance the pseudoscalar decay constant $`f_\pi `$ is related to the Debye mass of the gluons, $`m_D`$, via
$$f_\pi m_D/g,$$
(3)
and the velocity of the pseudoscalar mesons is determined from the ratio of Debye and Meissner masses,
$$v_\pi m_M/m_D.$$
(4)
The problem is thus reduced to computing these masses in the underlying theory. In this case, one has two choices. First, one may use QCD in the color-superconducting ground state. This theory has quasiparticle as well as quasi-antiparticle excitations, and is valid at all energy scales. Second, one may start from the effective theory proposed by Hong . This theory contains only quasiparticle excitations around the Fermi surface, quasi-antiparticle excitations have already been integrated out. It is valid at energy scales which are much smaller than the chemical potential, $`\mu `$, but which can be larger than the gap, $`\varphi _0`$.
Son and Stephanov used the latter theory to compute the Debye and Meissner masses. Their results are
$$m_D^2=m_g^2\frac{218\mathrm{ln}2}{18},m_M^2=m_g^2\frac{218\mathrm{ln}2}{54}.$$
(5)
Consequently, the velocity of the pseudoscalar mesons is $`v_\pi =1/\sqrt{3}`$. Note that the (square of the) Debye mass in the three-flavor superconductor, $`m_D^20.859m_g^2`$, is reduced by a factor 3.5 as compared to its value in a normal, cold medium, $`m_D^2=3m_g^2`$.
The result (5) is not undisputed throughout the literature. For instance, Rho, Shuryak, Wirzba, and Zahed computed the Debye and Meissner masses from the gluon self-energy in the full theory, including quasi-antiparticle excitations. They obtained \[see Eqs. (A.72) and (A.75) of \],
$$m_D^2=\frac{1}{2}m_g^2,m_M^2=\frac{5}{6}m_g^2.$$
(6)
This is quite surprising, as it implies that the velocity of the pseudoscalar mesons is superluminous, $`v_\pi m_M/m_D=\sqrt{5/3}`$. Other results that can be found in the literature are those of Zarembo , which agree with Son and Stephanov’s calculation. Beane, Bedaque, and Savage agree with Son and Stephanov on the Debye and Meissner masses up to a factor of 2.
The second goal of this paper is to resolve this ambiguity in the literature. The Debye and Meissner masses will be computed in the full theory, i.e., QCD in the color-superconducting ground state. The framework for such a computation was already established in . As shown in the following section, the results are found to agree with those of Son and Stephanov, Eq. (5).
The Debye and Meissner masses are not only important for the nonlinear version (2) of the effective low-energy theory in a three-flavor color superconductor. As outlined in they also determine the coefficients of the kinetic terms in the linear version of the effective theory,
$$_{\mathrm{l}\mathrm{\Sigma }}^{\mathrm{kin}}=\alpha _\mathrm{e}\underset{h=r,\mathrm{}}{}\mathrm{Tr}\left[\left(D_0\mathrm{\Phi }_h\right)^{}D^0\mathrm{\Phi }_h\right]+\alpha _\mathrm{m}\underset{h=r,\mathrm{}}{}\mathrm{Tr}\left[\left(D_i\mathrm{\Phi }_h\right)^{}D^i\mathrm{\Phi }_h\right].$$
(7)
Since there is no reason why right- and left-handed terms should differ in normalization, I assumed $`\alpha _{\mathrm{e}r}\alpha _\mathrm{e}\mathrm{}\alpha _\mathrm{e}`$, and similarly for the coefficient of the space-like terms, $`\alpha _\mathrm{m}`$. For color-flavor locking, the order parameter is a $`3\times 3`$ matrix with expectation value $`\mathrm{\Phi }_h=\mathrm{diag}(\varphi _{0h},\varphi _{0h},\varphi _{0h})`$ . Consequently, $`g^2\alpha _\mathrm{e}(\varphi _{0r}^2+\varphi _0\mathrm{}^2)m_D^2`$, $`g^2\alpha _\mathrm{m}(\varphi _{0r}^2+\varphi _0\mathrm{}^2)m_M^2`$. Due to explicit symmetry breaking by nonzero quark masses (and, at less than asymptotically high densities, instantons), the true ground state of the color-flavor locked phase corresponds to the $`J^P=0^+`$ channel where $`\varphi _{0r}\varphi _0\mathrm{}\varphi _0`$. Then, $`\alpha _\mathrm{e}m_D^2/(2g^2\varphi _0^2)`$, $`\alpha _\mathrm{m}^2m_M^2/(2g^2\varphi _0^2)`$.
I use natural units, $`\mathrm{}=c=k_B=1`$, and work in Euclidean space-time $`𝐑^4V/T`$, where $`V`$ is the volume and $`T`$ the temperature of the system. Nevertheless, I find it convenient to retain the Minkowski notation for 4-vectors, with a metric tensor $`g^{\mu \nu }=\mathrm{diag}(+,,,)`$. For instance, the space-time coordinate vector is $`x^\mu (t,𝐱)`$, $`ti\tau `$, where $`\tau `$ is Euclidean time. 4-momenta are denoted as $`K^\mu (k_0,𝐤)`$, $`k_0i\omega _n`$, where $`\omega _n`$ is the Matsubara frequency, $`\omega _n2n\pi T`$ for bosons and $`\omega _n(2n+1)\pi T`$ for fermions, $`n=0,\pm 1,\pm 2,\mathrm{}`$. The absolute value of the 3-momentum $`𝐤`$ is denoted as $`k|𝐤|`$, and its direction as $`\widehat{𝐤}𝐤/k`$.
## II Explicit computation of Debye and Meissner masses
A convenient starting point to compute the gluon self-energy in the color-flavor locked phase is Eq. (68) of ,
$`\mathrm{\Pi }_{ab}^{\mu \nu }(P)`$ $`=`$ $`{\displaystyle \frac{1}{2}}g^2{\displaystyle \frac{T}{V}}{\displaystyle \underset{K}{}}\mathrm{Tr}_{s,c,f}[\mathrm{\Gamma }_a^\mu G^+(K)\mathrm{\Gamma }_b^\nu G^+(KP)+\overline{\mathrm{\Gamma }}_a^\mu G^{}(K)\overline{\mathrm{\Gamma }}_b^\nu G^{}(KP)`$ (9)
$`+\mathrm{\Gamma }_a^\mu \mathrm{\Xi }^{}(K)\overline{\mathrm{\Gamma }}_b^\nu \mathrm{\Xi }^+(KP)+\overline{\mathrm{\Gamma }}_a^\mu \mathrm{\Xi }^+(K)\mathrm{\Gamma }_b^\nu \mathrm{\Xi }^{}(KP)].`$
Here, the trace is over color, flavor, and spinor space. The vertices are $`\mathrm{\Gamma }_a^\mu \gamma ^\mu T_a`$ and $`\overline{\mathrm{\Gamma }}_a^\mu \gamma ^\mu T_a^T`$. $`G^\pm `$ and $`\mathrm{\Xi }^\pm `$ are the diagonal and off-diagonal elements of the Nambu–Gor’kov propagator for quasiparticle excitations,
$$G^\pm (G_0^\pm \mathrm{\Sigma }^\pm )^1,\mathrm{\Xi }^\pm G_0^{}\mathrm{\Phi }^\pm G^\pm .$$
(10)
$`G_0^\pm (K)\left(\gamma K\pm \gamma _0\mu \right)^1`$ is the propagator for massless, non-interacting quarks (charge-conjugate quarks), and $`\mathrm{\Sigma }^\pm \mathrm{\Phi }^{}G_0^{}\mathrm{\Phi }^\pm `$ is the quark self-energy generated by exchanging particles or charge-conjugate particles with the condensate. In mean-field approximation, the condensate $`\mathrm{\Phi }^+`$ is computed from the gap equation discussed in , and $`\mathrm{\Phi }^{}`$ can be obtained from
$$\mathrm{\Phi }^{}(K)\gamma _0\left[\mathrm{\Phi }^+(K)\right]^{}\gamma _0.$$
(11)
In Eq. (9), the first line corresponds to the diagram in Fig. 1(a), where only the diagonal components of the Nambu–Gor’kov propagator appear, while the second line corresponds to the diagram in Fig. 1(b), formed from the off-diagonal components. Note that, at small temperatures $`T\varphi _0`$, and in weak coupling, $`\varphi _0\mu `$, the fermion loops of Fig. 1 constitute the dominant contribution to the gluon self-energy, since contributions from gluon (or ghost) loops are relatively suppressed by a factor $`T^2/\mu ^2\varphi _0^2/\mu ^2`$ .
In the two-flavor case, the trace over color and flavor in Eq. (9) could be performed independently . In the three-flavor case, due to color-flavor locking this is no longer possible. The traces over color and flavor must be performed simultaneously. The most elegant way to do this is to utilize the color-flavor space projectors introduced by Shovkovy and Wijewardhana, Eq. (7) of ,
$`𝒞_{}^{(1)}{}_{rs}{}^{ij}`$ $``$ $`{\displaystyle \frac{1}{3}}\delta _r^i\delta _s^j,`$ (13)
$`𝒞_{}^{(2)}{}_{rs}{}^{ij}`$ $``$ $`{\displaystyle \frac{1}{2}}\left(\delta _{rs}\delta ^{ij}\delta _r^j\delta _s^i\right),`$ (14)
$`𝒞_{}^{(3)}{}_{rs}{}^{ij}`$ $``$ $`{\displaystyle \frac{1}{2}}\left(\delta _{rs}\delta ^{ij}+\delta _r^j\delta _s^i\right){\displaystyle \frac{1}{3}}\delta _r^i\delta _s^j.`$ (15)
(To avoid proliferation of the symbol $`𝒫`$, I denote them here as $`𝒞`$.) All projectors are symmetric under simultaneous exchange of color, $`i,j`$, and flavor, $`r,s`$, indices. Note that $`𝒞^{(1)}`$ is the singlet projector $`𝐏_\mathrm{𝟏}`$ introduced by Zarembo, Eq. (3.11) in . Furthermore, $`𝒞^{(2)}+𝒞^{(3)}\mathrm{𝟏}𝐏_\mathrm{𝟏}𝐏_\mathrm{𝟖}`$ is Zarembo’s octet projector, Eq. (3.12) in .
With the projectors (1), the gap matrices $`\mathrm{\Phi }^\pm `$ can be written as
$$\mathrm{\Phi }^\pm \underset{n=1}{\overset{3}{}}𝒞^{(n)}\mathrm{\Phi }_n^\pm .$$
(16)
Here,
$`\mathrm{\Phi }_1^\pm `$ $``$ $`2\left(\mathrm{\Phi }_{\overline{\mathrm{𝟑}}}^\pm +2\mathrm{\Phi }_\mathrm{𝟔}^\pm \right),`$ (18)
$`\mathrm{\Phi }_2^\pm `$ $``$ $`\mathrm{\Phi }_{\overline{\mathrm{𝟑}}}^\pm \mathrm{\Phi }_\mathrm{𝟔}^\pm ,`$ (19)
$`\mathrm{\Phi }_3^\pm `$ $``$ $`\mathrm{\Phi }_2^\pm ,`$ (20)
are gap matrices in spinor space,
$$\mathrm{\Phi }_n^+(K)\underset{h=r,\mathrm{}}{}\underset{e=\pm }{}\varphi _{n,h}^e(K)𝒫_h\mathrm{\Lambda }_𝐤^e,\mathrm{\Phi }_n^{}(K)\underset{h=r,\mathrm{}}{}\underset{e=\pm }{}\left[\varphi _{n,h}^e(K)\right]^{}𝒫_h\mathrm{\Lambda }_𝐤^e,$$
(21)
where $`𝒫_{r,\mathrm{}}(1\pm \gamma _5)/2`$ are chirality projectors, $`h=\mathrm{}`$ when $`h=r`$ and $`h=r`$ when $`h=\mathrm{}`$, $`\mathrm{\Lambda }_𝐤^\pm (1\pm \gamma _0𝜸\widehat{𝐤})/2`$ are energy projectors, and $`\varphi _{n,h}^e(K)`$ are simple functions of 4-momentum $`K^\mu `$.
In Eqs. (16), $`\mathrm{\Phi }_{\overline{\mathrm{𝟑}}}^\pm `$ is the gap matrix in the antitriplet channel and $`\mathrm{\Phi }_\mathrm{𝟔}^\pm `$ is the gap matrix in the sextet channel,
$$\mathrm{\Phi }_{}^{\pm }{}_{rs}{}^{ij}\mathrm{\Phi }_{\overline{\mathrm{𝟑}}}^\pm \left(\delta _r^i\delta _s^j\delta _s^i\delta _r^j\right)+\mathrm{\Phi }_\mathrm{𝟔}^\pm \left(\delta _r^i\delta _s^j+\delta _s^i\delta _r^j\right).$$
(22)
Why condensation in the (repulsive) sextet channel is possible in the color-flavor locked phase was explained in . Antitriplet and sextet gaps are related to the gap functions $`\kappa _1`$ and $`\kappa _2`$ of by $`\varphi _{\overline{\mathrm{𝟑}}\mathrm{}}^+=\varphi _{\overline{\mathrm{𝟑}}r}^+(\kappa _1\kappa _2)/2`$, $`\varphi _\mathrm{𝟔}\mathrm{}^+=\varphi _{\mathrm{𝟔}r}^+(\kappa _1+\kappa _2)/2`$. For future purpose, it will be convenient to define a singlet and an octet gap matrix according to
$$\mathrm{\Phi }_\mathrm{𝟏}^\pm \mathrm{\Phi }_1^\pm ,\mathrm{\Phi }_\mathrm{𝟖}^\pm \mathrm{\Phi }_2^\pm \mathrm{\Phi }_3^\pm .$$
(23)
The quasiparticle propagators take the form
$$G^\pm (K)\underset{n=1}{\overset{3}{}}𝒞^{(n)}G_n^\pm (K),$$
(24)
where
$$G_n^\pm (K)=\underset{h=r,\mathrm{}}{}\underset{e=\pm }{}𝒫_{\pm h}\mathrm{\Lambda }_𝐤^{\pm e}\frac{1}{k_0^2[ϵ_𝐤^e(\varphi _{n,h}^e)]^2}\left[G_0^{}(K)\right]^1.$$
(25)
The quasiparticle energies are
$$ϵ_𝐤^e(\varphi _{n,h}^e)\sqrt{(\mu ek)^2+|\varphi _{n,h}^e|^2},$$
(26)
where $`\varphi _{n,h}^e`$ is the gap function for pairing of quarks ($`e=+1`$) or antiquarks ($`e=1`$) with chirality $`h`$.
The right-hand side of (25) does not depend on the sign of the gap functions, i.e., the difference in sign between $`\varphi _{2,h}^e`$ and $`\varphi _{3,h}^e`$, Eq. (20), is irrelevant, $`G_2^\pm G_3^\pm `$. Then, Eq. (24) has the alternative representation
$$G^\pm (K)𝐏_\mathrm{𝟏}G_\mathrm{𝟏}^\pm (K)+𝐏_\mathrm{𝟖}G_\mathrm{𝟖}^\pm (K),$$
(27)
where $`𝐏_{\mathrm{𝟏},\mathrm{𝟖}}`$ are the singlet and octet projectors introduced by Zarembo , see above, and, following Eq. (23), $`G_\mathrm{𝟏}^\pm G_1^\pm `$, $`G_\mathrm{𝟖}^\pm G_2^\pm =G_3^\pm `$.
The off-diagonal components of the quasiparticle propagators are similarly computed as
$$\mathrm{\Xi }^\pm (K)\underset{n=1}{\overset{3}{}}𝒞^{(n)}\mathrm{\Xi }_n^\pm (K),$$
(28)
where
$$\mathrm{\Xi }_n^+(K)=\underset{h=r,\mathrm{}}{}\underset{e=\pm }{}\frac{\varphi _{n,h}^e(K)}{k_0^2[ϵ_𝐤^e(\varphi _{n,h}^e)]^2}𝒫_h\mathrm{\Lambda }_𝐤^e,\mathrm{\Xi }_n^{}(K)=\underset{h=r,\mathrm{}}{}\underset{e=\pm }{}\frac{\left[\varphi _{n,h}^e(K)\right]^{}}{k_0^2[ϵ_𝐤^e(\varphi _{n,h}^e)]^2}𝒫_h\mathrm{\Lambda }_𝐤^e.$$
(29)
Since $`\varphi _{2,h}^e=\varphi _{3,h}^e`$, there is no simple representation in terms of singlet and octet projectors for $`\mathrm{\Xi }^\pm `$. Nevertheless, in line with (23) let us define for future purpose $`\mathrm{\Xi }_\mathrm{𝟏}^\pm \mathrm{\Xi }_1^\pm `$, $`\mathrm{\Xi }_\mathrm{𝟖}^\pm \mathrm{\Xi }_2^\pm \mathrm{\Xi }_3^\pm `$.
Inserting Eqs. (24) and (28) into Eq. (9), one straightforwardly performs the trace over color and flavor space to obtain
$`\mathrm{\Pi }_{ab}^{\mu \nu }(P)`$ $`=`$ $`\delta _{ab}\mathrm{\Pi }^{\mu \nu }(P),`$ (31)
$`\mathrm{\Pi }^{\mu \nu }(P)`$ $`=`$ $`{\displaystyle \frac{g^2}{12}}{\displaystyle \frac{T}{V}}{\displaystyle \underset{K}{}}\mathrm{Tr}_s[\gamma ^\mu G_\mathrm{𝟏}^+(K)\gamma ^\nu G_\mathrm{𝟖}^+(KP)+\gamma ^\mu G_\mathrm{𝟖}^+(K)\gamma ^\nu G_\mathrm{𝟏}^+(KP){\displaystyle \frac{}{}}`$ (37)
$`+\gamma ^\mu G_\mathrm{𝟏}^{}(K)\gamma ^\nu G_\mathrm{𝟖}^{}(KP)+\gamma ^\mu G_\mathrm{𝟖}^{}(K)\gamma ^\nu G_\mathrm{𝟏}^{}(KP){\displaystyle \frac{}{}}`$
$`+\mathrm{\hspace{0.33em}7}\gamma ^\mu G_\mathrm{𝟖}^+(K)\gamma ^\nu G_\mathrm{𝟖}^+(KP)+7\gamma ^\mu G_\mathrm{𝟖}^{}(K)\gamma ^\nu G_\mathrm{𝟖}^{}(KP){\displaystyle \frac{}{}}`$
$`+\gamma ^\mu \mathrm{\Xi }_\mathrm{𝟏}^{}(K)\gamma ^\nu \mathrm{\Xi }_\mathrm{𝟖}^+(KP)+\gamma ^\mu \mathrm{\Xi }_\mathrm{𝟖}^{}(K)\gamma ^\nu \mathrm{\Xi }_\mathrm{𝟏}^+(KP){\displaystyle \frac{}{}}`$
$`+\gamma ^\mu \mathrm{\Xi }_\mathrm{𝟏}^+(K)\gamma ^\nu \mathrm{\Xi }_\mathrm{𝟖}^{}(KP)+\gamma ^\mu \mathrm{\Xi }_\mathrm{𝟖}^+(K)\gamma ^\nu \mathrm{\Xi }_\mathrm{𝟏}^{}(KP){\displaystyle \frac{}{}}`$
$`+\mathrm{\hspace{0.33em}2}\gamma ^\mu \mathrm{\Xi }_\mathrm{𝟖}^{}(K)\gamma ^\nu \mathrm{\Xi }_\mathrm{𝟖}^+(KP)+2\gamma ^\mu \mathrm{\Xi }_\mathrm{𝟖}^+(K)\gamma ^\nu \mathrm{\Xi }_\mathrm{𝟖}^{}(KP){\displaystyle \frac{}{}}].`$
From (1) one learns two things. First, unlike the two-flavor case , the gluon self-energy is diagonal in the adjoint colors $`a,b`$. Second, there are no terms where both quasiparticle propagators involve singlet gaps. The reason is that such terms are proportional to $`\mathrm{Tr}_cT_a\mathrm{Tr}_cT_b0`$.
The evaluation of the spin traces proceeds in complete analogy to the two-flavor case . Assuming $`\varphi _{n,r}^e=\varphi _{n,\mathrm{}}^e\varphi _n^e𝐑`$, the result is \[cf. Eq. (96a) of \]
$`\mathrm{\Pi }^{\mu \nu }(P)={\displaystyle \frac{g^2}{12}}{\displaystyle }{\displaystyle \frac{d^3𝐤}{(2\pi )^3}}{\displaystyle \underset{e_1,e_2=\pm }{}}({\displaystyle \frac{}{}}𝒯_+^{\mu \nu }(𝐤_1,𝐤_2)`$ (38)
$`\times `$ $`[({\displaystyle \frac{\widehat{n}_1(1n_2)}{p_0+\widehat{ϵ}_1+ϵ_2}}{\displaystyle \frac{(1\widehat{n}_1)n_2}{p_0\widehat{ϵ}_1ϵ_2}})(1\widehat{N}_1N_2)+({\displaystyle \frac{(1\widehat{n}_1)(1n_2)}{p_0\widehat{ϵ}_1+ϵ_2}}{\displaystyle \frac{\widehat{n}_1n_2}{p_0+\widehat{ϵ}_1ϵ_2}})(\widehat{N}_1N_2)`$ (39)
$`+`$ $`\left({\displaystyle \frac{n_1(1\widehat{n}_2)}{p_0+ϵ_1+\widehat{ϵ}_2}}{\displaystyle \frac{(1n_1)\widehat{n}_2}{p_0ϵ_1\widehat{ϵ}_2}}\right)(1N_1\widehat{N}_2)+\left({\displaystyle \frac{(1n_1)(1\widehat{n}_2)}{p_0ϵ_1+\widehat{ϵ}_2}}{\displaystyle \frac{n_1\widehat{n}_2}{p_0+ϵ_1\widehat{ϵ}_2}}\right)(N_1\widehat{N}_2)`$ (40)
$`+`$ $`7({\displaystyle \frac{n_1(1n_2)}{p_0+ϵ_1+ϵ_2}}{\displaystyle \frac{(1n_1)n_2}{p_0ϵ_1ϵ_2}})(1N_1N_2)+7({\displaystyle \frac{(1n_1)(1n_2)}{p_0ϵ_1+ϵ_2}}{\displaystyle \frac{n_1n_2}{p_0+ϵ_1ϵ_2}})(N_1N_2)]`$ (42)
$`+𝒯_{}^{\mu \nu }(𝐤_1,𝐤_2)`$
$`\times `$ $`[({\displaystyle \frac{(1\widehat{n}_1)n_2)}{p_0+\widehat{ϵ}_1+ϵ_2}}{\displaystyle \frac{\widehat{n}_1(1n_2)}{p_0\widehat{ϵ}_1ϵ_2}})(1\widehat{N}_1N_2)+({\displaystyle \frac{\widehat{n}_1n_2}{p_0\widehat{ϵ}_1+ϵ_2}}{\displaystyle \frac{(1\widehat{n}_1)(1n_2)}{p_0+\widehat{ϵ}_1ϵ_2}})(\widehat{N}_1N_2)`$ (43)
$`+`$ $`\left({\displaystyle \frac{(1n_1)\widehat{n}_2}{p_0+ϵ_1+\widehat{ϵ}_2}}{\displaystyle \frac{n_1(1\widehat{n}_2)}{p_0ϵ_1\widehat{ϵ}_2}}\right)(1N_1\widehat{N}_2)+\left({\displaystyle \frac{n_1\widehat{n}_2}{p_0ϵ_1+\widehat{ϵ}_2}}{\displaystyle \frac{(1n_1)(1\widehat{n}_2)}{p_0+ϵ_1\widehat{ϵ}_2}}\right)(N_1\widehat{N}_2)`$ (44)
$`+`$ $`7({\displaystyle \frac{(1n_1)n_2}{p_0+ϵ_1+ϵ_2}}{\displaystyle \frac{n_1(1n_2)}{p_0ϵ_1ϵ_2}})(1N_1N_2)+7({\displaystyle \frac{n_1n_2}{p_0ϵ_1+ϵ_2}}{\displaystyle \frac{(1n_1)(1n_2)}{p_0+ϵ_1ϵ_2}})(N_1N_2)]`$ (46)
$`\left[𝒰_+^{\mu \nu }(𝐤_1,𝐤_2)+𝒰_{}^{\mu \nu }(𝐤_1,𝐤_2)\right]`$
$`\times `$ $`\{{\displaystyle \frac{\widehat{\varphi }_1\varphi _2}{4\widehat{ϵ}_1ϵ_2}}[({\displaystyle \frac{1}{p_0+\widehat{ϵ}_1+ϵ_2}}{\displaystyle \frac{1}{p_0\widehat{ϵ}_1ϵ_2}})(1\widehat{N}_1N_2)({\displaystyle \frac{1}{p_0\widehat{ϵ}_1+ϵ_2}}{\displaystyle \frac{1}{p_0+\widehat{ϵ}_1ϵ_2}})(\widehat{N}_1N_2)]`$ (47)
$`+`$ $`{\displaystyle \frac{\varphi _1\widehat{\varphi }_2}{4ϵ_1\widehat{ϵ}_2}}\left[\left({\displaystyle \frac{1}{p_0+ϵ_1+\widehat{ϵ}_2}}{\displaystyle \frac{1}{p_0ϵ_1\widehat{ϵ}_2}}\right)(1N_1\widehat{N}_2)\left({\displaystyle \frac{1}{p_0ϵ_1+\widehat{ϵ}_2}}{\displaystyle \frac{1}{p_0+ϵ_1\widehat{ϵ}_2}}\right)(N_1\widehat{N}_2)\right]`$ (48)
$`+`$ $`2{\displaystyle \frac{\varphi _1\varphi _2}{4ϵ_1ϵ_2}}[({\displaystyle \frac{1}{p_0+ϵ_1+ϵ_2}}{\displaystyle \frac{1}{p_0ϵ_1ϵ_2}})(1N_1N_2)({\displaystyle \frac{1}{p_0ϵ_1+ϵ_2}}{\displaystyle \frac{1}{p_0+ϵ_1ϵ_2}})(N_1N_2)]\}).`$ (49)
Here, I denoted the octet and singlet gaps by
$$\varphi _i\varphi _\mathrm{𝟖}^{e_i},\widehat{\varphi }_i\varphi _\mathrm{𝟏}^{e_i}.$$
(50)
Correspondingly,
$$ϵ_iϵ_{𝐤_i}^{e_i}(\varphi _i),\widehat{ϵ}_iϵ_{𝐤_i}^{e_i}(\widehat{\varphi }_i)$$
(51)
are the excitation energies for quasiparticles with octet and singlet gaps, $`𝐤_1𝐤`$, $`𝐤_2𝐤𝐩`$,
$$n_in_{𝐤_i}^{e_i}\frac{ϵ_i\xi _i}{2ϵ_i},\widehat{n}_i\widehat{n}_{𝐤_i}^{e_i}\frac{\widehat{ϵ}_i\xi _i}{2\widehat{ϵ}_i}$$
(52)
are the occupation numbers for quasiparticles with octet and singlet gaps, $`\xi _ie_ik_i\mu `$, and
$$N_iN_{𝐤_i}^{e_i}[\mathrm{exp}\left(\frac{ϵ_i}{T}\right)+1]^1,\widehat{N}_i\widehat{N}_{𝐤_i}^{e_i}[\mathrm{exp}\left(\frac{\widehat{ϵ}_i}{T}\right)+1]^1$$
(53)
are the corresponding thermal occupation numbers. The spin traces are \[see Eqs. (45) and (97) of \]
$`𝒯_\pm ^{00}`$ $`=`$ $`𝒰_\pm ^{00}=1+e_1e_2\widehat{𝐤}_1\widehat{𝐤}_2,`$ (55)
$`𝒯_\pm ^{0i}`$ $`=`$ $`𝒯_\pm ^{i0}=𝒰_\pm ^{0i}=𝒰_\pm ^{i0}=\pm e_1\widehat{k}_1^i\pm e_2\widehat{k}_2^i,i=x,y,z,`$ (56)
$`𝒯_\pm ^{ij}`$ $`=`$ $`𝒰_\pm ^{ij}=\delta ^{ij}\left(1e_1e_2\widehat{𝐤}_1\widehat{𝐤}_2\right)+e_1e_2\left(\widehat{k}_1^i\widehat{k}_2^j+\widehat{k}_1^j\widehat{k}_2^i\right),i,j=x,y,z.`$ (57)
In Eq. (38), the terms proportional to $`𝒯_\pm ^{\mu \nu }`$ correspond to the diagram in Fig. 1(a), while those proportional to $`𝒰_\pm ^{\mu \nu }`$ correspond to that in Fig. 1(b).
In order to compute the Debye and Meissner mass, it is sufficient to consider the time-like, $`\mu =\nu =0`$, and space-like, $`\mu =i,\nu =j`$, components of the self-energy, since the Debye and Meissner masses are defined as
$$m_D^2\underset{p0}{lim}\mathrm{\Pi }^{00}(0,p),m_M^2\underset{p0}{lim}\mathrm{\Pi }^{ii}(0,p).$$
(58)
The sign in the first equation is due to the choice of metric. The self-energy of electric gluons is
$`\mathrm{\Pi }^{00}(P)`$ $`=`$ $`{\displaystyle \frac{g^2}{12}}{\displaystyle \frac{d^3𝐤}{(2\pi )^3}\underset{e_1,e_2=\pm }{}(1+e_1e_2\widehat{𝐤}_1\widehat{𝐤}_2)}`$ (66)
$`\times [({\displaystyle \frac{1}{p_0+\widehat{ϵ}_1+ϵ_2}}{\displaystyle \frac{1}{p_0\widehat{ϵ}_1ϵ_2}})(1\widehat{N}_1N_2){\displaystyle \frac{\widehat{ϵ}_1ϵ_2\xi _1\xi _2\widehat{\varphi }_1\varphi _2}{2\widehat{ϵ}_1ϵ_2}}`$
$`+\left({\displaystyle \frac{1}{p_0+ϵ_1+\widehat{ϵ}_2}}{\displaystyle \frac{1}{p_0ϵ_1\widehat{ϵ}_2}}\right)(1N_1\widehat{N}_2){\displaystyle \frac{ϵ_1\widehat{ϵ}_2\xi _1\xi _2\varphi _1\widehat{\varphi }_2}{2ϵ_1\widehat{ϵ}_2}}`$
$`+\mathrm{\hspace{0.17em}7}\left({\displaystyle \frac{1}{p_0+ϵ_1+ϵ_2}}{\displaystyle \frac{1}{p_0ϵ_1ϵ_2}}\right)(1N_1N_2){\displaystyle \frac{ϵ_1ϵ_2\xi _1\xi _22\varphi _1\varphi _2/7}{2ϵ_1ϵ_2}}`$
$`+\left({\displaystyle \frac{1}{p_0\widehat{ϵ}_1+ϵ_2}}{\displaystyle \frac{1}{p_0+\widehat{ϵ}_1ϵ_2}}\right)(\widehat{N}_1N_2){\displaystyle \frac{\widehat{ϵ}_1ϵ_2+\xi _1\xi _2+\widehat{\varphi }_1\varphi _2}{2\widehat{ϵ}_1ϵ_2}}`$
$`+\left({\displaystyle \frac{1}{p_0ϵ_1+\widehat{ϵ}_2}}{\displaystyle \frac{1}{p_0+ϵ_1\widehat{ϵ}_2}}\right)(N_1\widehat{N}_2){\displaystyle \frac{ϵ_1\widehat{ϵ}_2+\xi _1\xi _2+\varphi _1\widehat{\varphi }_2}{2ϵ_1\widehat{ϵ}_2}}`$
$`+7({\displaystyle \frac{1}{p_0ϵ_1+ϵ_2}}{\displaystyle \frac{1}{p_0+ϵ_1ϵ_2}})(N_1N_2){\displaystyle \frac{ϵ_1ϵ_2+\xi _1\xi _2+2\varphi _1\varphi _2/7}{2ϵ_1ϵ_2}}].`$
On the other hand, the self-energy of magnetic gluons is
$`\mathrm{\Pi }^{ij}(P)`$ $`=`$ $`{\displaystyle \frac{g^2}{12}}{\displaystyle \frac{d^3𝐤}{(2\pi )^3}\underset{e_1,e_2=\pm }{}\left[\delta ^{ij}\left(1e_1e_2\widehat{𝐤}_1\widehat{𝐤}_2\right)+e_1e_2\left(\widehat{k}_1^i\widehat{k}_2^j+\widehat{k}_1^j\widehat{k}_2^i\right)\right]}`$ (73)
$`\times [({\displaystyle \frac{1}{p_0+\widehat{ϵ}_1+ϵ_2}}{\displaystyle \frac{1}{p_0\widehat{ϵ}_1ϵ_2}})(1\widehat{N}_1N_2){\displaystyle \frac{\widehat{ϵ}_1ϵ_2\xi _1\xi _2+\widehat{\varphi }_1\varphi _2}{2\widehat{ϵ}_1ϵ_2}}`$
$`+\left({\displaystyle \frac{1}{p_0+ϵ_1+\widehat{ϵ}_2}}{\displaystyle \frac{1}{p_0ϵ_1\widehat{ϵ}_2}}\right)(1N_1\widehat{N}_2){\displaystyle \frac{ϵ_1\widehat{ϵ}_2\xi _1\xi _2+\varphi _1\widehat{\varphi }_2}{2ϵ_1\widehat{ϵ}_2}}`$
$`+\mathrm{\hspace{0.17em}7}\left({\displaystyle \frac{1}{p_0+ϵ_1+ϵ_2}}{\displaystyle \frac{1}{p_0ϵ_1ϵ_2}}\right)(1N_1N_2){\displaystyle \frac{ϵ_1ϵ_2\xi _1\xi _2+2\varphi _1\varphi _2/7}{2ϵ_1ϵ_2}}`$
$`+\left({\displaystyle \frac{1}{p_0\widehat{ϵ}_1+ϵ_2}}{\displaystyle \frac{1}{p_0+\widehat{ϵ}_1ϵ_2}}\right)(\widehat{N}_1N_2){\displaystyle \frac{\widehat{ϵ}_1ϵ_2+\xi _1\xi _2\widehat{\varphi }_1\varphi _2}{2\widehat{ϵ}_1ϵ_2}}`$
$`+\left({\displaystyle \frac{1}{p_0ϵ_1+\widehat{ϵ}_2}}{\displaystyle \frac{1}{p_0+ϵ_1\widehat{ϵ}_2}}\right)(N_1\widehat{N}_2){\displaystyle \frac{ϵ_1\widehat{ϵ}_2+\xi _1\xi _2\varphi _1\widehat{\varphi }_2}{2ϵ_1\widehat{ϵ}_2}}`$
$`+7({\displaystyle \frac{1}{p_0ϵ_1+ϵ_2}}{\displaystyle \frac{1}{p_0+ϵ_1ϵ_2}})(N_1N_2){\displaystyle \frac{ϵ_1ϵ_2+\xi _1\xi _22\varphi _1\varphi _2/7}{2ϵ_1ϵ_2}}].`$
In Eq. (58), terms proportional to $`(ϵ_1ϵ_2\pm \xi _1\xi _2)/(2ϵ_1ϵ_2)`$ (and similar terms involving the singlet gaps, $`\widehat{\varphi }_i`$) arise from the diagram in Fig. 1(a), while terms proportional to $`\varphi _1\varphi _2/(2ϵ_1ϵ_2)`$ (and similar terms involving the singlet gaps, $`\widehat{\varphi }_i`$) arise from that in Fig. 1(b).
To proceed I treat quasi-antiparticles as free antiparticles, as in ,
$`\varphi _\mathrm{𝟖}^{}0,\varphi _\mathrm{𝟏}^{}0,ϵ_{𝐤_i}^{}(\varphi _\mathrm{𝟖}^{})k_i+\mu ,ϵ_{𝐤_i}^{}(\varphi _\mathrm{𝟏}^{})k_i+\mu ,`$ (75)
$`n_{𝐤_i}^{}1,\widehat{n}_{𝐤_i}^{}1,\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}1}n_{𝐤_i}^{}0,\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}1}\widehat{n}_{𝐤_i}^{}0,N_{𝐤_i}^{}0,\widehat{N}_{𝐤_i}^{}0.`$ (76)
Setting $`p_0=0`$ and sending $`p0`$, there are then no antiparticle contributions to the electric part of the gluon self-energy. Abbreviating $`\varphi _\mathrm{𝟖}^+\varphi `$, $`\varphi _\mathrm{𝟏}^+\widehat{\varphi }`$, $`ϵ_𝐤^+ϵ`$, $`\widehat{ϵ}_𝐤^+\widehat{ϵ}`$, $`N_𝐤^+N`$, $`\widehat{N}_𝐤^+\widehat{N}`$,
$`\mathrm{\Pi }^{00}(0)`$ $``$ $`\mathrm{\Pi }_\mathrm{e}^{(a)}(0)+\mathrm{\Pi }_\mathrm{e}^{(b)}(0),`$ (78)
$`\mathrm{\Pi }_\mathrm{e}^{(a)}(0)`$ $``$ $`{\displaystyle \frac{g^2}{24\pi ^2}}{\displaystyle _0^{\mathrm{}}}dkk^2[8{\displaystyle \frac{1\widehat{N}N}{\widehat{ϵ}+ϵ}}{\displaystyle \frac{\widehat{ϵ}ϵ\xi ^2}{2\widehat{ϵ}ϵ}}+7(12N){\displaystyle \frac{\varphi ^2}{ϵ^3}}`$ (80)
$`8{\displaystyle \frac{N\widehat{N}}{ϵ\widehat{ϵ}}}{\displaystyle \frac{\widehat{ϵ}ϵ+\xi ^2}{2\widehat{ϵ}ϵ}}28{\displaystyle \frac{\mathrm{d}N}{\mathrm{d}ϵ}}(1{\displaystyle \frac{\varphi ^2}{2ϵ^2}})],`$
$`\mathrm{\Pi }_\mathrm{e}^{(b)}(0)`$ $``$ $`{\displaystyle \frac{g^2}{24\pi ^2}}{\displaystyle _0^{\mathrm{}}}dkk^2[\mathrm{\hspace{0.17em}8}{\displaystyle \frac{1\widehat{N}N}{\widehat{ϵ}+ϵ}}{\displaystyle \frac{\widehat{\varphi }\varphi }{2\widehat{ϵ}ϵ}}2(12N){\displaystyle \frac{\varphi ^2}{ϵ^3}}`$ (82)
$`8{\displaystyle \frac{N\widehat{N}}{ϵ\widehat{ϵ}}}{\displaystyle \frac{\widehat{\varphi }\varphi }{2\widehat{ϵ}ϵ}}4{\displaystyle \frac{\mathrm{d}N}{\mathrm{d}ϵ}}{\displaystyle \frac{\varphi ^2}{ϵ^2}}].`$
In order to facilitate comparison with the results of Son and Stephanov, and to elucidate the origin of the various terms, I have separated the self-energy into the contributions from the diagram in Fig. 1(a), $`\mathrm{\Pi }_\mathrm{e}^{(a)}`$, and that in Fig. 1(b), $`\mathrm{\Pi }_\mathrm{e}^{(b)}`$.
For the self-energy of magnetic gluons one obtains with $`𝑑\mathrm{\Omega }\widehat{k}^i\widehat{k}^j/(4\pi )=\delta ^{ij}/3`$, $`n_𝐤^+n`$, $`\widehat{n}_𝐤^+\widehat{n}`$,
$`\mathrm{\Pi }^{ij}(0)`$ $``$ $`\delta ^{ij}\left[\mathrm{\Pi }_\mathrm{m}^{(a1)}(0)+\mathrm{\Pi }_\mathrm{m}^{(a2)}(0)+\mathrm{\Pi }_\mathrm{m}^{(b)}(0)\right],`$ (83)
$`\mathrm{\Pi }_\mathrm{m}^{(a1)}(0)`$ $``$ $`{\displaystyle \frac{g^2}{72\pi ^2}}{\displaystyle _0^{\mathrm{}}}dkk^2[8{\displaystyle \frac{1\widehat{N}N}{\widehat{ϵ}+ϵ}}{\displaystyle \frac{\widehat{ϵ}ϵ\xi ^2}{2\widehat{ϵ}ϵ}}+7(12N){\displaystyle \frac{\varphi ^2}{ϵ^3}}`$ (85)
$`8{\displaystyle \frac{N\widehat{N}}{ϵ\widehat{ϵ}}}{\displaystyle \frac{\widehat{ϵ}ϵ+\xi ^2}{2\widehat{ϵ}ϵ}}28{\displaystyle \frac{\mathrm{d}N}{\mathrm{d}ϵ}}(1{\displaystyle \frac{\varphi ^2}{2ϵ^2}})],`$
$`\mathrm{\Pi }_\mathrm{m}^{(a2)}(0)`$ $``$ $`{\displaystyle \frac{g^2}{72\pi ^2}}{\displaystyle _0^{\mathrm{}}}dkk^2[16{\displaystyle \frac{(1\widehat{N})(1\widehat{n})}{k+\mu +\widehat{ϵ}}}+128{\displaystyle \frac{(1N)(1n)}{k+\mu +ϵ}}`$ (87)
$`+\mathrm{\hspace{0.17em}16}{\displaystyle \frac{\widehat{N}\widehat{n}}{k+\mu \widehat{ϵ}}}+128{\displaystyle \frac{Nn}{k+\mu ϵ}}72{\displaystyle \frac{1}{k}}],`$
$`\mathrm{\Pi }_\mathrm{m}^{(b)}(0)`$ $``$ $`{\displaystyle \frac{g^2}{72\pi ^2}}{\displaystyle _0^{\mathrm{}}}dkk^2[8{\displaystyle \frac{1\widehat{N}N}{\widehat{ϵ}+ϵ}}{\displaystyle \frac{\widehat{\varphi }\varphi }{2\widehat{ϵ}ϵ}}+2(12N){\displaystyle \frac{\varphi ^2}{ϵ^3}}`$ (89)
$`+8{\displaystyle \frac{N\widehat{N}}{ϵ\widehat{ϵ}}}{\displaystyle \frac{\widehat{\varphi }\varphi }{2\widehat{ϵ}ϵ}}+4{\displaystyle \frac{\mathrm{d}N}{\mathrm{d}ϵ}}{\displaystyle \frac{\varphi ^2}{ϵ^2}}].`$
Again, I have separated contributions from Fig. 1(a), $`\mathrm{\Pi }_\mathrm{m}^{(a1)}+\mathrm{\Pi }_\mathrm{m}^{(a2)}`$, from those of Fig. 1(b), $`\mathrm{\Pi }_\mathrm{m}^{(b)}`$. Obviously, $`\mathrm{\Pi }_\mathrm{e}^{(a)}3\mathrm{\Pi }_\mathrm{m}^{(a1)}`$, $`\mathrm{\Pi }_\mathrm{e}^{(b)}3\mathrm{\Pi }_\mathrm{m}^{(b)}`$. There are two contributions from Fig. 1(a). The first, $`\mathrm{\Pi }_\mathrm{m}^{(a1)}`$, arises from quasiparticle-quasiparticle excitations, while the second, $`\mathrm{\Pi }_\mathrm{m}^{(a2)}`$, originates from quasiparticle-antiparticle excitations. The latter is UV-divergent, and thus requires renormalization, which is achieved by adding the last term in Eq. (87). As we shall see shortly, $`\mathrm{\Pi }_\mathrm{m}^{(a2)}`$ gives rise to the “bare” Meissner mass in Hong’s effective theory , where contributions involving antiparticles are integrated out.
Let us now consider the case of zero temperature, where $`N\widehat{N}0`$. One may restrict the $`k`$ integration to the region $`0k2\mu `$, the contribution from $`k2\mu `$ can be shown to be negligible. The various components of the self-energies are
$`\mathrm{\Pi }_\mathrm{e}^{(a)}(0)`$ $``$ $`3\mathrm{\Pi }_\mathrm{m}^{(a1)}(0){\displaystyle \frac{g^2\mu ^2}{12\pi ^2}}{\displaystyle _0^\mu }d\xi \left[{\displaystyle \frac{4}{\widehat{\varphi }^2\varphi ^2}}\left({\displaystyle \frac{\widehat{ϵ}^2+\xi ^2}{\widehat{ϵ}}}{\displaystyle \frac{ϵ^2+\xi ^2}{ϵ}}\right)+7{\displaystyle \frac{\varphi ^2}{ϵ^3}}\right],`$ (91)
$`\mathrm{\Pi }_\mathrm{e}^{(b)}(0)`$ $``$ $`3\mathrm{\Pi }_\mathrm{m}^{(b)}(0){\displaystyle \frac{g^2\mu ^2}{12\pi ^2}}{\displaystyle _0^\mu }d\xi \left[{\displaystyle \frac{4\widehat{\varphi }\varphi }{\widehat{\varphi }^2\varphi ^2}}\left({\displaystyle \frac{1}{ϵ}}{\displaystyle \frac{1}{\widehat{ϵ}}}\right)2{\displaystyle \frac{\varphi ^2}{ϵ^3}}\right],`$ (92)
$`\mathrm{\Pi }_\mathrm{m}^{(a2)}(0)`$ $``$ $`{\displaystyle \frac{g^2}{72\pi ^2}}{\displaystyle _0^{2\mu }}dkk\left[8{\displaystyle \frac{\mu (\widehat{ϵ}k+\mu )+\widehat{\varphi }^2}{\widehat{ϵ}(\widehat{ϵ}+k+\mu )}}+64{\displaystyle \frac{\mu (ϵk+\mu )+\varphi ^2}{ϵ(ϵ+k+\mu )}}\right].`$ (93)
The integral appearing in (93) was already computed in , Eq. (122). The result is
$$\mathrm{\Pi }_\mathrm{m}^{(a2)}(0)m_g^2.$$
(94)
As advertised above, this is the “bare” Meissner mass appearing in Hong’s effective theory .
To obtain the expressions (91) and (92) for $`\mathrm{\Pi }_\mathrm{e}^{(a)}`$ and $`\mathrm{\Pi }_\mathrm{e}^{(b)}`$, I substituted $`\xi k\mu `$ and exploited the symmetry of the integrand around $`\xi =0`$. Furthermore, contributions $`\xi ^2/\mu ^2`$ in the integrands were neglected, because they give rise to terms of order $`\varphi ^2/\mu ^2`$ relative to the leading terms. Neglecting the momentum dependence of the gap function, all remaining integrals are exactly solvable. First note that, to leading order, $`_0^\mu d\xi \varphi ^2/ϵ^31`$. This takes care of the last term in Eqs. (91) and (92). To compute the first, substitute $`y\mathrm{ln}[(\xi +ϵ)/\varphi ]`$ for $`\xi `$ in the first term in parentheses, and $`\widehat{y}\mathrm{ln}[(\xi +\widehat{ϵ})/\widehat{\varphi }]`$ in the second. Evaluating the $`y`$ integral, note that one must not approximate $`\mathrm{ln}[(\mu +\sqrt{\mu ^2+\varphi ^2})/\varphi ]\mathrm{ln}(2\mu /\varphi )`$ for the upper boundary of the $`y`$ integral, and similarly for the $`\widehat{y}`$ integral. The reason is that leading-order terms cancel between these two integrals. The subleading terms conspire to cancel the denominator $`\widehat{\varphi }^2\varphi ^2`$. To leading order, one then obtains
$$\mathrm{\Pi }_\mathrm{e}^{(a)}(0)\frac{3}{2}m_g^2,\mathrm{\Pi }_\mathrm{e}^{(b)}(0)\frac{1}{3}m_g^2(1+\frac{2\widehat{\varphi }\varphi }{\widehat{\varphi }^2\varphi ^2}\mathrm{ln}\frac{\widehat{\varphi }}{\varphi }).$$
(95)
Neglecting the sextet gap, the singlet gap is twice the octet gap, $`\widehat{\varphi }=2\varphi `$, cf. Eq. (16), and
$$\mathrm{\Pi }_\mathrm{e}^{(b)}(0)\frac{1}{3}m_g^2\left(1+\frac{4}{3}\mathrm{ln}2\right).$$
(96)
Using the definitions of the Debye and Meissner masses, Eq. (58), this confirms the results of Son and Stephanov, Eq. (5), q.e.d.
## Acknowledgements
I thank R. Pisarski, T. Schäfer, and D.T. Son for discussions. My thanks go to RIKEN, BNL and the U.S. Dept. of Energy for providing the facilities essential for the completion of this work, and to Columbia University’s Nuclear Theory Group for continuing access to their computing facilities.
|
warning/0003/physics0003051.html
|
ar5iv
|
text
|
# SYMMETRY BREAKING of VIBRATING INTERFACES: A MECHANISM for MORPHOGENESIS
## Abstract
We show that very small-amplitude oscillations of a highly symmetric, spheric or cylindrical, interface (thin membrane) between two fluids can result in inhomogeneous instability and breaking of the interface symmetry: the frequency of the breathing vibration selects the spatial symmetry. This mechanism may govern morphogenesis.
The nature of spontaneous symmetry breaking remains one of the most enigmatic questions of modern science. This problem emerges in connection with the equilibrium phase transitions, self-organization in nonequilibrium systems and many other areas in physics, chemistry and biology (see, e.g., ), as well as with cell fission and morphogenesis, i.e., the development and spatial differentiation of complex body structures during their growth .
In 1952 Turing showed that the homogeneous state of some specific chemical reactions can lose stability with regard to a spontaneous increase of perturbations of certain form . Since then the chemical basis is the prevalent idea of phenomenological theory of morphogenesis (see, e.g., ). Turing’s model is based on chemical or biological processes of local self-reproduction of some chemical agent (the activator) and far-ranging inhibition. As a consequence of such processes a very small increase of the activator concentration in a local region results in a global redistribution of the substance concentrations and formation of more complex structure . However, the Turing’s chemical reactions are uncommon, unique and very complex processes.
In this work we develop a new mechanism, without complexity, that breaks the symmetry by creating an instability in an oscillating interface, thin membrane, separating two different fluids. In other words, we show that if, for example, a spherical or cylindrical structure vibrates with a breathing symmetric mode for a given set of the frequencies the symmetry breaks with respect to bimodal, trimodal, pentagonal, etc. modes, i.e., the vibration frequency selects the spatial symmetry of the interface.
We consider a thin symmetric membrane, spherical or cylindrical interface, with the radius $`R_0`$ which separates two fluids with densities $`\rho _1`$and $`\rho _2`$ ($`\rho _1`$ $`\rho _2\rho `$) respectively. Owing to Archimed’s force the effective gravity acceleration operating on the internal fluid is $`g=g_e(1\rho _1/\rho _2)<<g_e.`$ We propose $`R_0`$ is small enough, so the condition $`\gamma k_m^2/\rho =\gamma m^2/R_0^2\rho >>g`$ is valid. Here $`\gamma `$ is the surface tension and $`k=m/R_0`$ is the typical wave vector of the increasing deformation of the symmetric interface, $`m=1,2,3,\mathrm{}`$. This is the condition when we can neglect the gravity and consider only the effect due to the surface tension of the interface.
Let us take, at first, for definiteness, a spherically symmetric interface $`𝐒`$ whose radius, $`R`$, oscillates with a frequency $`\omega :`$ $`R=R_0d\mathrm{cos}\omega t.`$ From the incompressibility of the fluid it follows that its radial velocity is $`v_{r0}=v_{R0}(t)R_0^2/r^2,`$ where $`v_{R0}(t)=dR/dt=d\omega \mathrm{sin}\omega t`$. (This means that some source, for example, a small pulsating ball has to be inside the interface.) The vortex-free motion of an ideal liquid (we consider the effect of the viscosity below) is described by the Euler and the continuity equations:
$$\frac{d𝐯}{dt}=\frac{𝐯}{t}+(𝐯)𝐯=\frac{1}{\rho }p,$$
(1)
$$^2\mathrm{\Phi }=_r^2\mathrm{\Phi }+_{}^2\mathrm{\Phi }=0$$
(2)
where $`\mathrm{\Phi }`$ is the velocity potential, $`𝐯=\mathrm{\Phi }`$ and $`_{}^2`$ is the part of Laplacian depending only on coordinates of the surface $`𝐒`$. For the undistorted spherical surface, from the symmetry of the problem, it follows that $`𝐯_{}=0,`$ i.e., $`_{}^2\mathrm{\Phi }=0.`$ Then, from Eq.2, we can write that $`_r^2\mathrm{\Phi }=_r(v_{r0})=v_{r0}/r+2v_{r0}/r=0,`$ in accord with $`v_{r0}=v_{R0}(t)R_0^2/r^2.`$ In the presence of a distortion, $`\varsigma `$, of the spherical surface $`𝐒`$ the interface radial velocity is $`v_r=v_{R0}(t)+\varsigma /t.`$ Using this, we find from Eq.1 that near the interface
$$\frac{dv_r}{dt}=F(t)+^2\varsigma /t^2=\frac{1}{\rho }\frac{p}{r},$$
(3)
$$\frac{𝐯_{}}{t}=\frac{1}{\rho }_{}p$$
(4)
where $`F(t)=d\omega ^2\mathrm{cos}\omega t`$ is the acceleration of the interface and we neglect the term $`(𝐯_{})𝐯_{}`$ in Eq. 4 by virtue of smallness of $`\varsigma `$ . Owing to smallness of $`\varsigma `$ we can write the pressure near the surface as
$$p=\rho F(t)(rR\varsigma )+\gamma _{}^2\varsigma +p_o(t)$$
(5)
Here we took into account that the pressure at the interface (when $`r=R+\varsigma )`$ is $`p=\gamma (\sigma _1+\sigma _2)+p_o(t)`$ where $`\sigma _1`$ and $`\sigma _2`$ are the principal curvatures of the interface : $`(\sigma _1+\sigma _2)=_{}^2\varsigma `$ since $`_{}^2\varsigma >R_0^1`$. Substituting Eq.5 into Eq.4 we obtain
$$\frac{𝐯_{}}{t}=F(t)_{}\varsigma \frac{\gamma }{\rho }_{}^3\varsigma \text{ or }\frac{}{t}_{}^2\mathrm{\Phi }=F(t)_{}^2\varsigma \frac{\gamma }{\rho }_{}^4\varsigma $$
(6)
We will seek solutions of the problem in the following form
$$\zeta =\underset{m=0}{\overset{\mathrm{}}{}}a_m(t)S_m\text{ and }\mathrm{\Phi }=\underset{m=0}{\overset{\mathrm{}}{}}c_m(t)\mathrm{\Psi }_m(r)S_mv_{R0}(t)R_0^2/r$$
(7)
where $`S_m`$ is the complete orthogonal set of eigenfunctions depending only on the coordinates of the undisturbed surface $`𝐒`$ and satisfying the following equation
$$(_{}^2+k_m^2)S_m=0$$
(8)
for $`r=R_0`$ and the boundary conditions corresponding to the symmetry of the problem. In the spherical case $`S_m=C_{l,m}P_l^m(\mathrm{cos}\theta )\mathrm{exp}(im\phi )`$ are the spherical functions of angles $`\phi `$ and $`\theta `$ and $`k_m^2=l(l+1)R_0^2`$ where $`m=l,l1,\mathrm{},l`$ and $`l=0,1,2,..`$. Substituting $`\mathrm{\Phi }`$ from Eq.7 into Eq.2, using Eq.8 and the condition $`_r(v_{r0})=0`$ cited above, we obtain the equation for $`\mathrm{\Psi }_m(r):`$
$$(_r^2k_m^2)\mathrm{\Psi }_m(r)=0$$
(9)
with the boundary conditions $`_r\mathrm{\Psi }_m0`$ when $`r0`$ and $`\mathrm{\Psi }_m(r)=A`$ at $`r=R_0`$ where $`A`$ is some constant which does not reveal itself in the final results. Near the interface $`_r\mathrm{\Phi }=v_r=v_{R0}(t)+\varsigma /t`$ and so from Eq.7 it follows that $`c_m(t)=da_m/dt(_r\mathrm{\Psi }_m)_{r=R_0}^1.`$ Substituting $`\mathrm{\Phi }`$ from Eq.7 into Eq.2 and using Eq.9 and $`c_m(t),`$ we find that
$$_{}^2\mathrm{\Phi }=\underset{m=0}{\overset{\mathrm{}}{}}k_m^2\varkappa _m^1S_mda_m/dt$$
(10)
where $`\varkappa _m=[_r\mathrm{\Psi }_m/\mathrm{\Psi }_m(r)]_{r=R_0}`$ does not depend on the constant $`A`$ . Then from Eq.10 and Eq.6, we obtain
$$d^2a_m/dt^2+[\gamma k_m^2\varkappa _m\rho ^1\varkappa _mF(t)]a_m=0.$$
(11)
Using $`T=\omega t/2`$ we can rewrite Eq.11 as
$$d^2a_m/dT^2+(p_m2q_m\mathrm{cos}\omega t)a_m=0,$$
(12)
where
$$q_m=2\varkappa _md\text{ and }p_m=\mathrm{\Omega }_m^2\omega ^2\text{ where }\mathrm{\Omega }_m^2=4k_m^2\varkappa _m\gamma \rho ^1.$$
(13)
For the spherical interface $`k_mR_0>1`$ and $`\varkappa _mk_m=`$ $`[l(l+1)]^{1/2}R_0^1`$ and so $`\mathrm{\Omega }_m^2=4[l(l+1)]^{3/2}R_0^3\gamma \rho ^1`$ and $`q_m=2d[l(l+1)]^{1/2}R_0^1.`$
These results can be extended easily to other cases. For example, when the interface have a form of a cylinder with vibrating radius, then $`S_m=\mathrm{cos}(k_lz)\mathrm{exp}(im\phi )`$ and in Eq.13 $`\varkappa _mk_m`$ and $`k_m^2=m^2/R_0^2+\pi ^2l^2/h_0^2`$ where $`h_0`$ is height of the cylinder. This vibrating cylindrical body can spontaneously distort in the axis $`z`$ or with respect the azimuthal perturbations.
We emphasize that Eq.12 coincides with Eq.(2.12) of Ref. to describe the Faraday’s instability of the plane free surface of an ideal liquid under vertical periodic vibrations. These equations differ in the values of the parameters $`p_m`$ and $`q_m`$. Moreover, in contrast to the Faraday’s instability when the vibrations are reduced to trivial renormalization of the gravity, in this work we consider spherical or cylindrical oscillating interfaces when the vertical direction, axial gravity, is not distinguished from other directions. Benjamin and Ursell have constructed the stability diagram for Eq.12 with respect to the universal parameters $`p_m`$ and $`q_m`$ using the analogy between Eq.12 and the Mathieu’s equation . From this diagram it follows that the instability is realized only in regions near the points $`p_m=n^2`$ where $`n=1,2,3,4,\mathrm{}`$. In other words, the condition
$$\omega =\omega _{n,m}n^1\mathrm{\Omega }_m=2n^1k_m(\varkappa _m\gamma /\rho )^{1/2}$$
(14)
determines the resonant vibration frequencies when the symmetric interface spontaneously deforms with respect to the standing wave with the azimuthal number $`m`$. However, the greater is $`n`$, the narrower is the width $`E_g^{(n)}(q_m)`$ of the $`n`$-th region of the instability for given $`q_m`$ . For the widest instability region, with $`n=1,`$ the value $`E_g^{(1)}(q_m)2q_m`$ for $`q_m<1`$. It means that the instability takes place when $`(1q_m)<p_m<(1+q_m),`$ i.e., the symmetry breaking is realized for the vibration frequency lying within the spectral range:
$$\mathrm{\Omega }_m(1\varkappa _md)<\omega <\mathrm{\Omega }_m(1\varkappa _md).$$
(15)
The threshold of the vibration amplitude $`d`$ is limited by the fluid viscosity. For real fluid Eqs.11 and 12 include the additional terms $`\gamma _mda_m(t)/dt`$ and $`\mathrm{\Gamma }_mda_m(T)/dT`$, respectively, where $`\gamma _m=2\nu k_m^2C_{1m}`$ and $`\mathrm{\Gamma }_m=4\nu k_m^2C_{1m}/\omega `$ are proportional to the kinematic viscosity $`\nu `$ and $`C_m`$ is some constant of the order of unity . The threshold vibration amplitude, $`d=d_t,`$ for the instability region with $`n=1,`$ can be estimated from the condition $`E_g^{(1)}(q_m)>2\mathrm{\Gamma }_m`$, i.e., $`q_m=2\varkappa _md>\mathrm{\Gamma }_m.`$ This condition follows practically from results of Refs. and is obtained in . Using Eq.15, this condition can be written as
$$d>d_t=2\nu C_mk_m^2/\varkappa _m\omega \nu (\rho /\gamma \varkappa _m)\nu (\rho R_0/\gamma m).$$
(16)
For parameters of water $`d_t4\mu m`$, i.e., the threshold vibration is a very small flutter of the interface.
We propose that the results above may be used as a basis for a simple, without complexity, mechanism to trigger the fanciful morphogenesis appearing in nature. The frequency of homogeneous interface vibrations self-selects the space symmetry. If the interface oscillates with a characteristic frequency the germ symmetry will break when its radius $`R_0`$ amounts to the quantity satisfied by Eq.15 or Eq.14 for $`n>1`$. After the new symmetry appears the growth rate increases with surface curvature as is usual for many of Stefan-like problems .
The mathematical results reported here will be applied in a forthcoming paper to explain the morphogesis of the acetabularia, equinoderms and cell fision.
This work has been supported by the Spanish DGCIYT and by a NATO fellowship Grant.
|
warning/0003/gr-qc0003037.html
|
ar5iv
|
text
|
# Second post-Newtonian radiative evolution of the relative orientations of angular momenta in spinning compact binaries
## I Introduction
Ongoing theoretical studies of the gravitational radiation emitted by compact binaries and of the radiation reaction on the orbit are strongly motivated by the desire to detect gravitational waves using the LIGO / VIRGO interferometers. In order to extract the signal from the experimental data one needs templates of sufficient accuracy. Based on the expectation that during the adiabatic regime the radiation reaction will circularize the orbit , templates for circular orbits are computed. If the orbits are eccentric, the detectability of the gravitational waves by using circular templates decreased . Quinlan and Shapiro and Hills and Bender argue for a significant number of eccentric binaries in galactic nuclei for which the time before plunging is insufficient for circularization. A faithful description of the behavior of these binaries in the adiabatic regime requires not only high post-Newtonian orders, but also eccentric orbits. Such a generic treatment, valid up to the second post-Newtonian order was provided by Gopakumar and Iyer .
If the spins of the neutron stars/black holes which form the binary are equally considered, the number of kinematical variables increases significantly and the dynamics complicates. Currently a second post-Newtonian order accurate description of the motion of the spinning binary - is available. Spin-orbit and spin-spin type contributions appear at the $`3/2`$ and second post-Newtonian orders, respectively. There is one notable exception over this rule: the spin-orbit and spin-spin terms in the spin precession equations
$`\dot{𝐒}_\mathrm{𝟏}`$ $`=`$ $`{\displaystyle \frac{G}{c^2r^3}}\left({\displaystyle \frac{4m_1+3m_2}{2m_1}}𝐋_𝐍𝐒_\mathrm{𝟐}+{\displaystyle \frac{3}{r^2}}(𝐫𝐒_\mathrm{𝟐})𝐫\right)\times 𝐒_\mathrm{𝟏},`$ (1)
$`\dot{𝐒}_\mathrm{𝟐}`$ $`=`$ $`{\displaystyle \frac{G}{c^2r^3}}\left({\displaystyle \frac{4m_2+3m_1}{2m_2}}𝐋_𝐍𝐒_\mathrm{𝟏}+{\displaystyle \frac{3}{r^2}}(𝐫𝐒_\mathrm{𝟏})𝐫\right)\times 𝐒_\mathrm{𝟐}`$ (2)
generate first and $`3/2`$ post-Newtonian order changes in the angles between the angular momenta vectors <sup>*</sup><sup>*</sup>* We evaluate the post-Newtonian order of a correction $`\delta f`$ to some quantity $`f`$ through the following prescription. Count the powers of $`ϵGm/rc^2v^2/c^2`$, representing one post-Newtonian order, in both $`\delta f`$ and $`f`$. ($`G`$ is the gravitational constant, $`c`$ the speed of light, $`m`$ the mass of binary, $`r`$ and $`v`$ the orbital radius and velocity; the characteristic radius for a compact object being $`RGm/c^2.`$) The post-Newtonian order is their ratio $`PN𝒪(\delta f)=𝒪_ϵ(\delta f)/𝒪_ϵ(f)`$. When the time derivative $`\dot{f}`$ is given, we can evaluate the order of the change $`\delta f`$ generated during a characteristic time-interval of the quasi-periodic motion: $`PN𝒪(\delta f)=𝒪_ϵ(\dot{f})𝒪_ϵ(T)/𝒪_ϵ(f)`$. Here $`T`$ is the radial period and $`𝒪_ϵ(T)=rϵ^{1/2}/c`$. For a dimensionless quantity, typically an angle variable $`\alpha `$, the simple estimate $`PN𝒪(\delta \alpha )=𝒪_ϵ(\dot{\alpha })𝒪_ϵ(T)`$ holds. A numerical example shows that for a binary consisting of massive black holes with mass $`100M_{}`$ and orbital radius of $`10^4`$ km, i.e. , one-hundredth of the Hulse-Taylor system, the post-Newtonian parameter is $`ϵGm/rc^210^2`$. We also assume rapid rotation, thus the characteristic velocity of rotation $`V`$ is commesurable with the velocity of light $`𝒪(V/c)=1`$. If this assumption is lifted, the spin-spin effects appear at higher post-Newtonian orders than is assumed here. . Consequences of the spin precession equations (2) and of the conservation of the total angular momentum $`𝐉=𝐋+𝐒_\mathrm{𝟏}+𝐒_\mathrm{𝟐}`$ holding in the second post-Newtonian order are summarized in Table 1. In the present paper the angular momenta, their directions, magnitudes and the angles, all enlisted in Table 1 will be needed to a precision at which precessional effects do not contribute.
All spin terms in the radiation reaction up to the $`3/2`$ post-Newtonian order were evaluated in , starting from the expressions of the radiated power and total angular momentum loss derived by Kidder and from the Burke-Thorne potential . The employed averaging method, based on a suitably introduced radial parameter $`\chi `$ and on the residue theorem was described in detail in .
The second-order spin contribution to the radiation reaction was considered recently in . There we have developed the toolchest for computing secular radiative effects of spin-spin origin, essentially by proposing a description in terms of constants of the motion and the angular average $`\overline{L}`$ of the magnitude of the orbital angular momentum $`L(\chi )`$. The latter depends on the Keplerian true anomaly $`\chi `$ through a spin-spin term:
$`L(\chi )`$ $`=`$ $`\overline{L}{\displaystyle \frac{G\mu ^2}{2c^2\overline{L}^3}}S_1S_2\mathrm{sin}\kappa _1\mathrm{sin}\kappa _2\{2\overline{A}\mathrm{cos}[\chi +2(\psi _0\overline{\psi })]`$ (4)
$`+(3Gm\mu +2\overline{A}\mathrm{cos}\chi )\mathrm{cos}2(\chi +\psi _0\overline{\psi })\}.`$
The quantity $`\overline{A}=(G^2m^2\mu ^2+2E\overline{L}^2/\mu )^{1/2}`$ is the magnitude of the Laplace-Runge-Lenz vector for a Keplerian motion characterized by the energy $`E`$ and magnitude of orbital angular momentum $`\overline{L}`$. The angles $`\kappa _i=\mathrm{cos}^1(\widehat{𝐒}_𝐢\widehat{𝐋}),(i=1,2)`$ together with $`\gamma =\mathrm{cos}^1(\widehat{𝐒}_\mathrm{𝟏}\widehat{𝐒}_\mathrm{𝟐}).`$ characterize the relative orientation of the angular momentum vectors $`𝐒_\mathrm{𝟏},𝐒_\mathrm{𝟐}`$ and $`𝐋`$.
This description enabled us to introduce the generalized true anomaly parametrization $`r(\chi )`$ for the radial motion: (Eq. (3.7) of Ref.). By employing $`\chi `$ as radial variable, the average of a generic function $`f`$ over a radial period $`T`$
$$f=\frac{1}{T}_0^{2\pi }f(\chi )\frac{1}{\dot{r}}\frac{dr}{d\chi }𝑑\chi $$
(5)
is computed by Taylor expanding $`\dot{r}^2`$ taken from the radial equation, by substituting
$$\frac{dr}{d\chi }=\frac{1}{2}\left(\frac{1}{r_{min}}\frac{1}{r_{max}}\right)r^2\mathrm{sin}\chi $$
(6)
($`r_{_{_{min}^{max}}}`$ being the turning points) and finally, by switching to the complex parameter $`z=\mathrm{exp}(i\chi )`$. Then the residue theorem is employed in the computation of the integral (5), which in the majority of cases is given by the residues in the origin .
As application, we have computed the secular losses of the energy $`E`$ and of $`\overline{L}`$. They depend on $`E,`$ $`\overline{L},`$ on the angles$`\psi _i`$ subtended by the node line with the projections of the spins in the plane of the orbit and on the angles $`\kappa _i`$ and $`\gamma `$. The enlisted five angles are constrained by two algebraic relations :
$`0`$ $`=`$ $`S_1\mathrm{sin}\kappa _1\mathrm{cos}\psi _1+S_2\mathrm{sin}\kappa _2\mathrm{cos}\psi _2`$ (7)
$`\mathrm{cos}\gamma `$ $`=`$ $`\mathrm{cos}\kappa _1\mathrm{cos}\kappa _2+\mathrm{sin}\kappa _1\mathrm{sin}\kappa _2\mathrm{cos}\mathrm{\Delta }\psi ,`$ (8)
where $`\mathrm{\Delta }\psi =\psi _2\psi _1`$. We also introduce here the notation $`\overline{\psi }=(\psi _1+\psi _2)/2.`$ In order to have a closed system of differential equations, beside $`dE/dt`$ and $`dL/dt`$ computed in the radiative change of the angles $`\kappa _i`$ and $`\gamma `$ is required. This was given to $`3/2`$ post-Newtonian order accuracy in . It is the purpose of the present paper to compute the radiative evolution of the angles $`\kappa _i`$ and $`\gamma `$ up to the second post-Newtonian order by the inclusion of the spin-spin terms.
We complete the task by passing through the following steps. A well-known result states that no leading-order gravitational radiation leaves an axisymmetric body, as its quadrupole moment is a constant. It is also known that the radiative loss $`d𝐒_𝐢/dt`$ in the spins of the axisymmetric neutron stars or black holes forming a binary is two post-Newtonian orders higher that the leading radiative loss $`d𝐋/dt`$ of the orbital angular momentum . In Sec. 2 we derive a related generic result, which holds for binaries with axisymmetric constituents. We show that axisymmetry implies that the leading order secular contribution to the radiative change of the spins (which would appear at the second post-Newtonian order after the leading radiation terms) averages out during one cycle of the quasi-periodic motion. This result simplifies the study of the secular evolution of the angles. Indeed we will drop all terms of the form $`d𝐒_𝐢/dt`$ multiplied with any quantity which is a constant to the required accuracy.
The rest of the paper is organized as follows. In Sec. 3 we discuss the radiative change in the angles $`\kappa _i`$ and $`\gamma `$ by considering the Newtonian ($`N`$), first post-Newtonian ($`PN`$), second post-Newtonian ($`2PN`$), spin-orbit ($`SO`$), spin-spin ($`SS`$) and leading order tail contributions. We show that the $`N,PN,2PN`$ and tail terms drop out as a consequence of the property of the respective parts in $`d𝐉/dt`$ of being aligned with the Newtonian part of the orbital angular momentum. The $`SO`$ terms were given previously in . We complete the list by computing in Sec. 4 the averaged $`SS`$ contributions. They contain both two-body and self-interaction terms.
As a result we obtain the system of equations which govern the radiative secular evolution of the relative orientations of the angular momentum vectors $`𝐒_\mathrm{𝟏},𝐒_\mathrm{𝟐}`$ and $`𝐋`$ up to the second post-Newtonian order.
## II Radiative loss in the spins
In we have derived the radiative loss in the spin $`𝐒_𝐢`$ of the $`i^{th}`$ axisymmetric body, following . The integral over the volume of the body of the moment of the reaction force (the sign swapped gradient of the Burke-Thorne potential) gave:
$$\frac{d\left(𝐒_𝐢\right)_\mu }{dt}=\frac{2G}{5c^5\mathrm{\Omega }_i}\left(\frac{\mathrm{\Theta }_i}{\mathrm{\Theta }_i^{}}1\right)ϵ_{\mu \nu \rho }I_N^{(5)\nu \sigma }S_i(\widehat{𝐒}_𝐢)_\rho (\widehat{𝐒}_𝐢)_\sigma .$$
(9)
(There is no summation over $`i.`$) Each body is characterized by its principal moments of inertia $`\mathrm{\Theta }_i`$ and $`\mathrm{\Theta }_i^{}`$ and by the angular velocity $`\mathrm{\Omega }_i=S_i/\mathrm{\Theta }_i^{}`$, while $`I_N^{(5)\nu \sigma }`$ is the $`5^{th\text{ }}`$time derivative of the system’s Newtonian symmetric trace-free mass quadrupole moment. As the loss in the spin vector described by Eq. (9) is of second post-Newtonian order (above the leading radiation term in the orbital angular momentum loss), Eq. (9) is still valid to this order for approximately axisymmetric bodies, with the deviation from axisymmetry being of any post-Newtonian order. We emphasize that Eq. (9) implies $`d𝐒_𝐢/dt=`$ $`S_id\widehat{𝐒}_𝐢/dt`$, therefore to this order the radiation reaction will change the orientation but not the magnitude of the spin vectors. We also stress that $`d\widehat{𝐒}_𝐢/dt`$ generates $`3/2`$ post-Newtonian order radiative changes $`\delta _{rad}\widehat{𝐒}_𝐢`$ above the leading order radiation losses (in $`L`$ for example).
In coordinates $`(x,y,z)=r(\mathrm{cos}\psi ,\mathrm{sin}\psi ,0)`$ the spins are expressed as $`𝐒_𝐢=𝐒_i(\mathrm{sin}\kappa _i\mathrm{cos}\psi _i,\mathrm{sin}\kappa _i\mathrm{sin}\psi _i,\mathrm{cos}\kappa _i).`$ By employing the Keplerian equation of motion $`𝐚_N=Gm𝐫/r^3`$ and the radial equation $`\dot{r}^2=2E/\mu +2Gm/r\overline{L}^2/\mu ^2r^2`$ in Eq. (9), we obtain the instantaneous spin-loss equation:
$`{\displaystyle \frac{1}{S_i}}{\displaystyle \frac{d\left(𝐒_𝐢\right)_\mu }{dt}}`$ $`=`$ $`{\displaystyle \frac{2G^2m\mathrm{sin}\kappa _i}{5c^5\mu ^2r^7\mathrm{\Omega }_i}}\left({\displaystyle \frac{\mathrm{\Theta }_i}{\mathrm{\Theta }_i^{}}}1\right)\mathrm{\Lambda }_{i\mu }`$ (10)
$`\mathrm{\Lambda }_{i1}`$ $`=`$ $`\mathrm{cos}\kappa _i[a_1\mathrm{sin}(2\psi \psi _i)+a_2\mathrm{cos}(2\psi \psi _i)+a_3\mathrm{sin}\psi _i]`$ (11)
$`\mathrm{\Lambda }_{i2}`$ $`=`$ $`\mathrm{cos}\kappa _i[a_2\mathrm{sin}(2\psi \psi _i)a_1\mathrm{cos}(2\psi \psi _i)a_3\mathrm{cos}\psi _i]`$ (12)
$`\mathrm{\Lambda }_{i3}`$ $`=`$ $`\mathrm{sin}\kappa _i[a_1\mathrm{sin}2(\psi \psi _i)+a_2\mathrm{cos}2(\psi \psi _i)],`$ (13)
with the coefficients
$`a_1`$ $`=`$ $`\mu r\dot{r}(12\mu Er^2+20Gm\mu ^2r+45\overline{L}^2)`$ (14)
$`a_2`$ $`=`$ $`4\overline{L}(18\mu Er^2+20Gm\mu ^2r15\overline{L}^2)`$ (15)
$`a_3`$ $`=`$ $`\mu r\dot{r}(12\mu Er^2+20Gm\mu ^2r15\overline{L}^2).`$ (16)
In order to obtain the secular radiative changes of the spins we insert the Newtonian expressions
$$\dot{r}=\frac{\overline{A}}{\overline{L}}\mathrm{sin}\chi ,r=\frac{\overline{L}^2}{\mu (Gm\mu +\overline{A}\mathrm{cos}\chi )},\psi =\chi +\psi _0.$$
(17)
in Eq. (13).
Then we average by the method described in the Introduction (see also ) employing the Newtonian expression $`(1/\dot{r})(dr/d\chi )=dt/d\chi =\mu r^2/\overline{L}.`$ We obtain the generic result that the radiative change in the spins averages to zero in the second post-Newtonian approximation:
$$\frac{d𝐒_𝐢}{dt}=0.$$
(18)
Therefore the averaging-out property of some particular projections of the spin losses found in and also emerge.
As the first correction to both the acceleration and the Burke-Thorne potential is one post-Newtonian order higher than the respective leading terms , the averaging-out property of the radiative change in the spins derived above holds at $`5/2`$ post-Newtonian orders as well.
## III Radiative evolution of the angles $`\kappa _i`$ and $`\gamma `$
We obtain the equations for the radiative evolution of $`\kappa _i=\mathrm{cos}^1(\widehat{𝐒}_𝐢\widehat{𝐋}),(i=1,2)`$ and $`\gamma =\mathrm{cos}^1(\widehat{𝐒}_\mathrm{𝟏}\widehat{𝐒}_\mathrm{𝟐})`$ by differentiating the defining relations of these angles.
The simplest is the evolution of the angle $`\gamma `$ spanned by the two spin vectors:
$$\frac{d}{dt}\mathrm{cos}\gamma =\frac{d}{dt}\left(\widehat{𝐒}_\mathrm{𝟏}\widehat{𝐒}_\mathrm{𝟐}\right)=\widehat{𝐒}_\mathrm{𝟏}\frac{d\widehat{𝐒}_\mathrm{𝟐}}{dt}+\widehat{𝐒}_\mathrm{𝟐}\frac{d\widehat{𝐒}_\mathrm{𝟏}}{dt}.$$
(19)
The terms of the right hand side were already evaluated in the framework of the spin-orbit contributions in . Though they do not contain the spin explicitly, the angular velocities $`\mathrm{\Omega }_i`$ are present in their expression, Eq. (3.13) of . An order of magnitude estimate shows that the ratio of the usual spin-orbit terms (taken for example from the losses of $`\kappa _i`$, Eqs. (3.10)-(3.12) of ) and such terms is of the order unity for rapidly rotating objects. It was computed in , and one can check by simple inspection of Eq. (18) that $`\gamma `$ will receive no secular radiative change.
The angles spanned by the spins with the orbital angular momentum evolve in a more complicated fashion:
$$\frac{d}{dt}\mathrm{cos}\kappa _i=\frac{d}{dt}\left(\widehat{𝐒}_𝐢\frac{𝐋}{L(\chi )}\right)=\frac{1}{L(\chi )}\left[\widehat{𝐒}_𝐢\frac{d(𝐉𝐒_\mathrm{𝟏}𝐒_\mathrm{𝟐})}{dt}\frac{dL}{dt}\mathrm{cos}\kappa _i+𝐋\frac{d\widehat{𝐒}_𝐢}{dt}\right].$$
(20)
In the forthcoming expressions we replace $`L(\chi )\overline{L}`$ in all post-Newtonian terms. This is possible in the required second order accuracy as $`L(\chi )`$ differs from the constant $`\overline{L}`$ in spin-spin terms (see Eq. (4)).
By inserting the following expression for the loss of the magnitude of orbital angular momentum
$$\frac{dL}{dt}=\widehat{𝐋}\frac{d𝐋}{dt}=\widehat{𝐋}\frac{d𝐉}{dt}\widehat{𝐋}\frac{d𝐒_\mathrm{𝟏}}{dt}\widehat{𝐋}\frac{d𝐒_\mathrm{𝟐}}{dt}$$
(21)
into Eq. (20), and from $`\widehat{𝐒}_𝐢d𝐒_𝐢/dt=0`$ (no summation) we obtain:
$$\frac{d}{dt}\mathrm{cos}\kappa _i=\frac{1}{L(\chi )}(\widehat{𝐒}_𝐢\widehat{𝐋}\mathrm{cos}\kappa _i)\frac{d𝐉}{dt}+\frac{1}{\overline{L}}\left[\widehat{𝐒}_𝐢\frac{d𝐒_𝐣}{dt}+\widehat{𝐋}\left(\frac{d𝐒_\mathrm{𝟏}}{dt}+\frac{d𝐒_\mathrm{𝟐}}{dt}\right)\mathrm{cos}\kappa _i+𝐋\frac{d\widehat{𝐒}_𝐢}{dt}\right],$$
(22)
where $`ji.`$ Remarkably all terms in the square bracket average out due to Eq.(18). Therefore we will not consider them further.
For the evaluation of the first term of Eq. (22) we need the loss in the total angular momentum $`d𝐉/dt`$. It has the following structure To avoid confusion in notation, we denote the post-Newtonian, $`2^d`$ post-Newtonian, spin-orbit and spin-spin terms from a decomposition of a quantity expressed in terms of $`𝐫,𝐯,𝐒_𝐢`$ and $`\dot{r}`$ by lower-case indices $`pn,2pn,so`$ and $`ss`$, respectively. This notation applies to all terms of $`d𝐉/dt`$ computed in and . After inserting the expressions for $`v^2`$, $`\dot{r}^2`$ and $`\dot{r}`$, Eqs. (2.32), (2.33) of and (17), the terms of the new decomposition will be denoted by the respective capital letters. :
$$\frac{d𝐉}{dt}=\left(\frac{d𝐉}{dt}\right)_N+\left(\frac{d𝐉}{dt}\right)_{pn}+\left(\frac{d𝐉}{dt}\right)_{2pn}+\left(\frac{d𝐉}{dt}\right)_{so}+\left(\frac{d𝐉}{dt}\right)_{ss}+\left(\frac{d𝐉}{dt}\right)_{tail}.$$
(23)
The leading order contribution was computed long time ago by Peters , the $`pn`$ contribution by Junker and Schāfer , the $`so`$ contribution by Kidder , and the $`2pn`$ contribution by Gopakumar and Iyer in terms of $`𝐫,𝐯,𝐒_𝐢`$ and $`\dot{r}`$. The method of computing the $`ss`$ contribution was also indicated in . The leading tail contribution appears at $`3/2`$ post-Newtonian order and its time average was given by Rieth and Schäfer in the form of a Fourier series.
A property we want to stress is the following:
$$\left(\frac{d𝐉}{dt}\right)_{N,pn,2pn}=\mathrm{\Gamma }_{0,1,2}𝐋_N$$
(24)
where the coefficients $`\mathrm{\Gamma }_{0,1,2}`$ can be read from Eq. (3.28.a,b) of and from (3.9.d) of . The same property applies to the averaged tail contribution, given by Eq. (84) of . To see this we note that all coefficients (denoted there by $`{}_{n}{}^{}𝒮_0`$) of the spectral decomposition of the angular momentum loss (denoted there by $`𝒮`$) have only $`z`$components, thus:
$$\frac{d𝐉}{dt}_{tail}=\mathrm{\Gamma }_{tail}𝐋_N.$$
(25)
The next essential remark is that from
$$𝐋=𝐋_N+𝐋_{PN}+𝐋_{2PN}+𝐋_{SO}=(1+\gamma _1+\gamma _2)𝐋_N+𝐋_{SO}$$
(26)
we can express the orbital angular momentum as
$$𝐋_N=(1\gamma _1\gamma _2+\gamma _1^2)𝐋𝐋_{SO}.$$
(27)
Here the coefficients $`\gamma _{1,2}`$ are given by Eqs. (2.9.b,d) of and they are of first and second post-Newtonian order, respectively. Therefore up to an accuracy of second post-Newtonian order we can write
$$\left(\frac{d𝐉}{dt}\right)_{N+pn+2pn}=(\mathrm{\Gamma }_0+\eta _1+\eta _2)𝐋\mathrm{\Gamma }_0𝐋_{SO}$$
(28)
and
$$\frac{d𝐉}{dt}_{tail}=\mathrm{\Gamma }_{tail}𝐋,$$
(29)
where the coefficients of the expansion (28) are given by
$`\eta _1`$ $`=`$ $`\mathrm{\Gamma }_1\mathrm{\Gamma }_0\gamma _1`$ (30)
$`\eta _2`$ $`=`$ $`\mathrm{\Gamma }_2\mathrm{\Gamma }_1\gamma _1\mathrm{\Gamma }_0(\gamma _2\gamma _1^2).`$ (31)
We will not need their explicit expressions. By inserting Eq. (28) in Eq. (23), then the resulting expression in Eq. (22), due to the remark
$$(\widehat{𝐒}_𝐢\widehat{𝐋}\mathrm{cos}\kappa _i)𝐋=0$$
(32)
we obtain for the first term of Eq. (22):
$$\frac{1}{L(\chi )}(\widehat{𝐒}_𝐢\widehat{𝐋}\mathrm{cos}\kappa _i)\left(\frac{d𝐉}{dt}\right)=\frac{1}{\overline{L}}(\widehat{𝐒}_𝐢\widehat{𝐋}\mathrm{cos}\kappa _i)\left[\left(\frac{d𝐉}{dt}\right)_{tail}+\left(\frac{d𝐉}{dt}\right)_{so}\mathrm{\Gamma }_0𝐋_{SO}+\left(\frac{d𝐉}{dt}\right)_{ss}\right].$$
(33)
The tail term averages out due to Eq. (29) and property (32).
As suggested by the notation, the second and third terms of Eq. (33) are those spin-orbit contributions to the radiative loss of the angles $`\kappa _i`$, which do not originate in the Burke-Thorne potential. Indeed, modulo Burke-Thorne type terms they can be put into the concise form ($`\widehat{𝐒}_𝐢d\widehat{𝐋}/dt)_{SO}`$ which was computed previously (Eqs. (3.8)-(3.10) of ). The averaged expression for the spin-orbit type loss in $`\kappa _1`$ was also given by Eq. (4.4) of , and a similar expression can be found for the loss of $`\kappa _2`$ by interchanging the indices $`12`$ and the ratios $`\eta =m_2/m_1\eta ^1`$.
The last term of Eq. (33), modulo Burke-Thorne type terms, is the spin-spin part of the radiative evolution of the angles $`\kappa _i`$:
$$\left(\frac{d}{dt}\mathrm{cos}\kappa _i\right)_{SS}\frac{1}{\overline{L}}(\widehat{𝐒}_𝐢\widehat{𝐋}\mathrm{cos}\kappa _i)\left(\frac{d𝐉}{dt}\right)_{ss},$$
(34)
which will be computed in the next section. (We have denoted by $``$ the equality modulo Burke-Thorne type terms.)
## IV Secular radiative evolution equations
We start from the expression of $`(d𝐉/dt)_{ss}`$ given by Kidder in terms of the time derivatives of the mass quadrupole and velocity quadrupole moments (see also Eq. (4.18) of for the respective expression with $`c1G`$). The required $`SO`$ part of the velocity quadrupole moment was computed first by Kidder and later verified by several authors. (Rieth and Schāfer presented a derivation based on a different spin supplementary condition while Owen, Tagoshi and Ohasha have employed a $`\delta `$-function type energy-momentum tensor.)
After computing the expression $`(d𝐉/dt)_{ss}`$ in detail, we rewrite the last term of Eq. (33) as function of the radial variables $`r(\chi )`$ and $`\chi .`$ The procedure we follow was described in detail in . As a result both self-interaction and two-body spin-spin terms emerge:
$`\left[{\displaystyle \frac{1}{\overline{L}}}(\widehat{𝐒}_𝐢\widehat{𝐋}\mathrm{cos}\kappa _i)\left({\displaystyle \frac{d𝐉}{dt}}\right)_{ss}\right]_{SSself}`$ $`=`$ $`{\displaystyle \frac{2G^3m^2\mu \mathrm{sin}\kappa _i}{5c^7r^6}}\left[\left({\displaystyle \frac{S_i}{m_i}}\right)^2\mathrm{sin}\kappa _i\mathrm{cos}\kappa _i+\left({\displaystyle \frac{S_j}{m_j}}\right)^2\mathrm{sin}\kappa _j\mathrm{cos}\kappa _j\mathrm{cos}\mathrm{\Delta }\psi \right]`$ (35)
$`\left[{\displaystyle \frac{1}{\overline{L}}}(\widehat{𝐒}_𝐢\widehat{𝐋}\mathrm{cos}\kappa _i)\left({\displaystyle \frac{d𝐉}{dt}}\right)_{ss}\right]_{S_1S_2}`$ $`=`$ $`{\displaystyle \frac{2G^2S_1S_2\mathrm{sin}\kappa _i}{5c^7\mu ^2\overline{L}^2r^7}}\{u_1[\mathrm{sin}\kappa _i\mathrm{cos}\kappa _j+\mathrm{sin}\kappa _j\mathrm{cos}\kappa _i\mathrm{cos}\mathrm{\Delta }\psi ]`$ (39)
$`+u_2[\mathrm{sin}\kappa _i\mathrm{cos}\kappa _j\mathrm{cos}2(\chi +\psi _0\psi _i)+\mathrm{sin}\kappa _j\mathrm{cos}\kappa _i\mathrm{cos}2(\chi +\psi _0\overline{\psi })]`$
$`+\mathrm{sin}\chi \{u_3\mathrm{sin}\kappa _j\mathrm{cos}\kappa _i\mathrm{sin}(\psi _j\psi _i)`$
$`+u_4[\mathrm{sin}\kappa _i\mathrm{cos}\kappa _j\mathrm{sin}2(\chi +\psi _0\psi _i)+\mathrm{sin}\kappa _j\mathrm{cos}\kappa _i\mathrm{sin}2(\chi +\psi _0\overline{\psi })]\}\}.`$
Here $`ji`$ and the coefficients are
$`u_1`$ $`=`$ $`\overline{L}^2(12\mu Er^2+4Gm\mu ^2r15\overline{L}^2)`$ (40)
$`u_2`$ $`=`$ $`3\overline{L}^2(Gm\mu ^2r+3\overline{L}^2)`$ (41)
$`u_3`$ $`=`$ $`3\mu \overline{A}r(2\mu Er^2+5\overline{L}^2)`$ (42)
$`u_4`$ $`=`$ $`3\mu \overline{A}r(2\mu Er^2+3\overline{L}^2)`$ (43)
The averaging procedure based on the parametrization $`r(\chi )`$ and on the residue theorem yields the self-interaction and two-body spin-spin terms in the secular radiative loss of the angles $`\kappa _i`$ and $`\gamma `$. We enlist them together with the spin-orbit contributions:
$`{\displaystyle \frac{d\gamma }{dt}}`$ $`=`$ $`0`$ (44)
$`{\displaystyle \frac{d\kappa _i}{dt}}`$ $`=`$ $`{\displaystyle \frac{d\kappa _i}{dt}}_{SO}+{\displaystyle \frac{d\kappa _i}{dt}}_{SSself}+{\displaystyle \frac{d\kappa _i}{dt}}_{SS}`$ (45)
$`{\displaystyle \frac{d\kappa _i}{dt}}_{SO}`$ $`givenbyEq.(4.4)of\text{[12]}`$ (46)
$`{\displaystyle \frac{d\kappa _i}{dt}}_{SSself}`$ $`=`$ $`{\displaystyle \frac{G^2m\mu (2\mu E)^{3/2}}{20c^7\overline{L}^9}}V_1\left[\left({\displaystyle \frac{S_i}{m_i}}\right)^2\mathrm{sin}\kappa _i\mathrm{cos}\kappa _i+\left({\displaystyle \frac{S_j}{m_j}}\right)^2\mathrm{sin}\kappa _j\mathrm{cos}\kappa _j\mathrm{cos}\mathrm{\Delta }\psi \right]`$ (47)
$`{\displaystyle \frac{d\kappa _i}{dt}}_{S_1S_2}`$ $`=`$ $`{\displaystyle \frac{G^2(2\mu E)^{3/2}S_1S_2}{20c^7\overline{L}^9\mathrm{sin}\kappa _i}}\{V_2(\mathrm{sin}\kappa _i\mathrm{cos}\kappa _j+\mathrm{sin}\kappa _j\mathrm{cos}\kappa _i\mathrm{cos}\mathrm{\Delta }\psi )`$ (49)
$`+V_3[\mathrm{sin}\kappa _i\mathrm{cos}\kappa _j\mathrm{cos}2(\psi _0\psi _i)+\mathrm{sin}\kappa _j\mathrm{cos}\kappa _i\mathrm{cos}2(\psi _0\overline{\psi })]\},`$
where $`ji`$ and the coefficients $`V_{13}`$ are:
$`V_1`$ $`=`$ $`12E^2\overline{L}^4+60G^2m^2\mu ^3E\overline{L}^2+35G^4m^4\mu ^6`$ (50)
$`V_2`$ $`=`$ $`564E^2\overline{L}^4+1620G^2m^2\mu ^3E\overline{L}^2+805G^4m^4\mu ^6`$ (51)
$`V_3`$ $`=`$ $`60(8E^2\overline{L}^4+18G^2m^2\mu ^3E\overline{L}^2+7G^4m^4\mu ^6).`$ (52)
Note that $`V_1=D_1`$, because both arose from the average of $`r^6.`$ (The coefficient $`D_1`$, given by Eq. (4.32) of , governs the self-interaction term in the radiative loss of $`L`$). In order the $`3/2`$ post-Newtonian order-accurate spin-orbit contribution to hold at $`2PN`$, the replacements $`L\overline{L}`$ and $`A_0\overline{A}`$ should be carried on in Eq. (4.4) of .
We stress that Eqs. (44)-(49) contain all radiative terms in the angular evolution up to the second post-Newtonian order above the leading radiative effects.
## V Concluding remarks
By proving the $`5/2`$ post-Newtonian accurate result that axisymmetric objects do not radiate away any fraction of their initial spins, the computation of the secular angular evolutions induced by radiation reaction became simpler. Then from the analysis of $`d𝐉/dt`$ we could conclude that there are no $`N,PN,2PN`$ and tail contributions to the secular radiative angular evolutions. The $`SO`$ part of the equations was given in . We have derived the $`SS`$ terms here.
No spin-spin terms appear in the radiative evolution of the angle $`\gamma .`$ As the spin-orbit terms average out , we found the remarkable result that the angle spanned by the spins receives no radiative secular change up to the second post-Newtonian order. (However the angle $`\gamma ,`$ together with all other angles is subjected to precessional \- both instantaneous and secular - evolution. Although the second post-Newtonian order precessional (nonradiative) evolution of the angles is not developed in the paper, inspection of Table 1 shows that the precessions do not contribute to the 2PN radiative evolution of the angles $`\kappa _i`$ and $`\gamma `$, which contain only 3/2 PN and 2PN parts.)
The angles $`\kappa _i`$ evolve under the influence of both spin-orbit and spin-spin terms. The radiative angular evolution equations (44)-(49) together with the algebraic constraints (7)-(8) and the expressions for $`dE/dt`$ and $`dL/dt`$ derived in form a closed system of first order differential equations.
The angle $`\psi _0`$ appearing in this system is an integration constant. It is interpreted as the angle subtended by the node line with $`𝐫`$ at $`\chi =0`$ (with the periastron line). Each precession modifies the value of $`\psi _0`$ by a small amount. According to the arguments of Ryan in , where this angle first appeared, the terms containing periodic functions of $`\psi _0`$ average to zero whenever the precession time scale is short compared to the radiation reaction time scale.
As it happened with the loss of $`E`$ and $`L`$ , the spin-spin terms of the radiative $`\kappa _i`$-evolution could be decomposed into two-body and self-interaction terms. In the one-spin limit $`𝐒_\mathrm{𝟐}=0`$ the terms proportional to $`S_1^2`$ from Eq.(47) represent the second post-Newtonian correction to the radiative evolution of the angle $`\kappa _1`$ derived earlier in the Lense-Thirring approximation .
## VI Acknowledgments
This work has been supported by the Hungarian Scholarship Board. The algebraic package REDUCE was employed in some of the computations.
|
warning/0003/cond-mat0003265.html
|
ar5iv
|
text
|
# Bak–Sneppen model near zero dimension
## Abstract
We consider the Bak–Sneppen model near zero dimension where the avalanche exponent $`\tau `$ is close to $`1`$ and the exponents $`\mu `$ and $`\sigma `$ are close to $`0`$. We demonstrate that $`\tau 1=\mu \sigma =\mathrm{exp}\{\mu ^1\gamma +\mathrm{}\}`$ in this limit, where $`\gamma `$ is the Euler’s constant. The avalanche hierarchy equation is rewritten in a form that makes possible to find the relation between the critical exponents $`\sigma `$ and $`\mu `$ with high accuracy. We obtain new, more precise, values of the critical exponents for the 1D and 2D Bak–Sneppen model and for the 1D anisotropic Bak–Sneppen model.
PACS numbers: 05.40.+j, 64.60.Fr, 64.60.Lx, 87.10.+e
Perhaps the most simply formulated model showing avalanche behavior is the Bak–Sneppen model : ”What could be simpler than replacing some random numbers with some other random numbers?” . Nevertheless, the exact solution of the Bak–Sneppen model is unknown even in one dimension. The value of the most fundamental quantity, i.e. of the upper critical dimension is also still under discussion .
The formulation of the model is short indeed. A random number $`f_i`$ from some distribution $`𝒫(f)`$ is placed at an each site $`i`$ of a lattice. One replaces simultaneously the smallest number $`f_{min}`$ of them and the random numbers at its nearest neighbour sites by new random numbers from the distribution $`𝒫(f)`$, and afterwards the process is repeated.
Avalanches in the Bak–Sneppen model are defined in the following way. The $`f`$-avalanche is defined as the sequence of the steps at which $`f_{min}`$ remains smaller than the given parameter $`f`$. (One may find a more detailed definition in .)
A very significant step to the understanding of the nature of the avalanches in this model was made in by S. Maslov who introduced the so called avalanche hierarchy equation for the distribution $`P(s,f)`$ of $`f`$-avalanche sizes $`s`$ (i.e. of temporal durations). From this exact equation, one may get the additional relation between the critical exponents of the model. Unfortunately, the exact solution of the equation is known only for the mean field situation. Two first terms of the expansion from the mean field solution, i.e. from the upper critical dimension, were calculated in . The precision of the results obtained by direct numerical integration of the avalanche hierarchy equation in its original form is only comparable with the precision of the Monte Carlo simulations .
There is another way to get analytical results. It seems natural to start from the lower critical dimension which equals zero for the Bak–Sneppen model to make something similar to the well known $`2+ϵ`$ expansion. Here, we derive from the avalanche hierarchy equation some convenient relations which enable us find the singular relation between the critical exponents near zero dimension and get values of the exponents at integer dimensions. Traditionally, one relates the exponents $`\tau `$ and $`\mu `$ (see the definition of these exponents below). The total curve $`\tau (\mu )`$ with the particular points for the integer dimensions and the areas of applicability of the approaches of and ours is depicted in Fig. 1.
For the distribution $`𝒫(f)=e^f,f>0`$ , the avalanche hierarchy equation is of the form
$$\frac{P(s,f)}{f}=\underset{t=1}{\overset{s1}{}}t^\mu P(t,f)P(st,f)s^\mu P(s,f).$$
(1)
Here, in the scaling region, $`s^\mu `$ gives the average number of the distinct sites updating during an avalanche of the size $`s`$, where $`\mu =d/D_f`$, $`d`$ is the dimension of the lattice, $`D_f`$ is an avalanche fractal dimension . The physical meaning of the equation describing the hierarchical nature of avalanches in the Bak–Sneppen model is the following. The distribution $`P(s,f)`$ changes while $`f`$ grows because of two reasons. First, two consecutive avalanches of size $`t`$ and $`st`$ contribute to the avalanche of size $`s`$ (the second avalanche starts from one of the sites changed during the first avalanche that gives the factor $`t^\mu `$ in the sum). Second, some avalanches of size $`s`$ merge into a larger avalanche.
For the problem under consideration, the exponent $`\mu `$, which equals $`0`$ at the lower critical dimension and equals $`1`$ at the upper critical dimension (simulation in and , gave different values for it, $`d_{uc}=8`$ and $`4`$, correspondingly), is the given parameter. All other exponents are related with $`\mu `$ by Eq. (1). Below the threshold $`f_c`$ (i.e. in the symmetric phase) the solution of Eq. (1) has the following scaling form:
$$P(s,f)=s^\tau F\left(s^\sigma (f_cf)\right).$$
(2)
Eq. (1) resembles nonlinear differential equations with a peaking regime . For such equations, it is possible to find both exponents included in Eq. (2).
In , it was proposed to search for the Laplace transform of the distribution $`P(s,f)`$:
$$p(\alpha ,f)=\underset{s=1}{\overset{\mathrm{}}{}}P(s,f)e^{\alpha s}.$$
(3)
Then Eq. (1) gives
$`{\displaystyle \frac{1}{1p(\alpha ,f)}}{\displaystyle \frac{p(\alpha ,f)}{f}}={\displaystyle \underset{s=1}{\overset{\mathrm{}}{}}}P(s,f)s^\mu e^{\alpha s}=`$ (4)
$`(1)^\mu {\displaystyle \frac{^\mu p(\alpha ,f)}{\alpha ^\mu }}={\displaystyle \frac{1}{\mathrm{\Gamma }(1\mu )}}{\displaystyle _0^{\mathrm{}}}𝑑tt^\mu {\displaystyle \frac{p(\alpha +t,f)}{\alpha }},`$ (5)
where $`^\mu /\alpha ^\mu `$ denotes the fractional partial derivative ($`\mu `$ is certainly noninteger) and the last expression is its integral representation. $`\mathrm{\Gamma }()`$ is the gamma-function. The scaling relation for the solution of Eq. (4) is
$$p(\alpha ,f)=1\alpha ^{\tau 1}h\left(\frac{f_cf}{\alpha ^\sigma }\right).$$
(6)
Inserting Eq. (6) into Eq. (4) one obtains the usual relation between the critical exponents
$$\tau =1+\mu \sigma $$
(7)
and the following integral-differential equation for the scaling function $`h(x)`$:
$$\mathrm{\Gamma }(1\mu )x\frac{h^{}(x)}{h(x)}=_0^x𝑑y\frac{yh^{}(y)\frac{\mu \sigma }{\sigma }h(y)}{[1(y/x)^{1/\sigma }]^\mu }.$$
(8)
Here $`h^{}(x)dh(x)/dx`$.
In , Eq. (8) was used to obtain the expansion from the mean field solution but it seems to be inconvenient. Let us show that one may transfer it to a purely integral form. The following lines demonstrate how that integration may be done.
$`\mathrm{\Gamma }(1\mu ){\displaystyle \frac{d\mathrm{ln}h(x)}{dx}}`$ (9)
$`={\displaystyle _0^1}{\displaystyle \frac{dz}{(1z^{1/\sigma })^\mu }}\left[xz{\displaystyle \frac{dh(xz)}{d(xz)}}{\displaystyle \frac{\mu \sigma }{\sigma }}h(xz)\right]={\displaystyle _0^1}{\displaystyle \frac{dz}{(1z^{1/\sigma })^\mu }}(xz)^{(\mu \sigma )/\sigma +1}{\displaystyle \frac{d}{d(xz)}}\left[(xz)^{(\mu \sigma )/\sigma }h(xz)\right]`$ (10)
$`=x^{(\mu \sigma )/\sigma +1}{\displaystyle _0^1}{\displaystyle \frac{dz}{(1z^{1/\sigma })^\mu }}z^{(\mu \sigma )/\sigma +1}{\displaystyle \frac{1}{z}}{\displaystyle \frac{d}{dx}}\left[x^{(\mu \sigma )/\sigma }z^{(\mu \sigma )/\sigma }h(xz)\right]`$ (11)
$`=x^{(\mu \sigma )/\sigma +1}{\displaystyle \frac{d}{dx}}x^{(\mu \sigma )/\sigma }{\displaystyle _0^1}𝑑z{\displaystyle \frac{h(xz)}{(1z^{1/\sigma })^\mu }}.`$ (12)
Applying $`_0^x𝑑x`$ to the first and the last lines of Eq. (9) and then integrating by parts (one may chose $`h(x=0)=1`$ ) we get
$$\mathrm{\Gamma }(1\mu )\mathrm{ln}h(x)=_0^1\frac{dz}{(1z^{1/\sigma })^\mu }\left[xh(xz)\left(\frac{\mu \sigma }{\sigma }+1\right)_0^x𝑑uh(uz)\right]$$
(13)
and finally we obtain the equation for the scaling function $`h(x)`$ in the most convenient form:
$$h(x)=\mathrm{exp}\left\{\frac{1}{\mathrm{\Gamma }(1\mu )}_0^x\frac{dy}{[1(y/x)^{1/\sigma }]^\mu }\left[h(y)\frac{\mu }{\sigma }\frac{1}{y}_0^y𝑑zh(z)\right]\right\}$$
(14)
(if one does not demand $`h(0)=1`$, $`h(x)`$ in the left parts of Eqs. (13) and (14) is $`h(x)/h(0)`$).
Asymptotic form of $`h(x)`$ for large $`x`$ follows from the expansion of Eq. (3) in small $`\alpha `$. Below the threshold, $`h(x)`$ has to be
$$h(x)x^{(\mu \sigma )/\sigma 1/\sigma }(c_0+c_1x^{1/\sigma }+c_2x^{2/\sigma }+\mathrm{}).$$
(15)
This particular asymptotic behavior fixes the solution of Eq. (14) and the value of $`\sigma `$ for any given $`\mu `$. Substituting Eq. (15) into Eq. (8), Eq. (13), or Eq. (14), one gets the sume rule
$$_0^{\mathrm{}}𝑑xh(x)=\frac{1(\mu \sigma )}{\mu }\mathrm{\Gamma }(1\mu ).$$
(16)
Note that if $`h(x,\mu ,\sigma )`$ is a solution of Eq. (14) then $`ch(cx,\mu ,\sigma )`$ is also a solution for any constant $`c`$.
Eqs. (14) and (16) are the set of equations that lead to the scaling function $`h(x,\mu )`$ and $`\sigma (\mu )`$. Instead of Eq. (16), one may use equally the condition on the value of the exponent of the asymptote
$$[xh^{}(x)/h(x)](x\mathrm{})=\frac{(\mu \sigma )1}{\sigma }.$$
(17)
Hence, the problem is reduced to the eigen value problem for the nonlinear equation .
Let us study the solution of the system for small $`\mu `$. The expansion of the solution of Eq. (14) in $`x`$ looks as
$`\mathrm{ln}h(x)=\left(1{\displaystyle \frac{\mu }{\sigma }}\right)B(\sigma ,1\mu )\left({\displaystyle \frac{\sigma x}{\mathrm{\Gamma }(1\mu )}}\right)+\left(1{\displaystyle \frac{\mu }{\sigma }}\right)B(\sigma ,1\mu )\left(1{\displaystyle \frac{1}{2}}{\displaystyle \frac{\mu }{\sigma }}\right)B(2\sigma ,1\mu )\left({\displaystyle \frac{\sigma x}{\mathrm{\Gamma }(1\mu )}}\right)^2+`$ (18)
$`\left(1{\displaystyle \frac{\mu }{\sigma }}\right)B(\sigma ,1\mu )\left[{\displaystyle \frac{1}{2}}\left(1{\displaystyle \frac{\mu }{\sigma }}\right)B(\sigma ,1\mu )+\left(1{\displaystyle \frac{1}{2}}{\displaystyle \frac{\mu }{\sigma }}\right)B(2\sigma ,1\mu )\right]\left(1{\displaystyle \frac{1}{3}}{\displaystyle \frac{\mu }{\sigma }}\right)B(3\sigma ,1\mu )\left({\displaystyle \frac{\sigma x}{\mathrm{\Gamma }(1\mu )}}\right)^3+\mathrm{}`$ (19)
$`={\displaystyle \frac{\mu \sigma }{\sigma }}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n}}B(\sigma ,1\mu )\mathrm{}B(n\sigma ,1\mu )\left({\displaystyle \frac{\sigma x}{\mathrm{\Gamma }(1\mu )}}\right)^n+\mathrm{},`$ (20)
where $`B(,)`$ is the beta-function. We shall see that the quantity $`(\mu \sigma )/\sigma `$ is the smallest parameter of the problem near the lower critical dimension. If one tends formally $`\mu `$ to $`0`$, the last line of Eq. (18) tends to
$$\mathrm{ln}h(x)=\frac{\mu \sigma }{\sigma }\underset{n=1}{\overset{\mathrm{}}{}}\frac{\mathrm{\Gamma }(1+n\mu )}{nn!}\left(\frac{x}{\mathrm{\Gamma }(1\mu )}\right)^n$$
(21)
and afterwards to
$$\mathrm{ln}h(x)=\frac{\mu \sigma }{\sigma }\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{nn!}\left(\frac{x}{\mathrm{\Gamma }(1\mu )}\right)^n.$$
(22)
Thus, for small $`\mu `$, the solution $`h(x)`$ behaves as the following. For low enough $`x`$, the solution very slowly decreases from the value $`h(0)=1`$ and in some crossover region, $`x1/\mu `$, it comes to the asymptotic power tail, Eq. (15). We have to stress that three last limit equalities may be justified only for small $`x`$ and that the omitted terms in Eq. (18) do contribute to the power-law tail. Nevertheless, one may try to estimate $`\sigma (\mu )`$ for small $`\mu `$ inserting Eqs. (21) and (22) into the sum rule (16). The solution of the first of the equations obtained in such a way is $`\mu \sigma =\mathrm{exp}\{\mu ^12\gamma +O(\mu )\}`$, where $`\gamma =0.5772\mathrm{}`$ is the Euler’s constant, and the solution of the second one is of the form
$$\mu \sigma =\tau 1=\mathrm{exp}\{\mu ^1\gamma +O(\mu )\}.$$
(23)
Thus, the dependence is non analytical but even the second term of the expansion can not be defined by such estimation.
In fact we failed to obtain the value of the constant analytically. Nevertheless, Eqs. (14) and (16) are very convenient for numerics, since iterations of Eq. (14) converge. (One may start, for instance, from the functions (21) or (22).) We checked the validity of the relation (23) for small $`\mu `$. The value of the constant in Eq. (23) obtained in such a way is $`0.5771(5)`$, i. e. that is the Euler’s constant indeed.
Solving Eq. (14) with the constraint, Eq. (16) or Eq. (17), and the initial condition, one may easily get $`\sigma (\mu )`$ and $`\tau (\mu )`$ for any given $`\mu `$. The values of the exponent $`\mu `$ are known from simulation at integer dimensions with much higher precision than the values of $`\tau `$ because of the better available statistics . Therefore we can improve essentially the precision of the known value of $`\tau `$. For the 1D Bak–Sneppen model, we get $`\tau =1.0637(5)+0.4(\mu 0.4114)`$, where $`\mu =0.4114(2)`$ is the value obtained from the Monte-Carlo simulation . For the 2D Bak–Sneppen model, we obtain $`\tau =1.229(1)+0.77(\mu 0.685)`$, where $`\mu =0.685(5)`$ is the value obtained in . (The last relations may be used to obtain better values of $`\tau `$ when the more precise values of the exponent $`\mu `$ will be available.) Now the precision of $`\tau `$ coincides with that one of $`\mu `$. Note that these values are below the values of $`\tau `$ previously obtained from the simulation, $`\tau (1D)=1.073(3)`$ and $`\tau (2D)=1.245(10)`$ but are in accordance with the less precise values found in by direct numerical solution of the avalance hierarchy equation, Eq. (1). The value of the exponent $`\tau `$ of the 3D Bak–Sneppen model may be obtained from the expansion from the upper critical dimension . In Fig. 1, we show the curve $`\tau (\mu )`$ together with the points for the integer dimensions and the low-$`\mu `$ asymptote, (23), and the expansion from the upper critical dimension .
Of course, the relation (23) is valid only for $`\mu 1`$. Nevertheless, let us compare the value of $`\tau `$ at $`\mu =0.4114`$ obtained from Eq. (23), $`\tau =1.0494`$, with the calculated above $`\tau (1D)`$. One may see that these values are in qualitative agreement.
In the case of the 1D anisotropic Bak–Sneppen model (i. e. for the update of the extremal site and only one its, for instance, right neighbour) the exponents $`\sigma `$ and $`\mu `$ are coupled by the following additional relation: $`\sigma +\mu =1`$ . Hence, one can find all the exponents of the problem. We obtained from Eq. (14) the following value, $`\mu =0.5779(5)`$. In , two different values for $`\mu `$ were obtained, $`\mu =0.58`$ from Eq. (1) and $`\mu =0.588`$ found in another way. The Monte-Carlo simulation made in and gave $`\mu =0.60(1)`$ and $`\mu =0.59(3)`$ correspondingly. Therefore, we had to check our result. For that we solved numerically Eq. (8) with the initial condition $`h(0)=1`$ and the constraint, Eq. (16) or Eq. (17). The result is $`\mu =0.5778(5)`$. Thus, the value $`\mu =0.578`$ seems to be more reliable but the question is still open.
In summary, we have demonstrated that the simple transformation of the avalanche hierarchy equation made it convenient for analysis and numerics. We have obtained the non trivial singular relation, $`\tau 1=\mu \sigma =\mathrm{exp}\{\mu ^1\gamma +\mathrm{}\}`$ with the Euler’s constant $`\gamma `$, between the scaling exponents of the Bak–Sneppen model near zero dimension. Using the known from simulation values of the exponent $`\mu `$, we have found, in fact, all other exponents of the Bak–Sneppen model in 1D and 2D with the same high precision. We have got also the exponents of the anisotropic 1D Bak–Sneppen model.
Nevertheless, one should note that the main problem of obtaining of the last independent critical exponent of the Bak–Sneppen model in a regular way remains open.
SND thanks PRAXIS XXI (Portugal) for a research grant PRAXIS XXI/BCC/16418/98. JFFM was partially supported by the project PRAXIS/2/2.1/FIS/299/94 and YGP was partially supported by the project PRAXIS/2/2.1/FIS/302/94. We also thank M.C. Marques for reading the manuscript and A.V. Goltsev and A.N. Samukhin for many useful discussions.
Electronic address: sdorogov@fc.up.pt
Electronic address: jfmendes@fc.up.pt
Electronic address: ypogorel@fc.up.pt
|
warning/0003/math0003041.html
|
ar5iv
|
text
|
# 1.Introduction
## 1.Introduction
It is well known that Virasoro algebra and affine Lie algebras play a central role in conformal field theories (CFT) in two dimensional spacetime, which in quantum theory corresponds to fields of massless particles, or massive quantum field theories (QFT) at the critical points, i.e. when the correlation length becomes infinity (scale invariant).
It is an interesting problem to describe the symmetry of the QFT off the critical points. Indeed, impressive progresses have been made in this direction, part of which is represented by the quantum affine algebra , which describes the symmetries of the certain lattice models and, in massive QFT, acts as the role of affine Lie algebra in CFT. For example, the spin $`1/2`$ XXZ chain in the region $`\mathrm{\Delta }<1`$ (off-critical) was exactly solved with the help of the bosonic representation of $`U_q(\widehat{sl_2})`$ at level $`1`$ ; in the same way, the higher spin $`XXZ`$ chain cannot be exactly solved without involving the higher level representation of the quantum affine algebra $`U_q(\widehat{sl_2})`$ . It is known that a quantum affine Lie algebra corresponds to certain trigonometric solution of the quantum Yang-Baxter equation (QYBE); while for a rational solution of QYBE, the corresponding algebraic structure is the the quantum double of Yangian with central extension . Contrary to the quantum affine case, the Yangian double with center is obtained relatively rather late due to some technical difficulty. Before its explicit formulation was given in , the Yangian double with center was expected to be the symmetry algebra of quantum non-local currents of massive field theories .
Integrable massive QFTs, which are QFTs away from the critical phase, can be obtained either by perturbations of CFT with relevant fields, or in terms of free field realizations, following Lykyanov . In fact the free field realization is a common approach used in both massive quantum field theories and the representation theory of their dynamical symmetry algebras. The free boson representations of $`U_q(\widehat{sl_2})`$ with an arbitrary level have been obtained in Refs.. For the Yangian double with center, the free field representation of $`DY_{\mathrm{}}(sl_2)`$ with level $`k(0,2)`$ was constructed in , the level-$`1`$ and level-$`k`$ representation of $`DY_{\mathrm{}}(sl_N)`$ are obtained in and , respectively. The level-$`k`$ free field representation of $`DY_{\mathrm{}}(gl_N)`$ is also given in . However, all of these are deformations of the Wakimoto modules. Another kind of module has been known in the representation theories of the classical and quantum affine algebras. In the case of Yangian double with center a free field realization, which corresponds to the Feigin-Fuchs representation of usual affine Lie algebras in the classical limit, is given in . However, different from the corresponding quantum affine case, the coset structure in that realization is absent. So, explicit expression of rational quantum currents in the coset space $`SU(2)/U(1)`$ is not known yet. In the classical case, the currents defined on the coset space $`SU(2)/U(1)`$ subject to a $`SU(2)`$ nonlocal currents algebra, which is also referred to as parafermion algebra . Naively speaking, we would expect that the rational quantum (or $`\mathrm{}`$)-deformed parafermion are relevant to the off-critical coset WZNW model, or equivalently, the massive WZNW model in homogeneous space $`SU(2)/U(1)`$. The $`\mathrm{}`$-deformed parafermion is also relevant to the $`SU(2)`$ Yangian double with center $`DY_{\mathrm{}}(sl_2)_k`$, which plays an important role in condensed state physics . In fact the relation with the $`DY_{\mathrm{}}(sl_2)_k`$ is our main criterion to define the rational coset quantum currents. One of the main properties of the massive current is nonlocality. It is interesting to study the rational quantum deformation of coset currents, which is nonlocal even in the classical case. On the other hand, to study the rational quantum deformation of the nonlocal algebra is helpful for the investigating of the quantum deformation of the chiral vertex operator in more general scheme.
The manuscript is arranged as follows. First, we briefly review the definition of the $`SU(2)`$ nonlocal currents in Section 2. Then we propose a nonlocal quantum currents $`\mathrm{\Psi }(u)`$ and $`\mathrm{\Psi }^{}(u)`$, and the central extension of Yangian double for the Drinfeld new realization are given in terms of these rational quantum nonlocal currents and a $`U(1)`$ current. In section $`3`$, we give a bosonic representation for a nonlocal currents through two sets of bosonic fields, the operator product expansion (OPE) of this nonlocal currents satisfies the definition of the proposed quantum nonlocal currents after a Wick rotation.
## 2.Quantum nonlocal currents
In this section we propose a quantum nonlocal currents (QNC) $`\mathrm{\Psi }(u)`$ and $`\mathrm{\Psi }^{}(u)`$, in the classical limit $`\mathrm{}0`$, this nonlocal quantum currents become the nonlocal currents $`\psi (z)`$ and $`\psi ^{}(z)`$ of the coset space $`SU(2)/U(1)`$ respectively. Then in terms of such quantum nonlocal currents and a quantum $`U(1)`$ bosonic current, we obtain a new kind of representation for the Yangian double with center in the Drinfeld new realization. Before we going to the details, we first review briefly the theory of parafermionic currents in CFT.
Parafermionic currents are primary fields of the 2d CFT. The general parafermion defined for root lattices are proposed in . For the case of $`SU(2)`$, the affine Kac-Moody currents and parafermionic currents are related by the relations,
$`\chi _+(z)`$ $`=`$ $`\sqrt{k}:\psi (z)\mathrm{exp}(i\varphi (z)/\sqrt{k}):,`$
$`\chi _{}(z)`$ $`=`$ $`\sqrt{k}:\psi ^{}(z)\mathrm{exp}(i\varphi (z)/\sqrt{k}):,`$
$`h(z)`$ $`=`$ $`i\sqrt{k}_z\varphi (z),`$ (2.1)
where $`\chi _\pm (z)`$ and $`h(z)`$ are currents of $`SU(2)`$ affine algebra, and the radial ordering of the parafermions are given by
$$R\left(\psi _\alpha (z)\psi _\beta (w)\right)(zw)^{2\alpha \beta /k}=R\left(\psi _\beta (w)\psi _\alpha (z)\right)(wz)^{2\alpha \beta /k},$$
(2.2)
in which $`\alpha ,\beta =\pm `$. We will drop the $`R`$ symbol in the following without of any confusion. The OPE of the parafermionic fields defined by
$`\psi _\pm (z)\psi _\pm (w)(zw)^{2/k}`$ $`=`$ $`reg.`$
$`\psi _+(z)\psi _{}(w)(zw)^{2/k}`$ $`=`$ $`{\displaystyle \frac{1}{(zw)^2}}+reg,`$ (2.3)
For rational quantum deformation the $`SU(2)`$ currents algebra becomes the central extension of Yangian double for the Drinfeld new realization. From this point view we propose rational quantum deformation of the above nonlocal currents. The results are:
###### Proposition 1
The rational quantum deformation of nonlocal currents for $`SU(2)/U(1)`$ can be defined as:
$`((uv)+\mathrm{}){\displaystyle \frac{\mathrm{\Gamma }(\frac{(uv)}{k\mathrm{}}+\frac{1}{k}1)}{\mathrm{\Gamma }(\frac{(uv)}{k\mathrm{}}\frac{1}{k}1)}}\mathrm{\Psi }(u)\mathrm{\Psi }(v)`$ $`=`$ $`((uv)\mathrm{}){\displaystyle \frac{\mathrm{\Gamma }(\frac{(vu)}{k\mathrm{}}\frac{1}{k}1)}{\mathrm{\Gamma }(\frac{(vu)}{k\mathrm{}}\frac{1}{k}1)}}\mathrm{\Psi }(v)\mathrm{\Psi }(u),`$
$`((uv)\mathrm{}){\displaystyle \frac{\mathrm{\Gamma }(\frac{(uv)}{k\mathrm{}}+\frac{1}{k})}{\mathrm{\Gamma }(\frac{(uv)}{k\mathrm{}}\frac{1}{k})}}\mathrm{\Psi }^{}(u)\mathrm{\Psi }^{}(v)`$ $`=`$ $`((uv)+\mathrm{}){\displaystyle \frac{\mathrm{\Gamma }(\frac{(vu)}{k\mathrm{}}+\frac{1}{k})}{\mathrm{\Gamma }(\frac{(vu)}{k\mathrm{}}\frac{1}{k})}}\mathrm{\Psi }^{}(v)\mathrm{\Psi }^{}(u),`$
$`{\displaystyle \frac{\mathrm{\Gamma }(\frac{(uv)}{k\mathrm{}}\frac{1}{k}\frac{1}{2})}{\mathrm{\Gamma }(\frac{(uv)}{k\mathrm{}}+\frac{1}{k}\frac{1}{2})}}\mathrm{\Psi }(u)\mathrm{\Psi }^{}(v)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }(\frac{(vu)}{k\mathrm{}}\frac{1}{k}\frac{1}{2})}{\mathrm{\Gamma }(\frac{(vu)}{k\mathrm{}}+\frac{1}{k}\frac{1}{2})}}\mathrm{\Psi }^{}(v)\mathrm{\Psi }(u),`$ (2.4)
At first glance the above relations may seem very strange. Their actual meaning in the theory of quantum deformed algebras will soon be clear if we resort to the construction of Yangian double analogous to that of SU(2) current algebra in terms of ordinary parafermionic currents, and indeed, such a construction exists, and its classical limit $`\mathrm{}0`$ gives precisely the parafermionic construction of SU(2) currents.
To make the above statements more concrete, we need to introduce another bosonic currents.
$$[\widehat{c}(t),\widehat{c}(t^{})]=\frac{\mathrm{sinh}\mathrm{}t\mathrm{sinh}\frac{k}{2}\mathrm{}t}{\mathrm{}^2t}\delta (t+t^{}).$$
(2.5)
Defining the bosonic currents
$`C^+(u)=\mathrm{exp}\{\mathrm{}{\displaystyle _{\mathrm{}}^0}𝑑te^{\frac{k}{4}\mathrm{}t}{\displaystyle \frac{e^{iut}}{\mathrm{sinh}\frac{k}{2}\mathrm{}t}}\widehat{c}(t)\mathrm{}{\displaystyle _0^{\mathrm{}}}𝑑te^{\frac{k}{4}\mathrm{}t}{\displaystyle \frac{e^{iut}}{\mathrm{sinh}\frac{k}{2}\mathrm{}t}}\widehat{c}(t)\},`$
$`C^{}(u)=\mathrm{exp}\{\mathrm{}{\displaystyle _{\mathrm{}}^0}𝑑te^{\frac{k}{4}\mathrm{}t}{\displaystyle \frac{e^{iut}}{\mathrm{sinh}\frac{k}{2}\mathrm{}t}}\widehat{c}(t)+\mathrm{}{\displaystyle _0^{\mathrm{}}}𝑑te^{+\frac{k}{4}\mathrm{}t}{\displaystyle \frac{e^{iut}}{\mathrm{sinh}\frac{k}{2}\mathrm{}t}}\widehat{c}(t)\},`$ (2.6)
we have the following exchange relations,
$`{\displaystyle \frac{\mathrm{\Gamma }(\frac{i(uv)}{k\mathrm{}}\frac{1}{k}1)}{\mathrm{\Gamma }(\frac{i(uv)}{k\mathrm{}}+\frac{1}{k}1)}}C^+(u)C^+(v)`$ $`=`$ $`C^+(v)C^+(u){\displaystyle \frac{\mathrm{\Gamma }(\frac{i(vu)}{k\mathrm{}}\frac{1}{k}1)}{\mathrm{\Gamma }(\frac{i(vu)}{k\mathrm{}}+\frac{1}{k}1)}},`$
$`{\displaystyle \frac{\mathrm{\Gamma }(\frac{i(uv)}{k\mathrm{}}\frac{1}{k})}{\mathrm{\Gamma }(\frac{i(uv)}{k\mathrm{}}+\frac{1}{k})}}C^{}(u)C^{}(v)`$ $`=`$ $`C^{}(v)C^{}(u){\displaystyle \frac{\mathrm{\Gamma }(\frac{i(vu)}{k\mathrm{}}\frac{1}{k})}{\mathrm{\Gamma }(\frac{i(vu)}{k\mathrm{}}+\frac{1}{k})}},`$
$`{\displaystyle \frac{\mathrm{\Gamma }(\frac{i(uv)}{k\mathrm{}}+\frac{1}{k}+\frac{1}{2})}{\mathrm{\Gamma }(\frac{i(uv)}{k\mathrm{}}\frac{1}{k}+\frac{1}{2})}}C^+(u)C^{}(v)`$ $`=`$ $`C^{}(v)C^+(u){\displaystyle \frac{\mathrm{\Gamma }(\frac{i(vu)}{k\mathrm{}}+\frac{1}{k}+\frac{1}{2})}{\mathrm{\Gamma }(\frac{i(vu)}{k\mathrm{}}\frac{1}{k}+\frac{1}{2})}}.`$ (2.7)
If we perform the Wick rotation for currents $`C^\pm `$, namely replacing $`\mathrm{}`$ by $`i\mathrm{}`$, and defining the quantum currents as
$`E(u)=:\mathrm{\Psi }(u)C^+(u):,`$
$`F(u)=:\mathrm{\Psi }^{}(u)C^{}(u):,`$ (2.8)
we obtain the Drinfeld currents of Yangian double with center, with relations
$`[H^\pm (u),H^\pm (v)]=0,[c,\text{everything}]=0,c=k,`$
$`(uv+\mathrm{}{\displaystyle \frac{k}{2}}\mathrm{})(uv\pm {\displaystyle \frac{k}{2}}\mathrm{}\mathrm{})H^\pm (u)H^{}(v)`$
$`=(uv\mathrm{}{\displaystyle \frac{k}{2}}\mathrm{}\mathrm{})(uv\pm {\displaystyle \frac{k}{2}}\mathrm{}+\mathrm{})H^{}(v)H^\pm (u),`$
$`(uv+\mathrm{}{\displaystyle \frac{k}{4}}\mathrm{})H^\pm (u)E(v)=(uv\mathrm{}{\displaystyle \frac{k}{4}}\mathrm{})E(v)H^\pm (u),`$
$`(uv\mathrm{}\pm {\displaystyle \frac{k}{4}}\mathrm{})H^\pm (u)F(v)=(uv+\mathrm{}\pm {\displaystyle \frac{k}{4}}\mathrm{})F(v)H^\pm (u),`$
$`(uv+\mathrm{})E(u)E(v)=(uv\mathrm{})E(v)E(u),`$
$`(uv\mathrm{})F(u)F(v)=(uv+\mathrm{})F(v)F(u),`$
$`[E(u),F(v)]={\displaystyle \frac{1}{\mathrm{}}}\left(\delta (uv+{\displaystyle \frac{k}{2}}\mathrm{})H^+(u+{\displaystyle \frac{k}{2}}\mathrm{})\delta (uv{\displaystyle \frac{k}{2}}\mathrm{})H^{}(v+{\displaystyle \frac{k}{2}}\mathrm{})\right),`$ (2.9)
wherein
$`H^+(u)=\text{exp}\{2\mathrm{}{\displaystyle _0^+\mathrm{}}𝑑te^{iut}\widehat{c}(t)\},`$
$`H^{}(u)=\text{exp}\{2\mathrm{}{\displaystyle _{\mathrm{}}^0}𝑑te^{iut}\widehat{c}(t)\}.`$ (2.10)
Since the fields $`C^\pm (u)`$ and $`H^\pm (u)`$ are only involved in the quantum $`U(1)`$ current, the coset structure of classical currents are preserved, and the coset structure could not be preserved by another kind of free field representation .
## 3.Bosonization of the quantum nonlocal currents
In order to see that the strange rational deformed quantum currents are well-defined, i.e. their definition is not empty, in this section we give a free field realization of the quantum nonlocal currents. First we introduce two kinds of Heisenberg algebras,
$`[\widehat{b}(t),\widehat{b}(t^{})]={\displaystyle \frac{\mathrm{sinh}\mathrm{}t\mathrm{sinh}\frac{k}{2}\mathrm{}t}{\mathrm{}^2t}}\delta (t+t^{}),`$
$`[\widehat{\lambda }(t),\widehat{\lambda }(t^{})]={\displaystyle \frac{\mathrm{sinh}\mathrm{}t\mathrm{sinh}\frac{k+2}{2}\mathrm{}t}{\mathrm{}^2t}}\delta (t+t^{}).`$ (3.1)
Then defining the following intermediate fields
$`B^+(u)=\text{exp}\{\mathrm{}{\displaystyle _{\mathrm{}}^0}𝑑te^{\frac{k}{4}\mathrm{}t}{\displaystyle \frac{e^{iut}}{\mathrm{sinh}\frac{k}{2}\mathrm{}t}}\widehat{b}(t)\mathrm{}{\displaystyle _0^{\mathrm{}}}𝑑te^{\frac{k}{4}\mathrm{}t}{\displaystyle \frac{e^{iut}}{\mathrm{sinh}\frac{k}{2}\mathrm{}t}}\widehat{b}(t)\},`$
$`B^{}(u)=\text{exp}\{\mathrm{}{\displaystyle _{\mathrm{}}^0}𝑑te^{\frac{k}{4}\mathrm{}t}{\displaystyle \frac{e^{iut}}{\mathrm{sinh}\frac{k}{2}\mathrm{}t}}\widehat{b}(t)+\mathrm{}{\displaystyle _0^{\mathrm{}}}𝑑te^{+\frac{k}{4}\mathrm{}t}{\displaystyle \frac{e^{iut}}{\mathrm{sinh}\frac{k}{2}\mathrm{}t}}\widehat{b}(t)\},`$ (3.2)
$`\mathrm{\Lambda }_+(u)=\text{exp}\{2\mathrm{}{\displaystyle _0^+\mathrm{}}𝑑t{\displaystyle \frac{\mathrm{sinh}\frac{\mathrm{}}{2}t}{\mathrm{sinh}\mathrm{}t}}e^{iut}\widehat{\lambda }(t)\},`$
$`\mathrm{\Lambda }_{}(u)=\text{exp}\{2\mathrm{}{\displaystyle _{\mathrm{}}^0}𝑑t{\displaystyle \frac{\mathrm{sinh}\frac{\mathrm{}}{2}t}{\mathrm{sinh}\mathrm{}t}}e^{iut}\widehat{\lambda }(t)\},`$
$`\beta _+(u)=\text{exp}\{2\mathrm{}{\displaystyle _0^+\mathrm{}}𝑑t{\displaystyle \frac{\mathrm{sinh}\frac{\mathrm{}}{2}t}{\mathrm{sinh}\mathrm{}t}}e^{iut}\widehat{b}(t)\},`$
$`\beta _{}(u)=exp\{2\mathrm{}{\displaystyle _{\mathrm{}}^0}𝑑t{\displaystyle \frac{\mathrm{sinh}\frac{\mathrm{}}{2}t}{\mathrm{sinh}\mathrm{}t}}e^{iut}\widehat{b}(t)\},`$ (3.3)
we have the following formal commutation relations, i.e. relations to be understood in the sense of analytic continuation,
$`\beta _+(u)B^\pm (v)=`$ $`{\displaystyle \frac{i(uv)\pm \frac{k}{4}\mathrm{}\pm \frac{1}{2}\mathrm{}}{i(uv)\pm \frac{k}{4}\mathrm{}\frac{1}{2}\mathrm{}}}B^\pm (v)\beta _+(u),`$
$`B^\pm (u)\beta _{}(v)=`$ $`{\displaystyle \frac{i(uv)\pm \frac{k}{4}\mathrm{}\pm \frac{1}{2}\mathrm{}}{i(uv)\pm \frac{k}{4}\mathrm{}\frac{1}{2}\mathrm{}}}\beta _{}(v)B^\pm (u),`$
$`\mathrm{\Lambda }_+(u)\mathrm{\Lambda }_{}(v)=`$ $`{\displaystyle \frac{\mathrm{\Gamma }(\frac{i(uv)}{2\mathrm{}}+\frac{k+2}{4})\mathrm{\Gamma }(\frac{i(uv)}{2\mathrm{}}+\frac{k+6}{4})}{\mathrm{\Gamma }(\frac{i(uv)}{2\mathrm{}}\frac{k+2}{4})\mathrm{\Gamma }(\frac{i(uv)}{2\mathrm{}}\frac{k2}{4})}}`$
$`\times {\displaystyle \frac{\mathrm{\Gamma }^2(\frac{i(uv)}{2\mathrm{}}\frac{k}{4})}{\mathrm{\Gamma }^2(\frac{i(uv)}{2\mathrm{}}+\frac{k+4}{4})}}\mathrm{\Lambda }_{}(v)\mathrm{\Lambda }_+(u),`$
$`\beta _+(u)\beta _{}(v)=`$ $`{\displaystyle \frac{\mathrm{\Gamma }(\frac{i(uv)}{2\mathrm{}}\frac{k}{4})\mathrm{\Gamma }(\frac{i(uv)}{2\mathrm{}}\frac{k4}{4})}{\mathrm{\Gamma }(\frac{i(uv)}{2\mathrm{}}+\frac{k}{4})\mathrm{\Gamma }(\frac{i(uv)}{2\mathrm{}}+\frac{k+4}{4})}}`$ (3.4)
$`\times {\displaystyle \frac{\mathrm{\Gamma }^2(\frac{i(uv)}{2\mathrm{}}+\frac{k+2}{4})}{\mathrm{\Gamma }^2(\frac{i(uv)}{2\mathrm{}}\frac{k2}{4})}}\beta _{}(v)\beta _+(u),`$
$`{\displaystyle \frac{\mathrm{\Gamma }(\frac{i(uv)}{k\mathrm{}}+\frac{1}{k}1)}{\mathrm{\Gamma }(\frac{i(uv)}{k\mathrm{}}\frac{1}{k}1)}}B^+(u)B^+(v)`$ $`=`$ $`B^+(v)B^+(u){\displaystyle \frac{\mathrm{\Gamma }(\frac{i(vu)}{k\mathrm{}}+\frac{1}{k}1)}{\mathrm{\Gamma }(\frac{i(vu)}{k\mathrm{}}\frac{1}{k}1)}},`$
$`{\displaystyle \frac{\mathrm{\Gamma }(\frac{i(uv)}{k\mathrm{}}+\frac{1}{k})}{\mathrm{\Gamma }(\frac{i(uv)}{k\mathrm{}}\frac{1}{k})}}B^{}(u)B^{}(v)`$ $`=`$ $`B^{}(v)B^{}(u){\displaystyle \frac{\mathrm{\Gamma }(\frac{i(vu)}{k\mathrm{}}+\frac{1}{k})}{\mathrm{\Gamma }(\frac{i(vu)}{k\mathrm{}}\frac{1}{k})}},`$
$`{\displaystyle \frac{\mathrm{\Gamma }(\frac{i(uv)}{k\mathrm{}}\frac{1}{k}+\frac{1}{2})}{\mathrm{\Gamma }(\frac{i(uv)}{k\mathrm{}}+\frac{1}{k}+\frac{1}{2})}}B^+(u)B^{}(v)`$ $`=`$ $`B^{}(v)B^+(u){\displaystyle \frac{\mathrm{\Gamma }(\frac{i(vu)}{k\mathrm{}}\frac{1}{k}+\frac{1}{2})}{\mathrm{\Gamma }(\frac{i(vu)}{k\mathrm{}}+\frac{1}{k}+\frac{1}{2})}}.`$ (3.5)
Using these results, we can get a realization of the quantum nonlocal currents as follows,
$`\mathrm{\Psi }(u)={\displaystyle \frac{1}{\mathrm{}}}:\{\beta _+(u+i{\displaystyle \frac{k+2}{4}}\mathrm{})\mathrm{\Lambda }_+(u+i{\displaystyle \frac{k}{4}}\mathrm{})\beta _{}(ui{\displaystyle \frac{k+2}{4}}\mathrm{})\mathrm{\Lambda }_{}(ui{\displaystyle \frac{k}{4}}\mathrm{})\}B^+(u):,`$
$`\mathrm{\Psi }^{}(u)={\displaystyle \frac{1}{\mathrm{}}}:\{\beta _+(ui{\displaystyle \frac{k+2}{4}}\mathrm{})\mathrm{\Lambda }_+^1(ui{\displaystyle \frac{k}{4}}\mathrm{})\beta _{}(u+i{\displaystyle \frac{k+2}{4}}\mathrm{})\mathrm{\Lambda }_{}^1(u+i{\displaystyle \frac{k}{4}}\mathrm{})\}B^{}(u):.`$ (3.6)
By direct calculation, we can show that the OPE of the above nonlocal currents coincide with the defining relations (2.4).
In this paper a set of rational deformed quantum nonlocal $`SU(2)`$ currents are proposed, bosonization and their relation with central extension of Yangian double in the Drinfeld new realization are also discussed. A related interesting object, i.e. the screening currents for these quantum nonlocal currents, which is important for the calculation of the correlation functions, will be considered in a separate paper.
Acknowledgments: One of the authors (Ding) would like to thanks Profs. S. K. Wang, K. Wu and Z. Y. Zhu for fruitful discussion. Ding and Zhao are supported in part by the ”Natural Science Foundation of China”.
|
warning/0003/nucl-th0003049.html
|
ar5iv
|
text
|
# Invariant Manifolds and Collective Coordinates
## I Introduction
Dynamical systems with invariant manifolds in phase space have been the subject of ongoing research in recent years. Many authors have considered the case of two or more coupled identical systems that are chaotic. On invariant manifolds the subsystems display identical or synchronized motion, and the manifold’s transverse stability is investigated .
An alternative approach is based on the observation that any rotationally invariant system of identical interacting particles possesses low dimensional invariant manifolds in classical phase space . On such manifolds, the classical motion displays largely collective behavior and decouples from more complex single-particle behavior. The importance of a given invariant manifold depends crucially on its stability properties. If the manifold under consideration is sufficiently stable in transverse directions, the quantum system may exhibit wave function scarring or display a strong revival for wave-packets localized to the vicinity of the manifold . These findings may be directly associated with the slow decay of collective motion due to the coupling between collective and single-particle motion.
In this paper we propose suitably adapted coordinate systems that separate collective and single-particle motion on the invariant manifolds mentioned above. Such coordinates clarify the separation of collective and single particle motion and may be useful in several applications. We have in mind (i) the problem of damping and dissipation of collective excitations and the interplay of collective and chaotic motion in atomic nuclei , which is often addressed in the framework of single-particle motion in a time dependent mean field; (ii) multi particle fragmentation of atoms at threshold which evolves over highly symmetric configurations corresponding to invariant manifolds ; (iii) the structural stability of invariant manifolds .
There is a traditional way to introduce collective and single-particle coordinates in interacting many-body systems. Aiming at the description of nuclear vibrations and rotations, Zickendraht introduced a system of collective coordinates in a self-bound many-body system. Three of these coordinates describe the center of mass motion, and six collective degrees of freedom govern the dynamics of the inertia ellipsoid. The remaining coordinates are of single-particle nature. We shall establish the relation of coordinates of invariant manifolds to those defined by Zickendraht. Furthermore we shall show that more complicated collective motion, e.g. shearing modes can be described
This article is divided as follows. In the next section we introduce suitable coordinate systems for interacting many-body systems with invariant manifolds. We give a construction recipe and present a detailed example calculation. As an application we give a potential expansion around an invariant manifold and discuss stability properties. In section IV we make a connection with the Zickendraht coordinates. We present examples where the motion of the inertia ellipsoid corresponds to the motion on an invariant manifold. For such initial conditions the traditional collective motion decouples completely from the single-particle degrees of freedom. We also find that collective coordinates as defined here are capable of other types of motion. Therefore we finally discuss how motion on or near such invariant manifolds could be interpreted as collective motion of a system.
## II Coordinates for invariant manifolds
In this section we present a transformation from Cartesian single particle coordinates in position and momentum space to Cartesian coordinates that are adapted to invariant manifolds. The new coordinates consist of “collective coordinates” that govern the motion on the invariant manifold and of coordinates transversal to this manifold that represent the single-particle aspects.
Consider rotationally invariant systems of $`N`$ identical particles in $`d`$ spatial dimensions ($`d=2`$ or $`d=3`$). The Hamiltonian is invariant under both, the action of the rotation group O($`d`$) and the group of permutations S<sub>N</sub>. One may now take a finite subgroup $`𝒢`$ O($`d`$) with elements $`g`$ and properly chosen permutations $`P(g)`$ such that
$$gP^1(g)(\stackrel{}{p},\stackrel{}{q})=(\stackrel{}{p},\stackrel{}{q}),g𝒢\stackrel{}{p}(p_1,\mathrm{},p_{Nd}),\stackrel{}{q}(q_1,\mathrm{},q_{Nd})$$
(1)
for points $`(\stackrel{}{p},\stackrel{}{q})`$ on some invariant submanifold of phase space. On such a manifold, the action of certain rotations $`g`$ can be canceled by permutations. These permutations clearly form a subgroup isomorphic to $`𝒢`$.
Fig. 1 shows a configuration of four particles in two spatial dimensions that corresponds to a point on an invariant manifold. The operations of elements from the discrete symmetry group $`𝒢=C_{2v}`$ can be undone by suitable permutations of particles. This leads to a collective motion with two degrees of freedom which we shall identify with vibrations.
Fig. 2 shows two spatial configurations of eight (2a) and six (2b) particles, respectively which display a $`D_{4h}`$ symmetry. If initial momenta display the same symmetry the motion on the invariant manifold will have two degrees of freedom. For eight particles the radii of the two circles will oscillate synchronously, and the two circles will vibrate against each other. For the six particles we will have a vibration of the radius of the circle and of the two particles along the vertical axis. We may choose initial momenta to reduce the symmetry group to $`C_{4h}`$ which will allow rotations around the vertical axis and thus add an additional degree of freedom. For eight particles we could alternatively choose initial conditions that are limited to a $`D_4`$ symmetry. Besides the vibrations discussed above this would allow for a shearing motion of the two circles thus yielding again three degrees of freedom. We could also reduce the fourfold rotation axis to a twofold one and obtain $`D_{2h}`$, $`D_2`$ or $`C_{2v}`$ as remaining symmetry groups yielding more collective degrees of freedom. Adding two particles symmetrically onto the principal axis of rotation would also increase the number of degrees of freedom by one. Other reductions of symmetry will yield different invariant manifolds with varying degrees of freedom. We will see this exemplified by explicit construction of coordinates.
We may use the definition (1) directly for the construction of coordinate systems where invariant manifolds correspond to coordinate axis or planes, i.e. non collective coordinates vanish for motion on the invariant manifold. To this purpose we consider the many-body system in Cartesian coordinates in momentum and position space. In what follows we will introduce orthogonal transformations in configuration space only; momenta will be subject to the same transformation.
In a Cartesian coordinate system each element $`g𝒢`$ and each permutation $`P(g)`$ can be represented by an orthogonal matrix $`𝐌_g`$ and $`𝐏_g`$ of dimension $`Nd`$. It is clear that the products $`𝐌_g𝐏_g^\mathrm{T}`$ form a matrix group $``$ that acts onto position and momentum space, respectively. The construction of the coordinate system is now straightforward. Every vector $`\stackrel{}{p}`$ and $`\stackrel{}{q}`$ may be expanded in basis vectors of the irreducible representations (IRs) of $``$ by means of projectors .
$$\mathrm{\Pi }_\nu =\underset{g𝒢}{}\chi _g^{(\nu )}𝐌_g𝐏_g^\mathrm{T}.$$
(2)
Here $`\chi _g^{(\nu )}`$ denotes the character of $`g`$ in the $`\nu `$’th IR. Similar formulae hold for momentum space. The projection onto the identical IR defines the invariant manifold. Note that the identical representation is one-dimensional while the invariant manifolds of interest typically have higher dimensionality. We can find independent vectors on the manifold by projecting from different vectors, but in practice the construction of the independent vectors seems to be unproblematic as we shall see in the example.
A comment on the rotation symmetry is in order. Like any Cartesian coordinates, the coordinates introduced in this article do not explicitly reflect the invariance under rotations. Acting on an invariant manifold, rotations generate a continuous family of equivalent manifolds. Our coordinates, however, single out one particular manifold. For quantum systems, the rotation operator may easily be constructed and used for projection onto subspaces of definite angular momentum.
## III A simple example
We now illustrate the proposed construction explicitly for four particles in two dimensions and a quartic potential, considering the invariant manifold shown in Fig. 1. We shall also expand the potential near the invariant manifold to second order in the transversal coordinates.
The invariant manifold is defined by those points which are invariant under $`=\{E,\sigma _xP_{(12)(34)},\sigma _yP_{(14)(23)},C_2P_{(13)(24)}\}`$, where $`E`$ denotes the identity, $`P`$ a permutation of particles as indicated, $`\sigma `$ a reflection at the axis indicated, and $`C_2`$ a rotation about $`\pi `$. Thus, $`=C_{2v}`$ with four IRs labeled by $`\nu =A_1,B_1,A_2,B_2`$ . Let $`\stackrel{}{q}=(x_1,x_2,x_3,x_4,y_1,y_2,y_3,y_4)`$ denote a coordinate vector in position space ($`x_i,y_i`$ denote the coordinates of the $`i`$’th particle). We have
$`E\stackrel{}{q}`$ $`=`$ $`(x_1,x_2,x_3,x_4,y_1,y_2,y_3,y_4),`$ (3)
$`\sigma _xP_{(12)(34)}\stackrel{}{q}`$ $`=`$ $`(x_2,x_1,x_4,x_3,y_2,y_1,y_4,y_3)`$ (4)
$`C_2P_{(13)(24)}\stackrel{}{q}`$ $`=`$ $`(x_3,x_4,x_1,x_2,y_3,y_4,y_1,y_2)`$ (5)
$`\sigma _yP_{(14)(23)}\stackrel{}{q}`$ $`=`$ $`(x_4,x_3,x_2,x_1,y_4,y_3,y_2,y_1).`$ (6)
Using the character table of $`C_{2v}`$ and the projectors (2) one constructs the following basis vectors corresponding to the IR labeled by
$`A_1`$ $`:`$ $`e_1^{}=(1,1,1,1,0,0,0,0)/2,e_2^{}=(0,0,0,0,1,1,1,1)/2,`$ (7)
$`B_1`$ $`:`$ $`e_3^{}=(1,1,1,1,0,0,0,0)/2,e_4^{}=(0,0,0,0,1,1,1,1)/2,`$ (8)
$`A_2`$ $`:`$ $`e_5^{}=(1,1,1,1,0,0,0,0)/2,e_6^{}=(0,0,0,0,1,1,1,1)/2,`$ (9)
$`B_2`$ $`:`$ $`e_7^{}=(1,1,1,1,0,0,0,0)/2,e_8^{}=(0,0,0,0,1,1,1,1)/2.`$ (10)
The vectors associated with the identical IR $`A_1`$ span the two-dimensional invariant manifold and the vectors associated with the IRs $`B_1,A_2,B_2`$ span the transverse directions.
We now present the orthogonal transformation that transforms the single particle coordinates $`\stackrel{}{q}`$ into the coordinates adapted to the invariant manifold. In our example $`x`$ and $`y`$-components do not mix and we have
$`\left[\begin{array}{c}x_1^{}\\ x_2^{}\\ x_3^{}\\ x_4^{}\end{array}\right]={\displaystyle \frac{1}{2}}\left[\begin{array}{cccc}\hfill 1& \hfill 1& \hfill 1& \hfill 1\\ \hfill 1& \hfill 1& \hfill 1& \hfill 1\\ \hfill 1& \hfill 1& \hfill 1& \hfill 1\\ \hfill 1& \hfill 1& \hfill 1& \hfill 1\end{array}\right]\left[\begin{array}{c}x_1\\ x_2\\ x_3\\ x_4\end{array}\right],\left[\begin{array}{c}y_1^{}\\ y_2^{}\\ y_3^{}\\ y_4^{}\end{array}\right]={\displaystyle \frac{1}{2}}\left[\begin{array}{cccc}\hfill 1& \hfill 1& \hfill 1& \hfill 1\\ \hfill 1& \hfill 1& \hfill 1& \hfill 1\\ \hfill 1& \hfill 1& \hfill 1& \hfill 1\\ \hfill 1& \hfill 1& \hfill 1& \hfill 1\end{array}\right]\left[\begin{array}{c}y_1\\ y_2\\ y_3\\ y_4\end{array}\right].`$ (35)
To illustrate the example and to further demonstrate the usefulness of the newly introduced coordinate system we want to consider the the interacting four-body system with Hamiltonian
$$H=\underset{i=1}{\overset{4}{}}\left((p_{x_i}^2+p_{y_i}^2)/2+16(x_i^2+y_i^2)^2\right)\underset{i<j}{}\left[(x_ix_j)^2+(y_iy_j)^2\right].$$
(36)
This Hamiltonian has been studied previously . In particular, the stability of the invariant manifold displayed in Fig. 1 has been studied by computing the full phase space monodromy matrix of several periodic orbits that are inside the invariant manifold. It was found that several orbits are linearly stable in transverse directions or possess rather small stability exponents. Qualitatively, this may also be understood by studying the Hamiltonian (36) close to the invariant manifold. We therefore use the transformation (35) and expand the potential of Hamiltonian (36) to second order in the transverse directions labeled by ($`ϵ_1,\mathrm{}ϵ_6`$) while keeping the full dependence of the coordinates ($`x,y`$) inside the invariant manifold. One obtains the quadratic form $`\stackrel{}{ϵ}^\mathrm{T}𝐕\stackrel{}{ϵ}`$ where
$`𝐕=\left[\begin{array}{cccccc}12x^2+4y^2& 0& 0& 16xy& 0& 0\\ 0& 12x^2& 0& 0& 8xy& 0\\ 0& 0& 24x^2+8y^2& 0& 0& 16xy\\ 16xy& 0& 0& 8x^2+24y^2& 0& 0\\ 0& 8xy& 0& 0& 12y^2& 0\\ 0& 0& 16xy& 0& 0& 4x^2+12y^2\end{array}\right].`$ (43)
A diagonalization of $`𝐕`$ yields the eigenvalues
$`\lambda _{1,2}`$ $`=`$ $`10x^2+14y^2\pm 2\sqrt{x^4+25y^4+54x^2y^2},`$ (44)
$`\lambda _{3,4}`$ $`=`$ $`14x^2+10y^2\pm 2\sqrt{y^4+25x^4+54x^2y^2},`$ (45)
$`\lambda _{5,6}`$ $`=`$ $`6(x^2+y^2)\pm 2\sqrt{9x^42x^2y^2+9y^4}.`$ (46)
All eigenvalues are non-negative and vanish at the origin $`(x=y=0)`$. Thus, instability may occur only in its vicinity. Though the expansion of a potential around an invariant manifold is no substitute for the computations of Lyapunov exponents or monodromy matrices, it is a first step when estimating stability properties of such manifolds.
## IV Zickendraht’s coordinates and invariant manifolds
Almost thirty years ago Zickendraht introduced a set of collective coordinates to describe nuclear vibrations and rotations, as well as their coupling with single particle motion. We shall discuss to what extent these coordinates correspond to the ones we introduced in the previous sections. On one hand this will allow to identify certain vibrational modes of a many-body system with invariant manifolds. On the other hand we shall also see that our procedure proposes collective movements that are not of the type described easily in Zickendraht’s coordinates.
Following Zickendraht we write the coordinates $`\stackrel{}{r}_i`$ of the $`i^{th}`$ particle in the center of mass system as
$$\stackrel{}{r}_i=s_{i1}\stackrel{}{y_1}+s_{i2}\stackrel{}{y_2}+s_{i3}\stackrel{}{y_3},i=1,\mathrm{},N$$
(47)
where the $`\stackrel{}{y}_i`$ span the inertia ellipsoid and $`s_{ik}`$ are non-collective coordinates which for simplicity we shall call single-particle coordinates. The newly introduced coordinates $`\stackrel{}{y}_i`$ and $`s_{ij}`$ are not independent. The constraints are
$`\stackrel{}{y}_i\stackrel{}{y}_j`$ $`=`$ $`y_iy_j\delta _{ij},i,j=1,2,3`$ (48)
$`{\displaystyle \underset{i=0}{\overset{N}{}}}s_{ij}`$ $`=`$ $`0,j=1,2,3`$ (49)
$`{\displaystyle \underset{i=0}{\overset{N}{}}}s_{ij}s_{ik}`$ $`=`$ $`\delta _{jk},j,k=1,2,3.`$ (50)
The first six equations ensure the orthogonality and normalization of the principal axis of the inertia ellipsoid whereas the next three equations fix the origin at the center of mass system. The last six equations are orthogonality relations of the single-particle coordinates. In the center of mass system, one may therefore characterize the $`N`$-body system by its inertia ellipsoid (e.g. three Euler angles of the principle axis and three moments of inertia) and $`3N9`$ single particle coordinates. The moments of inertia $`I_i`$ are related to the coordinates $`y_i`$ by
$$I_1=m(y_2^2+y_3^2),I_2=m(y_1^2+y_3^2),I_3=m(y_1^2+y_2^2),$$
(51)
where $`m`$ denotes the mass of the particles.
It is interesting to determine those configurations, where the motion of the many-body system may be described in terms of the collective coordinates $`y_i`$ only. While such motion would be restricted to some invariant manifold in phase space it would not obviously be one of those defined by eq. (1). We may however determine invariant manifolds (1) such that the motion on the manifold changes only the inertia ellipsoid of the system and hence may be described entirely by Zickendraht’s collective coordinates $`y_i`$. Two necessary conditions for this a situation are easily stated. First, the number of coordinates on such invariant manifold may not exceed six in the general case and three in the case of pure vibrations. Second, every motion on such an invariant manifold has to change the inertia ellipsoid of the many-body system.
For simplicity let us start with the a system of four particles in two spatial dimensions and the invariant manifold displayed in Fig. 1, i.e.
$`\stackrel{}{r}_1=\left[\begin{array}{c}x\\ y\end{array}\right],\stackrel{}{r}_2=\left[\begin{array}{c}x\\ y\end{array}\right],\stackrel{}{r}_3=\left[\begin{array}{c}x\\ y\end{array}\right],\stackrel{}{r}_4=\left[\begin{array}{c}x\\ y\end{array}\right],`$ (60)
and the momenta are chosen by replacing $`xp_x,yp_y`$. Computation of the moments of inertia yield the collective Zickendraht coordinates $`y_1=2x`$, $`y_2=2y`$. On the invariant manifold the remaining coordinates are given by $`s_{11}=s_{12}=s_{21}=s_{22}=s_{31}=s_{32}=s_{41}=s_{42}=1/2`$. This shows that every motion on the invariant manifold only changes the moments of inertia and therefore decouples from the single-particle motion.
We next consider the example of an eight-body system in three dimensions. Let
$`\stackrel{}{r}_1=\left[\begin{array}{c}x\\ y\\ z\end{array}\right],\stackrel{}{r}_2=\left[\begin{array}{c}y\\ x\\ z\end{array}\right],\stackrel{}{r}_3=\left[\begin{array}{c}x\\ y\\ z\end{array}\right],\stackrel{}{r}_4=\left[\begin{array}{c}y\\ x\\ z\end{array}\right],\stackrel{}{r}_{4+i}=\stackrel{}{r}_i(zz)`$ (73)
denote a configuration restricted to the invariant manifold displayed in Fig. 2 (a) with $`C_{4h}`$ symmetry. (The momenta are chosen by replacing $`xp_x,yp_y,zp_z`$ in eq.(73).) The moments of inertia are $`I_1=I_2=4m(x^2+y^2)+8mz^2,I_3=8m(x^2+y^2)`$ and yield collective coordinates (51) $`y_1^2=y_2^2=4(x^2+y^2),y_3^2=8z^2`$. Since the inertia ellipsoid is symmetric we have a freedom in choosing two of its principle axis. Using
$`\stackrel{}{y}_1=2\left[\begin{array}{c}x\\ y\\ 0\end{array}\right],\stackrel{}{y}_2=2\left[\begin{array}{c}y\\ x\\ 0\end{array}\right],\stackrel{}{y}_3=\sqrt{8}\left[\begin{array}{c}0\\ 0\\ z\end{array}\right].`$ (83)
one obtains constant single-particles coordinates $`s_{11}=s_{31}=s_{51}=s_{71}=s_{22}=s_{42}=s_{62}=s_{82}=1/2,s_{13}=s_{23}=s_{33}=s_{43}=s_{53}=s_{63}=s_{73}=s_{83}=1/\sqrt{8}`$ for the motion on the invariant manifold. Thus, the single-particle motion decouples from the collective motion on the invariant manifold. Similar results hold for the six particle configuration displayed in Fig. 2.
It is also instructive to consider one counterexample. The configuration
$`\stackrel{}{r}_1=\left[\begin{array}{c}x\\ y\\ z\end{array}\right],\stackrel{}{r}_2=\left[\begin{array}{c}y\\ x\\ z\end{array}\right],\stackrel{}{r}_3=\left[\begin{array}{c}x\\ y\\ z\end{array}\right],\stackrel{}{r}_4=\left[\begin{array}{c}y\\ x\\ z\end{array}\right],`$ (96)
$`\stackrel{}{r}_5=\left[\begin{array}{c}x\\ y\\ z\end{array}\right],\stackrel{}{r}_6=\left[\begin{array}{c}y\\ x\\ z\end{array}\right],\stackrel{}{r}_7=\left[\begin{array}{c}x\\ y\\ z\end{array}\right],\stackrel{}{r}_8=\left[\begin{array}{c}y\\ x\\ z\end{array}\right],`$ (109)
displays $`D_4`$ symmetry and differs from configuration (73) by a shearing motion. Like in the previous example, the moments of inertia are given by $`I_1=I_2=4m(x^2+y^2)+8mz^2,I_3=8m(x^2+y^2)`$ and the ellipsoid of inertia is symmetric. However, no choice of the principal axis allows to fulfill eqs. (47) with constant single-particle coordinates $`s_{ij}`$. Therefore, single-particle degrees of freedom depend on collective degrees of freedom and a decoupling does not exist using Zickendraht’s coordinate system. A decoupling is obtained by using the coordinates introduced in this work. However, the collective motion on the appropriate invariant manifold does not correspond to pure vibrations or rotations of the inertia ellipsoid. These findings are interesting e.g. in relation with with the magnetic dipole mode in nuclei since this type of collective behavior is associated with a shearing motion.
## V Discussion
We constructed an orthogonal transformation that maps the Cartesian single particle coordinates of a many-body system to a new Cartesian coordinate system that distinguishes collective and single-particle motion. The collective degrees of freedom govern the motion that is restricted to a low-dimensional invariant manifold and are decoupled from single-particle degrees of freedom on this manifold. We have demonstrated that there are several configurations of few-body systems, where the motion on the invariant manifold corresponds to a vibration or rotation and may be described in terms of Zickendraht’s collective coordinates, but differs when the collective motion goes beyond that. These results are independent of the details of the Hamiltonian of the $`N`$-body system, and are entirely determined by rotational and permutational symmetry.
Using the results of this article as well as those of refs. we can draw the following picture: First it is possible that an invariant manifold is spanned exactly by the vibrational and rotational modes of a few-body system; second such manifolds may be stable or have small instability exponents in transversal directions; third the revival probabilities of wave packets launched on such manifolds are large; last, as a conclusion of these points we may have a collective motion near the manifold whose damping is characterized by the decay rate in transversal direction. We also found that there may be other collective motions; this was displayed in an example of shearing motion, but there can be others such as breathing modes etc. The coincidence of Zickendraht coordinates with our collective coordinates depends on particle numbers; typically they do not span an invariant manifold. This confirms the well-known fact that in general the collective motion in these coordinates does not separate rigorously, but only in some adiabatic approximation.
As we found more general invariant configurations which in turn induce collective coordinates we may hope that these are useful for approximate considerations for larger particle numbers that do include the corresponding invariant manifold in a non-trivial fashion. The construction of appropriate coordinates is an open problem.
## Acknowledgments
T.P. acknowledges the warm hospitality at the Centro Internacional de Ciencias, Cuernavaca, Mexico, where part of this work has been done. This work was partially supported by Dept. of Energy (USA), by DGAPA (UNAM) and by CONACyT (Mexico).
|
warning/0003/hep-ph0003277.html
|
ar5iv
|
text
|
# 𝒪(𝛼²Γ,𝛼³Γ) Binding Effects in Orthopositronium Decay
## I The Effective Hamiltonian
We now proceed to the effective Hamiltonian. We work in the center-of-momentum frame of the electron and positron, and consider only states of orbital angular momentum $`0`$ ($`S`$-states), and spin $`1`$ (triplet- or ortho- states). In order to make use of the threshold results in , we give the photon a small mass, $`\lambda `$, which is taken to $`0`$ after the determination of the local operator coefficients; the three final-state photons are not given a mass.
We begin with Hamiltonian
$$H\frac{p^2}{m}\frac{p^4}{4m^3}+V+iW$$
(1)
where $`V`$ and $`W`$ are hermitian. Potential $`W`$ accounts for the effects of three-photon decay, while $`V`$ accounts for all other interactions. To determine corrections through $`𝒪(\alpha ^2\mathrm{\Gamma })`$, or $`𝒪(\alpha ^3\mathrm{\Gamma })`$, we need retain only non-annihilation terms that contribute to the positronium $`S`$ -state binding energy through $`𝒪(\alpha ^4m)`$, or $`𝒪(\alpha ^5m)`$, respectively. We write $`V`$ as the sum of three terms,
$`V(E)`$ $`=`$ $`V_0+V_{\mathrm{rel}}+V_{\mathrm{rad}}(E),`$ (2)
where the Coulomb potential, $`V_0`$, and the leading relativistic corrections, $`V_{\mathrm{rel}}`$, are given by:
$`l|V_0|k`$ $`=`$ $`{\displaystyle \frac{4\pi \alpha }{|lk|^2+\lambda ^2}}e^{\frac{|lk|^2+\lambda ^2}{2\mathrm{\Lambda }^2}}`$ (3)
$`l|V_{\mathrm{rel}}|k`$ $`=`$ $`\left[{\displaystyle \frac{4\pi \alpha }{|lk|^2+\lambda ^2}}\left({\displaystyle \frac{2}{3m^2}}|lk|^2{\displaystyle \frac{1}{2m^2}}(l^2+k^2)\right){\displaystyle \frac{4\pi \alpha }{2\mathrm{\Lambda }^2}}+{\displaystyle \frac{4\pi \alpha }{4m^2}}{\displaystyle \frac{(l^2k^2)^2}{(|lk|^2+\lambda ^2)^2}}\right]e^{\frac{|lk|^2+\lambda ^2}{2\mathrm{\Lambda }^2}}`$ (5)
$`+{\displaystyle \frac{4\pi \alpha }{2m^2}}e^{\frac{|lk|^2}{2\mathrm{\Lambda }^2}}.`$
We have introduced Gaussian factors to suppress the potentials at high momentum. Typically we take the ultraviolet cutoff $`\mathrm{\Lambda }m`$, although our final decay rate is (and must be) $`\mathrm{\Lambda }`$-independent. With potentials $`V_0`$ and $`V_{\mathrm{rel}}`$ in place, our Hamiltonian correctly reproduces the QED amplitude for one-photon exchange and one-photon virtual annihilation to lowest and first order in $`(v/c)^2`$ (of the electron). $`V_{\mathrm{rad}}(E)`$ accounts for one-loop radiative corrections:
$`l|V_{\mathrm{rad}}(E)|k`$ $`=`$ $`{\displaystyle \frac{8\alpha }{3\pi m^2}}l|e^{\frac{p^2}{4\mathrm{\Lambda }^2}}p^i(p^2/m+V_0E)\mathrm{ln}\left({\displaystyle \frac{m/4}{p^2/m+V_0E}}\right)p^ie^{\frac{p^2}{4\mathrm{\Lambda }^2}}|k`$ (7)
$`+{\displaystyle \frac{14\alpha ^2}{3m^2}}\mathrm{ln}{\displaystyle \frac{|lk|}{m/2}}e^{\frac{|lk|^2}{2\mathrm{\Lambda }^2}}+{\displaystyle \frac{\alpha ^2}{m^2}}\left\{{\displaystyle \frac{74}{15}}+{\displaystyle \frac{2}{3}}\mathrm{ln}2+D\right\}e^{\frac{|lk|^2}{2\mathrm{\Lambda }^2}}.`$
$`V_{\mathrm{rad}}(E)`$ gives in positronium the analogue of the non-recoil Lamb shift and the recoil Salpeter correction in hydrogen. The final term is a local operator accounting for effects at high momentum; parameter $`D`$ is a counterterm which will be determined shortly.
The annihilation of the electron and positron occurs over distances of order $`\mathrm{\Delta }x1/m`$, which are much smaller than wavelengths typical of the electron and positron in the atom, $`\lambda 1/\alpha m`$. Thus the annihilation potential $`W`$ consists entirely of short-distance interactions, which, to the order of interest, can be parameterized for $`S`$-states as follows:
$$l|\mathrm{W}|k=A^{(0)}\left[(1+\alpha A^{(1)}+\alpha ^2A^{(2)}+\alpha ^3A^{(3)})+(B^{(0)}+\alpha B^{(1)})\frac{E}{m}\right]e^{\frac{|lk|^2}{2\mathrm{\Lambda }^2}}.$$
(8)
We will adjust parameters $`A^{(0)}`$, $`A^{(1)}`$, $`A^{(2)}`$ and $`B^{(0)}`$ so that our Hamiltonian reproduces QED results for electron-positron annihilation into three photons. Since determination of $`A^{(3)}`$ and of $`B^{(1)}`$ requires the as-yet unknown $`𝒪(m\alpha ^6)`$ threshold annihilation rate, and the leading term in the momentum expansion for the $`𝒪(m\alpha ^4)`$ rate, respectively, we simply set $`A^{(3)}=B^{(1)}=0`$. Doing so introduces an error in the decay rate of $`𝒪(\alpha ^3\mathrm{\Gamma })`$.
The $`p^4`$ operator in Eq. (1) is ill-defined at high momenta. To regulate this operator, we replace it by an energy dependent potential :
$$p^4m^2(E(V+iW))^2.$$
(9)
Keeping only the relevant parts of this expression, our final effective Hamiltonian is therefore:
$$H(E)=H_0(E)+i\overline{W}(E)$$
(10)
where
$`H_0(E)`$ $`=`$ $`{\displaystyle \frac{p^2}{m}}+V(E){\displaystyle \frac{1}{4m}}(EV_0)^2,`$ (11)
$`\overline{W}(E)`$ $`=`$ $`W+{\displaystyle \frac{1}{2m}}(EV_0)(W_0+W_1).`$ (12)
Potentials $`V`$ (Eqs.(2),(3),(5),(7)), $`V_0`$ (Eq.(3)) and $`W`$ (Eq.(8)) are given above, and
$`l|W_0|k`$ $``$ $`l|W_1|k/A^{(1)}A^{(0)}e^{\frac{|lk|^2}{2\mathrm{\Lambda }^2}}.`$ (13)
Parameter $`D`$ in Eq.(7) is tuned to correctly reproduce the one-loop contribution to $`e\overline{e}e\overline{e}`$ threshold scattering, and is found to be
$$D=\sqrt{\pi }\left[\frac{121}{36}\frac{\mathrm{\Lambda }}{m}9\frac{m}{\mathrm{\Lambda }}+\frac{5}{3}\left(\frac{m}{\mathrm{\Lambda }}\right)^3\right].$$
(14)
The hermitian Hamiltonian $`H_0(E)`$ is accurate through $`𝒪(\alpha ^5m)`$.
The parameters in $`W`$ are determined by considering the imaginary part of the $`S`$-wave scattering amplitude for $`e\overline{e}3\gamma e\overline{e}`$, with electron momentum $`k`$ in the center-of-momentum frame; the optical theorem relates this amplitude to the free-particle annihilation rate. For small $`k`$, the imaginary part of this amplitude can be parameterized as:
$`𝒯(k)`$ $`=`$ $`𝒯_0\{[1+\beta {\displaystyle \frac{k^2}{m^2}}+𝒪\left({\displaystyle \frac{k^4}{m^4}}\right)]+\alpha [2{\displaystyle \frac{m}{\lambda }}+a_0+a_1{\displaystyle \frac{\lambda }{m}}+𝒪\left({\displaystyle \frac{\lambda ^2}{m^2}}\right)]`$ (16)
$`+\alpha ^2[(1+2\mathrm{ln}2){\displaystyle \frac{m^2}{\lambda ^2}}+2a_0{\displaystyle \frac{m}{\lambda }}{\displaystyle \frac{1}{3}}\mathrm{ln}{\displaystyle \frac{m}{\lambda }}+b_0+𝒪\left({\displaystyle \frac{\lambda }{m}}\right)]+𝒪(\alpha ^3)\}.`$
We define $`𝒯(k)`$ using nonrelativistic normalization for the external particles . Parameters $`𝒯_0,\beta ,a_0,a_1,b_0`$ are determined using QED perturbation theory .
We now calculate $`𝒯`$ in our Hamiltonian theory, and adjust the unknown parameters to reproduce the QED result order by order in $`\alpha `$ and $`k^2`$. Matching at lowest order in $`\alpha `$ implies that
$$A^{(0)}=𝒯_0,B^{(0)}=\beta +\frac{m^2}{\mathrm{\Lambda }^2}.$$
(17)
$`A^{(1)}`$ and $`A^{(2)}`$ are determined by matching the $`𝒪(\alpha )`$ and $`𝒪(\alpha ^2)`$ contributions in $`𝒯(0)`$, respectively; we obtain
$`A^{(1)}`$ $`=`$ $`a_0+{\displaystyle \frac{1}{\sqrt{\pi }}}\left[{\displaystyle \frac{4}{3}}\left({\displaystyle \frac{\mathrm{\Lambda }}{m}}\right)+3\left({\displaystyle \frac{\mathrm{\Lambda }}{m}}\right)^1\right]+\left(a_1{\displaystyle \frac{7}{12}}+\left({\displaystyle \frac{\mathrm{\Lambda }}{m}}\right)^2\right){\displaystyle \frac{\lambda }{m}}+𝒪\left({\displaystyle \frac{\lambda ^2}{m^2}}\right),`$ (18)
$`A^{(2)}`$ $`=`$ $`b_02a_1+{\displaystyle \frac{1}{\sqrt{\pi }}}[{\displaystyle \frac{4}{3}}\left({\displaystyle \frac{\mathrm{\Lambda }}{m}}\right)+3\left({\displaystyle \frac{\mathrm{\Lambda }}{m}}\right)^1]a_0+{\displaystyle \frac{1}{3}}\mathrm{ln}{\displaystyle \frac{\mathrm{\Lambda }}{m}}+{\displaystyle \frac{1}{\pi \sqrt{\pi }}}\{[{\displaystyle \frac{44\sqrt{6}}{81}}(\gamma \mathrm{ln}{\displaystyle \frac{2\mathrm{\Lambda }^2}{3m^2}}2)]\left({\displaystyle \frac{\mathrm{\Lambda }}{m}}\right)^3`$ (21)
$`+[{\displaystyle \frac{7}{3}}\mathrm{ln}{\displaystyle \frac{\mathrm{\Lambda }}{m}}+{\displaystyle \frac{56\sqrt{6}}{27}}(\gamma \mathrm{ln}{\displaystyle \frac{2\mathrm{\Lambda }^2}{3m^2}}{\displaystyle \frac{2}{7}}){\displaystyle \frac{37}{15}}+{\displaystyle \frac{1}{3}}\mathrm{ln}2{\displaystyle \frac{7}{6}}\gamma ]\left({\displaystyle \frac{\mathrm{\Lambda }}{m}}\right)\}+({\displaystyle \frac{83}{24\pi }}{\displaystyle \frac{11\sqrt{3}}{12\pi }}+{\displaystyle \frac{11}{48}}\left)\right({\displaystyle \frac{\mathrm{\Lambda }}{m}})^2`$
$`+\left({\displaystyle \frac{25}{2\pi }}{\displaystyle \frac{4\sqrt{3}}{3\pi }}+{\displaystyle \frac{17}{18}}{\displaystyle \frac{5}{6}}\mathrm{ln}2{\displaystyle \frac{1}{3}}\gamma +{\displaystyle \frac{2}{\sqrt{\pi }}}\kappa \right)+\left({\displaystyle \frac{49}{6\pi }}{\displaystyle \frac{3\sqrt{3}}{2\pi }}{\displaystyle \frac{1}{4}}\right)\left({\displaystyle \frac{\mathrm{\Lambda }}{m}}\right)^2+𝒪\left({\displaystyle \frac{\lambda }{m}}\right).`$
Here $`\gamma =\psi (1)=0.577216`$ is the Euler constant and $`\kappa 𝑑x\mathrm{ln}x^1\mathrm{exp}(x^2)\mathrm{erf}(x)^2=0.051428`$. Having determined all the necessary parameters, we can now safely set $`\lambda =0`$.
## II The Decay Rate
Now that our Hamiltonian is completely specified we finally solve for the decay rate, given by the imaginary part of the ground state energy eigenvalue. Note that due to the presence of the cutoff, no divergences occur when calculating matrix elements, and no intricate limiting procedures are necessary to solve the eigenvalue problem — renormalization is automatic. To avoid dealing with nonhermitian matrices we choose to work only to first order in the annihilation potential $`\overline{W}`$; higher-order terms are suppressed by several powers of $`\alpha `$ beyond the precision of interest.
We first solve the eigenvalue problem for $`H_0`$,
$$H_0(E_0)|\psi _0=E_0|\psi _0,$$
(22)
to obtain the ground-state energy and wavefunction. The energy dependence of $`H_0`$ is easily handled by iterating the eigenvalue equation, starting with an approximate energy in $`H_0`$; the answer converges to adequate precision after only a few iterations. The eigenfunctions for our energy dependent Hamiltonian must be normalized so that
$$\psi _0|1\frac{H_0}{E}|\psi _0|_{E=E_0}=1.$$
(23)
Then the decay rate, to first order in $`\overline{W}`$, is
$`\mathrm{\Gamma }`$ $`=`$ $`2\psi _0|\overline{W}(E_0)|\psi _0`$ (24)
$`=`$ $`\mathrm{\Gamma }_0\left\{\left[1+\alpha A^{(1)}+\alpha ^2A^{(2)}+\alpha ^3A^{(3)}+\left({\displaystyle \frac{1}{2}}\left(1+\alpha A^{(1)}\right)+B^{(0)}+\alpha B^{(1)}\right){\displaystyle \frac{E_0}{m}}\right]\alpha ^2M_1+{\displaystyle \frac{1}{2}}(1+\alpha A^{(1)})\alpha ^2M_2\right\}`$ (25)
where $`𝒪\psi _0|𝒪|\psi _0`$, $`\delta _\mathrm{\Lambda }^3(r)(m^3\alpha ^5/8\pi )M_1`$, $`\alpha v_\mathrm{\Lambda }(r)\delta _\mathrm{\Lambda }^3(r)(m^4\alpha ^5/8\pi )M_2`$ and $`\mathrm{\Gamma }_0`$ is the lowest order $`1S`$ decay rate. The cutoff operators $`v_\mathrm{\Lambda }(r)`$ and $`\delta _\mathrm{\Lambda }^3(r)=^2v_\mathrm{\Lambda }(r)/4\pi `$ are defined by Fourier transform:
$$\frac{4\pi }{q^2}e^{\frac{q^2}{2\mathrm{\Lambda }^2}}v_\mathrm{\Lambda }(r)\frac{1}{r}\mathrm{erf}\left(\frac{\mathrm{\Lambda }r}{\sqrt{2}}\right).$$
(26)
The matrix elements $`M_1`$, $`M_2`$, in Eq. (25) can be evaluated for any $`S`$-state of $`H_0`$ and the corresponding decay rate computed using this equation. We used bases consisting of 30 to 60 gaussians, with varying widths, to diagonalize $`H_0`$. The numerical eigenvalues accurately reproduce the $`{}_{}{}^{3}S_{1}^{}`$ spectrum through $`𝒪(\alpha ^5m)`$ . Our ground-state results, for several values of $`\mathrm{\Lambda }`$, are shown in Table I; there we introduce the dimensionless parameter $`X_\mathrm{\Gamma }`$, defined for any $`S`$-state by:
$$\mathrm{\Gamma }(nS)\frac{\mathrm{\Gamma }_0}{n^3}\left[1+\alpha a_0+\alpha ^2\left(\frac{1}{3}(1+\alpha a_0)\mathrm{ln}\alpha +b_02a_1\frac{\beta }{4n^2}\frac{3}{2\pi }\alpha \mathrm{ln}^2\alpha +X_\mathrm{\Gamma }(nS)\right)\right].$$
(27)
This definition anticipates the leading $`\alpha ^2\mathrm{ln}\alpha `$ and $`\alpha ^3\mathrm{ln}^2\alpha `$ contributions, which are correctly reproduced by our numerical analysis. The final results for $`X_\mathrm{\Gamma }`$ are almost independent of $`\mathrm{\Lambda }`$, while the changes in the matrix elements from one $`\mathrm{\Lambda }`$ to the next are two orders of magnitude larger than $`X_\mathrm{\Gamma }`$ itself. Renormalization theory guarantees that $`\mathrm{\Lambda }`$ dependence due to the matrix elements cancels, in the final answer, against $`\mathrm{\Lambda }`$ dependence due to coefficients $`A^{(1)}`$, $`A^{(2)}`$ and $`B^{(0)}`$ in the annihilation potential $`W`$. The residual $`\mathrm{\Lambda }`$ dependence in $`X_\mathrm{\Gamma }`$ is due to our approximations in potentials $`V`$ and $`W`$, the dominant effect being at $`𝒪(\alpha ^3\mathrm{\Gamma })`$ where we have left out the contributions from $`A^{(3)}`$ and $`B^{(1)}`$.
Our calculation is nonperturbative in potential $`V`$ and so automatically includes order $`\alpha ^3\mathrm{\Gamma }`$ (and higher-order) corrections to the decay rate. To facilitate comparison with other calculations, we suppressed these higher-order effects by extrapolating our analysis to $`\alpha =0`$. Our results for $`\mathrm{\Lambda }=m`$ and a range of $`\alpha `$’s are shown in Table II. Upon extrapolation to $`\alpha =0`$, we obtain:
$$X_\mathrm{\Gamma }(n)|_{\alpha 0}=\{\begin{array}{cc}0.913,& n=1S\hfill \\ 2.588,& n=2S\hfill \\ 2.936,& n=3S\hfill \end{array}$$
(28)
where all results are accurate to within $`1`$ in the last digit. The 1$`S`$ result agrees well with the analytic result in .
Since the analysis would be complete through $`𝒪(\alpha ^3\mathrm{\Gamma })`$ with the inclusion of the operators parameterized by $`A^{(3)}`$ and $`B^{(1)}`$, and since neither of these operators generates factors of $`\mathrm{ln}\alpha `$, the complete contribution in the decay rate of $`𝒪(\alpha ^3\mathrm{ln}\alpha \mathrm{\Gamma })`$ is already present, along with that of $`𝒪(\alpha ^3\mathrm{ln}^2\alpha \mathrm{\Gamma })`$. By examining the $`\alpha `$ dependence of $`X_\mathrm{\Gamma }`$, we find:
$$X_\mathrm{\Gamma }(1S)=0.9130.665\alpha \mathrm{ln}\alpha +𝒪(\alpha ).$$
(29)
The coefficient of $`\alpha \mathrm{ln}\alpha `$ is independent of the cutoff $`\mathrm{\Lambda }`$. The non-logarithmic term of $`𝒪(\alpha ^3\mathrm{\Gamma })`$ is cutoff-dependent, since we have neglected contributions from the cutoff-dependent $`A^{(3)}`$ and $`B^{(1)}`$. Using the values for $`a_0`$ and $`\mathrm{\Gamma }_0`$ in Ref. , the $`𝒪(\alpha ^3\mathrm{ln}\alpha \mathrm{\Gamma })`$ contribution amounts to $`2.4\times 10^5\mu s^1`$. Together with a small contribution from five-photon decay of $`0.73\times 10^5\mu s^1`$ , this brings the current theoretical prediction for the decay rate to $`\mathrm{\Gamma }=7.039967(10)\mu s^1`$.
After completing our calculation, we learned of an independent analysis of $`𝒪(\alpha ^3\mathrm{ln}\alpha \mathrm{\Gamma })`$ contributions to positronium decay by Kniehl and Penin . Our result, $`0.665`$ for the coefficient independent of $`a_0`$, disagrees with their result, $`(4/5+8\mathrm{ln}2/3)/\pi =0.8430`$. To verify our analysis, we compared it with published results on $`𝒪(\alpha ^3\mathrm{ln}\alpha )`$ contributions to muonium hyperfine splitting (HFS). The HFS results involve the same operators, with different coefficients, as those contributing at the same relative order in the decay rate. Terms in the orthopositronium decay rate not involving $`a_0`$ arise from second order perturbations of $`V_{\mathrm{rad}}`$ with the leading decay operator:
$$\frac{\delta \mathrm{\Gamma }}{\mathrm{\Gamma }_0}=2\frac{V_{\mathrm{rad}}\stackrel{~}{G}\delta ^3(r)}{\delta ^3(r)}+\frac{V_{\mathrm{rad}}}{E},$$
(30)
where $`\stackrel{~}{G}`$ is the Coulomb Green’s function with the ground state pole removed, and expectation values are taken between unperturbed Coulomb eigenfunctions. $`V_{\mathrm{rad}}`$ can be expressed as
$$V_{\mathrm{rad}}=\frac{2}{3}\alpha ^2𝒪_1+\frac{7}{6}\alpha ^2𝒪_2+\left(\frac{1}{6}\mathrm{ln}2\frac{37}{30}\right)\alpha ^2𝒪_3,$$
(31)
where
$$𝒪_1=\frac{1}{\pi \alpha m_r^2}p^i\left(\frac{p^2}{2m_r}\frac{\alpha }{r}E\right)\mathrm{ln}\frac{m_r/2}{\frac{p^2}{2m_r}\frac{\alpha }{r}E}p^i,l|𝒪_2|k=\frac{1}{m_r^2}\mathrm{ln}\frac{|lk|}{m_r},l|𝒪_3|k=\frac{1}{m_r^2}.$$
(32)
The reduced mass $`m_r`$ equals $`m/2`$ in positronium. The logarithmic parts of the matrix elements for the second-order perturbations can be inferred, for the most part, directly from the HFS papers :
$$\left(2\frac{𝒪_i\stackrel{~}{G}\delta ^3(r)}{\delta ^3(r)}+\frac{𝒪_i}{E}\right)\frac{\alpha }{\pi }\times \{\begin{array}{cc}4\mathrm{ln}^2\alpha 8(\mathrm{ln}2+3/4)\mathrm{ln}\alpha & ,i=1\hfill \\ \mathrm{ln}^2\alpha +(2\mathrm{ln}21)\mathrm{ln}\alpha & ,i=2\hfill \\ 2\mathrm{ln}\alpha & ,i=3\hfill \end{array}$$
(33)
The only exception is the coefficient $`(2\mathrm{ln}21)`$ of $`\mathrm{ln}\alpha `$ for $`𝒪_2`$. A partial analysis in Ref. gives $`(2\mathrm{ln}2+2)`$ in place of $`(2\mathrm{ln}21)`$. We have calculated the full contribution analytically, and find the result shown in Eq. (33. We have also verified the results in Eq.(33) by direct numerical evaluation, using our Gaussian basis set. Non-logarithmic terms of $`𝒪(\alpha ^0,\alpha ^1)`$ are cutoff dependent, but the logarithmic terms were cutoff independent, as expected. Combining our analytic results with those in Ref., our Eq. (29) becomes:
$`X_\mathrm{\Gamma }(1S)`$ $`=`$ $`{\displaystyle \frac{11}{8}}{\displaystyle \frac{2}{3}}\mathrm{ln}2+\left(8\mathrm{ln}2{\displaystyle \frac{229}{30}}\right){\displaystyle \frac{\alpha }{\pi }}\mathrm{ln}\alpha `$ (34)
$`=`$ $`0.91290.6647\alpha \mathrm{ln}\alpha ,`$ (35)
in complete agreement with our numerical analysis.
We thank the authors of for sharing their results with us before they were published. We also thank Patrick Labelle for several discussions. This work was supported by a grant from the National Science Foundation.
|
warning/0003/cond-mat0003158.html
|
ar5iv
|
text
|
# Quantum Phase Diagram of the 𝑡-𝐽_𝑧 Chain Model
0pt0.4pt 0pt0.4pt 0pt0.4pt
## Abstract
We present the quantum phase diagram of the one-dimensional $`t`$-$`J_z`$ model for arbitrary spin (integer or half-integer) and sign of the spin-spin interaction $`J_z`$, using an exact mapping to a spinless fermion model that can be solved exactly using the Bethe ansatz. We discuss its superconducting phase as a function of hole doping $`\nu `$. Motivated by the new paradigm of high temperature superconductivity, the stripe phase, we also consider the effect the antiferromagnetic background has on the $`t`$-$`J_z`$ chain intended to mimic the stripe segments.
Introduction. Phase diagrams of generic models of strongly interacting quantum particles are considered fundamental to understanding the complex physical behavior of cuprate superconductors, heavy fermion, and related compounds. It is rare to encounter situations where unambiguously these diagrams can be completely determined and only a few exceptional cases are exactly solvable. It is a purpose of this paper to show that the $`t`$-$`J_z`$ chain belongs to this latter class of models.
A new paradigm in superconductivity springs up as a consequence of the growing body of experimental evidence suggesting that the quantum state of high temperature superconductors is a striped phase. Unlike conventional conductors where the charge carriers distribute in an spatially homogeneous way, the stripe paradigm assumes that carriers cluster into quasi one-dimensional (1$`d`$) channels. These channels act as domain walls separating different antiferromagnetic (AF) domains. It is remarkable that experiments are consistent with a spin ordering that is $`\pi `$-shifted across the wall , indicating the topological character of these extended defects .
Motivated by this new paradigm Ref. argued that planar models, with appropriate inhomogeneous magnetic terms, breaking translational and local spin $`SU(2)`$ symmetries are appropriate to understand neutron scattering and angle-resolved photoemission spectroscopy experiments in cuprates . It is interesting to understand why spin anisotropies are relevant to obtain substantial pair hole binding and whether the stripes themselves have important superconducting fluctuations. The simplest representation of a stripe segment is realized by a $`t`$-$`J_z`$ chain model.
In this Letter we study the quantum phase diagram of the $`t`$-$`J_z`$ chain for arbitrary spin and sign of $`J_z`$ by using an exact mapping to an attractive spinless fermion model, and solve this problem by using the Bethe ansatz integral equations. We then consider the effect of the AF boundaries on the stripe as an effective confining potential and determine the resulting phase diagram. While a superconducting phase exists in both cases, the superconducting region is more prominent in the latter.
Model Hamiltonian. The Hamiltonian representing the 1$`d`$ $`t`$-$`J_z`$ model with $`L`$ sites (equal to the length of the chain, i.e., lattice constant $`a`$=1) and $`M`$ holes with open boundary conditions (BC) (the thermodynamic limit, $`L,M\mathrm{}`$ with $`\nu =M/L`$ finite, is performed at the end of the calculation), for arbitrary half-integer spin $`S`$, is $`\widehat{H}=\widehat{T}+\widehat{H}_{J_z}`$ with
$`\widehat{T}`$ $`=`$ $`{\displaystyle \frac{t{\displaystyle }}{\alpha =1}}`$ (1)
$`\sigma [S,S]^{L1}\widehat{T}_{\alpha ,\sigma },\widehat{T}_{\alpha ,\sigma }=c_{\alpha \sigma }^{}c_{\alpha +1\sigma }^{}+\mathrm{H}.\mathrm{c}.,`$ (2)
$`\widehat{H}_{J_z}`$ $`=`$ $`J_z{\displaystyle \underset{\alpha =1}{\overset{L1}{}}}S_\alpha ^zS_{\alpha +1}^z,S_\alpha ^z={\displaystyle \underset{\sigma [S,S]}{}}\sigma c_{\alpha \sigma }^{}c_{\alpha \sigma }^{}.`$ (3)
Here, $`c_{\alpha \sigma }^{}`$($`c_{\alpha \sigma }^{}`$) creates(annihilates) a fermion of $`z`$ spin-component $`\sigma `$ in a Wannier orbital centered at $`\alpha `$. The Hilbert space of states corresponds to a constrained space with no doubly occupied sites .
Consider the set of parent states, labeled by the string configuration $`\stackrel{}{𝝈}`$, with $`M`$ holes and $`LM`$ quantum particles, $`|\mathrm{\Phi }_0(\stackrel{}{𝝈})`$, defined as
$$|\mathrm{\Phi }_0(\stackrel{}{𝝈})=|\underset{LM}{\underset{}{\sigma _1\sigma _2\sigma _3\mathrm{}\sigma _{LM}}}\underset{𝑀}{\underset{}{_{}\mathrm{}}},$$
(4)
where $`\sigma _\alpha `$ indicates the $`z`$-component of the spin of the particle at site $`\alpha `$. These states are eigenstates of $`\widehat{H}_{J_z}`$ with energy $`E_M(\stackrel{}{𝝈})=J_z_{\alpha =1}^{LM1}\sigma _\alpha \sigma _{\alpha +1}`$, and $`z`$-component of the total spin $`S_z=_{\alpha =1}^{LM}\sigma _\alpha `$.
¿From a given parent state one can generate a subspace of the Hilbert space $`(\stackrel{}{𝝈})`$ by applying the hopping operators $`\widehat{T}_{\alpha ,\sigma }`$ to the parent state and its descendants, $`|\mathrm{\Phi }_i(\stackrel{}{𝝈})`$,
$$|\mathrm{\Phi }_1(\stackrel{}{𝝈})=\widehat{T}_{LM,\sigma }|\mathrm{\Phi }_0(\stackrel{}{𝝈})$$
(5)
or, in general,
$$|\mathrm{\Phi }_i(\stackrel{}{𝝈})=\widehat{T}_{\alpha ,\sigma }|\mathrm{\Phi }_j(\stackrel{}{𝝈}).$$
(6)
The dimension $`𝒟`$ of the subspace $`(\stackrel{}{𝝈})`$ is $`\left(\genfrac{}{}{0pt}{}{L}{M}\right)`$. Moreover, these different subspaces are orthogonal and are not mixed by the Hamiltonian $`\widehat{H}`$, although they can be degenerate. In the following we will only consider the AF case $`J_z>0`$. At the end we will discuss two important generalizations: the ferromagnetic (FM) $`J_z<0`$ and the arbitrary integer spin hard-core boson cases.
Among all possible initial configurations the one corresponding to the Néel string $`\stackrel{}{𝝈}_0`$ (i.e., $`\sigma _\alpha =(1)^\alpha S`$), which is two-fold degenerate, turns out to be special. We want to show now that for a given number of holes $`M`$ the subspace generated by the Néel parent state, $`(\stackrel{}{𝝈}_0)_0`$, contains the ground state. To this end, one has to realize that the kinetic energy matrix elements $`\mathrm{\Phi }_i(\stackrel{}{𝝈})|\widehat{T}|\mathrm{\Phi }_j(\stackrel{}{𝝈})`$ are the same for the different subspaces $``$. Nonetheless, the magnetic matrix elements $`\mathrm{\Phi }_i(\stackrel{}{𝝈})|\widehat{H}_{J_z}|\mathrm{\Phi }_j(\stackrel{}{𝝈})=\delta _{ij}A_i(\stackrel{}{𝝈})`$ are different for the different subspaces, with $`A_i(\stackrel{}{𝝈}_0)A_i(\stackrel{}{𝝈})`$ . Notice that if one assigns $`\sigma _\alpha =0`$ to the presence of a hole at site $`\alpha `$, then $`A_i(\stackrel{}{𝝈})=J_z_{\alpha =1}^{L1}\sigma _\alpha \sigma _{\alpha +1}`$. Therefore, the Hamiltonian matrices H$`{}_{}{}^{}{}_{i,j}{}^{}`$ (of dimension $`𝒟\times 𝒟`$) in each subspace $``$, consists of identical off-diagonal matrix elements (H$`{}_{}{}^{}{}_{i,j}{}^{}`$ = H$`{}_{i,j}{}^{^{}},ij`$) and different diagonal ones. These hermitian matrices can be ordered according to the increasing value of the energy $`E_M`$ of their parent states (for fixed $`L`$ and $`M`$). For any $`E_M(\stackrel{}{𝝈})<E_M(\stackrel{}{𝝈}^{})`$, $`\mathrm{H}^{^{}}=\mathrm{H}^{}+B`$, where $`B`$ is a positive semidefinite matrix. Then, the monotonicity theorem tells us that
$$E_k(\stackrel{}{𝝈})E_k(\stackrel{}{𝝈}^{})k=1,\mathrm{},𝒟,$$
(7)
where $`E_k(\stackrel{}{𝝈})`$’s are the eigenvalues of $`\mathrm{H}^{}`$ arranged in increasing order ($`E_k(\stackrel{}{𝝈})E_{k+1}(\stackrel{}{𝝈})`$). Therefore, we conclude that the lowest eigenvalue of $`\widehat{H}`$ must be in $`_0`$, is $`E_1(\stackrel{}{𝝈}_0)`$, and is two-fold degenerate.
Spinless Fermion Mapping. The next step consists in showing, within the ground state subspace $`_0`$, that the Hamiltonian $`\widehat{H}`$ maps into an attractive spinless fermion model. If one makes the following identification
$$|\underset{LM}{\underset{}{\mathrm{}}}\underset{𝑀}{\underset{}{_{}\mathrm{}}}|\underset{LM}{\underset{}{_{}\mathrm{}}}\underset{𝑀}{\underset{}{_{}\mathrm{}}},$$
(8)
i.e., any spin ($`c_{\alpha S}^{}`$ or $`c_{\alpha S}^{}`$) maps into a single spinless fermion ($`b_\alpha ^{}`$) in $`_0`$, it is straightforward to show that all matrix elements of H$`^_0`$ are identical to the matrix elements of the interacting quantum lattice gas
$$H=t\underset{\alpha =1}{\overset{L1}{}}(b_\alpha ^{}b_{\alpha +1}^{}+\mathrm{H}.\mathrm{c}.)J_zS^2\underset{\alpha =1}{\overset{L1}{}}n_\alpha n_{\alpha +1}$$
(9)
in the corresponding new basis. In Eq. 9, $`n_\alpha =b_\alpha ^{}b_\alpha ^{}`$.
Quantum Phase Diagram. The attractive spinless fermion model of Eq. 9 certainly has a superconducting phase (i.e., correlation exponent $`K_\rho >1`$) . For arbitrary values of $`J_zS^2`$, $`t`$, and hole density $`\nu `$, the spinless model is equivalent (via a Jordan-Wigner transformation) to a Heisenberg-Ising spin-$`\frac{1}{2}`$ chain (also known as XXZ model) with Hamiltonian
$$H_{\mathrm{xxz}}=\underset{\alpha =1}{\overset{L}{}}J_{}(s_\alpha ^xs_{\alpha +1}^x+s_\alpha ^ys_{\alpha +1}^y)+J_{}(s_\alpha ^zs_{\alpha +1}^z+s_\alpha ^z)$$
(10)
and $`J_{}=2t`$, $`J_{}=J_zS^2`$. In this new representation the spin up(down) density is $`\nu `$($`1\nu `$). $`J_{}=0`$ represents the classical Ising limit while $`J_{}=0`$ is the extreme quantum limit (XY model). In general, the exchange anisotropy parameter $`\mathrm{\Delta }=J_{}/J_{}<0`$ determines the physical nature of the correlations while $`J_{}`$ defines the energy scale.
There is a vast literature on this model that was exactly solved by the Bethe ansatz . For $`|\mathrm{\Delta }|<1`$, solutions of this model belong to the universality class called “Luttinger liquids” , with correlation functions characterized by power laws with non-universal exponents continuously depending on $`\mathrm{\Delta }`$. The correlation exponent $`K_\rho `$ is determined from
$$K_\rho =\frac{\pi }{2}\nu ^2\kappa v_\rho =\pi \sqrt{\frac{\nu ^2\kappa D_c}{2}},$$
(11)
where the Drude weight (or charge stiffness) $`D_c`$ is related to the velocity of charge excitations $`v_\rho `$ by $`D_c=v_\rho K_\rho /\pi `$ with $`D_c=\frac{1}{2}^2e(\mathrm{\Phi })/(\mathrm{\Phi }/L)^2`$ as usual , and $`\kappa `$ is the isothermal compressibility. On the other hand, $`\kappa `$ can be calculated from the variation of the ground state energy per site $`e(\nu )`$ as $`(\nu ^2\kappa )^1=^2e/\nu ^2`$.
At $`\nu =1/2`$, several quantities and properties are known in closed analytic form. There is a Mott transition at $`\mathrm{\Delta }=1`$ (Umklapp scattering becomes relevant at $`\mathrm{\Delta }=1`$, while it is irrelevant for $`|\mathrm{\Delta }|<1`$). Moreover, the exact expressions for $`K_\rho `$ and $`v_\rho `$ can be determined from the Bethe ansatz
$$K_\rho =\frac{\pi }{4(\pi \mu )},v_\rho =\frac{\pi t\mathrm{sin}\mu }{\mu },$$
(12)
which implies $`D_c=\pi t\mathrm{sin}\mu /[4\mu (\pi \mu )]`$ with $`\mathrm{cos}\mu =\mathrm{\Delta }`$. The energy per site ($`|\mathrm{\Delta }|<1`$) is given by
$$e(1/2)=\frac{\mathrm{\Delta }t}{2}2t\mathrm{sin}\mu _0^{\mathrm{}}𝑑x\frac{\mathrm{sinh}(\pi \mu )x}{\mathrm{cosh}\mu x\mathrm{sinh}\pi x}.$$
(13)
Thus, one obtains $`\frac{1}{2}K_\rho 1`$ in the region $`0\mathrm{\Delta }1/\sqrt{2}`$, and $`K_\rho >1`$ (superconducting correlations dominate at large distances) for $`1<\mathrm{\Delta }<1/\sqrt{2}`$. At $`\mathrm{\Delta }=1`$, there is a transition to a phase segregated state ($`\kappa =2\mu /[\pi t\mathrm{sin}\mu (\pi \mu )]`$ diverges).
For $`\mathrm{\Delta }=0`$, the system reduces to a free spinless fermion system with energy per site $`e=\frac{2t}{\pi }\mathrm{sin}(\pi \nu )`$, stiffness $`D_c=e/2`$, and $`\kappa ^1=\pi ^2\nu ^2e`$. This trivially corresponds to $`K_\rho =\frac{1}{2}`$. Also the cases $`\nu 0`$ and $`\nu 1`$ map to free fermions independently of the value of $`\mathrm{\Delta }`$, therefore, $`K_\rho =\frac{1}{2}`$. The value $`\mathrm{\Delta }=1`$ is also special: After the unitary transformation $`s_\alpha ^{x,y}(1)^\alpha s_\alpha ^{x,y}`$, the Hamiltonian $`H_{\mathrm{xxz}}`$ maps into the FM Heisenberg model in a magnetic field $`J_{}`$ (full $`SU(2)`$ symmetry is recovered). Here $`e=\frac{t}{2}+\left(\frac{1}{2}\nu \right)`$, implying the opening of a gap with a diverging $`\kappa `$, i.e., $`\mathrm{\Delta }=1`$ determines the line of phase segregation for all densities $`\nu `$.
Away from $`\nu =1/2`$ and the special limiting cases discussed above, the quantities $`K_\rho `$ and $`v_\rho `$ are obtained from the numerical solution of the Bethe ansatz integral equations . To calculate $`v_\rho `$ one needs to determine hole excitations with a well-defined momentum $`q`$ and energy $`\mathrm{\Delta }e`$ with respect to the ground state $`e`$. We find the velocity of this elementary excitation from $`v_\rho =lim_{q0}d\mathrm{\Delta }e/dq`$, and together with the numerical second order derivative of $`e(\nu )`$, we determine the correlation exponent $`K_\rho `$. For $`|\mathrm{\Delta }|<1`$ the excitations are gapless. The resulting quantum phase diagram is shown in Fig. 1(a). Notice that the largest superconducting region corresponds to $`\nu =1/2`$.
It is interesting to determine the influence of the antiphase domain wall (ADW), associated with each charge, on the spin-spin correlations of the metallic and the superconducting phases. It is well-known that the charge structure factor of the spinless model has a peak at $`k=2k_F=2\pi \nu `$. Since each charge of the effective model carries an ADW, the spin structure factor, $`𝐒(k)`$, will peak at $`k=\pi \pm 2\pi \nu `$ in the metallic phase. In the superconducting phase a broad peak at $`k=\pi `$ is obtained for $`𝐒(k)`$ since the pairs do not carry an ADW.
Effect of AF Boundaries on the Chain. It is also interesting to study the effect of the AF background in which our stripe segments are embedded. This background provides a strong BC that results in an additional attractive (confining) potential for the holes on the stripe. In this way, an enhanced superconducting region is expected. In fact, the influence of the AF background on the stripe is equivalent to the effect a staggered magnetic field (STM) $`B_s`$ (see Fig. 2). Since each hole carries an ADW, the staggered field gives rise to a confining linear potential between the $`\alpha `$ and the $`\alpha +1`$ holes for even $`\alpha `$. Therefore, the effect of a STM is to bind pairs of holes tightly by a string potential. These pairs interact as hard core bosons. In the very dilute limit one expects these bosons to condense at $`T=0`$ for any value of $`B_s/t`$ and $`J_z/t`$ lower than some critical value which gives rise to phase segregation.
In the limit $`B_s/t1`$ the model can be solved analytically for any concentration of hole pairs. In this limit the problem reduces to nearest-neighbors (NN) pairs of holes moving into an AF background (see Fig. 2). Each pair can hop to its NN with an effective hopping $`t_{\mathrm{eff}}=2t^2/(B_sS+J_zS^2)`$. In addition, there is an attractive $`J_zS^2`$ interaction between NN pairs which comes from the second term of Eq. 9. If we map each pair into a spinless particle and each spin into an empty site
$$|\underset{(L,M)}{\underset{}{\mathrm{}}}|\underset{(\stackrel{~}{L},\stackrel{~}{M})}{\underset{}{_{}\mathrm{}}},$$
(14)
the problem reduces to the spinless Hamiltonian Eq. 9. But as each pair is replaced by an effective particle, the number of particles and the length of the chain for the effective spinless problem are
$$\stackrel{~}{L}=L\frac{M}{2}\stackrel{~}{M}=\frac{M}{2},$$
(15)
where $`\stackrel{~}{\nu }=\stackrel{~}{M}/\stackrel{~}{L}=\nu /(2\nu )`$. Here, as in the FM case, we can use closed or open BC. The sign that arises after a cyclic permutation of fermions is absorbed in the BC. For an odd(even) number of fermions these BC are periodic(anti-periodic). This corresponds to $`M=4n+2`$ ($`M=4n`$) of the original problem. The relations between energies and charge velocities are
$$e(\nu )=\stackrel{~}{e}(\stackrel{~}{\nu })\left(1\frac{\nu }{2}\right),v_\rho (\nu )=\stackrel{~}{v}_\rho (\stackrel{~}{\nu })\frac{2}{2\nu },$$
(16)
where $`\stackrel{~}{e}(\stackrel{~}{\nu })`$ and $`\stackrel{~}{v}_\rho (\stackrel{~}{\nu })`$ are the energy per site and charge velocity of the corresponding spinless model of concentration $`\stackrel{~}{v}`$. Therefore, simple algebraic manipulations lead to
$$K_\rho (\nu )=\stackrel{~}{K}_\rho (\stackrel{~}{\nu })(2\nu )^2,$$
(17)
and the phase diagram is depicted in Fig. 1(b).
For completeness, we would like to mention that the mapping of the low energy spectra of the $`t`$-$`J_z`$ model into the spinless Hamiltonian, Eq. 9, is also valid for the FM case, i.e., $`J_z<0`$. In this case, the magnetic background, which is replaced by empty sites in the spinless model, is FM. Notice, however, that the effective spinless model is also attractive ($`\mathrm{\Delta }<0`$). This implies that the dynamics of the charge degrees of freedom in an AF background is the same as in the FM one. But in the latter case, the charges do not carry an ADW. Moreover, the mapping does not depend upon the statistics of the quantum particles. In other words, we could also apply these concepts to constrained quantum particles with integer spin $`S`$, i.e., hard-core bosons. In the large spin $`S`$ limit, the quantum phase diagram of the $`t`$-$`J_z`$ model approaches the one of isotropic $`t`$-$`J`$ Hamiltonian. Notice, however, the qualitative similarity between the phase diagram in Fig. 1(a) and the one for the isotropic spin-$`\frac{1}{2}`$ $`t`$-$`J`$ model obtained numerically . Finally, the solution can be trivially extended to the $`t`$-$`J_z`$-$`V`$ model, where $`V`$ represents a NN density-density interaction. The effect of $`V`$ is simply to renormalize the spinless fermion interaction in Eq. 9. Furthermore, it is simple to prove that there is a family of bilinear-biquadratic spin-1 chain Hamiltonians that can be mapped onto a $`t`$-$`J_z`$-$`V`$ model and, therefore, its low energy physics is exactly solvable .
In summary, we presented the exact quantum phase diagram of the $`t`$-$`J_z`$ chain model for arbitrary spin $`S`$, particle statistics, and sign of the magnetic interaction $`J_z`$. We also exactly determined the phase diagram of a modified $`t`$-$`J_z`$ chain that includes the effects of a strong antiferromagnetic background. A metallic, superconducting and segregated phases characterize these two phase diagrams.
We thank A.A. Aligia, S. Sondhi, and J.E. Gubernatis for useful discussions, and A.A. Aligia for helping us with the Bethe ansatz equations. Work at Los Alamos is sponsored by the US DOE under contract W-7405-ENG-36.
|
warning/0003/hep-ph0003205.html
|
ar5iv
|
text
|
# References
KANAZAWA-99-15
March, 2000
Mass bound of the lightest neutral Higgs scalar in the extra U(1) models
Y. Daikoku <sup>1</sup><sup>1</sup>1e-mail: daikoku@hep.s.kanazawa-u.ac.jp and D. Suematsu <sup>2</sup><sup>2</sup>2e-mail: suematsu@hep.s.kanazawa-u.ac.jp
Institute for Theoretical Physics, Faculty of Science, Kanazawa University,
Kanazawa 920-1192, Japan
The upper mass bound of the lightest neutral Higgs scalar is studied in the $`\mu `$ problem solvable extra U(1) models by using the analysis of the renormalization group equations. In order to restrict the parameter space we take account of a condition of the radiative symmetry breaking and some phenomenological constraints. We compare the bound obtained based on this restricted parameter space with the one of the next to the minimal supersymmetric standard model (NMSSM). Features of the scalar potential and renormalization group equations of the Yukawa couplings among Higgs chiral supermultiplets are rather different between them. They can reflect in this bound.
1. Introduction
Low energy supersymmetry is one of the main subjects of present particle physics. It is considered to solve a weak scale stability problem called the gauge hierarchy problem in the standard model (SM). Although we donot have any direct evidence for it, it has been stressed that the gauge coupling unification occuring in a rather precise way in the minimal supersymmetric extension (MSSM) of the SM may be an encouraging sign for the presence of the low energy supersymmetry. In the MSSM its phenomenology crucially depends on the soft supersymmetry breaking parameters and then it seems to be difficult to make useful predictions unless we know how the supersymmetry breaks down. However, there is an important exception that the lightest neutral Higgs scalar mass cannot be so heavy and it is mainly controled by the feature of the weak scale symmetry breaking . This is not heavily dependent on the feature of the soft supersymmetry breaking parameters at least at the tree level. Thus the knowledge of its possible upper bound is crucial to judge the validity of the low energy supersymmetry from a viewpoint of the energy front of the accelerator experiment. This aspect has been extensively studied taking account of a radiative correction mainly due to a large top Yukawa coupling .
It is well known that there still remains a hierarchy problem called $`\mu `$ problem in the MSSM. Why a supersymmetric Higgsino mixing term parametrized by $`\mu `$ is a weak scale cannot be explained in the MSSM . A simple and promising candidate for its solution is an extension of the MSSM by the introduction of an extra U(1) gauge symmetry and a SM singlet field $`S`$ with a nonzero charge of this extra U(1) . The essential feature of this model is described by the following superpotential
$$W_{U(1)^{}}=\lambda SH_1H_2+kS\overline{g}g+h_tQH_2\overline{T}+\mathrm{},$$
(1)
where $`H_1`$ and $`H_2`$ are usual doublet Higgs chiral superfields and the ellipses stand for the remaining terms in the MSSM superpotential other than the $`\mu `$ term and the top Yukawa coupling. In the second term $`g`$ and $`\overline{g}`$ stand for the extra color triplet chiral superfields which are important to induce the $`\mu `$ scale. In the superpotential $`W_{U(1)^{}}`$ we also explicitly write the top Yukawa coupling because of its importance in the electroweak radiative symmetry breaking as in the case of the MSSM .
The vacuum of these models is parametrized by the vacuum expectation values (VEVs) of Higgs scalar fields such as
$$H_1=\left(\begin{array}{c}v_1\\ 0\end{array}\right),H_2=\left(\begin{array}{c}0\\ v_2\end{array}\right),S=u,$$
(2)
where $`v_1`$ and $`v_2`$ are assumed to be positive and $`v_1^2+v_2^2=v^2((174\mathrm{G}\mathrm{e}\mathrm{V})^2)`$ should be satisfied<sup>1</sup><sup>1</sup>1 In the following discussion we donot consider the spontaneous CP violation. Under this assumption the sign of $`u`$ cannot be fixed freely but it should be dynamically determined by finding the potential minimum.. The vacuum in this model is parametrized by $`\mathrm{tan}\beta =v_2/v_1`$ and $`u`$. The extra U(1) symmetry is assumed to be broken at the region not far from the weak scale by a VEV of the scalar component of $`S`$ because of the radiative effect caused by the second term in $`W_{U(1)^{}}`$ and then the $`\mu `$ scale is induced as $`\mu =\lambda u`$. Thus in this model the sign of $`\mu `$ is fixed as the one of $`u`$ automatically.
This extra U(1) symmetry forbids a bare $`\mu `$ term in the superpotential and simultaneously makes the model free from the massless axion and tadpole problems. These features seem to make this model more promising than the next to the minimal supersymmetric standard model (NMSSM) which is similar to this extra U(1) model but is extended only by a SM singlet chiral superfield $`S`$ with the superpotential
$$W_{\mathrm{NMSSM}}=\lambda SH_1H_2+\frac{1}{3}\kappa S^3+h_tQH_2\overline{T}+\mathrm{}.$$
(3)
It is also interesting that this kind of extra U(1) models can be often obtained as the effective models of a lot of superstring models . Various interesting features of this type of models have been studied in many works by now \[9$``$12\]. Among the phenomenology of these models the lightest neutral Higgs scalar mass is also an important target for the detailed investigation. Of course, also in these models the lightest neutral Higgs scalar can be expected to be generally not so heavy. The interesting point is that its upper bound can be calculable with no dependence on the soft supersymmetry breaking parameters at least at the tree level as in the case of the NMSSM \[13$``$17\]. The dependence on the soft supersymmetry breaking parameters comes in through the loop correction mainly due to the large top Yukawa coupling and the second term of Eq. (1).
In this paper we estimate the upper bound $`m_{h^0}`$ of this lightest neutral Higgs scalar mass on the correct vacuum. The correct vacuum is determined as the radiatively induced minimum of the effective potential in the suitable parameter space. In this approach we use the one-loop effective potential and solve the relevant renormalization group equations (RGEs) numerically. We will pay our attention on the comparison of this upper bound with the one of the NMSSM within the phenomenologically allowable parameter region. In the NMSSM it has been known through many works that the triviality bound of a Yukawa coupling $`\lambda `$ of Higgs chiral superfields strictly control the upper bound of the lightest neutral Higgs mass . Our approach is somehow different from this usual one. We find the phenomenologically acceptable parameter subspace in the rather wide parameter space by taking account of the radiative symmetry breaking condition and some phenomenological conditions such as the chargino mass and the charged Higgs scalar mass etc.. The estimation of the upper mass bound of the lightest neutral Higgs scalar is carried out in this restricted parameter subspace. Although the result of this approach is necessarily dependent on the assumption for the soft supersymmetry breaking parameters, we consider that it is possible to obtain the useful results by studying the wide region of the parameter space.
2. Extra U(1) models
In this section we discuss more detailed features of the extra U(1) models and give the basis of the present study. Since the NMSSM is well known and discussed in many papers, it is cinvenient to explain the points by using the extra U(1) models. The superpotential of our considering extra U(1) models is defined by Eq. (1). Soft supersymmetry breaking parameters are introduced as
$`_{\mathrm{soft}}`$ $`=`$ $`{\displaystyle \underset{i}{}}m_{\varphi _i}^2|\varphi _i|^2+({\displaystyle \underset{a}{}}{\displaystyle \frac{1}{2}}M_a\overline{\lambda }_a\lambda _a+\mathrm{h}.\mathrm{c}.)`$ (4)
$`+`$ $`(A_\lambda \lambda SH_1H_2+A_kkS\overline{g}g+A_th_tQH_2\overline{T}+\mathrm{h}.\mathrm{c}.),`$
where the first two terms are mass terms of the scalar component $`\varphi _i`$ of each chiral supermultiplet and of gauginos $`\lambda _a`$. We use the same notation for the scalar component as the one of the chiral superfield to represent the trilinear scalar couplings in the last parentheses. Other freedoms remaining in the models are extra matter contents and a type of extra U(1). On these points we confine our study into the typical extra U(1) models derived from E<sub>6</sub>, which are listed in Table 1. At the TeV region they are assumed to have only one extra U(1) symmetry which is broken only by the VEV of $`S`$ and give a solution to the $`\mu `$ problem . As discussed in Ref. for the case of NMSSM, the extra matter contents affect indirectly the low energy value of the Yukawa coupling $`\lambda `$ through the influence on the running of the top Yukawa coupling. This is rather important to estimate the Higgs mass bound. In the present model such kind of effects on the Yukawa couplings may also be expected but its effect is more complicated than the NMSSM as discussed later. If we introduce the extra field contents arbitrarily, the cancellation of the gauge anomaly may require to introduce the additional fields which again affect the running of Yukawa coupling $`\lambda `$ and so on. Thus for the estimation of the Higgs mass bound it is important to fix the matter contents in the anomaly free way in the present study.
As the matter contents we assume the MSSM contents and additional extra matter fields
$$\left[3(Q,\overline{U},\overline{D},L,\overline{E})+(H_1,H_2)\right]_{\mathrm{MSSM}}+3(g,\overline{g})+2(H_1,H_2)+3(S)+3(N),$$
which can be derived from three 27s of E<sub>6</sub> shown in Table 1. This set satisfies the anomaly free conditions. We can also add extra fields to these in the form of vector representations constructed from the fields listed in Table 1. Here we consider the following two cases as the additional extra chiral superfields
$$(\mathrm{A})(H_a)+(H_a^{}),(\mathrm{B})(g+H_a+H_b)+(g^{}+H_a^{}+H_b^{}),$$
where $`a,b=1`$ or 2 and the fields in the second parentheses come from 27 of $`E_6`$. At least on the sector of SU(3)$`{}_{C}{}^{}\times `$SU(2)$`{}_{L}{}^{}\times `$U(1)<sub>Y</sub> these matter contents are the same as the one of \[MSSM + $`n(\mathrm{𝟓}+\mathrm{𝟓}^{})`$\] where $`\mathrm{𝟓}`$ and $`\mathrm{𝟓}^{}`$ are the representations of the usual SU(5). The case (A) corresponds to $`n=3`$ and (B) to $`n=4`$. The $`n=3`$ is the critical value for the one-loop $`\beta `$-function of SU(3). It makes this one-loop $`\beta `$-function be zero. The interesting point of these field contents is that the unification scale of SU(3)$`{}_{C}{}^{}\times `$SU(2)$`{}_{L}{}^{}\times `$U(1)<sub>Y</sub> is not shifted from the MSSM one. The $`n=4`$ case saturates the $`\beta `$-function for the pertubative running of gauge couplings up to the unification scale $`3\times 10^{16}`$ GeV. Although this addition seems to be artificial, this type of spectrum can be expected in the Wilson line breaking scenario of the E<sub>6</sub> type superstring model. We use these contents to compare the feature between the NMSSM and our extra U(1) models<sup>2</sup><sup>2</sup>2We should note that if every extra U(1) is broken near the unification scale, these models are equal to the NMSSM with the equivalent extra matters as discussed in Ref. ..
The existence of multi-generation extra fields brings an ambiguity in Eq. (1). The coupling $`\lambda `$ and $`k`$ can have generation indices for extra fields such as $`S`$, $`H_1`$, $`H_2`$, $`g`$ and $`\overline{g}`$. On this point we make the following assumption to make the argument simple.<sup>3</sup><sup>3</sup>3 As far as we use Eq. (7) for the upper bound of the lightest neutral Higgs mass, this assumption seems to be reasonable. This assumption affects the RGEs of some parameters and also one-loop correction to the Higgs scalar mass. Other cases will be discussed later. Only one $`S`$ can have the couplings in Eq. (1) and one pair of $`(H_1,H_2)`$ corresponding to the one of the MSSM alone gets the VEVs. Extra colored singlets $`(g_i,\overline{g}_i)`$ have a diagonal coupling to this $`S`$ as $`k_iSg_i\overline{g}_i`$, where all the coupling constant $`k_i`$ show the same behavior in the RGEs because $`(g_i,\overline{g}_i)`$ are completely symmetric for the generation index $`i`$ in the models. The fermion components of $`g_i`$ and $`\overline{g}_i`$ can get mass through this coupling.<sup>4</sup><sup>4</sup>4 If we change the charge assignment for some fields, some extra fields can be heavy at the intermediate scale as discussed in . Although this kind of possibility can be realized in the $`\xi _\pm `$ model, we donot consider it in this paper. On the other hand, the fermion components of the remaining $`S`$ which donot couple to the usual Higgsinos in $`H_1`$ and $`H_2`$ can get their masses through the one-loop correction. From a viewpoint of the model construction, the serious phenomenological problem will be how the fermion components in other remaining extra matter fields can get their masses. Although they can be generally massive through the gaugino mediated one-loop diagrams, their magnitude seems not to be enough to satisfy the phenomenological constraints. In the $`\xi _{}`$ model given in Table 1 we can introduce the intermediate scale through the D-flat direction of $`N`$ and $`N^{}`$ whose exsistence does not affect our discussion in the later part of this paper. If this is the case, they can have the weak scale masses through the nonrenormalizable interactions in the superpotential such as $`{\displaystyle \frac{1}{M_{\mathrm{pl}}}}NN^{}gg^{}`$. Although this is phenomenologically important, it can be improved by the suitable extension without changing the following results and thus we donot get involved in this point further here.
In our considering models the tree level scalar potential including the soft supersymmetry breaking terms can be written as
$`V_0`$ $`=`$ $`{\displaystyle \frac{1}{8}}\left(g_2^2+g_1^2\right)\left(|H_1|^2|H_2|^2\right)^2+\left(|\lambda SH_1|^2+|\lambda SH_2|^2\right)`$ (5)
$`+`$ $`m_1^2|H_1|^2+m_2^2|H_2|^2(A_\lambda \lambda SH_1H_2+\mathrm{h}.\mathrm{c}.)`$
$`+`$ $`{\displaystyle \frac{1}{8}}g_E^2\left(Q_1|H_1|^2+Q_2|H_2|^2+Q_S|S|^2\right)^2+\lambda ^2|H_1H_2|^2+m_S^2|S|^2,`$
where $`Q_1`$, $`Q_2`$ and $`Q_S`$ are the extra U(1) charges of $`H_1`$, $`H_2`$ and $`S`$, respectively. The first two lines are found to have the corresponding terms in the MSSM if we remind the fact that $`\mu `$ is realized as $`\mu =\lambda u`$. The third line contains new ingredients. Its first term is a D-term contribution of the extra U(1) and $`g_E`$ stands for its gauge coupling constant.
Potential minimum condition for Eq. (5) can be written as,
$`m_1^2={\displaystyle \frac{1}{4}}(g_2^2+g_1^2)(v_1^2v_2^2){\displaystyle \frac{1}{4}}g_E^2Q_1(Q_1v_1^2+Q_2v_2^2+Q_Su^2)\lambda ^2(u^2+v_2^2)+\lambda A_\lambda u{\displaystyle \frac{v_2}{v_1}},`$
$`m_2^2={\displaystyle \frac{1}{4}}(g_2^2+g_1^2)(v_1^2v_2^2){\displaystyle \frac{1}{4}}g_E^2Q_2(Q_1v_1^2+Q_2v_2^2+Q_Su^2)\lambda ^2(u^2+v_1^2)+\lambda A_\lambda u{\displaystyle \frac{v_1}{v_2}},`$
$`m_S^2={\displaystyle \frac{1}{4}}g_E^2Q_S(Q_1v_1^2+Q_2v_2^2+Q_Su^2)\lambda ^2(v_1^2+v_2^2)+\lambda A_\lambda {\displaystyle \frac{v_1v_2}{u}}.`$ (6)
This constrains the soft SUSY breaking masses of Higgs scalars around the weak scale. As the second derivative of $`V_0`$ in Eq. (5) we can derive the mass matrix of the CP-even neutral Higgs scalar sector which is composed of three neutral components $`H_1^0,H_2^0`$ and $`S`$. The goodness of this treatment has been discussed in the MSSM case and we follow this argument. If we note the fact that the smallest eigenvalue of any matrix is always smaller than any diagonal elements, we can obtain the tree level upper bound of this lightest Higgs scalar in an independent way of the soft supersymmetry breaking parameters by transforming the basis into the suitable one. This upper bound can be written as
$$m_{h^0}^{(0)2}m_Z^2\left[\mathrm{cos}^22\beta +\frac{2\lambda ^2}{g_1^2+g_2^2}\mathrm{sin}^22\beta +\frac{g_E^2}{g_1^2+g_2^2}\left(Q_1\mathrm{cos}^2\beta +Q_2\mathrm{sin}^2\beta \right)^2\right],$$
(7)
where we used the potential minimization condition (6). The first two terms correspond to the ones of the NMSSM in which their behavior has been studied in many works . The running of a coupling constant $`\lambda `$ and its triviality bound have been shown to be crucially dependent on the extra matters . The extra U(1) effect appears through the last term which is its D-term contribution. Equation (7) can show the different $`\mathrm{tan}\beta `$ dependence from the one in the MSSM depending on the value of $`\lambda `$ and also the type of extra U(1). In the case of MSSM the upper bound of the lightest neutral Higgs mass always increases with $`\mathrm{tan}\beta `$ in the region of $`\mathrm{tan}\beta >1`$. If $`\lambda {}_{}{}^{<}0.6`$, the present model also shows the same behavior. On the other hand, for the same region of $`\mathrm{tan}\beta `$ its upper bound can decrease with increasing $`\mathrm{tan}\beta `$ when $`\lambda {}_{}{}^{>}0.6`$ which does not depend on the model so heavily. The NMSSM shows the similar feature, which can be seen in Ref. . Although this may be potentially altered by the radiative correction, it is one of the typical features coming from the $`\lambda SH_1H_2`$ in these models different from the MSSM.
We should note that the bound formula Eq. (7) is applicable only in the case of $`uv_1,v_2`$ <sup>5</sup><sup>5</sup>5 In the case of $`u<v_1,v_2`$ the diagonal element corresponding to $`S`$ can be smaller than the right-hand side of Eq. (7). In such a case we cannot use Eq. (7) as the bound of the lightest neutral Higgs mass. We will exclude it from our study.. In the extra U(1) model the value of $`u`$ can be constrained from below by the conditions on the mass of this extra U(1) gauge boson and its mixing with ordinary $`Z^0`$. As far as we donot consider the special situation such as $`\mathrm{tan}^2\beta Q_1/Q_2`$ under which the mixing with the ordinary $`Z^0`$ is negligible, the hierarchical condition $`u>v_1,v_2`$ should be imposed to satisfy the phenomenological constraints on the extra $`Z^{}`$ mass and its mixing with ordinary $`Z^0`$ . In the sufficiently large $`u`$ case $`\lambda `$ may be constrained into a limited range required by the successful radiative symmetry breaking at the weak scale so that $`\lambda u(\mu )`$ takes a suitable value. In the NMSSM this kind of constraint on $`\lambda `$ is expected to be weaker than the one of the extra U(1) model since $`u`$ has no phenomenological constraint at this stage. Anaway, we need the RGE study to check whether $`\lambda `$ can be constrained in a substantial way by this condition.
3. The comparison of extra U(1) models and the NMSSM
It is useful to discuss some qualitative features of the extra U(1) models and the NMSSM in more detail before comparing the mass bound of the lightest neutral Higgs in both models. Although the extra U(1) models and the NMSSM have the similar feature related to the $`\mu `$ term, they are expected to show rather different behavior in the running of Yukawa couplings $`k`$, $`\kappa `$ and $`\lambda `$. The top Yukawa coupling has the same one-loop RGE in both models as,
$$\frac{dh_t}{d\mathrm{ln}\mu }=\frac{h_t}{16\pi ^2}\left(6h_t^2+\lambda ^2\frac{16}{3}g_3^2\right).$$
(8)
In the present field contents $`g_3`$ takes larger value at $`M_X`$ than the one of the MSSM. Even if the initial value of $`h_t`$ takes the large value like $`O(1)`$, the $`\beta `$-function in Eq. (8) can be small due to the cancellation between a $`h_t`$ term and a $`g_3`$ term. As a result, $`h_t`$ tends to stay near its initial value at the intermediate scale independently whether it starts from a large value or a small value. This feature is shared by both models. On the other hand, the one-loop RGEs of $`\kappa `$, $`k`$ and $`\lambda `$ are largely different from each other. They are witten as, in the NMSSM,
$`{\displaystyle \frac{d\kappa }{d\mathrm{ln}\mu }}={\displaystyle \frac{\kappa }{16\pi ^2}}\left(6\kappa ^2+6\lambda ^2\right),`$
$`{\displaystyle \frac{d\lambda }{d\mathrm{ln}\mu }}={\displaystyle \frac{\lambda }{16\pi ^2}}\left(3h_t^2+2\kappa ^2+4\lambda ^2\right).`$ (9)
and in the extra U(1) model,
$`{\displaystyle \frac{dk}{d\mathrm{ln}\mu }}={\displaystyle \frac{k}{16\pi ^2}}\left((3N_g+2)k^2+2\lambda ^2{\displaystyle \frac{16}{3}}g_3^2\right),`$
$`{\displaystyle \frac{d\lambda }{d\mathrm{ln}\mu }}={\displaystyle \frac{\lambda }{16\pi ^2}}\left(3h_t^2+3N_gk^2+4\lambda ^2\right),`$ (10)
where $`N_g`$ is a number of the pair of the singlet colored fields $`g`$ and $`\overline{g}`$ which have a coupling to $`S`$. In these RGEs we neglect the effect of gauge couplings $`g_2`$, $`g_1`$ and $`g_E`$ <sup>6</sup><sup>6</sup>6 In these equations we cannot find the fixed ratio point other than $`k=0`$ or $`\lambda =0`$ as far as $`N_g0`$ even if we ignore $`g_3`$. This is very different situation from the NMSSM which has been discussed in Ref. .. At first we consider the running behavior of $`\kappa `$ and $`k`$. Since $`k`$ has an effect of $`g_3`$, it can be rather larger at the intermediate scale than $`\kappa `$ which has no such effect and rapidly decreases according to lowering energy. This is important to determine the value of $`u`$ realized in both models, which are mainly determined by $`m_S^2`$ at the low energy region. They are controled by the one-loop RGE as
$$\frac{dm_S^2}{d\mathrm{ln}\mu }=\frac{1}{8\pi ^2}\left(2\kappa ^2(3m_S^2+A_k^2)+2\lambda ^2(m_S^2+m_{H_1}^2+m_{H_2}^2+A_\lambda ^2)\right)$$
(11)
in the NMSSM and
$$\frac{dm_S^2}{d\mathrm{ln}\mu }=\frac{1}{8\pi ^2}\left(3N_gk^2(m_S^2+m_g^2+m_{\overline{g}}^2+A_k^2)+2\lambda ^2(m_S^2+m_{H_1}^2+m_{H_2}^2+A_\lambda ^2)\right),$$
(12)
in the extra U(1) models. The larger $`k`$ compared with $`\kappa `$ makes $`m_S^2`$ much more negative in the extra U(1) models. The larger value of $`u`$ is expected in the extra U(1) models if we remind Eq. (6).
As easily seen from the RGE of $`\lambda `$ in the extra U(1) model, the running of $`\lambda `$ is made fast by the existence of the second term of Eq. (1) which is needed for the successful radiative symmetry breaking of these models . The one-loop $`\beta `$-function of the coupling $`k`$ has a contribution of $`g_3`$ differently from the case of $`\kappa `$ in the NMSSM. If we start $`k`$ and $`\kappa `$ from the large values at the unification scale, this feature can keep $`k`$ rather large at the intermediate region and then the running of $`\lambda `$ can be made fast by its effect compared with the one of $`\kappa `$ in the NMSSM. This feature tends to make the value of $`\lambda `$ at the low energy scale smaller compared with the NMSSM case if the same initial value is adopted at least. However, the initial value of $`k`$ and $`\kappa `$ should be controled from the requirement of the radiative symmetry breaking from our view point since they play an important role in this phenomenon. We need the numerical analysis to study this aspect in more quantitative way. The extra matter effects on the RGEs are also rather different between the NMSSM and the extra U(1) models. As far as all the couplings are within the perturbative regime, the larger number of extra matter fields make the gauge couplings at the unification scale larger. As pointed out in , in the NMSSM this indirectly makes the low energy value of $`\lambda `$ larger through the smallness of $`h_t`$ at the intermediate scale whose $`\beta `$-function in Eq. (8) is kept small there. On the other hand, in the extra U(1) models the runnings of $`k`$ and $`\lambda `$ are simultaneously affected by the extra matters in both direct and indirect manner, as is easily seen in Eq. (10).
We know from these considerations that the resulting low energy values of $`\lambda `$ and $`u`$ are rather different in both models. We should note that these values affect the upper bound of the lightest neutral Higgs scalar mass. Although Eq. (7) shows $`\lambda `$ is crucial to determine the tree level bound, $`u`$ is essential to determine the magnitude of the one-loop effect, especially in the extra U(1) models. The radiative correction to Eq. (7) can be taken into account based on the one-loop effective potential. It is well-known that the one-loop contribution to the effective potential can be written as
$$V_1=\frac{1}{64\pi ^2}\mathrm{Str}^4\left(\mathrm{ln}\frac{^2}{\mathrm{\Lambda }^2}\frac{3}{2}\right),$$
(13)
where $`^2`$ is a matrix of the squared mass of the fields contributing to the one-loop correction and $`\mathrm{\Lambda }`$ is a renormalization point. In the usual estimation of the lightest neutral Higgs mass in the NMSSM the top and stop contributions to $`V_1`$ are mainly considered as the relevant fields because of their large Yukawa coupling. However, in the study of the extra U(1) models $`k`$ is rather large and then we should also take account of the effect on $`^2`$ from the extra singelt colored chiral superfields $`g`$ and $`\overline{g}`$ which have a coupling with $`S`$. A mass matrix of the stops is written as
$$\left(\begin{array}{cc}\stackrel{~}{m}_Q^2+h_t^2v_2^2& h_tv_2(A_t+\lambda u\mathrm{cot}\beta )\\ h_tv_2(A_t+\lambda u\mathrm{cot}\beta )& \stackrel{~}{m}_{\overline{T}}^2+h_t^2v_2^2\end{array}\right),$$
(14)
and the one of the s-$`g`$quarks is expressed as
$$\left(\begin{array}{cc}\stackrel{~}{m}_g^2+k^2u^2& A_kku+\lambda kv_1v_2\\ A_kku+\lambda kv_1v_2& \stackrel{~}{m}_{\overline{g}}^2+k^2u^2\end{array}\right),$$
(15)
where $`\stackrel{~}{m}_{Q,\overline{T}}^2,\stackrel{~}{m}_{g,\overline{g}}^2`$ and $`A_t,A_k`$ are soft supersymmetry breaking parameters. Here a D-term contribution is neglected as it has been done in many previous investigations of the MSSM . Mass eigenvalues of these mass matrices are respectively expressed as,
$`\stackrel{~}{m}_{t_i}^2={\displaystyle \frac{1}{2}}(\stackrel{~}{m}_Q^2+\stackrel{~}{m}_{\overline{T}}^2)+h_t^2v_2^2\pm \sqrt{{\displaystyle \frac{1}{4}}(\stackrel{~}{m}_Q^2\stackrel{~}{m}_{\overline{T}}^2)+h_t^2v_2^2(A_t+\lambda u\mathrm{cot}\beta )^2},`$
$`\stackrel{~}{m}_{g_i}^2={\displaystyle \frac{1}{2}}(\stackrel{~}{m}_g^2+\stackrel{~}{m}_{\overline{g}}^2)+k^2u^2\pm \sqrt{{\displaystyle \frac{1}{4}}(\stackrel{~}{m}_g^2\stackrel{~}{m}_{\overline{g}}^2)+(A_kku+\lambda kv_1v_2)^2}.`$ (16)
If we estimate the upper bound of the lightest Higgs mass in the same procedure as the one used to obtain Eq. (7) by minimizing the one-loop effective potential $`V_{\mathrm{eff}}=V_0+V_1`$, the following one-loop correction should be added to the right-hand side of Eq. (7):
$$\mathrm{\Delta }m_{h^0}^2=\frac{1}{2}\left(\frac{^2V_1}{v_1^2}\frac{1}{v_1}\frac{V_1}{v_1}\right)\mathrm{cos}^2\beta +\frac{1}{2}\frac{^2V_1}{v_1v_2}\mathrm{sin}2\beta +\frac{1}{2}\left(\frac{^2V_1}{v_2^2}\frac{1}{v_2}\frac{V_1}{v_2}\right)\mathrm{sin}^2\beta .$$
(17)
From these we find that $`u`$ can crucially affect to the mass bound through the one-loop effect of $`g`$quark sector in the extra U(1) models. This addtional effect cannot be escapable as far as the occurence of the radiative symmetry breaking is required.
It may also be important to take account of the difference in both models coming from some phenomenological constraints, in particular, the ones related to $`\lambda `$ and $`u`$. Although this kind of constraints depend on the values of soft supersymmetry breaking parameters, it may be useful to improve the upper bound estimation based on the triviality bound of $`\lambda `$. We should remind the fact that the chargino mass, the charged Higgs mass and squark masses are dependent on $`\lambda `$ and $`u`$ . The chargino and the charged Higgs scalar have the same constituents as the MSSM. However, they have a different mass formulas from the MSSM . In both models the chargino mass is expressed as
$`m_{\chi ^\pm }={\displaystyle \frac{1}{2}}\left(\lambda ^2u^2+2m_W^2+M_2^2\right)`$
$`\pm \sqrt{{\displaystyle \frac{1}{4}}\left(2m_W^2\mathrm{cos}2\beta +\lambda ^2u^2M_2^2\right)^2+2m_W^2\left(\lambda u\mathrm{sin}\beta +M_2\mathrm{cos}\beta \right)^2},`$ (18)
where $`m_W`$ and $`M_2`$ represent the W boson and the gaugino $`\lambda _2^\pm `$ masses. The charged Higgs scalar mass has the different mass formula between both models. In the extra U(1) models it is expressed as
$$m_{H^\pm }^2=m_W^2\left(1\frac{2\lambda ^2}{g_2^2}\right)+\frac{2A_\lambda \lambda u}{\mathrm{sin}2\beta },$$
(19)
while in the NMSSM it is written as
$$m_{H^\pm }^2=m_W^2\left(1\frac{2\lambda ^2}{g_2^2}\right)+\frac{2(A_\lambda \lambda u\kappa \lambda u^2)}{\mathrm{sin}2\beta }.$$
(20)
Recently the lower bounds of these masses become larger and we may use these to put some constraints on $`\lambda `$ and $`u`$. Another important point to use Eq. (7) is that it must be smaller than other two diagonal mass matrix elements of the $`3\times 3`$ neutral Higgs scalar mass matrix. Especially the diagonal mass for the singlet Higgs scalar $`S`$ can give a substantial constraint on $`u`$. Its tree level formula is
$$m_{H_3^0}^2=\frac{1}{2}g_E^2Q_S^2u^2+\frac{A_\lambda \lambda v_1v_2}{u}$$
(21)
in the extra U(1) models, while it is expressed as
$$m_{H_3^0}^2=4\kappa ^2u^2+\frac{A_\lambda \lambda v_1v_2}{u}A_\kappa \kappa u$$
(22)
in the NMSSM. This constraint may be substantial in the NMSSM where there is no other clear constraint on the small $`u`$.
4. Numerical analysis and its results
In this section we numerically estimate the bound of $`m_{h^0}^2(m_{h^0}^{(0)2}+\mathrm{\Delta }m_{h^0}^2)`$ by solving the RGEs and taking account of the phenomenological constraints presented above. In order to improve the one-loop effective potential We use two-loop RGEs for dimensionless coupling constants and one-loop ones for dimensional SUSY breaking parameters, for simplicity. In this estimation we adopt the following procedure. As the initial conditions for the SUSY breaking parameters we take
$$\stackrel{~}{m}_{\varphi _i}^2=(\gamma _i\stackrel{~}{m})^2,M_a=M,A_t=A_k=A_\kappa =A_\lambda =A,$$
(23)
where $`\stackrel{~}{m}^2`$ is the universal soft scalar mass and we introduce the nonuniversality represented by $`\gamma _i`$ only among soft scalar masses of $`H_1,H_2`$ and $`S`$. We comment on this point later. These initial conditions are assumed to be applied at the scale where the coupling unification of SU(2)<sub>L</sub> and U(1)<sub>Y</sub> occurs. We donot require the regolous coupling unification of SU(3)<sub>C</sub> but only impose the realization of the low energy experimental value following Ref. . For the extra U(1) coupling $`g_E`$ we use the same initial value as the one of U(1)<sub>Y</sub> at the unification scale $`M_X`$. The initial values of these parameters are surveyed through the following region,
$`0h_t1.2(0.1),0k,|\kappa |2.0(0.2),0\lambda 3.0(0.2),`$
$`0M/M_S0.8(0.2),0\stackrel{~}{m}/M_S,|A|/M_S3.0(0.3),`$ (24)
where in the parentheses we give the interval which we use in the survey of these parameter regions. Since the sign of $`\kappa `$ and $`A`$ affect the scalar potential, we need to investigate both sign of them. We also assume that the RGEs of the model are changed from the ones of the supersymmetric extra U(1) models to the nonsupersymmetric ones at a supersymmetry breaking scale $`M_S`$ for which we take $`M_S=1`$ TeV as a typical numerical value <sup>7</sup><sup>7</sup>7In principle we should solve the RGEs of soft supersymmetry breaking parameters under the initial values given in Eq. (15) in order to estimate this scale $`M_S`$. However, we donot take such a way here, for simplicity. It is beyond the present scope to study the dependence of our results on the supersymmetry breaking scale $`M_S`$..
As a criterion for the choice of the correct vacuum, we impose that the radiative symmetry breaking occurs correctly. We check whether the potential minimum satisfying the conditions such as Eq. (6) improved by the one-loop effective potential can satisfy the phenomenologically required conditions such as $`v=174`$ GeV and $`m_t=174`$ GeV starting from the above mentioned initial conditions. It is not so easy to find this solution under the completely universal soft breaking parameters so that in our RGEs analysis we allow the nonuniversality in the region $`0.8\gamma _i1.2`$ among soft supersymmetry breaking masses of Higgs scalars. The nonuniversality of soft scalar masses are generally expected in the superstring models . This treatment seems to be good enough for our purpose such as to estimate the upper mass bound of Higgs scalar. We also additionally impose the following phenomenological conditions.
(i) $`m_{h^0}^2`$ should be smaller than other diagonal components of the Higgs mass matrix (see also footnote 3 and the discussion related to Eqs. (13) and (14)).
(ii) the experimental mass bounds on the charged Higgs bosons, charginos, stops, gluinos and $`Z^{}`$ should be satisfied. Here we require the following values:
$`m_{H^\pm }67\mathrm{GeV},m_{\chi ^\pm }72\mathrm{GeV},\stackrel{~}{m}_{t_{1,2}}>67\mathrm{GeV},`$
$`M_3173\mathrm{GeV},m_Z^{}500\mathrm{GeV}.`$ (25)
(iii) the vacuum should be a color conserving one .
We adopt only the parameters set satisfying these criterions as the candidates of the correct vacua and calculate the Higgs mass bound $`m_{h^0}^2`$ for them.
At first in order to see the difference in the allowed vacuum between the NMSSM and the extra U(1) models we plot the radiative symmetry breaking solutions for the present parameter settings in the $`(\mathrm{tan}\beta ,u)`$ plane in Fig. 1. Solutions are classified by the initial value of $`h_t`$ at $`M_X`$ into three classes which show rather different qualitative features. As an example of the extra U(1) models we take the $`\xi _{}`$ model here but the $`\eta `$ model has been checked to show the similar feature to the $`\xi _{}`$ model. We take the case (A) as the extra matter contents. Throught the present calculation an effect of the translation of the running mass to the pole mass is taken into account to determine $`\mathrm{tan}\beta `$. We take $`\mathrm{tan}\beta 15`$ and neglect the large $`\mathrm{tan}\beta `$ solutions since the bottom Yukawa coupling is assumed to be small in the RGEs so that in the present analysis the large $`\mathrm{tan}\beta `$ solutions cannot be recognized as the appropriate ones. Figure 1 shows that the $`\xi _{}`$ model can have solutions in the larger $`u`$ region of the $`(\mathrm{tan}\beta ,u)`$ plane compared with the NMSSM. As mentioned in the previous section, this is a result that $`k`$ can be larger than $`\kappa `$ at the $`m_t`$ scale due to the SU(3) effect. This is shown in Fig.2, where the values of $`k(m_t)`$, $`\kappa (m_t)`$ and $`\lambda (m_t)`$ corresponding to each solution are plotted for $`\mathrm{tan}\beta `$. The soft scalar mass $`m_S^2`$ of the singlet Higgs scalar $`S`$ becomes much more negative in the extra U(1) models than in the NMSSM. In the sufficiently large $`u`$ region the potential minimum condition for $`u`$ reduces to
$$u^2=\frac{4m_S^2}{g_E^2Q_S^2}\mathrm{for}\mathrm{extra}\mathrm{U}(1),u^2=\frac{m_S^2}{2\kappa ^2}\mathrm{for}\mathrm{NMSSM}.$$
(26)
In the NMSSM $`u`$ depends not only on $`m_S^2`$ but also on $`\kappa `$ and as a result $`u`$ can take a rather large value. In the $`\xi _{}`$ model the smaller $`u`$ region such as $`u_{}^<`$ 1 TeV is cut due to the experimental extra Z mass bound. Also in the NMSSM very small $`u`$ seems to be forbidden. This seems to be a result of the phenomenological conditions (i) and (ii).
The big qualitative difference of the vacuum in both models is that there can be large $`u`$ solutions for $`\mathrm{tan}\beta {}_{}{}^{>}5`$ corresponding to $`h_t(M_X)=0.3`$ in the extra U(1) model. One reason of this is that the smaller $`\lambda (m_t)`$ is realized in the extra U(1) models than in the NMSSM. This is clearly shown in Fig. 2(b). The discussion on this aspect has already given based on the RGE in the previous section. On this point we should also note that in the $`\mathrm{tan}\beta {}_{}{}^{>}5`$ region the small $`\lambda (m_t)`$ is allowed. Thus $`\mu =\lambda u`$ can be in the suitable range even if $`u`$ is large. However, the boundary value of $`u`$ seems not to have so strong dependence on $`\lambda (m_t)`$ in both models and the value of $`\lambda u`$ does not seem to be strictly restricted by the radiative symmetry breaking at least within the parameter region searched in this paper.
In Figs. 3$``$5 we give the results of our numerical estimations of $`m_h^0`$ for each model. In these figures we plot the boundary values of $`m_h^0`$ for the parameters obtained as the solutions of our radiative symmetry breaking study. In each figure (a) the upper and lower boundaries of $`m_h^0`$ are drawn by using the all solutions obtained under the intial values shown in (24). In oder to show the $`h_t(M_X)`$ dependence of $`m_h^0`$ we classify the solutions into three classes and draw them separately. In figures (b) we plot the upper and lower boundaries of $`m_h^0`$ for the remaining solutions after imposing the additional condition $`\stackrel{~}{m}=1`$ TeV. We also add the scatter plots of the all solutions corresponding to 2.4 TeV $`<u<`$ 2.6 TeV in the same figures. They are represented by three kinds of triangles corresponding to each $`h_t(M_X)`$. As a common feature in all models, we find that the larger $`h_t(M_X)`$ realizes the smaller $`\mathrm{tan}\beta `$ and then brings the larger contribution of the second term of Eq. (7). Thus the largest $`\lambda (m_t)`$ in the small $`\mathrm{tan}\beta `$ in Fig. 2(b) gives the largest $`m_h^0`$. Although $`\lambda (m_t)`$ in the extra U(1) models can be smaller than the one of the NMSSM as shown in Fig. 2(b), the boundary values of $`m_h^0`$ is larger in the extra U(1) models than in the NMSSM by a few to ten GeV. This is mainly due to the extra contribution to Eq. (17) coming from the singlet colored fields $`(g_i,\overline{g}_i)`$. Since the existence of this contribution is the basic feature of the present extra U(1) models, the boundary value of $`m_h^0`$ is generally expected to be larger than the one of the NMSSM inspite of the running feature of the Yukawa coupling $`\lambda `$. This one-loop effect is large enough to cancel the difference of $`\lambda (m_t)`$ in the second term of Eq. (7). In our studying parameters space the largest value of $`m_h^0`$ is
$$m_h^0{}_{}{}^{<}156\mathrm{GeV}(\mathrm{NMSSM}),m_h^0{}_{}{}^{<}164\mathrm{GeV}(\mathrm{NMSSM}),m_h^0{}_{}{}^{<}158\mathrm{GeV}(\mathrm{NMSSM}).$$
(27)
By comparing (a) and (b) in Figs. 3$``$5 we can get the tendency how the solutions are restricted when we reduce the parameter space. The change of $`\stackrel{~}{m}`$ and $`u`$ mainly affect the one loop contribution through the mass matrices (14) and (15).
In Fig. 6 we plot the boundary value of $`m_h^0`$ for $`u`$ in the $`\xi _{}`$ model. This shows the tendency that the larger $`u`$ gives the larger value of $`m_h^0`$. This is expected from the one-loop contribution of the extra singlet colored fields $`(g_i,\overline{g}_i)`$. From this figure we can read off the relation between $`m_Z^{}`$ and $`m_h^0`$ by using $`m_Z^{}^2g_E^2Q_S^2u^2/2`$. The lower bound of $`m_Z^{}`$ in Fig. 6 is about 600 GeV where we used $`g_E(m_t)=0.36`$. The conditions (i) and (ii) also determine the lower bound of $`u`$ in the extra U(1) models.
Finally we give a few comments on some points related to the extra matters. We also studied the case (B) of the extra matter contents for the same parameter settings as the above study. In that case, as a common feature we can find, it becomes rather difficult to satisfy both of the radiative symmetry breaking conditions and the phenomenological conditions (i) to (iii) compared with the case (A). The number of solutions in the case (B) is drastically less than in the case (B). Since the value of $`g_3(M_X)`$ increases, $`h_t(m_t)`$ and $`k(m_t)`$ becomes larger. In fact, the initial value of $`h_t`$ in the wide region such as $`0.2h_t(m_t)0.9`$ results in only the small $`\mathrm{tan}\beta `$ (larger $`h_t(m_t)`$) solution such as $`\mathrm{tan}\beta {}_{}{}^{<}1.8`$. This also makes $`\lambda (m_t)`$ smaller. The larger $`\mathrm{tan}\beta `$ solutions disappear and the value of $`|u|`$ is shifted upward. However, the upper boundary value of $`m_h^0`$ behaves in the different way between the NMSSM and the extra U(1) models. Although in both models $`m_{h^0}`$ becomes smaller in the region of $`\mathrm{tan}\beta >_{}2`$, the behavior is different at $`\mathrm{tan}\beta <_{}2`$. In the NMSSM it is a little bit larger than the one of case (A). On the other hand, it becomes smaller than the one of case (A) by a several GeV in the extra U(1) models. Here we should remind the fact that even if $`\lambda (m_t)`$ is smaller $`m_{h^0}`$ can be larger in the case that corresponding $`tan\beta `$ is smaller. The difference in the RGE of $`\lambda `$ in both model is also important in this behavior. To have more confident quantitative results in this case we need to search the parameter space in the finer way. We also changed the number of $`(g_i,\overline{g}_i)`$ which couples to $`S`$ in the superpotential (1) in the case (A). If we decrease this number from three to one, the boundary values of the allowed $`m_h^0`$ become larger. This reason is considered as follows. Although this decrease reduces the number of fields contributing to the one-loop effective potential, this also decreases the $`N_g`$ value in Eq. (10). As a result the larger $`k`$ and $`\lambda `$ are realized at the low energy region. The larger $`k`$ also brings the larger $`u`$. The contribution to the one-loop effect per a field can be larger. Thus the decrease of the number of $`(g_i,\overline{g}_i)`$ which couples to $`S`$ causes the increase of $`m_h^0`$ at not only the tree level but also the one-loop level.
5. Summary
There are two well-known low energy candidates to solve the $`\mu `$ problem in the MSSM. These are the NMSSM and the extra U(1) models. We have estimated the upper bound of the lightest neutral Higgs mass in both models. Apart from a Higgs coupling $`\lambda SH_1H_2`$, there is a typical coupling $`\kappa S^3`$ in the NMSSM and $`kSg\overline{g}`$ in the extra U(1) models. In the NMSSM $`\kappa `$ plays a crucial role in the evolution of $`\lambda `$ which dominantly determines the tree level mass bound of the lightest neutral Higgs scalar and in the radiative symmetry breaking. In the extra U(1) models the introduction of the extra colored fields $`g,\overline{g}`$ and its coupling with the singlet Higgs $`S`$ are crucial to cause the radiative symmetry breaking at the weak scale successfully. This coupling can also affect the running of the coupling constant $`\lambda `$. We focussed our attention on these points and estimated the the upper bound of the lightest neutral Higgs mass in both models. In this estimation we additionally imposed some phenomenological constraints related to $`\lambda `$ and the VEV of $`S`$ coming from, for example, the mass bounds of the charginos, the charged Higgs scalars and the $`Z^{}`$ boson. We solved the minimum conditions of the one-loop effective potential improved by the RGEs for the couplings and soft supersymmetry breaking parameters whose initial conditions are taken in the suitable region. We estimated the upper bound of the lightest neutral Higgs scalar for the parameters which bring the phenomenologically correct potential minimum. Its tree level contribution due to $`\lambda `$ can be smaller in the extra U(1) models than in the NMSSM. However, there is the extra one-loop contribution originated from the Yukawa coupling $`kSg\overline{g}`$ and this makes its upper bound larger in the extra U(1) than in the NMSSM by a few to ten GeV. It is interesting enough that the upper bound of the lightest neutral Higgs scalar in the extra U(1) models is not so different from the one of the NMSSM. The extra U(1) models may be an equal candidate to the NMSSM for the experimental Higgs search.
This work has been partly supported by the a Grant-in-Aid for Scientific Research from the Ministry of Education, Science and Culture(#11640267 and #11127206).
|
warning/0003/cond-mat0003422.html
|
ar5iv
|
text
|
# Quarter-filled spin density wave states with long-range Coulomb interaction**footnote *submitted to J. Phys. Chem. Solids
## 1 Introduction
Bechgaard salts ( (TMTSF)<sub>2</sub>X and (TMTTF)<sub>2</sub>X ), which are known as low dimensional organic conductors, exhibit spin density wave (SDW) states at low temperatures. $`^{\text{?}\text{}\text{?}\text{)}}`$ The SDW states exhibit several unconventional properties associated with charge density wave (CDW). The recent X-ray experiment on (TMTSF)<sub>2</sub>PF<sub>6</sub> salt $`^{\text{?}\text{}\text{?}\text{)}}`$ has shown that 2$`k_\mathrm{F}`$-SDW coexists with 2$`k_\mathrm{F}`$-CDW at temperatures just below the onset temperature of SDW where $`k_\mathrm{F}`$ denotes a Fermi momentum. The coexistence of SDW and CDW has been studied theoretically in terms of the mean-field theory at the absolute zero temperature. By taking into account repulsive interactions of both on-site and nearest-neighbor sites and dimerization, it has been demonstrated that 2$`k_\mathrm{F}`$-SDW coexists with 4$`k_\mathrm{F}`$-CDW.$`^{\text{?}\text{)}}`$ Further the coexistence of 2$`k_\mathrm{F}`$-SDW and 2$`k_\mathrm{F}`$-CDW has been found by adding the next-nearest-neighbor repulsive interaction. $`^{\text{?}\text{)}}`$
In the present paper, by extending the previous calculations, $`^{\text{?, ?, }\text{?}\text{)}}`$ we study if such a long range Coulomb interaction results in the coexistence even at the onset temperature of the SDW state, as found in the experiment. $`^{\text{?, ?)}}`$
## 2 Formulation
We examine a one-dimensional extended Hubbard model with interactions of on-site ($`U`$), nearest-neighbor ($`V`$), next-nearest-neighbor ($`V_2`$) sites and dimerization energy ($`t_\mathrm{d}`$). The Hamiltonian is expressed as $`^{\text{?)}}`$
$`H`$ $`=`$ $`{\displaystyle \underset{\sigma =,}{}}{\displaystyle \underset{j=1}{\overset{N}{}}}(t(1)^jt_\mathrm{d})(C_{j\sigma }^{}C_{j+1,\sigma }+h.c.)`$ (1)
$`+U{\displaystyle \underset{j=1}{\overset{N}{}}}n_jn_j+{\displaystyle \underset{j=1}{\overset{N}{}}}V_{1j}n_jn_{j+1}+V_2{\displaystyle \underset{j=1}{\overset{N}{}}}n_jn_{j+2}.`$
The quantity $`C_{j\sigma }^{}`$ denotes the creation operator of the electron at the $`j`$-th site and $`V_{1j}=V(1)^j\delta V`$ where $`\delta V`$ also comes from dimerization and $`n_j=n_j+n_j`$ with $`n_{j\sigma }=C_{j\sigma }^{}C_{j\sigma }`$. Quantities $`t`$, $`k_\mathrm{B}`$ and lattice constant are taken as unity. Due to the quarter-filled band with the Fermi wave vector $`k_\mathrm{F}=\pi /4`$, order parameters with $`m=1,2,3`$ are calculated self-consistently by
$`S_{mQ_0}`$ $`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{\sigma =,}{}}{\displaystyle \underset{\pi <k\pi }{}}\mathrm{sgn}(\sigma )C_{k\sigma }^{}C_{k+mQ_0,\sigma }_{\mathrm{MF}},`$
$`D_{mQ_0}`$ $`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{\sigma =,}{}}{\displaystyle \underset{\pi <k\pi }{}}C_{k\sigma }^{}C_{k+mQ_0,\sigma }_{\mathrm{MF}},`$ (3)
where $`Q_0=2k_\mathrm{F}`$, $`S_0=0`$, $`D_0=1/2`$, $`S_{Q_0}=S_{3Q_0}^{}`$, $`D_{Q_0}=D_{3Q_0}^{}`$, $`S_{2Q_0}=S_{2Q_0}^{}S_2`$ and $`D_{2Q_0}=D_{2Q_0}^{}D_2`$. In Eqs. (2) and (3), $`S_1(|S_{Q_0}|)`$, $`S_2`$, $`D_1(|D_{Q_0}|)`$ and $`D_2`$ correspond to the amplitudes for 2$`k_\mathrm{F}`$-SDW, 4$`k_\mathrm{F}`$-SDW, 2$`k_\mathrm{F}`$-CDW and 4$`k_\mathrm{F}`$-CDW respectively. From eqs. (2) and (3), the free energy per site with the quantity $`U/16+V/4+V_2/4`$ subtracted is given by
$`F_{\mathrm{MF}}`$ $`=`$ $`{\displaystyle \frac{T}{N}}{\displaystyle \underset{\sigma }{}}{\displaystyle \underset{0<kQ_0}{}}{\displaystyle \underset{n=1}{\overset{4}{}}}\mathrm{ln}\left(1+\mathrm{exp}[(E_{n\sigma }(k)\mu )/T]\right)`$ (4)
$`+U\left[{\displaystyle \frac{1}{8}}{\displaystyle \frac{1}{2}}\left(|D_{Q_0}|^2|S_{Q_0}|^2\right){\displaystyle \frac{1}{4}}\left(D_{2Q_0}^2S_{2Q_0}^2\right)\right]`$
$`+V\left({\displaystyle \frac{1}{2}}+D_{2Q_0}^2\right)+\mathrm{i}\delta V\left(D_{Q_0}^2D_{Q_0}^2\right)`$
$`+V_2\left({\displaystyle \frac{1}{2}}+2|D_{Q_0}|^2D_{2Q_0}^2\right)+{\displaystyle \frac{\mu }{2}},`$
where $`T`$ is temperature and $`\mu `$ is a chemical potential determined by $`D_0=1/2`$. In eq. (4), $`E_{n\sigma }`$ is the eigen value for the mean-filed Hamiltonian, $`^{\text{?)}}`$
$`H_{\mathrm{MF}}`$ $`=`$ $`{\displaystyle \underset{\sigma =,}{}}{\displaystyle \underset{\pi <k\pi }{}}[(\epsilon _k+{\displaystyle \frac{U}{4}}+V+V_2)C_{k\sigma }^{}C_{k\sigma }`$ (5)
$`+(\mathrm{\Delta }_{Q_0\sigma }C_{k+Q_0,\sigma }^{}C_{k\sigma }+h.c.)+\mathrm{\Delta }_{2Q_0\sigma }C_{k\sigma }^{}C_{k+2Q_0,\sigma }],`$
where $`\epsilon _k=2t\mathrm{cos}k`$, $`\mathrm{\Delta }_{Q_0\sigma }=\left(U/22V_2\right)D_{Q_0}+2\mathrm{i}\delta VD_{Q_0}^{}\mathrm{sgn}(\sigma )US_{Q_0}/2`$ and $`\mathrm{\Delta }_{2Q_0\sigma }=\left(U/22V+2V_2\right)D_{2Q_0}\mathrm{sgn}(\sigma )US_{2Q_0}/22\mathrm{i}t_\mathrm{d}\mathrm{sin}k`$. Here we note an excess free energy, $`\delta F`$, which is obtained by expanding $`F_{\mathrm{MF}}`$ in terms of order parameters. Since there are two kinds of coupling for $`D_2S_1^2`$ and $`S_1S_2D_1`$ as found in the next section, the relevant expressions for $`\delta F`$ are written as
$`\delta F`$ $`=`$ $`B_{S_1}S_1^2+B_{D_2}D_2^2+G_1D_2S_1^2+\mathrm{},`$ (6)
$`\delta F`$ $`=`$ $`B_{S_1}S_1^2+B_{S_2}S_2^2+B_{D_1}D_1^2+G_2S_1S_2D_1+\mathrm{},`$
where $`G_1`$ and $`G_2`$ are coupling constants. For small $`B_X`$, one obtains $`B_X(TT_X^0)`$ with $`T_X^0`$ ($`X=S_1,D_2`$ and $`D_1`$) being the onset temperature for $`S_1`$-state, $`D_2`$-state and $`D_1`$-state.
## 3 Phase diagram at finite temperatures
We examine SDW states at finite temperatures, by calculating eqs. (2), (3) and (4). The numerical calculation is performed by using $`t_b=tt_\mathrm{d}`$ and $`V_a=V+\delta V`$ with the fixed $`t_\mathrm{d}/t`$= 1/21, $`U/t=3.81`$ and $`\delta V/V`$ = 1/9. $`^{\text{?)}}`$
In Fig. 1, $`T`$-dependence of order parameters for $`V_2=0`$ and $`V_a/t_b`$ =0, 1.6 and 3 is shown where solid curve and dashed curve correspond to $`S_1`$ and $`D_2`$, respectively and $`S_2=D_1=0`$. The dotted curve with $`V_a`$ =0 denotes a conventional 2$`k_\mathrm{F}`$-SDW state. There are two kinds of phase transitions for $`V_a/t_b`$ = 1.6 (3) where pure $`S_1`$ state ($`D_2`$ state) is obtained at high temperatures with $`0.314<T/t_b<0.341`$ ($`0.688<T/t_b<0.924`$) while a coexistent state of $`S_1`$ and $`D_2`$ is obtained at low temperature with $`T/t_b<0.314`$ ($`T/t_b<0.688`$). The inset denotes a first order transition into a coexistent state of $`S_1`$ and $`D_2`$, which occurs at $`T/t_b`$=0.488 for $`V_a/t_b`$ = 2.5. The first order transition temperature is estimated by comparing $`F_{\mathrm{MF}}`$ of eq.(4).
Based on these results, phase diagram on the plane of $`V_a/t_b`$ and $`T/t_b`$ for $`V_2=0`$ is shown in Fig. 2, where $`S_10`$ ($`D_20`$) at temperatures below the solid curve (dashed curve). With decreasing temperature, the second order transition into $`S_1`$ state occurs for $`V_a/t_b<0.75`$ while the second order transition into $`S_1`$ state is followed by the second (first) order transition into the coexistent state of $`S_1`$ and $`D_2`$ for $`0.75<V_a/t_b<1.67`$ ($`1.67<V_a/t_b<1.89`$). It is noticeable that the first order transition from the normal ($`N`$) state into the coexistent state of $`S_1`$ and $`D_2`$ takes place on the dash-dotted curve with $`1.89<V_a/t_b<2.78`$. Further the second order transition into $`D_2`$ state is followed by the first (second) order transition into the coexistent state of $`S_1`$ and $`D_2`$ for $`2.78<V_a/t_b<3.07`$ ($`3.07<V_a/t_b`$ ). Equation (6) indicates a fact that the first order transition originates in the third term with a coefficient $`G_1`$ and that the second order transition from $`S_1`$ state to the coexistent state of $`S_1`$ and $`D_2`$ is due to $`G_1=0`$ in the pure $`S_1`$ state. Thus the first order transition for $`V_2=0`$ is attributable to the third term of eq. (6).
Next we examine another case of $`V_2/V_a=1`$. In Fig. 3, $`T`$-dependence of order parameters is shown with some choices of $`V_a/t_b`$ where the solid curve, dashed curve and dash-dotted curve correspond to $`S_1`$, $`D_1`$ and $`S_2`$, respectively. For $`V_a/t_b`$ = 1.4 (2.5), pure $`S_1`$ state ($`D_1`$ state) is obtained at high temperatures with $`0.223<T/t_b<0.341`$ ($`0.746<T/t_b<1.00`$) while a coexistent state of $`S_1`$, $`D_1`$ and $`S_2`$ is obtained at low temperature with $`T/t_b<0.223`$ ($`T/t_b<0.746`$). The inset denotes a first order transition, which occurs at $`T/t_b`$=0.4 for $`V_a/t_b`$ = 1.8.
For $`V_2/V_a=1.0`$, the phase diagram on the plane of $`V_a/t_b`$ and $`T/t_b`$ is shown in Fig. 4, where $`S_10`$ ($`D_10`$) below the solid curve (dashed curve). For $`V_a/t_b<1.26`$, the second order transition into $`S_1`$ state occurs while the second order transition into $`S_1`$ state is followed by the second (first) order transition into the coexistent state of $`S_1`$, $`D_1`$ and $`S_2`$ for $`1.26<V_a/t_b<1.43`$ ($`1.43<V_a/t_b<1.66`$). A salient feature is the first order transition on the dash-dotted curve, where the normal ($`N`$) state moves into the coexistent state of $`S_1`$, $`D_1`$ and $`S_2`$ in the interval region of $`1.66<V_a/t_b<2.08`$. The second order transition into $`D_1`$ state is followed by the first (second) order transition into the coexistent state of $`S_1`$, $`D_1`$ and $`S_2`$ for $`2.08<V_a/t_b<2.32`$ ($`2.32<V_a/t_b`$ ). The first order transition originates in the third term with a coefficient $`G_2`$ of eq. (2). It should be noticed that the existence of $`S_2`$ is crucial to obtain a direct transition from $`N`$ state to the coexistent of $`S_1`$ and $`D_1`$. Actually, for $`V_a/t_b=2.5`$ and $`V_2/V_a=0.7`$, the mean-field calculation with $`S_2=0`$ leads to a second order transition from $`N`$ state to pure $`S_1`$ state, $`^{\text{?}\text{)}}`$ while the corresponding calculation in the presence of $`S_2`$ exhibits a first order transition into a coexistent state.
Base on these results, a phase diagram on a plane of $`V_2/V_a`$ and $`V_a/t_b`$ is shown in Fig. 5. With decreasing temperatures, we obtain the following phase transitions. In the region I, a transition from $`N`$ state into pure $`S_1`$ state appears while the region II<sub>a</sub> (II<sub>c</sub>) shows the successive transition given by $`N`$ $``$ $`D_2`$ $``$ $`S_1`$ and $`D_2`$ ( $`N`$ $``$ $`S_1`$ $``$ $`S_1`$ and $`D_2`$ ). The first order transition from $`N`$ state into a coexistent state of $`S_1`$ and $`D_2`$ is obtained in the region II<sub>b</sub>. In the region III<sub>a</sub> (III<sub>c</sub>), there is the successive transition given by $`N`$ $``$ $`D_1`$ $``$ $`S_1`$, $`D_1`$ and $`S_2`$ ( $`N`$ $``$ $`S_1`$ $``$ $`S_1`$, $`D_1`$ and $`S_2`$). The first order transition from $`N`$ state into a coexistent state of $`S_1`$, $`D_1`$ and $`S_2`$ is obtained in the region III<sub>b</sub> while the region III<sub>d</sub> corresponds to the transition given by $`NS_1S_1,D_2S_1,D_1`$ and $`S_2`$. The region III<sub>d</sub> is understood by an example shown in the inset of Fig. 4 on the plane of $`V_2/V_a`$ and $`T/t_b`$, where the corresponding transition is obtained for $`0.5<V_2/V_a<0.54`$ for $`V_a/t_b=2.5`$.
By studying a model at quarter-filling with long range Coulomb interactions, we have found that, with decreasing temperature, a first order transition from $`N`$ state into the coexistent state of SDW and CDW occurs. Since there is a reasonable range of parameters, it is considered that the present result of the first order transition could be relevant to the coexistent state of $`2k_\mathrm{F}`$-SDW and 2$`k_\mathrm{F}`$-CDW found in the X-ray experiment. $`^{\text{?, ?)}}`$
Acknowledgment
The authors thank M. Ogata and K. Yonemitsu for useful discussions. This work was partially supported by a Grant-in-Aid for Scientific Research from the Ministry of Education, Science, Sports and Culture (Grant No.09640429), Japan.
|
warning/0003/math0003235.html
|
ar5iv
|
text
|
# Chapter 0 Navier-Stokes Equations and Fluid Turbulence
## Chapter 0 Navier-Stokes Equations and Fluid Turbulence
Incompressible fluids are described by the Navier-Stokes equation. Turbulence (, , ) experiments provide measurements that correspond to averages of certain quantities associated to the variables appearing in the Navier-Stokes equation. The present mathematical knowledge about the Navier-Stokes equations is incomplete. Some of the quantities measured in experiments are accessible to mathematical theory. They are usually low order, one-point bulk averages like the time average of integrals of squares of gradients. Most other measured quantities are not amenable to rigorous quantitative a priori analysis. Turbulence is concerned with statistical or collective properties of fluids. Nevertheless, the main impediment to progress in the rigorous analysis of turbulence is the present lack of understanding of possible blow up in individual solutions of the Euler and Navier-Stokes systems.
I will discuss briefly the blow up problem and present an Eulerian-Lagrangian approach to fluids. I will also give examples of low order one-point bulk quantities that can be treated with present knowledge and discuss results on certain two-point quantities.
### 1 An Eulerian-Lagrangian Approach to Fluids
I will start by recalling that the Navier-Stokes-Euler system can be written as an evolution equation for the three-component velocity vector $`u=u(x,t)`$,
$$\frac{u}{t}+uu+p=\nu \mathrm{\Delta }u+f;$$
the pressure $`p=p(x,t)`$ preserves incompressibility
$$u=0.$$
The Euler system is obtained if the kinematic viscosity vanishes, $`\nu =0`$; the Navier-Stokes system if $`\nu >0`$. Boundary conditions are different for the two systems. The blow up question can be stated in its simplest form for the pure initial value problem: are there any smooth initial data with finite energy that lead to solutions that diverge in finite time? The answer is not known.
For a blow up in the Navier-Stokes system one would have to have a finite time divergence of an eddy-viscosity-like quantity:
$$_0^T\underset{x,r}{sup}|u(x+r,t)u(x,t)|^2dt<\mathrm{}u𝐂^{\mathrm{}}.$$
By contrast, it is known () that
$$_0^T\underset{x,r}{sup}|u(x+r,t)u(x,t)|dt<\mathrm{}.$$
In a situation in which all velocities are finite no singularities can appear in the Navier-Stokes equations. One could accept the finiteness of velocities as a physical hypothesis. For the Euler system this hypothesis would not be sufficient to ensure smoothness of solutions (). A sufficient condition for regularity in the Euler equations is the finiteness of the time integral of the maximum magnitude of vorticity (). The vorticity (the curl of velocity or anti-symmetric part of the velocity gradient), $`\omega =\times u`$ obeys a quadratic evolution equation. The magnitude of vorticity obeys
$$D_t|\omega |=\alpha |\omega |$$
with $`D_t=_t+u`$ the material derivative along particle paths. The scalar stretching term $`\alpha `$ is related to the magnitude of vorticity by a principal value singular integral ():
$$\alpha (x)=P.V.D(\widehat{y},\xi (x),\xi (x+y))|\omega (x+y)|\frac{dy}{|y|^3}$$
$$\xi (x)=\frac{\omega (x)}{|\omega (x)|}.$$
The smooth, mean zero function of three unit vectors $`D`$ vanishes when two of its arguments are on the same line. Consequently, if the direction of vorticity $`\xi `$ is Lipschitz then the singular integral representing $`\alpha `$ is mild and the solutions remain smooth (, , ). This is a generalization of the two dimensional situation where $`\xi =0`$ and the solutions remain smooth. The possibility of blow up due to strong vortex stretching is not removed by the previous result; the result only precludes blow up in a smooth vortex line field. Blow up can occur also because of strain intensification (the strain matrix is the symmetric part of the gradient of vorticity). I will present now a description of the Euler equations that is convenient for analysis () and allows for a clearer geometric picture of the possible singularity formation. I will present the results only in the periodic case for simplicity of exposition. The case of decay at infinity is almost identical.
###### Theorem 1
A function $`u(x,t)`$ solves the incompressible Euler equations if and only if it can be represented in the form $`u=u_A`$,
$$u_A^i(x,t)=u_0^m(A(x,t))\frac{A^m(x,t)}{x_i}\frac{n_A(x,t)}{x_i}$$
where $`A(x,t)`$ solves the active vector equation
$$\frac{A(x,t)}{t}+u_A(x,t)A(x,t)=0$$
with initial data
$$A(x,0)=x.$$
$`u_0`$ is the initial velocity and $`n_A(x,t)`$ is determined up to additive constants by the requirement of incompressibility, $`u_A=0`$.
Let us denote by
$$𝐏=(\mathrm{𝟏}\mathrm{\Delta }^1)$$
the Leray-Hodge projector on divergence-free vectors. The local existence theorem requires just one derivative to be Hölder continuous:
###### Theorem 2
Let $`u_0`$ be a divergence free $`C^{1,ϵ}`$ periodic vector valued function of three variables. There exists a time interval $`[0,T]`$ and a unique $`C([0,T];C^{1,ϵ})`$ spatially periodic vector valued function $`\delta (x,t)`$ such that
$$A(x,t)=x+\delta (x,t)$$
solves the active vector formulation of the Euler equations,
$$\frac{A}{t}+u_AA=0,$$
$$u_A=𝐏\left\{(A(x,t))^{}u_0(A(x,t))\right\}$$
with initial datum $`A(x,0)=x`$.
The proof of this result is based on an identity that removes the apparent ill-posedness, on singular integral calculus, and on the use of the method of characteristics. In the active vector formulation, the ”back-to-labels” map $`A`$ has conserved distribution. Its time evolution is a smooth, volume-preserving rearrangement. Singularities can occur only if the gradient map $`A`$ diverges rapidly in finite time:
$$_0^T\underset{x}{sup}|A(x,t)|^2dt<\mathrm{}A𝐂^{\mathrm{}}.$$
Thus, would-be singularities are gradient singularities in a conserved quantity, similar to shocks in conservation laws, but with the significant difference that the characteristic flow is measure preserving. The formula relating the velocity to the value of the spatial gradient at the same instance of time (, , , ) may have important mathematical consequences. On one hand, conservation of kinetic energy confers a constraint to the growth of $`A`$. On the other hand, the formula suggests that near regions of high gradient of $`A`$ the velocity is exceedingly high, making the regions of high gradient difficult to track and perhaps unstable. Dynamical stability or instability of blow up modalities is a difficult subject. There are obvious space and time symmetries (for instance, a minute delay of blow up), that clearly should not be categorized as instabilities. Nevertheless, even relatively simple PDEs can exhibit the coexistence of a variety of dynamical behaviors, including several stable blow up modalities, stable time independent solutions, unstable blow up modalities and dynamical connections between the unstable behaviors and the stable ones ().
I will pass now from the blow up problem to some more tractable questions about average properties. One can obtain rigorous upper bounds for certain bulk averages of solutions of Navier-Stokes equations. Lower bounds are harder to obtain. Upper bounds for bulk averages of low order moments for Rayleigh-Bénard convection will be described further below; I start with a lower bound for the burning rate in a simple model of turbulent combustion.
### 2 Bulk Burning Rate
Mixtures of reactants may interact in a burning region that has a rather complicated spatial structure but is thin across (). This reaction region moves towards the unburned reactants leaving behind the burned ones. When the reactants are carried by an ambient fluid then the burning rate is enhanced. The physical reason for the observed speed-up is believed to be that fluid advection tends to increase the area available for reaction. What characteristics of the ambient fluid flow are responsible for burning rate enhancement? The question needs first to be made precise, because the reaction region may be complicated and, in general, may move with an ill-defined velocity. An unambiguous quantity $`V`$ representing the bulk burning rate is defined in () and explicit estimates of $`V`$ in terms of the magnitude of the advecting velocity and the geometry of streamlines are derived. In situations where traveling waves are known to exist, $`V`$ coincides with the traveling wave speed and the estimates thus provide automatically bounds for the speed of the traveling waves. The main result of () is the identification of a class of flows that are particularly effective in speeding up the bulk burning rate. The main feature of these “percolating flows” is the presence of tubes of streamlines connecting distant regions of burned and unburned material. For such flows we obtained an optimal linear enhancement bound $`VKU`$ where $`U`$ represents the magnitude of the advecting velocity and $`K`$ is a proportionality factor that depends on the geometry of streamlines but not the speed of the flow. Other flows and in particular cellular flows, which have closed streamlines, on the other hand, may produce a weaker enhancement. The instantaneous bulk burning rate is defined by the formula
$$V(t)=_D\frac{T}{t}(x,y,t)𝑑x𝑑y$$
where the integral extends over the spatial domain $`D`$, taken here for simplicity of exposition to be a two-dimensional strip of unit width and infinite length
$$0y1,\mathrm{}<x<\mathrm{}.$$
The temperature $`T`$ is assumed to obey Neumann boundary conditions at the finite boundaries and to obey
$$T(\mathrm{},y)=1,T(\mathrm{},y)=0.$$
The simplified model is a passive reactive scalar with a KPP nonlinearity
$$T_t+uT\kappa \mathrm{\Delta }T=\frac{v_0^2}{4\kappa }T(1T).$$
with prescribed velocity $`u`$ that satisfies
$$_0^1u(x,y,t)𝑑y=0,u=0.$$
The constant $`v_0`$ represents the speed of a stable one-dimensional laminar ($`u=0`$) traveling wave. We start with a very general lower bound:
###### Theorem 3
For arbitrary initial data obeying
$$0T_0(x,y)1.$$
one has the general lower bound
$$V(t)Cv_0\left(1e^{\frac{v_0^2t}{2\kappa }}\right).$$
The proof is based on the product lemma
###### Lemma 2.1.
There exists a constant $`C>0`$ such that
$$0T(x,y)1,$$
$$T(\mathrm{},y)=1,T(\mathrm{},y)=0\mathrm{for}\mathrm{any}y[0,1].$$
implies
$$\left(_DT(1T)𝑑x𝑑y\right)\left(_D|T|^2𝑑x𝑑y\right)C.$$
Although information about the velocity is not present in the general result, it nevertheless shows that this model does not permit quenching. Also, the general lower bound applies to the homogenized version of the equations as well.
For a very general class of velocities $`u(x,y,t)`$ it can be shown that the bulk burning rate may not exceed a linear bound in the amplitude of the advecting velocity. For a large class of flows we proved lower bounds on the bulk burning rate that are linear in the magnitude of advection. We denote by
$$V_\tau =\frac{1}{\tau }\underset{0}{\overset{\tau }{}}V(t)𝑑t$$
the time average of the instantaneous bulk burning rate. The main result of () is too technical to state here precisely but its meaning is that presence of coherent tubes of streamlines connecting unburned and burned regions enhances the burning rate
$$V_\tau KU$$
as long as the velocity spatial scales are not too small compared to the reaction length scale $`\frac{\kappa }{v_0}`$, and the time scale of change of the advecting velocity is not too small compared to $`\tau _0=\mathrm{max}[\frac{\kappa }{v_0^2},\frac{\stackrel{~}{H}}{v_0}]`$ where $`\stackrel{~}{H}`$ is associated to the width of the coherent tubes of streamlines. For instance, a result concerning mean zero shear flow of the form
$$u(x,y)=(u(y),0),_0^1u(y)𝑑y=0$$
can be stated as
###### Theorem 4.
Let us consider an arbitrary partition of the interval $`[0,1]`$ into subintervals $`I_j=[c_jh_j,c_j+h_j]`$ on which $`u(y)`$ does not change sign. Denote by $`D_{},`$ $`D_+`$ the unions of intervals $`I_j`$ where $`u(y)>0`$ and $`u(y)<0`$ respectively. Then there exist constants $`C_\pm >0`$, independent of the partition, $`u(y),`$ and the initial data $`T_0(x,y)`$, so that the average burning rate $`V_\tau `$ satisfies the following estimate:
$$V_\tau C_+c_+\underset{I_jD_+}{}\left(1+\frac{l^2}{h_j^2}\right)^1\underset{c_j\frac{h_j}{2}}{\overset{c_j+\frac{h_j}{2}}{}}|u(y)|𝑑y$$
$$+C_{}c_{}\underset{I_jD_{}}{}\left(1+\frac{l^2}{h_j^2}\right)^1\underset{c_j\frac{h_j}{2}}{\overset{c_j+\frac{h_j}{2}}{}}|u(y)|𝑑y$$
for any $`\tau \tau _0=\text{max}[\frac{\kappa }{v_0^2},\frac{H}{v_0}].`$ ($`H=1`$). Here $`l=\kappa /v_0.`$ The constants $`c_\pm `$ are defined by
$$c_\pm =\left(\underset{I_jD_{}}{}\frac{h_j^3}{h_j^2+l^2}\right)\left(\underset{I_j}{}\frac{h_j^3}{h_j^2+l^2}\right)^1.$$
The main result of () applies to a large class of flows that are not necessarily spatially periodic, nor shears, and can have completely arbitrary features outside the tubes of streamlines. The bulk burning rate is still linear in the magnitude of the advecting velocity, no matter what kind of behavior (closed streamlines, areas of still fluid, etc.) the flow has outside the tubes. The proportionality coefficient depends on the geometry of the flow in a rather complex manner. These bounds can be extended to larger classes of chemistries. The lower bounds, however, are not yet available for models in which there is a feedback coupling of temperature on the velocity. For such models upper bounds can be derived. In the next section we will discuss upper bounds for heat transfer. We will concentrate on the simplest coupling, mediated by gravity in a Boussinesq approximation and discuss the canonical case of Rayleigh-Bénard convection.
### 3 Bulk Heat Transfer
The equations for Rayleigh-Bénard convection in the Boussinesq approximation are
$$\frac{u}{t}+uu+p=\sigma \mathrm{\Delta }u+\sigma Ra\widehat{e}T,$$
$$u=0$$
$$\frac{T}{t}+uT=\mathrm{\Delta }T.$$
The vertical direction $`\widehat{e}`$ of gravity is singled out. We consider as spatial domain a box of height $`1`$ and lateral side $`L`$. The velocity vanishes at the boundary, $`T=1`$ at the bottom boundary and $`T=0`$ at top. The Nusselt number is the space-time average of the flux of temperature across horizontal cross-section planes. From the equations of motion it follows that
$$|T|^2=N.$$
and also
$$|u|^2=Ra(N1).$$
Here $`<\mathrm{}>`$ is global space-time average. The general rigorous result here is (, )
###### Theorem 5.
There exists an absolute constant $`C`$, independent of Rayleigh number $`Ra`$, aspect ratio $`L`$ and Prandtl number $`\sigma `$ such that
$$N1+C\sqrt{Ra}$$
holds for all solutions of the Boussinesq equations.
The experimental data () point however more to results of the type $`NRa^{\frac{1}{3}}`$ or $`NRa^{\frac{2}{7}}`$. The exponent $`\frac{1}{3}`$ can be obtained rigorously for a simplified model. Consider the infinite Prandtl number equations for rotating convection
$$\left(_t+u\right)T=\mathrm{\Delta }T$$
$$\begin{array}{c}\mathrm{\Delta }uE^1v+p_x=0\hfill \\ \mathrm{\Delta }v+E^1u+p_y=0\hfill \\ \mathrm{\Delta }w+p_z=RT.\hfill \end{array}$$
$$𝐮=0$$
with boundary conditions: ($`(u,v,w)`$, $`p`$, $`T`$) periodic in $`x`$ and $`y`$ with period $`L`$; $`u`$, $`v`$, and $`w`$ vanish for $`z=0,1`$, $`T=0`$ at $`z=1`$, $`T=1`$ at $`z=0`$. One can prove ()
###### Theorem 6.
There exist absolute constants $`c_1,\mathrm{},c_4`$ so that the Nusselt number for rotating infinite Prandtl-number convection is bounded by
$$N1$$
$$\mathrm{min}\{c_1R^{\frac{2}{5}};(c_2E^2+c_3E)R^2;c_4R^{1/3}(E^1+\mathrm{log}_+R)^{\frac{2}{3}}\}.$$
This coincides, in the limit of no rotation $`E\mathrm{}`$ with a logarithmic correction () to the $`\frac{1}{3}`$ exponent. The bound also shows that strong rotation $`E0`$ stabilizes the system and that increasing rotation may result in a non-monotonic behavior of the Nusselt number, as observed in experiments.
Consider now the horizontal average $`\overline{T}(z,t)`$ of $`T(x,y,z,t)`$ and define
$$n=|(T\overline{T})|^2.$$
Note that
$$nN$$
###### Theorem 7.
For the full Boussinesq system
$$N1+c(nRa)^{\frac{1}{3}}$$
holds. For the infinite Prandtl number system
$$N1+C\left(Ra(\mathrm{log}_+Ra)^2\sqrt{n}\right)^{\frac{2}{7}}$$
holds.
Note that, if $`nN`$ is used then the result recovers the exponents $`\frac{1}{2}`$ for general Rayleigh - Bénard and logarithmically corrected $`\frac{1}{3}`$ for the infinite Prandtl number case. But the rigorous appearance of the exponent $`\frac{2}{7}`$ is perhaps not coincidental.
One of the technical ingredients for the proof of these results concerns zero order operators that are not translation invariant
$$B=\frac{^2}{z^2}(\mathrm{\Delta }_{DN}^2)^1\mathrm{\Delta }_h$$
where $`w=(\mathrm{\Delta }_{DN}^2)^1f`$ is the solution of
$$\mathrm{\Delta }^2w=f$$
with horizontally periodic and vertically Dirichlet and Neumann boundary conditions $`w=w^{}=0`$.
###### Theorem 8.
For any $`\alpha (0,1)`$ there exists a positive constant $`C_\alpha `$ such that every Hölder continuous function $`\theta `$ that is horizontally periodic and vanishes at the vertical boundaries satisfies
$$B\theta _L^{\mathrm{}}C_\alpha \theta _L^{\mathrm{}}\left(1+\mathrm{log}_+\theta _{C^{0,\alpha }}\right)^2.$$
The proof of this result is based on a pointwise bound on exponential-oscillatory sums of the type:
$$K(x,z)=\underset{k𝐙^2}{}e^{\frac{2\pi }{L}ikx}m_k^pe^{ϵm_k}$$
where $`ϵ=ϵ(z)0`$, $`m_k=\frac{2\pi }{L}|k|`$ and $`ϵ(z)=0z=z_0`$. The sum is singular at $`z=z_0`$ and the pointwise bound
$$|K(x,z)|C_p\left[|x|^2+ϵ^2(z)\right]^{\frac{p+2}{2}}$$
is obtained using the Poisson summation formula.
The bounds on bulk one-point quantities presented above are among the most successful areas of mathematical and experimental agreement. The reason is perhaps that the quantities involved are numbers, albeit numbers depending on a parameter. The next step beyond the description of bulk one-point averages is the description of power spectra. These are asserted to have some universal features in physical turbulence theory; we present some mathematical results in the next section.
### 4 Spectra
Unlike bulk one-point quantities, spectra are averages of functions. There are some well-established spectra in the physical literature associated to small scale turbulence: the Kraichnan spectrum in two dimensions and the Kolmogorov spectrum in three dimensions.
The energy spectrum $`E(k)`$ is a function that has the property that
$$_0^{\mathrm{}}E(k)𝑑k=|u|^2.$$
(The convention is that $`<\mathrm{}>`$ is normalized space integral followed by long time average). The Kolmogorov spectrum for 3D turbulence is
$$E(k)=C_{Kl}ϵ^{\frac{2}{3}}k^{\frac{5}{3}}.$$
The 2D Kraichnan spectrum is
$$E(k)=C_{Kr}\eta ^{\frac{2}{3}}k^3.$$
Here $`ϵ=\nu |u|^2`$ is the rate of dissipation of energy and $`\eta =\nu |\omega |^2`$ is the rate of dissipation of enstrophy. The spectra are supposed to be valid in a range of scales $`k[k_i,k_d]`$ where $`k_d`$ is the dissipation scale and is determined by viscosity and $`ϵ`$ (respectively viscosity and $`\eta `$) alone. Their expressions are then determined by dimensional analysis. We consider the Littlewood-Paley decomposition of the velocity associated to a mollifier $`\varphi (\xi )`$. This is a smooth function in $`R^d`$ that is non-increasing, smooth, radially symmetric, satisfying $`\varphi (\xi )=1`$ for $`|\xi |\frac{5}{8}`$, $`\varphi (\xi )=0`$ for $`|\xi |\frac{3}{4}`$. One sets $`\psi _{(0)}(\xi )=\varphi (\frac{\xi }{2})\varphi (\xi )`$, $`\psi _{(m)}(\xi )=\psi _{(0)}(2^m\xi )`$ and
$$\varphi (\xi )=_{𝐑^d}e^{i(\xi z)}\mathrm{\Phi }(z)𝑑z,$$
$$\psi _{(m)}(\xi )=_{𝐑^d}e^{i(\xi \dot{z})}\mathrm{\Psi }_{(m)}(z)𝑑z.$$
The Littlewood-Paley decomposition is
$$u(x,t)=u_{(\mathrm{})}(x,t)+\underset{m=0}{\overset{\mathrm{}}{}}u_{(m)}(x,t)$$
where
$$u_{(\mathrm{})}(x,t)=L^d_{𝐑^d}\mathrm{\Phi }\left(\frac{y}{L}\right)u(xy,t)𝑑y,$$
$$u_{(m)}(x,t)=L^d_{𝐑^d}\mathrm{\Psi }_{(m)}\left(\frac{y}{L}\right)u(xy,t)𝑑y$$
and $`L`$ is a length (the integral scale).
We define the Littlewood-Paley spectrum to be
$$E_{LP}(k)=k_m^1|u_{(m)}|^2$$
for $`k_{m1}k<k_m`$, $`m1`$ with $`k_m=2^mL^1`$.
We start with the two dimensional Navier-Stokes equation
$$\left(_t+u\nu \mathrm{\Delta }\right)\omega =f$$
with
$$u(x,t)=\frac{1}{2\pi }_{𝐑^2}\frac{y^{}}{|y|^2}\omega (xy,t)𝑑y$$
One can prove ()
###### Theorem 9.
Assume that the source of vorticity in the Navier-Stokes equations has spectrum localized in the region of wave numbers $`k\frac{1}{L}`$ for some $`L>0`$. Then there exists a constant $`C`$ such that the Littlewood-Paley energy spectrum of solutions of two dimensional forced Navier-Stokes equations obeys the bound
$$E_{LP}(k)Ck^3\left\{\tau ^2\left(\frac{k_d}{k}\right)^6\right\}$$
for $`kL^1`$. Here $`\tau ^1=<u_L^{\mathrm{}}>`$.
The corresponding three-dimensional energy spectrum result requires a significant assumption:
$$|u|^3<\mathrm{}.$$
Denoting
$$\widehat{ϵ}=\nu \left\{|u|^3\right\}^{\frac{2}{3}}$$
$$\widehat{\eta }=\nu ^{\frac{3}{4}}\left(\widehat{ϵ}\right)^{\frac{1}{4}}$$
$$\widehat{k}_d=\nu ^{\frac{3}{4}}(\widehat{ϵ})^{\frac{1}{4}}$$
and setting
$$C_\psi =|\mathrm{\Psi }_{(0)}(a)||a|^2|\mathrm{\Psi }_{(0)}(b)||b|𝑑a𝑑b.$$
we have ():
###### Theorem 10.
Consider three-dimensional body forces that satisfy
$$\widehat{f}(k)=0$$
for all $`|k|\frac{C}{L}`$ and some $`C>0`$. Consider solutions of the three dimensional Navier-Stokes equation that satisfy $`\widehat{ϵ}<\mathrm{}`$. Then
$$E_{LP}(k)C_\psi \left(\widehat{ϵ}\right)^{\frac{2}{3}}k^{\frac{5}{3}}\left(\frac{k}{\widehat{k}_d}\right)^{\frac{10}{3}}$$
holds for $`|k|\frac{C}{L}`$.
Acknowledgment. Research partially supported by NSF DMS-9802611, by AIM, and by the ASCI Flash Center at the University of Chicago under DOE contract B341495.
|
warning/0003/math0003214.html
|
ar5iv
|
text
|
# Geometry of Quaternionic Kähler connections with torsion
## 1 Introduction and statement of the results
An almost hyper complex structure on a 4n-dimensional manifold $`M`$ is a triple $`H=(J_\alpha ),\alpha =1,2,3`$, of almost complex structures $`J_\alpha :TMTM`$ satisfying the quaternionic identities $`J_\alpha ^2=id`$ and $`J_1J_2=J_2J_1=J_3`$. When each $`J_\alpha `$ is a complex structure, $`H`$ is said to be a hyper complex structure on $`M`$.
An almost quaternionic structure on $`M`$ is a rank-3 subbundle $`QEnd(TM)`$ which is locally spanned by almost hypercomplex structure $`H=(J_\alpha )`$; such a locally defined triple $`H`$ will be called an admissible basis of $`Q`$. A linear connection $``$ on $`TM`$ is called quaternionic connection if $``$ preserves $`Q`$, i.e. $`_X\sigma \mathrm{\Gamma }(Q)`$ for all vector fields $`X`$ and smooth sections $`\sigma \mathrm{\Gamma }(Q)`$. An almost quaternionic structure is said to be a quaternionic if there is a torsion-free quaternionic connection. A $`Q`$-hermitian metric is a Riemannian metric which is Hermitian with respect to each almost complex structure in $`Q`$. An almost quaternionic (resp. quaternionic) manifold with Q-hermitian metric is called an almost quaternionic Hermitian (resp. quaternionic hermitian) manifold
For $`n=1`$ an almost quaternionic structure is the same as an oriented conformal structure and it turns out to be always quaternionic. When $`n2`$, the existence of torsion-free quaternionic connection is a strong condition which is equivalent to the 1-integrability of the associated GL(n,H)SP(1) structure . If the Levi-Civita connection of a quaternionic hermitian manifold $`(M,g,Q)`$ is a quaternionic connection then $`(M,g,Q)`$ is called Quaternionic Kähler (briefly QK). This condition is equivalent to the statement that the holonomy group of $`g`$ is contained in SP(n).SP(1) . If on a QK manifold there exist an admissible basis $`(H)`$ such that each almost complex structure $`(J_\alpha )(H),\alpha =1,2,3`$ is parallel with respect to the Levi-Civita connection then the manifold is called hyper Kähler (briefly HK). In this case the holonomy group of $`g`$ is contained in SP(n).
The notions of quaternionic manifolds arise in a natural way from the theory of supersymmetric sigma models. The geometry of the target space of two-dimensional sigma models with extended supersymmetry is described by the properties of a metric connection with torsion . The geometry of (4,0) supersymmetric two-dimensional sigma models without Wess-Zumino term (torsion) is a hyper Kähler manifold. In the presence of torsion the geometry of the target space becomes hyper Kähler with torsion (briefly HKT) . This means that the complex structures $`J_\alpha ,\alpha =1,2,3`$, are parallel with respect to a metric quaternionic connection with totally skew-symmetric torsion . Local (4,0) supersymmetry requires that the target space of two dimensional sigma models with Wess-Zumino term be either HKT or quaternionic Kähler with torsion (briefly QKT) which means that the quaternionic subbundle is parallel with respect to a metric linear connection with totally skew-symmetric torsion and the torsion 3-form is of type (1,2)+(2,1) with respect to all almost complex structures in $`Q`$. The target space of two-dimensional (4,0) supersymmetric sigma models with torsion coupled to (4,0) supergravity is a QKT manifold . If the torsion of a QKT manifold is a closed 3-form then it is called strong QKT manifold. The properties of HKT and QKT geometries strongly resemble those of HK and QK ones, respectively. In particular, HKT and QKT manifolds admit twistor constructions with twistor spaces which have similar properties to those of HK and QK .
The main object of interest in this article is the differential geometric properties of QKT manifolds. We find necessary and sufficient conditions to the existence of a QKT connection in terms of the Kähler 2-forms and show that the QKT-connection is unique if dimension is at least 8 (see Theorem 2.2 below). We prove that the QKT manifolds are invariant under conformal transformations of the metric. This allows us to present a lot of (compact) examples of QKT manifolds. In particular, we show that the compact quaternionic Hopf manifolds studied in , which do not admit a QK structure, are QKT manifolds. In the compact case we show the existence of Gauduchon metric i.e. the unique conformally equivalent QKT structure with co-closed torsion 1-form.
It is shown in that the twistor space of a QKT manifold is always complex manifold provided the dimension is at least 8. It admits complex contact (resp. Kähler) structure if the torsion 4-form is of type (2,2) and some additional nondegeneratity (positivity) conditions are fulfilled . Most of the known examples of QKT manifolds are homogeneous constructed in. However, there are no homogeneous proper QKT manifolds (i.e. QKT which is not QK or HKT) with torsion 4-form of type (2,2) in dimensions greater than four by the result of . We generalise this result showing that there are no proper QKT manifolds with torsion 4-form of type (2,2) provided that the torsion is parallel and dimension is at least 8.
In dimension 4 a lot of examples of QKT manifolds are known . In particular, examples of homogeneous QKT manifolds are constructed in . We notice that there are many (even strong) QKT structures in dimension 4, all depending on an arbitrary 1-form. We give a local description of 4-dimensional QKT manifolds with parallel torsion; namely such a QKT manifold is a Riemannian product of a real line and a 3-dimensional Riemannian manifold. We observe that homogeneous QKT manifolds are precisely naturally reductive homogeneous Riemannian manifolds, the objects which are well known. We present a complete local description (up to an isometry) of 4-dimensional homogeneous QKT which was known in the setting of naturally reductive homogeneous 4-manifold . In the last section we consider 4-dimensional Einstein-like QKT manifold and find a closed relation with Einstein-Weyl geometry in dimension four. In particular, we show that every 4-dimensional HKT manifold is of this type.
Acknowledgements. The research was done during the author’s visit at the Abdus Salam International Centre for Theoretical Physics, Trieste Italy. The author thanks the Abdus Salam ICTP for support and the excellent environment. The author also thanks to G. Papadopoulos for his interest, useful suggestions and remarks. He is grateful to S.Marchiafava for pointing out some incorrect statements and L.Ornea for finding the time to read and comment on a first draft of the manuscript.
## 2 Characterisations of QKT connection
Let $`(M,g,(J_\alpha )Q,\alpha =1,2,3)`$ be a 4n-dimensional almost quaternionic manifold with $`Q`$-hermitian Riemannian metric $`g`$ and an admissible basis $`(J_\alpha )`$. The Kähler form $`F_\alpha `$ of each $`J_\alpha `$ is defined by $`F_\alpha =g(.,J\alpha .)`$. The corresponding Lee forms are given by $`\theta _\alpha =\delta F_\alpha J_\alpha `$.
For an $`r`$-form $`\psi `$ we denote by $`J_\alpha \psi `$ the $`r`$-form defined by
$`J_\alpha \psi (X_1,\mathrm{},X_r):=(1)^r\psi (J_\alpha X_1,\mathrm{},J_\alpha X_r),\alpha =1,2,3`$. Then $`(d^c\psi )_\alpha =(1)^rJ_\alpha d\psi `$. We shall use the notations $`d_\alpha F_\beta :=(d^cF_\beta )_\alpha `$, i.e. $`d_\alpha F_\beta (X,Y,Z)=dF_\beta (J_\alpha X,J_\alpha Y,J_\alpha Z),\alpha ,\beta =1,2,3.`$
We recall the decomposition of a skew-symmetric tensor $`P\mathrm{\Lambda }^2T^{}MTM`$ with respect to a given almost complex structure $`J_\alpha `$. The (1,1), (2,0) and (0,2) part of $`P`$ are defined by $`P^{1,1}(J_\alpha X,J_\alpha Y)=P^{1,1}(X,Y),P^{2,0}(J_\alpha X,Y)=J_\alpha P^{2,0}(X,Y),P^{0,2}(J_\alpha X,Y)=J_\alpha P^{0,2}(X,Y)`$, respectively.
For each $`\alpha =1,2,3`$, we denote by $`dF_\alpha ^+`$ (resp. $`dF_\alpha ^{}`$) the $`(1,2)+(2,1)`$-part (resp. $`(3,0)+(0,3)`$-part) of $`dF_\alpha `$ with respect to the almost complex structure $`J_\alpha `$. We consider the following 1-forms
$$\theta _{\alpha ,\beta }=\frac{1}{2}\underset{i=1}{\overset{4n}{}}dF_\alpha ^+(X,e_i,J_\beta e_i),\alpha ,\beta =1,2,3.$$
Here and further $`e_1,e_2,\mathrm{},4n`$ is an orthonormal basis of the tangential space.
Note that $`\theta _{\alpha ,\alpha }=\theta _\alpha `$.
The Nijenhuis tensor $`N_\alpha `$ of an almost complex structure $`J_\alpha `$ is given by
$`N_\alpha (X,Y)=[J_\alpha X,J_\alpha Y][X,Y]J_\alpha [J_\alpha X,Y]J_\alpha [X,J_\alpha Y].`$
The celebrated Newlander-Nirenberg theorem states that an almost complex structure is a complex structure if and only if its Nijenhuis tensor vanishes.
Let $``$ be a quaternionic connection i.e.
(2.1)
$$J_\alpha =\omega _\beta J_\gamma +\omega _\gamma J_\beta ,$$
where the $`\omega _\alpha ,\alpha =1,2,3`$ are 1-forms.
Here and henceforth $`(\alpha ,\beta ,\gamma )`$ is a cyclic permutation of $`(1,2,3)`$.
Let $`T(X,Y)=_XY_YX[X,Y]`$ be the torsion tensor of type (1,2) of $``$. We denote by the same letter the torsion tensor of type (0,3) given by $`T(X,Y,Z)=g(T(X,Y),Z)`$. The Nijenhuis tensor is expressed in terms of $``$ as follows
$`N_\alpha (X,Y)`$ $`=`$ $`4T_\alpha ^{0,2}(X,Y)`$
$`+`$ $`(_{J_\alpha X}J_\alpha )(Y)(_{J_\alpha Y}J_\alpha )(X)(_YJ_\alpha )(J_\alpha X)+(_XJ_\alpha )(J_\alpha Y),`$
where the (0,2)-part $`T_\alpha ^{0,2}`$ of the torsion with respect to $`J_\alpha `$ is given by
(2.3)
$$T_\alpha ^{0,2}(X,Y)=\frac{1}{4}\left(T(X,Y)T(J_\alpha X,J_\alpha Y)+J_\alpha T(J_\alpha X,Y)+J_\alpha T(X,J_\alpha Y)\right).$$
We recall that if a 3-form $`\psi `$ is of type (1,2)+(2,1) with respect to an almost complex structure $`J`$ then it satisfies the equality
(2.4)
$$\psi (X,Y,Z)=\psi (JX,JY,Z)+\psi (JX,Y,JZ)+\psi (X,JY,JZ).$$
Definition. An almost quaternionic hermitian manifold $`(M,g,(H_\alpha )Q)`$ is QKT manifold if it admits a metric quaternionic connection $``$ with totally skew symmetric torsion which is (1,2)+(2,1)-form with respect to each $`J_\alpha ,\alpha =1,2,3`$. If the torsion 3-form is closed then the manifold is said to be strong QKT manifold.
It follows that the holonomy group of $``$ is a subgroup of SP(n).SP(1).
By means of (2.1), (2) and (2.4), the Nijenhuis tensor $`N_\alpha `$ of $`J_\alpha ,\alpha =1,2,3`$, on a QKT manifold is given by
(2.5)
$$N_\alpha (X,Y)=A_\alpha (Y)J_\beta XA_\alpha (X)J_\beta YJ_\alpha A_\alpha (Y)J_\gamma X+J_\alpha A_\alpha (X)J_\gamma Y,$$
where
(2.6)
$$A_\alpha =\omega _\beta +J_\alpha \omega _\gamma .$$
Remark 1. The definition of QKT manifolds given above is equivalent to that given in because the requirement the torsion to be (1,2)+(2,1)-form with respect to each $`J_\alpha ,\alpha =1,2,3`$, is equivalent, by means of (2.5), to the fourth condition of (4) in . The torsion of $``$ is (1,2)+(2,1)-form with respect to any (local) almost complex structure $`JQ`$ . This follows also from (2.5) and the general formula (6) in which expresses $`N_J`$ in terms of $`N_{J_1},N_{J_2},N_{J_3}`$. In fact, it is sufficient that the torsion is a (1,2)+(2,1)-form with respect to the only two almost complex structures of $`(H)`$ since the formula (3.4.4) in ) gives the necessary expression of $`N_{J_3}`$ by $`N_{J_1}`$ and $`N_{J_2}`$. Indeed, it is easy to see that the formula (3.4.4) in holds for the (0,2)-part $`T_\alpha ^{0,2},\alpha =1,2,3`$, of the torsion. Hence, the vanishing of the (0,2)-part of the torsion with respect to any two almost complex structures in $`(H)`$ implies the vanishing of the (0,2)-part of $`T`$ with respect to the third one.
On a QKT manifold there are three naturally associated 1-forms to the torsion defined by
(2.7)
$$t_\alpha (X)=\frac{1}{2}\underset{i=1}{\overset{4n}{}}T(X,e_i,J_\alpha e_i),\alpha =1,2,3.$$
We have
###### Proposition 2.1
On a QKT manifold $`J_1t_1=J_2t_2=J_3t_3.`$
Proof. Applying (2.4) with respect to $`J_\beta `$ we obtain
$`t_\alpha (X)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{4n}{}}}T(X,e_i,J_\alpha e_i)={\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{4n}{}}}T(X,J_\beta e_i,J_\gamma e_i)`$
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{4n}{}}}T(J_\beta X,e_i,J_\gamma e_i){\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{4n}{}}}T(J_\beta X,J_\beta e_i,J_\alpha e_i)+{\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{4n}{}}}T(X,e_i,J_\alpha e_i).`$
The last equality implies $`t_\alpha =J_\beta t_\gamma `$ which proves the assertion. Q.E.D.
The 1-form $`t=J_\alpha t_\alpha `$ is independent of the chosen almost complex structure $`J_\alpha `$ by Proposition 2.1. We shall call it the torsion 1-form of a given QKT manifold.
Remark 2. Every QKT manifold is a quaternionic manifold. This is an immediate consequence of (2.5) and Proposition 2.3 in .
However, the converse to the above property is not always true. In fact, we have
###### Theorem 2.2
Let $`(M,g,(J_\alpha Q)`$ be a 4n-dimensional ($`n>1`$) quaternionic manifold with $`Q`$-hermitian metric $`g`$. Then $`M`$ admits a QKT structure if and only if the following conditions hold
(2.8)
$$(d_\alpha F_\alpha )^+(d_\beta F_\beta )^+=\frac{1}{2}\left(K_\alpha F_\beta J_\beta K_\beta F_\alpha (K_\beta J_\alpha K_\alpha )F_\gamma \right),$$
where $`(d_\alpha F_\alpha )^+`$ denotes the (1,2)+(2,1) part of $`(d_\alpha F_\alpha )`$ with respect to the $`J_\alpha ,\alpha =1,2,3`$. The 1-forms $`K_\alpha ,\alpha =1,2,3`$, are given by
(2.9)
$$K_\alpha =\frac{1}{1n}\left(J_\beta \theta _\alpha +\theta _{\alpha ,\gamma }\right).$$
The metric quaternionic connection $``$ with torsion 3-form of type (1,2)+(2,1) is unique and is determined by
(2.10)
$$=^g+\frac{1}{2}\left((d_\alpha F_\alpha )^+\frac{1}{2}\left(J_\alpha K_\alpha F_\gamma +K_\alpha F_\beta \right)\right),$$
where $`^g`$ is the Levi-Civita connection of $`g`$.
Proof. To prove the ’if’ part, let $``$ be a metric quaternionic connection satisfying (2.1) which torsion $`T`$ has the required properties. We follow the scheme in . Since $`T`$ is skew-symmetric we have
(2.11)
$$=^g+\frac{1}{2}T.$$
We obtain using (2.1) and (2.11) that
$`{\displaystyle \frac{1}{2}}(T(X,J_\alpha Y,Z)+(T(X,Y,J_\alpha Z))`$ $`=`$ $`g((_X^gJ_\alpha )Y,Z)`$
$`+`$ $`\omega _\beta (X)F_\gamma (Y,Z)\omega _\gamma (X)F_\beta (Y,Z).`$
The tensor $`^gJ_\alpha `$ is decomposed by parts according to $`J_\alpha =(J_\alpha )^{2,0}+(J_\alpha )^{0,2}`$, where
(2.13)
$$g((_X^gJ_\alpha )^{2,0}Y,Z)=\frac{1}{2}\left((dF_\alpha )^+(X,J_\alpha Y,J_\alpha Z)(dF_\alpha )^+(X,Y,Z)\right)$$
(2.14)
$$g((_X^gJ_\alpha )^{0,2}Y,Z)=\frac{1}{2}\left(g(N_\alpha (X,Y),J_\alpha Z)g(N_\alpha (X,Z),J_\alpha Y)g(N_\alpha (Y,Z),J_\alpha X)\right)$$
Taking the (2,0) part in (2) we obtain using (2.13) that
$`T(X,J_\alpha Y,Z)+T(X,Y,J_\alpha Y)`$ $`=`$ $`(dF_\alpha ^+(X,J_\alpha Y,J_\alpha Z)(dF_\alpha ^+(X,Y,Z)`$
$`+`$ $`C_\alpha (X)F_\gamma (Y,Z)+C_\alpha (J_\alpha X)F_\beta (Y,Z),`$
where
(2.16)
$$C_\alpha =\omega _\beta J_\alpha \omega _\gamma .$$
The cyclic sum of (2) and the fact that $`T`$ and $`(dF_\alpha )^+`$ are (1,2)+(2,1)-forms with respect to each $`J_\alpha `$, gives
(2.17)
$$T=(d_\alpha F_\alpha )^+\frac{1}{2}\left(J_\alpha C_\alpha F_\gamma +C_\alpha F_\beta \right).$$
Further, we take the contractions in (2.17) to get
$`J_\alpha t_\alpha =\theta _\alpha J_\beta C_\alpha ,`$
(2.18) $`J_\alpha t_\alpha =J_\gamma \theta _{\beta ,\alpha }nJ_\gamma C_\beta ,`$
$`J_\alpha t_\alpha =J_\beta \theta _{\gamma ,\alpha }nJ_\alpha C_\gamma `$
Using Proposition 2.1, (2.6) and (2.16), we obtain consequently from (2) that
(2.19)
$$A_\alpha =J_\alpha C_\beta +J_\gamma C_\gamma =J_\beta \left(\theta _\gamma \theta _\beta \right),$$
(2.20)
$$(n1)J_\beta C_\alpha =\theta _\alpha J_\beta \theta _{\alpha ,\gamma }.$$
Then (2.8) and (2.9) follow from (2.17) and (2.20).
For the converse, we define $``$ by (2.10). To complete the proof we have to show that $``$ is a quaternionic connection. We calculate
$`g((_XJ_\alpha )Y,Z)`$ $`=`$ $`g((_X^gJ_\alpha )Y,Z)+{\displaystyle \frac{1}{2}}\left(T(X,J_\alpha Y,Z)+T(X,Y,J_\alpha Z)\right)`$
$`=`$ $`\omega _\beta (X)F_\gamma (Y,Z)\omega _\gamma (X)F_\beta (Y,Z),`$
where we used (2.13), (2.14), (2.19), (2.9), (2.6), (2.16) and the compatibility condition (2.8) to get the last equality. The uniqueness of $``$ follows from (2.10) as well as from Theorem 10.3 in which states that any quaternionic connection is entirely determined by its torsion (see also ). Q.E.D.
In the case of HKT manifold, $`K_\alpha =dF_\alpha ^{}=0`$ and Theorem 2.2 is a consequence of the general results in (see also ) which imply that on a hermitian manifold there exists a unique linear connection with totally skew-symmetric torsion preserving the metric and the complex structure, the Bismut connection. This connection was used by Bismut to prove a local index theorem for the Dolbeault operator on non-Kähler manifold. The geometry of this connection is referred to KT-geometry by physicists. Obstructions to the existence of (non-trivial) Dolbeault cohomology groups on a compact KT-manifold are presented in .
We note that (2.19) and (2.20) are also valid in the case $`n=1`$.
We get, as a consequence of the proof of Theorem 2.2, the following integrability criterion which is discovered in dimension 4 in .
###### Proposition 2.3
The Nijenhuis tensors of a QKT manifold depend only on the difference between the Lie forms. In particular, the almost complex structures $`J_\alpha `$ on a QKT manifold $`(M,(J_\alpha )Q,g,)`$ are integrable if and only if
$$\theta _\alpha =\theta _\beta =\theta _\gamma $$
Proof. The Nijenhuis tensors are given by (2.5) and (2.19). Q.E.D.
###### Corollary 2.4
On a 4n-dimensional QKT manifold the following formulas hold
$$J_\beta \theta _{\alpha ,\gamma }=J_\gamma \theta _{\alpha ,\beta },$$
(2.21)
$$(n^2+n)\theta _\alpha n\theta _\beta n^2\theta _\gamma +J_\gamma \theta _{\beta ,\alpha }+nJ_\alpha \theta _{\gamma ,\beta }(n+1)J_\beta \theta _{\alpha ,\gamma }=0.$$
If $`n=1`$ then $`\theta _\alpha =J_\beta \theta _{\alpha ,\gamma }=J_\gamma \theta _{\alpha ,\beta }`$.
Proof. The first formula follows directly from the system (2). Solving the system (2) with respect to $`C_\alpha `$ we obtain
(2.22)
$$(n^31)J_\beta C_\alpha =(\theta _\alpha J_\gamma \theta _{\beta ,\alpha })+n(\theta _\beta J_\alpha \theta _{\gamma ,\beta })+n^2(\theta _\gamma J_\beta \theta _{\alpha ,\gamma }).$$
Then (2.21) is a consequence of (2.22) and (2.20). The last assertion follows from (2.20) . Q.E.D.
###### Corollary 2.5
On a 4n-dimensional ($`n>1`$) QKT manifold the $`sp(1)`$-connection 1-forms are given by
(2.23)
$$\omega _\beta =\frac{1}{2}J_\beta \left(\theta _\gamma \theta _\beta +\frac{1}{1n}\theta _\alpha \right)+\frac{1}{2(1n)}\theta _{\alpha ,\gamma }.$$
Proof. The proof follows in a straightforward way from (2.19), (2.20), (2.6) and (2.16). Q.E.D.
Theorem 2.2 and the above formulas lead to the following criterion
###### Proposition 2.6
Let $`(M,g,(H))`$ be a 4n-dimensional ($`n>1`$) QKT manifold. The following conditions are equivalent:
i) $`(M,g,(H))`$ is a HKT manifold;
ii) $`d_\alpha F_\alpha ^+=d_\beta F_\beta ^+=d_\gamma F_\gamma ^+`$;
iii) $`\theta _\alpha =J_\beta \theta _{\gamma ,\alpha }.`$
Proof. If $`(M,g,(H))`$ is a HKT manifold, the connection 1-forms $`\omega _\alpha =0,\alpha =1,2,3`$. Then ii) and iii) follow from (2.16), (2.20), (2.9) and (2.8).
If iii) holds, then (2.20) and (2.19) yield $`C_\alpha =A_\alpha =0,\alpha =1,2,3`$, since $`n>1`$. Consequently, $`2\omega _\alpha =J_\beta C_\beta J_\beta A_\beta =0`$ by (2.16) and (2.6). Thus the equivalence of i) and iii) is proved.
Let ii) holds. Then we compute that $`\theta _\alpha =J_\gamma \theta _{\beta ,\alpha }`$. Since $`n>1`$, the equality (2.22) leads to $`C_\alpha =0,\alpha =1,2,3`$, which forces $`\omega _\alpha =0,\alpha =1,2,3`$ as above. This completes the proof. Q.E.D.
The next theorem shows that QKT manifolds are stable under a conformal transformations.
###### Theorem 2.7
Let $`(M,g,(J_\alpha ),)`$ be a 4n-dimensional QKT manifold. Then every Riemannian metric $`\overline{g}`$ in the conformal class \[g\] admits a QKT connection. If $`\overline{g}=fg`$ for a positive function $`f`$ then the QKT connection $`\overline{}`$ corresponding to $`\overline{g}`$ is given by
$`\overline{g}(\overline{}_XY,Z)=fg(_XY,Z)`$ $`+`$ $`{\displaystyle \frac{1}{2}}\left(df(X)g(Y,Z)+df(Y)g(X,Z)df(Z)g(X,Y)\right)`$
$`+`$ $`{\displaystyle \frac{1}{2}}\left(J_\alpha dfF_\alpha +J_\beta dfF_\beta +J_\gamma dfF_\gamma \right)(X,Y,Z).`$
If $`M`$ is compact then there exists a unique (up to homotety) metric $`g_G[g]`$ with co-closed torsion 1-form.
Proof. First we assume $`n>1`$. We shall apply Theorem 2.2 to the quaternionic Hermitian manifold $`(M,\overline{g}=fg,(J_\alpha )Q)`$. We denote the objects corresponding to the metric $`\overline{g}`$ by a line above the symbol e.g. $`\overline{F_\alpha }`$ denotes the Kähler form of $`J_\alpha `$ with respect to $`\overline{g}`$. An easy calculation gives the following sequence of formulas
(2.25)
$$d_\alpha \overline{F}_\alpha ^+=J_\alpha dfF_\alpha +fd_\alpha F_\alpha ^+;\overline{\theta }_\alpha =\theta _\alpha +(2n1)d\mathrm{ln}f;\overline{\theta }_{\alpha ,\gamma }=\theta _{\alpha ,\gamma }J_\beta d\mathrm{ln}f.$$
We substitute (2.25) into (2.9), (2.19) and (2.23) to get
(2.26)
$$\overline{K}_\alpha =K_\alpha 2J_\beta d\mathrm{ln}f,\overline{A}=A,\overline{\omega }_\alpha =\omega _\alpha J_\beta d\mathrm{ln}f.$$
Using (2.25) and (2.26) we verify that the conditions (2.8) with respect to the metric $`\overline{g}`$ are fulfilled. Theorem 2.2 implies that there exists a QKT connection $`\overline{}`$ with respect to $`(\overline{g},Q)`$. Using the well known relation between the Levi-Civita connections of conformally equivalent metrics, (2.25) and (2.26), we obtain (2.7) from (2.10).
If $`n=1`$ we define the new QKT connection with respect to $`(\overline{g},Q)`$ by (2.7).
Using (2.7), we find that the torsion tensors $`T`$ and $`\overline{T}`$ of $``$ and $`\overline{}`$ are related by
(2.27)
$$\overline{T}=fT+J_\alpha dfF_\alpha +J_\beta dfF_\beta +J_\gamma dfF_\gamma .$$
Consequently, we obtain from (2.27) for the torsion 1-forms $`t`$ and $`\overline{t}`$ that
(2.28)
$$\overline{t}=t(2n+1)d\mathrm{ln}f.$$
If $`M`$ is compact, we may apply to (2.28) the theorem of Gauduchon for the existence of a Gauduchon metric on a compact Weyl manifold to obtain the desired metric $`g_G`$. Q.E.D.
We shall call the unique metric with co-closed torsion 1-form on a compact QKT manifold the Gauduchon metric.
###### Corollary 2.8
On a compact QKT manifold with closed (non exact) torsion 1-form the Gauduchon metric $`g_G`$ cannot have positive definite Riemannian Ricci tensor. In particular, if it is an Einstein manifold then it is of non-positive scalar curvature.
Further, if the Gauduchon metric is Ricci flat then the corresponding torsion 1-form $`t_G`$ is parallel with respect to the Levi-Civita connection of $`g_G`$.
Proof. The two form $`dt`$ is invariant under conformal transformations by (2.28). Then the Gauduchon metric has harmonic torsion 1-form i.e $`dt=\delta t=0`$. The claim follows from the Weitzenböeck formula (see e.g. ) $`_M\{|dt|^2+|\delta t|^2\}𝑑V=_M\{|^gt|^2+Ric^g(t^\mathrm{\#},t^\mathrm{\#})\}𝑑V=0,`$ where $`t^\mathrm{\#}`$ is the dual vector field of $`t`$, $`|.|`$ is the usual tensor norm and $`dV`$ is the volume form. Q.E.D.
Theorem 2.7 allows us to supply a large class of (compact) QKT manifold. Namely, any conformal metric of a QK, HK or HKT manifold will give a QKT manifold. This leads to the notion of locally conformally QK (resp. locally conformally HK, resp. locally conformally HKT) manifolds (briefly l.c.QK (resp. l.c.HK, resp. l.c.HKT) manifolds) in the context of QKT geometry.
The l.c.QK and l.c.HK manifolds have already appeared in the context of Hermitian-Einstein-Weyl structures and of 3-Sasakian structures . These two classes of quaternionic manifolds are studied in detail (mostly in the compact case) in .
We recall that a quaternionic Hermitian manifold $`(M,g,Q)`$ is said to be l.c.QK (resp. l.c.HK, resp. l.c.HKT) manifold if each point $`pM`$ has a neighbourhood $`U_p`$ such that $`g|_{U_p}`$ is conformally equivalent to a QK (resp.HK, resp.HKT) metric. There are compact l.c.QK manifold which do not admit any QK structure . Typical examples of compact l.c. QK manifolds without any QK structure are the quaternionic Hopf spaces $`H=(^n\{0\})/\mathrm{\Gamma }`$, where $`\mathrm{\Gamma }`$ is an appropriate discrete group acting diagonally on the quaternionic coordinates in $`^n`$ (see ).
We recall that on a l.c.QK manifold the 4-form $`\mathrm{\Omega }=_{\alpha =1}^3F_\alpha F_\alpha `$ satisfies $`d\mathrm{\Omega }=\omega \mathrm{\Omega },d\omega =0`$, where $`\omega `$ is locally defined by $`\omega =2d\mathrm{ln}f`$. On a l.c.QK manifold viewed as a QKT manifold by Theorem 2.7 the torsion 1-form is equal to $`t=(2n+1)\omega `$ by (2.28). The QK manifolds are Einstein provided the dimension is at least 8 . Then, the Gauduchon Theorem applied to l.c.QK manifold in can be stated in our context as follows
###### Corollary 2.9
Let $`(M,g)`$ be a compact 4n-dimensional ($`n>1`$) QKT manifold which is l.c.QK and assume that no metric in the conformal class \[g\] of g is QK. Then the torsion 1-form of the Gauduchon metric $`g_G`$ is parallel with respect to the Levi-Civita connection of $`g_G`$.
Theorem 2.7, Theorem 2.2 together with Proposition 2.3 and Proposition 2.6 imply the following
###### Corollary 2.10
Every l.c.QK manifold admits a QKT structure.
Further, if $`(M,g,(J_\alpha ),)`$ is a 4n-dimensional $`n>1`$ QKT manifold then:
i) $`(M,g,(J_\alpha ),)`$ is a l.c.QK manifold if and only if
(2.29)
$$T=\frac{1}{2n+1}\left(t_\alpha F_\alpha +t_\beta F_\beta +t_\gamma F_\gamma \right),dt=0;$$
ii) $`(M,g,(J_\alpha ),)`$ is a l.c.HKT manifold if and only if the 1-form $`\theta _\alpha J_\beta \theta _{\alpha ,\gamma }`$ is closed i.e.
$$d(\theta _\alpha J_\beta \theta _{\alpha ,\gamma })=0;$$
iii) $`(M,g,(J_\alpha ),)`$ is a l.c.HK manifold if an only if (2.29) holds and
$$\theta _\alpha J_\beta \theta _{\alpha ,\gamma }=\frac{2(1n)}{2n+1}t.$$
## 3 Curvature of a QKT space
Let $`R=[,]_[,]`$ be the curvature tensor of type (1,3) of $``$. We denote the curvature tensor of type (0,4) $`R(X,Y,Z,V)=g(R(X,Y)Z,V)`$ by the same letter. There are three Ricci forms given by
$$\rho _\alpha (X,Y)=\frac{1}{2}\underset{i=1}{\overset{4n}{}}R(X,Y,e_i,J_\alpha e_i),\alpha =1,2,3.$$
###### Proposition 3.1
The curvature of a QKT manifold $`(M,g,(J_\alpha ),)`$ satisfies the following relations
(3.30)
$$R(X,Y)J_\alpha =\frac{1}{n}\left(\rho _\gamma (X,Y)J_\beta \rho _\beta (X,Y)J_\gamma \right),$$
(3.31)
$$\rho _\alpha =d\omega _\alpha +\omega _\beta \omega _\gamma .$$
Proof. We follow the classical scheme (see e.g. ). Using (2.1) we obtain
$$R(X,Y)J_\alpha =(d\omega _\beta +\omega _\gamma \omega _\alpha )(X,Y)J_\gamma +(d\omega _\gamma +\omega _\alpha \omega _\beta )(X,Y)J_\beta .$$
Taking the trace in the last equality, we get
$`\rho _\alpha (X,Y)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{4n}{}}}R(X,Y,e_i,J_\alpha e_i)={\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{4n}{}}}R(X,Y,J_\beta e_i,J_\gamma e_i)`$
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{4n}{}}}R(X,Y,e_i,J_\alpha e_i)+2n(d\omega _\alpha +\omega _\beta \omega _\gamma )(X,Y)J_\beta .`$
Q.E.D.
Using Proposition 3.1 we find a simple necessary and sufficient condition a QKT manifold to be a HKT one, i.e. the holonomy group of $``$ to be a subgroup of Sp(n).
###### Proposition 3.2
A 4n-dimensional $`(n>1)`$ QKT manifold is a HKT manifold if and only if all the three Ricci forms vanish, i.e $`\rho _1=\rho _2=\rho _3=0`$.
Proof. If a QKT manifold is a HKT manifold then the holonomy group of $``$ is contained in Sp(n). This implies $`\rho _\alpha =0,\alpha =1,2,3`$.
For the converse, let the three Ricci forms vanish. The equations (3.31) mean that the curvature of the Sp(1) connection on $`Q`$ vanish. Then there exists a basis $`(I_\alpha ,\alpha =1,2,3)`$ of almost complex structures on $`Q`$ and each $`I_\alpha `$ is $``$-parallel i.e. the corresponding connection 1-forms $`\omega _{I_\alpha }=0,\alpha =1,2,3`$. Then each $`I_\alpha `$ is a complex structure, by (2.5) and (2.6). This implies that the QKT manifold is a HKT manifold. Q.E.D.
We denote by $`Ric,Ric^g`$ the Ricci tensors of the QKT connection and of the Levi-Civita connection, respectively. In fact $`Ric(X,Y)=_{i=1}^{4n}R(e_i,X,Y,e_i)`$.
Our main technical result is the following
###### Proposition 3.3
Let $`(M,g,(J_\alpha ),)`$ be a 4n-dimensional QKT manifold. The following formulas hold
(3.32) $`n\rho _\alpha (X,J_\alpha Y)+\rho _\beta (X,J_\beta Y)+\rho _\gamma (X,J_\gamma Y)=`$
$`nRic(XY)+{\displaystyle \frac{n}{4}}(dT)_\alpha (X,J_\alpha Y)+{\displaystyle \frac{n}{2}}(T)_\alpha (X,J_\alpha Y);`$
(3.33) $`(n1)\rho _\alpha (X,J_\alpha Y)={\displaystyle \frac{n(n1)}{n+2}}Ric(X,Y)`$
$`+{\displaystyle \frac{n}{4(n+2)}}\left\{(n+1)(dT)_\alpha (X,J_\alpha Y)(dT)_\beta (X,J_\beta Y)(dT)_\gamma (X,J_\gamma Y)\right\}`$
(3.34) $`+{\displaystyle \frac{n}{2(n+2)}}\left\{(n+1)(T)_\alpha (X,J_\alpha Y)(T)_\beta (X,J_\beta Y)(T)_\gamma (X,J_\gamma Y)\right\},`$
where
$`(dT)_\alpha (X,Y)=_{i=1}^{4n}dT(X,Y,e_i,J_\alpha e_i),(T)_\alpha (X,Y)=_{i=1}^{4n}(_XT)(Y,e_i,J_\alpha e_i)`$.
Proof. Since the torsion is a 3-form, we have
(3.35)
$$(_X^gT)(Y,Z,U)=(_XT)(Y,Z,U)+\frac{1}{2}\genfrac{}{}{0pt}{}{\sigma }{XYZ}\{g(T(X,Y),T(Z,U)\},$$
where $`\genfrac{}{}{0pt}{}{\sigma }{XYZ}`$ denote the cyclic sum of $`X,Y,Z`$.
The exterior derivative $`dT`$ is given by
$`dT(X,Y,Z,U)`$ $`=`$ $`{\displaystyle \genfrac{}{}{0pt}{}{\sigma }{XYZ}}\{(_XT)(Y,Z,U)+g(T(X,Y),T(Z,U)\}`$
$``$ $`(_UT)(X,Y,Z)+{\displaystyle \genfrac{}{}{0pt}{}{\sigma }{XYZ}}\{g(T(X,Y),T(Z,U)\}.`$
The first Bianchi identity for $``$ states
(3.37)
$$\genfrac{}{}{0pt}{}{\sigma }{XYZ}R(X,Y,Z,U)=\genfrac{}{}{0pt}{}{\sigma }{XYZ}\{(_XT)(Y,Z,U)+g(T(X,Y),T(Z,U)\}.$$
We denote by $`B`$ the Bianchi projector i.e. $`B(X,Y,Z,U)=\genfrac{}{}{0pt}{}{\sigma }{XYZ}R(X,Y,Z,U)`$.
The curvature $`R^g`$ of the Levi-Civita connection is connected by $`R`$ in the following way
(3.38) $`R^g(X,Y,Z,U)`$ $`=`$ $`R(X,Y,Z,U){\displaystyle \frac{1}{2}}(_XT)Y,Z,U)+{\displaystyle \frac{1}{2}}(_YT)X,Z,U)`$
$``$ $`{\displaystyle \frac{1}{2}}g(T(X,Y),T(Z,U)){\displaystyle \frac{1}{4}}g(T(Y,Z),T(X,U)){\displaystyle \frac{1}{4}}g(T(Z,X),T(Y,U)).`$
Define $`D`$ by $`D(X,Y,Z,U)=R(X,Y,Z,U)R(Z,U,X,Y)`$, we obtain from (3.38)
(3.39) $`D(X,Y,Z,U)=`$
$`{\displaystyle \frac{1}{2}}(_XT)(Y,Z,U){\displaystyle \frac{1}{2}}(_YT)(X,Z,U){\displaystyle \frac{1}{2}}(_ZT)(U,X,Y)+{\displaystyle \frac{1}{2}}(_UT)(Z,X,Y),`$
since $`D^g`$ of $`R^g`$ is zero.
Using (3.30) and (3.37) we find the following relation between the Ricci tensor and the Ricci forms
$`\rho _\alpha (X,Y)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{4n}{}}}\left(R(Y,e_i,X,J_\alpha e_i)+R(e_i,X,Y,J_\alpha e_i)\right)+{\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{4n}{}}}B(X,Y,e_i,J_\alpha e_i)`$
$`=`$ $`{\displaystyle \frac{1}{2}}Ric(Y,J_\alpha X)+{\displaystyle \frac{1}{2}}Ric(X,J_\alpha Y)+{\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{4n}{}}}B(X,Y,e_i,J_\alpha e_i)`$
$`+`$ $`{\displaystyle \frac{1}{2n}}\left\{\rho _\beta (J_\gamma Y,X)\rho _\beta (J_\gamma X,Y)+\rho _\gamma (J_\beta X,Y)\rho _\gamma (J_\beta Y,X)\right\}.`$
On the other hand, using (3.30), we calculate
$`{\displaystyle \underset{i=1}{\overset{4}{}}}D(X,e_i,J_\alpha e_i,Y)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{4n}{}}}\{R(X,e_i,J_\alpha e_i,Y)+R(Y,e_i,J_\alpha e_iX)\}`$
$`=`$ $`Ric(Y,J_\alpha X)Ric(X,J_\alpha Y)`$
$`+`$ $`{\displaystyle \frac{1}{n}}\left\{\rho _\beta (X,J_\gamma Y)+\rho _\beta (Y,J_\gamma X)\rho _\gamma (Y,J_\beta X)\rho _\gamma (X,J_\beta Y)\right\}.`$
Combining (3) and (3), we derive
(3.42)
$$n\rho _\alpha (X,J_\alpha Y)+\rho _\beta (X,J_\beta Y)+\rho _\gamma (X,J_\gamma Y)=$$
$$nRic(XY)+\frac{n}{2}B_\alpha (X,J_\alpha Y)+\frac{n}{2}D_\alpha (X,J_\alpha Y),$$
where the tensors $`B_\alpha `$ and $`D_\alpha `$ are defined by $`B_\alpha (X,Y)=_{i=1}^{4n}B(X,Y,e_i,J_\alpha e_i)`$ and
$`D_\alpha (X,Y)=_{i=1}^{4n}D(X,e_i,J_\alpha e_i,Y)`$. Taking into account (3.39), we get the expression
(3.43)
$$D_\alpha (X,Y)=\frac{1}{2}\underset{i=1}{\overset{4n}{}}(_XT)(Y,e_i,J_\alpha e_i)+\frac{1}{2}\underset{i=1}{\overset{4n}{}}(_YT)(X,e_i,J_\alpha e_i)\alpha =1,2,3.$$
To calculate $`B_\alpha +D_\alpha `$ we use (3) twice and (3.43). After some calculations, we derive
(3.44)
$$B_\alpha (X,Y)+D_\alpha (X,Y)=\frac{1}{2}\underset{i=1}{\overset{4n}{}}dT(X,Y,e_i,J_\alpha e_i)+\underset{i=1}{\overset{4n}{}}(_XT)(Y,e_i,J_\alpha e_i),\alpha =1,2,3.$$
We substitute (3.44) into (3.42). Solving the obtained system, we obtain
(3.45) $`(n1)\left\{\rho _\alpha (X,J_\alpha Y)\rho _\beta (X,J_\beta Y)\right\}=`$
$`{\displaystyle \frac{n}{2}}\left\{(dT)_\alpha (X,J_\alpha Y)(dT)_\beta (X,J_\beta Y)\right\}+{\displaystyle \frac{n}{2}}\left\{(T)_\alpha (X,J_\alpha Y)(T)_\beta (X,J_\beta Y)\right\}.`$
Finally, (3.42) and (3.45) imply (3.33). Q.E.D.
Remark 3. The Ricci tensor of a QKT connection is not symmetric in general. From (3.37), (3.35) and the fact that $`T`$ is a 3-form we get the formula $`Ric(X,Y)Ric(Y,X)=_{i=1}^{4n}(_{e_i}^gT)(e_i,X,Y)=\delta T(X,Y).`$ Hence, the Ricci tensor of a linear connection with totally skew-symmetric torsion is symmetric if and only if the torsion 3-form is co-closed.
## 4 QKT manifolds with parallel torsion and homogeneous QKT structures
Let $`(G/K,g)`$ be a reductive (locally) homogeneous Riemannian manifold. The canonical connection $``$ is characterised by the properties $`g=T=R=0`$ ,p.193. A homogeneous quaternionic Hermitian manifold (resp. homogeneous hyper Hermitian) manifold $`(G/K,g,Q)`$ is a homogeneous Riemannian manifold with an invariant quaternionic Hermitian subbundle $`Q`$ (resp. three invariant anti commuting complex structures ). This means that the bundle $`Q`$ (resp. each of the three complex structures) is parallel with respect to the canonical connection $``$. The torsion of $``$ is totally skew-symmetric if and only if the homogeneous Riemannian manifold is naturally reductive (see also . Homogeneous QKT (resp. HKT) manifolds are homogeneous quaternionic Hermitian (resp. homogeneous hyper Hermitian) manifold which are naturally reductive. Examples of homogeneous HKT and QKT manifolds are presented in . The homogeneous QKT manifolds in are constructed from homogeneous HKT manifolds.
In this section we generalise the result of which states that there are no homogeneous QKT manifold with torsion 4-form $`dT`$ of type (2,2) in dimensions greater than four. First, we prove the following technical result
###### Proposition 4.1
Let $`(M,g,(J_\alpha ),)`$ be a 4n-dimensional $`(n>1)`$ QKT manifold with 4-form $`dT`$ of type (2,2) with respect to each $`J_\alpha ,\alpha =1,2,3`$. Suppose that the torsion is parallel with respect to the QKT-connection. Then the Ricci forms $`\rho _\alpha `$ are given by
(4.46)
$$\rho _\alpha (X,J_\alpha Z)=\lambda g(X,Y),\alpha =1,2,3,$$
where $`\lambda `$ is a smooth function on $`M`$.
Proof. Let the torsion be parallel i.e. $`T=0`$. Remark 3 shows that the Ricci tensor is symmetric. The equalities (3) and (3.37) imply
(4.47)
$$B(X,Y,Z,U)=\genfrac{}{}{0pt}{}{\sigma }{XYZ}\{g(T(X,Y),T(Z,U)\}=\frac{1}{2}dT(X,Y,Z,U).$$
We get $`D=0`$ from (3.39).
Suppose now that the 4-form $`dT`$ is of type (2,2) with respect to each $`J_\alpha ,\alpha =1,2,3.`$. Then it satisfies the equalities
(4.48)
$$dT(X,Y,Z,U)=dT(J_\alpha X,J_\alpha Y,Z,U)+dT(J_\alpha X,Y,J_\alpha Z,U)+dT(X,J_\alpha Y,J_\alpha Z,U).$$
The similar arguments as we used in the proof of Proposition 2.1 but applying (4.48) instead of (2.4), yield
###### Lemma 4.2
On a QKT manifold with 4-form $`dT`$ of type (2,2) with respect to each $`J_\alpha ,\alpha =1,2,3`$, the following equalities hold:
(4.49)
$$(dT)_1(X,J_1Y)=(dT)_2(X,J_2Y)=(dT)_3(X,J_3Y),$$
(4.50)
$$(dT)_\alpha (X,J_\alpha Y)=(dT)_\alpha (J_\alpha X,Y),\alpha =1,2,3.$$
We substitute (4.49), (4.47) and $`D=0`$ into (3.45) and (3.33) to get
(4.51)
$$\rho _1(X,J_1Y)=\rho _2(X,J_2Y)=\rho _3(X,J_3Y),$$
(4.52)
$$\rho _\alpha (X,J_\alpha Y)=\frac{n}{n+2}Ric(X,Y)+\frac{n}{4(n+2)}(dT)_\alpha (X,J_\alpha Y),\alpha =1,2,3.$$
The equality (4.50) shows that the 2-form $`dT_\alpha `$ is a (1,1)-form with respect to $`J_\alpha `$. Hence, the $`dT_\alpha `$ is (1,1)-form with respect to each $`J_\alpha ,\alpha =1,2,3`$, because of (4.49). Since the Ricci tensor $`Ric`$ is symmetric, (4.52) shows that the Ricci tensor $`Ric`$ is of hybrid type with respect to each $`J_\alpha `$ i.e. $`Ric(J_\alpha X,J_\alpha Y)=Ric(X,Y),\alpha =1,2,3`$ and the Ricci forms $`\rho _\alpha ,\alpha =1,2,3`$ are (1,1)-forms with respect to all $`J_\alpha ,\alpha =1,2,3`$. Taking into account (3.30), we obtain
$`R(X,J_\alpha X,Z,J_\alpha Z)+R(X,J_\alpha X,J_\beta Z,J_\gamma Z)`$
$`+`$ $`R(J_\beta X,J_\gamma X,Z,J_\alpha Z)+R(J_\beta X,J_\gamma X,J_\beta Z,J_\gamma Z)`$
$`=`$ $`{\displaystyle \frac{1}{n}}\left(\rho _\alpha (X,J_\alpha X)+\rho _\alpha (J_\beta X,J_\gamma X)\right)g(Z,Z)={\displaystyle \frac{2}{n}}\rho _\alpha (X,J_\alpha X)g(Z,Z),`$
where the last equality of (4) is a consequence of the following identity
$$\rho _\alpha (J_\beta X,J_\gamma X)=\rho _\beta (J_\beta X,X)=\rho _\alpha (X,J_\alpha X).$$
The left side of (4) is symmetric with respect to the vectors $`X,Z`$ because $`D=0`$. Hence, $`\rho _\alpha (X,J_\alpha X)g(Z,Z)=\rho _\alpha (Z,J_\alpha Z)g(X,X),\alpha =1,2,3`$. The last equality together with (4.51) implies (4.46). Q.E.D.
###### Theorem 4.3
Let $`(M,g,(J_\alpha ))`$ be a 4n-dimensional ($`n>1`$) QKT manifold with 4-form $`dT`$ of type (2,2) with respect to each $`J_\alpha ,\alpha =1,2,3`$. Suppose that the torsion is parallel with respect to the QKT-connection. Then $`(M,g,(J_\alpha ))`$ is either a HKT manifold with parallel torsion or a QK manifold.
Proof. We apply Proposition 4.1. If the function $`\lambda =0`$ then $`\rho _\alpha =0,\alpha =1,2,3`$, by (4.46) and Proposition 3.2 implies that the QKT manifold is actually a HKT manifold.
Let $`\lambda 0`$. The condition (4.46) determines the torsion completely. We proceed involving (3.31) into the computations as in . We calculate using (2.1) and (4.46) that
(4.54)
$$(_Z\rho _\alpha )(X,Y)=\lambda \left\{\omega _\beta (Z)F_\gamma (X,Y)\omega _\gamma (Z)F_\beta (X,Y)\right\}d\lambda (Z)F_\alpha (X,Y).$$
Applying the operator $`d`$ to (3.30), we get taking into account (4.46) that
(4.55)
$$d\rho _\alpha =\lambda (F_\beta \omega _\gamma \omega _\beta F_\gamma )$$
On the other hand, we have
(4.56)
$$d\rho _\alpha =\genfrac{}{}{0pt}{}{\sigma }{XYZ}\{(_Z\rho _\alpha )(X,Y)+\lambda (T(X,Y,J_\alpha Z)\},\alpha =1,2,3.$$
Comparing the left-hand sides of (4.55) and (4.56) and using (4.54), we derive
$$\lambda \genfrac{}{}{0pt}{}{\sigma }{XYZ}\left\{(T(X,Y),J_\alpha Z)\right\}=d\lambda F_\alpha (X,Y,Z),\alpha =1,2,3.$$
The last equality implies $`\lambda T=J_\alpha d\lambda F_\alpha ,\alpha =1,2,3.`$ If $`\lambda `$ is a non zero constant then $`T=0`$ and we recover the result of . If $`\lambda `$ is not a constant then there exists a point $`pM`$ and a neighbourhood $`V_p`$ of $`p`$ such that $`\lambda |_{V_p}0`$. Then
(4.57)
$$T=J_\alpha d\mathrm{ln}\lambda F_\alpha ,\alpha =1,2,3.$$
We take the trace in (4.57) to obtain
(4.58)
$$4(n1)J_\alpha d\mathrm{ln}\lambda =0,\alpha =1,2,3.$$
The equation (4.58) forces $`d\lambda =0`$ since $`n>1`$ and consequently $`T=0`$ by (4.57). Hence, the QKT space is a QK manifold which completes the proof. Q.E.D.
On a locally homogeneous QKT manifold the torsion and curvature are parallel and Theorem 4.3 leads to the following
###### Theorem 4.4
A (locally) homogeneous 4n-dimensional $`(n>1)`$ QKT manifold with torsion 4-form $`dT`$ of type (2,2) is either (locally) homogeneous HKT space or a (locally) symmetric QK space.
Theorem 4.4 shows that there are no homogeneous (proper) QKT manifolds with torsion 4-form of type (2,2) in dimensions greater than four which is proved in by different methods using the Lie algebra arguments.
## 5 Four dimensional QKT manifolds
In dimension 4 the situation is completely different from that described in Theorem 2.2 and Theorem 4.3 in higher dimensions. For a given quaternionic structure on a 4-dimensional manifold $`(M,g(H))`$ (or equivalently, given an orientation and a conformal class of Riemannian metrics ) there are many QKT structures . More precisely, all QKT structures associated with $`(g,(H))`$ depend on a 1-form $`\psi `$ due to the general identity
(5.59)
$$\psi =J\psi F,$$
where $``$ is the Hodge $``$-operator, $`J`$ is an $`g`$-orthogonal almost complex structure with Kähler form $`F`$ (see ). Indeed, for any given 1-form $`\psi `$ we may define a QKT-connection $``$ as follows: $`=^g+\frac{1}{2}\psi `$. Conversely, any 3-form $`T`$ can be represented by $`T=(T)`$ and the connection given above is a quaternionic connection with torsion $`T=\psi `$. Hence, a QKT structure on a 4-dimensional oriented manifold is a pair $`(g,t)`$ of a Riemannian metric $`g`$ and an 1-form $`t`$. The choice of $`g`$ generates three almost complex structures $`(J_\alpha ),\alpha =1,2,3`$, satisfying the quaternionic identities . The torsion 3-form T is given by
(5.60)
$$T=t=t_\alpha F_\alpha =t_\beta F_\beta =t_\gamma F_\gamma .$$
As consequence of (5.59), we obtain $`dT=dt=\delta t`$. The last identity means that the torsion 3-form $`T`$ is closed if and only if the 1-form $`t`$ is co-closed. Thus, in dimension 4 there are many strong QKT structures.
In higher dimensions the conformal change of the metric induces a unique QKT structure by Theorem 2.7. We may define a QKT connection corresponding to a conformally equivalent metric $`\overline{g}=fg`$ in dimension 4 by (2.7) and call this conformal QKT transformation. In the compact case, taking the Gauduchon metric of Theorem 2.7, we obtain
###### Proposition 5.1
Let $`(M,g,(H),)`$ be a compact 4-dimensional QKT manifold. In the conformal class \[g\] there exists a unique (up to homotety) strong QKT structure conformally equivalent to the given one.
Further, we consider QKT structures with parallel torsion. We have
###### Theorem 5.2
A 4-dimensional QKT manifold $`M`$ with parallel torsion 3-form is a strong QKT manifold, the torsion 1-form is parallel with respect to the Levi-Civita connection and $`M`$ is locally isometric to the product $`N^3\times `$, where $`N^3`$ is a three dimensional Riemannian manifold admitting a Riemannian connection $``$ with totally skew-symmetric torsion, parallel with respect to $``$.
Proof. The proof is based on the following
###### Lemma 5.3
A 4-dimensional QKT manifold has parallel torsion 3-form if and only if it has parallel torsion 1-form with respect to the Levi-Civita connection.
Proof of Lemma 5.3. We calculate using (5.60) and (2.1) that
$`(_ZT)(X,Y,U)`$ $`=`$ $`t_\alpha (U)\left(\omega _\beta (Z)F_\gamma (Y,X)\omega _\gamma (Z)F_\beta (Y,X)\right)`$
$``$ $`t_\alpha (X)\left(\omega _\beta (Z)F_\gamma (Y,U)\omega _\gamma (Z)F_\beta (Y,U)\right)`$
$`+`$ $`t_\alpha (Y)\left(\omega _\beta (Z)F_\gamma (X,U)\omega _\gamma (Z)F_\beta (X,U)\right)`$
$`+`$ $`F_\alpha (Y,U)(_Zt_\alpha )X+F_\alpha (X,Y)(_Zt_\alpha )UF_\alpha (X,U)(_Zt_\alpha )Y.`$
Taking the trace in (5), we obtain
(5.62)
$$\underset{i=1}{\overset{4}{}}(_ZT)(X,e_i,J_\alpha e_i)=2(_Zt_\alpha )X2\left(\omega _\beta (Z)t_\gamma (X)\omega _\gamma (Z)t_\beta (X)\right).$$
Using (2.1), we get
(5.63)
$$(_Zt_\alpha )X=(_Zt)J_\alpha X\left(\omega _\beta (Z)t_\gamma (X)\omega _\gamma (Z)t_\beta (X)\right).$$
The equation (5.63) and (5.62) yield
(5.64)
$$\underset{i=1}{\overset{4}{}}(_ZT)(J_\alpha X,e_i,J_\alpha e_i)=2(_Zt)X,\alpha =1,2,3.$$
Then $`t=0`$ since the torsion is parallel. But $`^gt=t`$ by (2.11) and (5.60). Hence, $`^gt=0`$.
For the converse, we insert (5.63) into (5) to get
(5.65)
$$(_ZT)(X,Y,U)=F_\alpha (Y,U)(_Zt)J_\alpha X+F_\alpha (X,Y)(_Zt)J_\alpha U+F_\alpha (U,X)(_Zt)J_\alpha Y,$$
since the dimension is equal to four. If $`^gt=0`$ then $`t=0`$ and (5.65) leads to $`T=0`$ which proves the lemma. Q.E.D.
Lemma 5.3 shows that $`(M,g)`$ is locally isometric to the Riemannian product $`\times N^3`$ of a real line and a 3-dimensional manifold $`N^3`$ (see e.g. ). Using (5.60) we see that $`T(t^\mathrm{\#},X^{},Y^{})=0`$ for every vector fields $`X^{},Y^{}`$ orthonormal to the vector field $`t^\mathrm{\#}`$ dual to the torsion 1-form $`t`$. Hence, the torsion $`T`$ and therefore the connection $``$ descend to $`N^3`$.
In particular, $`\delta t=0`$ and therefore the QKT structure is strong. Q.E.D.
As a consequence of Theorem 5.2, we recover the following two results proved in in the setting of naturally reductive homogeneous 4-manifolds
###### Theorem 5.4
A (locally) homogeneous 4-dimensional QKT manifold is locally isometric to the Riemannian product $`\times N^3`$ of a real line and a naturally reductive homogeneous 3-manifold $`N^3`$.
###### Theorem 5.5
Let $`(M,g)`$ be a 4-dimensional compact homogeneous QKT manifold. Then the universal covering space $`\stackrel{~}{M}`$ of $`M`$ is isometric to the Riemannian product $`\times N^3`$ of a real line and the three dimensional space $`N^3`$ is one of the following
i) $`R^3,S^3,^3`$;
ii) isometric to one of the following Lie groups with a suitable left invariant metric:
1. $`SU(2)`$;
2. $`\stackrel{~}{SL(2,)}`$, the universal covering of $`SL(2,)`$;
3. the Heisenberg group.
Theorem 5.5 is based on the classification of 3-dimensional simply connected naturally reductive homogeneous spaces given in .
### 5.1 Einstein-like QKT 4-manifolds
It is well known that a 4n-dimensional ($`n>1)`$ QK manifold is Einstein and the Ricci forms satisfy $`\rho _\alpha (X,J_\alpha Y)=\rho _\beta (X,J_\beta Y)=\rho _\gamma (X,J_\gamma Y)=\lambda g(X,Y)`$, where $`\lambda `$ is a constant. However, the assumptions that these properties hold on a QKT manifold ($`n>1)`$ force the torsion to be zero and the QKT manifold is a QK manifold. Actually, we have already generalised this result proving that if $`\lambda `$ is not a constant the torsion has to be zero (see the proof of Theorem 4.3).
If the dimension is equal to 4 the situation is different. In this section we show that there exists a 4-dimensional (proper) QKT manifold satisfying similar curvature properties as those mentioned above.
We denote by $`K`$ the following (0,2) tensor
$`K(X,Y):=\rho _\alpha (X,J_\alpha Y)+\rho _\beta (X,J_\beta Y)+\rho _\gamma (X,J_\gamma Y)`$.
The tensor $`K`$ is independent of the chosen local almost complex structures $`(J_\alpha )`$ because of the following
###### Proposition 5.6
Let $`(M,g,(J_\alpha ),)`$ be a 4-dimensional QKT manifold. Then:
(5.66)
$$K=Ric+^gt\frac{\delta t}{2}g;$$
(5.67)
$$Skew(Ric)=\frac{1}{4}<dt,F_\alpha >F_\alpha +\frac{1}{2}(d^ct)_\alpha ,\alpha =1,2,3;$$
(5.68)
$$Ric^g=Sym(Ric)+\frac{1}{2}(|t|^2gtt),$$
where $`<,>`$ is the scalar product of tensors induced by $`g`$, $`Skew`$ (resp. $`Sym`$) denotes the skew-symmetric (resp. symmetric) part of a tensor.
In particular, the Ricci tensor is symmetric if and only if the torsion 1-form is closed.
Proof. We use (3.42). From (5.64) and (3.43), we obtain
(5.69)
$$D_\alpha (X,J_\alpha Y)=(_Xt)Y(_{J_\alpha Y}t)J_\alpha X,\alpha =1,2,3.$$
To compute $`B_\alpha `$ we need the following general identity
###### Lemma 5.7
On a 4-dimensional QKT manifold we have $`\genfrac{}{}{0pt}{}{\sigma }{XYZ}g(T(X,Y),T(Z,U))=0.`$
Proof of Lemma 5.7. Since $`\genfrac{}{}{0pt}{}{\sigma }{XYZ}g(T(X,Y),T(Z,U))`$ is a 4-form it is sufficient to check the equality for a basis of type $`\{X,J_\alpha X,J_\beta X,J_\gamma X\}`$. The last claim is obvious because of (5.60).
For each $`\alpha \{1,2,3\}`$, Lemma 5.7, (5.65) and (5.64) yield
(5.70)
$$B_\alpha (X,J_\alpha Y)=\underset{i=1}{\overset{4}{}}\genfrac{}{}{0pt}{}{\sigma }{XJ_\alpha Ye_i}(_XT)(J_\alpha Y,e_i,J_\alpha e_i)=(_Xt)Y+(_{J_\alpha Y}t)J_\alpha X\delta tg(X,Y).$$
Substituting (5.69), (5.70) into (3.42) and putting $`n=1`$, we derive (5.66) since $`^gt=t`$. Taking the trace in (5.65), we get $`_{i=1}^4(_{e_i}T)(e_i,X,Y)=\frac{1}{2}_{i=1}^4dt(e_i,J_\alpha e_i)F_\alpha (X,Y)+dt(J_\alpha X,J_\alpha Y),\alpha =1,2,3.`$ Then (5.67) follows from the last equality and Remark 3. The equation (5.68) is a direct consequence of (3.38) and (5.60). Q.E.D.
A $`4n`$-dimensional QKT manifold $`(M,g,(J_\alpha ),)`$ is said to be a Einstein QKT manifold if the symmetric part $`Sym(Ric)`$ of the Ricci tensor of $``$ is a scalar multiple of the metric $`g`$ i.e. $`Sym(Ric)=\frac{Scal}{4n}g,`$ where $`Scal=tr_gRic`$ is the scalar curvature of $``$.
We note that the scalar curvature $`Scal`$ of an Einstein QKT manifold may not be a constant.
We shall say that a 4-dimensional QKT manifold is sp(1)-Einstein if the symmetric part $`Sym(K)`$ of the tensor $`K`$ is a scalar multiple of the metric $`g`$ since the tensor $`K`$ is determined by the sp(1)-part of the curvature. On a sp(1)-Einstein QKT manifold $`Sym(K)=\frac{Scal^K}{4}g,`$ where $`Scal^K=tr_gK`$.
For a given QKT manifold with torsion 1-form $`t`$ we consider the corresponding Weyl structure $`^W`$, i.e. the unique torsion-free linear connection determined by the condition
(5.71)
$$^Wg=tg.$$
Conversely, in dimension 4, to a given Weyl structure $`^Wg=\psi g`$ we associate the QKT connection with torsion $`T=(\psi )`$. Note that a given Weyl structure on a conformal manifold (M,\[g\]) does not depend on the particularly chosen metric $`g[g]`$ but depends on the conformal class \[g\]. A Weyl structure is said to be Einstein-Weyl if the symmetric part $`Sym(Ric^W)`$ of its Ricci tensor is a scalar multiple of the metric $`g`$. Weyl structures and especially Einstein-Weyl structures have been much studied. For a nice overview of Einstein-Weyl geometry see . The next theorem shows the link between Einstein-Weyl geometry and sp(1)-Einstein QKT manifolds in dimension 4.
###### Theorem 5.8
Let $`(M,g,(J_\alpha ),)`$ be a 4-dimensional QKT manifold with torsion 1-form $`t`$. The following conditions are equivalent:
i) $`(M,g,(J_\alpha ),)`$ is a sp(1)-Einstein QKT manifold.
ii) The corresponding Weyl structure is an Einstein-Weyl structure.
Proof. The Weyl connection $`^W`$ determined by (5.71) is given explicitly by
$$_X^WY=_X^gY+\frac{1}{2}t(X)Y+\frac{1}{2}t(Y)X\frac{1}{2}g(X,Y)t^\mathrm{\#}.$$
The symmetric part of its Ricci tensor is equal to
(5.72)
$$Sym(Ric^W)=Ric^gSym(^gt)\frac{1}{2}(|t|^2gtt)+\frac{\delta t}{2}g.$$
Keeping in mind that $`^gt=t`$, we get from (5.66), (5.68) and (5.72) that
$`Sym(Ric^W)=Sym(K)`$. The theorem follows from the last equality. Q.E.D.
It is well known that on a 4-dimensional conformal manifold there exists a hypercomplex structure iff the conformal structure has anti-self-dual Weyl tensor (see also ). Every 4-dimensional hypercomplex manifold $`(M,g,(H_\alpha ))`$, i.e. (an oriented anti-self-dual 4-manifold) carries a unique HKT structure in view of the results in . Indeed, let $`\theta =\theta _\alpha =\theta _\beta =\theta _\gamma `$ be the common Lee form. The unique HKT structure is defined by $`=^g\frac{1}{2}\theta `$ (the uniqueness is a consequence of a general result in , see also ). The HKT structure on a 4-dimensional hypercomplex manifold is sp(1)-Einstein since the tensor $`K`$ vanishes. The corresponding Weyl structure to the given HKT structure on a 4-dimensional hyperhermitian manifold is the Obata connection , i.e. the unique torsion-free linear connection which preserves each of the three hypercomplex structures. As a consequence of Theorem 5.8, we recover the result in which states that the Obata connection of a hyper complex 4-manifold is Einstein-Weyl and the symmetric part of its Ricci tensor is zero.
Theorem 5.8 and (5.66) show that every Einstein-Weyl structure determined by (5.71) on a 4-dimensional conformal manifold whose vector field dual to the 1-form $`t`$ is Killing, induces an Einstein and sp(1)-Einstein QKT structure.
###### Corollary 5.9
Let $`(M,[g],^W)`$ be a compact 4-dimensional Einstein-Weyl manifold. Then the corresponding QKT structure to the Gauduchon metric of $`^W`$ is Einstein and sp(1)-Einstein.
Proof. On a compact Einstein-Weyl manifold the vector field dual to the Lee form of the Gauduchon metric is Killing by the result of Tod . Then the claim follows from Theorem 5.8 and (5.66). Q.E.D.
The Ricci tensor of a 4-dimensional QKT manifold is symmetric iff the torsion 1-form is closed by Proposition 5.6. Applying Theorem 3 in and using Theorem 5.8, we obtain
###### Corollary 5.10
Let $`(M,g,(J_\alpha ),)`$ be a 4-dimensional compact sp(1)-Einstein QKT manifold with symmetric Ricci tensor. Suppose that the torsion 1-form is not exact. Then the torsion 1-form corresponding to the Gauduchon metric $`g_G`$ of $`(M,g,(J_\alpha ),)`$ is parallel with respect to the Levi-Civita connection of $`g_G`$ and the universal cover of $`(M,g_G)`$ is isometric to $`\times S^3`$. In particular, the quaternionic bundle $`(J_\alpha )`$ admits hypercomplex structure.
A lot is known about Einstein-Weyl manifolds (see a nice survey ). There are many (compact) Einstein-Weyl 4-manifolds (e.g. $`S^2S^2`$). Among them there are (anti)-self-dual as well as non (anti)-self-dual. We mention here the Einstein-Weyl examples of Bianchi IX type metric . All these Einstein-Weyl 4-manifolds admit sp(1)-Einstein QKT structures by Theorem 5.8.
It is also known that there are obstructions to the existence of Einstein-Weyl structures on compact 4-manifold . If the manifold $`M`$ is finitely covered by $`T^2S^2`$ which cannot be Einstein-Weyl then $`M`$ does not admit Einstein-Weyl structure and therefore there are no sp(1)-Einstein structures on $`M`$.
Authors’ address:
Stefan Ivanov
University of Sofia, Faculty of Mathematics and Informatics, Department of Geometry,
5 James Bourchier blvd, 1126 Sofia, BULGARIA.
E-mail: ivanovsp@fmi.uni-sofia.bg
|
warning/0003/hep-ph0003036.html
|
ar5iv
|
text
|
# I Introduction
## I Introduction
Hints of new physics have recently been advocated through the neutrino sector. The recent observation by the Super-Kamiokande Collaboration of the atmospheric zenith angle dependent deficit has strengthened this conclusion. Also the long term puzzle of the solar neutrino deficit has been a strong demonstration on the existence of new physics. Recent results on neutrino oscillation have been also reported by the Liquid Scintillation Neutrino Detector (LSND) experiment . The observed anomalies in the neutrino data can be naturally understood in terms of massive neutrino oscillations.
The theoretical picture for neutrino oscillation poses a real challenge in understanding the form of the lepton mass matrix as derived from the neutrino oscillation data. A definite picture is obscure due to the large number of free parameters in the neutrino and charged lepton mass matrices. The phenomenological solution is not unique and clearly the need for more data and theoretical breakthrough is essential.
Neutrino oscillation can explain both the solar and atmospheric data in terms of three-generation neutrinos (ignoring the LSND results .) In the simplest explanation picture, the solar neutrino data can be understood in terms of $`\nu _e\nu _\mu `$ oscillation with a small mass splitting not to influence atmospheric data. On the other hand, atmospheric data can be explained in terms of $`\nu _\mu \nu _\tau `$ large mixing with a large mass splitting compared to the solar case . However, if we combine the LSND result with the solar and atmospheric data then we have to include at least an extra light neutrino. The full oscillation data requires the existence of three different scales of neutrino mass-squared differences. The different scales can be accommodated only if at least a light fourth neutrino exists. Such a light neutrino has to be sterile, i.e., to decouple from the low-energy observables as indicated by the low-energy experiments . Some recent phenomenological studies indicate that in the minimal scheme the dominant transition of solar neutrinos is due to $`\nu _e\nu _s`$ mixing, while the dominant transition of atmospheric neutrinos in long-baseline (LBL) experiments is due to $`\nu _\mu \nu _\tau `$ mixing .
The inclusion of a sterile neutrino still poses another theoretical challenge. Namely, to understand both the origin of the sterile fermion and the very low mass it has. The small mass is probably the most difficult issue in introducing such a particle. Any successful scenario has to explain the tiny mass of the sterile neutrino in a natural way. A possible scenario would be to generate the light mass through radiative corrections . Another interesting scenario would be to postulate that the extra neutrino is active at a relatively high-energy scale. At that scale the extra neutrino is massless as the assumed dynamics, due to some symmetry, forbids its mass generation. Once the high-energy symmetry is broken (probably in the TeV region), a Dirac mass can be generated while the neutrino decouples from the low-energy regime, i.e., becomes a sterile. Finally, by invoking the seesaw mechanism we can understand the highly suppressed Majorana mass of such a sterile neutrino.
In this work we consider the possibility of understanding the origin of a light sterile neutrino through the second scenario. A similar model has been discussed in Ref. , however, the model suffers from theoretical anomalies. Furthermore, an explicit formulation of the model is highly complicated. The model we discuss in this work does not suffer from theoretical drawbacks. It is based on the gauge nonuniversal symmetry $`SU(3)_c\times SU(2)_l\times SU(2)_h\times U(1)_Y`$ discussed extensively in Refs. . We refer to this model as the TopFlavor model which is anomaly free and phenomenologically well motivated. Several recent phenomenological studies have been published in the literature . To account for the existence of the sterile neutrino, we modify the standard fermion content by the inclusion of few extra fermions. The extra fermion spectrum does not appreciably affect the low-energy regime because of the heavy mass of the extra active fermions, as discussed later. Only one-neutral fermion emerges with a small mass while decouples from the low-energy regime which we then call the sterile neutrino.
The rest of this paper is organized as follows. In Sec. II, we briefly review the model. In Sec. III, we enlarge the fermion spectrum by introducing extra fermions and discuss the mechanism for generating the mass of the sterile and active fermions. Finally, in Sec. IV we provide some numerical discussion of the model.
## II Structure of the Model
The TopFlavor model is based on the gauge symmetry G= $`SU(3)_c\times SU(2)_l\times SU(2)_h\times U(1)_Y`$. In this model, the third generation of fermions (top quark $`t`$, bottom quark $`b`$, tau lepton $`\tau `$, and its neutrino $`\nu _\tau `$) is subjected to a new gauge interaction at the high energy scale, instead of the usual weak interaction advocated by the standard model (SM) of the electroweak interaction. On the contrary, the first and second generations only feel the weak interaction supposedly equivalent to the SM case. The new gauge dynamics is attributed to the $`SU(2)_h`$ symmetry under which the left-handed fermions of the third generation transform in the fundamental representation (doublets), while they remain to be singlets under the $`SU(2)_l`$ symmetry. On the other hand, the left-handed fermions of the first and second generation transform as doublets under the $`SU(2)_l`$ group and singlets under the $`SU(2)_h`$ group. The $`U(1)_Y`$ group is the SM hypercharge group. The right-handed fermions only transform under the $`U(1)_Y`$ group as assigned by the SM. Finally, the QCD interactions and the color symmetry $`SU(3)_c`$ are the same as that in the SM.
The spontaneous symmetry-breaking of the group G=$`SU(3)_c\times SU(2)_l\times SU(2)_h\times U(1)_Y`$ is accomplished by introducing the complex scalar fields $`\mathrm{\Sigma }`$, $`\mathrm{\Phi }_1`$, and $`\mathrm{\Phi }_2`$, where $`\mathrm{\Sigma }(1,2,2,0)`$, $`\mathrm{\Phi }_1(1,2,1,1)`$, and $`\mathrm{\Phi }_2(1,1,2,1)`$. For the $`\mathrm{\Sigma }`$ field we explicitly write
$$\mathrm{\Sigma }=\left(\begin{array}{cc}\pi _1^0& \pi _1^+\\ \pi _2^{}& \pi _2^0\end{array}\right),$$
(1)
where all scalar fields are complex. The group G is then broken in three different stages. The first stage of symmetry breaking is accomplished once the $`\mathrm{\Sigma }`$ field acquires a vacuum expectation value (vev) $`u`$, i.e., $`\mathrm{\Sigma }=\left(\begin{array}{cc}u& 0\\ 0& u\end{array}\right),`$ where $`u`$ is taken to be real. The form of $`\mathrm{\Sigma }`$ guarantees the breakdown of $`SU(2)_l\times SU(2)_hSU(2)`$. Therefore, the unbroken symmetry is essentially the SM gauge symmetry $`SU(3)_c\times SU(2)_w\times U(1)_Y`$, where $`SU(2)_w`$ is the usual SM weak group. At this stage, three of the gauge bosons acquire a mass of the order $`u`$, while the other gauge bosons remain massless. The phenomenology of the model imposes the constraint $`u\mathrm{\Gamma }>\mathrm{\hspace{0.17em}1}`$ TeV The second and third stage of symmetry breaking (the electroweak symmetry-breaking) is accomplished through the scalar fields $`\mathrm{\Phi }_1`$ and $`\mathrm{\Phi }_2`$ by acquiring their vacuum expectation values $`\mathrm{\Phi }_1=\left(\begin{array}{c}0\\ v_1\end{array}\right)`$, and $`\mathrm{\Phi }_2=\left(\begin{array}{c}0\\ v_2\end{array}\right)`$, respectively. The electroweak symmetry-breaking scale $`v`$ is defined $`v\sqrt{v_1^2+v_2^2}=246`$ GeV. Since the third generation of fermions is heavier than the first two generations, it is suggestive to conclude that $`v_2v_1`$. The surviving symmetry at low energy is $`SU(3)_c\times U(1)_{\mathrm{em}}`$ gauge symmetry. Because of this pattern of symmetry breaking, the gauge couplings $`g_l`$, $`g_h`$, and $`g_Y`$ of $`SU(2)_l`$, $`SU(2)_h`$, and $`U(1)_Y`$, respectively, are related to the $`U(1)_{\mathrm{em}}`$ gauge coupling $`e`$ by the relation $`1/e^2=1/g_l^2+1/g_h^2+1/g_Y^2`$ . We define
$$g_l=\frac{e}{\mathrm{sin}\theta \mathrm{cos}\varphi },g_h=\frac{e}{\mathrm{sin}\theta \mathrm{sin}\varphi },g_Y=\frac{e}{\mathrm{cos}\theta },$$
(2)
where $`\theta `$ is the usual weak mixing angle and $`\varphi `$ is a new parameter of the model.
For $`g_h>g_l`$ (equivalently $`\mathrm{tan}\varphi <1`$), we require $`g_h^24\pi `$ (which implies $`\mathrm{sin}^2\varphi g^2/(4\pi )0.04`$) so that the perturbation theory is valid. Similarly, for $`g_h<g_l`$, we require $`\mathrm{sin}^2\varphi 0.96`$. For simplicity, we focus on the region where $`xu^2/v^21`$, and ignore the corrections which are suppressed by higher powers of $`1/x`$. The light gauge boson masses are found to be $`M_{W^\pm }^2=M_Z^2\mathrm{cos}^2\theta =M_0^2(1+O(1/x)),`$ where $`M_0ev/2\mathrm{sin}\theta `$. While for the heavy gauge bosons $`W^\pm `$ and $`Z^{}`$, one finds
$$M_{W_{}^{}{}_{}{}^{\pm }}^2=M_Z^{}^2=M_0^2\left(\frac{x}{\mathrm{sin}^2\varphi \mathrm{cos}^2\varphi }+O(1)\right).$$
(3)
Up to this order, the heavy gauge bosons are degenerate in mass because they do not mix with the hypercharge gauge boson field. The SM fermions acquire their masses through their Yukawa interaction via the $`\mathrm{\Phi }_1`$ and $`\mathrm{\Phi }_2`$ scalar fields. For instance, the leptonic Yukawa sector is given by
$`_{\text{Yukawa}}^{\mathrm{}}`$ $`=`$ $`\overline{\mathrm{\Psi }_L}^1\mathrm{\Phi }_1\left[g_{11}^ee_R+g_{12}^e\mu _R+g_{13}^e\tau _R\right]+`$ (4)
$`\overline{\mathrm{\Psi }_L}^2\mathrm{\Phi }_1\left[g_{21}^ee_R+g_{22}^e\mu _R+g_{23}^e\tau _R\right]+`$
$`\overline{\mathrm{\Psi }_L}^3\mathrm{\Phi }_2\left[g_{31}^ee_R+g_{32}^e\mu _R+g_{33}^e\tau _R\right]+h.c.,`$
where
$$\mathrm{\Psi }_L^1=\left(\begin{array}{c}\nu _{eL}\\ e_L\end{array}\right),\mathrm{\Psi }_L^2=\left(\begin{array}{c}\nu _{\mu L}\\ \mu _L\end{array}\right),\mathrm{and}\mathrm{\Psi }_\mathrm{L}^3=\left(\begin{array}{c}\nu _{\tau \mathrm{L}}\\ \tau _\mathrm{L}\end{array}\right).$$
(5)
The phenomenology of the model has been studied extensively in Refs. <sup>2</sup><sup>2</sup>2Although some differences in the scalar sector exist among those references.. Comparisons with the Large Electron Positron (LEP) and other low-energy data have been investigated and constraints on the heavy gauge bosons mass are reported as $`M_W^{}\mathrm{\Gamma }>\mathrm{\hspace{0.17em}1}`$ TeV. The parameter $`xu^2/v^2`$ is constrained by LEP data to be larger than 20. Other low-energy data such as the $`\tau `$ life time imposes a higher constraint on $`x`$ for specific scenarios of lepton mixing. Flavor changing neutral current (FCNC) effects in the lepton and quark sectors have been explored and contributions to different processes have been calculated . Therefore, in this work we only concentrate on the leptonic mass matrix and refer the reader to Refs. for a detailed study of the phenomenology of the model.
## III Extra Fermions
To explain the existence of the light sterile neutrino in our scenario we enlarge the particle spectrum by the inclusion of extra fermions. We consider the minimal number of particles needed to account for the existence of the light sterile neutrino and without introducing anomalies in the structure of the model. Furthermore, consistency with low-energy data should be maintained and therefore the extra active fermions should decouple from the low-energy regime.
At the high-energy scale the gauge symmetry is assumed to be G=$`SU(3)_c\times SU(2)_l\times SU(2)_h\times U(1)_Y`$. At a lower scale, $`\mathrm{\Sigma }`$ a few TeV, the gauge symmetry is broken into the SM symmetry group $`H_1=SU(3)_c\times SU(2)_w\times U(1)_Y`$. At the electroweak scale the final stage of symmetry breaking occurs and the surviving symmetry group is $`H_2=SU(3)_c\times U(1)_{\mathrm{em}}`$. The fermion spectrum includes the standard three fermion generations with the transformation, under G, as explained in Sec. II.
We enlarge the fermion spectrum by the inclusion of three sets of extra fermions as follow:
* To explain the extremely light neutrino masses we invoke the seesaw mechanism. Therefore, four right-handed neutrinos are introduced, $`\nu _{eR}`$, $`\nu _{\mu R}`$, $`\nu _{\tau R}`$, and $`\nu _{sR}`$. The four right-handed neutrinos are singlets under G and are assumed to be Majorana fermions with masses of the order of the Grand Unified Theory (GUT) scale.
* We introduce a bi-doublet fermion field $`S_L`$ with the transformation $`S_L(1,2,2,0)`$. Explicitly, $`S_L=\nu _{sL}+S_L^a\tau ^a`$ transforms, under G, as
$$S_Lg_1S_Lg_2^{},$$
(6)
where $`g_1SU(2)_l`$ and $`g_2SU(2)_h`$. Once the symmetry group G is broken down into the symmetry group $`H_1`$, the field $`S_L`$ decomposes into two parts, with transformation under $`H_1`$ as $`(1,3,0)+(1,1,0)`$. The neutral field with the transformation $`(1,1,0)`$ corresponds to the sterile neutrino. The triplet field remains an active and thus must acquire a heavy Dirac mass in order to be consistent with the low-energy data . To prevent the uncontrolled Majorana mass we assume that $`S_L`$ carries a conserved quantum number $`z_L`$ at the high-energy scale. The new quantum number is due to a global $`U(1)`$ symmetry which causes no harmful anomaly to spoil the foundation of the model .
* We introduce another triplet fermion field $`S_R=S_R^a\tau ^a`$ with the transformation, under G, as $`(1,1,3,0)`$ just for the purpose of giving a Dirac mass to the active triplet field of $`S_L`$. Similar to $`S_L`$, the field $`S_R`$ is assumed to carry a quantum number $`z_R`$ to prevent the dangerous Majorana mass term. Hence, a Dirac mass for the extra active fermions is generated through the Yukawa interaction term
$$_a=\frac{g_a}{2}\mathrm{Tr}\left[\overline{\mathrm{S}_\mathrm{L}}\mathrm{\Sigma }\mathrm{S}_\mathrm{R}\right],$$
(7)
where $`g_a1`$ is a Yukawa coupling. Once the scalar field $`\mathrm{\Sigma }`$ acquires its vev $`u\mathrm{\Gamma }>`$ 1 TeV, a Dirac mass for the triplet fermions is generated, $`m_a=g_au\mathrm{\Gamma }>\mathrm{\hspace{0.17em}1}`$ TeV. In order for the Yukawa term in Eq. (7) to conserve the assumed global $`U(1)`$ symmetry, we require $`\mathrm{\Sigma }`$ to carry a quantum number $`z_0`$ such that $`z_L=z_R+z_0`$.
The sterile neutrino $`\nu _{sL}`$ acquires its Dirac mass through the Yukawa term
$$_s=\frac{g_s}{2}\mathrm{Tr}\left[\overline{\mathrm{S}_\mathrm{L}}\tau _2\mathrm{\Sigma }^{}\tau _2\right]\nu _{\mathrm{sR}},$$
(8)
where $`g_s1`$ is a Yukawa coupling. The Yukawa terms in Eqs. (7,8) conserve the new quantum number provided that we demand $`z_L=z_0=z_R/2`$. It is important to notice that the form of $`\mathrm{\Sigma }`$ is as given in Eq. (1). A simple choice would be
$$\mathrm{\Sigma }=\pi ^0+i\pi ^a\tau _a=\left(\begin{array}{cc}\pi ^0+i\pi ^3& i\sqrt{2}\pi ^+\\ i\sqrt{2}\pi ^{}& \pi ^0i\pi ^3\end{array}\right),$$
(9)
where $`\pi ^0`$ and $`\pi ^a`$ are taken to be real fields. However, for this particular choice $`\tau _2\mathrm{\Sigma }^{}\tau _2=\mathrm{\Sigma }`$ and therefore, the Yukawa terms in Eqs. (7,8) are not simultaneously invariant under the assumed global $`U(1)`$ symmetry. To conclude, the extra fields we introduce are the minimal number of fields required to account for the existence of a light sterile neutrino without spoiling the accuracy of the low-energy data, and also without introducing any theoretical anomalies into the model.
The full neutrino Yukawa interaction terms are given as
$`_{\text{Yukawa}}^\nu `$ $`=`$ $`\overline{\mathrm{\Psi }_L}^1\stackrel{~}{\mathrm{\Phi }}_1\left[g_{11}^\nu \nu _{eR}+g_{12}^\nu \nu _{\mu R}+g_{13}^\nu \nu _{\tau R}+g_{14}^\nu \nu _{sR}\right]+`$ (10)
$`\overline{\mathrm{\Psi }_L}^2\stackrel{~}{\mathrm{\Phi }}_1\left[g_{21}^\nu \nu _{eR}+g_{22}^\nu \nu _{\mu R}+g_{23}^\nu \nu _{\tau R}+g_{24}^\nu \nu _{sR}\right]+`$
$`\overline{\mathrm{\Psi }_L}^3\stackrel{~}{\mathrm{\Phi }}_2\left[g_{31}^\nu \nu _{eR}+g_{32}^\nu \nu _{\mu R}+g_{33}^\nu \nu _{\tau R}+g_{34}^\nu \nu _{sR}\right]+`$
$`{\displaystyle \frac{1}{2}}\mathrm{Tr}[\overline{S_L}\tau _2\mathrm{\Sigma }^{}\tau _2]\left[g_{41}^\nu \nu _{eR}+g_{42}^\nu \nu _{\mu R}+g_{43}^\nu \nu _{\tau R}+g_{44}^\nu \nu _{sR}\right]+h.c.,`$
where $`\stackrel{~}{\mathrm{\Phi }}_{1,2}i\tau _2\mathrm{\Phi }_{1,2}^{}`$. The Dirac mass matrix derived from Eq. (10) is written as
$$M_D=\left(\begin{array}{cccc}g_{11}^\nu v_1& g_{12}^\nu v_1& g_{13}^\nu v_1& g_{14}^\nu v_1\\ g_{21}^\nu v_1& g_{22}^\nu v_1& g_{23}^\nu v_1& g_{24}^\nu v_1\\ g_{31}^\nu v_2& g_{32}^\nu v_2& g_{33}^\nu v_2& g_{34}^\nu v_2\\ g_{41}^\nu u& g_{42}^\nu u& g_{43}^\nu u& g_{44}^\nu u\end{array}\right).$$
(11)
The right-handed neutrino Majorana mass matrix $`M_R`$ is assumed to have a common mass scale of the order of the GUT scale, $`M_X10^{15}`$ GeV. Therefore, the full neutrino mass matrix forms a $`8\times 8`$ matrix which can be written as
$$M_\nu =\left(\begin{array}{cc}0& M_D\\ M_D^T& M_R\end{array}\right).$$
(12)
By invoking the seesaw mechanism the left-handed neutrino Majorana mass matrix is then given as
$$M_L=M_DM_R^1M_D^T.$$
(13)
Due to the seesaw mechanism all elements of $`M_L`$ are highly suppressed by the GUT scale $`M_X`$ of the right-handed Majorana mass matrix $`M_R`$. Therefore, we can introduce the sterile neutrino with a natural mechanism for generating its light mass as required by the neutrino oscillation data. In the next section we give further discussion of the derived neutrino mass matrix.
## IV Discussion and Conclusions
The mass matrix in Eq. (11) is written in its most general form. A quantitative analysis is attainable only if the structure of the mass matrix is fully determined which requires further theoretical assumptions to be incorporated in the structure of the model. Nevertheless, the structure of the mass matrix already suggests an interesting behavior. There are three hierarchical energy scales in the mass matrix $`M_D`$, namely, $`uv_2v_1`$ which could be connected to the observed three hierarchical mass scales in the neutrino data, namely, $`\mathrm{\Delta }m_{\mathrm{LSND}}^2\mathrm{\Delta }m_{\mathrm{atm}.}^2\mathrm{\Delta }m_{\mathrm{solar}}^2`$. It is suggestive to conclude that
$`\mathrm{\Delta }m_{\mathrm{LSND}}^2\left({\displaystyle \frac{u^2}{M_X}}\right)^21\mathrm{eV}^2,`$
$`\mathrm{\Delta }m_{\mathrm{atm}.}^2\left({\displaystyle \frac{v_2^2}{M_X}}\right)^210^3\mathrm{eV}^2,`$
$`\mathrm{\Delta }m_{\mathrm{solar}}^2\left({\displaystyle \frac{v_1^2}{M_X}}\right)^210^5\mathrm{eV}^2,`$ (14)
as indicated by the neutrino oscillation data . In fact from the LEP data, we already know that $`u\mathrm{\Gamma }>v\sqrt{20}1.2`$ TeV . The observed mass scales can be obtained if we simply choose $`u1.2`$ TeV and $`v_2230`$ GeV. From which we conclude that $`v_175`$ GeV, $`M_X(110)\times 10^{15}`$ GeV, and
$`{\displaystyle \frac{\mathrm{\Delta }m_{\mathrm{LSND}}^2}{\mathrm{\Delta }m_{\mathrm{atm}.}^2}}{\displaystyle \frac{u^4}{v_2^4}}10^{+3},`$
$`{\displaystyle \frac{\mathrm{\Delta }m_{\mathrm{LSND}}^2}{\mathrm{\Delta }m_{\mathrm{solar}}^2}}{\displaystyle \frac{u^4}{v_1^4}}10^{+5},`$
$`{\displaystyle \frac{\mathrm{\Delta }m_{\mathrm{atm}.}^2}{\mathrm{\Delta }m_{\mathrm{solar}}^2}}{\displaystyle \frac{v_2^4}{v_1^4}}10^{+2}.`$ (15)
In the simplest scheme where oscillation data can be explained in terms of two flavor mixing, it has been argued that the dominant transition of solar neutrino is due to $`\nu _e\nu _s`$ mixing . Such a picture can hardly be satisfied by our model with the above choice of parameters and without the need for fine tuning. In the case of $`\nu _e\nu _s`$ mixing one can show that the effective $`2\times 2`$ mass matrix is given as
$$M_L=\frac{1}{M_X}\left(\begin{array}{cc}g_1v_1^2& g_2v_1u\\ g_2v_1u& g_3u^2\end{array}\right),$$
(16)
where the couplings $`g_{1,2,3}1`$. Therefore, one can show that the solar mass splitting is given as $`\mathrm{\Delta }m_{\mathrm{solar}}^2u^4/M_X^2`$ which is many orders of magnitude larger than the experimental fit . If we take $`M_X`$ to be of order $`10^{18}`$ GeV, which is close to the Planck scale, we get a result consistent with the experimental fit as shown in Table 1, where we consider the numerical values $`u=1200`$ GeV, $`M_X=10^{18}`$ GeV, and $`v_1=75`$ GeV. In Table 1, the numerical values of the Yukawa couplings as well as their solar neutrino solution are given.
However, such a solution is not favored as we can not explain the apparent hierarchy among the solar, atmospheric, and LSND data.
In conclusion we have provided a scenario based on the TopFlavor model to explain the existence of a light sterile neutrino. The scenario is anomaly free and phenomenologically compatible with all existing low-energy data. The scenario can also qualitatively explain the hierarchy in the observed mass scales of the neutrino oscillation data. Quantitative results are obtained for special cases.
## Acknowledgement
The authors would like to thank G. Senjanovic for useful discussion and comments. Also they would like to thank ICTP for the kind hospitality where some part of this work was done.
|
warning/0003/nucl-ex0003003.html
|
ar5iv
|
text
|
# Inhomogeneous isospin distribution in the reactions of 28Si + 112Sn and 124Sn at 30 and 50 MeV/nucleon.
##
|
warning/0003/math0003205.html
|
ar5iv
|
text
|
# Untitled Document
Exponentially decaying eigenvectors for certain
almost periodic operators
Norbert Riedel
Abstract. For every point $`\chi `$ in the spectrum of the operator
$$(h(\delta )\xi )=\xi _{n+1}+\xi _{n1}+\beta \left(\delta e^{2\pi \alpha ni}+\delta ^1e^{2\pi \alpha ni}\right)\xi _n$$
on $`\mathrm{}^2(\text{})`$ there exists a complex number $`x`$ of modulus one such that the equation
$$\xi _{n+1}+\xi _{n1}+\beta \left(x\delta e^{2\pi \alpha ni}+\overline{x}\delta ^1e^{2\pi \alpha ni}\right)\xi _n=\chi \xi _n$$
has a non-trivial solution satisfying the condition
$$\underset{|n|\mathrm{}}{\overline{\mathrm{lim}}}|\xi _n|^{1/|n|}\delta ^1\beta ^1$$
provided that $`\beta ,\delta >1`$ and $`\alpha `$ satisfies the diophantine condition
$$\underset{n\mathrm{}}{lim}|\mathrm{sin}\pi \alpha n|^{\frac{1}{n}}=1.$$
The parameters $`x\delta `$ and $`\chi `$ are in the range of analytic functions which are defined on a Riemann surface covering the resolvent set of the operator $`h(1)`$.
Introduction The spectral properties of almost periodic operators have been investigated extensively over the past 25 years, both in mathematics as well as in physics. A particularly intriguing problem in this area this author has been occupied with for some time is the question under what conditions such operators have point spectrum and how the prevalence of point spectrum affects the topological nature of their spectrum. In certain cases the second part of this question appears to be intimately related to a problem that deserves some consideration in its own right, namely, simplistically put, under what conditions is every point in the spectrum an eigenvalue?
In the sequel a family of almost periodic operators will be considered for which this question has a satisfactory answer. The operators to be considered are complex perturbations of bounded self-adjoint operators, known as almost Mathieu operators or Harper’s operators. The approach chosen is $`C^{}`$-algebraic. The key to the proofs are certainautomorphisms $`\rho _\beta `$ of the irrational rotation $`C^{}`$-algebra $`𝒜_\alpha `$ associated with an irrational number $`\alpha `$. If $`u`$ and $`v`$ are unitary generators of $`𝒜_\alpha `$ satisfying the defining relation $`uv=e^{2\pi \alpha i}vu`$, the operators of interest are of the form
$$h(\delta )=u+u^{}+\beta (\delta v+\delta ^1v^{})$$
where $`\beta >1`$ is a fixed constant and $`|\delta |>1`$. The said automorphism $`\rho _\beta `$ flips a defining parameter when applied to a slightly enlarged family of operators. This can then be used to generate exponentially decaying eigenvectors for the operators $`h(\delta )`$ represented on the Hilbert space $`\mathrm{}^2(\text{})`$, provided $`\alpha `$ satisfies a suitable diophantine condition. From a dynamical systems point of view, the automorphism $`\rho _\beta `$ is related to a skew translation, extending the irrational rotation underlying the dynamics of the operators in question, in a sense to be made precise below. It is noteworthy that as $`\beta `$ approaches $`1`$, $`\rho _\beta `$ approaches an automorphism of period $`4`$, a so-called “Fourier transform”. This shows, in particular, that the extension of the dynamical systems picture, which is so vital for the case $`\beta >1`$, is no longer available in the case $`\beta =1`$.
Introducing the automorphism $`\rho _\beta `$ and presenting a brief discussion of the extended dynamical systems picture will be taken up in the first paragraph. In the second paragraph the existence of exponentially decaying eigenvectors for points in the spectrum of the operators $`h(\delta )`$ will be proved. It will then be shown that a far more specific formulation of the eigenvalue problem for the operators $`h(\delta )`$ can be obtained through parametrization in a suitable Riemann surface covering the resolvent set $``$ of the operator $`h(1)`$. More specifically, since the spectrum of $`h(1)`$ turns out to be a regular compactum in the sense of potential theory, and since the spectra of the operators $`h(\delta )`$ are exactly the level curves of the corresponding conductor potential, there exists a Riemann surface $`\stackrel{~}{}`$ covering $``$ and an analytic function $`G`$ which maps $`\stackrel{~}{}`$ onto the complement of the closed unit disk, such that
$$h(G(z))\xi =p(z)\xi $$
has an exponentially decaying solution $`\xi `$ for every $`z`$ in $`\stackrel{~}{}`$. Here $`p`$ denotes the canonical mapping from $`\stackrel{~}{}`$ onto $``$. As $`z`$ ranges over $`\stackrel{~}{}`$, $`\xi `$ ranges over all possible eigenvectors for the operators $`h(\delta )`$. Due to the basic $`K`$-theory for the $`C^{}`$-algebra $`𝒜_\alpha `$, one can see that the group of covering transformations of $`\stackrel{~}{}`$ over $``$ is infinite cyclic. Moreover, translation by one of the two generators of this group, $`\omega `$ say, corresponds to shifting the eigenvector $`\xi `$.
The automorphism $`\rho _\beta `$ gives rise to an eigenvalue problem in its own right which is intimately interconnected with the one expounded above. The eigenvalues are given through an analytic function $`\mathrm{\Gamma }`$ on $`\stackrel{~}{}`$ which has the property
$$\mathrm{\Gamma }(\omega (z))=G(z)^2\mathrm{\Gamma }(z).$$
In paragraph 3 it will be shown that the two eigenvalue problems are essentially equivalent. To this end, the latter eigenvalue problem will be transformed into a question regarding the kernel of certain Fredholm operators of index zero. The problem then boils down to the question whether these kernels are one-dimensional. With the aid of analytic perturbation theory, it will be shown that this is indeed the case.
Finally, in paragraph 4, the case $`|\delta |=1`$ will be discussed. Since $`h=h(1)`$ is a fixed-point of $`\rho _\beta `$, no exponentially decaying eigenvectors can be generated along the lines this was possible for the case $`|\delta |>1`$. Nevertheless, the extended dynamical systems picture shows that the eigenvalue problem for $`h`$ is intimately related to similar questions about Schrödinger type difference operators with unbounded potentials, such as
$$n\mathrm{tan}\pi (\alpha n^2+2\theta n+\nu ).$$
Even though these operators are distinctly non almost periodic, they are related to, and in fact extensions of, a family of operators which was a focus of research in the early 1980’s (let $`\alpha =0`$ and let $`\theta `$ be irrational). In the restricted case ($`\alpha =0`$) it can be shown that the said operators have pure point spectrum. This is being accomplished by relating these operators to certain bounded operators which can be diagonalized by solving a specific cocycle equation, and then by translating this information back to the original context. In the extended case ($`\alpha 0`$) there still exist those related bounded operators, which are actually derived from the automorphism $`\rho _\beta `$, and which are the point of departure, rather than an auxiliary device, as it happens to be the case when $`\alpha =0`$. But it is not possible anymore to diagonalize these related bounded operators, due to the complications brought about by switching from an irrational rotation to a skew translation extending it.
1. The automorphism $`\rho _\beta `$ For an irrational number $`\alpha `$ we consider the $`C^{}`$-algebra $`𝒜=𝒜_\alpha `$ generated by two unitaries $`u`$ and $`v`$ satisfying the relation $`uv=\lambda ^2vu`$, where $`\lambda =e^{\pi \alpha i}`$. For $`\beta 1`$, let
$$\begin{array}{cc}\hfill \rho _\beta (u)& =vuv(uv+\beta )^1(v^{}u^{}+\beta )\hfill \\ \hfill \rho _\beta (v)& =v(uv+\beta )^1(v^{}u^{}+\beta ).\hfill \end{array}$$
Then $`\rho _\beta (u)`$ and $`\rho _\beta (v)`$ are unitaries satisfying again the relation
$$\rho _\beta (u)\rho _\beta (v)=\lambda ^2\rho _\beta (v)\rho _\beta (u).$$
Therefore, $`\rho _\beta `$ extends to an automorphism of $`𝒜`$ which we will also denote by $`\rho _\beta `$. The significance of this automorphism for what is to follow rests with the identities
$$\rho _\beta (u+\beta v)=u^{}+\beta v,\rho _\beta (u^{}+\beta v^{})=u+\beta v^{}.$$
$`(1.1)`$
Let $`GL(2,\text{})`$ be the group of $`2\times 2`$ matrices with integer entries and a determinant of modulus 1. The assignment
$$\begin{array}{cc}\hfill w_{m,n}& w_{pq}\hfill \\ \hfill \left(\genfrac{}{}{0pt}{}{p}{q}\right)& =A\left(\genfrac{}{}{0pt}{}{m}{n}\right)\hfill \end{array}\}AGL(2,\text{}),$$
where $`w_{pq}=\lambda ^{pq}u^pv^q`$, is known to define a linear isometry of $`𝒜`$. This isometry is an automorphism if $`\mathrm{det}A=1`$, and it is an antiautomorphism if $`\mathrm{det}A=1`$. The mapping assigning to every $`AGL(2,\text{})`$ the corresponding isometry is a faithful homomorphism from the group $`GL(2,\text{})`$ into the group of isometries of $`𝒜`$. For convenience we will denote matrices in $`GL(2,\text{})`$ and their corresponding isometries by the same symbol. For instance $`\left(\genfrac{}{}{0pt}{}{01}{10}\right)`$ represents the antiautomorphism $`uv`$, $`vu`$.
Returning to the automorphism $`\rho _\beta `$ we are now going to list a number of useful identities.
$$\begin{array}{cc}& \rho _\beta =\left(\genfrac{}{}{0pt}{}{10}{11}\right)\rho _\beta ^{(0)}\left(\genfrac{}{}{0pt}{}{10}{11}\right),\text{where}\hfill \\ & \rho _\beta ^{(0)}(u)=u\hfill \\ & \rho _\beta ^{(0)}(v)=v(\lambda u+\beta )^1(\overline{\lambda }u^{}+\beta )\hfill \end{array}\}$$
$`(1.2)`$
$$\rho _{\beta ^1}=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)\rho _\beta ^1\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)$$
$`(1.3)`$
This identity says essentially that $`\rho _{\beta ^1}`$ and $`\rho _\beta ^1\left(\genfrac{}{}{0pt}{}{10}{01}\right)`$ are conjugates of each other via a “Fourier transform” (also known as “duality”).
$$\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)\rho _\beta =\rho _\beta \left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)$$
$`(1.4)`$
If $`\beta `$ approaches $`1`$ then $`\rho _\beta ^{(0)}`$ approaches $`\left(\genfrac{}{}{0pt}{}{11}{01}\right)`$ on $`w_{pq}`$, hence $`\rho _\beta `$ approaches
$$\left(\begin{array}{cc}1& 0\\ 1& 1\end{array}\right)\left(\begin{array}{cc}1& 1\\ 0& 1\end{array}\right)\left(\begin{array}{cc}1& 0\\ 1& 1\end{array}\right)=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)$$
on $`w_{pq}`$. Thus $`\rho _\beta `$ approaches a “Fourier transform”. Throughout in the subsequent discussion we will limit our attention to irrational numbers $`\alpha `$ only which satisfy a diophantine condition.
$$\underset{n\mathrm{}}{lim}|\mathrm{sin}\pi \alpha n|^1=1$$
$`(1.5)`$
For such numbers $`\alpha `$ and $`\beta >1`$ the automorphism $`\rho _\beta ^{(0)}`$ becomes an inner automorphism. Namely,
$$\{\begin{array}{cc}& \rho _\beta ^{(0)}(a)=e^{ig(u)}ae^{ig(u)},a𝒜;\text{or}\rho _\beta ^{(0)}=Ad(e^{ig(u)}),\hfill \\ & \text{where}\hfill \\ & g(z)=\underset{n=1}{\overset{\mathrm{}}{}}\frac{(1)^n}{n}(\mathrm{sin}\pi \alpha n)^1\beta ^n(z^n+z^n).\hfill \end{array}$$
$`(1.6)`$
We turn now to a dynamical systems interpretation of the automorphisms $`\rho _\beta `$ in the framework of $`C^{}`$-algebras, under the assumption that (1.5) holds. Consider the crossed product $``$ of $`𝒜`$ by the automorphism $`\left(\genfrac{}{}{0pt}{}{10}{11}\right)`$
$$=𝒜_{\left({\scriptscriptstyle \genfrac{}{}{0pt}{}{10}{11}}\right)}\text{}.$$
This $`C^{}`$-algebra is generated by three unitaries $`u`$, $`v`$ and $`w`$ satisfying the defining relations
$$\begin{array}{cc}\hfill w^{}uw& =\overline{\lambda }uv\hfill \\ \hfill vw& =wv\hfill \\ \hfill uv& =\lambda ^2vu\hfill \end{array}\}$$
Since $`\rho _\beta =\left(\genfrac{}{}{0pt}{}{10}{11}\right)Ad(e^{ig(u)})\left(\genfrac{}{}{0pt}{}{10}{11}\right)`$, $`\rho _\beta `$ extends to an inner automorphism of $``$. So $``$ provides a natural framework where all our manipulations so far take place. Manipulating the first of those relations we get
$$uwu^{}=\lambda vw.$$
This, in conjunction with the other two relations, suggests that $``$ can also be realized as the crossed-product of the $`C^{}`$-algebra of continuous functions on the two dimensional torus $`C(\text{𝕋}^2)`$ by a skew translation followed by a translation. More specifically, let
$$\varphi (x,y)=\lambda (\lambda x,xy),x,y\text{𝕋}^2.$$
In dynamical systems theory it is a well known fact that $`\varphi `$ is a uniquely ergodic homeomorphism of $`\text{𝕋}^2`$, the unique invariant probability measure being the Haar measure on the compact group $`\text{𝕋}^2`$, which has of course full support. It is then a well known fact in $`C^{}`$-algebra theory, that the crossed-product
$$\stackrel{~}{}=C(\text{𝕋}^2)_\varphi \text{}$$
is a simple $`C^{}`$-algebra (i.e., it has no non-trivial ideals) with a unique tracial stateextending the Haar measure on $`C(\text{𝕋}^2)`$. Letting $`\stackrel{~}{u}`$ be the unitary in $`\stackrel{~}{}`$ corresponding to $`\varphi `$, $`\stackrel{~}{v}`$ the projection from $`\text{𝕋}^2`$ onto the first component, and $`\stackrel{~}{w}`$ the projection from $`\text{𝕋}^2`$ onto the second component, then it is easily seen that $`\stackrel{~}{u}`$, $`\stackrel{~}{v}`$ and $`\stackrel{~}{w}`$ satisfy the same three relations stated above for $`u`$, $`v`$ and $`w`$. Using this information it is not difficult to see that the assignments $`\stackrel{~}{u}u`$, $`\stackrel{~}{v}v`$, $`\stackrel{~}{w}w`$ extend to an isomorphism from $`\stackrel{~}{}`$ onto $``$.
2. Spectrum and point spectrum for a family of non self-adjoint almost periodic operators
In this paragraph we will employ the automorphism $`\rho _\beta `$ to investigate the spectrum of the operators
$$h(\delta )=u+u^{}+\beta (\delta v+\delta ^1v^{})$$
for $`|\delta |>1`$, provided that $`\beta >1`$ and $`\alpha `$ satisfies the property (1.5). To this end we consider the extended family
$$h_\gamma (\delta )=\gamma u+\gamma ^1u^{}+\beta (\delta v+\delta ^1v^{}),$$
where $`\beta ^1\delta ^1<|\gamma |<\beta \delta `$.
The assignments
$$\{\begin{array}{cc}\hfill (u\xi )_n& =\xi _{n1}\hfill \\ \hfill (v\xi )_n& =\overline{\lambda }^{2n}\xi _n\hfill \\ \hfill (w\xi )_n& =\lambda ^{n^2}\xi _n\hfill \end{array}$$
$`(2.1)`$
define linear operators on the vector space $`\text{}^{\mathrm{}}`$ of (two-sided) sequences $`\xi `$. When restricted to square summable sequences, these assignments extend to a (faithful) representation of the $`C^{}`$-algebra $``$ on the Hilbert space $`\mathrm{}^2(\text{})`$ introduced in paragraph 1. For $`\beta ^1<|\gamma |<\beta `$, let
$$k_\gamma =we^{ig(\gamma u)}w,k=k_1.$$
Then (1.1), (1.2) and (1.6) yield
$$h_{\gamma \delta ^1}(\delta )k_\gamma =k_\gamma h_{\gamma \delta }(\delta ).$$
$`(2.2)`$
For any complex number $`x`$, let
$$(D_x\xi )_n=x^n\xi ,\xi \text{}^{\mathrm{}}.$$
2.1 Lemma. Let $`|\delta |>1`$; $`\beta ^1<|\gamma |<1`$ or $`1<|\gamma |<\beta `$.
$`(+)`$ Suppose that
$$h_{\gamma \delta }(\delta )\eta =z\eta \text{for some}z\text{},\eta \mathrm{}^{\mathrm{}}(\text{}).$$
Then
$$h(\delta )\xi =z\xi ,\underset{n\mathrm{}}{\overline{\mathrm{lim}}}|\xi _n|^{1/|n|}|\gamma |\delta ^1,\underset{n\mathrm{}}{\overline{\mathrm{lim}}}|\xi _n|^{1/n}|\gamma |^1\delta ^1$$
where $`\xi =D_{\gamma ^1\delta ^1}\eta `$.
$`()`$ Suppose that
$$h_{\gamma \delta ^1}(\delta )\eta =z\eta \text{for some}z\text{},\eta \mathrm{}^{\mathrm{}}(\text{}).$$
Then
$$h(\delta )\xi =z\xi ,\underset{n\mathrm{}}{\overline{\mathrm{lim}}}|\xi _n|^{1/|n|}|\gamma |^1\delta ^1,\underset{n\mathrm{}}{\overline{\mathrm{lim}}}|\xi _n|^{1/n}|\gamma |\delta ^1,$$
where $`\xi =D_{\gamma ^1\delta }\eta `$.
Proof: We will deal with the case $`(+)`$ only. The case $`()`$ can be handled in a similar fashion.
The identity (2.2) yields
$$h_{\gamma \delta ^1}k_\gamma \eta =zk_\gamma \eta ,\text{where}k_\gamma \eta \mathrm{}^{\mathrm{}}(\text{}).$$
Let $`\stackrel{~}{\xi }=D_{\gamma ^1\delta }\eta `$. Then
$$h(\delta )\stackrel{~}{\xi }=z\stackrel{~}{\xi },\underset{n\mathrm{}}{\overline{\mathrm{lim}}}|\stackrel{~}{\xi }_n|^{1/|n|}|\gamma |\delta ^1,\underset{n\mathrm{}}{\overline{\mathrm{lim}}}|\stackrel{~}{\xi }_n|^{1/n}|\gamma |^1\delta .$$
On the other hand,
$$h(\delta )\xi =z\xi ,\underset{n\mathrm{}}{\overline{\mathrm{lim}}}|\stackrel{~}{\xi }_n|^{1/|n|}|\gamma |\delta ,\underset{n\mathrm{}}{\overline{\mathrm{lim}}}|\xi _n|^{1/n}|\gamma |^1\delta ^1.$$
Therefore
$$\underset{n\mathrm{}}{\overline{\mathrm{lim}}}\left|\xi _n\stackrel{~}{\xi }_{n+1}\right|^{1/|n|}|\gamma |^2,\underset{n\mathrm{}}{\overline{\mathrm{lim}}}\left|\xi _n\stackrel{~}{\xi }_{n+1}\right||\gamma |^2$$
and
$$\underset{n\mathrm{}}{\overline{\mathrm{lim}}}\left|\xi _{n+1}\stackrel{~}{\xi }_n\right|^{1/|n|}|\gamma |^2,\underset{n\mathrm{}}{\overline{\mathrm{lim}}}\left|\xi _{n+1}\stackrel{~}{\xi }_n\right||\gamma |^2.$$
Since $`|\gamma |1`$ by assumption, this entails
$$\underset{n\mathrm{}}{\overline{\mathrm{lim}}}\left(\xi _n\stackrel{~}{\xi }_{n+1}\xi _{n+1}\stackrel{~}{\xi }_n\right)=0$$
or
$$\underset{n\mathrm{}}{\overline{\mathrm{lim}}}\left(\xi _n\stackrel{~}{\xi }_{n+1}\xi _{n+1}\stackrel{~}{\xi }_n\right)=0.$$
Either way, it follows that $`\xi `$ and $`\stackrel{~}{\xi }`$ are linearly dependent. For, if this were not the case, the expression following the $`\overline{\mathrm{lim}}`$ in the last two identities would have to be constant and non-zero for all $`n\text{}`$. Therefore, observing that $`\delta >1`$ by assumption, we conclude
$$\underset{n\mathrm{}}{lim}|\xi _n|^{1/|n|}|\gamma |\delta ^1,\underset{n\mathrm{}}{lim}|\xi _n|^{1/n}|\gamma |^1\delta ^1,$$
as claimed.
2.2 Lemma. Let $``$ be a set of real numbers and let $`b>0`$, $`d>0`$. Suppose that $``$ has a non-empty intersection with at least one of the two open intervals $`(db,d)`$ or $`(d,d+b)`$. Suppose in addition that $``$ has the following properties:
$`(+)`$$`0<|t|<b`$ and $`t+d`$ implies $`(td,t+d)`$.
$`()`$$`0<|t|<b`$ and $`td`$ implies $`(td,t+d)`$.
Then $`(db,d+b)`$.
Proof: Consider the case that $`(db,d)\mathrm{}`$. Then there exists a $`c_0(b,0)`$ such that $`c_0+d`$.
Suppose first that $`bd`$.
Let $`s>db`$, but close to $`db`$. Since $`(+)`$ ensures that $`(c_0d,c_0+d)`$, we can find $`c_1(b,0)`$ such that $`c_1+d`$ and $`c_1d<s`$. Again, $`(+)`$ ensures that $`(c_1d,c_1+d)`$. Now let $`t<d+b`$ but close to $`d+b`$. Then we can find $`c_2(0,b)`$ such that $`c_2d(c_1d,c_1+d)`$ and $`c_1+d>t`$. Then $`()`$ ensures that $`(c_2d,c_2+d)`$. By construction $`(s,t)(c_1d,c_1+d)(c_2d,c_2+d)`$. Since $`s`$ and $`t`$ can be chosen arbitrarily close to $`db`$ and $`d+b`$, respectively, we conclude that $`(db,d+b)`$.
Now suppose that $`b>d`$.
Let $`s>db`$, but close to $`db`$. Through induction we can generate a (possibly empty) chain $`c_1>\mathrm{}>c_n`$ such that $`c_0>c_1`$ if $`n>0`$; $`c_k(b,0)`$, $`c_k(c_{k1}d,c_{k1}+d)`$, for $`1kn`$, and $`c_nd<b`$. Repeated applications of $`(+)`$ show that $`(c_kd,c_k+d)`$ for $`k=0,\mathrm{},n`$. Now choose $`c_{n+1}(b,c_n)`$ such that $`c_{n+1}d<s`$. Once again, $`(+)`$ ensures that $`(c_{n+1}d,c_{n+1}+d)`$. By construction we have $`s(c_{n+1}d,c_0+d)`$.
Next, let $`t<d+b`$ but close to $`d+b`$. Using $`()`$ instead of $`(+)`$, we can construct in the same fashion a chain $`c_0<\stackrel{~}{c}_1<\mathrm{}<\stackrel{~}{c}_{m+1}`$ such that $`t(c_0d,c_{m+1}+d)`$. This means $`(s,t)`$, and again we conclude that $`(db,d+b)`$.
The case that $`(d,d+b)\mathrm{}`$ can be handled in a similar fashion.
2.3 Lemma. Let $`|\delta |>1`$; $`\beta ^1<|\gamma _0|<1`$ or $`1<|\gamma _0|<\beta `$. Suppose that
$$h_{\gamma _0\delta ^{\pm 1}}(\delta )\eta =z\eta \text{for some}z\text{},\eta \mathrm{}^{\mathrm{}}(\text{}).$$
Then
$$h(\delta )\xi =z\xi ,\underset{|n|\mathrm{}}{\overline{\mathrm{lim}}}|\xi _n|^{1/|n|}\beta ^1|\delta |^1,$$
where $`\xi =D_{\gamma _0^1\delta 1}\eta `$.
Proof: Let
$$\begin{array}{c}\text{ }b=\mathrm{log}\beta ,d=\mathrm{log}|\delta |,c_0=\mathrm{log}|\gamma _0|,\text{ }\hfill \\ \text{ }=\left\{\mathrm{log}|\gamma |\right|D_\gamma \xi \mathrm{}^{\mathrm{}}(\text{})\}.\text{ }\hfill \end{array}$$
Then $`(+)`$ and $`()`$ in Lemma 2.1 translate into the namesake properties of Lemma 2.2, which then yields the desired conclusion.
Let $`Sp_0(h_\gamma (\delta ))`$ be the spectrum of $`h_\gamma (\delta )`$ considered as a bounded linear operator on the Banach space $`c_0(\text{})`$ of bounded two-sided sequences which vanish at infinity.
2.4 Lemma. For every $`zSp_0(h_\gamma (\delta ))`$ there exists $`x\text{𝕋}`$ and $`\eta \mathrm{}^{\mathrm{}}(\text{})\backslash \{0\}`$ such that
$$h_\gamma (x\delta )\xi =z\xi .$$
Proof: We choose an approximate eigenvector for $`z`$ in $`c_0(\text{})`$, $`\eta ^{(1)},\eta ^{(2)},\mathrm{}`$ in $`c_0(\text{})`$, $`\eta ^{(m)}_{\mathrm{}}=1`$,
$$\underset{m\mathrm{}}{lim}(h_\gamma (\delta )z))\eta ^{(m)}_{\mathrm{}}=0.$$
For each $`m`$ there is a $`j_m`$ such that $`\left|\eta _{j_m}^{(m)}\right|\frac{1}{2}`$. Let $`\xi ^{(m)}=u^{j_m}\eta ^{(m)}`$. Then
$$\underset{m\mathrm{}}{lim}(h_\gamma (x_m\delta )z)\xi ^{(m)}_{\mathrm{}}=0,$$
where $`x_m=\overline{\lambda }^{2j_m}`$. Also,
$$\left|\xi _0^{(m)}\right|\frac{1}{2}.$$
Since the unit ball of $`\mathrm{}^{\mathrm{}}(\text{})`$ is weakly compact and metrizable, there exists a subsequence of $`\xi ^{(1)},\xi ^{(2)},\mathrm{}`$ which converges weakly to some $`\xi \mathrm{}^{\mathrm{}}(\text{})`$. Moreover, we can arrange for the corresponding subsequence of $`x_1,x_2,\mathrm{}`$ to converge to some $`x\text{𝕋}`$. Since $`|\xi _0|\frac{1}{2}`$, $`\xi `$ is non-zero. By construction
$$h_\gamma (x\delta )\xi =z\xi $$
2.5 Lemma. $`Sp_0(h_\gamma (\delta ))=Sp_0(h_{|\gamma |}(|\delta |))`$.
Proof: Since $`h_\gamma (\delta )=D_{\gamma /|\gamma |}h_{|\gamma |}(\delta )D_{\gamma /|\gamma |}^1`$, we have
$$Sp_0\left(h_\gamma (\delta )\right)=Sp_0\left(h_{|\gamma |}(\delta )\right).$$
Now let $`zSp_0(h_\gamma (\delta ))`$ and $`x\text{𝕋}`$. Then there exists an approximate eigenvector for $`z`$ in $`c_0(\text{})`$:
$$\eta ^{(1)},\eta ^{(2)},\mathrm{};\eta ^{(m)}_{\mathrm{}}=1,$$
$$\underset{m\mathrm{}}{lim}(h_\gamma (\delta )z)\eta ^{(m)}_{\mathrm{}}=0.$$
For each $`m`$ choose $`j_m\text{}`$ such that $`\lambda ^{2j_m}x`$. Then $`u^{j_1}\eta ^{(1)},u^{j_2}\eta ^{(2)},\mathrm{}`$ is seen to be an approximate eigenvector of $`h_\gamma (x\delta )`$ for $`z`$, which entails $`zSp_0(h_\gamma (x\delta ))`$.
2.6 Lemma. Let $`|\delta |>1`$, $`|\delta |\beta ^1<|\gamma |<|\delta |\beta `$. Then
$$Sp_0\left(h_\gamma (\delta )\right)=Sp_0\left(h_{\gamma ^1}(\delta )\right).$$
Proof: Let $`zSp_0(h_\gamma (\delta ))`$. By Lemma 2.4 there exist $`x\text{𝕋}`$, $`\xi \mathrm{}^{\mathrm{}}(\text{})\backslash \{0\}`$ such that
$$h_\gamma (x\delta )\xi =z\xi .$$
Lemma 2.3 implies that $`\xi `$ and $`\eta =D_{\gamma ^2}\xi `$ are in $`c_0(\text{})`$. But
$$h_{\gamma ^1}(x\delta )\eta =z\eta .$$
Hence
$$zSp_0(h_{\gamma ^1}(x\delta )),$$
which, by Lemma 2.5, entails
$$zSp_0(h_{\gamma ^1}(\delta )).$$
so we have
$$Sp_0(h_\gamma (\delta ))Sp_0(h_{\gamma ^1}(\delta )).$$
The opposite inclusion is shown in exactly the same way.
2.7 Lemma. Let $`|\delta |>1`$, $`|\delta |\beta ^1<|\gamma |<|\delta |\beta `$. Then
$$Sp(h_\gamma (\delta ))=Sp_0(h_\gamma (\delta )),$$
where $`Sp(h_\gamma (\delta ))`$ denotes the spectrum of $`h_\gamma (\delta )`$ considered as an operator on the Hilbert space $`\mathrm{}^2(\text{})`$.
Proof: First we are going to show that $`Sp_0(h_\gamma (\delta ))Sp(h_\gamma (\delta ))`$. Let $`z\text{}\backslash Sp(h_\gamma (\delta ))`$, and let
$$\left(h_\gamma (\delta )z\right)^1=\underset{p,q=\mathrm{}}{\overset{\mathrm{}}{}}c_{pq}(z)w_{pq}$$
be the “Fourier expansion” of $`(h_\gamma (\delta )z)^1`$. Since this expansion decays exponentially in $`p`$ and $`q`$, it defines a bounded linear operator on $`c_0(\text{})`$, which is an inverse of $`h_\gamma (\delta )z`$ on $`c_0(\text{})`$. Hence $`z\text{}\backslash Sp_0(h(\delta ))`$. So, $`Sp_0(h_\gamma (\delta ))Sp(h_\gamma (\delta ))`$.
To establish the opposite inclusion, let $`z\text{}\backslash Sp_0(h_\gamma (\delta ))`$. Then $`h_\gamma (\delta )z`$ has an inverse $`S`$ on $`c_0(\text{})`$. Since $`Sp_0(h_\gamma (\delta ))=Sp_0(h_{\gamma ^1}(\delta ))`$ by Lemma 2.6, $`h_{\gamma ^1}(\delta )z`$ also has an inverse on $`c_0(\text{})`$, which we denote by $`T`$. Let $`(h_{\gamma ^1}(\delta )z)^t`$ and $`T^t`$ be the transposed operators of $`h_{\gamma ^1}(\delta )z`$ and $`T`$, respectively, on the dual Banach space of $`c_0(\text{})`$, which happens to be $`\mathrm{}^{}(\text{})`$. Since $`(h_{\gamma ^1}(\delta )z)^t`$ equals the restriction of $`h_\gamma (\delta )z`$ on $`\mathrm{}^1(\text{})c_0(\text{})`$ and $`T^t`$ is an inverse of $`(h_{\gamma ^1}(\delta )z)^t`$ on $`\mathrm{}^1(\text{})`$, $`T^t`$ must be the restriction of $`S`$ on $`\mathrm{}^1(\text{})`$. Now consider the double transposition $`(h_\gamma (\delta )z)^{tt}`$ and $`S^{tt}`$ of $`h_\gamma (\delta )z`$ and $`S`$, respectively, on the dual Banach space of $`\mathrm{}^1(\text{})`$, which happens to be $`\mathrm{}^{\mathrm{}}(\text{})`$. Then $`(h_\gamma (\delta )z)^{tt}`$ restricted on $`c_0(\text{})`$ equals $`h_\gamma (\delta )z`$, while $`S^{tt}`$ restricted on $`c_0(\text{})`$ equals $`S`$. Thus, dispensing with the double $`t`$’s, we can summarize the situation as follows:
Considering the linearly embedded Banach spaces $`\mathrm{}^1(\text{})\mathrm{}^2(\text{})\mathrm{}^{\mathrm{}}(\text{})`$, we have two linear operators on $`\mathrm{}^{\mathrm{}}(\text{})`$, namely, $`h_\gamma (\delta )z`$ and $`S`$, which are inverses of each other. Both map $`\mathrm{}^1(\text{})`$ into $`\mathrm{}^1(\text{})`$, and both are continuous on $`\mathrm{}^{\mathrm{}}(\text{})`$ as well as on $`\mathrm{}^1(\text{})`$ with regard to the respective Banach space norms. In addition, $`h_\gamma (\delta )z`$ maps $`\mathrm{}^2(\text{})`$ (continuously) into $`\mathrm{}^2(\text{})`$. On account of the Riesz-Thorin theorem (\[Kn\], Ch. IV), we conclude that $`S`$ maps $`\mathrm{}^2(\text{})`$ into $`\mathrm{}^2(\text{})`$ and that $`S`$ restricted on $`\mathrm{}^2(\text{})`$ is continuous with respect to the Hilbert space norm. Moreover, $`S`$ and $`h_\gamma (\delta )z`$, restricted on $`\mathrm{}^2(\text{})`$, are inverses of each other. Hence $`z\text{}\backslash Sp(h(\delta ))`$, and therefore $`Sp(h_\gamma (\delta ))Sp_0(h_\gamma (\delta ))`$.
2.8 Lemma. Let $`|\delta |>1`$. Then $`Sp(h_\gamma (\delta ))`$ is constant for $`|\delta |^1\beta ^1<|\gamma |<|\delta |\beta `$.
Proof: The Lemmas 2.3, 2.4 and 2.7 show that $`Sp(h_\gamma (\delta ))`$ is constant for $`|\delta |^1\beta ^1<|\gamma |<|\delta |\beta `$, with the possible exception when $`|\gamma |=|\delta |`$ and $`|\gamma |=|\delta |^1`$. Let $`Sp(\delta )`$ denote the constant spectrum covered by those cases. The lemmas listed also show that
$$Sp(\delta )Sp(h_\gamma (\delta ))\text{for}|\delta |^1\beta ^1<|\gamma |<|\delta |\beta .$$
Now consider the “Fourier expansion” of the resolvent of $`h(\delta )`$
$$(h(\delta )z)^1=\underset{p,q=\mathrm{}}{\overset{\mathrm{}}{}}c_{pq}(z)w_{pq}.$$
We have $`c_{pq}(z)=c_{|p|,q}(z)`$. From \[R2\], paragraph 4, we know that
$$(h_\gamma (\delta )z)^1=\underset{p,q=\mathrm{}}{\overset{\mathrm{}}{}}\gamma ^pc_{pq}(z)w_{pq}\text{for}z\text{}\backslash Sp(h_\gamma (\delta )).$$
These two facts taken together show that $`Sp(h_\gamma (\delta ))`$ increases as $`\left|\mathrm{log}|\gamma |\right|`$ increases. In conjunction with the inclusion established above, this shows that $`Sp(h_\gamma (\delta ))=Sp(\delta )`$ for $`|\delta |^1\beta ^1<|\gamma |<|\delta |\beta `$, as claimed.
Remark. The proof of Lemma 2.8 looks a bit twisted. Since it is not this author’s aspiration to state and prove the relevant facts in their utmost generality, he feels that the line of reasoning chosen in the context of the objective in this exposition is somewhat adequate. However, it should be mentioned that the claim of Lemma 2.8 is valid for arbitrary irrational numbers $`\alpha `$. Here is the brief sketch of a proof.
Consider the set $`𝒜^{(\omega )}`$ of analytic elements in $`𝒜`$. An element is called analytic if it has an exponentially decaying “Fourier expansion”. $`𝒜^{(\omega )}`$ is seen to be a $``$-subalgebra of $`𝒜`$. Even though the concept of the operator $`k_\gamma `$, as defined at the beginning of this paragraph, cannot be extended to the general setup, the concept of the algebra automorphism $`Ad(k_\gamma )`$ can. To this end, one needs to show that
$$Sp\left((e^{\pi \alpha i}\gamma u+\beta )^1(e^{\pi \alpha i}\gamma ^1u+\beta )\right)=\text{𝕋},\text{whenever}\beta ^1<|\gamma |<\beta ,$$
for all irrational numbers $`\alpha `$. This can be fairly easily proved by exploiting the spectral radius formula as well as the ergodicity of the irrational rotation on 𝕋 associated with $`\alpha `$, in conjunction with Ascoli’s theorem. One can use this information to show that the assignments
$$\begin{array}{cc}\hfill u& u\hfill \\ \hfill v& v(e^{\pi \alpha i}\gamma +\beta )^1(e^{\pi \alpha i}\gamma ^1u^{}+\beta )\hfill \end{array}$$
yield an algebra automorphism $`\rho _{(\gamma )}^{(0)}`$ of $`𝒜^{(\omega )}`$ (which does of course not preserve the involution $``$ unless $`|\gamma |=1`$).
Multiplying $`\rho _{(\gamma )}^{(0)}`$ from both sides by the automorphism $`\left(\genfrac{}{}{0pt}{}{10}{11}\right)`$ yields an automorphism $`\rho _{(\gamma )}`$ of $`𝒜^{(\omega )}`$ which has the property
$$\rho _{(\gamma )}\left(h_{\gamma \delta }(\delta )\right)=h_{\gamma \delta ^1}(\delta )\text{for}\beta ^1<|\gamma |<\beta .$$
This identity extends (2.2). Since an element in $`𝒜^{(\omega )}`$ is invertible in $`𝒜^{(\omega )}`$ if and only if it is invertible in $`𝒜`$, $`\rho _{(\gamma )}`$ preserves the spectrum of an element in $`𝒜^{(\omega )}`$.
Now consider the resolvent $`(h_\gamma (\delta )z)^1`$ of $`h_\gamma (\delta )`$. There is a power series
$$\underset{p,q=0}{\overset{\mathrm{}}{}}c_{pq}(z)x^py^q$$
in two variables such that
$$\left(h_\gamma (\delta )z\right)^1=\underset{p,q=\mathrm{}}{\overset{\mathrm{}}{}}\stackrel{~}{c}_{pq}(z)\gamma ^p\delta ^qw_{pq},$$
where throughout in each of the four $`(p,q)`$-quadrants either
$$\stackrel{~}{c}_{pq}(z)=c_{|p|,|q|}(z)\text{ or }\stackrel{~}{c}_{pq}(z)=0$$
(cf. \[R2\], paragraph 4). This power series depends on $`z`$ only and the type of the component in the resolvent set of $`h_\gamma (\delta )`$ to which $`z`$ belongs (i.e. bounded component vs. unbounded component). Since $`\rho _{(\gamma )}`$ preserves the spectrum of any element in $`𝒜^{(\omega )}`$, one can show that the domain of convergence of the said power series, which is known to be a logarithmically convex complete Reinhardt domain, must contain the polydisc
$$\{x\text{}||x|d\beta \}\times \{y\text{}||y|d\},$$
where $`d=\mathrm{max}\{|\delta |,|\delta |^1\}`$. This is valid whenever $`|\delta |>0`$ and $`d^1\beta ^1<|\gamma |<d\beta `$. It follows immediately that $`Sp(h_\gamma (\delta ))`$ is constant for $`d^1\beta ^1<|\gamma |<d\beta `$. (Notice that no exception needs to be made for the case $`|\delta |=1`$.)
We are now in a position to record a preliminary conclusion to our discussion regarding the existence of eigenvectors.
2.9 Proposition. For every $`\delta \text{}\backslash (\text{𝕋}\{0\})`$ and every $`zSp(h(\delta ))`$ there exist $`x\text{𝕋}`$, $`\xi \mathrm{}^2(\text{})\backslash \{0\}`$ such that
$$h(x\delta )\xi =z\xi ,\underset{|n|\mathrm{}}{lim}\left|\xi _n\right|^{1/|n|}\mathrm{min}\{|\delta |\beta ^1,|\delta |^1\beta ^1\}.$$
Proof: If $`|\delta |>1`$, the claim follows from Lemmas 2.3, 2.4, 2.7 and 2.8. Since
$$𝒥h(\delta )𝒥=h(\delta ^1),\text{where}(𝒥\xi )_n=\xi _n,$$
the case $`0<|\delta |<1`$ can be reduced to the case $`|\delta |>1`$.
The translation by $`x`$ in Proposition 2.9, and hence $`\xi `$, is (essentially) uniquely determined by $`z`$.
2.10 Proposition. If $`\delta >1`$; $`x,y\text{𝕋}`$ and
$$h(x\delta )\xi =z\xi ,h(y\delta )\eta =z\eta ;\xi ,\eta \mathrm{}^2(\text{}),$$
then $`y=\lambda ^2\mathrm{}x`$ for some $`\mathrm{}\text{}`$ and $`u^{\mathrm{}}\xi `$, $`\eta `$ are linearly dependent.
Proof: We have
$$h_x(x\delta )D_x\xi =zD_x\xi ,h_{\overline{y}}(y\delta )D_{\overline{y}}\eta =zD_{\overline{y}}\eta .$$
So, if we let
$$\phi (a)=aD_x\xi ,D_y\overline{\eta },a𝒜,$$
then
$$\phi (ah_x(x\delta ))=\phi (h_y(y\delta )a)=z\phi (a).$$
Now let
$$\phi _{pq}=s^{p+q}\delta ^q\phi (w_{pq}),\text{where}s^2=\overline{x}\overline{y}.$$
Then the double sequence $`\left\{\phi _{pq}\right\}`$ is seen to solve the system of difference equations
$$\begin{array}{cc}\hfill \mathrm{cos}(\pi \alpha q+\theta )(X_{p1,q}+X_{p+1,q})& +\beta \mathrm{cos}(\pi \alpha p+\theta )(X_{p,q1}+X_{p,q+1})=zX_{pq}\hfill \\ \hfill \mathrm{sin}(\pi \alpha q+\theta )(X_{p1,q}X_{p+1,q})& \beta \mathrm{sin}(\pi \alpha p+\theta )(X_{p,q1}X_{p,q+1})=0,\hfill \end{array}$$
$`()`$
where $`e^{i\theta }=s`$. Starting over again with $`\overline{\xi }`$ and $`\overline{\eta }`$ in place of $`\xi `$ and $`\eta `$, respectively, we get
$$h_x(x\delta ^1)D_x\overline{\xi }=\overline{z}D_x\overline{\xi },h_{\overline{y}}(y\delta ^1)D_{\overline{y}}\overline{\eta }=\overline{z}D_{\overline{y}}\overline{\eta }.$$
Let
$$\psi (a)=aD_x\overline{\xi },D_y\eta ,a𝒜.$$
Now let
$$\psi _{pq}=s^{p+q}\delta ^q\psi (w_{pq}),s\text{ as before.}$$
Replacing $`z`$ by $`\overline{z}`$ in $`()`$ it is seen that the double sequence $`\{\psi _{pq}\}`$ solves $`()`$. So both, $`\{\phi _{pq}\}`$ and $`\{\overline{\psi }_{pq}\}`$ solve $`()`$ for the same parameter $`z`$. Moreover, the sequence
$$\left\{\phi _{pp}\psi _{p+1,p+1}\phi _{p+1,p+1}\psi _{pp}\right\}$$
is bounded.
Now suppose that the claim of the proposition is not true. Then it follows from appendix A1 that $`\{\phi _{pq}\}`$ and $`\{\overline{\psi }_{pq}\}`$ are linearly dependent. This however implies that $`\{\phi _{pq}\}`$ decays exponentially uniformly in $`p`$ of order at least $`\delta ^2`$ as $`q\mathrm{}`$. Since the sequence $`\{\psi _{pq}\}_q\text{}`$ is almost periodic for every $`p\text{}`$, and hence does not approach zero as $`q\mathrm{}`$, we have reached a contradiction.
Our next objective is to show that the spectrum of $`h=u^{}+u+\beta (v+v^{})`$ is a regular compactum in the sense of potential theory. In view of \[R3\], Theorem 2.2 we are going to prove a stronger statement. In preparation of this, we need the following.
2.11 Lemma. $`Sp(h)\text{}\backslash Sp(h(\delta ))`$ whenever $`|\delta |1`$.
Proof: Suppose this were not true. Then there exists $`\chi Sp(h)`$ and $`\delta _0>1`$ such that $`\chi Sp(h(\delta ))`$ for $`\delta _0^1\delta \delta _0`$. According to Proposition 2.9, for any $`\delta [\delta _0^1,\delta _0]\backslash \{1\}`$, there exist $`x\text{𝕋}`$ and $`\xi \mathrm{}^2(\text{})\backslash \{0\}`$ such that
$$h(x\delta )\xi =\chi \xi ,\underset{|n|\mathrm{}}{\overline{\mathrm{lim}}}|\xi |^{1/|n|}\delta ^1\beta ^1,$$
For every $`y\text{𝕋}`$, let
$$\eta _n^{(y)}=\underset{p=\mathrm{}}{\overset{\mathrm{}}{}}y^px^n\delta ^nw_{pn}\xi ,𝒥\xi w_{pn}e^{(0)},$$
where ($`𝒥`$ as in 2.9),
$$e_n^{(0)}=\{\begin{array}{cc}1,\hfill & n=0\hfill \\ 0,\hfill & n0.\hfill \end{array}$$
Let $`\widehat{h}(y)=u+u^{}+\beta ^1(yv+\overline{y}v^{})`$. Then
$$\widehat{h}(y)\eta ^{(y)}=\chi \eta ^{(y)},\underset{n\mathrm{}}{\overline{\mathrm{lim}}}\left|\eta _n^{(y)}\right|^{1/n}\delta ^1,\underset{n\mathrm{}}{\overline{\mathrm{lim}}}\left|\eta _n^{(y)}\right|^{1/|n|}\delta .$$
Exempting for every $`\delta [\delta _0^1,\delta _0]\backslash \{1\}`$ a possibly non-empty, but countable subset of values for the parameter $`y`$, for which $`\eta ^{(y)}`$ might be zero, this implies the following: There is a countable subset $`M\text{𝕋}`$ and a dense countable subset $`N[\delta _0^1,\delta _0]\backslash \{1\}`$ such that for every $`y\text{𝕋}\backslash M`$ and every $`\delta N`$ there exists $`\eta \text{}^{\mathrm{}}\backslash \{0\}`$ such that
$$\widehat{h}(y)\eta =\chi \eta ,\underset{n\mathrm{}}{\overline{\mathrm{lim}}}\left|\eta _n\right|^{1/n}\delta ^1,\underset{n\mathrm{}}{\overline{\mathrm{lim}}}\left|\eta _n\right|^{1/|n|}\delta .$$
This in turn entails that for every $`y\text{𝕋}\backslash M`$, the operator $`\widehat{h}(y)`$ has an exponentially decaying eigenvector for $`\chi `$. Since $`\widehat{h}(y)`$ is known to have no eigenvalues for any $`y\text{𝕋}`$, we have reached a contradiction.
Letting $`\mu `$ be the probability measure of $`Sp(h)`$ obtained through restricting the canonical trace $`\tau `$ of $`𝒜`$ on the $`C^{}`$-algebra generated by $`h`$, Lemma 2.11 entails, on account of \[R3\], Theorem 2.2:
2.12 Proposition. $`Sp(h(\delta ))=\{z\text{}\mathrm{log}|zs|d\mu (s)=\mathrm{log}(\beta \delta )\}\mathrm{for}\delta 1`$. In particular $`Sp(h)`$ is a regular compactum and $`\mu `$ is its equilibrium distribution.
2.13 Corollary. $`Sp(h)`$ is not connected.
Proof: Consider the moments $`\tau (h^n)`$ of $`h`$. While $`\tau (h^{2n1})=0`$ for every $`n\text{}`$, elementary calculations show that
$$\tau (h^2)=2\beta ^2+2,\tau (h^4)=6\beta ^4+(24+16\mathrm{cos}2\pi \alpha )\beta ^2+6.$$
Now consider $`v+v^{}=\underset{\beta \mathrm{}}{lim}\beta ^1(u+u^{}+\beta (v+v^{}))`$. This element has a connected spectrum and $`\tau `$ restricted to the $`C^{}`$-algebra generated by $`v+v^{}`$ is nothing but the equilibrium distribution for $`Sp(v+v^{})=[2,2]`$. Again we have $`\tau (v+v^{})^{2n1})=0`$ for every $`n\text{}`$. Moreover,
$$\tau ((v+v^{})^2)=2,\tau ((v+v^{})^4)=6.$$
In order to identify possible values for $`\beta `$ for which $`Sp(h)`$ is connected, we have to solve the equations
$$2x=\tau (h^2),6x^2=\tau (h^4),$$
for $`x`$ and $`\beta `$. Eliminating $`x`$, we obtain
$$3\tau (h^2)^2=2\tau (h^4),$$
which in turn yields
$$(1+\mathrm{cos}2\pi \alpha )\beta ^2=0.$$
Since $`\alpha `$ is irrational this is valid only for $`\beta =0`$. Note that $`\beta =0`$ does indeed correspond to an element with a connected spectrum, namely $`u+u^{}`$, which is the image of $`v+v^{}`$ under a “Fourier transform”.
We will show now that the eigenvectors of the operators $`h(\delta )`$ depend continuously on the spectral parameter $`z`$ in a sense to be made precise below.
2.14 Proposition. Let $`d>1`$, and for every $`m\text{}`$ let $`|\delta _m|d`$, $`z_mSp(h(\delta _m))`$, $`\xi ^{(m)}\mathrm{}^2(\text{})\backslash \{0\}`$ such that
$$h(\delta _m)\xi ^{(m)}=z_m\xi ^{(m)},\underset{|n|\mathrm{}}{\overline{\mathrm{lim}}}\left|\xi _n^{(m)}\right|^{1/|n|}|\delta _m|^1\beta ^1$$
and $`\underset{m\mathrm{}}{lim}z_m=z`$.
Then there exist $`\delta \text{}`$, $`\xi \mathrm{}^2(\text{})\backslash \{0\}`$ and $`j_m\text{}`$, $`c_m\text{}\backslash \{0\}`$ such that $`h(\delta )\xi =z\xi `$, $`\underset{m\mathrm{}}{lim}\lambda ^{2j_m}\delta _m=\delta `$ and
$$\underset{m\mathrm{}}{lim}D_\gamma \left(c_mu^{j_m}\xi ^{(m)}\xi \right)_{\mathrm{}}=0\text{uniformly for}d^1\gamma d.$$
Proof: It follows from (2.2) that there exists $`t_m\text{}`$ such that
$$kD_{\delta _m}\xi ^{(m)}=t_mD_{\delta _m}^1\xi ^{(m)}.$$
Let
$$\eta ^{(m)}=D_{\delta _m}\xi ^{(m)}_{\mathrm{}}^1D_{\delta _m}\xi ^{(m)}.$$
Then $`\eta ^{(m)}_{\mathrm{}}=1`$ and there exist $`j_m\text{}`$ such that
$$\left|\eta _{j_m}^{(m)}\right|\frac{1}{2}.$$
Adjusting $`\eta ^{(m)}`$ suitably through shifting and multiplication by scalars, we may assume that
$$\eta _0^{(m)}\frac{1}{2}.$$
Finally, switching to a subsequence if necessary, we may assume that the sequence $`\eta ^{(1)},\eta ^{(2)},\mathrm{}`$ converges weakly to some $`\eta \mathrm{}^{\mathrm{}}(\text{})\backslash \{0\}`$. Adjusting the $`t_m`$ accordingly we have
$$k\eta ^{(m)}=t_mD_{\delta _m}^2\eta ^{(m)}.$$
Since $`k\eta 0`$, it follows that $`inf\{|t_m|m\text{}\}>0`$. Also, since Proposition 2.12 implies that $`|\delta _1|,|\delta _2|,\mathrm{}`$ is convergent, $`sup\{|t_m|m\text{}\}<\mathrm{}`$. For, if this were not the case, then we could find some point of density $`\delta `$ for the $`\delta _m`$ such that $`D_\delta ^2\eta =0`$.
Switching once again to a subsequence if necessary, we may assume that the sequences $`\{\delta _m\}`$ and $`\{t_m\}`$ are convergent. Let $`\delta =\underset{m\mathrm{}}{lim}\delta _m`$. Since $`k`$ defines a bounded operator on $`\mathrm{}^{\mathrm{}}(\text{})`$, it follows that the corresponding sequence of the $`\xi ^{(m)}`$, after having been suitably scaled and adjusted through shifts, does indeed converge uniformly to some eigenvector $`\xi `$ of $`h(\delta )`$ for $`z`$, as claimed in the proposition. Since eigenvectors of $`h(\delta )`$ are essentially unique by Proposition 2.10, the conclusion of the proposition applies to the original sequence of the $`\xi ^{(m)}`$ as well.
Remark: Tracing the steps in the proof of Lemmas 2.1 and 2.2 one can improve the convergence of eigenvectors in Proposition 2.14 to the effect that $`\gamma `$ may range over any closed interval contained in $`(d^1\beta ^1,d\beta )`$.
Our next objective is to refine Propositions 2.9 and 2.13 by showing that there is a natural parametrization of the eigenvalue problem for the operators $`h(\delta )`$, $`|\delta |>1`$, through a Riemann surface $`\stackrel{~}{}`$ covering the resolvent set $`=\text{}\backslash Sp(h)`$. By doing so we will discover that all the eigenvectors of these operators correspond to orbits of a cyclic group of covering transformations on $`\stackrel{~}{}`$, while the square of their components can be obtained through evaluation of a single analytic function on $`\stackrel{~}{}`$. For the basic concepts of Riemann surfaces we are going to use and for the terminology that comes with it, we refer to \[AS\], Chapters I and II.
Let $`I`$ be the smallest interval containing $`Sp(h)`$. There exists a unique analytic function $`G`$ on $`\text{}\backslash I`$ such that
$$\mathrm{log}|G(z)|=\mathrm{log}|zs|d\mu (s)\mathrm{log}\beta ,G(\text{}^+\backslash Sp(h))\text{}^+.$$
If $`f`$ is a closed arc surrounding a component $`K`$ of $`Sp(h)`$, then any two analytic continuations of $`G`$ along $`f`$ over the same point $`z`$ on $`f`$ differ by a multiplicative constant of the form
$$e^{2\pi \mu (K)ni}.$$
A well-known fact in the $`K`$-theory of the irrational rotation $`C^{}`$-algebra provides us with the information that
$$\mu (K)\left(\text{}+\alpha \text{}\right)[0,1],$$
so that the said multiplicative factor takes the form $`\lambda ^{2n}`$. Now $`G`$ can be continued analytically along any arc in $`=\text{}\backslash Sp(h)`$. So, if we let $``$ be the universal covering of $``$, $`𝒢`$ the corresponding group of covering transformations, and finally $`\stackrel{~}{G}`$ the analytic function on $``$ obtained through analytic continuation of $`G`$, then we have for every $`g𝒢`$
$$\stackrel{~}{G}g=\lambda ^{2n}\stackrel{~}{G}\text{for some}n\text{}.$$
Let
$$𝒢_0=\{g𝒢\stackrel{~}{G}g=\stackrel{~}{G}\}.$$
Then $`𝒢_0`$ is a normal subgroup of $`𝒢`$ with a cyclic quotient group $`𝒢/𝒢_0`$. The $`𝒢_0`$-orbits in $``$ form a Riemann surface $`\stackrel{~}{}`$ which covers $``$, and whose group of covering transformations corresponds to $`𝒢/𝒢_0`$ in a natural way. Let $`p`$ be the covering map of $`\stackrel{~}{}`$ over $``$. Finally, $`\stackrel{~}{G}`$ defines an analytic map on $`\stackrel{~}{}`$ which we denote by $`G`$ again. We assume $`\omega `$ to be that generator of the group of covering transformations of $`\stackrel{~}{}`$ over $``$ for which $`G(\omega (z))=\lambda ^2G(z)`$.
We are now in a position to state the major claim in this paragraph. For its proof, we need another technical lemma.
2.15 Lemma. Suppose $`f:\stackrel{~}{}\text{𝕋}`$ has the following properties
(i) Every subsequence of $`\{f^nn\text{}\}`$ has a subsequence which converges uniformly on compact subsets of $`\stackrel{~}{}`$.
(ii) If $`\underset{n\mathrm{}}{lim}z_n=z`$ and $`\underset{n\mathrm{}}{lim}f(z_n)=c`$, then $`c\{f(z),\lambda ^2f(z)\}`$.
Then $`f(\stackrel{~}{})\{x\lambda ^{2n}n\text{}\}`$ for some $`x\text{𝕋}`$.
Proof: Let $`K\stackrel{~}{}`$ be compact and connected. It suffices to show that $`f(K)\{x\lambda ^{2n}n\text{}\}`$ for some $`x\text{𝕋}`$. Let $`\mathrm{\Omega }`$ be the uniform closure in $`\mathrm{}^{\mathrm{}}(K)`$ of $`\{g^nn\text{}\}`$, where $`g=f/k`$. It follows from (i) that $`\mathrm{\Omega }`$ is a metrizable and compact group. To every point $`y`$ in $`g(K)`$ there corresponds a character $`\phi _y`$ of $`\mathrm{\Omega }`$, a continuous homomorphism from $`\mathrm{\Omega }`$ into 𝕋, such that $`\phi _y(g)=y`$. Let $`𝒴=\{\phi _yyg(K)\}`$. Since the uniform structure on $`\mathrm{\Omega }`$ is inherited from $`\mathrm{}^{\mathrm{}}(K)`$, the set $`𝒴`$ is equicontinuous on $`\mathrm{\Omega }`$. Hence $`\overline{𝒴}`$, the closure of $`𝒴`$ in the topology of uniform convergence on $`\mathrm{\Omega }`$, is compact. In other words, $`\overline{𝒴}`$ is a compact subset of the dual group $`\widehat{\mathrm{\Omega }}`$ of $`\mathrm{\Omega }`$. Since $`\mathrm{\Omega }`$ is compact,$`\widehat{\mathrm{\Omega }}`$ is discrete. Therefore $`\overline{𝒴}`$, being a compact subset of $`\widehat{\mathrm{\Omega }}`$, must be finite. It follows that $`g(K)`$ is a finite subset of 𝕋.
For $`x\text{𝕋}`$, let $`𝒪_x=\{x\lambda ^{2n}n\text{}\}`$. Now let $`x\text{𝕋}`$ be such that
$$M=g(K)𝒪_x\mathrm{}.$$
Let $`N=g(K)\backslash 𝒪_x`$. Since $`g(K)`$ is finite, $`M`$ and $`N`$ are finite as well and hence, closed. Therefore, since $`K`$ is compact, property (ii) implies
$$\overline{g^1(M)}M\lambda ^2M,\overline{g^1(N)}N\lambda ^2N.$$
By the definition of $`M`$ and $`N`$
$$(M\lambda ^2M)(N\lambda ^2N)=\mathrm{}.$$
Hence,
$$\overline{g^1(M)}\overline{g^1(N)}=\mathrm{},\overline{g^1(M)}\overline{g^1(N)}=K.$$
Since $`K`$ is connected by assumption,
$$g^1(M)=\overline{g^1(M)}=K,$$
which settles the claim.
2.16 Theorem. For every $`z\stackrel{~}{}`$, there exists $`\xi ^{(z)}\mathrm{}^2(\text{})\backslash \{0\}`$ such that
$$h(G(z))\xi ^{(z)}=p(z)\xi ^{(z)},\underset{|n|\mathrm{}}{\overline{\mathrm{lim}}}\left|\xi _n^{(z)}\right|^{1/|n|}|G(z)|^1\beta ^1,$$
and for every $`p\text{}`$, $`\vartheta _p(z)=\xi _0^{(z)}\xi _p^{(z)}`$ is an analytic function on $`\stackrel{~}{}`$ with the property
$$\vartheta _p(\omega ^n(z))=\xi _n^{(z)}\xi _{n+p}^{(z)}.$$
Proof: Propositions 2.9 and 2.12 taken together imply that for every $`z\stackrel{~}{}`$, there exist $`x_z\text{𝕋}`$, $`\xi ^{(z)}\mathrm{}^2(\text{})\backslash \{0\}`$ such that
$$h(x_zG(z))\xi ^{(z)}=p(z)\xi ^{(z)},\overline{\underset{|n|\mathrm{}}{lim}}\left|\xi _n^{(z)}\right|^{1/|n|}|G(z)|^1\beta ^1.$$
On account of (2.2), there exists for every $`z\stackrel{~}{}`$ a $`t_z\text{}`$ such that
$$kD_{x_zG(z)}\xi ^{(z)}=t_zD_{x_zG(z)}^1\xi ^{(z)}.$$
$`()`$
Applying $`u^m`$ to both sides of this identity yields
$$kD_{x_zG(z)}u^m\xi ^{(z)}=(x_zG(z))^2t_zD_{x_zG(z)}^1u^m\xi ^{(z)}.$$
Therefore, we can adjust $`x_z`$, $`\xi ^{(z)}`$ and $`t_z`$ in $`()`$ such that
$$1|t_z|<|G(z)|^2.$$
It follows from Proposition 2.10 that
$$\underset{n\mathrm{}}{lim}z_n=z,\underset{n\mathrm{}}{lim}|t_{z_n}|=s\text{and}\underset{n\mathrm{}}{lim}x_{z_n}=y$$
implies
$$s\{|t_z|,|G(z)|^2|t_z|\},y\{x_z,\lambda ^2x_z\}.$$
Now let
$$\phi ^{(z)}(a)=a\xi ^{(z)},\overline{\xi ^{(z)}},a𝒜.$$
Then
$$\phi ^{(z)}(h(x_zG(z))a)=\phi (ah(x_zG(z))=p(z)\phi ^{(z)}.$$
So, if we define
$$\psi _{pq}^{(z)}=\overline{x}_z^qG(z)^q\phi (w_{pq});p,q\text{},$$
then the double sequence $`\left\{\psi _{pq}^{(z)}\right\}`$ is a solution to the system of difference equations in the Appendix A2. Also, $`\left\{\psi _{pq}^{(z)}\right\}`$ decays exponentially uniformly in $`p`$ with an order less than or equal to $`\beta ^1`$. Therefore, it follows from A2 that, after scaling the eigenvectors $`\xi ^{(z)}`$ suitably,
$$\lambda ^{pq}\overline{x}_z^q\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\overline{\lambda }^{2qn}\xi _n^{(z)}\xi _{n+p}^{(z)}=(c_{pq}(p(z))d_{pq}(p(z))G(z)^q;p,q\text{},$$
$`()`$
where $`c_{pq}`$ and $`d_{pq}`$ are analytic functions on $``$ which are determined by the identities
$$\begin{array}{cc}\hfill hx)^1& =\underset{p,q=\mathrm{}}{\overset{\mathrm{}}{}}c_{pq}(x)w_{pq}\text{and}\hfill \\ \multicolumn{2}{c}{}\\ \hfill (h(\delta )x)^1& =\underset{p,q=\mathrm{}}{\overset{\mathrm{}}{}}d_{pq}(x)\delta ^qw_{pq},|\delta |>|G(z)|.\hfill \end{array}$$
If $`d>1`$ and $`K\{z\stackrel{~}{}|G(z)|d\}`$ is compact, then Proposition 2.11 ensures that
$$D_\gamma \xi ^{(z)}_{\mathrm{}}\text{is uniformly bounded in }zK,d^1|\gamma |d.$$
Since the functions on the right hand side of (2.4) are analytic on $``$ as well as uniformly bounded on $`K`$ in $`p`$ and $`q`$, we conclude that if we let $`f(z)=\overline{x}_z`$, then any subsequence of $`\{f^q\}`$ has a subsequence which is locally uniformly convergent. Thus we have show that $`f`$ has the properties (i) and (ii) in Lemma 2.15. It follows that
$$f(\stackrel{~}{})\{\overline{y}\lambda ^{2n}n\text{}\}\text{for some}y\text{𝕋}.$$
By applying suitable powers of $`u`$ to the vectors $`\xi ^{(z)}`$, we can adjust $`()`$ to the effect that $`x_z=y`$ for every $`z\stackrel{~}{}`$.
Let $`t>\mathrm{max}Sp(h)`$. Then there exists $`z_t\stackrel{~}{}`$ such that $`G(z)\text{}^+`$. Hence
$$h(yG(z_t))\xi ^{(z_t)}=t\xi ^{(z_t)},h(\overline{y}G(z_t)𝒥\overline{\xi ^{(z_t)}}=t𝒥\overline{\xi ^{(z_t)}}$$
and by Proposition 2.10, $`\overline{y}=\lambda ^{2m}y`$ for some $`m\text{}`$. Thus, adjusting the eigenvectors through a suitable power of the shift we may assume that $`y\{1,1,\lambda ,\lambda \}`$. In order to show that $`y`$ equals actually $`1`$, one can proceed as follows: Since $`\beta ^1G(z_t)h(yG(z_t))`$ approaches $`yv`$ as $`t\mathrm{}`$, one can use $`()`$ to show that $`\xi ^{(z_t)}`$ approaches an eigenvector $`\eta `$ of $`yv`$ (in $`\mathrm{}^2(\text{})`$, say) for some positive eigenvalue $`s`$. Thus $`s\{\overline{y}\lambda ^{2n}n\text{}\}`$. The only way this can happen is when $`y=1`$.
Proposition 2.13 in conjunction with $`()`$ entails that the function
$$\psi _p:\text{𝕋}\times \stackrel{~}{}\text{};\psi _p(x,z)=\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}x^n\xi _n^{(z)}\xi _{n+p}^{(z)}$$
is continuous and analytic in $`z`$ for every $`x\text{𝕋}`$. It follows that the functions
$$\vartheta _p(z)=\psi _p(x,z)𝑑x=\xi _0^{(z)}\xi _p^{(z)}$$
are analytic. Moreover, since $`G(\omega (z))=\lambda ^2G(z)`$,
$$\vartheta (\omega ^n(z))=\xi _n^{(z)}\xi _{n+p}^{(z)}$$
by $`()`$, as claimed.
Remarks. 1) The second part of Theorem 2.16 says (take $`p=0`$), that given $`\xi _0^{(z)}`$ in a neighborhood of $`z`$, one can generate $`\left(\xi _n^{(z)}\right)^2`$ through analytic continuation of $`\left(\xi _0^{(z)}\right)^2`$.
2) Theorem 2.16 also shows that we can identify the points of the Riemann surface $`\stackrel{~}{}`$ with the one-dimensional eigenspaces of the operators $`h(G(z))`$ for $`p(z)`$ in such a way that the covering transformation $`\omega `$ corresponds to the two-sided shift on $`\mathrm{}^2(\text{})`$.
3) From $`()`$ in the proof of Theorem 2.16 we can extract the identity
$$\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\left(\xi _n^{(z)}\right)^2=(tp(z))^1𝑑\mu (t).$$
It follows that $`\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\left(\xi _n^{(z)}\right)^2`$ equals zero if and only if $`z`$ is a critical point for the conductor potential of $`Sp(h)`$. These critical points are simple; and exactly one is located in every gap of $`Sp(h)`$ and nowhere else. This shows that the eigenvectors $`\xi ^{(z)}`$ carry the relevant information regarding the gap structure for $`Sp(h)`$ in a very explicit form.
4) Elaborating on the comments made above, there are some implications for the case $`|\delta |=1`$. One might be tempted to try to generate exponentially decaying eigenvectors for the operators $`h(x)`$, $`|x|=1`$, as (uniform) limits of the eigenvectors $`\xi ^{(z)}`$. Assuming that $`Sp(h)`$ has infinitely many gaps, one is confronted with the following impediment: Let $`z_1,z_2,\mathrm{}`$ be a sequence of critical points converging to $`\chi Sp(h)`$. Suppose that $`h(x)\xi =\chi \xi `$ for some $`\xi \mathrm{}^2(\text{})\backslash \{0\}`$. Then
$$\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\xi _n^2=y\xi _20\text{for some}y\text{𝕋}.$$
It follows that $`\xi `$ cannot be approximated by the (suitably scaled) sequence $`\xi ^{(z_1)},\xi ^{(z_2)},\mathrm{}`$ in the Hilbert space norm.
2.17 Corollary. There exist homeomorphisms $`\sigma `$ and $`\iota `$ on $`\stackrel{~}{}`$ such that
$$\begin{array}{cc}\hfill \xi ^{(\sigma (z))}& =𝒥\overline{\xi ^{(z)}},p(\sigma (z))=\overline{p(z)},G(\sigma (z))=\overline{G(z)}\hfill \\ \hfill \xi ^{(\iota (z))}& =D_1\xi ^{(z)},p(\iota (z))=p(z),G(\iota (z))=G(z)\hfill \end{array}.$$
In particular, $`\sigma ^2=\iota ^2=Id`$, $`\sigma \omega =\omega ^1\sigma `$, $`\omega \iota =\iota \omega `$.
Proof: This follows immediately from Theorem 2.16, noting that
$$h(\overline{G(z)})𝒥\overline{\xi ^{(z)}}=\overline{p(z)}𝒥\overline{\xi ^{(z)}}$$
and
$$h\left(G(z)\right)D_1\xi =\overline{p}(z)D_1\xi ^{(z)}.$$
The following is in preparation for the discussion in paragraph 3.
2.18 Corollary. There exists an analytic function $`\mathrm{\Gamma }`$ on $`\stackrel{~}{}`$ which has the following properties:
$$kD_{G(z)}\xi ^{(z)}=\mathrm{\Gamma }(z)D_{G(z)^1}\xi ^{(z)},\mathrm{\Gamma }(\omega (z))=G(z)^2\mathrm{\Gamma }(z),\mathrm{\Gamma }(\sigma (z))=\overline{\mathrm{\Gamma }(z)}^1$$
$$\text{and}\mathrm{\Gamma }(\iota (z))=\mathrm{\Gamma }(z).$$
Proof: There clearly exists a function satisfying the stated identities. To show that it is actually analytic, we consider scalar products
$$kD_{G(z)}\xi ^{(z)},a\overline{\xi ^{(z)}}=\mathrm{\Gamma }(z)D_{G(z)^1}\xi ^{(z)},a\overline{\xi ^{(z)}},$$
where $`a`$ is any linear combination of the elements $`w_{pq}`$. Written out in components, one can see that the function on the left as well as the second factor on the right are algebraically generated by $`G`$, $`G^1`$ and functions of the form $`\vartheta _p\omega ^n`$, all of which are analytic. Therefore, $`\mathrm{\Gamma }`$ has to be analytic as well.
3. The kernel of a family of related operators The first identity in Corollary 2.18 gives rise to an eigenvalue problem in its own right, involving the operator $`k`$ and the function $`\mathrm{\Gamma }`$ as an eigenvalue parameter. The question arises whether the two related eigenvalue problems are actually equivalent. It will be shown that this is essentially the case. More specifically, it will be shown that to any two distinct points $`z_1`$ and $`z_2`$ in the complement of a (possibly empty) discrete subset $`\stackrel{~}{}`$, there correspond distinct pairs of parameters $`(G(z_1),\mathrm{\Gamma }(z_1))`$ and $`(G(z_2),\mathrm{\Gamma }(z_2))`$. To this end, we introduce new operators whose kernel will hold all the relevant information.
Let $`g`$ be the function defined in (1.6) and let $`\gamma \text{}`$ be close to $`1`$ having the property that
$$g(\gamma \text{𝕋})(\left\{\pi (n+\frac{1}{2})|n\text{}\}\right)\{\pi (m+i(n+\frac{\pi }{4}))|m,n\text{}\}=\mathrm{}.$$
Let $`\mathrm{\Gamma }`$ be as in Corollary 2.18 and let
$$\mathrm{\Omega }=\{z\stackrel{~}{}1+\mathrm{\Gamma }(\omega ^n(z))0\text{for all}n\text{}\}.$$
Since $`\mathrm{\Gamma }`$ is analytic, $`\mathrm{\Omega }`$ is a discrete subset of $`\stackrel{~}{}`$. In the following, $`\mathrm{\Omega }`$ will be augmented as needed, but it will always be discrete. Now we define for every $`z\stackrel{~}{}\backslash \mathrm{\Omega }`$ a bounded operator $`H_\gamma (z)`$ on the Hilbert space $`\mathrm{}^2(\text{})`$,
$$\left(H_\gamma (z)\xi \right)_n=\underset{j=\mathrm{}}{\overset{\mathrm{}}{}}a_j\xi _{n+j}+i\frac{1\mathrm{\Gamma }(\omega ^n(z))}{1+\mathrm{\Gamma }(\omega ^n(z))}\xi _n,$$
where
$$\mathrm{tan}g(\gamma u)=\underset{j=\mathrm{}}{\overset{\mathrm{}}{}}a_ju^j.$$
The sequence $`\{a_j\}`$ decays exponentially as $`|j|\mathrm{}`$. Obviously, $`H(z)`$ is analytic in $`z`$.
Moreover,
$$\{\begin{array}{cc}\hfill H_\gamma (\omega (z))& =uH_\gamma (z)u^{}\hfill \\ \hfill H_\gamma (\sigma (z))& =𝒥H_{\gamma ^1}(z)^{}𝒥.\hfill \end{array}$$
$`(3.1)`$
3.1 Proposition. The operator $`H(z)`$ is Fredholm with index zero. Any two of these operators are compact perturbations of each other, and their essential spectrum equals
$$(\mathrm{tan}g(\gamma \text{𝕋})+i)(\mathrm{tan}g(\gamma \text{𝕋})i).$$
Proof: Let $`T(\mathrm{}^2(\text{}))`$ be defined as follows:
$$(T\xi )_n=\underset{j=\mathrm{}}{\overset{\mathrm{}}{}}a_j\xi _{n+j}+(n)\xi _n,$$
where
$$(n)=\{\begin{array}{cc}i,\hfill & n0\hfill \\ \multicolumn{2}{c}{}\\ i,\hfill & n<0\text{.}\hfill \end{array}$$
Since $`|\mathrm{\Gamma }(\omega ^n(z))|=|G(z)^{2n}\mathrm{\Gamma }(z)|`$ and $`|G(z)|>1`$,
$$\underset{n\mathrm{}}{lim}i\frac{1\mathrm{\Gamma }(\omega ^n(z))}{1+\mathrm{\Gamma }(\omega ^n(z))}=i,\underset{n\mathrm{}}{lim}i\frac{1\mathrm{\Gamma }(\omega ^n(z))}{1+\mathrm{\Gamma }(\omega ^n(z))}=i.$$
It follows that $`H(z)`$ is a compact perturbation of $`T`$. Let
$$T_0=T,$$
and let $`\stackrel{~}{T}_0`$, $`\stackrel{~}{}`$ be the range of $`T_0`$, $``$, respectively, in the Calkin algebra $`(\mathrm{}^2(\text{}))/𝒦`$, where $`𝒦`$ denotes the $`C^{}`$-algebra of compact operators on $`\mathrm{}^2(\text{})`$. Then $`\stackrel{~}{T}_0`$ and $`\stackrel{~}{}`$ are normal operators which commute and their joint spectrum equals
$$\mathrm{tan}g(\text{𝕋})\times \{i,i\}.$$
It follows that
$$Sp\left(\stackrel{~}{T}_0+\stackrel{~}{}\right)=(\mathrm{tan}g(\gamma \text{𝕋})+i)(\mathrm{tan}g(\gamma \text{𝕋})i)$$
which equals the essential spectrum of $`T`$. By our choice of $`\gamma `$, $`\stackrel{~}{T}_0+\stackrel{~}{}`$ is invertible, which entails that $`H(z)`$ is Fredholm with index zero.
In the following we will constantly make use of Theorem 2.16. In particular, the eigenvectors which occur will be assumed to decay of a sufficiently high order, so that all manipulations make sense, unless specified otherwise. The significance of the operators $`H(z)`$ rests with the following statement.
3.2 Proposition. Suppose $`h_\gamma (G(z))\xi =p(z)\xi `$ for some $`z\stackrel{~}{}\backslash \mathrm{\Omega }`$. Then there exist $`t_n\text{}`$ such that
$$t_n^2=\mathrm{\Gamma }(\omega ^n(z)),H_\gamma (z)\eta =0,$$
where
$$\eta _n=\left(t_n+t_n^1\right)\xi _n.$$
Proof: By Corollary 2.18
$$k_\gamma D_{G(z)}\xi =\mathrm{\Gamma }(z)D_{G(z)^1}\xi .$$
Rearranging the diagonal operators involved, we obtain
$$Te^{ig(\gamma u)}\xi =T^1\xi $$
where $`T`$ is a diagonal operator such that $`(T\xi )_n=t_n\xi _n`$ with $`t_n^2=\mathrm{\Gamma }(\omega ^n(z))`$. Using the identity
$$e^{ig(\gamma u)}=\frac{1+i\mathrm{tan}\frac{g(\gamma u)}{2}}{1i\mathrm{tan}\frac{g(\gamma u)}{2}}$$
we have
$$\left(1+i\mathrm{tan}\frac{g(\gamma u)}{2}\right)T\xi =\left(1i\mathrm{tan}\frac{g(\gamma u)}{2}\right)T^1\xi $$
or
$$t_n\xi _n+i\underset{j=\mathrm{}}{\overset{\mathrm{}}{}}a_jt_{n+j}\xi _{n+j}=t_n^1\xi _ni\underset{j=\mathrm{}}{\overset{\mathrm{}}{}}a_jt_{n+j}^1\xi _{n+j}$$
which turns into
$$\underset{j=\mathrm{}}{\overset{\mathrm{}}{}}a_j\eta _{n+j}+i\frac{t_n^1t_n}{t_n^1+t_n}\eta _n=0.$$
Since $`t_n^2=\mathrm{\Gamma }(z)`$, the second term equals
$$i\frac{1\mathrm{\Gamma }(\omega ^n(z))}{1+\mathrm{\Gamma }(\omega ^n(z))}.$$
Remark: If we replace $`\xi `$ by $`\stackrel{~}{\xi }=D_1\xi `$ in Proposition 3.2, then
$$h(G(\iota (z))\stackrel{~}{\xi }=p(z)\stackrel{~}{\xi }.$$
But if we transform $`\stackrel{~}{\xi }`$ analogous to $`\xi `$, thus obtaining a vector $`\stackrel{~}{\eta }`$, we have $`\eta =\stackrel{~}{\eta }`$. Succinctly put, one can say that to every non-zero element in the kernel of $`H_\gamma (z)`$ there correspond eigenvectors of $`h(G(z))`$ and $`h(G(\iota (z))`$, respectively, whose eigenvalues differ by a negative sign. Also note that $`H(\iota (z))=H(z)`$, since $`\mathrm{\Gamma }(\iota (z))=\mathrm{\Gamma }(z)`$. This means that the operators $`H(z)`$ are more appropriately parametrized through the Riemann surface obtained from $`\stackrel{~}{}`$ by identifying $`z`$ and $`\iota (z)`$.
Propositions 3.1 and 3.2 taken together say that $`H(z)`$ is a Fredholm operator of index zero with a non-trivial kernel and that $`0`$ is an isolated point in the spectrum of $`H_\gamma (z)`$ for every $`z\stackrel{~}{}\backslash \mathrm{\Omega }`$. This situation provides the proper setting for the employment of analytic perturbation theory of linear operators as expounded in \[Ko\], for instance.
First we enlarge $`\mathrm{\Omega }`$ by the set of points $`z\stackrel{~}{}`$ for which the following is true: For every neighborhood $`𝒰\stackrel{~}{}`$ of $`z`$ and for every neighborhood $`𝒱\text{}`$ of $`0`$, there exists $`\stackrel{~}{z}𝒰`$ such that $`H_\gamma (\stackrel{~}{z})`$ has a non-zero eigenvalue in $`𝒱`$. Since $`H_\gamma (z)`$ is analytic in $`z`$, the inclusion of those branch-points in $`\mathrm{\Omega }`$ still yields a discrete subset of $`\stackrel{~}{}`$.
For every point $`z\stackrel{~}{}\backslash \mathrm{\Omega }`$, let
$$P(z)=\frac{1}{2\pi i}\left(tH_\gamma (z)\right)^1𝑑t,$$
where the integral is taken over a positively-oriented circle enclosing $`0`$ but no other point in the spectrum of $`H_\gamma (z)`$. Then $`P(z)`$ is a projection of finite rank. Moreover, $`P(z)`$ is analytic in $`z`$ and the rank of $`P(z)`$ is constant. For every $`z\stackrel{~}{}\backslash \mathrm{\Omega }`$, let
$$N(z)=H_\gamma (z)P(z).$$
Then $`N(z)`$ is a nilpotent operator which is also analytic in $`z`$.
Our next goal is to show that the kernel of $`H_\gamma (z)`$ is one-dimensional for every$`z\stackrel{~}{}\backslash \mathrm{\Omega }`$. In preparation of this, we settle a number of technical questions first.
3.3 Lemma. There exists a discrete subset $`\mathrm{\Omega }_0\stackrel{~}{}`$ such that the following holds true: If $`z\stackrel{~}{}\backslash \mathrm{\Omega }_0`$ and $`\mathrm{}^2(\text{})`$ is a finite-dimensional subspace which is invariant under $`h(G(z))`$, then $``$ contains a linear basis of eigenvectors of $`h(G(z))`$.
Proof: Let $`\xi ^{(z)}`$ and $`\vartheta _p`$ be as in Theorem 2.16. Let
$$\mathrm{\Omega }_0=\{z\stackrel{~}{}\text{There exists }\stackrel{~}{z}G^1(\{p(z)\})\text{ such that }\vartheta _0(\stackrel{~}{z})=0\}.$$
Since $`\vartheta _0`$ is analytic, the set $`M_0=\{\stackrel{~}{z}\vartheta _0(\stackrel{~}{z})=0\}`$ is a discrete subset of $`\stackrel{~}{}`$. For every compact subset $`K\stackrel{~}{}`$, the set $`\{zKzG^1(\{p(\stackrel{~}{z})\})`$for some $`\stackrel{~}{z}M_0\}`$ is finite. Hence $`\mathrm{\Omega }_0`$ is discrete.
Now let $`z\stackrel{~}{}\backslash \mathrm{\Omega }_0`$ and consider the Jordan canonical form of $`h(G(z))`$ restricted on $``$. We need to show that the nilpotent components in this decomposition are trivial. To this end,we need to show that if $`\stackrel{~}{z}\stackrel{~}{}`$ such that $`p(\stackrel{~}{z})`$ is an eigenvalue for $`h(G(z))`$ with eigenvector $`\xi ^{(\stackrel{~}{z})}`$, then there does not exist $`\eta \mathrm{}^2(\text{})`$ such that $`(h(G(z))p(\stackrel{~}{z}))\eta =\xi ^{(\stackrel{~}{z})}`$. Suppose that the opposite is true: There exists $`\eta `$ with the said property. Then
$$\begin{array}{cc}\hfill \vartheta _0(z)& =\xi ^{(\stackrel{~}{z})},\overline{\xi ^{(\stackrel{~}{z})}}=\left(h(G(z))p(\stackrel{~}{z})\right)\eta ,\overline{\xi ^{(\stackrel{~}{z})}}\hfill \\ & =\eta ,\left(h(G(z))p(\stackrel{~}{z})\right)^{}\overline{\xi ^{(\stackrel{~}{z})}}\hfill \\ & =\eta ,\overline{\left(h(G(z))p(\stackrel{~}{z})\right)\xi ^{(\stackrel{~}{z})}}=0,\hfill \end{array}$$
which means that $`z\mathrm{\Omega }_0`$, thus contradicting our assumption on $`z`$.
We augment $`\mathrm{\Omega }`$, if necessary, to include the set $`\mathrm{\Omega }_0`$.
3.4 Corollary. Let $`z\stackrel{~}{}\backslash \mathrm{\Omega }`$ and for every $`n\text{}`$, let $`t_n`$ be chosen as in Proposition 3.2. Let
$$=\{\xi \xi _n=(t_n+t_n^1)^1\eta _n,\text{where}\eta \mathrm{}^2(\text{}),H_\gamma (z)\eta =0\}.$$
Then $``$ contains a linear basis consisting of eigenvectors of $`h_\gamma (G(z))`$.
Proof: Since the kernel of $`H_\gamma (z)`$ is finite-dimensional, in view of Lemma 3.3 all that needs to be shown is that $`h_\gamma (G(z))`$. Let $`\eta \mathrm{}^2(\text{})`$, $`H_\gamma (z)\eta =0`$ and let $`\xi _n=(t_n+t_n^1)\eta `$, $`\stackrel{~}{\xi }=h_\gamma (G(z))\xi `$, $`\stackrel{~}{\eta }_n=(t_n+t_n^1)\stackrel{~}{\xi }_n`$. Since $`|t_n||G_n(z)|^n`$ is constant,
$$\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}|t_n+t_n^1|^2|\xi _n|^2<\mathrm{}\text{implies}\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}|t_n+t_n^1|^2|\stackrel{~}{\xi }_n|^2<\mathrm{},$$
which implies $`\stackrel{~}{\eta }\mathrm{}^2(\text{})`$.
Next, since $`D_{G(z)}k_\gamma D_{G(z)}`$ and $`h_\gamma (G(z))`$ commute, and since $`D_{G(z)}k_\gamma D_{G(z)}\xi =\mathrm{\Gamma }(z)\xi `$,
$$D_{G(z)}k_\gamma D_{G(z)}\stackrel{~}{\xi }=\mathrm{\Gamma }(z)\stackrel{~}{\xi },$$
which in turn implies that $`H_\gamma (z)\stackrel{~}{\eta }=0`$. In conclusion, $`\stackrel{~}{\xi }=h_\gamma (G(z))\xi `$.
Since $`N`$ is analytic, the rank of $`N(z)`$ is constant on $`\stackrel{~}{}\backslash \mathrm{\Omega }`$ with the possible exception of a discrete subset. We now enlarge $`\mathrm{\Omega }`$ by this set of exceptional points.
3.5 Lemma. For every $`z\stackrel{~}{}\backslash \mathrm{\Omega }`$, let $`Q(z)`$ be the projection whose kernel equals the range of $`H_\gamma (z)`$ and whose range equals the kernel of $`H_\gamma (z)`$. Then $`Q`$ is analytic.
Proof: Let $`z_0\stackrel{~}{}\backslash \mathrm{\Omega }`$ and let $`𝒰\stackrel{~}{}\backslash \mathrm{\Omega }`$ be a simply connected open neighborhood of $`z_0`$. We shall shrink $`𝒰`$ when needed. By \[Ko\], II–§4.2 and VII–§1.3, there exists an analytic function $`T`$ from $`𝒰`$ into $`()`$ such that $`T(z)`$ is invertible for every $`z𝒰`$ and
$$T(z)P(z_0)T(z)^1=P(z).$$
Let
$$\stackrel{~}{N}(z)=T^1(z)N(z)T(z).$$
Then $`\stackrel{~}{N}`$ is analytic on $`𝒰`$ and $`\stackrel{~}{N}(z)`$ is nilpotent.
We may consider $`\stackrel{~}{N}(z)`$ as an operator on the finite-dimensional subspace $`=P(z_0)\mathrm{}^2(\text{})`$. We choose a linear basis in $``$, and we denote the matrix representation of $`\stackrel{~}{N}(z)`$ with respect to this basis once again by $`\stackrel{~}{N}(z)`$. Let $`m`$ be the dimension of $``$ and let $``$ be the canonical basis of $`\text{}^m`$. Then we choose a subset $`_1`$ of $``$ such that $`\{\stackrel{~}{N}(z_0)ee\}`$ is a linear basis for $`\stackrel{~}{N}(z_0)\text{}^m`$. Since the functions $`z\stackrel{~}{N}(z)e`$ are analytic in $`𝒰`$ for all $`e_1`$, $`_1(z)=\{\stackrel{~}{N}(z)ee_1\}`$ is a basis for $`\stackrel{~}{N}(z)\text{}^m`$ for every $`z`$ in an open neighborhood of $`z_0`$ contained in $`𝒰`$. We replace $`𝒰`$ by that neighborhood. Since $`\stackrel{~}{N}`$ is nothing but a matrix whose entries are analytic functions, we can choose a submatrix $`M`$ of $`\stackrel{~}{N}`$ of type $`r\times m`$, where $`r`$ equals the rank of $`\stackrel{~}{N}(z_0)`$, such that $`M(z_0)`$ has rank $`r`$. Again, since the entries of $`M`$ are analytic in $`z`$, $`r=\mathrm{rank}(\mathrm{M}(\mathrm{z}))=\mathrm{rank}(\stackrel{~}{\mathrm{N}}(\mathrm{z}))`$ for every $`z`$ in an open neighborhood of $`z_0`$ contained in $`𝒰`$. Once again we replace $`𝒰`$ by that neighborhood. We can now use Cramer’s rule to solve a system of linear equations for analytic functions $`f_1,\mathrm{},f_r`$ on $`𝒰`$ with values in $`\text{}^m`$ such that $`_2(z)=\{f_1(z),\mathrm{},f_r(z)\}`$ is a basis for the kernel of $`\stackrel{~}{N}(z)`$. Let $`\stackrel{~}{Q}(z)`$ be the projection whose range equals the kernel of $`\stackrel{~}{N}(z)`$ and whose kernel equals the range of $`\stackrel{~}{N}(z)`$. It follows that $`\stackrel{~}{Q}`$ is analytic on $`𝒰`$. By construction $`Q(z)=T(z)(\stackrel{~}{Q}(z)P(z_0))T(z)^1`$. In conclusion, $`Q`$ is locally analytic and hence analytic on $`\stackrel{~}{}\backslash \mathrm{\Omega }`$ as claimed.
3.6 Lemma. Let $`F`$ be a mapping from $`\backslash p(\mathrm{\Omega })`$ into the subsets of $``$ such that $`F(z)`$ contains exactly $`r`$ elements. Suppose that $`F`$ has the following properties:
(I) If $`zSp(h(\delta ))`$, then $`F(z)Sp(h(\delta ))`$.
(II) For every $`z_0`$, there exists an open neighborhood $`𝒰`$ of $`z_0`$ and for every $`j`$, $`1jr`$, there exists an analytic function $`s_j^{(z_0)}`$ on $`𝒰`$ such that
$$F(z)=\{s_1^{(z_0)}(z),\mathrm{},s_r^{(z_0)}(z)\}\text{for every }z𝒰.$$
Then $`F(z)\{z,z\}`$. (A mapping $`F`$ having the property (II) is called an analytic multifunction).
Proof: For every $`z\backslash p(\mathrm{\Omega })`$, we form the polynomial
$$\left(Xs_1^{(z)}(z)\right)\mathrm{}\left(Xs_r^{(z)}(z)\right)=X^r+\underset{j=0}{\overset{r1}{}}f_j(z)X^j.$$
Then (II) entails that the coefficients $`f_j`$ are analytic functions on $`\backslash p(\mathrm{\Omega })`$. By (I) the sets $`F(z)`$ are uniformly bounded in a punctured neighborhood of any point in $`p(\mathrm{\Omega })`$. It follows that the same is true for the functions $`f_j`$. Thus, all points in $`p(\mathrm{\Omega })`$ are removable singularities for $`f_j`$, and we may consider $`f_j`$ as an analytic function on $``$. It follows that the roots of the above polynomial are branches of analytic functions with exceptional point located in $`p(\mathrm{\Omega })`$ (see \[Ko\], II–§1.2).
Let $`d`$ be the smallest number such that the logarithmic potential associated with the equilibrium distribution $`\mu `$ of $`Sp(h)`$ has a critical point $`y`$ on the corresponding level curve. More precisely,
$$\mathrm{log}|yt|d\mu (t)=d,$$
while
$$(yt)^1𝑑\mu (t)=0.$$
By Corollary 2.13, we know that there always exists a critical point. Let
$$_d=\{z\text{}|\mathrm{log}|zt|d\mu (t)>d\}.$$
Then there exists an analytic function $`G_d`$ on $`_d`$ such that
$$\mathrm{log}|G_d(z)|=\mathrm{log}|zt|d\mu (t)$$
and
$$G_d(z)\text{}^+\text{whenever}z\text{}^+.$$
Also, $`G_d`$ has a simple pole at infinity. Let $`\overline{}_d`$ be the closure of $`_d`$ and let $`M`$ be the set of critical points in $`\overline{}_d`$. Then $`G_d`$ has a continuous extension on $`\overline{}_d\backslash M`$ and it approaches two distinct points as $`z`$ approaches a point in $`M`$. Let $`z_0_d\backslash p(\mathrm{\Omega })`$. Then (I) entails that
$$\begin{array}{cc}& s_j^{(z_0)}(z)_d,\hfill \\ & \left|G_d\left(s_j^{(z_0)}(z)\right)\right|=\left|G_d(z)\right|\hfill \end{array}\}\text{in a neighborhood }U\text{ of }z_0$$
for $`1jr`$.
It follows that there exist $`c_j\text{𝕋}`$ such that
$$G_ds_j^{(z_0)}=c_jG_d\text{on}𝒰,1jr.$$
Since this holds true for every $`z_0_d\backslash p(\mathrm{\Omega })`$, (II) implies that the set of those scaling factors $`c_j`$ is the same for all points in $`_d\backslash p(\mathrm{\Omega })`$. Hence
$$G_dF(z)=\{c_1,\mathrm{},c_r\}G_d(z),$$
for $`z_d\backslash p(\mathrm{\Omega })`$, where $`G_dF`$ denotes the multifunction obtained by applying $`G_d`$ to every element in $`F(z)`$. Since $`G_d`$ is a conformal map from $`_d\{\mathrm{}\}`$ onto $`\{z||z|>e^d\}\{\mathrm{}\}`$, it follows that there exist analytic functions $`s_1,\mathrm{},s_r`$ on $`_d`$ such that
$$F(z)=\{s_1(z),\mathrm{},s_r(z)\},z_d\backslash p(\mathrm{\Omega }).$$
Every $`s_j`$ is transformed via $`G_d`$ into a rotation by $`c_j`$ on $`\{z|z|>e^d\}`$. Moreover, the analytic continuation of $`s_j`$ across the boundary of $`\overline{}_d`$ does not result in the occurrence of exceptional points (branching-point) in $`\overline{}_d\backslash _d`$ or points of order larger than one. For, the existence of such points would conflict with the fact that all the critical points in $`M`$ are simple. It follows that $`s_j`$ has a differentiable extension on $`\overline{}_d`$ with non-vanishing derivatives. So, $`s_j`$ has an inverse transformation with the same properties. Let $`𝒦`$ be the group of transformations on $`\overline{}_d`$ which is generated by $`s_1,\mathrm{},s_r`$ and their inverses. Since the derivative $`G_d^{}(z)`$ approaches zero as $`z`$ approaches a point in $`M`$, it follows that every $`s𝒦`$ maps $`M`$ into $`M`$. Moreover, if $`zM`$ is a fixed point for some transformation in $`𝒦`$, that transformation must be the identity. Hence, since $`M`$ is finite, $`𝒦`$ is finite too. In particular, the group of scaling factors
$$𝒞=\{c\text{𝕋}|G_ds=cG_d\text{for some }s𝒦\}$$
is finite, which in turn means that $`c^m=1`$ for every $`c𝒞`$, where $`m`$ is the cardinality of $`𝒞`$. We want to show that $`𝒞\{1,1\}`$.
To this end, we first note that $`G_d(z)G(p^1(z))`$, if we let $`G_d(z)`$ be the set of the two limit points of $`G_d`$ at $`z`$ in case $`z`$ is in $`M`$. So, let $`zM`$. Since $`z`$ is real, since $`G\sigma =\overline{G}`$, and since $`\sigma \omega =\omega ^1\sigma `$, there exists a point $`\stackrel{~}{z}p^1(z)`$ such that either
$$\sigma (\stackrel{~}{z})=\stackrel{~}{z}\mathrm{or}\sigma (\stackrel{~}{z})=\omega (z).$$
In either case, $`G(p^1(z))\{\lambda ^nn\text{}\}\{\lambda ^nn\text{}\}`$. Now let $`s𝒦`$ and let $`z^{}=s(z)`$. Then
$$G_d(z),G_d(z^{})\{\lambda ^n|n\text{})\{\lambda ^n|n\text{}\}.$$
Let $`G_ds=cG_d`$ for some $`c𝒞`$. Since $`c^m=1`$ on the one hand, while $`\lambda `$ is non-periodic on the other hand, $`c\{1,1\}`$ as claimed. Finally, since $`z`$ whenever $`z`$, we conclude that $`r2`$ and $`F(z)\{z,z\}`$ for every $`z`$, by virtue of analytic continuation.
3.7 Proposition. The kernel of $`H_\gamma (z)`$is one-dimensional for every $`z\stackrel{~}{}\backslash \mathrm{\Omega }`$.
Proof: For every $`z\stackrel{~}{}\backslash \mathrm{\Omega }`$ choose $`t_n(z)`$ as in Proposition 3.2, that is,
$$H_\gamma (z)\eta =0\text{if and only if}k_\gamma D_{G(z)}\xi =\mathrm{\Gamma }(z)D_{G(z)^1}\xi ,$$
whenever $`\eta \mathrm{}^2(\text{})`$ and $`\xi _n=(t_n(z)+t_n(z)^1)^1\eta _n`$.
Let $`\stackrel{~}{h}_\gamma (G(z))`$ be defined as follows:
$$\left(\stackrel{~}{h}_\gamma (G(z))\eta \right)_n=s_n(z)\gamma \eta _{n+1}+s_n(z)^1\gamma ^1\eta _{n1}+\beta \left(G(z)\lambda ^n+G(z)^1\lambda ^n\right)\eta _n,$$
where
$$s_n(z)=\left(t_{n+1}(z)+t_{n+1}(z)^1\right)\left(t_n(z)+t_n(z)^1\right)^1.$$
Then $`\stackrel{~}{h}_\gamma (G(z))`$ is a bounded operator. Since $`\mathrm{\Gamma }`$ is analytic, $`\stackrel{~}{h}_\gamma (G(z))`$ is analytic in $`z`$. Let
$$B(z)=Q(z)\stackrel{~}{h}_\gamma (G(z))Q(z),z\stackrel{~}{}\backslash \mathrm{\Omega }.$$
Then $`B(z)`$ is a finite rank operator and $`B`$ is analytic by Lemma 3.5. By Corollary 3.4 the cardinality of the spectrum of $`B(z)`$ equals the rank of $`B(z)`$. The rank of $`B(z)`$, however, equals the dimension of the kernel of $`H_\gamma (z)`$. So, in order to settle the claim it has to be shown that the spectrum of $`B(z)`$ contains exactly one point. Since $`Sp(B(\omega (z))=Sp(B(z))`$, we can define a multi-function $`F`$ on $`\backslash p(\mathrm{\Omega })`$,
$$F(z)=Sp(B(\stackrel{~}{z})),p(\stackrel{~}{z})=z.$$
Then $`F`$ enjoys the properties stated in Lemma 3.6. It follows that $`F(z)\{z,z\}`$ for every $`z\backslash p(\mathrm{\Omega })`$. Since $`F(p(\stackrel{~}{z}))Sp(h_\gamma (G(\stackrel{~}{z}))`$ for every $`\stackrel{~}{z}\stackrel{~}{}\backslash \mathrm{\Omega }`$, and since $`z`$ cannot be an eigenvalue of $`h_\gamma (G(\stackrel{~}{z}))`$ if $`z`$ is one, we conclude that $`Sp(B(\stackrel{~}{z}))`$ contains exactly one element, as claimed. (Compare this with the remark following Proposition 3.2.)
As an immediate consequence of Proposition 3.7, we obtain the following corollary, announced in the introductory remarks to this paragraph:
3.8 Corollary. If $`z_1.z_2\stackrel{~}{}\backslash \mathrm{\Omega }`$ have the property that $`(G(z_1),\mathrm{\Gamma }(z_1))=(G(z_2),\mathrm{\Gamma }(z_2))`$, then $`z_1=z_2`$.
4. The case $`|\delta |=1`$ In this final paragraph we will describe how certain features from the preceding two paragraphs carry over to the operator $`h=h(1)`$.
First, $`h`$ is a fixed point of the automorphism $`\rho _\beta `$, and therefore the operators $`h`$ and $`k`$ commute. This information can be used to show the following:
$$\text{The spectrum of }h_\gamma =h_\gamma (1)\text{ is constant for }\beta ^1|\gamma |\beta .$$
$`(4.1)`$
Moreover, if $`\xi `$ is a solution of
$$h(x)\xi =\chi \xi ;\chi Sp(h),x\text{𝕋},$$
and $`|\xi _n|`$ grows moderately as $`|n|\mathrm{}`$, then
$$D_xkD_x\xi =c\xi $$
for some complex number $`c`$. Hence
$$D_{\overline{x}}k^1D_{\overline{x}}\overline{\xi }=\overline{D_xkD_x\xi }=\overline{c}\overline{\xi },$$
or
$$D_xkD_x\overline{\xi }=\overline{c}^1\overline{\xi }.$$
So, if $`\xi `$ is real, then $`|c|=1`$. This information, in conjunction with (4.1), can be used to show the following:
$$Sp(D_xk_\gamma D_x)=\text{𝕋}\text{for }\beta ^1|\gamma |\beta \text{ and }x\text{𝕋}.$$
$`(4.2)`$
In analogy to the operators $`H_\gamma (z)`$ considered in paragraph 3, we define for $`\gamma \text{}`$ close to $`1`$ and $`\theta ,\nu \text{}`$;
$$\left(H_\gamma (\theta ,\nu )\eta \right)_n=\underset{j=\mathrm{}}{\overset{\mathrm{}}{}}a_j\eta _{n+j}+\mathrm{tan}\pi (\alpha n^2+2\theta n+\nu )\eta _n,$$
where $`|\eta _n|`$ grows moderately as $`|n|\mathrm{}`$, and
$$\mathrm{tan}g(\gamma u)=\underset{j=\mathrm{}}{\overset{\mathrm{}}{}}a_ju^j,$$
with $`g`$ as in (1.6). Then $`H_\gamma (\theta ,\nu )`$ determines an unbounded operator on $`\mathrm{}^2(\text{})`$ with a dense domain, as long as
$$(\theta ,\nu )\mathrm{\Omega }_0=\left\{(\stackrel{~}{\theta },\stackrel{~}{\nu })\right|\alpha n^2+2\stackrel{~}{\theta }n+\stackrel{~}{\nu }\frac{1}{2}\text{}\}.$$
As $`\mathrm{\Gamma }(z)`$ and $`G(z)`$ approach suitable complex numbers of modulus one, respectively, $`H_\gamma (z)`$ is seen to approach an operator of this type, on the linear subspace of all vectors with finitely many non-vanishing components only, say.
As in Proposition 3.2, one can now show the following: Let $`\xi \mathrm{}^2(\text{})`$ be an eigenvector of $`h_\gamma (x)`$ and suppose that
$$D_xk_\gamma D_x\xi =c\xi .$$
Let
$$\eta _n=\mathrm{cos}\pi (\alpha n^2+2\theta n+\nu )\xi _n,$$
where $`e^{2\pi \theta i}=x,e^{2\pi \nu i}=c`$. Then
$$H_\gamma (\theta ,\nu )\eta =0,$$
in case $`(\theta ,\nu )\mathrm{\Omega }_0`$. Thus, information about the eigenvalue problem for the operator $`h_\gamma (\theta )`$ can be transformed into information regarding the kernel of an operator of the form $`H_\gamma (\theta ,\nu )`$, and to some degree, this works in the opposite direction as well.
If $`\beta `$ is sufficiently large, then we may choose $`\gamma =1`$. In this case, $`H(\theta ,\nu )=H_1(\theta ,\nu )`$ becomes an essentially self-adjoint operator. A complete analysis of the spectral properties of operators of this type can be obtained in case $`\alpha =0`$ and $`\theta `$ is an irrational number satisfying a stronger diophantine condition than the one that has been assumed to be valid for $`\alpha `$ in this paper (cf. \[PF\], Chapter VII, §18).
Appendix Two items related to a system of difference equations which have been used explicitly in the text will be assembled. For more details, we refer to \[R1\] and \[R2\].
A1. Consider the system
$$\{\begin{array}{cc}\hfill \mathrm{cos}(\pi \alpha q+\theta )(X_{p1,q}& +X_{p+1,q})+\beta \mathrm{cos}(\pi \alpha p+\theta )(X_{p,q1}+X_{p,q+1})=zX_{pq}\hfill \\ \hfill \mathrm{sin}(\pi \alpha q+\theta )(X_{p1,q}& X_{p+1,q})\beta \mathrm{sin}(\pi \alpha p+\theta )(X_{p,q1}X_{p,q+1})=0.\hfill \end{array}$$
$`(A1.1)`$
As in \[R1\], one can construct a recursion
$$\left(\begin{array}{c}X_{p+2,p+2}\\ X_{p+2,p+1}\\ X_{p+1,p+1}\end{array}\right)=F_p\left(\begin{array}{c}X_{p+1,p+1}\\ X_{p+1,p}\\ X_{pp}\end{array}\right),p0,$$
where $`\{X_{pq}\}`$ is any solution of (A1.1) and
$$F_p=E_pD_pC_p,$$
$$\begin{array}{cc}\hfill C_p& =\left[\begin{array}{ccc}\frac{\beta \mathrm{sin}[2\pi \alpha (p+1)+2\theta ]}{\mathrm{sin}\pi \alpha }& \frac{\chi \mathrm{sin}[\pi \alpha (p+1)+\theta ]}{\mathrm{sin}\pi \alpha }& \frac{\mathrm{sin}[\pi \alpha (2p+1)+2\theta ]}{\mathrm{sin}\pi \alpha }\\ & \\ 1& 0& 0\\ & \\ 0& 1& 0\end{array}\right]\hfill \\ \multicolumn{2}{c}{}\\ \hfill D_p& =\left[\begin{array}{ccc}0& \frac{\chi }{2\mathrm{cos}[\pi \alpha (p+1)+\theta ]}& \beta \\ & \\ 0& 1& 0\\ & \\ 1& 0& 0\end{array}\right]\hfill \\ \multicolumn{2}{c}{}\\ \hfill E_p& =\left[\begin{array}{ccc}\frac{\chi \mathrm{sin}[\pi \alpha (p+1)+\theta ]}{\beta \mathrm{sin}[\pi \alpha (2p+3)+2\theta ]}& \frac{\mathrm{sin}[2\pi \alpha (p+1)+2\theta ]}{\beta \mathrm{sin}[\pi \alpha (2p+3)+2\theta ]}& \frac{\mathrm{sin}\pi \alpha }{\mathrm{sin}[\pi \alpha (2p+3)+2\theta ]}\\ & \\ 1& 0& 0\\ & \\ 0& 1& 0\end{array}\right].\hfill \end{array}$$
If $`\{X_{pq}\}`$ and $`\{Y_{pq}\}`$ are solutions of (A1.1), then
$$X_{p+1,p+1}Y_{p+2,p+2}X_{p+2,p+2}Y_{p+1,p+1}=\beta \frac{\mathrm{sin}(\pi \alpha +\theta )}{\mathrm{sin}[\pi \alpha (2p+3)+2\theta ]}\left(X_{00}Y_{11}X_{11}Y_{00}\right).$$
Now suppose $`\theta \pi \alpha \text{}`$. Then the following holds true: If $`(X_{00},X_{11})(0,0)`$, $`(Y_{00},Y_{11})(0,0)`$ and the sequence
$$\{X_{pp}Y_{p+1,p+1}X_{p+1,p+1}Y_{pp}\}_{p0}$$
is bounded, while $`\{\underset{p\text{}}{sup}|X_{pq}|\}_{q0}`$ and $`\{\underset{p\text{}}{sup}|Y_{pq}|\}_{q0}`$ are bounded as well, then $`\{X_{pq}\}`$ and $`\{Y_{pq}\}`$ must be linearly dependent.
A2. Now consider the system (A1.1) for $`\theta =0`$. Let
$$(hz)^1=\underset{p,q\text{}}{}c_{pq}(z)w_{pq},z=\text{}\backslash Sp(h).$$
Then $`c_{pq}`$ is analytic and $`\{c_{pq}(z)\}`$ solves the system (A1.1), except for $`p=q=0`$.
By \[R2\], Paragraph 4, there exist polynomials $`d_{pq}`$ which are either zero or of degree $`||p||q||1`$ such that $`\{d_{pq}(z)\}`$ solves the system (A1.1) except for $`p=q=0`$ and
$$d_{pq}(z)=0\text{for}q|p|;d_{p,p1}(z)=(1)^p\beta ^{p1}\text{for}p0.$$
If $`\{X_{pq}\}`$ is a solution of (A1.1) for some $`z`$ which decays exponentially uniformly in $`p`$ as $`q\mathrm{}`$, then $`\{X_{pq}\}`$ and $`\{c_{pq}(z)d_{pq}(z)\}`$ must be linearly independent.
References
\[AS\] L. V. Ahlfors, L. Sario, “Riemann Surfaces”, Princeton Mathematical Series, 1976.
\[Ko\] T. Kato, “Perturbation Theory for Linear Operators” Second Edition, Springer-Verlag 1976.
\[Kn\] Y. Katznelson, “An Introduction to Harmonic Analysis”, Second Corrected Edition, Dover 1976.
\[PF\] L. Pastur, A. Figotin, “Spectra of Random and Almost-Periodic Operators”, Springer-Verlag 1992.
\[R1\] N. Riedel, “Almost Mathieu operators and rotation $`C^{}`$–algebras”, Proc. London Math. Soc. 3 (56), (1988), 281–302.
\[R2\] N. Riedel, “The spectrum of a class of almost periodic operators”, submitted since July 1992.
\[R3\] N. Riedel, “Regularity of the spectrum for the almost Mathieu operator” (September 1992), Proc. Amer. Math. Soc., to appear.
Department of Mathematics
Tulane University
e-mail: nriedel@mailhost.tcs.tulane.edu
|
warning/0003/math0003116.html
|
ar5iv
|
text
|
# The generalized Chern character and Lefschetz numbers in W*-modules
## 1 Introduction
For an arbitrary von Neumann algebra $`A`$ we introduce an abelian group $`N_0(A)`$ in the following way. It is possible to define some equivalence relation between normal elements of the inductive limit $`M_{\mathrm{}}(A)=\underset{}{lim}M_n(A)`$ such that for projections it coinside with the usual stable equivalence relation. Then the set of all equivalence classes of normal elements from $`M_{\mathrm{}}(A)`$ is an abelian semigroup (with respect to the direct sum operation) and $`N_0(A)`$ is its symmetrization. The first part of our paper is devoted to the consideration of some propeties of $`N`$-groups. More detail on this subject can be found in .
Further, we introduce Banach cyclic homology of $`A`$ as some analogue of usual cyclic homology and construct the generalized Chern character as a map from $`N_0(A)`$ to even Banach cyclic homology. Furthermore, this map is an extension of the classic Chern character to the group $`N_0(A)K_0(A)`$ in some natural sence.
In the final section we define generalized Lefschetz numbers for an arbitrary unitary endomorphism $`U`$ of an $`A`$-elliptic complex. Besides, in the case when $`U`$ is an element of a representation of some compact Lie group we describe the connection of the generalized Lefschetz numbers with the W\*-Lefschetz numbers of the first and of the second types introduced in .
Some results of the present paper were formulated in .
The paper is organized as follows:
| 1. | Introduction |
| --- | --- |
| 2. | Some properties of $`N`$-groups |
| 3. | The group $`N_0(A)_{fin}`$ |
| 4. | Banach cyclic homology |
| 5. | The generalized Chern character |
| 6. | Generalized Lefschetz numbers |
| | References |
## 2 Some properties of $`N`$-groups
Suppose $`A`$ is a von Neumann algebra, $`M_r(A)`$ is the set of $`r\times r`$ matrices with entries in $`A`$, $`M_{\mathrm{}}(A)`$ is the inductive limit of the sequence $`\{M_r(A)\}_{r=1}^{\mathrm{}}`$, and $`M_{\mathrm{}}(A)_\nu `$ is the set of normal elements for $`M_{\mathrm{}}(A)`$. Denote by $`(𝐂)`$ the family of all Borel subsets of the complex plane. If $`aM_{\mathrm{}}(A)_\nu `$ and $`E(𝐂)`$, then by $`P_a(E)`$ we denote the spectral projection of $`a`$ corresponding to the set $`E`$. We remark that $`P_a(E)M_{\mathrm{}}(A)_\nu `$ since von Neumann algebras are closed with respect to the Borel calculus. We denote the stable equivalence relation (see ) of projections $`p,qM_{\mathrm{}}(A)`$ by $`pq`$. Finally, let $`sp(a)`$ denote the spectrum of an element $`a`$.
A Borel set $`E𝐂`$ is called admissible if zero does not belong to the closure of $`E`$. Denote by $`_{}(𝐂)`$ the family of all admissible Borel subsets of the complex plane.
###### Definition 1
Call elements $`a,bM_{\mathrm{}}(A)`$ equivalent (and denote by $`ab`$) if and only if $`P_a(E)P_b(E)`$ for all $`E_{}(𝐂)`$.
Note that this equivalence relation coinsides with the usual stable equivalence relation whenever $`a,b`$ are projections. It is easy to see that $`a0_ma`$, where $`0_m`$ is the zero $`m\times m`$ matrix and $`aM_{\mathrm{}}(A)_\nu `$. We put
$$𝒩(A)=M_{\mathrm{}}(A)_\nu /.$$
For $`aM_{\mathrm{}}(A)_\nu `$ let us denote by $`[a]`$ the equivalence class of $`a`$ in $`𝒩(A)`$. Since
$$P_{ab}(E)=P_a(E)P_b(E)$$
for all $`a,bM_{\mathrm{}}(A)_\nu ,E(𝐂)`$, this implies that the set $`𝒩(A)`$ is an abelian semigroup with respect to the direct sum operation.
###### Definition 2
The symmetrization of $`𝒩(A)`$ is called the $`N`$-group of $`A`$ and is denoted by $`N_0(A)`$.
Under the previous considerations, the following result is clear.
###### Proposition 1
$`K_0(A)`$ is a subgroup of the group $`N_0(A)`$. $`\mathrm{}`$
###### Proposition 2
If the group $`K_0(A)`$ is trivial, then the group $`N_0(A)`$ is trivial too.
Proof. Suppose $`[a][b]N_0(A)`$. Since $`K_0(A)`$ is trivial, this implies that $`P_a(E)0P_b(E)`$ for all $`E_{}(𝐂)`$. Whence, $`ab`$. Thus the group $`N_0(A)`$ is trivial. $`\mathrm{}`$
Note that $`𝒩(A)`$ is a cancellation semigroup, i.e., the condition $`[a]+[c]=[b]+[c]`$ implies $`[a]=[b]`$ for any $`[a],[b],[c]𝒩(A)`$. In particular, the symmetrization homomorphism $`s:𝒩(A)N_0(A),s([a])=[a][0]`$ is injective.
Our next aim is to establish a functorial property for $`N_0`$. We recall that an arbitrary \*-homomorphism of C\*-algebras is a contraction, i.e., its norm does not exeed $`1`$ \[9, Theorem 2.1.7\]. Besides, an arbitrary surjective \*-homomorphism of von Neumann algebras is continuous with respect to the ultra-strong topology \[1, Theorem 2.4.23\].
Now let $`A,B`$ be von Neumann algebras and $`\phi :AB`$ an ultra-strong continuous unital \*-homomorphism. By definition, put $`\phi (a)=(\phi (a_{ij}))`$ and $`\phi _{}([a])=[\phi (a)]`$ for each matrix $`a=(a_{ij})M_{\mathrm{}}(A)_\nu `$.
###### Theorem 1
The map $`\phi _{}:N_0(A)N_0(B)`$ is a well defined homomorphism of abelian groups.
The following lemma is the main ingredient of the proof of Theorem 1.
###### Lemma 1
$`P_{\phi (a)}(E)=\phi (P_a(E))`$ for each $`aM_{\mathrm{}}(A)_\nu `$ and for each Borel subset $`Esp(a)`$.
Proof. We can assume that $`aM_r(A)`$ for some $`r1`$. It is clear that $`sp(\phi (a))sp(a)`$. It can be directly verified that $`\phi (R(a))=R(\phi (a))`$ for an arbitrary polynomial $`R`$. Further, for any function $`f`$, which is continuous on $`sp(a)`$, there exists a sequence of polinomials $`\{R_n\}_{n=1}^{\mathrm{}}`$ such that it converges uniformly to $`f`$. Then
$$\begin{array}{ccc}\phi (f(a))& =& \phi (lim_nR_n(a))=lim_n\phi (R_n(a))\hfill \\ & =& lim_nR_n(\phi (a))=f(\phi (a)).\hfill \end{array}$$
Now let $`\chi _E`$ be the characteristic function of $`E`$. Then we can find a sequence $`\{f_n\}_{n=1}^{\mathrm{}}`$ of continuous functions on the compact space $`sp(a)`$ such that $`\{f_n\}_{n=1}^{\mathrm{}}`$ converges to $`\chi _E`$ pointwise, i.e., with respect to the strong topology. Moreover, we can assume that the family $`\{f_n\}_{n=1}^{\mathrm{}}`$ is norm-bounded. In this case the sequence $`\{f_n(a)\}_{n=1}^{\mathrm{}}`$ of elements of the von Neumann algebra $`M_r(A)`$ converges strongly to the element $`\chi _E(a)=P_a(E)`$ and the family $`\{f_n(a)\}_{n=1}^{\mathrm{}}`$ is norm-bounded. Since strong and ultra-strong topologies coinside on bounded sets, we have the convergence with respect to the ultra-strong topology: $`P_a(E)=\sigma \text{-}lim_nf_n(a)`$. Finally, using the ultra-strong continuity of the \*-homomorphism $`\phi `$, we obtain:
$$\begin{array}{ccc}\hfill \phi (P_a(E))& =& \phi (\sigma \text{-}lim_nf_n(a))=\sigma \text{-}lim_n\phi (f_n(a))\hfill \\ & =& \sigma \text{-}lim_nf_n(\phi (a))=\chi _E(\phi (a))=P_{\phi (a)}(E).\mathrm{}\hfill \end{array}$$
Proof of Theorem 1. We have to establish that $`\phi _{}`$ is well defined. Let elements $`a,bM_{\mathrm{}}(A)_\nu `$ be equivalent. Then $`\phi (P_a(E))\phi (P_b(E))`$ for any $`E_{}(𝐂)`$. Now it immediately follows from Lemma 1 that the elements $`\phi (a)`$ and $`\phi (b)`$ are equivalent too. $`\mathrm{}`$
## 3 The group $`N_0(A)_{fin}`$
Let $`M_{\mathrm{}}(A)_{fin}M_{\mathrm{}}(A)_\nu `$ be the subset of all elements $`aM_{\mathrm{}}(A)`$ such that their spectrum is finite. Suppose,
$$𝒩(A)_{fin}=\{[a]:aM_{\mathrm{}}(A)_{fin}\}.$$
We remark that $`𝒩(A)_{fin}`$ is a subsemigroup of $`𝒩(A)`$. By $`N_0(A)_{fin}`$ we denote the symmetrization of the abelian monoid $`𝒩(A)_{fin}`$. So $`N_0(A)_{fin}`$ is a subgroup of the group $`N_0(A)`$. Any element of the group $`N_0(A)_{fin}`$ we can represent in the form $`[_{i=1}^n\lambda _ip_i][_{i=1}^n\lambda _iq_i],`$ where $`\lambda _i𝐂`$ and $`p_i,q_i`$ are some (possibly zero) projections in $`M_{\mathrm{}}(A)`$.
For an arbitrary map $`f:𝐂\{0\}K_0(A)`$ let us put $`\mathrm{\Lambda }_f=\{\lambda 𝐂\{0\}:f(\lambda )0\}`$. Let us denote by
$$_{fin;A}=_{fin}(𝐂\{0\},K_0(A))$$
the set of all maps from $`𝐂\{0\}`$ to $`K_0(A)`$ such that $`\mathrm{\Lambda }_f`$ is finite (or empty). The set $`_{fin;A}`$ is an abelian group with respect to the pointwise addition of maps. Let us consider the map
$$\varphi :N_0(A)_{fin}_{fin;A},(\varphi ([a][b]))(\lambda )=[P_a(\{\lambda \})][P_b(\{\lambda \})].$$
###### Theorem 2
The map $`\varphi `$ is an isomorphism of groups.
Proof. It is clear that $`\varphi `$ is a well defined homomorphism. Let $`[a][b]N_0(A)_{fin}`$, and $`\varphi ([a][b])=0`$. Therefore for each $`\lambda 𝐂\{0\}`$ the projections $`P_a(\{\lambda \})`$ and $`P_b(\{\lambda \})`$ are stably equivalent. Since the spectra of elements $`a,b`$ are finite, we conclude that these elements are equivalent. Hence, $`\varphi `$ is injective.
Let us examine a map $`f_{fin;A}`$ such that $`\mathrm{\Lambda }_f=\{\lambda _i\}_{i=1}^n`$ and $`f(\lambda _i)=[p_i][q_i]`$ $`(1in)`$, where $`p_i,q_i`$ are projections from $`M_{\mathrm{}}(A)`$. Since $`p_i0_kp_i`$, we can assume that the projections $`\{p_i\}_{i=1}^n`$ (and $`\{q_i\}_{i=1}^n`$) are pairwise orthogonal. We put $`a=_{i=1}^n\lambda _ip_i`$ and $`b=_{i=1}^n\lambda _iq_i`$. Then the elements $`a,b`$ belong to $`M_{\mathrm{}}(A)_{fin}`$. Furtermore, $`\varphi ([a][b])=f`$. Hence, $`\varphi `$ is surjective. $`\mathrm{}`$
###### Corollary 1
The groups $`N_0(M_r(𝐂))`$ and $`_{fin}(𝐂\{0\},𝐙)`$ are isomorphic.
Proof. The spectrum of any element from $`M_{\mathrm{}}(𝐂)`$ is finite. Therefore, $`N_0(M_r(𝐂))=N_{fin}(M_r(𝐂))`$. To complete the proof, it remains to use Theorem 2. $`\mathrm{}`$
Assume $`h_{\lambda ,\mu ;p}:=[p(\lambda +\mu )][p\lambda p\mu ]`$ and $`g_{\lambda ;p}^{(n)}:=[\lambda p^n][n\lambda p]`$, where $`n𝐍`$, $`\lambda ,\mu 𝐂`$ and $`p`$ is a projection in $`M_{\mathrm{}}(A)`$. Then $`h_{\lambda ,\mu ;p}`$ and $`g_{\lambda ;p}^{(n)}`$ are elements from $`N_0(A)_{fin}`$.
Let us denote by $`H`$ the subgroup of $`N_0(A)_{fin}`$ with the following system of generators $`\{h_{\lambda ,\mu ;p}:\lambda ,\mu 𝐂,p\text{ is a projection in }M_{\mathrm{}}(A)\},`$ and by $`G`$ the subgroup of $`N_0(A)_{fin}`$ with the following system of generators $`\{g_{\lambda ;p}^{(n)}:n𝐍,\lambda 𝐂,p\text{ is a projection in }M_{\mathrm{}}(A)\}.`$
###### Lemma 2
The group $`G`$ is a subgroup of $`H`$.
Proof. Take $`\lambda 𝐂`$ and a projection $`p`$ from $`M_{\mathrm{}}(A)`$. We have to demonstrate that $`g_{\lambda ;p}^{(k)}`$ belongs to $`H`$ for all $`k1`$. Let us prove this statement by induction over $`k`$. The case $`k=1`$ is clear. Suppose $`g_{\lambda ;p}^{(k)}H`$ for all $`kn1`$. In particular,
$$g_{\lambda ;p}^{(n1)}=[\lambda p^{(n1)}][(n1)\lambda p]=y$$
for some $`yH`$. Therefore,
$$\begin{array}{ccc}g_{\lambda ;p}^{(n)}& =& [\lambda p^{(n1)}]+[\lambda p][((n1)\lambda +\lambda )p]\hfill \\ & =& [(n1)\lambda p]+y+[\lambda p][((n1)\lambda +\lambda )p]\hfill \\ & =& yh_{(n1)\lambda ,\lambda ;p}.\hfill \end{array}$$
Thus $`g_{\lambda ;p}^{(n)}H`$ and by induction we obtain the desired statement. $`\mathrm{}`$
Let us consider the map
$$h:N_0(A)_{fin}K_0(A)𝐂$$
(1)
defined as follows
$$h([a][b])=\underset{i=1}{\overset{n}{}}[P_a(\lambda _i)]\lambda _i\underset{j=1}{\overset{m}{}}[P_b(\mu _j)]\mu _j,$$
where $`a,bM_{\mathrm{}}(A)_{fin}`$ and $`sp(a)=\{\lambda _1,\mathrm{},\lambda _n\}`$, $`sp(b)=\{\mu _1,\mathrm{},\mu _m\}`$.
###### Proposition 3
The map $`h`$ is a surjective homomorphism of groups. Besides, the kernel of $`h`$ coinsides with the group $`H`$.
Proof. It is obvious that $`h`$ is a well defined surjective homomorphism. Also, it is clear that $`H`$ belongs to the kernel of $`h`$. To complete the proof let us construct the inverse for $`h`$ homomorphism
$$t:K_0(A)𝐂N_0(A)_{fin}/H.$$
We put
$$t(\underset{i=1}{\overset{n}{}}([p_i][q_i])\lambda _i)=\underset{i=1}{\overset{n}{}}[p_i\lambda _i]\underset{i=1}{\overset{n}{}}[q_i\lambda _i]+H.$$
Let us demonstrate that the map $`t`$ is well defined. In the other words, we have to verify that the homomorphism $`t`$ is trivial on the elements $`c_{\lambda ,\mu ;p}^{(1)}=[p](\lambda +\mu )[p]\lambda [p]\mu `$, $`c_{z;\lambda ;p}^{(2)}=[p]z\lambda [p]z\lambda `$, and $`c_{\lambda ;p,q}^{(3)}=([p]+[q])\lambda [p]\lambda [q]\lambda `$, where $`\lambda ,\mu 𝐂,z𝐙`$ and $`p,qM_{\mathrm{}}(A)`$ are projections. We derive $`t(c_{\lambda ,\mu ;p}^{(1)})=h_{\lambda ,\mu ;p}+H=H`$. Besides, $`c_{1;\lambda ;p}^{(2)}=c_{\lambda ,\lambda ;p}^{(1)}`$. Therefore it suffices to regard the elements $`c_{z;\lambda ;p}^{(2)}`$ provided $`z>0`$. In this case we conclude $`t(c_{z;\lambda ;p}^{(2)})=t([p^z]\lambda [p]z\lambda )=g_{\lambda ;p}^{(z)}+H=H`$, where we have used Lemma 2. Finally, it can be directly verified that $`t(c_{\lambda ;p,q}^{(3)})=0+H=H`$. To complete the proof, it remains to note that $`t=h^1`$. $`\mathrm{}`$
## 4 Banach cyclic homology
As above, let $`A`$ be a von Neumann algebra. First let us recall some concepts from noncommutative geometry (see, for example, ). Consider the complex vector space $`C_n(A)=A^{(n+1)}`$, where $`A^{(n+1)}=\underset{n+1}{\underset{}{AA\mathrm{}A}}`$. The cyclic operator $`\tau _n:C_n(A)C_n(A)`$ is defined on generators by the formula
$$\tau _n(a_0\mathrm{}a_n)=(1)^na_na_0\mathrm{}a_{n1}.$$
The cokernel of the endomorphism $`1\tau _n:C_n(A)C_n(A)`$ we denote by
$$CC_n(A)=A^{(n+1)}/Im(1\tau _n).$$
Further, we define the face operator $`b_n:C_n(A)C_{n1}(A)`$ by the formula
$$b_n(a_0a_1\mathrm{}a_n)=\underset{i=0}{\overset{n1}{}}(1)^ia_0a_1\mathrm{}a_ia_{i+1}\mathrm{}a_n+$$
$$(1)^na_na_0a_1\mathrm{}a_{n1}.$$
It is clear that
$$b=\underset{i=0}{\overset{n}{}}(1)^id_i,$$
(2)
where the linear maps $`d_i:A^{(n+1)}A^n`$ are defined as follows
$$\begin{array}{ccc}d_i(a_0\mathrm{}a_n)\hfill & =& a_0\mathrm{}a_ia_{i+1}\mathrm{}a_n,\\ & & 0in1,\\ d_n(a_0\mathrm{}a_n)\hfill & =& a_na_0a_1\mathrm{}a_{n1}.\end{array}$$
It can be verified by the direct calculation that the family of linear spaces $`CC_{}(A)=\{CC_n(A),b_n\}`$ is a chain complex. Homology of this complex is called cyclic homology of $`A`$ and is denoted by $`HC_n(A)=H_n(CC_{}(A))`$, $`n0`$.
The trace map $`Tr:M_r(A)^{(n+1)}A^{(n+1)}`$ is defined by the formula
$$Tr(\xi ^{(0)}\xi ^{(1)}\mathrm{}\xi ^{(n)})=\underset{i_0,\mathrm{},i_n=1}{\overset{r}{}}\xi _{i_0i_1}^{(0)}\xi _{i_1i_2}^{(1)}\mathrm{}\xi _{i_ni_0}^{(n)},$$
where $`\xi ^{(k)}=(\xi _{i,j}^{(k)})_{i,j=1}^rM_r(A)`$. It can be directly verified that the trace map $`Tr:CC_{}(M_r(A))CC_{}(A)`$ is a morphism of chain complexes. Furthermore, the induced map
$$Tr_{}:HC_{}(M_r(A))\stackrel{}{}HC_{}(A)$$
is an isomorphism.
Now let $`X,Y`$ be normed spaces, and $`xX,yY`$. A representative of the equivalence class $`xyXY`$ we shall denote by $`x\mathrm{}y`$. Assume
$$\underset{i}{}x_i\mathrm{}y_i:=\underset{i}{}x_iy_i(x_iX,y_iY).$$
Then the projective norm of an equivalence class $`\xi XY`$ is defined as follows
$$\xi =inf\{\underset{i}{}x_i\mathrm{}y_i:x_iX,y_iY\text{and}\underset{i}{}x_i\mathrm{}y_i\xi \}.$$
Below we assume that all tensor products of normed spaces are equipped with the projective norm.
Under the previous conventions, let us set
$$𝒞𝒞_n(A)=A^{(n+1)}/\overline{Im(1\tau _n)}.$$
Note that $`𝒞𝒞_n(A)`$ is a Banach space. For $`\xi A^{(n+1)}`$ we denote by $`[\xi ]_{CC_n(A)}`$ the quotient class of $`\xi `$ in $`CC_n(A)`$, and by $`[\xi ]_{𝒞𝒞_n(A)}`$ the quotient class of $`\xi `$ in $`𝒞𝒞_n(A)`$.
###### Lemma 3
The face operator $`b_n:A^{(n+1)}A^n`$ is a continuous map.
Proof. Under equality (2), it is sufficiently to prove that all maps $`d_i`$ $`(0in)`$ are continuous. Given any $`\epsilon >0`$. For each $`\xi A^{(n+1)}`$ we can find an element $`_ka_0^{(k)}\mathrm{}\mathrm{}\mathrm{}a_n^{(k)}\xi `$ such that
$$\underset{k}{}a_0^{(k)}\mathrm{}\mathrm{}\mathrm{}a_n^{(k)}\xi <\epsilon .$$
Therefore,
$$\begin{array}{ccc}d_i(\xi )\hfill & =& _kd_i(a_0^{(k)}\mathrm{}a_n^{(k)})\hfill \\ & =& _ka_0^{(k)}\mathrm{}a_i^{(k)}a_{i+1}^{(k)}\mathrm{}a_n^{(k)}\hfill \\ & & _ka_0^{(k)}\mathrm{}\mathrm{}\mathrm{}a_i^{(k)}a_{i+1}^{(k)}\mathrm{}\mathrm{}\mathrm{}a_n^{(k)}\hfill \\ & =& _ka_0^{(k)}\mathrm{}a_i^{(k)}a_{i+1}^{(k)}\mathrm{}a_n^{(k)}\hfill \\ & & _ka_0^{(k)}\mathrm{}a_i^{(k)}a_{i+1}^{(k)}\mathrm{}a_n^{(k)}\hfill \\ & =& _ka_0^{(k)}\mathrm{}\mathrm{}\mathrm{}a_n^{(k)}<\xi +\epsilon \hfill \end{array}$$
for all $`\epsilon >0`$. Hence, $`d_i(\xi )\xi `$. $`\mathrm{}`$
Let us define the map $`\beta _n:𝒞𝒞_n(A)𝒞𝒞_{n1}(A)`$ by the formula
$$\beta _n([x]_{𝒞𝒞_n(A)})=[b_n(x)]_{𝒞𝒞_{n1}(A)}.$$
From Lemma 3 we conclude that the family $`𝒞𝒞_{}(A)=\{𝒞𝒞_n(A),\beta _n\}`$ is a well defined chain complex. Let us put
$$𝒞_n(A)=Ker\beta _n/\overline{Im\beta _{n+1}}.$$
The quotient space $`𝒞_{}(A)`$ we shall call Banach (cyclic) homology of $`A`$ (cf. ). Note that $`𝒞_n(A)`$ is a Banach space (for each $`n0`$).
For $`\xi M_r(A)^{(n+1)}`$ let us denote by $`\xi _{HC}=\xi _{HC_{}(M_r(A))}`$ the cyclic homology class of $`\xi `$ and by $`\xi _𝒞=\xi _{𝒞_{}(M_r(A))}`$ the Banach cyclic homology class of $`\xi `$.
Let $`pM_r(A)`$ be projection. Then
$$b_{2l}(p^{(2l+1)})=\underset{k=0}{\overset{2l}{}}d_k(p^{(2l+1)})=\underset{k=0}{\overset{2l}{}}(1)^kp^{2l}=p^{2l}.$$
On the other hand, we have $`[p^{2l}]_{CC_{}(M_r(A))}=(1)^{(2l1)}[p^{2l}]_{CC_{}(M_r(A))}`$. So $`[p^{2l}]_{CC_{}(M_r(A))}=0`$. Therefore $`p^{(2l+1)}`$ is a cycle. Now let us define the Chern character
$$Ch_{2l}^0:K_0(A)HC_{2l}(A)$$
by $`Ch_{2l}^0([p])=Tr_{}(1)^lp^{(2l+1)}_{HC}`$. The Chern character is a well defined homomorphism of groups \[6, Theorem 8.3.2\].
Let us study the linear epimorphism $`\pi _n:CC_n(A)𝒞𝒞_n(A)`$, where
$$\pi _n([\xi ]_{CC_n(A)})=[\xi ]_{𝒞𝒞_n(A)},\xi A^{(n+1)}.$$
It is obvious that the family of the maps $`\{\pi _n\}:CC_{}(A)𝒞𝒞_{}(A)`$ is a chain homomorphism, i.e., $`\pi _{n1}b_n=\beta _n\pi _n`$ for all $`n1`$. So $`\pi _n(Kerb_n)Ker\beta _n`$ and $`\pi _n(Imb_{n+1})\overline{Im\beta _{n+1}}`$. Therefore,
$$\pi _{}:HC_{}(A)𝒞_{}(A)$$
(3)
is a well defined map of homology spaces.
Below we shall need the following result.
###### Lemma 4
The trace $`Tr:M_r(A)^{(n+1)}A^{(n+1)}`$ is a continuous map.
Proof. Given any $`\epsilon >0`$. For each quotient class $`\xi M_r(A)^{(n+1)}`$ of the tensor product we can find a representative $`_{k=1}^N\xi ^{(0),k}\mathrm{}\xi ^{(1),k}\mathrm{}\mathrm{}\mathrm{}\xi ^{(n),k}`$ of $`\xi `$ such that
$$\underset{k=1}{\overset{N}{}}\xi ^{(0),k}\mathrm{}\xi ^{(1),k}\mathrm{}\mathrm{}\mathrm{}\xi ^{(n),k}\xi <\epsilon .$$
Therefore,
$$\begin{array}{ccc}Tr(\xi )\hfill & =& _{k=1}^NTr(\xi ^{(0),k}\xi ^{(1),k}\mathrm{}\xi ^{(n),k})\hfill \\ & =& _{k=1}^N_{i_0,\mathrm{},i_n=1}^r\xi _{i_0i_1}^{(0),k}\xi _{i_1i_2}^{(1),k}\mathrm{}\xi _{i_ni_0}^{(n),k}\hfill \\ & & _{k=1}^N_{i_0,\mathrm{},i_n=1}^r\xi _{i_0i_1}^{(0),k}\mathrm{}\xi _{i_1i_2}^{(1),k}\mathrm{}\mathrm{}\mathrm{}\xi _{i_ni_0}^{(n),k}\hfill \\ & =& _{k=1}^N_{i_0,\mathrm{},i_n=1}^r\xi _{i_0i_1}^{(0),k}\xi _{i_1i_2}^{(1),k}\mathrm{}\xi _{i_ni_0}^{(n),k}\hfill \\ & & _{k=1}^N_{i_0,\mathrm{},i_n=1}^r\xi ^{(0),k}\xi ^{(1),k}\mathrm{}\xi ^{(n),k}\hfill \\ & =& r^{n+1}_{k=1}^N\xi ^{(0),k}\mathrm{}\xi ^{(1),k}\mathrm{}\mathrm{}\mathrm{}\xi ^{(n),k}<r^{n+1}(\xi +\epsilon )\hfill \end{array}$$
for all $`\epsilon >0`$. So, $`Tr(\xi )r^{n+1}\xi `$. $`\mathrm{}`$
From Lemma 4 we conclude that the map
$$Tr:𝒞𝒞_{}(M_r(A))𝒞𝒞_{}(A),[\xi ]_{𝒞𝒞_{}(M_r(A))}[Tr(\xi )]_{𝒞𝒞_{}(A)}$$
(4)
is well defined. Furthermore, it can be directly checked up that trace map (4) is a morphism of chain complexes. By Lemma 4, this implies that the induced homomorphism
$$Tr_{}:𝒞_{}(M_r(A))𝒞_{}(A)$$
of Banach homology is well defined.
## 5 The generalized Chern character
Given $`l0`$. We want to define a map
$$T:M_{\mathrm{}}(A)_\nu 𝒞_{2l}(A)$$
to the even Banach homology by the following construction. Let an element $`a`$ belong to $`M_{\mathrm{}}(A)_\nu `$. Then we can suppose that $`aM_r(A)`$ for some $`r1`$. For each natural number $`n`$ let us consider a cover $`^{(n)}=\{E_k^{(n)}\}_{k=1}^{k_n}`$ of the spectrum of $`a`$ by disjoint Borel sets such that the diameter of each of these sets does not exceed $`1/n`$. Moreover, we can assume that $`\{E_k^{(m)}\}_{k=1}^{k_m}`$ is a subdivision of the cover $`\{E_k^{(n)}\}_{k=1}^{k_n}`$ when $`mn`$. Therefore we can write
$$E_k^{(n)}=\underset{j=1}{\overset{j(k)}{}}E_{k,j}^{(m)},$$
where $`E_{k,j}^{(m)}`$ are some elements of the cover $`^{(m)}`$. Thus, $`^{(m)}=\{E_{k,j}^{(m)}\}_{k=1;j=1}^{k_n;j(k)}`$ and $`_{k=1}^{k_n}j(k)=k_m`$.
Also, for any $`\lambda _k^{(n)}E_k^{(n)}`$ let us consider a sequence $`\{a_n\}_{n=1}^{\mathrm{}}`$ from $`M_r(A)`$, where
$$a_n=\underset{k=1}{\overset{k_n}{}}P_a(E_k^{(n)})\lambda _k^{(n)}.$$
It follows from the spectral theorem that $`\{a_n\}_{n=1}^{\mathrm{}}`$ converges uniformly to $`a`$. Further, for each $`a_n`$ let us examine
$$\stackrel{~}{a}_n=\underset{k=1}{\overset{k_n}{}}P_a(E_k^{(n)})^{(2l+1)}\lambda _k^{(n)}M_r(A)^{(2l+1)}.$$
Given natural numbers $`n,m`$ $`(mn)`$. Under the previous notation, we see that
$$a_n=\underset{k=1}{\overset{k_n}{}}P_a(\underset{j=1}{\overset{j(k)}{}}E_{k,j}^{(m)})\lambda _k^{(n)}=\underset{k=1}{\overset{k_n}{}}\underset{j=1}{\overset{j(k)}{}}P_a(E_{k,j}^{(m)})\lambda _k^{(n)}$$
and $`a_m=_{k=1}^{k_n}_{j=1}^{j(k)}P_a(E_{k,j}^{(m)})\lambda _{k,j}^{(m)}`$, where $`\lambda _{k,j}^{(m)}E_{k,j}^{(m)}`$. This yields that
$$a_na_m=\underset{k=1}{\overset{k_n}{}}\underset{j=1}{\overset{j(k)}{}}P_a(E_{k,j}^{(m)})(\lambda _k^{(n)}\lambda _{k,j}^{(m)}).$$
(5)
Furthermore,
$`\stackrel{~}{a}_n`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{k_n}{}}}P_a(E_k^{(n)})^{(2l+1)}\lambda _k^{(n)}={\displaystyle \underset{k=1}{\overset{k_n}{}}}P_a({\displaystyle \underset{j=1}{\overset{j(k)}{}}}E_{k,j}^{(m)})^{(2l+1)}\lambda _k^{(n)}`$
$`=`$ $`{\displaystyle \underset{k=1}{\overset{k_n}{}}}({\displaystyle \underset{j_0=1}{\overset{j(k)}{}}}P_a(E_{k,j_0}^{(m)}))\mathrm{}({\displaystyle \underset{j_{2l}=1}{\overset{j(k)}{}}}P_a(E_{k,j_{2l}}^{(m)}))\lambda _k^{(n)}`$
$`=`$ $`{\displaystyle \underset{k=1}{\overset{k_n}{}}}({\displaystyle \underset{j_0=1}{\overset{j(k)}{}}}\mathrm{}{\displaystyle \underset{j_{2l}=1}{\overset{j(k)}{}}}P_a(E_{k,j_0}^{(m)})\mathrm{}P_a(E_{k,j_{2l}}^{(m)}))\lambda _k^{(n)}`$
and $`\stackrel{~}{a}_m=_{k=1}^{k_n}_{j=1}^{j(k)}P_a(E_{k,j}^{(m)})^{(2l+1)}\lambda _{k,j}^{(m)}`$.
For brevity we shall use the following notation
$$\stackrel{~}{}_{j_0,\mathrm{},j_{2l}=1}^{j(k)}:=\underset{\{1j_0\mathrm{}j_{2l}j(k):j_pj_q\}}{}$$
for the sum over $`j_0,\mathrm{},j_{2l}`$ from one to $`j(k)`$, where not all indices coinside.
So we obtain
$`\stackrel{~}{a}_n\stackrel{~}{a}_m`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{k_n}{}}}{\displaystyle \underset{j=1}{\overset{j(k)}{}}}P_a(E_{k,j}^{(m)})^{(2l+1)}(\lambda _k^{(n)}\lambda _{k,j}^{(m)})`$
$`+`$ $`{\displaystyle \underset{k=1}{\overset{k_n}{}}}(\stackrel{~}{{\displaystyle }}_{j_0,\mathrm{},j_{2l}=1}^{j(k)}P_a(E_{k,j_0}^{(m)})\mathrm{}P_a(E_{k,j_{2l}}^{(m)}))\lambda _k^{(n)}`$
$`=`$ $`\alpha _{n,m}+\gamma _{n,m},`$
where by $`\alpha _{n,m}`$ $`(\gamma _{n,m})`$ we denote the first (the second) summand in expression (5).
The following result is the main ingredient of our definition of the map $`T`$.
###### Theorem 3
Let $`\{p_i\}_{i=1}^NM_r(A)`$ be a family of pairwise orthogonal projections and $`\eta =\stackrel{~}{}_{j_0,\mathrm{},j_{2l}=1}^Np_{j_0}\mathrm{}p_{j_{2l}}`$. Then the element $`\eta `$ belongs to the kernel of the face operator $`\beta _{2l}`$. Besides, $`\eta _{𝒞_{2l}(M_r(A))}=0`$.
Proof. We have
$$\eta =(\underset{j=1}{\overset{N}{}}p_j)^{(2l+1)}\underset{j=1}{\overset{N}{}}p_j^{(2l+1)}.$$
So $`\eta `$ belongs to the kernel of the face operator $`\beta _{2l}`$ as a difference of elements from $`Ker\beta _{2l}`$.
Now let us examine the following element
$$\begin{array}{ccc}\alpha & \hfill :=& Tr_{}(_{j=1}^Np_j)^{(2l+1)}_{HC}\hfill \\ & \hfill =& (1)^lCh_{2l}^0(_{j=1}^N[p_j])=_{j=1}^N(1)^lCh_{2l}^0([p_j])\hfill \\ & \hfill =& _{j=1}^NTr_{}p_j^{(2l+1)}_{HC}=Tr_{}_{j=1}^Np_j^{(2l+1)}_{HC}.\hfill \end{array}$$
On the other hand,
$$\alpha =Tr_{}\underset{j_0,\mathrm{},j_{2l}=1}{\overset{N}{}}p_{j_0}\mathrm{}p_{j_{2l}}_{HC}.$$
Thus,
$$\begin{array}{ccc}0& =& Tr_{}_{j_0,\mathrm{},j_{2l}=1}^Np_{j_0}\mathrm{}p_{j_{2l}}_{HC}Tr_{}_{j=1}^Np_j^{(2l+1)}_{HC}\hfill \\ & =& Tr_{}_{j_0,\mathrm{},j_{2l}=1}^Np_{j_0}\mathrm{}p_{j_{2l}}_{j=1}^Np_j^{(2l+1)}_{HC}\hfill \\ & =& Tr_{}\stackrel{~}{}_{j_0,\mathrm{},j_{2l}=1}^Np_{j_0}\mathrm{}p_{j_{2l}}_{HC}.\hfill \end{array}$$
The trace map is an isomorphism. Therefore we conclude
$$\eta _{HC}=\stackrel{~}{}_{j_0,\mathrm{},j_{2l}=1}^Np_{j_0}\mathrm{}p_{j_{2l}}_{HC}=0.$$
Whence, $`\eta _𝒞=\pi _{}(\eta _{HC})=0`$, where $`\pi _{}`$ is map (3). $`\mathrm{}`$
Now let us return to expressions (5),(5). By the previous theorem, we see that $`\gamma _{n,m}_𝒞=0`$ so $`\stackrel{~}{a}_n\stackrel{~}{a}_m_𝒞=\alpha _{n,m}_𝒞`$.
We claim that $`\alpha _{n,m}=a_na_m`$. Indeed, it is clear that elements $`a_na_m`$ and $`\alpha _{n,m}`$ are normal. Besides, $`sp(a_na_m)\{0\}`$ and $`sp(\alpha _{n,m})\{0\}`$ coinside with the set $`\{\lambda _k^{(n)}\lambda _{k,j}^{(m)}\}_{k=1;j=1}^{k_n;j(k)}`$. This implies that spectral radii of these elements coinside too. Thus we obtain the desired statement.
So we can write
$$\stackrel{~}{a}_n\stackrel{~}{a}_m_𝒞=\alpha _{n,m}_𝒞\alpha _{n,m}=a_na_m.$$
(7)
By Lemma 4 and inequality (7), we conclude that $`\{Tr_{}(\stackrel{~}{a}_n_𝒞)\}_{n=1}^{\mathrm{}}`$ is a Cauchy sequence. Therefore it converges to some element $`T(a;\{a_n\})𝒞_{2l}(A)`$.
It remains to verify that the limit $`T(a;\{a_n\})`$ does not depend on $`\{a_n\}`$. Let us regard covers $`^{(n)}=\{E_k^{(n)}\}_{k=1}^{k_n}`$ and $`^{(n)}=\{F_j^{(n)}\}_{j=1}^{j_n}`$ of the spectrum of $`a`$ by disjoint Borel sets. Besides, we shall assume that the diameter of each of these sets does not exceed $`1/n`$. Also, for any $`\mu _j^{(n)}F_j^{(n)}`$ let us examine a sequence $`\{c_n\}_{n=1}^{\mathrm{}}`$, where $`c_n=_{j=1}^{j_n}P_a(F_j^{(n)})\mu _j^{(n)}`$.
###### Theorem 4
The elements $`T(a;\{a_n\})`$ and $`T(a;\{c_n\})`$ coinside. Thus the map
$$T:M_{\mathrm{}}(A)_\nu 𝒞_{2l}(A),aT(a;\{a_n\})=T(a)$$
is well defined.
Proof. Let us denote $`X_{k,j}^{(n)}=E_k^{(n)}F_j^{(n)}`$. Then
$$\begin{array}{ccc}\hfill a_n=_{k=1}^{k_n}P_a(E_k^{(n)})\lambda _k^{(n)}& =& _{k=1}^{k_n}P_a(_{j=1}^{j(n)}E_k^{(n)}F_j^{(n)})\lambda _k^{(n)}\hfill \\ & =& _{k=1}^{k_n}_{j=1}^{j_n}P_a(X_{k,j}^{(n)})\lambda _k^{(n)}\hfill \end{array}$$
and by the same reason $`c_n=_{k=1}^{k_n}_{j=1}^{j_n}P_a(X_{k,j}^{(n)})\mu _j^{(n)}`$. Hence,
$$a_nc_n=\underset{k=1}{\overset{k_n}{}}\underset{j=1}{\overset{j_n}{}}P_a(X_{k,j}^{(n)})(\lambda _k^{(n)}\mu _j^{(n)}).$$
(8)
If $`X_{k,j}^{(n)}=\mathrm{}`$, then $`P_a(X_{k,j}^{(n)})=0`$. Therefore we can assume that $`X_{k,j}^{(n)}\mathrm{}`$ in expression (8). In this case let us consider $`zX_{k,j}^{(n)}`$. We deduce that $`|\lambda _k^{(n)}\mu _j^{(n)}||\lambda _k^{(n)}z|+|z\mu _j^{(n)}|1/n+1/n=2/n`$ for all $`1kk_n`$, $`1jj_n`$. Note that $`a_nc_n`$ is a normal element. Therefore,
$$a_nc_n=sup\{|\lambda _k^{(n)}\mu _j^{(n)}|:X_{k,j}^{(n)}\mathrm{},1kk_n,1jj_n\}2/n.$$
(9)
On the other hand, we have
$$\begin{array}{ccc}\stackrel{~}{a}_n& =& _{k=1}^{k_n}P_a(E_k^{(n)})^{(2l+1)}\lambda _k^{(n)}=_{k=1}^{k_n}(P_a(_{j=1}^{j(n)}X_{k,j}^{(n)}))^{(2l+1)}\lambda _k^{(n)}\hfill \\ & =& _{k=1}^{k_n}_{j_0,\mathrm{}j_{2l}=1}^{j_n}P_a(X_{k,j_0}^{(n)})\mathrm{}P_a(X_{k,j_{2l}}^{(n)})\lambda _k^{(n)}\hfill \end{array}$$
and by the same reason
$$\stackrel{~}{c}_n=\underset{j=1}{\overset{j_n}{}}\underset{k_0,\mathrm{}k_{2l}=1}{\overset{k_n}{}}P_a(X_{j,k_0}^{(n)})\mathrm{}P_a(X_{j,k_{2l}}^{(n)})\mu _j^{(n)}.$$
Thus we obtain
$$\begin{array}{ccc}\stackrel{~}{a}_n\stackrel{~}{c}_n& =& _{k=1}^{k_n}_{j=1}^{j_n}P_a(X_{k,j}^{(n)})^{(2l+1)}(\lambda _k^{(n)}\mu _j^{(n)})\hfill \\ & +& _{k=1}^{k_n}(\stackrel{~}{}_{j_0,\mathrm{},j_{2l}=1}^{j(n)}P_a(X_{k,j_0}^{(n)})\mathrm{}P_a(X_{k,j_{2l}}^{(n)}))\lambda _k^{(n)}\hfill \\ & & _{j=1}^{j_n}(\stackrel{~}{}_{k_0,\mathrm{},k_{2l}=1}^{k(n)}P_a(X_{j,k_0}^{(n)})\mathrm{}P_a(X_{j,k_{2l}}^{(n)}))\mu _j^{(n)}\hfill \\ & =& \gamma _n^{(1)}+\gamma _n^{(2)}\gamma _n^{(3)}.\hfill \end{array}$$
From Theorem 3 we conclude that $`\gamma _n^{(2)}_𝒞=\gamma _n^{(3)}_𝒞=0`$. The elements $`a_nc_n`$ and $`\gamma _n^{(1)}`$ are normal. Furtermore, the sets $`sp(a_nc_n)\{0\}`$ and $`sp(\gamma _n^{(1)})\{0\}`$ coinside. Therefore, $`a_nc_n=\gamma _n^{(1)}`$. Using inequality (9), we obtain
$$\stackrel{~}{a}_n_𝒞\stackrel{~}{c}_n_𝒞=\gamma _n^{(1)}_𝒞\gamma _n^{(1)}=a_nc_n2/n.$$
Therefore, $`Tr_{}(\stackrel{~}{a}_n_𝒞)Tr_{}(\stackrel{~}{c}_n_𝒞)2Tr_{}/n`$. This estimate implies that
$$T(a,\{a_n\})=\underset{n}{lim}Tr_{}(\stackrel{~}{a}_n_𝒞)=\underset{n}{lim}Tr_{}(\stackrel{~}{c}_n_𝒞)=T(a,\{c_n\}).$$
The proof is complete. $`\mathrm{}`$
###### Proposition 4
Suppose $`a,bM_{\mathrm{}}(A)_\nu `$ are equivalent in the sence of Definition 1. Then $`T(a)=T(b)`$.
Proof. For each $`n𝐍`$ let us cover the space $`sp(a)sp(b)`$ by disjoint Borel sets and enlarge this system of sets to a disjoint cover $`\{E_k^{(n)}\}_{k=0}^{k_n}`$ of the space $`sp(a)sp(b)`$. As above, we suppose that $`\mathrm{diam}(E_k^{(n)})1/n`$ for all $`0kk_n`$. Also, let us consider $`\lambda _k^{(n)}E_k^{(n)}`$. If $`0sp(a)sp(b)`$, then we shall assume that $`0E_0^{(n)}`$ and $`\lambda _0^{(n)}=0`$. In the opposite case, we put $`E_0^{(n)}=\mathrm{}`$. Furtemore, for any $`1kk_n`$ we can assume that $`E_k^{(n)}`$ is an admissible Borel set.
Projections $`P_a(E_k^{(n)})`$ and $`P_b(E_k^{(n)})`$ are stably equivalent for all $`k,n1`$. Therefore, $`Ch_{2l}^0([P_a(E_k^{(n)})])=Ch_{2l}^0([P_b(E_k^{(n)})])`$. Since the trace map is an isomorphism, we conclude that $`P_a(E_k^{(n)})^{(2l+1)}_{HC}=P_b(E_k^{(n)})^{(2l+1)}_{HC}`$. Thus,
$$\begin{array}{ccc}P_a(E_k^{(n)})^{(2l+1)}_𝒞& =& \pi _{}(P_a(E_k^{(n)})^{(2l+1)}_{HC})\hfill \\ & =& \pi _{}(P_b(E_k^{(n)})^{(2l+1)}_{HC})=P_b(E_k^{(n)})^{(2l+1)}_𝒞,\hfill \end{array}$$
where $`\pi _{}`$ is map (3). So we obtain that
$$\stackrel{~}{a}_n_𝒞=\underset{k=0}{\overset{k_n}{}}P_a(E_k^{(n)})^{(2l+1)}_𝒞\lambda _k^{(n)}=\underset{k=0}{\overset{k_n}{}}P_b(E_k^{(n)})^{(2l+1)}_𝒞\lambda _k^{(n)}=\stackrel{~}{b}_n_𝒞$$
for all $`n1`$ so $`T(a)=T(b)`$. $`\mathrm{}`$
###### Definition 3
We define the generalized Chern character as the map
$$𝒞h_{2l}^0:N_0(A)𝒞_{2l}(A),[a][b](1)^l(T(a)T(b)).$$
It follows from Proposition 4 that the generalized Chern character is well defined. An immediate verification gives us
###### Proposition 5
The generalized Chern character is a homomorphism of groups. $`\mathrm{}`$
###### Theorem 5
For any $`l0`$ there is a commutative diagramm
$$\begin{array}{ccc}K_0(A)& & N_0(A)\\ Ch_0^{2l}& & 𝒞h_0^{2l}\\ HC_{2l}(A)& \stackrel{\pi _{}}{}& 𝒞_{2l}(A),\end{array}$$
where $`\pi _{}`$ is map (3).
Proof. Under the notation of the beginning of this section let us argue as follows. Let $`pM_{\mathrm{}}(A)`$ be a projection. In this case the cover $`^{(n)}`$ of the spectrum of $`p`$ coinside with the set $`\{\{1\},\{0\}\}`$ for all $`n1`$. Therefore for all $`n1`$ we have $`\stackrel{~}{p}_n=p^{(2l+1)}`$. Finally, we obtain
$$\begin{array}{ccc}𝒞h_0^{2l}([p])& =& lim_nTr_{}((1)^l\stackrel{~}{p}_n_{𝒞_{2l}})\hfill \\ & =& lim_nTr_{}((1)^lp^{(2l+1)}_{𝒞_{2l}})\hfill \\ & =& Tr_{}((1)^lp^{(2l+1)}_{𝒞_{2l}})\hfill \\ & =& (1)^lTr(p^{(2l+1)})_{𝒞_{2l}}\hfill \\ & =& \pi _{}((1)^lTr(p^{(2l+1)})_{HC_{2l}})=\pi _{}Ch_{2l}^0([p]).\hfill \end{array}$$
The proof is complete. $`\mathrm{}`$
###### Theorem 6
For any $`l0`$ there is a commutative diagramm
$$\begin{array}{ccc}N_0(A)_{fin}& \stackrel{h}{}& K_0(A)𝐂\\ 𝒞h_0^{2l}& & \stackrel{~}{Ch}_0^{2l}\\ 𝒞_{2l}(A)& \stackrel{\pi _{}}{}& HC_{2l}(A),\end{array}$$
where $`h`$ is map (1) and $`\stackrel{~}{Ch}_{2l}^0([p]\lambda )=Ch_{2l}^0([p])\lambda `$.
Proof. Let $`a`$ be an element of $`M_{\mathrm{}}(A)_{fin}`$ such that $`sp(a)=\{\lambda _1,\mathrm{},\lambda _n\}`$. Assume $`p_i:=P_a(\{\lambda _i\})`$. Then we have
$$\begin{array}{ccc}\pi _{}\stackrel{~}{Ch}_0^{2l}h([a])& =& \pi _{}\stackrel{~}{Ch}_0^{2l}(_{i=1}^n[p_i]\lambda _i)\hfill \\ & =& \pi _{}(_{i=1}^nCh_0^{2l}([p_i])\lambda _i)\hfill \\ & =& _{i=1}^nTr_{}((1)^lp_i^{(2l+1)}_𝒞)\lambda _i.\hfill \end{array}$$
On the other hand, we can suppose that $`^{(k)}=\{\{\lambda _1\},\mathrm{},\{\lambda _n\}\}`$ and $`\stackrel{~}{a}_k=_{i=1}^np_i^{(2l+1)}\lambda _i`$ for all $`k1`$. Whence,
$$\begin{array}{ccc}𝒞h_0^{2l}([a])=(1)^lT(a)& =& (1)^llim_kTr_{}\stackrel{~}{a}_k_𝒞\hfill \\ & =& (1)^l_{i=1}^nTr_{}(p_i^{(2l+1)}_𝒞)\lambda _i.\mathrm{}\hfill \end{array}$$
In particular, Theorem 6 implies that one can extend the generalized Chern character to the map from the quotient group $`N_0(A)/Kerh`$ to the even Banach homology.
## 6 Generalized Lefshetz numbers
Suppose $`A`$, as above, is a von Neumann algebra, $`G`$ is a compact Lie group, and $`X`$ is a compact $`G`$-manifold. Let us denote by $`𝒫(A)`$ the category of finitely generated projective modules over $`A`$.
Let us recall some notation from . The set of all $`G`$-$`A`$-bundles over $`X`$ is an abelian semigroup with respect to the direct sum operation. The symmetrization of this semigroup is denoted by $`K_G(X;A)`$. Assume $`K^G(A):=K_G(pt;A)`$. In this situation there is an isomorphism
$$K^G(A)K_0(A)R(G),$$
(10)
where $`R(G)`$ is the ring of representations for $`G`$.
Let us consider a sequence $`\{E^i\}`$ of $`G`$-$`A`$-bundles over $`X`$ together with equivariant pseudo-differential operators $`\{d_i:\mathrm{\Gamma }(E^i)\mathrm{\Gamma }(E^{i1})\}`$, where by $`\mathrm{\Gamma }(E^i)`$ we denote the Banach (with respect to the uniform topology) $`A`$-module of continuous sections of $`E^i`$. Besides, let us denote by $`\sigma _i`$ the symbol of $`d_i`$. Then this sequence of bundles and operators is called a $`G`$-$`A`$-elliptic complex (and is denoted by $`(E,d)`$) if it satisfies the following conditions:
(i) $`d_id_{i+1}=0,`$
(ii) the sequence of symbols
$$0\pi ^{}E^n\stackrel{\sigma _n}{}\pi ^{}E^{n1}\mathrm{}\stackrel{\sigma _1}{}\pi ^{}E^00$$
is exact out of some compact neighbourhood of the zero section $`XT^{}X`$. Here $`\pi :T^{}XX`$ is the natural projection.
The index of the elliptic operator $`F=d+d^{}:\mathrm{\Gamma }(E_{ev})\mathrm{\Gamma }(E_{od})`$ is an element of the group $`K^G(A)`$. Furthermore, for any $`gG`$ by computation of the character we can define the map $`g:R(G)𝐂`$. Whence, using isomorphism (10), we obtain the map
$$g:K^G(A)K_0(A)𝐂.$$
Then the Lefschetz number of the first type is defined as follows
$$L_1(E,g)=g(\mathrm{index}(F))K_0(A)𝐂.$$
Note that there exists a connection between these Lefschetz numbers and fixed points of $`g`$ (see ).
Now let us consider an $`A`$-elliptic complex $`(E,d)`$ and its unitary endomorphism $`U`$. Furthermore, we shall assume that $`U=U_g`$ for some representation $`U_g`$ of a compact Lie group $`G`$.
Let $``$ be a Hilbert $`A`$-module (see, for example, ). We denote by $`\mathrm{End}_A()`$ the Banach algebra of all bounded $`A`$-homomorphisms of $``$. Now let us consider some strongly continuous representation $`G\mathrm{End}_A()`$. Then this representation is called unitary, if $`gx,gy=x,y`$ for any $`gG,x,y`$. A Hilbert $`A`$-module together with a unitary representation of the group $`G`$ is called a Hilbert $`G`$-$`A`$-module. Besides, a set $`\{x\}_{\beta B}`$ is a system of generators for $``$, if finite sums $`\{_kx_ka_k:a_kA\}`$ are dense in $``$.
We need the following result of .
###### Theorem 7
Let $``$ be a countably generated Hilbert $`G`$-$`A`$-module. Besides, let $`\{V_\pi \}`$ be a full system of pairwise not isomorphic unitary finite-dimentional irreducible representations for $`G`$. Then there exists a $`G`$-$`A`$-isomorphism
$$\underset{\pi }{}\mathrm{Hom}_G(V_\pi ,)_𝐂V_\pi .$$
Here the algebra $`A`$ (the group $`G`$) acts on the first (on the second) multiplier of the space $`\mathrm{Hom}_G(V_\pi ,)_𝐂V_\pi `$. $`\mathrm{}`$
Let the $`G`$-$`A`$-module $``$ belong to the class $`𝒫(A)`$. Then it is clear that $`_\pi :=\mathrm{Hom}_G(V_\pi ,)𝒫(A)`$. Furthermore, only finite number of terms in the sum $`\underset{\pi }{}_\pi _𝐂V_\pi `$ is not equal to zero (see \[13, 1.3.49\]).
In particular, we obtain for $`=A^n`$ the following formula
$$A^n\underset{k=1}{\overset{M}{}}Q_kV_k,$$
where $`V_k𝐂^{L_k}`$ and $`Q_k𝒫(A)`$. Therefore,
$$U_g(\underset{k=1}{\overset{M}{}}x_kv_k)=\underset{k=1}{\overset{M}{}}x_ku_g^kv_k=\underset{k=1}{\overset{M}{}}\underset{s=1}{\overset{L_k}{}}x_ke^{i\phi _s^k}v_k^sf_s.$$
(11)
Here $`f_1,\mathrm{},f_{L_k}`$ is a basis for $`V_k`$ such that the operator $`u_g^k`$ is diagonal with respect to it; $`v_k=v_k^sf_s`$. In this case let us define
$$\tau (U_g)=\underset{k=1}{\overset{M}{}}Ch_{2l}^0[Q_k]\mathrm{Trace}(u_g^k)HC_{2l}(A).$$
The following result was proved in .
###### Lemma 5
For the $`A`$-Fredholm operator $`F=d+d^{}:\mathrm{\Gamma }(E_{ev})\mathrm{\Gamma }(E_{od})`$ there exists a decomposition
$$F:M_0\stackrel{~}{N_0}M_1\stackrel{~}{N_1},F:M_0M_1$$
(12)
such that
$$\stackrel{~}{N_0}=\underset{j=0}{\overset{T}{}}N_{2j},\stackrel{~}{N_1}=\underset{j=1}{\overset{T}{}}N_{2j1},N_m\mathrm{\Gamma }(E_m),$$
where $`N_m`$ are projective $`U`$-invariant Hilbert $`A`$-modules. $`\mathrm{}`$
Now the Lefschetz number of the second type is defined as follows
$$L_{2l}(E,U_g)=\underset{j}{}(1)^j\tau (U_g|N_j)HC_{2l}(A).$$
This definition is well.
For more detail about W\*-Lefschetz numbers we refer to the works .
Now let us consider an $`A`$-elliptic complex $`(E,d)`$ and an arbitrary unitary endomorphism $`U`$ of it ($`U`$ is not necessarily an element of some representation of $`G`$). In this situation let us formulate the following
###### Definition 4
We define the generalized Lefschetz number $`_1`$ as follows
$$_1(E,U)=\underset{j}{}(1)^j[U|N_j]N_0(A).$$
Note that generalized Lefschetz numbers are well defined. This follows by the same reason that for the Lefschetz numbers of the second type (see \[13, 5.2.21\]).
###### Theorem 8
Let $`U`$ be a unitary endomorphism of an $`A`$-elliptic complex $`(E,d)`$. Besides, suppose that $`U=U_g`$ for some representation $`U_g`$ of a compact Lie group $`G`$. Then $`_1(E,U_g)`$ belongs to the group $`N_0(A)_{fin}`$ and $`h(_1(E,U_g))=L_1(E,U_g)`$, where $`h`$ is map (1).
Proof. Let us examine decomposition (12) for $`F`$. We have shown above that there exist isomorphisms
$$N_{2j}\underset{k_j=1}{\overset{K_j}{}}P_{k_j}^{(j)}V_{k_j}^{(j)},N_{2j1}\underset{l_j=1}{\overset{L_j}{}}Q_{l_j}^{(j)}W_{l_j}^{(j)},$$
where $`V_{k_j}^{(j)}`$ and $`W_{l_j}^{(j)}`$ are comlex vector spaces of irreducible unitary representations of $`G`$, $`P_{k_j}^{(j)}`$ and $`Q_{l_j}^{(j)}`$ are $`G`$-trivial modules from $`𝒫(A)`$. Thus we get
$$\mathrm{index}(F)=\underset{j=0}{\overset{T}{}}\underset{k_j=1}{\overset{K_j}{}}[P_{k_j}^{(j)}]\chi (V_{k_j}^{(j)})\underset{j=1}{\overset{T}{}}\underset{l_j=1}{\overset{L_j}{}}[Q_{l_j}^{(j)}]\chi (W_{l_j}^{(j)})$$
and
$$L_1(E,g)=\underset{j=0}{\overset{T}{}}\underset{k_j=1}{\overset{K_j}{}}[P_{k_j}^{(j)}]\mathrm{Trace}(g|V_{k_j}^{(j)})\underset{j=1}{\overset{T}{}}\underset{l_j=1}{\overset{L_j}{}}[Q_{l_j}^{(j)}]\mathrm{Trace}(g|W_{l_j}^{(j)}).$$
Here we have denoted by $`\chi `$ the character of the representation of $`G`$.
On the other hand, using expression (11), we conclude that
$$sp(U_g|(P_{k_j}^{(j)}V_{k_j}^{(j)}))=sp(u_g^{(j),k_j}|V_{k_j}^{(j)}).$$
This implies that $`_1(E,U_g)`$ belongs to the group $`N_0(A)_{fin}`$. Furtermore, for any $`e^{i\phi }`$ from $`sp(u_g^{(j),k_j}|V_{k_j}^{(j)})`$ the spectral projection of $`U_g`$ corresponding to this point is equal to $`P_{k_j}^{(j)}`$. Thus we obtain
$$\begin{array}{ccc}h(_1(E,U_g))& =& h(_{j=0}^{2T}(1)^j[U_g|N_j])\hfill \\ & =& h(_{j=0}^T_{k_j=1}^{K_j}[Id_{P_{k_j}^{(j)}}u_g^{(j),k_j}|V_{k_j}^{(j)}]\hfill \\ & & _{j=1}^T_{l_j=1}^{L_j}[Id_{Q_{l_j}^{(j)}}u_g^{(j),l_j}|W_{l_j}^{(j)}])\hfill \\ & =& h(_{j=0}^T_{k_j=1}^{K_j}[P_{k_j}^{(j)}\mathrm{Trace}(u_g^{(j),k_j}|V_{k_j}^{(j)})]\hfill \\ & & _{j=1}^T_{l_j=1}^{L_j}[Q_{l_j}^{(j)}\mathrm{Trace}(u_g^{(j),l_j}|W_{l_j}^{(j)})])\hfill \\ & =& _{j=0}^T_{k_j=1}^{K_j}[P_{k_j}^{(j)}]\mathrm{Trace}(u_g^{(j),k_j}|V_{k_j}^{(j)})\hfill \\ & & _{j=1}^T_{l_j=1}^{L_j}[Q_{l_j}^{(j)}]\mathrm{Trace}(u_g^{(j),l_j}|W_{l_j}^{(j)}).\hfill \end{array}$$
Thus we establish the required statement. $`\mathrm{}`$
###### Theorem 9
Under the assumptions of the previous theorem, we have
$$\pi _{}(L_{2l}(E,U_g))=𝒞h_{2l}^0(_1(E,U_g)).$$
Here $`\pi _{}`$ is map (3).
Proof. From Theorems 68 we deduce that
$$\begin{array}{ccc}𝒞h_{2l}^0(_1(E,U_g))& =& \pi _{}\stackrel{~}{Ch}_{2l}^0h(_1(E,U_g))\hfill \\ & =& \pi _{}\stackrel{~}{Ch}_{2l}^0(L_1(E,g)).\hfill \end{array}$$
Furtermore, it follows from \[13, Theorem 5.2.22\] that $`\stackrel{~}{Ch}_{2l}^0(L_1(E,g))=L_{2l}(E,U_g)`$. The proof is complete. $`\mathrm{}`$
Acknowledgements. The author is grateful to Prof. E.V. Troitsky for the attention to the work and also to Dr. V.M. Manuilov and Prof. A.S. Mishchenko for helpful discussions.
Alexandre Pavlov
Chair of Higher Geometry and Topology
Dept. of Mech. and Mathematics
Moscow State University
Vorobjevi gori
Moscow 119899
Russia
email: pavlov@mech.math.msu.su
|
warning/0003/math-ph0003018.html
|
ar5iv
|
text
|
# References
IMSc/2000/03/09
AN INTRODUCTION TO QUANTUM ALGEBRAS
AND THEIR APPLICATIONS Article based on Narayani, N. and Kadayam S. Sankaran Memorial Lectures delivered by the author at the Department of Mathematics, Ramakrishna Mission Vivekananda College, Mylapore, Chennai, on 23 & 27 January 1999
Ramaswamy JAGANATHAN
The Institute of Mathematical Sciences
C.I.T. Campus, Tharamani, Chennai - 600 113, India
E-Mail :jagan@imsc.ernet.in
Abstract : A very elementary introduction to quantum algebras is presented and a few examples of their physical applications are mentioned.
I shall give a very elementary introduction to the topic of quantum algebras and mention a few physical applications. Quantum algebras, or quantum groups, extend the domain of classical group theory and constitute a new and growing field of mathematics with vast potential for applications in physics. In fact, the origins of quantum groups lie in physics : in the studies on the behaviour of integrable systems in quantum field theory and statistical mechanics, using the quantum inverse scattering method, by Sklyanin, Kulish, Reshetikhin, Takhtajan, and Faddeev in the 1980s. Mathematical abstraction from the observations on the common features of these systems led to the definition of the concept of a quantum group and the studies of Drinfeld, Jimbo, Manin, Woronowicz, Connes, Wess, Zumino, Macfarlane, Biedenharn, …, revealed many aspects of quantum groups from different mathematical and physical points of view. I shall neither go into the history of the developements of the various significant concepts nor give any specific references to their origins. I shall cite at the end some books and articles which would give these details and lead to the vast literature on the subject of quantum groups and their applications.
Let us consider a two dimensional classical vector space whose elements can be written as
$$\left(\begin{array}{c}x\\ y\end{array}\right),$$
(1)
where the coordinates $`x`$ and $`y`$ are real variables and commute with each other, i.e.,
$$xy=yx.$$
(2)
In the space of functions $`\left\{f(x,y)\right\}`$ defined in the above vector space the partial derivative operations $`\frac{}{x}`$ and $`\frac{}{y}`$ and the operations of multiplications by $`x`$ and $`y`$ satisfy the differential calculus
$`[x,y]=0,[{\displaystyle \frac{}{x}},{\displaystyle \frac{}{y}}]`$ $`=`$ $`0,[{\displaystyle \frac{}{x}},y]=0,[{\displaystyle \frac{}{y}},x]=0,`$
$`[{\displaystyle \frac{}{x}},x]`$ $`=`$ $`1,[{\displaystyle \frac{}{y}},y]=1,`$ (3)
where the commutator bracket $`[,]`$ is defined by
$$[A,B]=ABBA.$$
(4)
Let us make a linear transformation of the vector $`\left(\begin{array}{c}x\\ y\end{array}\right)`$ as
$$\left(\begin{array}{c}x^{}\\ y^{}\end{array}\right)=M\left(\begin{array}{c}x\\ y\end{array}\right)=\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)\left(\begin{array}{c}x\\ y\end{array}\right)$$
(5)
where the entries of the matrix $`M`$ are real and satisfy the condition
$$adbc=\text{det}M=1,$$
(6)
or in other words the transformation (5) is an element of the group $`SL(2,R)`$. Then, the new coordinates $`x^{}`$ and $`y^{}`$ and partial derivatives with respect to them, namely $`\frac{}{x^{}}`$ and $`\frac{}{y^{}}`$, also satisfy the same relations as (3). This is easily checked by noting
$$\left(\begin{array}{c}\frac{}{x^{}}\\ \\ \frac{}{y^{}}\end{array}\right)=\left(\begin{array}{ccc}\hfill d& & \hfill c\\ & & \\ \hfill b& & \hfill a\end{array}\right)\left(\begin{array}{c}\frac{}{x}\\ \\ \frac{}{y}\end{array}\right)=\stackrel{~}{M^1}\left(\begin{array}{c}\frac{}{x}\\ \\ \frac{}{y}\end{array}\right),$$
(7)
where $`\stackrel{~}{A}`$ means the transpose of the matrix $`A`$. We say that the differential calculus on the two dimensional $`(x,y)`$-plane is covariant under the group $`SL(2,R)`$.
In nature all physical systems are quantum mechanical. But, the quantum mechanical behaviour is generally revealed only at the microscopic (molecular and deeper) level. At the macroscopic level of everyday experience quantum physics becomes classical physics as an approximation. So, only classical physics was discovered first and observations of failure of classical physics at the atomic level led to the discovery of quantum physics in the 20th century. This passage from classical physics to quantum physics can be mathematically described as a process of deformation of the classical physics wherein the commuting classical observables of a physical system are replaced by noncommuting hermitian operators. This process is characterized by a very small deformation parameter known as the Planck constant $`\mathrm{}`$ and, roughly speaking, in the limit $`\mathrm{}0`$ classical physics is recovered from the structure of quantum physics.
In analogy with the process of quantizing the classical physics let us now quantize the classical vector space to get a quantum vector space by assuming that the coordinates do not commute with each other at any point. Let us model the noncommutativity of the coordinates $`X`$ and $`Y`$ of a two dimensional quantum plane by
$$XY=qYX,$$
(8)
where $`q`$ is the deformation parameter which we shall consider, in general, to be any nonzero complex number. Note that in the limit $`q1`$ the noncommuting quantum coordinates $`X`$ and $`Y`$ become commuting classical coordinates. To be specific, let us choose
$$q=e^{i\theta }.$$
(9)
It is easy to find an example of such noncommuting variables. Let $`T_\alpha `$ and $`G_{\theta /\alpha }`$ be operators acting on functions of a single real variable $`x`$ such that
$$T_\alpha \psi (x)=\psi (x+\alpha ),G_{\theta /\alpha }\psi (x)=e^{i\theta x/\alpha }\psi (x).$$
(10)
Then, for any $`\psi (x)`$,
$$T_\alpha G_{\theta /\alpha }\psi (x)=e^{i\theta (x+\alpha )/\alpha }\psi (x+\alpha )=e^{i\theta }G_{\theta /\alpha }T_\alpha \psi (x),$$
(11)
for a given fixed value of $`\theta `$. Thus, with the variation of $`\alpha `$, $`T_\alpha `$ and $`G_{\theta /\alpha }`$ become noncommuting variables obeying the relation
$$T_\alpha G_{\theta /\alpha }=e^{i\theta }G_{\theta /\alpha }T_\alpha ,$$
(12)
with fixed value for $`\theta `$.
Now the interesting questions are :
* Is it possible to define a differential calculus on the two dimensional $`(X,Y)`$-plane?
* If so, will it be covariant under some generalization of the classical group $`SL(2)`$?
The answers are yes, yes! First let us understand how one can give a meaning to partial derivatives with respect to $`X`$ and $`Y`$. These have to operate on the space of functions $`\left\{f(X,Y)\right\}`$ which we shall consider to be polynomials in $`X`$ and $`Y`$. We can write $`f(X,Y)=_{m,n}f_{mn}X^mY^n`$ since any polynomial in $`X`$ and $`Y`$, with coefficients commuting with $`X`$ and $`Y`$, can be brought to this form using the commutation relation (8). Then, if we take formally
$$\frac{}{X}X^m=mX^{m1},\frac{}{Y}Y^n=nY^{n1},$$
(13)
we would have a differential calculus in the $`(X,Y)`$-plane, as desired, once we prescribe consistently the remaining commutation relations between $`X`$, $`Y`$, $`\frac{}{X}`$ and $`\frac{}{Y}`$.
Without going into any further details, let me state the complete set of commutation relations defining the differential calculus in the $`(X,Y)`$-plane (with $`q=e^{i\theta }`$) :
$`XY=qYX,{\displaystyle \frac{}{X}}{\displaystyle \frac{}{Y}}=q^1{\displaystyle \frac{}{Y}}{\displaystyle \frac{}{X}},{\displaystyle \frac{}{X}}Y=qY{\displaystyle \frac{}{X}},{\displaystyle \frac{}{Y}}X=qX{\displaystyle \frac{}{Y}},`$
$`{\displaystyle \frac{}{X}}Xq^2X{\displaystyle \frac{}{X}}=1+(q^21)Y{\displaystyle \frac{}{Y}},{\displaystyle \frac{}{Y}}Yq^2Y{\displaystyle \frac{}{Y}}=1.`$ (14)
This noncommutative differential calculus on the two-dimensional quantum plane is seen to be covariant under the transformations
$`\left(\begin{array}{c}X^{}\\ Y^{}\end{array}\right)`$ $`=`$ $`T\left(\begin{array}{c}X\\ Y\end{array}\right)=\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)\left(\begin{array}{c}X\\ Y\end{array}\right)`$ (23)
$`\left(\begin{array}{c}\frac{}{X^{}}\\ \\ \frac{}{Y^{}}\end{array}\right)`$ $`=`$ $`\stackrel{~}{T^1}\left(\begin{array}{c}\frac{}{X}\\ \\ \frac{}{Y}\end{array}\right)=\left(\begin{array}{ccc}D& & qC\\ & & \\ q^1B& & A\end{array}\right)\left(\begin{array}{c}\frac{}{X}\\ \\ \frac{}{Y}\end{array}\right),`$ (36)
provided
$$A,B,C,\text{and}D\text{commute with}X\text{and}Y,$$
(37)
$`AB=qBA,CD=qDC,AC=qCA,BD=qDB,`$
$`BC=CB,ADDA=\left(qq^1\right)BC,`$ (38)
and
$$ADqBC=\text{det}_qT=1.$$
(39)
In other words $`X^{},Y^{},\frac{}{X^{}},\text{and}\frac{}{Y^{}}`$ defined by (36) satisfy the relations obtained from (14) by just replacing $`X`$ and $`Y`$ by $`X^{}`$ and $`Y^{}`$ respectively. Note that $`\text{det}_qT`$ defined in (38) commutes with all the matrix elements of $`T`$. Verify that
$$T^1=\left(\begin{array}{cc}D& q^1B\\ qC& A\end{array}\right)$$
(40)
is such that
$$TT^1=T^1T=1\mathrm{l}=\left(\begin{array}{cc}\hfill 1& \hfill 0\\ \hfill 0& \hfill 1\end{array}\right).$$
(41)
A matrix $`T=\left(\begin{array}{cc}\hfill A& \hfill B\\ \hfill C& \hfill D\end{array}\right)`$ is called a $`2\times 2`$ quantum matrix if its matrix elements $`\{A,B,C,D\}`$ satisfy the commutation relations in (38). In the limit $`q1`$ a quantum matrix $`T`$ becomes a classical matrix with commuting elements. Note that the identity matrix $`1\mathrm{l}=\left(\begin{array}{cc}\hfill 1& \hfill 0\\ \hfill 0& \hfill 1\end{array}\right)`$ is a quantum matrix.
Entries of a quantum matrix $`T`$ are noncommuting variables satisfying the commutation relations (38). Let $`T_1=\left(\begin{array}{cc}\hfill A_1& \hfill B_1\\ \hfill C_1& \hfill D_1\end{array}\right)`$ and $`T_2=\left(\begin{array}{cc}\hfill A_2& \hfill B_2\\ \hfill C_2& \hfill D_2\end{array}\right)`$ be any two quantum matrices; i.e., $`\{A_1,B_1,C_1,D_1\}`$ obey the relations (38), and $`\{A_2,B_2,C_2,D_2\}`$ also obey the relations (38). The matrix elements of $`T_1`$ and $`T_2`$ may be ordinary classical matrices satisfying the required relations (38). Define the product
$`\mathrm{\Delta }_{12}(T)`$ $`=`$ $`T_1\dot{}T_2=\left(\begin{array}{cc}\hfill A_1& \hfill B_1\\ \hfill C_1& \hfill D_1\end{array}\right)\dot{}\left(\begin{array}{cc}\hfill A_2& \hfill B_2\\ \hfill C_2& \hfill D_2\end{array}\right)`$ (46)
$`=`$ $`\left(\begin{array}{cc}\hfill A_1A_2+B_1C_2& \hfill A_1B_2+B_1D_2\\ \hfill C_1A_2+D_1C_2& \hfill C_1B_2+D_1D_2\end{array}\right)`$ (49)
$`=`$ $`\left(\begin{array}{cc}\hfill \mathrm{\Delta }_{12}(A)& \hfill \mathrm{\Delta }_{12}(B)\\ \hfill \mathrm{\Delta }_{12}(C)& \hfill \mathrm{\Delta }_{12}(D)\end{array}\right)`$ (52)
where $``$ denotes the direct product with the property $`(PR)(QS)=PQRS`$. Then one finds that the matrix elements of $`\mathrm{\Delta }_{12}(T)`$, namely,
$`\mathrm{\Delta }_{12}(A)`$ $`=`$ $`A_1A_2+B_1C_2,\mathrm{\Delta }_{12}(B)=A_1B_2+B_1D_2,`$
$`\mathrm{\Delta }_{12}(C)`$ $`=`$ $`C_1A_2+D_1C_2,\mathrm{\Delta }_{12}(D)=C_1B_2+D_1D_2,`$ (53)
also satisfy the commutation relations (38). In other words, $`\mathrm{\Delta }_{12}(T)`$ is also a quantum matrix. This product, $`\mathrm{\Delta }_{12}(T)=T_1\dot{}T_2`$, is called the coproduct or comultiplication. Note that there is no inverse for this coproduct. Under this coproduct the $`2\times 2`$ quantum matrices $`\left\{T=\left(\begin{array}{cc}\hfill A& \hfill B\\ \hfill C& \hfill D\end{array}\right)\right\}`$ form a pseudomatrix group, commonly called a quantum group, denoted by $`SL_q(2)`$. The algebra of functions over $`SL_q(2)`$, or the algebra of polynomials in $`\{A,B,C,D\}`$, is denoted by $`Fun_q(SL(2))`$. The coproduct operation ($`\mathrm{\Delta }`$) defined by (53), is symbolically written as
$`\mathrm{\Delta }(A)`$ $`=`$ $`AA+BC,\mathrm{\Delta }(B)=AB+BD,`$
$`\mathrm{\Delta }(C)`$ $`=`$ $`CA+DC,\mathrm{\Delta }(D)=CB+DD.`$ (54)
For any $`f(A,B,C,D)Fun_q(SL(2))`$ the definition of $`\mathrm{\Delta }`$ is extended as $`\mathrm{\Delta }f(A,B,C,D)=f(\mathrm{\Delta }(A),\mathrm{\Delta }(B),\mathrm{\Delta }(C),\mathrm{\Delta }(D))`$. The algebra $`Fun_q(SL(2))`$ is technically a Hopf algebra. To complete this algebraic structure two more operations called the coinverse (or antipode) denoted by $`S`$, and the counit denoted by $`\epsilon `$, are defined and these operations $`\mathrm{\Delta }`$, $`S`$, and $`\epsilon `$ are required to satisfy certain axioms. We shall not consider these details of the Hopf algebraic structure. From the point of view physical applications quantum groups provide a generalization of symmetry concepts and involve mainly two fundamental new ideas : deformation and noncommutative comultiplication.
The matrix $`T=\left(\begin{array}{cc}\hfill A& \hfill B\\ \hfill C& \hfill D\end{array}\right)`$ corresponds to the fundamental irreducible representation of $`SL_q(2)`$. Higher dimensional representations are defined as follows. An $`n\times n`$ matrix
$$T=\left(T_{ij}\right)=\left(\begin{array}{cccc}T_{11}& T_{12}& \mathrm{}& T_{1n}\\ T_{21}& T_{22}& \mathrm{}& T_{2n}\\ .& .& \mathrm{}& .\\ .& .& \mathrm{}& .\\ T_{n1}& T_{n2}& \mathrm{}& T_{nn}\end{array}\right),$$
(55)
is said to be an $`n`$-dimensional representation of $`SL_q(2)`$ if its matrix elements $`\left(T_{ij}\right)`$ are polynomials in $`\{A,B,C,D\}`$, or in other words elements of $`Fun_q(SL(2))`$, and satisfy the property
$`\left(T\dot{}T\right)_{ij}`$ $`=`$ $`{\displaystyle \underset{l=1}{\overset{n}{}}}T_{il}T_{lj}=\mathrm{\Delta }\left(T_{ij}\right)`$ (56)
$`=`$ $`T_{ij}(\mathrm{\Delta }(A),\mathrm{\Delta }(B),\mathrm{\Delta }(C),\mathrm{\Delta }(D)),i,j=1,2,\mathrm{},n.`$
Now, for example, look at
$$T^{(1)}=\left(\begin{array}{ccc}A^2& \sqrt{1+q^2}AB& B^2\\ \sqrt{1+q^2}AC& AD+q^1BC& \sqrt{1+q^2}BD\\ C^2& \sqrt{1+q^2}CD& D^2\end{array}\right).$$
(57)
It can be verified that $`T^{(1)}`$ is the $`3`$-dimensional representation of $`SL_q(2)`$. For instance, see that
$`\left(T^{(1)}\dot{}T^{(1)}\right)_{11}`$ $`=`$ $`{\displaystyle \underset{l=1}{\overset{3}{}}}T_{1l}^{(1)}T_{l1}^{(1)}`$ (58)
$`=`$ $`A^2A^2+\left(1+q^2\right)ABAC+B^2C^2`$
$`=`$ $`A^2A^2+ABAC+BACA+B^2C^2`$
$`=`$ $`(AA+BC)^2=(\mathrm{\Delta }(A))^2=\mathrm{\Delta }\left(T_{11}^{(1)}\right),`$
as required. Similarly, for other matrix elements of $`T^{(1)}`$ one can verify the property (56), namely,
$$\left(T\dot{}T\right)_{ij}=\mathrm{\Delta }\left(T_{ij}\right).$$
(59)
In the theory of classical Lie groups we know that an element of a Lie group $`G`$, say $`g`$, can be written as
$$g=e^{ϵ_1L_1}e^{ϵ_2L_2}\mathrm{}e^{ϵ_nL_n},$$
(60)
where the parameters $`\left\{ϵ_i\right\}`$ characterize the group element $`g`$ and $`\left\{L_i\right\}`$ are constant generators of the group $`G`$ satisfying a Lie algebra
$$[L_i,L_j]=\underset{k=1}{\overset{n}{}}C_{ij}^kL_k,i,j=1,2,\mathrm{},n,$$
(61)
with $`\left\{C_{ij}^k\right\}`$ as the structure constants. When the group element $`g`$ is close to the identity element ($`I`$) of the group, the parameters $`\left\{ϵ_i\right\}`$ are infinitesimals and one can write
$$gI+\underset{i=1}{\overset{n}{}}ϵ_iL_i.$$
(62)
Now, the interesting question is
* Is there an analogue of the Lie algebra in the case of a quantum group?
The answer is yes! To this end, first we have to recall some basic notions of the theory of $`q`$-series.
One defines the $`q`$-shifted factorial by
$$(x;q)_n=\{\begin{array}{cc}1,\hfill & n=0,\hfill \\ (1x)(1xq)(1xq^2)\mathrm{}\left(1xq^{n1}\right),\hfill & n=1,2,\mathrm{}.\hfill \end{array}$$
(63)
Then, with the notation
$$(x_1,x_2,\mathrm{},x_m;q)_n=(x_1;q)_n(x_2;q)_n\mathrm{}(x_m;q)_n,$$
(64)
an $`{}_{r}{}^{}\varphi _{s}^{}`$ basic hypergeometric series, or a general $`q`$-hypergeometric series, is given by
$`{}_{r}{}^{}\varphi _{s}^{}(a_1,a_2,\mathrm{},a_r;b_1,b_2,\mathrm{},b_s;q,z)`$
$`={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(a_1,a_2,\mathrm{},a_r;q)_n}{(b_1,b_2,\mathrm{},b_s;q)_n(q;q)_n}}\left((1)^nq^{n(n1)/2}\right)^{1+sr}z^n,`$
$`r,s=0,1,2,\mathrm{}.`$ (65)
Consider
$${}_{1}{}^{}\varphi _{0}^{}(0;;q,(1q)z)=e_q^z=\underset{n=0}{\overset{\mathrm{}}{}}\frac{z^n}{[n]_q!},$$
(66)
where
$`[n]_q={\displaystyle \frac{1q^n}{1q}},`$
$`[n]_q!=[n]_q[n1]_q[n2]_q\mathrm{}[2]_q[1]_q,n=1,2,\mathrm{},[0]_q!=1.`$ (67)
The $`q`$-number $`[n]_q`$, or the so-called basic number, was defined by Heine (1846). Much of Ramanujan’s work is related to these $`q`$-series. Note that
$$[n]_q\stackrel{q1}{}n,e_q^z\stackrel{q1}{}e^z.$$
(68)
Thus, $`e_q^z`$ defined by (66) is a $`q`$-generalization of the exponential function, called a $`q`$-exponential function ; note that there can be several generalizations of the exponential function satisfying the condition that in the limit $`q1`$ it should become the standard exponential function. In the theory of quantum groups a new definition of the $`q`$-number is often useful. It is
$$\left[\left[n\right]\right]_q=\frac{q^nq^n}{qq^1}.$$
(69)
Note that $`\left[\left[n\right]\right]_q`$ also becomes $`n`$ in the limit $`q1`$, and $`\left[\left[n\right]\right]_q`$ is symmetric with respect to the interchange of $`q`$ and $`q^1`$ unlike Heine’s $`[n]_q`$.
Now consider the $`2`$-dimensional $`T`$-matrix parametrized as
$$T=\left(\begin{array}{cc}\hfill A& \hfill B\\ \hfill C& \hfill D\end{array}\right)=\left(\begin{array}{cc}e^\alpha & e^\alpha \beta \\ \gamma e^\alpha & e^\alpha +\gamma e^\alpha \beta \end{array}\right),$$
(70)
which requires the variable parameters $`\{\alpha ,\beta ,\gamma \}`$ to satisfy a Lie algebra
$$[\alpha ,\beta ]=(\mathrm{log}q)\beta ,[\alpha ,\gamma ]=(\mathrm{log}q)\gamma ,[\beta ,\gamma ]=0,$$
(71)
so that $`\{A,B,C,D\}`$ obey the algebra (38). Then, one can write
$$T=e_{q^2}^{\gamma 𝒳_{}^{(1/2)}}e^{2\alpha 𝒳_0^{(1/2)}}e_{q^2}^{\beta 𝒳_+^{(1/2)}},$$
(72)
with
$$𝒳_0^{(1/2)}=\frac{1}{2}\left(\begin{array}{cc}\hfill 1& \hfill 0\\ \hfill 0& \hfill 1\end{array}\right),𝒳_{}^{(1/2)}=\left(\begin{array}{cc}\hfill 0& \hfill 0\\ \hfill 1& \hfill 0\end{array}\right),𝒳_+^{(1/2)}=\left(\begin{array}{cc}\hfill 0& \hfill 1\\ \hfill 0& \hfill 0\end{array}\right).$$
(73)
Of course, in this case, $`e_{q^2}^{\gamma 𝒳_{}^{(1/2)}}`$ and $`e_{q^2}^{\beta 𝒳_+^{(1/2)}}`$ are trivially the same as $`e^{\gamma 𝒳_{}^{(1/2)}}`$ and $`e^{\beta 𝒳_+^{(1/2)}}`$, respectively, since $`\left(𝒳_\pm ^{(1/2)}\right)^2=0`$.
Actually, equation (72) is the special case of a universal formula for the representations of the $`T`$-matrices of $`SL_q(2)`$ and corresponds to the fundamental representation. The generic form of the $`T`$-matrix is given by
$$T=e_{q^2}^{\gamma 𝒳_{}}e^{2\alpha 𝒳_0}e_{q^2}^{\beta 𝒳_+},$$
(74)
called the universal $`𝒯`$-matrix, where $`\{𝒳_0,𝒳_+,𝒳_{}\}`$ obey the algebra
$`[𝒳_0,𝒳_\pm ]=\pm 𝒳_\pm ,[𝒳_+,𝒳_{}]={\displaystyle \frac{q^{2𝒳_0}q^{2𝒳_0}}{qq^1}}=\left[\left[2𝒳_0\right]\right]_q,`$ (75)
called the quantum algebra $`sl_q(2)`$. In the limit $`q1`$, $`\{𝒳\}\{X\}`$ which generate $`sl(2)`$,
$$[X_0,X_\pm ]=\pm X_\pm ,[X_+,X_{}]=2X_0,$$
(76)
the Lie algebra of $`SL(2)`$. The matrices $`\{𝒳_0^{(1/2)},𝒳_+^{(1/2)},𝒳_{}^{(1/2)}\}`$ in (73) provide the fundamental $`2`$-dimensional irreducible representation of the algebra (75) (actually, they also provide the fundamental representation of the generators of $`sl(2)`$ algebra (76)). When a higher dimensional representation of (75) is plugged in the formula (74) $`𝒯`$ becomes a higher dimensional representation of the $`T`$-matrix. For example, the three dimensional representation (57) is obtained by substituting in (74)
$`𝒳_0^{(1)}`$ $`=`$ $`\left(\begin{array}{ccc}\hfill 1& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 1\end{array}\right),`$ (80)
$`𝒳_{}^{(1)}`$ $`=`$ $`\sqrt{\left[\left[2\right]\right]_q}\left(\begin{array}{ccc}0& 0& 0\\ \sqrt{q}& 0& 0\\ 0& 1/\sqrt{q}& 0\end{array}\right),`$ (84)
$`𝒳_+^{(1)}`$ $`=`$ $`\sqrt{\left[\left[2\right]\right]_q}\left(\begin{array}{ccc}0& 1/\sqrt{q}& 0\\ 0& 0& \sqrt{q}\\ 0& 0& 0\end{array}\right),`$ (88)
which provide the three dimensional irreducible representation of the $`sl_q(2)`$ algebra (75), and using the parametrization of $`\{A,B,C,D\}`$ in terms of $`\{\alpha ,\beta ,\gamma \}`$ as given by (70). Note that in the limit $`q1`$ these matrices (88) obey the $`sl(2)`$ algebra (76). As seen from (70), (74) and (75), one can say that for a quantum group the group-parameter space is noncommutative.
The algebra $`sl_q(2)`$ is also often called a quantum group. Actually, the relations in (75) define the generators of the $`q`$-deformation of the universal enveloping algebra of $`sl(2)`$. Hence, the relations (75) are also referred to, more properly, as $`U_q(sl(2))`$, the $`q`$-deformed universal enveloping algebra of $`sl(2)`$. The algebra $`U_q(sl(2))`$, generated by polynomials in $`\{𝒳_0,𝒳_+,𝒳_{}\}`$ obeying the relations (75), is also a Hopf algebra. We shall consider only the coproduct(s) for $`U_q(sl(2))`$. A coproduct for $`sl_q(2)`$ is
$`\mathrm{\Delta }_q\left(𝒳_0\right)=𝒳_01\mathrm{l}+1\mathrm{l}𝒳_0,\mathrm{\Delta }_q\left(𝒳_\pm \right)=𝒳_\pm q^{𝒳_0}+q^{𝒳_0}𝒳_\pm .`$ (89)
It can be easily verified that this comultiplication rule is an algebra isomorphism for $`sl_q(2)`$ :
$`[\mathrm{\Delta }_q\left(𝒳_0\right),\mathrm{\Delta }_q\left(𝒳_\pm \right)]=\pm \mathrm{\Delta }_q\left(𝒳_\pm \right),[\mathrm{\Delta }_q\left(𝒳_+\right),\mathrm{\Delta }_q\left(𝒳_{}\right)]=\left[\left[2\mathrm{\Delta }_q\left(𝒳_0\right)\right]\right]_q.`$ (90)
The most important property of this coproduct is its noncommutativity. Note that the algebra (75) is invariant under the interchange $`qq^1`$ since $`\left[\left[2𝒳_0\right]\right]_q=\left[\left[2𝒳_0\right]\right]_{q^1}`$. However, the comultiplication (89) is not invariant under such an interchange. This means that the comultiplication obtained from (89) by an interchange $`qq^1`$ should also be an equally good comultiplication. It can be verified that the coproduct so obtained, namely,
$`\mathrm{\Delta }_{q^1}\left(𝒳_0\right)=𝒳_01\mathrm{l}+1\mathrm{l}𝒳_0,\mathrm{\Delta }_{q^1}\left(𝒳_\pm \right)=𝒳_\pm q^{𝒳_0}+q^{𝒳_0}𝒳_\pm .`$ (91)
is indeed an algebra isomorphism for $`sl_q(2)`$. This coproduct (91), $`\mathrm{\Delta }_{q^1}`$, is called the opposite coproduct in view of the relation
$$\mathrm{\Delta }_{q^1}(𝒳)=\tau \left(\mathrm{\Delta }_q(𝒳)\right),\text{where}\tau (uv)=vu.$$
(92)
Since
$$\mathrm{\Delta }_{q^1}\mathrm{\Delta }_q,\text{or}\tau (\mathrm{\Delta })\mathrm{\Delta },$$
(93)
the comultiplications $`\mathrm{\Delta }_q`$ and $`\mathrm{\Delta }_{q^1}`$ of $`sl_q(2)`$ are noncommutative. In the limit $`q1`$, the classical $`sl(2)`$ has only a single comultiplication, $`\mathrm{\Delta }(X)=X1\mathrm{l}+1\mathrm{l}X`$ , which is commutative (i.e., $`\tau (\mathrm{\Delta })=\mathrm{\Delta }`$).
One can show that the two comultiplications of $`sl_2(q)`$, namely $`\mathrm{\Delta }_q`$ and $`\mathrm{\Delta }_{q^1}`$, are related to each other by an equivalence relation such that there exists an $`U_q(sl(2))U_q(sl(2))`$, called the universal $``$-matrix, satisfying the relation
$$\mathrm{\Delta }_{q^1}(𝒳)=\mathrm{\Delta }_q(𝒳)^1.$$
(94)
This universal $``$-matrix is the central object of the quantum group theory. In this case it can be shown that
$$=q^{2\left(𝒳_0𝒳_0\right)}\underset{n=0}{\overset{\mathrm{}}{}}\frac{\left(1q^2\right)^n}{\left[\left[n\right]\right]_q!}q^{n(n1)/2}\left(q^{𝒳_0}𝒳_+q^{𝒳_0}𝒳_{}\right)^n.$$
(95)
If we insert the matrix representations of $`\{𝒳\}`$ in this expression for $``$ we get numerical $`R`$-matrices. For example, substituting in (95) the $`2\times 2`$ representation of $`\{𝒳\}`$, given in (73), we get the fundamental $`4`$-dimensional $`R`$-matrix
$$R=\frac{1}{\sqrt{q}}\left(\begin{array}{cccc}q& 0& 0& 0\\ 0& 1& \left(qq^1\right)& 0\\ 0& 0& 1& 0\\ 0& 0& 0& 1\end{array}\right).$$
(96)
Let us now write any $`R`$ in the form
$$R=\underset{i}{}a_iu_iv_i.$$
(97)
It is clear from (95) that this can be done. Now, define
$`R_{12}=R1\mathrm{l},R_{13}={\displaystyle \underset{i}{}}a_iu_i1\mathrm{l}v_i,R_{23}=1\mathrm{l}R.`$ (98)
Then, these satisfy the remarkable relation
$$R_{12}R_{13}R_{23}=R_{23}R_{13}R_{12}.$$
(99)
known as the quantum Yang-Baxter equation, or simply the Yang-Baxter equation (YBE).
We have considered only the simplest example of a quantum group, namely $`SL_q(2)`$, associated with the classical group $`SL(2)`$. There exists a systematic theory of deformation of any classical group. It is also possible, in certain cases, to obtain deformations with several $`q`$-parameters. Actually, the study of quantum groups sheds more light on the structure of the classical group theory. I shall not go further into the details of the formalism of quantum group theory.
Now, I am in a position to mention a few applications of quantum groups and algebras. First, let us see how these things started. Define
$`T_1`$ $`=`$ $`T1\mathrm{l}=\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),`$ (104)
$`T_2`$ $`=`$ $`1\mathrm{l}T=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)\left(\begin{array}{cc}A& B\\ C& D\end{array}\right).`$ (109)
Note that
$$T_1T_2=\left(\begin{array}{cccc}A^2& AB& BA& B^2\\ AC& AD& BC& BD\\ CA& CB& DA& DB\\ C^2& CD& DC& D^2\end{array}\right)T_2T_1=\left(\begin{array}{cccc}A^2& BA& AB& B^2\\ CA& DA& CB& DB\\ AC& BC& AD& BD\\ C^2& DC& CD& D^2\end{array}\right),$$
(110)
because $`\{A,B,C,D\}`$ ae noncommutative. The relation between $`T_1T_2`$ and $`T_2T_1`$ turns out to be
$$RT_1T_2=T_2T_1R.$$
(111)
This type of relation is commonly encountered in the quantum inverse scattering method approach to integrable models in quantum field theory and statistical mechanics. Substituting in (111) $`R`$ from (96), and $`T_1`$ and $`T_2`$ from (109), it is found that equation (111) is a compact way of stating the commutation relations (38) defining the fundamental $`T`$-matrix of $`SL_q(2)`$. What about the commutation relations (75) defining $`sl_q(2)`$? Define
$`L^{(+)}`$ $`=`$ $`\left(\begin{array}{cc}q^{𝒳_0}& \sqrt{q}\left(qq^1\right)𝒳_{}\\ 0& q^{𝒳_0}\end{array}\right),`$ (114)
$`L^{()}`$ $`=`$ $`\left(\begin{array}{cc}q^{𝒳_0}& 0\\ q^{1/2}\left(qq^1\right)𝒳_+& q^{𝒳_0}\end{array}\right),`$ (117)
$`L_1^{(\pm )}`$ $`=`$ $`L^{(\pm )}1\mathrm{l},L_2^{(\pm )}=1\mathrm{l}L^{(\pm )}.`$ (118)
Then, the commutation relations (75), defining the generators of $`sl_q(2)`$, can be stated elegantly as
$`R^1L_1^{(\pm )}L_2^{(\pm )}=L_2^{(\pm )}L_1^{(\pm )}R^1,R^1L_1^{(+)}L_2^{()}=L_2^{()}L_1^{(+)}R^1.`$ (119)
Note that the $`L^{(\pm )}`$-matrices are special realizations of the $`T`$-matrices, i.e., the elements of $`L^{(\pm )}`$-matrices obey the commutation relations required of the $`T`$-matrix elements.
If we define for the $`R`$-matrix in (96),
$$S_1=\stackrel{ˇ}{R}1\mathrm{l},S_2=1\mathrm{l}\stackrel{ˇ}{R},$$
(120)
where
$$\stackrel{ˇ}{R}=PR,P=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 0& 1& 0\\ 0& 1& 0& 0\\ 0& 0& 0& 1\end{array}\right),$$
(121)
then, it is found that
$$S_1S_2S_1=S_2S_1S_2,$$
(122)
which is an alternative form of the YBE (99). For any general $`R`$-matrix the YBE (99) can be put in this form (122). This relation (122) represents a property of the generators of a braid group which is a generalization of the well known symmetric group $`S_n`$. The symmetric group $`S_n`$ is the group of all permutations of $`n`$ given objects. An element of the braid group $`B_n`$ can be depicted as a system of $`n`$ strings joining two sets of $`n`$ points, each set located on a line, the two lines, say top and bottom, being parallel, with over-crossings or under-crossings of the strings. The over-crossings and the under-crossings of the strings make $`B_n`$ an infinite group which will otherwise reduce to $`S_n`$. If $`i`$ and $`i+1`$ are two consecutive points on the top and bottom lines, the string starting at $`i`$ on the top line can reach $`i+1`$ on the bottom line by either under-crossing or over-crossing the string starting at $`i+1`$ on the top line and reaching $`i`$ on the bottom line. The corresponding elements of the braid group are usually denoted by $`\sigma _i`$ and $`\sigma _i^1`$, respectively. The elements $`\left\{\sigma _i\right|i=1,2,\mathrm{}n1\}`$ , generating the braid group $`B_n`$, satisfy two relations,
$$\sigma _i\sigma _j=\sigma _j\sigma _i\text{for}|ij|>1,$$
(123)
and
$$\sigma _i\sigma _{i+1}\sigma _i=\sigma _{i+1}\sigma _i\sigma _{i+1}.$$
(124)
Now, comparing the relations (122) and (124) it is obvious that the solutions of the YBE ($`R`$-matrices), or the quantum algebras, should play a central role in the theory of representations of braid groups. Braid groups have many applications. In mathematics they are useful in the study of complex functions of hypergeometric type having several variables. In physics they appear in knot theory, statistical mechanics, two-dimensional conformal field theory, and so on.
In quantum mechanics the notion of continuous space-time with commutative coordinates is taken over from classical mechanics. This is an assumption. What will happen if at some deeper microscopic level the space-time coordinates themselves are noncommutative? It is clear that to deal with such a situation one will have to use a noncommutative differential calculus and the theory of quantum groups provides the necessary framework as we have seen above. For example, consider the motion of a quantum particle in a two dimensional noncommutative plane with $`X`$ and $`Y`$ as the coordinates. If we take the corresponding conjugate momenta to be proportional to $`\frac{}{X}`$ and $`\frac{}{Y}`$, respectively, then the relations (14) indicate how the two-dimensional quantum mechanical phase-space would be deformed at that level. Thus, the theory of quantum groups would provide the mathematical framework for the future of quantum physics if it turns out that at some deeper microscopic level space-time manifold is noncommutative. Hence there has been a lot of interest in studying the fundamental modifications that would occur in the framework quantum mechanics, relativity theory, Poincaré group, …, etc., if the space-time manifold happens to be noncommutative.
Apart from applications of fundamental nature, such as those mentioned above, there have been many phenomenological applications of quantum algebras in nuclear physics, condensed matter physics, molecular physics, quantum optics, and elementary particle physics. In these applications either an existing model is identified with a quantum algebraic structure, or a standard model is deformed to have an underlying quantum algebraic structure and studied to reveal the new features emerging. To give an idea of such applications, let me mention, as the final example, the $`q`$-deformation of the quantum mechanical harmonic oscillator algebra, also known as the boson algebra. The algebraic treatment of the quantum mechanical harmonic oscillator involves a creation operator $`\left(a^{}\right)`$, an annihilation operator $`(a)`$, and a number operator $`(N)`$, obeying the commutation relations
$$[a,a^{}]=1,[N,a^{}]=a^{},$$
(125)
where $`N`$ is a hermitian operator and $`a^{}`$ is the hermitian conjugate of $`a`$. The energy spectrum of the harmonic oscillator is given by the eigenvalues of the Hamiltonian operator
$$H=\frac{1}{2}\left(aa^{}+a^{}a\right),$$
(126)
in the appropriate units. Taking two such sets of oscillator operators, $`\{a_1,a_1^{},N_1\}`$ and $`\{a_2,a_2^{},N_2\}`$, which are assumed to commute with each other, and defining
$$J_0=\frac{1}{2}\left(N_1N_2\right),J_+=a_1^{}a_2,J_{}=a_2^{}a_1,$$
(127)
it is found that
$$J_0^{}=J_0,J_+^{}=J_{},$$
(128)
and
$$[J_0,J_\pm ]=\pm J_\pm ,[J_+,J_{}]=2J_0.$$
(129)
This Lie algebra (129) is seen to be the same as the $`sl(2)`$ algebra (76) subject to the hermiticity conditions (128) and is known as $`su(2)`$ algebra, the Lie algebra of the group $`SU(2)`$. The $`su(2)`$ algebra is the algebra of three dimensional rotations, or the rigid rotator, with $`\{J_0,J_\pm \}`$ representing the angular momentum operators. The coproduct rule
$$\mathrm{\Delta }\left(J_0\right)=J_01\mathrm{l}+1\mathrm{l}J_0,\mathrm{\Delta }\left(J_\pm \right)=J_\pm 1\mathrm{l}+1\mathrm{l}J_\pm ,$$
(130)
for the algebra (129), obtained by setting $`q=1`$ in (89) (or (91)), represents the rule for addition of angular momenta. Correspondingly, the relations (75) rewritten as
$$[𝒥_0,𝒥_\pm ]=\pm 𝒥_\pm ,[𝒥_+,𝒥_{}]=\left[\left[2J_0\right]\right]_q,$$
(131)
with the hermiticity conditions
$$𝒥_0^{}=𝒥_0,𝒥_+^{}=𝒥_{},$$
(132)
represent the $`su_q(2)`$ (or $`U_q(su(2))`$) algebra or the $`q`$-deformed version of the $`su(2)`$ algebra (129). One can say that $`su_q(2)`$ is the algebra of the $`q`$-rotator. For the $`q`$-angular momentum operators there are two possible addition rules,
$$\mathrm{\Delta }_{q^{\pm 1}}\left(𝒥_0\right)=𝒥_01\mathrm{l}+1\mathrm{l}𝒥_0,\mathrm{\Delta }_{q^{\pm 1}}\left(𝒥_\pm \right)=𝒥_\pm q^{\pm 𝒥_0}+q^{𝒥_0}𝒥_\pm ,$$
(133)
as seen from (89) and (91). Now the interesting fact is that one has a realization of $`su_q(2)`$ generators given by
$$𝒥_0=\frac{1}{2}\left(𝒩_1𝒩_2\right),𝒥_+=A_1^{}A_2,𝒥_{}=A_2^{}A_1,$$
(134)
exactly analogous to the $`su(2)`$ case (127), where the two sets of operators $`\{A_1,A_1^{},𝒩_1\}`$ and $`\{A_2,A_2^{},𝒩_2\}`$ commute with each other and obey, within each set, the algebra
$$AA^{}qA^{}A=q^𝒩,[𝒩,A^{}]=A^{}.$$
(135)
Further, $`𝒩`$ is hermitian and $`\{A,A^{}\}`$ is a hermitian conjugate pair. The $`q`$-deformed oscillator algebra (135) is known as the $`q`$-oscillator or the $`q`$-boson algebra. When $`q1`$ the $`q`$-oscillator algebra (135) reduces to the canonical oscillator algebra (125). As is easy to guess, phenomenological applications of quantum algebras in nuclear and molecular spectroscopy involve the substitution of harmonic oscillator model by the $`q`$-oscillator model and the rigid rotator model by the $`q`$-rotator model. Such applications lead to impressive results showing that the actual vibrational-rotational spectra of nuclei and molecules can be fit into schemes in which the number of phenomenological $`q`$-parameters required are very much fewer than the number of traditional phenomenological parameters required to fit the same spectral data. Somehow such $`q`$-deformed models seem to take into account more efficiently the anharmonicity of vibrations and the nonrigidity of rotations in nuclear and molecular systems.
I wish to thank the Management, and Prof. Dr. K. V. Parthasarathy (Head, Department of Mathematics), of the Ramakrishna Mission Vivekananda College, for giving me the privilege of delivering the Narayani, N. and Kadayam S. Sankaran Memorial Lectures for 1999.
|
warning/0003/quant-ph0003120.html
|
ar5iv
|
text
|
# 1 Introduction
## 1 Introduction
Quantum entanglement , is a property of bipartite systems by virtue of which the spatially separated subsystems exhibit strong correlations among themselves that cannot be explained in classical terms. Often termed as quantum nonlocality, the subject had been under extensive study and debate since 1935 when Einstein, Podolsky and Rosen questioned the completeness of quantum mechanics . Later, in 1964 in his pioneering work John Bell showed that such nonlocal correlations are *inherently* quantum mechanical and any model admitting the description of local hidden variables (or one may prefer to say local realism) fails to explain such correlations. A consequence of quantum superposition principle, thus entanglement has been the most celebrated manifestation of quantum nonlocality .
In recent years there is a renewed interest in entanglement primarily because of its newly found applications in quantum information processing and quantum computing and as it turned out, many of the applications would have either been impossible or less efficient classically. A necessary condition for faithful implementation of a quantum information protocol is that the parties share a maximally entangled state (MES), each of them having access to their respective subsystem whereby they perform local quantum operations and communicate among themselves one way or both ways via classical channel to implement the concerned protocol. The local operations include unitary operations, von neuman measurements and generalized measurements which may involve ancillary systems Alice and Bob might prepare in their laboratories. For two qubit systems a class of maximally entangled states are called Bell states that are defined by
$`|\mathrm{\Phi }^\pm `$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(|00\pm |11\right)`$ (1)
$`|\mathrm{\Psi }^\pm `$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(|01\pm |10\right)`$ (2)
The four Bell states are all equally good for faithful quantum communication. For instance, quantum teleportation with fidelity one is achieved only with maximally entangled states, viz. by any one of the four Bell states. They also share the property of being local unitary equivalent (they are mapped onto each other by local pauli rotations).
As noted above maximally entangled states are crucial for faithful information processing. In practice, however, entanglement is susceptible to local interactions with the environment. Such dissipative interactions may take place during encoding/decoding processes, transmission and sharing of entanglement. This give rise to a mixed entangled state or a separable state depending on the nature, strength and duration of the interaction. The noisy states thus formed are of little or no use for information processing. However one can still generate a fewer number of maximally entangled states from an ensemble of mixed entangled states by applying the distillation protocols .
Despite being responsible for destroying quantum coherence there is a genuine positive side of dissipative effects as shown recently by Badziag *et al*. . They demonstrated that a family of non teleporting mixed entangled states can be made teleporting through dissipative intaeractions with the environment via amplitude damping channel. Here it may be worth mentioning that ADC is fairly rich in structure and was shown to be having many interesting properties. For instance Bennett *et al.* showed that if the noisy channel is ADC then by entangling the transmission bits one can increase the receiver’s capability of correct inference .
In the present work we uncover some curious features when a pure but maximal entanglement interacts with the environment via ADC. The problem that we consider in this paper is the following. Suppose we have a bipartite maximally entagled state, say, $`\rho `$ ($`\rho `$ being one of the four Bell states), and the qubits of the entangled pair undergo local interactions with their respective environment via amplitude damping channel. There are two possibilities:
*Case 1*: Only one qubit gets affected: $`\rho \rho _1`$ (a mixed entangled state).
*Case 2:* Both the qubits are affected: $`\rho \rho _2`$ (a mixed entangled state).
Typically, both the cases may arise while sending entanglement across a noisy channel-the so called transmission errors or due to some coding / decoding process via some noisy channel. Whatever be the source of errors we ask which of the two states $`\rho _1`$ and $`\rho _2`$, is more useful for teleportation ?
Since both qubits are affected in the second case, one may be tempted to speculate that $`\rho _2`$ is always qualitatively worse that $`\rho _1`$ in terms of efficiency to perform teleportation. To our surprise we find that the state $`\rho _2`$ is *sometimes* “better” than $`\rho _1`$ having a higher teleportaion fidelity ! Quite interestingly this effect is observed *only* if the Bell state that we wish to share is either of the two states $`|\mathrm{\Phi }^\pm `$. The effect is also shown to depend on the strength of the amplitude damping channel. Moreover we show that if $`\rho _1`$ is a non teleporting state then it can *always* be made teleporting by subjecting the other qubit (which didn’t interact with the environment before) to similar dissipative interaction. An interesting feature of this effect is that the teleportation fidelity of the noisy states that are formed while sharing the other two Bell states *$`|\mathrm{\Psi }^\pm `$, cannot* be enhanced by the same prescription. Here we note that since the process is trace preserving, we have thus been able to identify and parametrize the class of non teleporting states having fully entangled fraction $`f`$ less than 1/2, for which we can increase $`f`$ to greater than half so that they become suitable for distillation *without discarding any pair,* as opposed to filtering procedures . This shows that *an a priori knowledge of the noisy channel might be helpful to decide which Bell state needs to be shared between Alice and Bob.*
As noted earlier, Badziag *et al.* obtained a class of noisy states whose teleportation fidelity can be enhanced by subjecting one qubit to undergo dissipative interactions with its local environment via ADC. *What is the origin of such noisy states ?* Our results show that the class of states obtained in Ref. belong to the family of noisy states that result when any of the particular Bell states, namely $`|\mathrm{\Phi }^\pm `$ interacts with the environment via the ADC.
This paper is arranged in the following way. In Sec. II we review some known results relating fully entangled fraction to teleportation fidelity and distillation. In Sec. III the action of an ADC on a qubit is reviewed. The generation of noisy states while sharing a Bell state across ADC is described in Sec. IV. In Sec. V we discuss the improvement of noisy states via dissipation. We conclude with some remarks and discussion in Sec. VI.
## 2 Entanglement Fidelity, Fully Entangled Fraction, Teleportation Fidelity and Distillation threshold
In this paper reliability for teleportation will be the criterion for judging the quality the noisy entangled states. The following results will be useful .
(1) Fully entangled fraction $`f`$ of a bipartite entangled state $`\rho `$ in $`C^2C^2`$ is defined as
$$f\left(\rho \right)=\underset{\psi }{\mathrm{max}}\psi |\rho \psi $$
(3)
where the maximum is taken over all maximally entangled states $`|\psi `$ . For $`\rho `$ to be useful for quantum teleportation we must have $`f>\frac{1}{2}`$ . It was also shown that in the standard teleportation scheme the maximal fidelity $`F`$ achievable is related to $`f`$ by
$$F=\frac{2f+1}{3}$$
(4)
(2) For the states with $`f\frac{1}{2}`$ one cannot directly apply the distillation protocol by Bennett *et al.* . For those states one first resorts to filtering procedures to enhance the value of $`f`$ from $`f\frac{1}{2}`$ to $`f>\frac{1}{2}`$.
## 3 Action of the Amplitude Damping Channel on a Qubit
Here we briefly review the action of ADC on a qubit. Details can be found in Ref. . The amplitude damping channel describes the interaction of a two level atom with the electromagnetic field (environment). Specifically the decay of an excited state of a two level atom by spontaneous emission of a photon in presence of an e.m. field is what modelled by this channel. The unitary transformation that governs the evolution of the system and the environment (the environment can be always taken to be in some pure state without any loss of generality) is defined by,
$`|0_A|0_E`$ $``$ $`|0_A|0_E`$ (5)
$`|1_A|0_E`$ $``$ $`\sqrt{1p}|1_A|0_E+\sqrt{p}|0_A|1_E`$ (6)
Physically this implies that if an atom is in an excited state $`|1_A`$, with probability $`p`$ it makes a transition to the ground state $`|0_A`$ with the emission of a photon. The environment as a result also makes a transition from the “no-photon” state $`|0_E`$ to the “one-photon” state $`|1_E`$. Tracing out the environment we obtain the Kraus operators, with the Kraus operators $`K_i`$:
$$K_1=\left(\begin{array}{cc}1& 0\\ 0& \sqrt{1p}\end{array}\right),K_2=\left(\begin{array}{cc}0& \sqrt{p}\\ 0& 0\end{array}\right)$$
(7)
satisfying the completeness relation
$$\underset{i=1}{\overset{2}{}}K_i^{}K_i=1$$
(8)
The density matrix $`\rho `$ of the quantum system then evolves as
$$\rho \rho ^{}=S\left(\rho \right)=\underset{i=1}{\overset{2}{}}K_i\rho K_i^{}$$
(9)
The above equation basically defines a linear map which takes a density matrix to another density matrix (a superoperator). It is clear that the amplitude damping channel is characterized by the parameter $`p`$, which denotes the dissipation strength when a qubit interacts with the environment via this channel.
## 4 Sending Bell States across the Amplitude Damping Channel
In this section we consider the problem of sending the Bell states across the amplitude damping channel. We assume that Alice prepares one of the Bell states locally and send one of the qubits to Bob. This qubit interacts with the environment via ADC during the transmission. The initial Bell state is thus transformed to a mixed entangled state. Let $`\rho `$ be the density operator representing one of the four Bell states, then the interaction with the environment via ADC is described by the following transformation:
$$\rho \stackrel{~}{\rho }\left(p\right)=S\left(\rho \right)=W_0\rho W_0^{}+W_1\rho W_1^{}$$
(10)
where $`W_iIK_i`$ are given by
$$W_0=\left[\begin{array}{cccc}1& 0& 0& 0\\ 0& \sqrt{1p}& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& \sqrt{1p}\end{array}\right];W_1=\left[\begin{array}{cccc}0& \sqrt{p}& 0& 0\\ 0& 0& 0& 0\\ 0& 0& 0& \sqrt{p}\\ 0& 0& 0& 0\end{array}\right]$$
(11)
The following notation will be used henceforth to represent the Bell states
$$|\mathrm{\Phi }^\pm \mathrm{\Phi }^\pm |\rho _\pm ^{};|\mathrm{\Psi }^\pm \mathrm{\Psi }^\pm |\rho _\pm ^{}$$
(12)
It is clear that such an interaction turns the maximally entangled state, originally shared by Alice and Bob into a mixed entangled state,
$$\stackrel{~}{\rho }_\pm ^{}=\left[\begin{array}{cccc}1& 0& 0& \pm \sqrt{1p}\\ 0& 0& 0& 0\\ 0& 0& p& 0\\ \pm \sqrt{1p}& 0& 0& 1p\end{array}\right]:\stackrel{~}{\rho }_\pm ^{}=\frac{1}{2}\left[\begin{array}{cccc}p& 0& 0& 0\\ 0& 1p& \pm \sqrt{1p}& 0\\ 0& \pm \sqrt{1p}& 1& 0\\ 0& 0& 0& 0\end{array}\right]$$
(13)
It turns out that the fully entangled fraction $`f`$ of the four resulting mixed states is the same:
$$f\left(\stackrel{~}{\rho }_\pm ^{}\right)=f\left(\stackrel{~}{\rho }_\pm ^{}\right)=\frac{1}{4}\left(1+\sqrt{1p}\right)^2=f\left(\stackrel{~}{\rho }\right)$$
(14)
where
$$f\left(\stackrel{~}{\rho }_\pm ^{}\right)=\mathrm{\Phi }^\pm \left|\stackrel{~}{\rho }_\pm ^{}\right|\mathrm{\Phi }^\pm ;f\left(\stackrel{~}{\rho }_\pm ^{}\right)=\mathrm{\Psi }^\pm \left|\stackrel{~}{\rho }_\pm ^{}\right|\mathrm{\Psi }^\pm $$
(15)
Here we note that fully entangled fraction is same as that of entanglement fidelity . Of course this is not always the case, as we will show later when we subject the second qubit to dissipation for possible improvement of the noisy states. Since entanglement fidelity measures how well the entanglement has been preserved while interacting with a noisy channel we conclude that till this end all the four Bell states are equally corrupted. This is also expected as the noisy channel acts locally on one of the qubits whose state given by the reduced density matrix is the same for all the four Bell states.
We now consider the usefulness of the mixed states given by Eq. (15) to perform quantum teleportation. These family of mixed states are useful for teleportation when $`f\left(\stackrel{~}{\rho }\right)>1/2`$ , from which one easily obtains the restriction on the parameter $`p`$
$$\begin{array}{cc}f\left(\stackrel{~}{\rho }\right)>\frac{1}{2}& p<2\sqrt{2}2\end{array}$$
(16)
Eq. (16) simply states that for all values of $`p2\sqrt{2}2`$, $`f\left(\stackrel{~}{\rho }\right)1/2`$ the family of mixed states $`\stackrel{~}{\rho }_\pm ^{},\stackrel{~}{\rho }_\pm ^{}`$ so obtained are not useful for teleportation. Here “not useful” is understood as not better than what can be done classically.
## 5 Improving the Noisy Bell States by letting the unaffected Qubit to interact with the Local Environment via Amplitude Damping Channel
In this section we discuss the possible improvement of the noisy states given by Eq. (15) by subjecting the second qubit to similar dissipative interaction.
### 5.1 Damping Parameters same for the two qubits
We now allow Alice’s qubit to interact with the environment via the amplitude damping channel. Here we also assume that the strength of dissipation affecting qubits locally is the same for both. This is not a necessary assumption but only simplifies the degree of algebraic complexity though it does not capture all the intricacies involved. The general case where the damping parameter is different for the two qubits will be treated in the next subsection.
This interaction is described by the following transformation:
$$\stackrel{~}{\rho }(p)\stackrel{~}{\stackrel{~}{\rho }}(p)=S^{}\left(\stackrel{~}{\rho }\right)=W_0^{}\rho W_0^{}+W_1^{}\rho W_1^{}$$
(17)
where $`W_i^{}=K_iI`$ and $`\stackrel{~}{\rho }`$ is the density operator representing one of the noisy Bell states: $`\{\stackrel{~}{\rho }_\pm ^{},\stackrel{~}{\rho }_\pm ^{}\}`$.
The fully entangled fraction of the $`\stackrel{~}{\stackrel{~}{\rho }}`$ states are now given by:
$`f\left(\stackrel{~}{\stackrel{~}{\rho }}_\pm ^{}\right)`$ $`=`$ $`1p+{\displaystyle \frac{1}{2}}p^2`$ (18)
$`f\left(\stackrel{~}{\stackrel{~}{\rho }}_\pm ^{}\right)`$ $`=`$ $`\begin{array}{cc}1p& p\frac{2}{3}\\ \frac{p}{2}& p\frac{2}{3}\end{array}`$ (21)
First we check if there is a range of $`p`$ such that $`f\left(\stackrel{~}{\stackrel{~}{\rho }}\right)>1/2`$. If so then we would like to know if that range has an overlap with the range $`p2\sqrt{2}2`$ because only then that would imply that further application of a dissipative interaction can indeed enhance the teleportation fidelity of the “one qubit affected” noisy states that are formed via the interaction defined by Eq. (10). It is easy to check that
$`f\left(\stackrel{~}{\stackrel{~}{\rho }}_\pm ^{}\right)`$ $`>{\displaystyle \frac{1}{2}}`$ $`p,p1`$ (22)
$`f\left(\stackrel{~}{\stackrel{~}{\rho }}_\pm ^{}\right)=`$ $`\begin{array}{cc}1p& >\frac{1}{2}\\ \frac{p}{2}& <\frac{1}{2}\end{array}`$ $`\begin{array}{cc}p<\frac{1}{2}& \\ p,p1& \end{array}`$ (27)
Recall that for all $`p2\sqrt{2}2`$, $`f\left(\stackrel{~}{\rho }\right)1/2`$. Since from Eq. (20) it follows that $`f\left(\stackrel{~}{\stackrel{~}{\rho }}_\pm ^{}\right)>\frac{1}{2}p,p1\text{ }`$, therefore $`f\left(\stackrel{~}{\stackrel{~}{\rho }}_\pm ^{}\right)\text{ }>\frac{1}{2}p2\sqrt{2}2`$ which implies that teleportation fidelity of *all* non teleporting states $`\stackrel{~}{\rho }_\pm ^{}\left(p\right);p2\sqrt{2}2`$ can indeed be enhanced by subjecting the second qubit to dissipation. Here we also note that $`f\left(\stackrel{~}{\stackrel{~}{\rho }}_\pm ^{}\right)f\left(\stackrel{~}{\rho }_\pm ^{}\right)p0.80585`$. This implies that teleportation fidelity of the teleporting states $`\stackrel{~}{\rho }_\pm ^{}\left(p\right);\mathrm{\hspace{0.33em}0.80585}p<2\sqrt{2}2`$ can also be enhanced by this method. Thus the two-qubit affected state $`\stackrel{~}{\stackrel{~}{\rho }}_\pm ^{}`$ is *better* (i.e. improved) with a higher teleportation fidelity than single qubit affected state $`\stackrel{~}{\rho }_\pm ^{}`$ whenever $`p>0.80585`$. This is surprising in the sense that the dissipative interaction which spoils entanglement in the first step is utilised to improve the quality of the mixed state by applying it to the second subsystem.
However Eqs. (19) and (21) show that the same prescription though curious doesn’t work to turn non teleporting states into teleporting ones. Although $`f\left(\stackrel{~}{\stackrel{~}{\rho }}_\pm ^{}\right)=p/2f\left(\stackrel{~}{\rho }_\pm ^{}\right)p8/9`$, the enhancement is not sufficient to transform non teleporting states to teleporting states as $`f\left(\stackrel{~}{\stackrel{~}{\rho }}_\pm ^{}\right)=p/2`$ is always less than or equal to $`\frac{1}{2}`$.
It may be worth stressing that in showing the above effect we have taken the damping parameter to be the same for both the qubits. We will now consider the case where the damping parameters are different.
### 5.2 Damping parameters different for the two qubits
Let us now treat the problem in a more general way. Here, as before, we will assume that first Bob’s qubit undergoes dissipative interaction with the environment and the parameter *p* of the amplitude damping channel will now be denoted by $`p_b`$. After that we allow the qubit that belongs to Alice also to interact with the environment via amplitude damping channel and the channel parameter now taken to be different from $`p_b`$, is denoted by $`p_a`$. Nevertheless the conclusion remains the same except for the fact that in this case we have the freedom of varying the strength of dissipation for the second qubit which allows us to maximize the enhancement of teleportation fidelity.
The fully entangled fraction $`f\left(\stackrel{~}{\rho }\right)`$ of the resulting mixed state $`\stackrel{~}{\rho }(p_b)`$ formed via interaction given by Eq. (10) is,
$$f\left(\stackrel{~}{\rho }_\pm ^{}\right)=f\left(\stackrel{~}{\rho }_\pm ^{}\right)=\frac{1}{4}\left(1+\sqrt{1p_b}\right)^2=f\left(\stackrel{~}{\rho }\right)$$
(28)
This is same as Eq. (14), the only difference is that $`p`$ is now denoted by $`p_b`$. Now when we allow Alice’s qubit also to interact with the environment described by Eq. (17) , $`\stackrel{~}{\rho }(p_b)`$ states are transformed to $`\stackrel{~}{\stackrel{~}{\rho }}(p_a,p_b)`$.
Irrespective of the source state the fully entangled fraction is the same when only one qubit is affected. The remaining analysis involving subjection of the second qubit to dissipation for possible enhancement of teleportation fidelity is divided into two parts. The first one corresponds to the cases when the source states are $`|\mathrm{\Phi }^\pm \mathrm{\Phi }^\pm |\rho _\pm ^{}`$ and the second part corresponds to the cases when the source states are $`|\mathrm{\Psi }^\pm \mathrm{\Psi }^\pm |\rho _\pm ^{}`$.
#### 5.2.1 Part 1
Let us first recall that for $`p_b2\sqrt{2}2`$, $`f\left(\stackrel{~}{\rho }\right)1/2`$, i.e., the states $`\stackrel{~}{\rho }_\pm ^{}`$ are non teleporting. After the second interaction for which the damping parameter is $`p_a`$, the fully entangled fraction corresponding to the $`\stackrel{~}{\stackrel{~}{\rho }}_\pm ^{}(p_a,p_b)`$ states are given by:
$$f\left(\stackrel{~}{\stackrel{~}{\rho }}_\pm ^{}\right)=\frac{1}{4}\left[p_ap_b+\left(1+\sqrt{\left(1p_a\right)\left(1p_b\right)}\right)^2\right]$$
(29)
5.2.1.1 Condition for “two qubit affected” state having larger teleportation fidelity than the “one qubit affected” state
Our objective is to obtain the condition such that the inequality
$$f\left(\rho ^{\prime \prime }\right)>f\left(\rho ^{}\right)$$
(30)
is satisfied. To be precise we are now looking for those values of $`p_a`$ , such that for $`p_b2\sqrt{2}2`$, the previous inequality is satisfied. It turns out that for
$$p_a<\frac{4\left[\sqrt{1p_b}\left(2p_b1\right)\left(1p_b\right)\right]}{\left(2p_b1\right)^2}=g\left(p_b\right)$$
(31)
the inequality $`f\left(\rho ^{\prime \prime }\right)>f\left(\rho ^{}\right)`$ is satisfied. However it remains to be ensured that a range of values of $`p_a`$ is obtained for $`p_b2\sqrt{2}2`$. First observe that $`g(3/4)=0`$. This in turn implies we can always find some suitable range of $`p_a`$ such that the inequality (20) is satisfied provided $`p_b>3/4`$. Stated more explicitly this means that only when Bob’s qubit interacted with the environment via the amplitude damping channel characterized by the parameter $`p_b>3/4`$ , then the possibility of improving the corrupted state by allowing the other qubit to interact, arises. Hence the two qubit affected states are thereby made better than the one qubit affected state. One may also note that not only non teleporting states can be made teleporting but also teleporting states with poor fidelity are made better. We now look for the condition when maximum enhancement of teleportation fidelity is achieved.
5.2.1.2 Maximum achievable teleportation fidelity
Let us now assume $`p_b>3/4`$. Thus we have at our disposal a range of possible values of $`p_a`$ satisfying (24) and the inequality (25). Then another important question remains to be answered: For which value of $`p_a`$, $`f\left(\rho ^{\prime \prime }\right)`$ is the maximum? That is, for which value $`p_a`$ highest fidelity for teleportation can be achievd?
One can easily check that (23) is maximised for
$$p_a=\frac{p_b\left(4p_b3\right)}{\left(2p_b1\right)^2}$$
(32)
Here it is interesting to note that (26) immediately gives a lower bound on the value of the parameter $`p_b`$, which is 3/4 and this bound we have obtained before. Substituting (26) in (23) one obtains the maximum value of fully entangled fraction that can be achieved:
$$f_{max}\left(\rho ^{\prime \prime }\right)=\frac{p_b^2}{2\left(2p_b1\right)}$$
(33)
One should note that $`f_{max}>1/2`$.
5.2.1.3 Non teleporting states can always be made teleporting
To begin with Alice and Bob initially shared a maximally entangled state. If Bob’s qubit gets affected by the ADC , $`p_b>3/4`$ then the fidelity of the resulting noisy state can be improved by subjecting the unaffected qubit to similar dissipation characterized by the parameter $`p_a`$ in accordance with Eq. (25). The maximum “fully entangled fraction” or the maximum fidelity of teleportation is obtained by substituting (26) in the expression (23).
Since for $`p_b2\sqrt{2}1`$ the state $`\rho ^{}(p_b)`$ is not capable of teleportation better than what can be achieved classically. It is now clear that those non teleporting states can always be made teleporting by allowing the other qubit to interact with the environment with an appropriate choice of the parameter $`p_a`$. In particular the maximum fidelity can be obtained if we choose $`p_a`$ in accordance with (26). Even if the state $`\rho ^{}(p_b)`$ is a teleporting one (which it is, if $`3/4<p_b<2\sqrt{2}2`$) still its fidelity can be enhanced by allowing the other qubit to interact with the environment for any value of $`p_a`$ defined by (25).
#### 5.2.2 Part 2
We now analyze what happens when the source states are the superposition of antiparallel spins. First we note that when one qubit is affected the fully entangled fraction is given by Eq. (10) or Eq. (22). After the second qubit is affected, the interaction being governed by an appropriate superoperator, the density operator of the mixed state is given by
$$\stackrel{~}{\stackrel{~}{\rho }}_\pm ^{}=\left[\begin{array}{cccc}p_a+p_b& 0& 0& 0\\ 0& 1p_b& \pm \sqrt{\left(1p_a\right)\left(1p_b\right)}& 0\\ 0& \pm \sqrt{\left(1p_a\right)\left(1p_b\right)}& 1p_a& 0\\ 0& 0& 0& 0\end{array}\right]$$
Now recall the definition of fully entangled fraction. One can easily see that two possible candidates for fully entangled fraction are
$`f_1`$ $`=`$ $`{\displaystyle \frac{p_a+p_b}{4}}`$
$`f_2`$ $`=`$ $`{\displaystyle \frac{1}{4}}\left(\sqrt{1p_a}+\sqrt{1p_b}\right)^2`$
The fully entangled fraction is given by $`\mathrm{max}[f_1,f_2]`$. Whether $`f_1>f_2`$ or vice versa depends on the particular values of $`p_a`$ and $`p_b`$, i.e., whether $`f_1`$ is the fully entangled fraction or $`f_2`$ is depends on the specific choices of the damping parameters.
Let us first assume that $`f_1>f_2`$. Now observe that $`f_11/2`$. Therefore teleportation fidelity can never exceed 2/3. Hence non teleporting states can never become teleporting. Now if $`f_2>f_1`$, then note that $`f_2(p_a,p_b)f\left(\stackrel{~}{\rho }_\pm ^{}\right)p_a,p_b`$. Therefore whatever be the case the above two observations guarantee that by subjecting the second qubit to a similar dissipative interaction, the non teleporting states cannot be made teleporting.
## 6 Remarks and Conclusion
We know that given any mixed state the corresponding teleportation fidelity has contribution from two parts: (a) the quantum part i.e. the shared entanglement and (b) the classical part, i.e. the classical correlations. Since entanglement cannot be increased by TPLOCC, an enhancement in teleportation fidelity suggests that, for the states on the border line of non teleporting and teleporting, it is the classical part that contributes more heavily than the quantum part so that even a non teleporting state becomes teleporting in spite of losing entanglement. Thus TPLOCC which spoils entanglement but nevertheless can sometimes redistribute the classical correlations in a helpful way.
We now discuss the implications of our result in the context of distilling entanglement from the noisy states. For distillation of mixed entangled states having fully entangled fraction $`f\frac{1}{2}`$one needs to apply filtering methods so that $`f`$ exceeds $`\frac{1}{2}`$. This is important because only then the distillation protocol of Bennett et al could be applied. However filtering being non trace preserving involves measurement followed by post selection which necessarily discards a large number of available pairs. In view of this our results provide a significant improvement for a large class of states for which fully entangled fraction can be made greater than half without resorting to filtering implying that we need not discard pairs.
We have shown that the noisy entangled states constructed from the source states by similar (meaning same for all the source states) dissipative interactions, can indeed have very different properties even though the source states are local unitary equivalent. This observation might be helpful while trying to establish entanglement when the noisy quantum channel is the amplitude damping channel. If we are provided with either of the states $`|\mathrm{\Psi }^\pm `$ then instead of trying to share it we may first convert it to $`|\mathrm{\Phi }^\pm `$ state by applying pauli rotations locally. The reason is even if the entanglement gets spoiled at its worse so that the mixed states are non teleporting we can still make them suitable for teleportation by TPLOCC. Besides being made useful for teleportation these states also allow for direct application of the Bennett protocol for distillation . Thus our results do indicate that a priori knowledge of the noisy channel may decide sometimes which Bell state is to be shared. One of course should try to share the case specific Bell state for which the resulting mixed states allow for further useful manipulations.
I would like to thank Daniel Terno, Ujjwal Sen and Guruprasad Kar for helpful discussions.
|
warning/0003/hep-th0003044.html
|
ar5iv
|
text
|
# String-generated quartic scalar interactionsTalk delivered by R. Marotta
## 1 Tree scalar diagrams from strings
The starting point is the planar tree scattering amplitude of four on-shell bosonic open string tachyons with momenta $`p_1,\mathrm{},p_4`$ each satisfying the mass-shell condition $`p^2=m^2=\frac{1}{\alpha ^{}}`$:
$$A_4^{(0)}(p_1,p_2,p_3,p_4)=\text{Tr}\left[\lambda ^{a_1}\lambda ^{a_2}\lambda ^{a_3}\lambda ^{a_4}\right]C_0𝒩_t^4$$
$$\times _0^1dz(1z)^{2\alpha ^{}p_2p_3}z^{2\alpha ^{}p_3p_4}$$
(1.1)
where the Koba-Nielsen variables relative to the tachyons labelled by $`1,2,4`$ have been respectively fixed at $`+\mathrm{},1`$ and $`0`$.
According to the corner of moduli space where the low-energy limit of the amplitude (1.1) is performed, one can recover, for instance, $`\mathrm{\Phi }^3`$\- or $`\mathrm{\Phi }^4`$-scalar diagrams.
In order to understand which regions in moduli space lead to the different field theory diagrams, one can use the so-called sewing and cutting procedure. This consists in starting from a string diagram and in cutting it in three-point vertices; next we fix the legs of each three-point vertex at $`+\mathrm{}`$, $`1`$ and $`0`$. Then we reconnect the diagram by inserting between two three-point vertices a suitable propagator acting as a well specified projective transformation. This is chosen in such a way that its fixed points are just the Koba-Nielsen variables of the two legs that have to be sewn. The geometric role of the propagator is to identify the local coordinate systems defined around the punctures to be sewn.
The amplitude in Eq. (1.1) can be expressed also in terms of the Green functions $`𝒢^{(0)}(z_i,z_j)`$, defined on the world-sheet in the following way:
$$𝒢^{(0)}(z_i,z_j)=\text{log}\left(z_iz_j\right)$$
(1.2)
A necessary intermediate step for deriving tree scalar $`\mathrm{\Phi }^4`$-diagrams is to generate tree diagrams of $`\mathrm{\Phi }^3`$-theory. With reference to the four-tachyon tree string diagram, one can see that it can be obtained by sewing two three-point vertices as shown in Fig. (1). We sew the leg corresponding to the point $`0`$ in the vertex at the left hand in Fig. (1a) to the leg corresponding to the point $`+\mathrm{}`$ in the one at the right hand through a propagator corresponding to the projective transformation
$$S(z)=Az$$
(1.3)
which has $`0`$ and $`+\mathrm{}`$ as fixed points and the parameter $`A`$, with $`0A1`$, as multiplier. Performing the sewing means, in this procedure, to transform only the punctures of the three-point vertex at the right hand in Fig. (1a) through (1.3), hence the puncture $`z_3=1`$ transforms into $`S(1)=A`$ while the other two punctures remain unchanged.
In general, after the sewing has been performed, the Koba-Nielsen variables become functions of the parameter $`A`$ appearing in the projective transformation (1.3). It is possible to give a simple geometric interpretation to this parameter, if a correspondence is established between the projective transformation in Eq. (1.3) and the string propagator, written in terms of the operator $`e^{\tau (L_01)}`$. The latter indeed propagates an open string through imaginary time $`\tau `$ and creates a strip of length $`\tau `$. In fact the change of variable $`z=e^\tau `$ allows the string propagator to be written as
$$\frac{1}{L_01}=_0^1𝑑zz^{L_02}$$
$$=_0^{\mathrm{}}𝑑\tau \mathrm{exp}\left(\tau \alpha ^{}\left[p^2+\frac{1}{\alpha ^{}}(N1)\right]\right)$$
(1.4)
and to establish the following relation between $`\tau `$ and $`A`$:
$$\tau =\text{log}A.$$
(1.5)
The multiplier $`A`$ results to be therefore related to the length of the strip connecting two three-vertices.
On the other hand, since we want to reproduce tree $`\mathrm{\Phi }^3`$-theory diagrams we have to consider a low-energy limit of string amplitudes in which only tachyons propagate as intermediate states. This is achieved observing from (1.4) that the only surviving contribution in the limit $`\alpha ^{}0`$ with $`\tau \alpha ^{}`$ kept fixed is the one coming from the level $`N=0`$, i.e. from tachyons with fixed mass given by $`m^2=\frac{1}{\alpha ^{}}`$. It is obvious that this also corresponds to the limit $`\tau \mathrm{}`$ and hence, from (1.5), to $`A0`$. From these considerations it seems natural to introduce the variable $`x=\tau \alpha ^{}`$ in terms of which the string propagator (1.4), reproduces, in the above mentioned limit, the scalar propagator
$$\frac{1}{p^2+m^2}=_0^{\mathrm{}}𝑑xe^{x(p^2+m^2)}$$
with $`x`$ being interpreted as the Schwinger proper time.
From a geometrical point of view, one can imagine that the the strip connecting the two three-vertices, in this field theory limit, becomes “very long and thin”, so that only the lightest states propagate.
By rewriting the amplitude (1.1) in terms of the Schwinger parameter $`x`$ or, equivalently, in terms of the multiplier $`A`$ we finally get :
$$A_4^{(0)}=\frac{1}{8}\text{Tr}\left[\lambda ^{a_1}\lambda ^{a_2}\lambda ^{a_3}\lambda ^{a_4}\right]\frac{g_{\varphi ^3}^2}{\left[(p_1+p_2)^2+m^2\right]}$$
(1.6)
where it has been used the well-known relation between $`g_s`$ and $`g_{\varphi ^3}`$ :
$$g_{\varphi ^3}=16g_s(2\alpha ^{^{}})^{\frac{d6}{4}}.$$
(1.7)
We are going now to consider a suitable limit of the string four-tachyon amplitude which can reproduce the diagram corresponding to the tree four-point vertex of $`\mathrm{\Phi }^4`$-theory. With reference again to the Fig. (1), this diagram has to correspond to a limit in which the length of the tube connecting the two three-vertices composing the string diagram goes to zero in the limit $`\alpha ^{}0`$, i.e.
$$\tau =\frac{x}{\alpha ^{}}=\text{log}A0$$
This corresponds to the limit $`A1`$, and hence $`z1`$ or, equivalently, $`z_3z_2`$.
In this limit the Green function $`𝒢^{(0)}(z_2,z_3)`$ is divergent. We regularize it by introducing a cut-off $`ϵ`$ on the world-sheet so that
$$\underset{z_2z_3}{lim}\text{log}\left[(z_2z_3)+ϵ\right]=\text{log}ϵ$$
and
$$\underset{\alpha ^{}0}{lim}\alpha ^{}\text{log}ϵ=0$$
We consider therefore the amplitude $`A_4^{(0)}`$ in Eq. (1.1) in the field theory limit defined by:
$$A=z=1ϵ,$$
$`\alpha ^{}0`$ and $`x=\alpha ^{}\text{log}ϵ0.`$ (1.8)
in which it reduces to
$$2^4\text{Tr}\left[\lambda ^{a_1}\lambda ^{a_2}\lambda ^{a_3}\lambda ^{a_4}\right]g_s^2(2\alpha ^{})^{\frac{d4}{2}}.$$
(1.9)
The complete amplitude is obtained by performing the sum over non cyclic permutations, finally getting a result coincident with the color ordered vertex generated by the following scalar field theory:
$$=\text{Tr}\left[^\mu \varphi _\mu \varphi +m^2\varphi ^2\frac{g_{\varphi ^4}}{4!}\varphi ^4\right]$$
(1.10)
obtaining the matching condition
$$g_{\varphi ^4}=4g_s(2\alpha ^{})^{d/22}.$$
(1.11)
## 2 One-loop $`\mathrm{\Phi }^4`$-diagrams from strings
### 2.1 Tadople diagram
In this subsection we show how the tadpole diagram in $`\mathrm{\Phi }^4`$-theory can be derived from string theory. The starting point will be, this time, the color ordered two-tachyon planar amplitude at one-loop:
$$A_2(p_1,p_2)=N\text{Tr}\left[\lambda ^{a_1}\lambda ^{a_2}\right]C_1\left[2g_s(2\alpha ^{})^{\frac{d2}{4}}\right]^2$$
$$\times _0^1\frac{dk}{k^2}[\frac{1}{2\pi }\text{log}k]^{\frac{d}{2}}\underset{n=1}{\overset{\mathrm{}}{}}(1k^n)^{2d}$$
$$\times _k^1\frac{dz}{z}\left[\frac{\mathrm{exp}𝒢^{(1)}(1,z)}{\sqrt{z}}\right]^{2\alpha ^{}p_1p_2}.$$
(2.1)
We would like now to stress that if we want to reproduce diagrams of scalar field theories we have to ensure that only tachyon states propagate in the loops of string amplitudes. In fact this condition is fulfilled if small values of the multiplier $`k`$ are considered: indeed this parameter plays exactly the same role as the multiplier $`A`$ in the tree level amplitudes. Therefore an expansion in powers of $`k`$ is performed keeping the most divergent terms. In so doing we get
$$A_2^{(1)}(p_1,p_2)=\frac{N}{2}\frac{1}{(4\pi )^{d/2}}\frac{1}{(2\alpha ^{})^{d/2}}\left[2g_s(2\alpha ^{})^{\frac{d2}{4}}\right]^2$$
$$\times _0^1\frac{dk}{k^2}[\frac{1}{2}\mathrm{log}k]^{\frac{d}{2}}_k^1\frac{dz}{z}e^{2\alpha ^{}𝒢^{(1)}(1,z)}$$
(2.2)
where the Green function, in the limit we are considering, is
$$𝒢^{(1)}(z_1,z_2)=\mathrm{log}(z_1z_2)\frac{1}{2}\mathrm{log}z_1z_2+\frac{\mathrm{log}^2z_1/z_2}{2\mathrm{log}k}$$
(2.3)
Our aim is to identify the right limit to get the tadpole diagram in Fig. (2).
Starting from two three-vertices, we sew the leg $`0`$ with the leg $`+\mathrm{}`$ according to the Fig. (3).
Such a sewing is performed by considering again the projective transformation $`S(z)=Az`$, which has $`+\mathrm{}`$ and $`0`$ as fixed points and which transforms $`z_2=1`$ in the second vertex in Fig. (3a) in the multiplier $`A`$ getting the configuration shown in Fig. (3b).
The next step consists in performing a limit in which $`z_2z_1`$, i.e. in which $`A1`$ with $`\alpha ^{}\text{log}(1A)0`$, as said before. In this limit we should get the tadpole diagram. Indeed we have:
$$A_2^{(1)}(p_1,p_2)=\frac{2N}{(4\pi )^{d/2}}\frac{1}{2\alpha ^{}}g_s^2_0^1\frac{dk}{k^2}\left[\frac{1}{2}\mathrm{log}k\right]^{d/2}$$
(2.4)
By defining:
$$x=\alpha ^{}\mathrm{log}k$$
with $`0x+\mathrm{}`$, we can rewrite (2.4) as follows:
$$A_2^{(1)}(p_1,p_2)=\frac{N}{(4\pi )^{d/2}}\lambda _{\varphi ^4}_0^{\mathrm{}}𝑑xe^{xm^2}x^{d/2}$$
(2.5)
By using the matching condition established at the tree level (1.11) we get from string theory the tadpole diagram of $`\mathrm{\Phi }^4`$\- theory.
### 2.2 Candy diagram
Let us now derive the candy diagram from the four-tachyon one-loop amplitude:
$$A_4^{(1)}(p_1,p_2,p_3,p_4)=\frac{N}{(4\pi )^{d/2}}\text{Tr}\left[\lambda ^{a_1}\lambda ^{a_2}\lambda ^{a_3}\lambda ^{a_4}\right]$$
$$\times \frac{1}{(2\alpha ^{^{}})^{d/2}}\left[2g_s(2\alpha ^{^{}})^{\frac{d2}{4}}\right]^4_0^1\frac{dk}{k^2}\left[\frac{1}{2}\text{log}k\right]^{\frac{d}{2}}$$
$$\times _k^1\frac{dz_4}{z_4}_{z_4}^1\frac{dz_3}{z_3}_{z_3}^1\frac{dz_2}{z_2}$$
$$\underset{i<j=1}{\overset{4}{}}\left[\frac{\mathrm{exp}(𝒢(z_i,z_j))}{\sqrt{z_iz_j}}\right]^{2\alpha ^{}p_ip_j}$$
(2.6)
The diagram relative to this amplitude can be obtained by means of the sewing procedure illustrated in Fig. (5).
The four-particle vertices of the candy diagram can be generated by the corner of the moduli space where the Koba-Nielsen variables $`z_1z_2`$ and $`z_3z_4`$. This is performed by considering the limit in which the multipliers $`B_i`$ $`(i=1,2)1`$. We stress here that, in this limit, the Green functions $`𝒢(z_1,z_2)`$ and $`𝒢(z_3,z_4)`$ result to be divergent and we regularize them by introducing a cut-off $`ϵ`$ on the world-sheet so that $`B_i=1ϵ`$. In this limit the length of the strips connecting the three-vertices become very short and in this way the four-particle vertices of the diagram in consideration are generated. Furthermore, in order to select in the loop only the lightest states, we also take the limit in which the multiplier $`A0`$, and, after having performed both the limits, we send the cut-off to zero in all the regular expressions.
From these geometrical considerations that shed light on the different roles played by the multipliers $`A`$-like and $`B`$-like, we select the following corner of the moduli space reproducing the candy diagram of $`\mathrm{\Phi }^4`$-theory:
$$A0B_i=1ϵ1$$
(2.7)
Let us now evaluate the amplitude (2.6) in the corner (2.7).
The first step consists in rewriting, in this region of the moduli space, the measure and the integration region in the amplitude (2.6).
The ordering of the Koba-Nielsen variables determines the integration regions of the multipliers $`A`$ and $`B_i`$ in terms of which the whole amplitude is expressed, after the sewing. More precisely, in the limits (2.7), one gets:
$`{\displaystyle _0^1}{\displaystyle \frac{dk}{k^2}}{\displaystyle _k^1}{\displaystyle \frac{dz_2}{z_2}}{\displaystyle _k^{z_2}}{\displaystyle \frac{dz_3}{z_3}}{\displaystyle _k^{z_3}}{\displaystyle \frac{dz_4}{z_4}}`$
$`{\displaystyle _0^1}{\displaystyle \frac{dk}{k^2}}{\displaystyle _k^1}{\displaystyle \frac{dA}{A}}+O\left(k\right)`$ (2.8)
For this diagram, the proper times associated to the single propagators in the loop, are identified with the Schwinger parameters
$$t_1=\alpha ^{}\mathrm{log}k/At_2=\alpha ^{}\mathrm{log}A$$
(2.9)
where $`k`$ has to be understood as the proper time of the whole loop.
The Green functions defined in (2.3), in this limit, drastically simplify.
In particular the Green function $`2\alpha ^{}𝒢(z_1,z_3)`$, when written in terms of the Schwinger parameters, becomes
$$2\alpha ^{}𝒢(z_1,z_3)=t_2\frac{t_{2}^{}{}_{}{}^{2}}{t_1+t_2}.$$
(2.10)
By expressing the full amplitude in terms of $`t_1`$ and $`t_2`$ one gets:
$`A_4={\displaystyle \frac{N}{(4\pi )^{d/2}}}{\displaystyle \frac{1}{2}}d^{a_1a_2l}d^{a_3a_4l}\left[2^6g_s^4(2\alpha ^{})^{d4}\right]`$
$`\times {\displaystyle _0^{\mathrm{}}}dt_1{\displaystyle _0^{\mathrm{}}}dt_2(t_1+t_2)^{d/2}e^{m^2(t_1+t_2)}`$
$`\times e^{(p_1+p_2)^2\left[t_2\frac{t_{2}^{}{}_{}{}^{2}}{t_1+t_2}\right]}`$ (2.11)
Once again we have the right result in field theory by using the matching condition (1.11).
## 3 A two-loop diagram: double-candy
In this section we show how to get the double-candy diagram of $`\mathrm{\Phi }^4`$ theory, Fig. (6), starting from the two-loop four-tachyon amplitude in bosonic string theory:
$$A_4^{(2)}(p_1p_2p_3p_4)=N^2Tr\left[\lambda ^{a_1}\lambda ^{a_2}\lambda ^{a_3}\lambda ^{a_4}\right]C_2N_0^4$$
$$\times \left[dm\right]_2^4\underset{i<j}{}\left[\frac{\mathrm{exp}𝒢^{(2)}(z_i,z_j)}{\sqrt{V_i^{^{}}(0)V_j^{^{}}(0)}}\right]^{2\alpha ^{^{}}p_ip_j}$$
(3.12)
where the expressions for $`V_i^{^{}}(0)`$ are given by:
$$(V_i^{^{}}(0))^1=\left|\frac{1}{z_i\rho _a}\frac{1}{z_i\rho _b}\right|$$
(3.13)
with $`\rho _a`$ and $`\rho _b`$ depending on the position of $`z_i`$ and being the two fixed points on the left and on the right hand of $`z_i`$.
We expand the previous amplitude for small values of the multipliers $`k_\mu `$ keeping the most divergent contribution that is the one corresponding to the tachyon state and again the Green functions reduce to the following form :
$$𝒢^{(2)}(z_i,z_j)=\mathrm{log}(z_iz_j)$$
$$+\frac{\mathrm{log}^2T\mathrm{log}k_2+\mathrm{log}^2U\mathrm{log}k_12\mathrm{log}T\mathrm{log}U\mathrm{log}S}{2(\mathrm{log}k_1\mathrm{log}k_2\mathrm{log}^2S)}$$
with
$`S`$ $`=`$ $`{\displaystyle \frac{(\eta _1\eta _2)(\xi _1\xi _2)}{(\xi _1\eta _2)(\eta _1\xi _2)}}`$
$`T`$ $`=`$ $`{\displaystyle \frac{(z_j\eta _1)(z_i\xi _1)}{(z_j\xi _1)(z_i\eta _1)}}`$
$`U`$ $`=`$ $`{\displaystyle \frac{(z_j\eta _2)(z_i\xi _2)}{(z_i\eta _2)(z_j\xi _2)}}`$ (3.14)
The measure, once used the projective invariance to fix $`z_4=1`$, $`\xi _2=+\mathrm{}`$ and $`\eta _2=0`$, becomes:
$$\left[dm\right]_2^4=\underset{i=1}{\overset{3}{}}\frac{dz_i}{V_i^{^{}}(0)}\underset{\mu =1}{\overset{2}{}}\frac{dk_\mu }{k_\mu ^2}\frac{d\xi _1d\eta _1}{(\xi _1\eta _1)^2}$$
$$\times \left[det\left(i\tau _{\mu \nu }\right)\right]^{d/2}$$
(3.15)
where $`\tau _{\mu \nu }`$ is the period matrix in the limit of small multipliers .
Let us now identify the corner of the moduli space that, in the field theory limit, reproduces the two-loop candy diagram, according to our procedure. In Fig. (8) it is shown the final configuration that we reach applying the sewing procedure with the following projective transformations:
$$\begin{array}{ccc}S_i=B_iz& \widehat{S}_1=A_1z& \widehat{S}_2=A_2z\end{array}$$
(3.16)
with $`i=1,2,3`$.
Once the sewing procedure is completed, the Koba-Nielsen variables and the moduli of the surface, are expressed in terms of the multiplier of the transformations according the correspondence shown in Fig. (8).
If we want to obtain the four-particle vertices peculiar of the two-loop candy diagram, we have to take in consideration the corner of the Koba-Nielsen variables characterized by $`z_1z_2`$, $`z_3z_4`$ and by the modulo $`\xi _11`$. This configuration is achieved considering the limits in which $`B_i1`$ and introducing the suitable regularizators, when necessary.
Furthermore, considering also the limit $`A_i0`$, we select scalar particles in the other internal legs.
The corner of moduli space reproducing the $`\mathrm{\Phi }^4`$ scalar diagram illustrated in Fig.(6) is
$$\begin{array}{cc}A_i0\hfill & \hfill B_i=1ϵ.\end{array}$$
(3.17)
Let us now evaluate the amplitude (3.12) in this corner.
The Green functions (3) are then evaluated in the limit (3.17) where they take a simple form and the same is done for the local coordinates $`V_i^{^{}}(0)`$ and for the measure.
As regards the integration region, we observe that the sewing procedure determines an ordering of the Koba-Nielsen variables and of the fixed points as shown in Fig. (8).
In the field theory limit (3.17) we integrate the multipliers $`B`$-like between $`0`$ and $`1`$ and the multipliers $`A`$-like, between $`0`$ and $`\delta `$ being $`\delta `$, a positive infinitesimal quantity.
The Schwinger parameters in this case are related to the $`A_i`$’s by the following relations :
$$t_{i+2}=\alpha ^\mathrm{`}\mathrm{log}A_i,t_1=\alpha ^{^{}}\mathrm{log}\frac{k_1}{A_2},t_2=\alpha ^{^{}}\mathrm{log}\frac{k_2}{A_1}$$
(3.18)
with $`i=1,2`$.
Rewriting the Green functions and the measure in terms of the Schwinger parameters we get:
$`A_4^{(2)}(p_1\mathrm{}p_4)={\displaystyle \frac{N^2}{(4\pi )^2}}d^{a_1a_2l}d^{a_3a_4l}`$
$`\times {\displaystyle \frac{\left[2^4g_s^2(2\alpha ^{^{}})^{\frac{d4}{2}}\right]^3}{2^5}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \underset{i=1}{\overset{4}{}}}dt_ie^{m^2(t_1+t_2+t_3+t_4)}`$
$`\times (t_1+t_4)^{d/2}(t_2+t_3)^{d/2}`$
$`\times e^{(p_1+p_2)^2\left[\frac{t_1t_4}{t_1+t_4}+\frac{t_2t_3}{t_2+t_3}\right]}`$ (3.19)
where a sum over inequivalent permutations of the external particles has been done analogously as in the one-loop candy-diagram case.
Now using the matching condition (1.11), we get the same result, including the overall factor, as the one obtained in field theory.
In conclusion, we have used the sewing and cutting procedure in order to show how $`\mathrm{\Phi }^4`$-theory diagrams can be reproduced from string amplitudes, up to two loop-order. The whole information so obtained can be in principle extendible to Yang-Mills diagrams involving quartic interactions.
###### Acknowledgments.
We would like to thank P. Di Vecchia, A. Frizzo, L. Magnea, R. Russo and M.G. Schmidt for helpful discussions. R.M. and F.P. thank respectively Università di Napoli and NORDITA for their kind hospitality during different stages of their work.
|
warning/0003/math0003182.html
|
ar5iv
|
text
|
# Proof of the Rokhlin’s Conjecture on Arnold’s surfaces
## Part I Introduction
A real algebraic curve $`𝒜`$ of even degree is by definition a homogeneous real polynomial $`F(x_0,x_1,x_2)`$ of degree $`m=2k`$ in three variables. The set
$$𝐂𝒜=\{(x_0:x_1:x_2)𝐂P^2F(x_0,x_1,x_2)=0\}$$
is called complex point set of the curve. The complex conjugation $`conj`$ of $`𝐂P^2`$ induces an antiholomorphic involution on $`𝐂𝒜`$ with fixed point set $`𝐑𝒜=𝐂𝒜𝐑P^2`$, which is the set of real points of the curve. Given a real algebraic curve $`𝒜`$ with non-singular complex set point $`𝐂𝒜`$, using the fact that $`𝐂P^2/conj`$ is diffeomorphic to $`S^4`$ , one can associate to the triple $`(𝐂P^2,conj,𝐂𝒜)`$ two pairs $`(S^4,𝔄)`$ where $`𝔄`$ is a closed connected surface embedded in $`S^4`$ called Arnold Surface. These surfaces were introduced in . The present paper is devoted to studying the topology of these pairs. The Arnold Surfaces are constructed as follows. The real part $`𝐑𝒜`$ is a union of disjoint circles embedded in $`𝐑P^2`$ which divide $`𝐑P^2`$ into two parts:
$`𝐑P_+^2=\{x𝐑P^2F(x_0,x_1,x_2)0\}`$ and $`𝐑P_{}^2=\{x𝐑P^2F(x_0,x_1,x_2)0\}`$, which have common boundary $`𝐑𝒜`$. Gluing $`𝐑P_+^2`$, respectively $`𝐑P_{}^2`$, to the quotient $`𝐂𝒜/conj`$ along their common boundary $`𝐑𝒜`$, we get two closed surfaces in $`S^4`$:
$$𝔄_\pm =(𝐂𝒜/conj)(𝐑P^2)_\pm $$
One of these surfaces is obviously non orientable. Changing the sign of $`F`$ if necessary, we can always assume that the non-orientable part of $`𝐑P^2`$ lies entirely in $`𝐑P_{}^2`$. Then $`𝔄_{}`$ is non-orientable. Note that these surfaces are not smooth along $`𝐑𝒜`$ but can be easily smoothed since $`𝐂𝒜`$ is nowhere tangent to $`𝐑P^2`$. In this paper we consider only the smoothed Arnold surfaces and we keep for them the same name ”Arnold surfaces”, since no confusion is possible. We shall say that an embedded surface $`SS^4`$ is standard if it is a connected sum of standard tori and standard $`𝐑P^2`$. Moreover, when it does not lead to confusion, we shall call curve a real algebraic curve with non-singular complex point set. Rokhlin’s conjecture: Let $`𝒜`$ be a real algebraic curve with non-empty real part $`𝐑𝒜`$. <sup>1</sup><sup>1</sup>1In the case of empty real part, $`𝔄_{}`$ is trivially non-standard. Then both Arnold surfaces are standard. The proof of the Rokhlin’s Conjecture for the surface $`𝔄_+`$ associated to the maximal nest curves can be extracted from the paper of S.Akbulut. (Recall that the maximal nest curve is a curve whose real part is , up to isotopy in $`𝐑P^2`$, $`k`$ circles linearly ordered by inclusion.) In the paper , S.Finashin proved Rokhlin’s Conjecture for $`L`$ curves of even degree. (Recall that an $`L`$ curve is a perturbation of a curve which splits into a union of $`2k`$ real lines in general position.)
Here is a brief description of the methods of our proof.
Let $`k`$ be a positive integer. Consider an algebraic curve of degree $`2k`$ whose real part is described up to isotopy in $`𝐑P^2`$ as $`(k1)(2k1)+1`$ disjoint circles embedded in $`𝐑P^2`$ such that: $`\frac{(k1)(k2)}{2}`$ circles lie inside one circle and $`\frac{k(k3)}{2}`$ circles lie outside this circle where the inside (resp, outside) of a circle design respectively the part of $`𝐑P^2`$ homeomorphic to a disc (resp, a Möbius band).
Such a curve was initially obtained in 1876 from Harnack’s construction based on perturbations of singularities . It can be as well obtained by the Patchworking procedure due to O.Viro. We recall the Patchworking procedure in the Preliminary Section. The construction of Harnack curve using the Patchworking method was given by I.Itenberg in . We recall these construction of the so-called Harnack curve in the Chapter 1.
Our proof starts in the Chapter 2. In the Chapter 3, we prove that Arnold surfaces for Harnack curves are standard. The proof makes use of a desingularization procedure.
Afterwards (sections 5.2 of Chapter 5.2, 5.3 of Chapter 5.3) we define a procedure of modification of the Harnack curve which allows to obtain (up to conj-equivariant isotopy) all real algebraic curves with non-empty real part. In the chapter 6 we prove that a curve $`𝒜`$ obtained as the result of such modification has standard Arnold surfaces if $`𝐑𝒜`$ is non-empty.
Acknowledgement
I am very grateful to Ilia Itenberg for posing the problem and for useful conversations during the course of this work.
### Preliminaries-Patchworking Process
##### Introduction
The T-curves are constructed by a combinatorial procedure due to O.Viro . They are naturally introduced via the procedure of patchworking polynomials.
Here we introduce some notations and definitions. In what follows, $`𝐊`$ denotes either the real field $`𝐑`$ or the complex field $`𝐂`$; $`𝐙`$ denotes the ring of integers, $`𝐍`$ the set of positive integers, $`𝐑_+`$ the set of positive reals.
A linear combination of products $`\{x^iy^j(i,j)𝐙^2\}`$ with coefficients from $`𝐊`$ is called an $`L`$-polynomial over $`𝐊`$ in two variables. L-polynomials over $`𝐊`$ in two variables form a ring isomorphic to the ring of regular function of the variety $`(𝐊^{})^2`$. The variety $`(𝐊^{})^2`$ is known as an algebraic torus over $`𝐊`$.
Let us recall some properties of $`(𝐊^{})^2`$.
Let $`l`$ be the map $`(𝐊^{})^2𝐑^2`$ defined by: $`l(x,y)=(ln(|x|),ln(|y|)).`$ Let $`U_𝐊=\{x𝐊|x|=1\}`$ and $`a`$ be the map $`(𝐊^{})^2U_𝐊^2`$ defined by : $`a(x,y)=(\frac{x}{|x|},\frac{y}{|y|})`$. The map $`la=(l,a):(𝐊^{})^2𝐑^2\times U_𝐊^2`$ is a diffeomorphism.
Being an abelian group $`(𝐊^{})^2`$ acts on itself by translations.
Let us recall some translations involved into this action. For $`(i,j)𝐑^2`$ and $`t>0`$, denote by $`qh_{(i,j),t}:(𝐊^{})^2(𝐊^{})^2`$ the translation defined by formula : $`qh_{(i,j),t}(x,y)=(t^ix,t^jy)`$. For $`(x^{},y^{})U_𝐊^2`$ denote by $`S_{(x^{},y^{})}`$ the translation $`(𝐊^{})^2(𝐊^{})^2`$ defined by formula : $`S_{(x^{},y^{})}(x,y)=(x^{}.x,y^{}.y)`$.
Any usual real polynomial $`f=f(x,y)=_{(i,j)𝐍^2}a_{i,j}x^iy^j`$ can be considered as a Laurent polynomial with coefficients from $`𝐊`$. Let $`f=f(x,y)=_{(i,j)𝐍^2}a_{i,j}x^iy^j`$ be a real polynomial in two variables. Denote by $`\mathrm{\Delta }(f)`$ and call Newton polyhedron of $`f`$ the convex hull of $`\{w𝐍^2a_w0\}`$ . By polyhedron we mean a convex bounded polyhedron in $`𝐑^2`$ with integer vertices having positive coordinates, in other words a Newton polyhedron of a real polynomial in two variables of finite degree. The interior of a polyhedron $`\mathrm{\Delta }`$ is denoted $`\mathrm{\Delta }^0`$.
A polyhedron is non-degenerate if its dimension is two, in other words its interior is non-empty.
A subdivision of a non-degenerate polyhedron $`\mathrm{\Delta }`$ is a set of non-degenerate polyhedra $`\mathrm{\Delta }_1,\mathrm{}\mathrm{\Delta }_r`$ such that $`\mathrm{\Delta }=_i\mathrm{\Delta }_i`$, and any intersection $`\mathrm{\Delta }_i\mathrm{\Delta }_j`$ is either empty or is a common face of both $`\mathrm{\Delta }_i,\mathrm{\Delta }_j`$. We call proper face of a polyhedron $`\mathrm{\Delta }`$ a face of codimension 1 and denote by $`𝒢^{}(\mathrm{\Delta })`$ the set of proper faces of a polyhedron $`\mathrm{\Delta }`$. The set of faces of a polyhedron $`\mathrm{\Delta }`$ is denoted by $`𝒢(\mathrm{\Delta })`$. We denote $`C(\mathrm{\Delta })`$ the vector subspace of $`𝐑^2`$ which corresponds to the minimal affine subspace containing $`\mathrm{\Delta }`$.
##### Toric varieties
To each polyhedron one can associate a variety called
$`𝐊`$-toric variety. The following definition is deduced from ( p.166-168).
###### Definition 0.0.1.
Let $`\mathrm{\Delta }`$ be a polyhedron.
Consider the finite set $`A(\mathrm{\Delta })=𝐙^2\mathrm{\Delta }`$ and choose an ordering of this set so that
$`A(\mathrm{\Delta })=\{(i_1,j_1),\mathrm{},(i_n,j_n)\}`$. Consider the subset
$$𝐊\mathrm{\Delta }^0=\{(x^{i_1}y^{j_1}:\mathrm{}:x^{i_n}y^{j_n})(x,y)(𝐊^{})^2\}$$
of $`𝐊P^{n1}`$.
The closure of $`𝐊\mathrm{\Delta }^0`$ is called the $`𝐊`$-projective toric variety associated to $`\mathrm{\Delta }`$ and the ordering of $`A`$.
###### Remark 0.0.2.
1. Another choice of the ordering of $`A`$ leads to an isomorphic $`𝐊`$-projective toric variety, so we shall use the term $`𝐊`$-projective toric variety associated to $`\mathrm{\Delta }`$ whenever no confusion is possible.
2. Moreover, let $`B𝐙^2`$ and $`T:𝐙^2𝐙^m`$ be an affine integer injective transformation such that $`T(A)=B`$, the $`𝐊`$ projective variety associated to $`\mathrm{\Delta }`$ and the set $`B`$ is isomorphic to the $`𝐊`$-projective toric variety associated to $`\mathrm{\Delta }`$ and an ordering set $`A(\mathrm{\Delta })`$.
Hence, denote by $`𝐊\mathrm{\Delta }`$ the $`𝐊`$-projective toric variety associated to a polyhedron $`\mathrm{\Delta }`$.
###### Definition 0.0.3.
Consider the action of the torus $`(𝐊^{})^2`$ given by formula $`(x,y).(z_1:\mathrm{}:z_n)=(x^{i_1}y^{j_1}z_1:\mathrm{}:x^{i_n}y^{j_n}z_n)`$, the variety $`𝐊\mathrm{\Delta }`$ is the closure of the orbit of the point $`(1:\mathrm{}:1)`$ under this action.
Let us recall some properties of toric varieties. We refer to , for more details about toric varieties. The following Lemma is implicit in , Prop 1.9,p.171.
###### Lemma 0.0.4.
Let $`\mathrm{\Delta }`$ be a polyhedron. Let $`\mathrm{\Gamma }`$ be a face of $`\mathrm{\Delta }`$. Let $`A(\mathrm{\Delta })=𝐙^2\mathrm{\Delta }`$ and choose an ordering of this set so that $`A(\mathrm{\Delta })=\{(i_1,j_1),\mathrm{},(i_n,j_n)\}`$. Consider the subset $`A(\mathrm{\Gamma })=𝐙^2\mathrm{\Gamma }`$.
Let $`:A(\mathrm{\Delta })0,1`$ be the map such that $`(A(\mathrm{\Gamma }))=1`$ $`(A(\mathrm{\Delta })\backslash A(\mathrm{\Gamma }))=0`$. Consider the subset
$$𝐊^{}\mathrm{\Gamma }^0=\{(x^{i_1}y^{j_1}((i_1,j_1)):\mathrm{}:(x^{i_n}y^{j_n}((i_n,j_n))(x,y)(𝐊^{})^2\}$$
of $`𝐊P^{n1}`$. The closure of $`𝐊^{}\mathrm{\Gamma }^0`$ is a subvariety of $`𝐊\mathrm{\Delta }`$ which corresponds to the closure of the orbit of the point $`((i_1,j_1):\mathrm{}:(i_n,j_n))`$ under the action of the torus $`(𝐊^{})^2`$ given by formula $`(x,y).(z_1:\mathrm{}:z_n)=(x^{i_1}y^{j_1}z_1:\mathrm{}:x^{i_n}y^{j_n}z_n)`$. Denote it by $`𝐊^{}\mathrm{\Gamma }`$. The variety $`𝐊^{}\mathrm{\Gamma }`$ is isomorphic to $`𝐊\mathrm{\Gamma }`$ the $`𝐊`$-projective toric variety associated to $`\mathrm{\Gamma }`$. In the cases when it does not lead to confusion we shall identify $`𝐊^{}\mathrm{\Gamma }`$ with $`𝐊\mathrm{\Gamma }`$ and $`𝐊^{}\mathrm{\Gamma }^0`$ with $`𝐊\mathrm{\Gamma }^0`$.
###### Proposition 0.0.5.
Let $`\mathrm{\Delta }`$ a polyhedron. Let $`𝐊\mathrm{\Delta }`$ be the $`𝐊`$-projective toric variety associated to $`\mathrm{\Delta }`$. Translations $`qh_{(i,j),t}`$ and $`S_{(x^{},y^{})}`$ of $`(𝐊^{})^2`$ define naturally translations $`qh_{(i,j),t}`$ and $`S_{(x^{},y^{})}`$ on $`𝐊\mathrm{\Delta }^0`$ which can be extended to transformations on $`𝐊\mathrm{\Delta }`$. In what follows, we keep the same notations for translations of $`𝐊\mathrm{\Delta }^0`$ and their extensions to transformations of $`𝐊\mathrm{\Delta }`$.
proof: On the assumptions of definition 0.0.1 of Chapter 0.0.1, let $`𝐊\mathrm{\Delta }`$ be the $`𝐊`$-projective toric variety associated to $`\mathrm{\Delta }`$. Since $`𝐊^{}\mathrm{\Gamma }^0`$ and their closure $`𝐊^{}\mathrm{\Gamma }`$,with $`\mathrm{\Gamma }`$ a face of $`\mathrm{\Delta }`$, are contained in the affine parts of $`𝐊P^{n1}`$, extensions of transformations $`qh_{(i,j),t}`$ and $`S_{(x^{},y^{})}`$ on varieties $`𝐊^{}\mathrm{\Gamma }`$ give transformations on the affine parts of $`𝐊P^{n1}`$ and thus define transformations of $`𝐊\mathrm{\Delta }`$. Q.E.D
###### Definition 0.0.6.
Let $`\mathrm{\Delta }`$ be a polyhedron. On the same assumptions as in Definition 0.0.1 of Chapter 0.0.1, let $`𝐊\mathrm{\Delta }`$ be the closure of $`𝐊\mathrm{\Delta }^0`$
$$𝐊\mathrm{\Delta }^0=\{(x^{i_1}y^{j_1}:\mathrm{}:x^{i_n}y^{j_n})(x,y)(𝐊^{})^2\}$$
subset of $`𝐊P^{n1}`$. Let $`𝐑_+\mathrm{\Delta }`$ be the closure (in the classic topology) of
$$𝐑_+\mathrm{\Delta }^0=\{(x^{i_1}y^{j_1}:\mathrm{}:x^{i_n}y^{j_n})(x,y)(𝐑_+^{})^2\}$$
subset of $`𝐊\mathrm{\Delta }`$. The subvariety $`𝐑_+\mathrm{\Delta }`$ of $`𝐊\mathrm{\Delta }`$ is homeomorphic,as stratified space,to the polyhedron $`\mathrm{\Delta }`$ stratified by its faces. An explicit homeomorphism is given by the moment map $`\mu :𝐊\mathrm{\Delta }\mathrm{\Delta },\mu (x,y)=\frac{_{(i,j)A(\mathrm{\Delta })}|x^iy^j|.(i,j)}{_{(i,j)A(\mathrm{\Delta })}|x^iy^j|}`$
proof: It follows easily from the properties of the moment map $`\mu `$(; see also , p.198 Theorem1.11) $`\mu :𝐊\mathrm{\Delta }\mathrm{\Delta }`$.Q.E.D
Consider the action of the torus $`U_𝐊^2`$ on $`𝐊\mathrm{\Delta }`$ given by formula $`(x,y).(z_1:\mathrm{}:z_n)=(x^{i_1}y^{j_1}z_1:\mathrm{}:x^{i_n}y^{j_n}z_n)`$. For a face $`\mathrm{\Gamma }`$ of $`\mathrm{\Delta }`$, denote by $`U_\mathrm{\Gamma }`$ the subgroup of $`U_𝐊^2`$ consisting of elements $`(e^{i\pi k},e^{i\pi l})`$ with $`(k,l)`$ a vector orthogonal to $`C(\mathrm{\Gamma })`$.
###### Lemma 0.0.7.
Let $`\mathrm{\Delta }`$ be a polyhedron. The map $`\rho :𝐑_+\mathrm{\Delta }\times U_𝐊^2𝐊\mathrm{\Delta }`$ defined by formula $`\rho ((x,y),(x^{},y^{}))=S_{(x^{},y^{})}(x,y)`$ is a proper surjection. The variety $`𝐊\mathrm{\Delta }`$ is homeomorphic to the quotient space of $`𝐑_+\mathrm{\Delta }\times U_𝐊^2`$ with respect to the partition into sets $`x\times yU_\mathrm{\Gamma }`$, $`x𝐑_+\mathrm{\Delta }𝐊\mathrm{\Gamma }^0,yU_𝐊^2`$.
proof:
Let $`\mathrm{\Delta }`$ be a polyhedron. Let $`𝐑_+\mathrm{\Delta }`$ be the closure of
$$𝐑_+\mathrm{\Delta }^0=\{(x^{i_1}y^{j_1}:\mathrm{}:x^{i_n}y^{j_n})(x,y)(𝐑_+^{})^2\}$$
subset of $`𝐊\mathrm{\Delta }`$. Translations $`S_{(x^{},y^{})},(x^{},y^{})U_𝐊^2`$ define an action of $`U_𝐊^2`$ in $`𝐊\mathrm{\Delta }`$ such that intersection of $`𝐑_+\mathrm{\Delta }`$ with each orbit under this action consists of one point. Furthermore, for $`x`$ in the interior of $`\mathrm{\Gamma }`$ where $`\mathrm{\Gamma }`$ is a face of $`\mathrm{\Delta }`$, the stationary subgroup of action of $`U_𝐊^2`$ consists of transformations $`S_{(e^{i\pi k},e^{i\pi l})}`$ where $`(k,l)`$ is a vector orthogonal to $`C(\mathrm{\Gamma })`$. Since $`𝐊\mathrm{\Delta }`$ is locally compact and Hausdorff,it follows the Lemma. Q.E.D
###### Lemma 0.0.8.
The $`𝐊`$-projective toric variety associated to a non-degenerate triangle $`\mathrm{\Delta }`$ is $`𝐊P^2`$. The $`𝐊`$-projective toric variety associated to a degenerate triangle $`\mathrm{\Delta }`$ is $`𝐊P^1`$.
proof:
It can be easily deduced from the Lemma 0.0.7 of Chapter 0.0.7 above. Let $`𝐊=𝐑`$. Place $`\mathrm{\Delta }\times U_R^2`$ in $`𝐑^2`$ identifying $`(x,y)(x^{},y^{})\mathrm{\Delta }\times U_R^2`$ with $`S_{(x^{},y^{})}(x,y)`$. The surface $`𝐑\mathrm{\Delta }`$ can be obtained by an appropriate gluing of four copies of $`\mathrm{\Delta }`$.
Let $`𝐊=𝐂`$. Place $`\mathrm{\Delta }\times U_C^2`$ in $`𝐑^2`$ identifying $`(x,y)(x^{},y^{})\mathrm{\Delta }\times U_C^2`$ with $`S_{(x^{},y^{})}(x,y)`$. Let $`\mathrm{\Delta }`$ be a non-degenerate triangle, it is not difficult to get the usual handlebody decomposition of $`𝐂P^2`$ as the union of three 4-balls. In the same way, let $`\mathrm{\Delta }`$ be a degenerate triangle, it is not difficult to get $`𝐂P^1`$ as the union of two 2-discs.Q.E.D
##### Hypersurfaces of toric varieties
Let $`f`$ be a real polynomial and $`\mathrm{\Delta }(f)`$ its Newton polyhedron. Let $`\mathrm{\Delta }`$ be a polyhedron such that $`C(\mathrm{\Delta }(f))C(\mathrm{\Delta })`$. The equation $`f=0`$ defines in $`𝐊\mathrm{\Delta }^0`$ an hypersurface $`V_{𝐊\mathrm{\Delta }^0}(f)`$ of $`𝐊\mathrm{\Delta }^0`$. Denote by $`V_{𝐊\mathrm{\Delta }}(f)`$ its closure (in the Zarisky topology) in $`𝐊\mathrm{\Delta }`$. (In the case $`𝐊=𝐂`$ the classic topology gives the same result, but in the case $`𝐊=𝐑`$ the usual closure may be a non-algebraic set)
###### Lemma 0.0.9.
(, p.175) Let $`f`$ a real polynomial and $`\mathrm{\Delta }(f)`$ its Newton polyhedron. Let $`\mathrm{\Delta }`$ a non-degenerate polyhedron ( $`C(\mathrm{\Delta }(f))C(\mathrm{\Delta })`$). For any vector $`(k,l)C(\mathrm{\Delta })`$ orthogonal to $`C(\mathrm{\Delta }(f))`$, the hypersurface $`V_{𝐊_\mathrm{\Delta }}(f)`$ is invariant under transformations $`S_{(e^{i\pi k},e^{i\pi l})}`$ and $`qh_{(k,l),t}`$ of $`𝐊\mathrm{\Delta }`$ .
Let $`\mathrm{\Delta }`$ be a polyhedron. An homeomorphism $`\varphi :\mathrm{\Delta }\mathrm{\Delta }`$ is said admissible if it maps any face $`\mathrm{\Gamma }`$ of $`\mathrm{\Delta }`$ to itself. (Thus, an admissible homeomorphism $`\varphi :\mathrm{\Delta }\mathrm{\Delta }`$ extends to an homeomorphism $`\stackrel{~}{\varphi }:\mathrm{\Delta }\times U_𝐊^2\mathrm{\Delta }\times U_𝐊^2`$ by equivariance).
Hence, the following definition is naturally introduced.
###### Definition 0.0.10.
1. Let $`f`$ be a real polynomial in two variables, $`\mathrm{\Delta }(f)`$ its Newton polyhedron and $`𝐊\mathrm{\Delta }(f)`$ the toric variety associated to $`\mathrm{\Delta }(f)`$. Denote by $`h=\mu |_{𝐑_+\mathrm{\Delta }(f)}^1`$ the homeomorphism $`h:\mathrm{\Delta }(f)𝐑_+\mathrm{\Delta }(f)`$. Let $`\rho :𝐑_+\mathrm{\Delta }(f)\times U_𝐊^2𝐊\mathrm{\Delta }(f)`$ be the surjection defined by formula $`\rho ((x,y),(x^{},y^{}))=S_{(x^{},y^{})}(x,y)=(x.x^{},y.y^{})`$.
$$\mathrm{\Delta }(f)\times U_𝐊^2^{h\times id}𝐑_+\mathrm{\Delta }(f)\times U_𝐊^2^\rho 𝐊\mathrm{\Delta }(f)$$
A pair consisting of $`(\mathrm{\Delta }(f))\times U_𝐊^2`$ and its subset $`l`$ which is the pre-image of $`V_{𝐊\mathrm{\Delta }(f)}`$ under $`\rho (h\times id)`$ is a canonical $`𝐊`$-chart of $`f`$.
2. Let $`(\mathrm{\Delta }(f)\times U_𝐊^2,l)`$ be the canonical $`𝐊`$-chart of $`f`$, and
$`\varphi :\mathrm{\Delta }(f)\times U_𝐊^2\mathrm{\Delta }(f)\times U_𝐊^2`$ the extended homeomorphism of an admissible homeomorphism $`\varphi :\mathrm{\Delta }(f)\mathrm{\Delta }(f)`$, the pair $`(\varphi (\mathrm{\Delta }(f)\times U_𝐊^2,\varphi (l))`$ is a $`𝐊`$-chart of $`f`$.
3. Let $`(𝐑)_+^2`$ be set of positive integers $`\{(x,y)𝐑^2x0,y0\}`$ Consider the map: $`(𝐑^2)_+\times U_𝐊^2𝐂^2:((x,y)(x^{},y^{}))S_{(x^{},y^{})}(x,y)`$. Call reduced $`𝐊`$-chart of $`f`$ the image of a $`𝐊`$-chart of $`f`$ under this map.
Singular Hypersurfaces Let $`f=f(x,y)=_{(i,j)𝐍^2}a_{i,j}x^iy^j`$ be a real polynomial. The set
$$SV_{(𝐂^{})^2}(f)=\{(x,y)(𝐂^{})^2f(x,y)=\frac{f}{x}(x,y)=\frac{f}{y}(x,y)=0\}$$
defines the set of singular points of the variety defined in $`(𝐂^{})^2`$ by the equation $`f=0`$. For a set $`\mathrm{\Gamma }𝐍^2`$ consider the polynomial $`_{(i,j)\mathrm{\Gamma }}a_{i,j}x^iy^j`$. It is called $`\mathrm{\Gamma }`$-truncation of $`f`$ and is denoted by $`f^\mathrm{\Gamma }(x,y)`$.
The polynomial $`f`$ is completely non-degenerate if for any face $`\mathrm{\Gamma }`$ (including $`\mathrm{\Delta }(f)`$) of its Newton polyhedron $`\mathrm{\Delta }(f)`$ the variety defined in $`(𝐂^{})^2`$ by the equation $`f^\mathrm{\Gamma }=0`$ is non-singular.
###### Lemma 0.0.11.
Completely non-degenerate real polynomials form a Zarisky open subset of the space of real polynomials with a given Newton polyhedron.
#### 0.1. Patchworking polynomials and T-curves
In this section, we shall explain the procedure of patchworking polynomials and define the concept of $`T`$-curve.
##### Patchworking polynomials
A reduced $`𝐑`$-chart of a non-degenerate polynomial $`a(x,y)=a_{i_1,j_1}x^{i_1}y^{j_1}+a_{i_2,j_2}x^{i_2}y^{j_2}+a_{i_3,j_3}x^{i_3}y^{j_3}`$ may be constructed from the following method.
Let $`a=a(x,y)=a_{i_1,j_1}x^{i_1}y^{j_1}+a_{i_2,j_2}x^{i_2}y^{j_2}+a_{i_3,j_3}x^{i_3}y^{j_3}`$ be a real polynomial with non-degenerate Newton polyhedron $`\mathrm{\Delta }`$. Assign to each vertex $`a_{i,j}`$ of $`\mathrm{\Delta }`$ the sign
$$ϵ_{i,j}=\frac{a_{i,j}}{|a_{i,j}|}$$
Consider the union of $`\mathrm{\Delta }`$ and its symmetric copies $`\mathrm{\Delta }_x=s_x(\mathrm{\Delta })`$ $`\mathrm{\Delta }_y=s_y(\mathrm{\Delta })`$ $`\mathrm{\Delta }_{xy}=s(\mathrm{\Delta })`$ where $`s_x,s_y`$ are reflections with respect to the coordinate axes and $`s=s_xs_y`$. For each vector $`(w_1,w_2)`$ orthogonal to $`\mathrm{\Delta }`$ with integer relatively prime coordinates glue the points $`(x,y)`$ and $`((1)^{w_1}x,(1)^{w_2}y)`$ of the union $`\mathrm{\Delta }\mathrm{\Delta }_x\mathrm{\Delta }_y\mathrm{\Delta }_{xy}`$. Denote $`\mathrm{\Delta }_{}`$ the resulting space. Extend the distribution of signs to $`\mathrm{\Delta }_{}`$ so that $`g^{}(ϵ_{i,j}x^iy^j)=ϵ_{g(i,j)}x^iy^j`$ for $`g=s_x,s_y,s`$. Let $``$ be one of the four triangles $`\mathrm{\Delta },\mathrm{\Delta }_x,\mathrm{\Delta }_y,\mathrm{\Delta }_{xy}`$. If $``$ has vertices of different signs, consider the midline of $``$ separating them. Denote by $`L`$ the union of such midlines.
We shall say that the pair $`(\mathrm{\Delta }_{},L)`$ is obtained from $`a`$ by combinatorial patchworking.
###### Lemma 0.1.1.
”The smallest patch” Let $`a(x,y)=a_{i_1,j_1}x^{i_1}y^{j_1}+a_{i_2,j_2}x^{i_2}y^{j_2}+a_{i_3,j_3}x^{i_3}y^{j_3}`$ be a real polynomial with non-degenerate Newton polyhedron $`\mathrm{\Delta }`$. Let $`(\mathrm{\Delta }_{},L)`$, obtained from $`a(x,y)`$ by combinatorial patchworking. The pair $`(\mathrm{\Delta }_{},L)`$ is a reduced $`𝐑`$-chart of $`a`$.
It is easy to deduce the following characterization of the set of real points $`\{(x,y)𝐑^2,a(x,y)=0\}`$ and $`\{(x_0:x_1:x_2)𝐑P^2,A(x_0,x_1,x_2)=0\}`$ where $`A`$ is the homogenization of $`a`$.
1. Remove from $`\mathrm{\Delta }_{}`$ the sides of $`\mathrm{\Delta }\mathrm{\Delta }_x\mathrm{\Delta }_y\mathrm{\Delta }_{xy}`$ which are not glued in the construction of $`\mathrm{\Delta }_{}`$. It turns the polyhedron $`\mathrm{\Delta }_{}`$ to the polyhedron $`\mathrm{\Delta }^{}`$ homeomorphic to $`𝐑^2`$ and the set $`L`$ to a set $`L^{}`$ such that the pair $`(\mathrm{\Delta }^{},v^{})`$ is homeomorphic to $`(𝐑^2,\{(x,y)𝐑^2,a(x,y)=0\})`$. Glue by $`s`$ the opposite sides of $`\mathrm{\Delta }_{}`$. The resulting space $`\overline{\mathrm{\Delta }}_{}`$ is homeomorphic to the projective plane $`𝐑P^2`$. Denote $`\overline{L}`$ the image of $`L`$ in $`\overline{\mathrm{\Delta }}_{}`$. Let $`A(x_0,x_1,x_2)`$ the homogenization of $`a(x,y)`$. It defines a curve $`𝒜`$. Then there exists an homeomorphism $`(\overline{\mathrm{\Delta }}_{},\overline{L})(𝐑P^2,𝐑𝒜)`$.
One can generalize this description to more complicated polynomials . Let $`a_1,\mathrm{},a_r`$ be completely non-degenerate real polynomials in two variables with Newton polyhedra $`\mathrm{\Delta }_1,\mathrm{},\mathrm{\Delta }_r`$ in such a way that $`\mathrm{\Delta }_1,\mathrm{},\mathrm{\Delta }_r`$ form a subdivision of a non-degenerate polyhedron $`\mathrm{\Delta }`$ and $`a_i^{\mathrm{\Delta }_i\mathrm{\Delta }_j}=a_j^{\mathrm{\Delta }_i\mathrm{\Delta }_j}`$
###### Lemma 0.1.2.
Patchworking polynomials
Assume that the subdivision is regular, that is, there exists a convex non-negative function $`\nu :\mathrm{\Delta }𝐑`$, satisfying the following conditions:
1. all the restrictions $`\nu |\mathrm{\Delta }_i`$ are linear.
2. if the restriction of $`\nu `$ to an open set is linear, then this set is contained in one of $`\mathrm{\Delta }_i`$
3. $`\nu (\mathrm{\Delta }𝐍^2)𝐍`$.
Such a function is said convexifying the subdivision $`\mathrm{\Delta }_1,\mathrm{},\mathrm{\Delta }_r`$ of $`\mathrm{\Delta }`$.
There exists a unique polynomial $`a`$ with $`\mathrm{\Delta }(a)=\mathrm{\Delta },a^{\mathrm{\Delta }_i}=a_i`$ for $`i=1,\mathrm{},r`$. Let it be $`a(x,y)=_{(i,j)𝐍^2}a_{i,j}x^iy^j`$.
Then introduce the one-parameter family of polynomials
$$b_t=b_t(x,y)=\underset{(i,j)𝐍^2}{}a_{i,j}x^iy^jt^{\nu (i,j)}$$
We say that polynomials $`b_t`$ are obtained by Patchworking the polynomials $`a_1,\mathrm{},a_r`$ by $`\nu `$.
###### Definition 0.1.3.
Patchworking Charts A pair $`(\mathrm{\Delta }\times U_𝐊^2,L)`$ is said to be obtained by patchworking from $`𝐊`$-charts of polynomials $`a_1,\mathrm{},a_r`$ if $`\mathrm{\Delta }=_{i=1}^r\mathrm{\Delta }_i`$ (where $`\mathrm{\Delta }_i`$ denotes the Newton polyhedron of $`a_i`$) and one can choose $`𝐊`$-charts $`(\mathrm{\Delta }_i\times U_𝐊^2,l_i)`$ of polynomials $`a_i`$ such that $`L=_{i=1}^r(l_i)`$.
###### Theorem 0.1.4.
Patchwork Theorem<sup>2</sup><sup>2</sup>2This theorem is a simplified version of the Patchworking Theorem for $`L`$-polynomials Let $`a_1,\mathrm{},a_r`$ be completely non-degenerate polynomial in two variables with Newton polyhedra $`\mathrm{\Delta }_1`$,…,$`\mathrm{\Delta }_r`$. Assume furthermore that $`\mathrm{\Delta }_1`$,…,$`\mathrm{\Delta }_r`$ form a regular subdivision of a non-degenerate polyhedron $`\mathrm{\Delta }`$. Then there exists $`t_0>0`$ such that for any $`t]0,t_0]`$ the polynomial $`b_t`$ is completely non-degenerate and its $`𝐊`$-chart is obtained by patchworking $`𝐊`$-charts of the polynomials $`a_1,\mathrm{},a_r`$.
Further, we say that such polynomials $`b_t`$ are obtained by patchworking process.
We shall give the main ideas of the proof of the patchworking theorem in the section 0.1 of Chapter 0.1.
##### T-curves
Bringing together the Lemma 0.1.1 of Chapter 0.1.1 (”the smallest patch ”) and the patchworking theorem 0.1.4 of Chapter 0.1.4 we deduce a combinatorial method of construction of curves. These curves are called T-curves.
Let $`m`$ be a positive integer. Let $`\mathrm{\Delta }`$ be the triangle
$$\{(x,y)𝐑^2x0,y0,x+ym\}$$
(Up to linear change of coordinates $`(x_0:x_1:x_2)`$ of $`𝐂P^2`$, the convex hull of $`\mathrm{\Delta }`$ may be the Newton polyhedron of the affine polynomial $`f(x,y)=F(1,x,y)`$ associated to a homogeneous polynomial $`F(x_0,x_1,x_2)`$ of degree $`m`$.)
Let $`\tau `$ be a regular triangulation of $`\mathrm{\Delta }`$ whose vertices have integer coordinates. Suppose that some distribution of signs $`\chi `$ at the vertices of the triangulation is given. Denote the sign $`\pm `$ at the vertex with coordinates $`(i,j)`$ by $`ϵ_{(i,j)}`$. Take the square $`\mathrm{\Delta }_{}`$ made of $`\mathrm{\Delta }`$ and its symmetric copies $`\mathrm{\Delta }_x=s_x(\mathrm{\Delta }),\mathrm{\Delta }_y=s_y(\mathrm{\Delta }),\mathrm{\Delta }_{xy}=s_{xy}(\mathrm{\Delta })`$ where $`s_x,s_y,s=s_xs_y`$ are reflections with respect to the coordinate axes. The resulting space is homeomorphic to $`𝐑^2`$. Extend the triangulation $`\tau `$ of $`\mathrm{\Delta }`$ to a symmetric triangulation $`\tau _{}`$ of $`\mathrm{\Delta }_{}`$. Extend the distribution of signs to a distribution at the vertices of $`\mathrm{\Delta }_{}`$ which verifies the modular properties: $`g^{}(ϵ_{i,j})x^iy^j=ϵ_{g(i,j)}x^iy^j`$ for $`g=s_x,s_y,s`$.
If a triangle of $`\mathrm{\Delta }_{}`$ has vertices of different signs, consider a midline separating them. Denote by $`L`$ the union of such midlines. Glue by $`s`$ the opposite sides of $`\mathrm{\Delta }_{}`$. The resulting space $`\overline{\mathrm{\Delta }}_{}`$ is homeomorphic to the projective plane $`𝐑P^2`$. Denote $`\overline{L}`$ the image of $`L`$ in $`\overline{\mathrm{\Delta }}_{}`$. Call $`\mathrm{\Delta }_{}`$, $`L`$, $`\overline{\mathrm{\Delta }}_{}`$,$`\overline{L}`$ obtained from $`\mathrm{\Delta },\tau ,\chi `$ by combinatorial patchworking.
###### Theorem 0.1.5.
Patchwork Theorem from the real view point
and T-curves. Define the one-parameter family of polynomials
$$b_t=b_t(x,y)=\underset{(i,j)verticesof\tau }{}ϵ_{i,j}x^iy^jt^{\nu (i,j)}$$
where $`\nu `$ is a function convexifying the triangulation $`\tau `$ of $`\mathrm{\Delta }`$. Let, $`\overline{\mathrm{\Delta }}_{}`$,$`\overline{L}`$ obtained from $`\mathrm{\Delta },\tau ,\chi `$ by combinatorial patchworking.
Denote by $`B_t=B_t(x_0,x_1,x_2)`$ the corresponding homogeneous polynomials:
$$B_t(x_0,x_1,x_2)=x_0^mb_t(x_1/x_0,x_2/x_0)$$
Then there exists $`t_0>0`$ such that for any $`t]0,t_0]`$ the equation $`B_t(x_0,x_1,x_2)=0`$ defines in $`𝐑P^2`$ the set of real points of an algebraic curve $`C_t`$ such that the pair $`(𝐑P^2,𝐑C_t)`$ is homeomorphic to the pair $`(\overline{\mathrm{\Delta }}_{},\overline{L})`$.
Such a curve is called a $`T`$-curve .
proof : It is an immediate consequence of the Lemma 0.1.2 of Chapter 0.1.2 and the theorem 0.1.4 of Chapter 0.1.4. Let $`\mathrm{\Delta }_{}`$,$`L`$ obtained from $`\mathrm{\Delta },\tau ,\chi `$ by combinatorial patchworking. There exists $`t_0`$ such that for any $`t]0,t_0]`$ the pair $`(\mathrm{\Delta }_{},L)`$ is a reduced $`𝐑`$-chart of $`b_t`$. Q.E.D
###### Remark 0.1.6.
1. Remove the sides $`x>0,y>0,x+y=m`$ and its symmetric copies to $`\mathrm{\Delta }_{}`$. It turns the polyhedron $`\mathrm{\Delta }_{}`$ to the polyhedron $`\mathrm{\Delta }^{}`$ homeomorphic to $`𝐑^2`$ and the set $`L`$ to a set $`L^{}`$ such that for any $`t]0,t_0]`$ the pair $`(\mathrm{\Delta }^{},v^{})`$ is homeomorphic to $`(𝐑^2,\{(x,y)𝐑^2,b_t(x,y)=0\})`$.
2. It is possible to recover in some cases whether the set of real points of the $`T`$-curve $`B`$ divides the set of its complex points. Assume the triangulation $`\tau `$ of $`\mathrm{\Delta }`$ sufficiently fine such that each triangle $`\mathrm{\Delta }_i`$ in the triangulation $`\tau `$ is the Newton polyhedron of a polynomial $`a_i`$ which defines a curve with orientable real set of points. Denote $`L_i`$ the union of midlines in $`(\mathrm{\Delta }_i)_{}`$ homeomorphic to this set. If one can (resp, not) choose an orientation on each $`L_i`$ compatible with an orientation of the union of $`_i(L_i)`$ then the set of real points of the $`T`$-curve $`B`$ divides (resp does not divide) the set of its complex points.
##### Main Ideas of the Proof of the Patchworking Theorem
In what follows, we shall give definitions and statements related to the patchworking method. Afterwards, we shall give a sketch of proof of the Patchwork Theorem.
Preliminaries
Let $`a_1,\mathrm{},a_r`$ be completely non-degenerate polynomials in two variables with Newton polyhedra $`\mathrm{\Delta }_1`$,…,$`\mathrm{\Delta }_r`$ such that $`\mathrm{\Delta }_1`$,…,$`\mathrm{\Delta }_r`$ form a regular subdivision of a non-degenerate triangle $`\mathrm{\Delta }`$. Let $`b_t`$ be a polynomial obtained by patchworking process. Let $`B_t`$ be the homogenization of $`b_t`$. It defines a real algebraic curve $`𝒞_t`$.
In this section, we explain the construction of $`b_t`$ and give definitions we shall need in the next sections.
Recall that real polynomials in two variables belong to the ring of $`L`$-polynomials, which is isomorphic to the ring of regular function of the variety $`(𝐊^{})^2`$.
As already introduced, set $`la`$ the diffeomorphism
$$la=(l,a):(𝐊^{})^2𝐑^2\times U_𝐊^2$$
Call $`ϵ`$-tubular neighborhood $`N`$ of a smooth submanifold $`M`$ of $`(𝐊^{})^2`$ the normal tubular neighborhood of the smooth submanifold $`la(M)`$ of $`𝐑^2\times U_𝐊^2`$ whose fibers lie in fibers $`𝐑\times t\times U_𝐊\times s`$, consist of segments of geodesics which are orthogonal to intersection of $`la(M)`$ with these $`𝐑\times t\times U_𝐊\times s`$, and are contained in balls of radius $`ϵ`$ centered in the points of intersection of $`la(M)`$ with these $`𝐑\times t\times U_𝐊\times s`$. The intersection of such tubular neighborhood of $`M`$ with the fiber $`𝐑\times t\times U_𝐊\times s`$ is a normal tubular neighborhood of $`M𝐑\times t\times U_𝐊\times s`$ in $`R\times t\times U_𝐊\times s`$. (Obviously, such definition of $`ϵ`$ tubular neighborhood of $`M`$ in $`(𝐊^{})^2`$ requires $`la(M)`$ transversal to $`R\times t\times U_𝐊\times s`$.) The following statement describes cases for which such tubular neighborhood exists.
We refer to ( p.59) and also to( p.175, p178) for its proof.
###### Lemma 0.1.7.
Let $`f`$ be a polynomial and $`\mathrm{\Delta }(f)`$ its Newton polyhedron. Let $`\mathrm{\Gamma }`$ be a face of $`\mathrm{\Delta }(f)`$ such that $`f^\mathrm{\Gamma }`$ is completely non- degenerate. Let $`\rho :𝐑_+\mathrm{\Delta }(f)\times U_𝐊^2𝐊\mathrm{\Delta }(f)`$ the natural surjection . Then the hypersurface $`\rho ^1(V_{𝐊\mathrm{\Delta }(f)}(f))`$ is transversal to $`𝐑_+\mathrm{\Gamma }\times U_𝐊^2`$.
Let $`M`$ be a smooth manifold of $`(𝐊^{})^2`$ and $`\mathrm{\Delta }`$ be a non-degenerate polyhedron. Denote $`\rho :𝐑_+\mathrm{\Delta }\times U_𝐊^2𝐊\mathrm{\Delta }`$ the natural surjection.
Using the moment map $`\mu :𝐊\mathrm{\Delta }\mathrm{\Delta }`$ which maps the torus $`(𝐊^{})^2𝐊\mathrm{\Delta }^0`$ onto the interior $`\mathrm{\Delta }^0`$ of $`\mathrm{\Delta }`$, one can identify points of $`𝐑_+\mathrm{\Delta }^0`$ with points of $`\mathrm{\Delta }^0`$. In such a way, using this identification, we call $`ϵ`$-tubular neighborhood of $`M`$ defined from $`p`$ in $`\mathrm{\Delta }^0`$, the $`ϵ`$-tubular neighborhood of $`M`$ in $`\rho (D(p,ϵ)\times U_𝐊^2)(𝐊^{})^2`$ where $`D(p,ϵ)`$ is an open (euclidian) disc around $`p`$ of radius $`ϵ`$ inside $`\mathrm{\Delta }^0`$. Thereby, given a triangulation of $`\mathrm{\Delta }`$, one can choose $`D(p,ϵ)`$ sufficiently small such that $`D(p,ϵ)`$ intersects only one proper $`\mathrm{\Gamma }`$ face of the triangulation.
###### Definition 0.1.8.
Let $`M`$ a smooth manifold of $`(𝐊^{})^2`$. Let $`\mathrm{\Delta }`$ be a non-degenerate polyhedron and $`\rho :𝐑_+\mathrm{\Delta }\times U_𝐊^2𝐊\mathrm{\Delta }`$ be the natural surjection. Given $`ϵ>0`$. We call $`ϵ`$-tubular neighborhood of $`M(𝐊^{})^2`$ defined from a point $`p`$ in $`𝐑_+\mathrm{\Delta }^0`$, the $`ϵ`$-tubular neighborhood of $`M`$ in $`U(p)=\rho (D(p,ϵ)\times U_𝐊^2)(𝐊^{})^2`$ where $`D(p,ϵ)`$ is an open (euclidian) disc around $`p`$ of radius $`ϵ`$ contained in $`𝐑_+\mathrm{\Delta }^0`$.
Let $`D(p,ϵ)𝐑_+T_m^0`$ be an open (euclidian) disc around $`p`$ of radius $`ϵ`$ such that $`\mu (D(p,ϵ))`$ intersects only the face $`\mathrm{\Gamma }`$ of the triangulation of $`T_m`$ and contains $`\mathrm{\Gamma }`$. We call $`ϵ`$-tubular neighborhood of $`M`$ defined from points of $`\mathrm{\Gamma }^0`$, the $`ϵ`$-tubular neighborhood of $`M`$ in $`U(p)=\rho (D(p,ϵ)\times U_𝐊^2)(𝐊^{})^2`$. We call $`U(p)=\rho (D(p,ϵ)\times U_𝐊^2)(𝐊^{})^2(𝐊^{})^2`$ the $`ϵ`$-neighborhood of $`M`$ defined from $`\mathrm{\Gamma }^0`$.
###### Definition 0.1.9.
Let the norm in vector spaces of L-polynomials be:
$$a_wx^w=max\{|a_w||w𝐙^2\}$$
Let $`a`$ a L-polynomial over $`𝐊`$ and $`U`$ a subset of $`(𝐊^{})^2`$ we say that in $`U`$ the truncation $`a^\mathrm{\Gamma }`$ is $`ϵ`$-sufficient for $`a`$ if for any $`L`$-polynomial $`b`$ over $`K`$ such that $`\mathrm{\Delta }(b)\mathrm{\Delta }(a),b^\mathrm{\Gamma }=a^\mathrm{\Gamma },bb^\mathrm{\Gamma }aa^\mathrm{\Gamma }`$ the following conditions are verified:
1. $`USV_{(𝐊^{})^2}(b)=\mathrm{}`$
2. the set $`la(UV_{(𝐊^{})^2}(b))`$ lies in a tubular $`ϵ`$-neighborhood of
$`la(V_{(𝐊^{})^2}(a^\mathrm{\Gamma })\backslash SV_{(𝐊^{})^2}(a^\mathrm{\Gamma }))`$
3. $`la(UV_{(𝐊^{})^2}(b))`$ can be extended to the image of a smooth section of the tubular fibration $`Nla(V_{(𝐊^{})^2}(a^\mathrm{\Gamma })\backslash SV_{(𝐊^{})^2}(a^\mathrm{\Gamma }))`$
Let $`\mathrm{\Delta }`$ a polyhedron. Denote by $`C_\mathrm{\Delta }(\mathrm{\Gamma })`$ the cone $`C_\mathrm{\Delta }(\mathrm{\Gamma })=_{r𝐑_+}r.(\mathrm{\Delta }y)`$, where y is a point of $`\mathrm{\Gamma }\backslash \mathrm{\Gamma }`$. It contains the minimal affine subspace $`C(\mathrm{\Gamma })`$ of $`𝐑^2`$ containing $`\mathrm{\Gamma }`$, $`C(\mathrm{\Gamma })C_\mathrm{\Delta }(\mathrm{\Gamma })`$. Denote by $`DC_\mathrm{\Delta }^{}(\mathrm{\Gamma })`$ the set $`\{x𝐑^2,aC_\mathrm{\Delta }(\mathrm{\Gamma })a.x0\}`$. For $`A𝐑^2`$ and $`\rho >0`$, set $`𝒩_\rho (A)=\{x𝐑^2\backslash dist(x,A)<\rho \}`$. Let $`\varphi :𝒢(\mathrm{\Delta })𝐑`$ be a positive function, denote by $`DC_{\mathrm{\Delta },\varphi }(\mathrm{\Gamma })`$ the set $`𝒩_{\varphi (\mathrm{\Gamma })}(DC_\mathrm{\Delta }^{}(\mathrm{\Gamma }))\backslash _{\mathrm{\Sigma }𝒢(\mathrm{\Delta }),\mathrm{\Gamma }𝒢(\mathrm{\Sigma })}𝒩_{\varphi (\mathrm{\Sigma })}(DC_\mathrm{\Delta }^{}(\mathrm{\Sigma }))`$.
###### Definition 0.1.10.
Let $`f`$ be a Laurent polynomial over $`𝐊`$ and $`\mathrm{\Delta }`$ its Newton polyhedron. A positive function $`\varphi :𝒢(\mathrm{\Delta })𝐑`$ describes domain of $`ϵ`$-sufficiency for $`f`$ if for any proper face $`\mathrm{\Gamma }\mathrm{\Gamma }(\mathrm{\Delta })`$, for which truncation is completely non-degenerate and the hypersurface $`la(V_{\mathrm{𝐊𝐑}^2}(f^\mathrm{\Gamma }))`$ has an $`ϵ`$-tubular neighborhood, the truncation $`f^\mathrm{\Gamma }`$ is $`ϵ`$-sufficient for $`f`$ in some neighborhood of $`l^1(DC_{\mathrm{\Delta },\varphi }(\mathrm{\Gamma }))`$.
Main Ideas of the proof of the Patchworking Theorem We shall just give the main ideas of a proof of the Patchworking theorem. We refer the reader to for the complete proof.
First notice that from the Lemma 0.1.7 of Chapter 0.1.7 since polynomials $`a_1`$,…,$`a_r`$ are completely non-degenerate one can consider $`ϵ`$-tubular neighborhood around points of any hypersurface $`a_i=0`$ in $`𝐊\mathrm{\Delta }_i`$ which belong to $`𝐊\mathrm{\Gamma }`$ with $`\mathrm{\Gamma }𝒢^{}(\mathrm{\Delta }_i)`$.
Besides, recall that to each non-degenerate polyhedron $`\mathrm{\Delta }`$ one can associate projective toric variety $`𝐊\mathrm{\Delta }`$. The variety $`𝐊\mathrm{\Delta }`$ can be seen as the completion of $`(𝐊^{})^2`$ in such a way that the toric varieties associated to the proper faces of $`\mathrm{\Delta }`$ cover $`𝐊\mathrm{\Delta }\backslash (𝐊^{})^2`$.
Recall the following statement extracted from .
Lemma 0 For any polynomial $`f`$ and $`ϵ>0`$ there exists a function
$`\varphi :𝒢(\mathrm{\Delta }(f))𝐑`$ describing domain of $`ϵ`$-sufficiency for $`f`$. Denote $`𝒢=_{i=1}^r𝒢(\mathrm{\Delta }_i)`$. Define $`b(x,y,t)=b_t(x,y)`$ the polynomial in the three variables $`(x,y,t)`$. Denote $`\mathrm{\Delta }\mathrm{"}`$ the Newton polyhedron of $`b`$. $`\mathrm{\Delta }\mathrm{"}`$ is the convex hull of the graph of $`\nu `$. For $`\mathrm{\Gamma }𝒢`$ denote by $`\mathrm{\Gamma }\mathrm{"}`$ the face of $`\mathrm{\Delta }\mathrm{"}`$ which is the graph of $`\nu |_\mathrm{\Gamma }`$. For $`t>0`$, let $`j_t`$ be the embedding $`𝐑^2𝐑^3`$ given by formula $`j_t(x,y)=(x,y,lnt)`$. Let $`\psi :𝒢𝐑`$ be a positive function. For $`\mathrm{\Gamma }𝒢`$ denote by $`_{t,\psi }`$ the subset :
$$𝒩_{\psi (\mathrm{\Gamma })}j_t^1(DC_{\mathrm{\Delta }\mathrm{"}}^{}(\mathrm{\Gamma }\mathrm{"}))\backslash _{\mathrm{\Sigma }𝒢(\mathrm{\Delta }),\mathrm{\Gamma }𝒢(\mathrm{\Sigma })}𝒩_{\varphi (\mathrm{\Sigma })}j_t^1(DC_{\mathrm{\Delta }\mathrm{"}}^{}(\mathrm{\Sigma }\mathrm{"}))$$
The following Lemma is extracted from . Lemma 1 Let $`a_1`$,..$`a_r`$ be completely non-degenerate polynomials. For any $`ϵ>0`$, there exists a function $`\psi :𝒢𝐑`$ such that for any $`t]0,t_0]`$ and any face $`\mathrm{\Gamma }𝒢`$, there exists a neighborhood $`_{t,\psi }(\mathrm{\Gamma })`$ of $`\mathrm{\Gamma }`$ such that $`b_t^\mathrm{\Gamma }`$ is $`ϵ`$-sufficient for $`b_t`$ in $`l^1_{t,\psi }(\mathrm{\Gamma })`$. sketch of proof:
It follows from the existence of a positive function $`\psi :𝒢𝐑`$ such that for each proper face $`\mathrm{\Gamma }\mathrm{"}`$ of $`\mathrm{\Delta }\mathrm{"}`$, $`b^{\mathrm{\Gamma }\mathrm{"}}`$ is $`ϵ`$-sufficient for $`b`$ in a neighborhood which contains $`_{t,\psi }(\mathrm{\Gamma })`$ for any $`t`$ sufficiently small. Q.E.D. Let $`b_t`$ be a polynomial obtained by patchworking the polynomials $`a_1`$,…$`a_r`$. The Newton polyhedron $`\mathrm{\Delta }(b_t)`$ of the polynomial $`b_t`$ is the polyhedron $`\mathrm{\Delta }`$. The following lemma follows immediately from Lemma 0. Lemma 2 For any polynomial $`b_t`$ obtained by patchworking the polynomials $`a_1`$,…$`a_r`$ and $`ϵ>0`$ there exists a function $`\varphi :𝒢(\mathrm{\Delta })𝐑`$ describing domain of $`ϵ`$-sufficiency for $`b_t`$.
Besides (see ), for a well chosen function $`\varphi :𝒢(\mathrm{\Delta })𝐑`$, the $`ϵ`$-sufficiency of $`b_t^\mathrm{\Gamma }`$ for $`b_t`$ in $`l^1_{t,\psi }(\mathrm{\Gamma })`$ for $`\mathrm{\Gamma }𝒢(\mathrm{\Delta })\backslash \mathrm{\Delta }`$ implies $`ϵ`$-sufficiency for $`b_t`$ in some neighborhood of $`l^1(DC_{\mathrm{\Delta },\varphi }(\mathrm{\Gamma }))`$ and of $`l^1(DC_{\mathrm{\Delta },\varphi }(\mathrm{\Delta }))`$ in such a way that for any $`t]0,t_0]`$ a $`𝐊`$-chart of the polynomial $`b_t`$ is obtained by patchworking $`𝐊`$-chart of the polynomials $`a_1,\mathrm{},a_r`$. Hence, it follows the Patchworking Theorem. Furthermore, one can make the following remark.
Since polynomials $`a_{i,i\{1,\mathrm{},r\}}`$ and $`b_t^{\mathrm{\Delta }_i}`$ for $`t`$ sufficiently small are completely non-degenerate, from the Lemma 1 and the existence of a well chosen function $`\varphi :𝒢(\mathrm{\Delta })𝐑`$ as above, the following approximations of $`V_{𝐊\mathrm{\Delta }}(b_t)`$ can be easily deduced. Denote the gradient of the restriction of $`\nu `$ on $`\mathrm{\Gamma }`$ by $`(\nu |_\mathrm{\Gamma })`$. The truncation $`b_t^\mathrm{\Gamma }`$ equals $`a^\mathrm{\Gamma }qh_{(\nu |_\mathrm{\Gamma }),t}`$ where $`a`$ is the unique polynomial with $`\mathrm{\Delta }(a)=\mathrm{\Delta },a^{\mathrm{\Delta }_i}=a_i`$ for $`i=1,\mathrm{},r`$. From the Lemma 1, we get that for $`t]0,t_0]`$ the space $`\mathrm{𝐊𝐑}^2`$ is covered by regions in which $`V_{\mathrm{𝐊𝐑}^2}(b_t)`$ is approximated by $`qh_{(\nu |_{\mathrm{\Delta }_i}),t}^1(V_{\mathrm{𝐊𝐑}^2}(a_i))`$. Extending translations $`qh_{(\nu |_\mathrm{\Gamma }),t}`$ to $`\mathrm{\Gamma }𝒢(\mathrm{\Delta }_i)`$ and $`S_{(x^{},y^{})}(x,y)`$ to the whole space $`𝐊\mathrm{\Delta }`$, it follows the description of $`V_{𝐊\mathrm{\Delta }}(b_t)`$.
#### 0.2. Metric on $`𝐂P^2`$ and $`𝐑P^2`$
The projective space $`𝐂P^2`$ as any differentiable manifold admits a Riemannian Metric. In what follows, we shall present the Riemannian Metric of $`𝐂P^2`$ usually called Fubini-Study metric.
##### Local Euclidian metric
The topological space $`𝐂P^2`$, as complex manifold looks like locally as the $`2`$-dimensional affine space $`𝐂^2`$. In other words, $`𝐂P^2`$ admits a covering by open set $`𝐂U_i=\{(z_0:z_1:z_2)𝐂P^2|z_i0\}`$, $`i\{0,1,2\}`$ provided with charts $`\varphi _i:𝐂U_i𝐂^2`$, $`i\{0,1,2\}`$ ($`\varphi _1(z_0:z_1:z_2)=(\frac{z_0}{z_1},\frac{z_2}{z_1})`$; $`\varphi _0`$ and $`\varphi _2`$ are symmetrically defined), where all the $`\varphi _i\varphi _j^1`$ are complex analytic.
With their help, one can use the cartesian coordinates in $`𝐂^2`$ as local coordinates in $`𝐂U_i𝐂P^2`$ and carry out complex analysis on $`𝐂P^2`$.
In such a way, we may locally define $`4`$-ball of $`𝐂P^2`$ in an open $`𝐂U_i`$ as $`\varphi _i^1(B^4)𝐂U_i`$ where $`B^4`$ is a usual $`4`$-ball of $`𝐂^2`$. Passing from complex to real set of points, any $`2`$-disc of $`𝐑P^2𝐂P^2`$ is defined in an open $`𝐑U_i𝐂U_i`$, as $`\varphi _i^1(D^2)𝐑U_i`$ where $`D^2`$ is a usual $`2`$-disc of $`𝐑^2`$.
##### Riemannian metric
One can also provide $`𝐂P^2`$ with another structure by giving a euclidian metric on the tangent space of each point $`p𝐂P^2`$ that is by introducing a Riemannian metric on the complex projective space $`𝐂P^2`$. This metric on $`𝐂P^2`$ is called the Fubini-Study metric.
Let $`S^5=\{(z_0,z_1,z_2)𝐂^3|z_0.\overline{z_0}+z_1.\overline{z_1}+z_2.\overline{z_2}=1\}`$ and $`\pi _𝐂:S^5𝐂P^2`$ be the natural projection. For any point $`p𝐂P^2`$ choose a representative $`(p_0:p_1:p_2)`$ with $`p_0^2+p_1^2+p_2^2=1`$, $`\pi _𝐂^1(p)`$ is the circle $`\{p^{}|p^{}=\lambda .(p_0,p_1,p_2)with\lambda 𝐂,|\lambda |=1\}`$.
Given two points $`p=(p_0:p_1:p_2)𝐂P^2`$, $`q=(q_0:q_1:q_2)𝐂P^2`$ ,( $`p_0^2+p_1^2+p_2^2=1`$, $`q_0^2+q_1^2+q_2^2=1`$), we define the distance $`\delta `$ between $`p`$ and $`q`$ in the Fubini-Study metric as the length of the geodesic between the circles $`\pi _𝐂^1(p)`$ $`\pi _𝐂^1(q)`$ of $`S^5`$.
It is easy to verify that the Fubini-Study metric is invariant under the action of complex conjugation. Besides, when consider the restriction of $`S^5`$ to its real set of points $`S^2`$, one define the metric on $`𝐑P^2`$.
Denote $`\pi _𝐑:S^2𝐑P^2`$ the natural projection.
Given two points $`p=(p_0:p_1:p_2)𝐑P^2`$, $`q=(q_0:q_1:q_2)𝐑P^2`$, ($`p_0^2+p_1^2+p_2^2=1`$, $`q_0^2+q_1^2+q_2^2=1`$), $`\pi _𝐑^1(p)=\{(p_0,p_1,p_2),(p_0,p_1,p_2)\}`$
$`\pi _𝐑^1(q)=\{(q_0,q_1,q_2),(q_0,q_1,q_2)\}`$, noticing that the antipodal mapping $`A:S^2S^2`$ is an isometry, we define the distance $`\delta `$ between $`p`$ and $`q`$ in the Fubini-Study metric as the minimal arc length of the arcs of $`S^2`$ between $`(p_0,p_1,p_2)\pi _𝐑^1(p)`$, and $`\pi _𝐑^1(q)=\{(q_0,q_1,q_2),(q_0,q_1,q_2)\}`$.
Given two real points of $`S^2S^5`$, it is easy to see (since geodesic minimizes arc length between two of its points) that any point of the geodesic from these two points is real.
In such a way, in the Fubini-Study metric of $`𝐂P^2`$, we define any $`4`$-ball of $`𝐂P^2`$ by $`\pi _𝐂()`$ where $``$ is an an ellipsoid of $`S^5`$; and any $`2`$-disc of $`𝐑P^2𝐂P^2`$ by $`\pi _𝐂(E)=\pi _𝐑(E)`$ where $`E`$ is an ellipse of $`S^2S^5`$.
## Part II Arnold Surfaces of Harnack Curves
The maximal number of connected components of the real point set of curves of degree $`m`$ is $`\frac{(m1)(m2)}{2}+1`$. Curves with this maximal number are called $`M`$-curves. In this section, we study the construction of some $`M`$-curves called Harnack curves and prove that Arnold surfaces of the so-called Harnack curves of even degree are standard.
### Chapter 1 Combinatorial patchworking construction for Harnack curves
Preliminaries Recall that the pair $`(𝐑P^2,𝐑𝒜)`$ where $`𝒜`$ is a non-singular real plane curve is determined up to homeomorphism by the real components of the curve and their relative location. In the case of even degree curve $`𝒜`$, each connected component of the real point set $`𝐑𝒜`$ is situated in $`𝐑P^2`$ as the boundary of an embedded disc and is called an oval. In the case of odd degree curve $`𝒜`$, the real point set $`𝐑𝒜`$ has besides oval one connected component situated in $`𝐑P^2`$ as an embedded projective line. It is called the one side component of $`𝐑𝒜`$. An oval divides $`𝐑P^2`$ into two components. The orientable component (i.e the component homeomorphic to a disc) is called the inside of the oval. The non-orientable component (i.e the component homeomorphic to a Möbius strip) is called the outside of the oval. Given $`𝒜`$ a non-singular real plane curve and $`F(x_0,x_1,x_2)`$ its polynomial, denote by $`𝐑P_+^2`$ the subset of $`𝐑P^2`$ $`\{x𝐑P^2|F(x_0,x_1,x_2)0\}`$. An oval is said outer (resp inner) if it bounds a component of $`𝐑P_+^2`$ from the outside (resp, the inside). In the case of an even degree curve $`𝒜`$, we can always assume (changing the sign of the polynomial $`F(x_0,x_1,x_2)`$ giving $`𝒜`$ if necessary) that $`𝐑P_+^2=\{x𝐑P^2|F(x_0,x_1,x_2)0\}`$ is the orientable component of $`𝐑P^2`$. In such away, in case of an even degree curve, since ovals lying in an even number of consecutive ovals are outer while the ovals lying in an odd number of consecutive ovals are inner; one also calls even ovals the outer ovals and odd ovals the inner ovals.
In the case of an odd degree curve $`𝒜`$, the definition of outer and inner oval apply only to oval which does not intersect the line at infinity. Call zero oval an oval of a curve of odd degree intersecting the line at infinity.
The pair $`(𝐑P^2,𝐑𝒜)`$ where$`𝒜`$ is a non singular curve of degree $`m`$ is defined by the scheme of disposition of the real component $`𝐑𝒜`$. This scheme is called the real scheme of the curve $`𝒜`$. In what follows, we use the following usual system of notations for real scheme (see ). A single oval is denoted by $`1`$. The one side component is denoted by $`J`$. The empty real scheme is denoted by $`0`$. If one set of ovals is denoted by $`A`$, then the set of ovals obtained by addition of an oval which contained $`A`$ in its inside component is denoted by $`1A`$. If the real scheme of a curve consists of two disjoint sets of ovals denoted by $`A`$ and$`B`$ in such a way that no oval of one set contains an oval of the other set in its inside component, then the real scheme of this curve is denoted by : $`AB`$. If one set of ovals is denoted by $`A`$, then the set $`A\mathrm{}A`$ where $`A`$ occurs n times, is denoted by $`n\times A`$; a set $`n\times 1`$ is denoted by $`n`$.
Furthermore, if $`𝒜`$ is a curve of which the set of real points $`𝐑𝒜`$ divides the set of complex points $`𝐂𝒜`$ in two connected components ( which induces two opposite orientations on $`𝐑𝒜`$), then the curve $`𝒜`$ is said of type $`I`$. Otherwise, $`𝒜`$ is said of type $`II`$. Hence, the real scheme of a curve $`𝒜`$ of degree $`m`$ is of type $`I`$ (resp, of type $`II`$), if any curve of degree $`m`$ having this real scheme is of type $`I`$(resp, of type $`II`$). Otherwise (i.e if there exists both curves of type $`I`$ and curves of type $`II`$ with the given real scheme), we say that the real scheme is of indeterminate type.
Thus, isotopy of $`𝐂P^2`$ which connects complex points of curves and commutes with the complex conjugation of $`𝐂P^2`$ , called conj-equivariant isotopy of $`𝐂P^2`$, provides a convenient way to classify pairs $`(𝐂P^2,𝐂𝒜)`$. We can already notice that since $`M`$-curves are curves of type $`I`$, their real scheme is sufficient to classify their complex set of points in $`𝐂P^2`$ up to conj-equivariant isotopy.
#### 1.1. Constructing Harnack curves with Harnack’s initial method
In 1876, Harnack proposed a method for constructing $`M`$-curves. Recall briefly Harnack’s initial construction of Harnack’s $`M`$-curves. (The detailed method can be found in )
Start with a line as $`M`$-curve of degree $`1`$. Then, consider a line $`L`$ which intersects it in one point. Assume an $`M`$-curve of degree $`m`$ $`𝒜_m`$ has been constructed at the step $`m`$ of the construction. At the step $`m+1`$ of the construction, the $`M`$-curve of degree $`m+1`$ is obtained from classical deformation (see Classical Small Perturbation Theorem for the definition of classical deformation) of $`𝒜_mL`$. The resulting $`M`$-curve of degree $`m+1`$ intersects the line $`L`$ in $`m+1`$ real points.
Passing from curves to polynomials giving these curves, an $`M`$-curve $`\stackrel{~}{X}_{m+1}=0`$ is constructed from an $`M`$-curve $`\stackrel{~}{X}_m=0`$ of degree $`m`$ by the formula $`\stackrel{~}{X}_{m+1}=x_0.\stackrel{~}{X}_m+t.C_{m+1}`$ where $`x_0`$ is a line, $`C_{m+1}=0`$ is $`m+1`$ parallels lines which intersect $`x_0=0`$ in $`m+1`$ points. It is essential in the construction that the curves $`\stackrel{~}{X}_m=0`$, $`x_0=0`$ and $`C_{m+1}=0`$ do not have common points.
Thus, one can choose projective coordinates $`(x_0:x_1:x_2)`$ of $`𝐂P^2`$, in such a way that $`\stackrel{~}{X}_m=0`$ is an homogeneous polynomial of degree $`m`$ in the variables $`x_0`$, $`x_1`$,$`x_2`$ and $`C_{m+1}=0`$ is an homogeneous polynomial of degree $`m+1`$ in the variables $`x_2`$,$`x_1`$. For sufficiently small $`t>0`$, $`\stackrel{~}{X}_{m+1}=0`$ is an Harnack curve of degree $`m+1`$.
In particular, this method provides curves with real scheme:
for even $`m=2k`$
(1.1)
$$1\frac{(k1)(k2)}{2}\frac{3k(k1)}{2}$$
for odd $`m=2k+1`$
(1.2)
$$Jk(2k1)$$
Let us denote $`_m`$ and call Harnack curve of degree $`m`$, any curve of degree $`m`$ with the real scheme above. Besides, we shall call Harnack polynomial of degree $`m`$ any polynomial giving the Harnack curve of degree $`m`$. We shall call curve of type $``$ a curve $`_m`$ which intersects the line at infinity of $`𝐂P^2`$ in $`m`$ real distinct points. Without loss of generality, one can always assume that the line $`L`$ is the line at infinity. Therefore, curves $`_m`$ constructed by the Harnack’s method are curves of type $``$.
#### 1.2. Patchworking construction for Harnack curves
In what follows, we shall give a construction of Harnack curves provided by the patchworking construction of curves .
Recall briefly the patchworking construction procedure of curves due to Viro . (We refer the reader to and also to the Preliminary section for details.)
Initial data
1. Let $`m`$ be a positive integer. Let $`T_m`$ be the triangle
$$\{(x,y)𝐑^2x0,y0,x+ym\}$$
(Up to linear change of coordinates $`(x_0:x_1:x_2)`$ of $`𝐂P^2`$, the convex hull of $`T_m`$ may be the Newton polyhedron of the affine polynomial $`f(x,y)=F(1,x,y)`$ associated to a homogeneous polynomial $`F(x_0,x_1,x_2)`$ of degree $`m`$.)
2. Let $`\tau `$ be a triangulation of $`T_m`$ whose vertices have integer coordinates. Call regular triangulation of $`T_m`$ a triangulation of $`T_m`$ such that there exists with a convexifying function $`\nu :T_m𝐑`$ (that is a piecewise-linear function $`\nu :T_m𝐑`$ which is linear on each triangle of the triangulation $`\tau `$ and not linear on the union of two triangles.) Assume $`\tau `$ regular.
3. Suppose that some distribution of signs $`\chi `$ at the vertices of the triangulation is given. Denote the sign $`\pm `$ at the vertex with coordinates $`(i,j)`$ by $`ϵ_{(i,j)}`$.
Combinatorial procedure
Take the square $`T_m^{}`$ made of $`T_m`$ and its symmetric copies $`T_m^x=s_x(T_m),T_m^y=s_y(T_m),T_m^{xy}=s_{xy}(T_m)`$ where $`s_x,s_y,s=s_xs_y`$ are reflections with respect to the coordinate axes. The resulting space is homeomorphic to $`𝐑^2`$. Extend the triangulation $`\tau `$ of $`T_m`$ to a symmetric triangulations$`\tau _{}`$ of $`T_m^{}`$. Extend the distribution of signs to a distribution at the vertices of $`T_m^{}`$ which verifies the modular properties:
$`g^{}(ϵ_{i,j})x^iy^j=ϵ_{g(i,j)}x^iy^j`$ for $`g=s_x,s_y,s`$.
If a triangle of $`T_m^{}`$ has vertices of different signs, consider a midline separating them. Denote by $`L`$ the union of such midlines. Glue by $`s`$ the opposite sides of $`T_m^{}`$. The resulting space $`\overline{T}_m^{}`$ is homeomorphic to the projective plane $`𝐑P^2`$. Denote $`\overline{L}`$ the image of $`L`$ in $`\overline{T}_m^{}`$. We shall say that $`T_m^{}`$, $`L`$, $`\overline{T}_m^{}`$,$`\overline{L}`$ are obtained from $`T_m,\tau ,\chi `$ by combinatorial patchworking.
###### Theorem 1.2.1.
Polynomial Patchworking Define the one-parameter family of polynomials
$$b_t=b_t(x,y)=\underset{(i,j)verticesof\tau }{}ϵ_{i,j}x^iy^jt^{\nu (i,j)}$$
where $`\nu `$ is a function convexifying the triangulation $`\tau `$ of $`T_m`$. Let $`\overline{T}_m^{}`$,$`\overline{L}`$ obtained from $`T_m,\tau ,\chi `$ by combinatorial patchworking.
Denote by $`B_t=B_t(x_0,x_1,x_2)`$ the corresponding homogeneous polynomials:
$$B_t(x_0,x_1,x_2)=x_0^mb_t(x_1/x_0,x_2/x_0)$$
Then there exists $`t_0>0`$ such that for any $`t]0,t_0]`$ the equation $`B_t(x_0,x_1,x_2)=0`$ defines in $`𝐑P^2`$ the set of real points of an algebraic curve $`C_t`$ such that the pair $`(𝐑P^2,𝐑C_t)`$ is homeomorphic to the pair $`(\overline{T}_m^{},\overline{L})`$. Such a curve is called a $`T`$-curve .
Let $`m`$ be a positive integer.
Let $`T_m`$ be the triangle
$$\{(x,y)𝐑^2x0,y0,x+ym\}$$
Define the regular triangulation $`\tau `$ of $`T_m`$ as follows:
1. All integer points of $`T_m`$ are vertices of the triangulation $`\tau `$.
2. Each proper face of the triangulation is contained in one segment of the set:
$$\{(x+y=i)_{i=1,\mathrm{},m},(x=i,ymi)_{i=1,\mathrm{},m},(y=i,xmi)_{i=1,\mathrm{},m}\}$$
Choose the following distribution of signs called ” Harnack distribution” at each integer points $`(i,j)`$ of $`T_m`$ :
1. If i,j are both even, the integer point $`(i,j)`$ gets the sign $`ϵ_{(i,j)}=`$
2. If $`i`$ or $`j`$ is odd, the integer point $`(i,j)`$ gets the sign $`ϵ_{(i,j)}=+`$.
The following propositions are deduced from a more general statement due to I.Itenberg .
###### Proposition 1.2.2.
Let $`m=2k`$ be a positive even integer. The patchworking process applied to $`T_m`$ with regular triangulation $`\tau `$ and the Harnack distribution of signs at the vertices produces a $`T`$-curve which is the Harnack curve $`_{2k}`$ of degree $`2k`$.
###### Proposition 1.2.3.
Let $`m=2k+1`$ be a positive even integer. The patchworking process applied to $`T_m`$ with regular triangulation $`\tau `$ and the Harnack distribution of signs at the vertices produces a $`T`$-curve which is the Harnack curve $`_{2k+1}`$ of degree $`2k+1`$ .
Now we shall work out the polynomial entering in this patchworking construction (following essentially Viro’s initial method with some modifications necessary for the sequel). We want in particular to stress the recursive character of the patchworking construction of Harnack’ s curves.
Fix $`m>0`$ integer
1. An essential element in the construction is the so-called convexifying function. We choose the convexifying function $`\nu _m:T_m𝐑`$ in such a way that it would be the restriction of $`\nu _{m+1}`$. Namely, we extend the convexifying function $`\nu _m:T_m𝐑`$ to a convexifying function $`\nu _{m+1}`$ of the triangulation of $`T_{m+1}=T_mD_{m+1,m}`$ with
$$D_{m+1,m}=\{(x,y)𝐑^20xm,0ym,x+ym+1\}$$
We shall assume the same notation for $`\nu _m`$ and its extension $`\nu _{m+1}`$, and for any $`m>1`$ denote by $`\nu `$ a convexifying function of the triangulation of $`T_m`$.
2. Now, we construct a polynomial in two variables
$$x_{m,t}(x,y)=\underset{(i,j)verticesofT_m}{}ϵ_{i,j}a_{i,j}x^iy^jt^{\nu (i,j)}$$
where: $`ϵ_{(i,j)}`$ is the sign at the vertex $`(i,j)`$ given by the Harnack distribution of sign, $`\nu (i,j)`$ is the value of the convexifying function $`\nu :T_m𝐑`$ on $`(i,j)`$. The numbers $`t`$ and $`a_{i,j}`$ are parameters which shall be chosen. We shall denote the polynomial $`x_{m,t}(x,y)`$ by $`x_{m(x,y;\stackrel{}{a}_m,t)}`$ or $`x_{m;\stackrel{}{a}_m,t}`$. Let $`\stackrel{~}{X}_{m,t}(x_0,x_1,x_2)`$ be the homogenization of $`x_{m,t}`$ obtained by remaining $`x`$ as $`x_1`$, $`y`$ as $`x_2`$, and adding the third variables $`x_0`$ so as to obtain :
$$\stackrel{~}{X}_{m,t}(x_0,x_1,x_2)=\underset{(i,j)verticesofT_m}{}ϵ_{i,j}a_{i,j}x_1^ix_2^jx_0^{mij}t^{\nu (i,j)}$$
We shall denote the polynomial $`\stackrel{~}{X}_{m,t}(x_0,x_1,x_2)`$ by $`\stackrel{~}{X}_{m(x_0,x_1,x_2;\stackrel{}{a}_m,t)}`$ or $`\stackrel{~}{X}_{m;\stackrel{}{a}_m,t}`$.
Note that
$$\stackrel{~}{X}_{m,t}=x_0.\stackrel{~}{X}_{m1,t}+C_{m,t}$$
where the polynomial $`C_{m,t}`$ is an homogeneous polynomial of degree in only $`x_1`$ and $`x_2`$.
Using the patchworking method, one can give an inductive construction, we shall call T-inductive construction of some Harnack polynomials.
Assume $`\stackrel{}{a}_m,t_m`$ given in such a way that for any $`t]0,t_m[`$, $`\stackrel{~}{X}_{m;\stackrel{}{a}_m,t}`$ is a Harnack polynomial of degree $`m`$. Then, the Harnack polynomial of degree $`m+1`$ is constructed as follows:
1. Property 1 Choose $`\stackrel{}{c}_{m+1}`$ in such a way that for any $`t]0,t_m[`$, singularities of $`x_0.\stackrel{~}{X}_{m;\stackrel{}{a}_m,t}`$ do not belong to the curve $`𝒞_{m+1;\stackrel{}{c}_{m+1},t}`$.
2. Property 2 Set $`\stackrel{}{a}_{m+1}=(\stackrel{}{a}_m,\stackrel{}{c}_{m+1})`$. Then, $`t_{m+1}`$ is chosen as follows:
1. $`t_{m+1}t_m`$
2. $`t_{m+1}`$ is the biggest $`\tau t_m`$ such that for any $`t]0,\tau [`$ the polynomial $`\stackrel{~}{X}_{m+1;\stackrel{}{a}_{m+1},\tau }=x_0.\stackrel{~}{X}_{m;\stackrel{}{a}_m,\tau }+C_{m+1;\stackrel{}{c}_{m+1},\tau }`$ is a Harnack polynomial of degree $`m+1`$.
Following the preliminary section, consider $`\rho ^m:𝐑_+T_m\times U_𝐂^2𝐂T_m`$ the natural surjection. Denote $`l_m=\{x0,y0,x+y=m\}`$ the hypotenuse of $`T_m`$. The variety $`𝐂T_m`$ is homeomorphic to $`𝐂P^2`$ with line at infinity $`L`$ of $`𝐂P^2`$ such that $`𝐂l_mL`$.
(The real line $`𝐑L𝐑P^2`$ is obtained from the square $`\{(x,y)|x|+|y|=m\}`$ by gluing of the sides by $`s`$ the central symmetry with center $`0=(0,0)`$).
From the Patchworking construction of Harnack curves it follows that:
###### Corollary 1.2.4.
The Harnack curve $`_m`$ of degree $`m`$ obtained by patchworking process applied to $`T_m`$ with triangulation $`\tau `$ and Harnack distribution of signs at the vertices intersects the line at infinity $`L𝐂l_m`$ of $`𝐂P^2𝐂T_m`$ in $`m`$ real points. Thus, the resulting curve $`_m`$ is a curve of type $``$. Let us denote $`a_1,\mathrm{},a_m`$ the real points of the intersection
$$𝐑_m𝐑L=\{a_1,\mathrm{},a_m\}$$
proof: First of all, recall that in the patchworking construction of $`_m`$, the real line at infinity $`𝐑L`$ of $`𝐑P^2`$ is obtained from the square $`\{(x,y)𝐑^2||x|+|y|=m\}`$ by gluing of the sides by $`s`$ the central symmetry with center $`0=(0,0)`$. Furthermore, according to the patchworking construction, the real point set of the Harnack curve separates any two consecutive integer points of the square $`\{(x,y)𝐑^2||x|+|y|=m\}`$ with different signs. Thus, the corollary follows immediately from the Harnack distribution of signs on $`T_m`$ and its extension to its symmetric copies.Q.E.D
Moreover, one can make the following remark:
###### Remark 1.2.5.
Harnack curves are curves of type $`I`$, in other words their real set of points are orientable. Let $`_m`$ be the Harnack curve of degree $`m`$. Assume $`𝐑_m`$ given with an orientation. There is only one orientation of $`𝐑L`$ compatible with the deformation of $`_mL`$ to $`_{m+1}`$ and an orientation of $`𝐑_{m+1}`$. In such a way, we obtain the relative orientation of the connected components of the real point set of the Harnack curves.
### Chapter 2 Morse-Petrovskii’s Theory of Harnack curves
In this chapter, at first we shall study critical points of Harnack polynomials from a view-point initiated by Petrovskii. This investigation is analogous to those of Morse on critical points of functions.
Then, using rigid isotopy classification of Harnack curves, we shall construct deformation of Harnack polynomial and deduce a description up to conj-equivariant isotopy of $`𝐂P^2`$ of the complex set of points of Harnack curve (see proposition 2.3.7 of Chapter 2.3.7 and theorem 2.3.9 of Chapter 2.3.9).
We shall divide this chapter into three sections.
In the first section, we start with proposition 2.1.1 of Chapter 2.1.1 in which we prove that up to slightly modify coefficients of a Harnack polynomial, the number of its critical points depends only on its degree. Then, in proposition 2.1.4 of Chapter 2.1.4, we precise, up to change the system of projective coordinates if necessary, the number of critical points of index $`1`$ with positive and respectively negative critical value of such modified Harnack polynomials. Such results may be obtained from the Viro’s patchworking method for $`T`$-Harnack polynomials, it is explicitely explained in Shustin’s paper .
In the second section, we consider only real points set of Harnack curves. In Theorem 2.2.1, we prove that isotopy also implies rigid isotopy for Harnack curves $`_m`$ of degree $`m`$. In the third section, we define in Proposition 2.3.1 of Chapter 2.3.1, a deformation of $`_m`$, to a singular irreducible curve of which singular points are critical points of index 1 with positive critical value of $`_m`$.
The main result of this chapter is given in the third and last section. It provides link between properties of $`L`$-curves () and the Harnack curves. In Theorem 2.3.9 of Chapter 2.3.9, we present any Harnack curve $`_m`$ of degree $`m`$ up to conj-equivariant isotopy of $`𝐂P^2`$ as follows. Outside a finite number of $`4`$-balls $`B(a_i)`$ globally invariant by complex conjugation, $`_m`$ splits in $`m`$ non-intersecting projective lines minus their intersection with the $`B(a_i)`$; inside any $`4`$-ball $`B(a_i)`$ it is the perturbation of a crossing.
#### 2.1. Critical points of Harnack polynomials
Recall a Harnack polynomial of degree $`m`$ is a non-singular homogeneous polynomial in three variables such that its set of zeros has real scheme: for even $`m=2k`$:
$$1\frac{(k1)(k2)}{2}\frac{3k(k1)}{2}$$
for odd $`m=2k+1`$:
$$Jk(2k1)$$
In this section, we study the critical points and the level surfaces of Harnack polynomials.
Given $`x_0:=0`$ the line at infinity and $`R(x_0,x_1,x_2)`$ an homogeneous polynomial of degree $`m`$, we call affine restriction $`r(x,y)`$ of $`R(x_0,x_1,x_2)`$ the unique polynomial such that $`R(x_0,x_1,x_2)=x_0^m.r(x_1/x_0,x_2/x_0)`$. In such a way, we consider $`𝐂P^2`$ as the completion of $`𝐂^2`$ with $`x_0:=0`$ the line at infinity. We call critical point of $`r(x,y)`$,(and by misuse critical point of $`R(x_0,x_1,x_2)`$), any point $`(x_0,y_0)`$ (finite or not) such that $`r_x(x_0,y_0)=0`$, $`r_y(x_0,y_0)=0`$.
We shall consider generic polynomials. To be more precise, let us call an homogeneous non-singular real polynomial $`R(x_0,x_1,x_2)`$ with affine restriction $`r(x,y)`$ regular if :
1. none of critical points of $`r(x,y)`$ lies on the line at infinity
2. any two different real critical points $`r(x,y)`$ have different critical values.
From (, p.359 Lemma 2), given a polynomial of a smooth algebraic curve, one can always perturb its coefficients so as to have $`1)`$ and $`2)`$ as above.
Bringing together Lemma 0.0.11 of Chapter 0.0.11 and Petrovskii’s Lemma (, p.359 Lemma 2), it follows that the set of regular polynomials is open and dense in the set of all polynomials.
We shall call polynomial of type $``$ a regular Harnack polynomial of degree $`m`$ giving a Harnack curve $`_m`$ which intersects the line at infinity of $`𝐑P^2`$ in $`m`$ real distincts points. We shall call polynomial of type $`^0`$ a polynomial of type $``$ such that the line at infinity of $`𝐑P^2`$ intersects the Harnack curve $`_m`$ of degree $`m`$ in $`m`$ real distinct points which belong to the same connected component of $`𝐑_m`$: in case $`m=2k+1`$ the real connected component of $`_{2k+1}`$ homeomorphic to the projective line; in case $`m=2k`$ the non-empty oval of $`_{2k}`$.
Let us start by the proposition 2.1.1 of Chapter 2.1.1 in which we find the number of critical points of all indices of any regular Harnack polynomial of degree $`m`$.
For a regular polynomial of homogeneous degree $`m`$, $`R(x_0,x_1,x_2)`$ with affine restriction $`r(x,y)`$ ,$`R(x_0,x_1,x_2)=x_0^m.r(x_1/x_0,x_2/x_0)`$, let us denote by $`c_i(R)`$ the number critical points of index $`i`$ of $`r(x,y)`$ and set $`c(R)=(c_0(R),c_1(R),c_2(R))`$.
###### Proposition 2.1.1.
Let $`B_m(x_0,x_1,x_2)`$ be a Harnack polynomial of degree $`m`$ of type $`^0`$.
Then, up to change the sign of $`B_m`$:
1. $`c_0(B_m)+c_1(B_m)+c_2(B_m)=(m1)^2`$
2. 1. For even $`m=2k`$,
$$c(B_{2k})=(\frac{(k1)(k2)}{2},k(2k1),\frac{3k(k1)}{2})$$
2. For odd $`m=2k+1`$,
$$c(B_{2k+1})=(\frac{k(k1)}{2},k(2k+1),\frac{k(3k1)}{2})$$
proof: Let $`b_m(x,y)`$ be the affine restriction of the polynomial $`B_m(x_0,x_1,x_2)`$. For any oval which does not intersect the line at infinity, consider the disc with boundary the oval. In case of even oval, the gradient of the affine polynomial $`b_m(x,y)`$ points inward. Therefore, the maximum of the polynomial $`b_m`$ on the disc is in the interior of the oval and is necessary a critical point of index $`2`$ of $`b_m`$.
In case of odd oval, the gradient of the affine polynomial $`b_m(x,y)`$ points outward. Therefore, the minimum of the polynomial $`b_m`$ on the disc is in the interior of the oval and is necessarily a critical point of index $`0`$ of $`b_m`$. This implies the following inequalities (up to change the sign of $`B_m`$ and in case of even degree curve according to the convention. -ovals not lying within other ovals or lying inside an even number number of consecutive ovals are outer, while ovals lying within an odd number of ovals are inner-)
for even $`m=2k`$:
$$c_0\frac{(k1)(k2)}{2},c_2\frac{3k(k1)}{2}$$
for odd $`m=2k+1`$:
$$c_0\frac{k(k1)}{2},c_2\frac{k(3k1)}{2}$$
Furthermore, since $`_m`$ intersects the line at infinity in $`m`$ real points we have $`c_0c_1+c_2=1m`$ (see ). Thus, it follows from the inequalities above that for even $`m=2k`$ :
(2.1)
$$c_1k(2k1)$$
for odd $`m=2k+1`$ :
(2.2)
$$c_1k(2k+1)$$
and thus $`c_0+c_1+c_2(m1)^2`$.
From the Bezout’s theorem, the number of non-degenerate critical points of a homogeneous polynomial of degree $`m`$ does not exceed $`(m1)^2`$. Therefore $`c_0+c_1+c_2=(m1)^2`$ and the inequalities 2.1 and 2.2 are equalities.
Q.E.D.
Our next aim is to find the number of critical points of indice 1 with negative (resp, positive) critical value of a any polynomial of type $`^0`$. This will be done in Proposition 2.1.4 of Chapter 2.1.4.
For a regular polynomial, let us denote by $`c_1^{}(R)`$ (resp, $`c_1^+(R)`$) the number of critical points of index 1 of $`R`$ of negative (resp, positive) critical value. Besides, denote by $`c_1^{}(R)`$ the number of critical points of index 1 of $`R`$ with positive critical value $`c_0`$ such that as $`c`$ increases from $`c_0ϵ`$ to $`c_0+ϵ`$ the number of connected components of $`M_c=\{(x,y)𝐑^2r(x,y)>c\}`$ of which boundary contains the line at infinity increases by $`1`$. We shall denote $`S_{2k}^{}`$ the set constituted by the $`c_1^{}(B_{2k})`$ critical points of index 1 of a polynomial $`B_{2k}`$ of type $`^0`$.
Our proof uses Petrovskii’s theory introduced in . It is based on the consideration of the lines $`b_m(x,y)=c`$ when $`c`$ crosses the critical value of a polynomial $`B_m`$ of type $`^0`$.
These last investigations are analogous to Morse theory.
Let us recall Petrovskii’s Lemmas implicit in
(p.360-361). First, introduce definitions.
Define $`𝐑P^2`$ as a Möbius strip and a disc $`D^2`$ glued along their boundaries. The core of the Möbius strip intersects the line at infinity in a finite number of points. Call ”odd component of the affine plane $`𝐑^2`$” the core of the Möbius strip minus its intersection with the line at infinity.
###### Lemma 2.1.2.
Let $`A`$ a curve of degree $`m`$ given by a regular polynomial $`R(x_0,x_1,x_2)=x_0^mr(x_1/x_0,x_2/x_0)`$ with $`(m1)^2`$ critical points. Let $`(x_0,y_0)`$ be a finite real critical point of $`r(x,y)`$. Setting $`r(x_0,y_0)=c_0`$, assume that $`r^1[c_0ϵ,c_0+ϵ]`$ contains no critical point other than $`(x_0,y_0)`$.
Assume the Hessian $`H(x_0,y_0)<0`$, (i.e $`(x_0,y_0)`$ is a critical point of index $`1`$). (\*) Assume $`r(x,y)`$ is of even degree, and $`c`$ decreases from $`c_0+ϵ`$ to $`c_0ϵ`$. Then there are three possibilities:
1. When $`c`$ takes value $`c_0`$, one outer oval touches another oval (outer or inner), then one outer oval disappears.
2. When $`c`$ takes value $`c_0`$, one oval (outer or inner) touches itself.
1. If one of the components $`r(x,y)>r(x_0,y_0)+ϵ`$ ; $`r(x,y)<r(x_0,y_0)ϵ`$ contains the odd component of $`𝐑^2`$, then one inner oval appears.
2. Otherwise, one outer oval becomes inner. (This case may present itself at most once as $`c`$ varies from $`+\mathrm{}`$ to $`\mathrm{}`$.)
(\*) Assume r(x,y) is of odd degree and $`c`$ decreases from $`c_0+ϵ`$ to $`c_0ϵ`$. Then there are three possibilities:
1. When $`c`$ takes value $`c_0`$, one outer oval touches an other oval or the one-side component, then the outer oval disappears.
2. When $`c`$ takes value $`c_0`$, one oval (outer or inner) touches itself then one inner oval appears.
3. When $`c`$ takes value $`c_0`$, one zero oval or the odd-component of the curve touches itself or an other zero oval, then there are three possibilities:
Consider the regions $`G_i`$ $`i=1,\mathrm{},r`$ where $`r(x,y)>c_0`$ which contain segment of the line at infinity of $`𝐑P^2`$ on their boundary.
1. the boundary of one region $`G_i`$ touches itself (in one point which does not belong to the line at infinity), then one inner oval appears.
2. one region $`G_i`$ touches the boundary (not contained in the line at infinity) of another such regions $`G_j`$ then the regions coalesce and a zero-oval disappears or a zero oval appears.
The next Lemma is implicit in(, Lemma 4, p.361)
###### Lemma 2.1.3.
(line at infinity () Let $`A`$ a curve of degree $`m`$ given by a regular polynomial $`R(x_0,x_1,x_2)=x_0^m.r(x_1/x_0,x_2/x_0)`$. with $`(m1)^2`$ critical points.
Assume furthermore that $`A`$ intersects the real line at infinity in $`k`$ $`(k>1)`$ distinct points. Let $`M_c=\{(x,y)𝐑^2r(x,y)>c\}`$ Let $`c_M`$ be the maximal critical value, and $`c_m`$ be the minimal critical value.
1. For $`c>c_M`$, $`M_c`$ has $`k`$ connected components.
2. For $`c<c_m`$, $`M_c`$ is simply connected. Two of the connected components above can coalesce only when $`r(x,y)`$ passes a critical value $`c`$ of index 1.
###### Proposition 2.1.4.
Let $`m0`$ and $`B_m(x_0,x_1,x_2)`$ be a Harnack polynomial of type $`^0`$. Then, up to change the sign of $`B_m`$:
1. for even $`m=2k`$;
$$c_1^{}(B_{2k})=\frac{k(3k1)}{2}$$
$$c_1^+(B_{2k})=\frac{k(k1)}{2},c_1^{}(B_{2k})=k1$$
2. for odd $`m=2k+1`$
$$c_1^{}(B_{2k+1})=\frac{k(3k+3)}{2}$$
$$c_1^+(B_{2k+1})=\frac{k(k1)}{2},c_1^{}(B_{2k+1})=0$$
proof :
The proof of Proposition 2.1.4 of Chapter 2.1.4 is based on the Lemma 2.1.2 of Chapter 2.1.2 and Lemma 2.1.3 of Chapter 2.1.3.
Let $`_m`$ be the Harnack curve of degree $`m`$ given by the polynomial $`B_m(x_0,x_1,x_2)=x_0^mb_m(x_1/x_0,x_2/x_0)`$ of type $`^0`$. Consider the pencil of curves given by polynomials $`x_0^m(b_m(x_1/x_0,x_2/x_0)c)`$ with $`c𝐑`$.
1. Let $`c`$ decrease from $`0`$ to $`c_mϵ`$, $`ϵ>0`$, the inner ovals shrink and disappear, the outer disappear in a unique oval. It follows from the Lemma 2.1.3 of Chapter 2.1.3 that the set $`M_{c_mϵ}`$ is simply connected. Furthermore, from Lemma 2.1.2 of Chapter 2.1.2 : one outer oval can touch another outer oval and then disappears as $`c`$ decreases from $`c_0+ϵ`$ to $`c_0ϵ`$, if and only if $`c_0`$ is a critical value of index $`1`$ of the polynomial $`b_m`$. (\*)Let $`_{2k}`$ be the Harnack curve of even degree $`2k`$
When $`c`$ decreases from $`0`$ to $`c_mϵ`$, outer ovals expand, then coalesce and finally disappear in a unique oval. Moreover, from Lemma 2.1.2 of Chapter 2.1.2, to each coalescence of a outer oval $`𝒪`$ with an other outer oval is associated a critical point $`(x_0,y_0)`$ of index 1 with $`b_{2k}(x_0,y_0)=c_0`$; $`c_0]c_m,0[`$ such that as $`c`$ decreases from $`c_0+ϵ`$ to $`c_0ϵ`$ the oval $`𝒪`$ disappears. Denote $`S_{2k}^{}`$ the set of these critical points. The curve $`_{2k}`$ has $`\frac{3k(k1)}{2}`$ outer ovals and $`(k+1)`$ connected components of $`M_0=\{(x,y)𝐑^2|b_{2k}(x,y)>0\}`$ contain a segment of the line at infinity of $`𝐑P^2`$ on their boundary. Therefore, the number of critical points of $`S_{2k}^{}`$ is
(2.3)
$$c_1^{}(B_{2k})\frac{3k(k1)}{2}1+(k+1)=\frac{k(3k1)}{2}$$
(\*)Let $`_{2k+1}`$ be the Harnack curve of odd degree $`2k+1`$
When $`c`$ decreases from $`0`$ to $`c_mϵ`$, outer ovals and regions $`G_i`$ where $`b_{2k+1}>c_0`$ which intersect the line at infinity of $`𝐑P^2`$ on their boundary, expand, then coalesce, and finally disappear in a simply connected region.
Moreover, from Lemma 2.1.2 of Chapter 2.1.2, one can associate to each coalescence of one outer oval $`𝒪`$ with an other outer oval or with one of the regions $`G_i`$ a critical point $`(x_0,y_0)`$ of index 1 with $`b_{2k}(x_0,y_0)=c_0`$; $`c_0]c_m,0[`$ with the property that $`𝒪`$ disappears as $`c`$ decreases from $`c_0+ϵ`$ to $`c_0ϵ`$.
Denote $`S_{2k+1}^{}`$ the set of these critical points. The curve $`_{2k+1}`$ has $`\frac{k(3k1)}{2}`$ outer ovals and $`(2k+1)`$ connected components
$`M_0=\{(x,y)𝐑^2|b_{2k+1}(x,y)>0\}`$ which contain a segment of the line at infinity of $`𝐑P^2`$ on their boundary. Therefore, the number of critical points of $`S_{2k+1}^{}`$ is:
(2.4)
$$c_1^{}(B_{2k+1})\frac{k(3k1)}{2}+2k=\frac{k(3k+3)}{2}$$
2. Let $`c`$ increase from $`0`$ to $`c_M+ϵ`$, then the outer oval shrink and disappear. Moreover, from the Lemma 2.1.3 of Chapter 2.1.3, the set $`M_{c_M+ϵ}`$ has $`m`$ components. (\*)Let $`_{2k}`$ be the Harnack curve of even degree $`2k`$
The set $`M_0`$ has $`(2k(k1))`$ components intersecting the line at infinity; the curve $`_{2k}`$ has $`\frac{(k1)(k2)}{2}`$ inner ovals. It is easy to deduce from the Lemma 2.1.2 of Chapter 2.1.2 that when $`c`$ increases from $`0`$ to $`c_M+ϵ`$, the inner ovals expand, touch the non-empty outer oval or another inner oval and then disappear. In other words, when $`c`$ decreases from $`c_M+ϵ`$ to $`0`$, $`\frac{(k1)(k2)}{2}`$ inner ovals appear. Hence, from Lemma 2.1.2 of Chapter 2.1.2(2a), one can associate to each inner oval $`𝒪`$ a critical point $`(x_0,y_0)`$ of index 1 with $`b_{2k}(x_0,y_0)=c_0`$; $`c_0]0,c_M[`$ with the property that $`𝒪`$ disappears as $`c`$ increases from $`c_0ϵ`$ to $`c_0+ϵ`$.
Furthermore, the set $`M_{c_M+ϵ}`$ has $`2k`$ components. Thus, from the Lemma 2.1.3 of Chapter 2.1.3, $`c`$ has passed at least $`(k1)`$ critical values of index $`1`$ in his way from $`0`$ to $`c_M+ϵ`$. We denote by $`S_{2k}^{}`$ the set of critical points associated to these critical values.
Denote $`S_{2k}^+`$ the set of critical points of index 1 $`(x_0,y_0)`$ such that $`b_{2k}(x_0,y_0)=c_0`$; $`c_0]0,c_M[`$. Therefore, the number of critical points of $`S_{2k}^+`$ is:
(2.5)
$$c_1^+(B_{2k})\frac{(k1)(k2)}{2}+(k1)=\frac{k(k1)}{2}$$
(\*)Let $`_{2k+1}`$ be the Harnack curve of odd degree $`2k+1`$.
The set $`M_0`$ has $`2k+1`$ components intersecting the line at infinity; the curve $`_{2k+1}`$ has $`\frac{k(k1)}{2}`$ inner ovals. When $`c`$ increases from $`0`$ to $`c_M+ϵ`$; inner ovals expand, then touch the one-side component or another inner oval, and finally disappear.
Moreover, the set $`M_{c_M+ϵ}`$ has $`2k+1`$ components.
From Lemma 2.1.2 of Chapter 2.1.2 (3), one can associate to each inner oval $`𝒪`$ of $`_{2k+1}`$ a critical point $`(x_0,y_0)`$ of index 1 with $`b_{2k+1}(x_0,y_0)=c_0`$; $`c_0]0,c_M[`$ with the property that $`𝒪`$ disappears as $`c`$ increases from $`c_0ϵ`$ to $`c_0+ϵ`$.
Denote $`S_{2k+1}^{}`$ the set of these critical points. The curve $`_{2k+1}`$ has $`\frac{k(k1)}{2}`$ inner ovals. Hence, the number of critical points of $`S_{2k+1}^+`$ is:
(2.6)
$$c_1^+(B_{2k+1})\frac{k(k1)}{2}$$
Hence, according to Proposition 2.1.1 of Chapter 2.1.1, inequalities 2.3, 2.4, 2.5, 2.6 are equalities. This implies the Proposition 2.1.4 of Chapter 2.1.4. Q.E.D In particular, this method provides curves with real scheme:
for even $`m=2k`$
#### 2.2. Harnack Curves from a real viewpoint-Rigid Isotopy Classification-
Passing from polynomials to real set of points, it follows that real algebraic curves form a real projective space of dimension $`\frac{m(m+3)}{2}`$. We shall denote this space by the symbol $`𝐑𝒞_m`$ and by $`𝐑𝒟_m`$ the subset of $`𝐑𝒞_m`$ corresponding to real singular curves. We call a path in the complement $`𝐑𝒞_m\backslash 𝐑𝒟_m`$ of the discriminant hypersurface in $`𝐑𝒞_m`$ a rigid isotopy of real point set of nonsingular curves of degree $`m`$. The classification of real point set of curves of degree $`4`$ up to rigid isotopy is known since the $`19^{th}`$ century. It was completed for curves of degree $`5`$ and $`6`$ at the end of the seventies. Up to rigid isotopy a curve of degree $`4`$ is determined by its real scheme; up to rigid isotopy a curve of degree $`5`$ or $`6`$, is determined by its real scheme and its type. Let us recall that we denote by $`_m`$ and call Harnack curve any curve with real scheme:
\- for even $`m=2k`$
$$1\frac{(k1)(k2)}{2}\frac{3k(k1)}{2}$$
\- for odd $`m=2k+1`$
$$Jk(2k1)$$
The main result of this section is given in the Theorem 2.2.1 of Chapter 2.2.1 where we establish the rigid isotopy classification of real point set of Harnack curves $`𝐑_m`$. Harnack curves of degree $`i6`$, as any $`M`$-curve of degree $`6`$, are rigidly isotopic. In Theorem 2.2.1 of Chapter 2.2.1, we extend this property to Harnack curves of arbitrary degree. Precisely, we prove that isotopy also implies rigid isotopy for real point set of Harnack curves. This section until its end is devoted to the proof of this result.
###### Theorem 2.2.1.
Rigid Isotopy Classification Theorem
Harnack curve $`_m`$ of degree $`m`$ are rigidly isotopic.
The proof of Theorem 2.2.1 of Chapter 2.2.1 is based on a modification of Harnack polynomials. Let us call regular modification of a polynomial any modification on its coefficients with the property that real point set of curves of the modified polynomial and the initial polynomial are rigidly isotopic.
(In particular, given a polynomial of a smooth algebraic curve, the modification of its coefficients such that the modified polynomial is a regular polynomial (i.e none of its critical points lies on the line at infinity, any two different critical points have distinct critical values) is obviously regular.)
Let $`𝒜`$ and $``$, be two smooth curves such that the union $`𝒜`$ is a singular curve all of whose singular points are crossings. Denote by $``$ the set of curves which result from the classical deformation of $`𝒜`$. We shall say that a curve is deduced from deformation of $`𝒜`$ if it is rigidly isotopic to a curve $`𝒞`$ of the set $``$.
Any two Harnack curves $`_m`$ constructed from the Harnack’s method are rigidly isotopic.(We recall Harnack’s method and propose a proof of this statement in the Appendix) The proof of the Theorem 2.2.1 of Chapter 2.2.1 is based on a regular modification of Harnack polynomials. From properties of this regular modification, we shall deduce that any Harnack curve $`_m`$ is, up to rigid isotopy, constructed from the Harnack’s method. In this way, we get the rigid isotopy Theorem 2.2.1 of Chapter 2.2.1.
##### Harnack curve of type $`^0`$
Let us recall that we call polynomial of type $``$ a regular Harnack polynomial of degree $`m`$ giving a Harnack curve $`_m`$ which intersects the line at infinity $`𝐑L`$ of $`𝐑P^2`$ in $`m`$ real distinct points. We shall call polynomial of type $`^0`$ (relatively to $`L`$) a polynomial of type $``$ such that the line at infinity $`L`$ of $`𝐑P^2`$ intersects the Harnack curve $`_m`$ of degree $`m`$ in $`m`$ real distinct points which belong to the same connected component of $`𝐑_m`$: in case $`m=2k`$, the non-empty oval of $`_{2k}`$; in case $`m=2k+1`$, the odd component of $`_{2k+1}`$. The odd component of $`_{2k+1}`$ is divided into $`2k+1`$ arcs which delimit a region of which boundary contains a segment of the line at infinity. The line at infinity is chosen such that only one of this region contains ovals of $`_{2k+1}`$. We say that a curve is of type $`^0`$ if its polynomial is of type $`^0`$. In particular, Harnack curves constructed from the Harnack’s method (see Appendix) are curves of type $`^0`$.
The main result of this subsection is the Theorem 2.2.2 of Chapter 2.2.2 where we give the rigid isotopy classification of curves of type $`^0`$.
###### Theorem 2.2.2.
Harnack curve $`_m`$ of degree $`m`$ and type $`^0`$ are rigidly isotopic.
Let $`_m`$ be a Harnack curve with polynomial $`B_m(x_0,x_1,x_2)`$ of type $`^0`$. Denote by $`C_i(x_1,x_2)`$ the unique homogeneous polynomial of degree $`i`$ in the variables $`x_1,x_2`$ such that: $`B_m(x_0,x_1,x_2)=x_0^{m1}.B_1(x_0,x_1,x_2)+\mathrm{\Sigma }_{i=2}^mx_0^{mi}.C_i(x_1,x_2)`$.
Let $`b_m(x,y)`$ be the affine polynomial associated to $`B_m(x_0,x_1,x_2)`$,
$$B_m(x_0,x_1,x_2)=x_0^m.b_m(x_1/x_0,x_2/x_0)$$
Let $`b_i(x,y)`$ be the truncation of $`b_m(x,y)`$ on the monomials $`x^\alpha .y^\beta `$ with $`0\alpha +\beta i`$ and $`B_i(x_0,x_1,x_2)=x_0^i.b_i(x_1/x_0,x_2/x_0)`$ be the homogeneous polynomial associated to $`b_i`$. We shall denote by $`_i`$ the curve with polynomial $`B_i(x_0,x_1,x_2)`$. Consider the norm in the vector space of polynomials
$$a_{i,j}x_1^ix_2^j=max\{|a_{i,j}||(i,j)𝐍^2\}$$
(Given $`A_m(x_0,x_1,x_2)`$ an homogeneous polynomial,
$$A_m(x_0,x_1,x_2)=A_m(1,x_1,x_2)$$
)
The Theorem 2.2.2 of Chapter 2.2.2 may be also formulated as follows:
###### Proposition 2.2.3.
On the assumption that $`B_m(x_0,x_1,x_2)`$ is of type $`^0`$ relatively to the line at infinity $`x_0=0`$, up to regular modification of $`B_m(x_0,x_1,x_2)`$, for any $`B_i(x_0,x_1,x_2)`$, $`i1`$, we have:
1. $`B_i(x_0,x_1,x_2)`$ is smooth of type $`^0`$
2. none of the critical points of $`B_i(x_0,x_1,x_2)`$ belongs to the line at infinity;
for any critical point $`(1,x_{0,1},x_{0,2})𝐑P^2`$ of $`B_i(1,x_1,x_2)=b_i(x_1,x_2)`$ its representative in $`S^2`$ $`\frac{1}{(1+x_{0,1}^2+x_{0,2}^2)^{1/2}}(1,x_{0,1},x_{0,2})`$ is such that:
$`B_i(\frac{1}{(1+x_{0,1}^2+x_{0,2}^2)^{1/2}}(1,x_{0,1},x_{0,2}))[\mathrm{\Sigma }_{j=i+1}^mC_j,\mathrm{\Sigma }_{j=i+1}^mC_j]`$
On the assumption that $`B_m`$ is of type $`^0`$, according to proposition 2.1.1 of Chapter 2.1.1, the following equalities are verified:
for even $`m=2k`$: $`c(B_{2k})=(\frac{(k1)(k2)}{2},k.(2k1),\frac{3k.(k1)}{2})`$
and for odd $`m=2k+1`$: $`c(B_{2k+1})=(\frac{k.(k1)}{2},k.(2k+1),\frac{k.(3k1)}{2})`$
proof:
Our proof is based on Morse Lemma and Petrovskii’s theory. We shall proceed by descending induction on the degree $`i`$ of $`B_i`$. Let us assume that assumptions (1) and (2) of Proposition 2.2.3 of Chapter 2.2.3 are satisfied for $`B_n`$ with $`i+1nm`$, and prove that $`B_i`$ also satisfies (1) and (2).
On these assumptions, we shall prove that, the Harnack curve $`_{i+1}`$ with polynomial $`B_{i+1}`$ is deduced from deformation of $`_iL`$ ,$`B_{i+1}(x_0,x_1,x_2)=x_0.B_i(x_0,x_1,x_2)+C_{i+1}(x_1,x_2)`$, where $`L:=\{x_0=0\}`$ and $`B_i`$ is of type $`^0`$ relatively to $`L`$. I) Let us in a first part study curves $`_i`$ of degree $`4`$. Up to regular modification of the Harnack polynomial $`B_{i+1}`$ one can always assume that $`B_i`$ is smooth.
Introduction Consider a real projective line $`𝐑L𝐑P^2`$; its tubular neighborhood in $`𝐑P^2`$ is homeomorphic to a Möbius band. The core of the Möbius band, i.e the real projective line $`𝐑L`$, is the circle with framing $`\pm 1`$. In such a way, one can identify points of $`𝐑L`$ with points of two halves (oriented) circles. Halves intersect each other in two points we shall call extremities. We shall denote $`𝐑L^+`$ and $`𝐑L^{}`$ the halves of $`𝐑L`$. On each half $`𝐑L^\pm `$ of $`𝐑L`$ one can consider an oriented tubular fibration. We shall denote by $`^\pm `$ the half of the Möbius band which is the tubular neighborhood of $`L^\pm `$. In such a way, the boundary of the half of Möbius band $`^\pm `$ contains the union of two real projective lines $`𝐑L_1^\pm `$, $`𝐑L_2^\pm `$. The half $`𝐑L_i^\pm `$ of $`𝐑L_i`$, $`i\{1,2\}`$, is the image of a smooth section of the tubular oriented fibration of a tubular neighborhood of $`𝐑L^\pm `$; $`^{}^+=`$, $`𝐑L_1^\pm 𝐑L_2^\pm `$. Let $`_m`$ be the Harnack curve of degree $`m`$ with polynomial $`B_m`$ of type $`^0`$ relatively to $`L`$. We shall study $`𝐑_m`$ in a tubular neighborhood $``$ of the line at infinity $`𝐑L`$.
Let $`𝒜`$ be the set of arcs of $`𝐑_m`$ bounding regions which contain a segment of the line at infinity on their boundary and do not contain ovals of $`𝐑_m`$. We shall say that a critical point of index $`1`$ is associated to an arc $`\gamma 𝒜`$ of $`𝐑_m`$ if there exists $`p`$ a critical point of $`b_m`$ with critical value $`c_0`$ such that as $`c`$ varies from $`0`$ to $`c_0`$, the region which contains $`\gamma `$ on its boundary varies in such a way that for $`c=c_0`$ it touches an other arc $`\gamma ^{}`$. The arc $`\gamma ^{}`$ is said associated to the arc $`\gamma `$ and the pair $`(\gamma ,\gamma ^{})`$ is said associated to the critical point $`p`$. We shall denote $`𝒜^{}`$ the set of arcs $`\gamma ^{}`$. From the study of the Petrovskii’s pencil $`x_0^m.(b_m(x_1/x_0,x_2/x_0)c)`$, $`c𝐑`$ over $`_m`$ with $`L:=x_0=0`$, we shall define a set $`𝒫`$ of critical points of $`b_m`$ and a set of arcs of $`𝐑_m`$ associated to $`𝒜`$.
The set $`𝒜`$ is the image of a smooth section $`s`$ of a tubular fibration of $`𝐑L^+`$ minus 2 points. The proof of Theorem 2.2.2 of Chapter 2.2.2 is based on the characterization (up to regular deformation of $`B_m`$) of a subset of $`𝐑_m`$ which is the image of the extension of $`s`$ to a smooth section of a tubular fibration of $`𝐑L`$ minus a finite number of points.
We shall proceed as follows. Using Morse Lemma and Petrovskii’s theory, we shall define a set of arcs $`𝒟`$ of $`𝐑_m`$ with the property that up to regular modification of $`B_m`$ the line at infinity $`𝐑L`$ divides any arc $`\xi 𝒟`$ into two halves which belong respectively to $`𝐑_{m1}`$ and $`𝐑_m\backslash 𝐑_{m1}`$ where $`_{m1}`$ is the Harnack curve of degree $`m1`$. In such a way, we shall deduce that the Harnack curve $`_m`$ is deduced from deformation of $`_{m1}L`$ $`B_m(x_0,x_1,x_2)=x_0.B_{m1}(x_0,x_1,x_2)+C_m(x_1,x_2)`$, where $`L:=\{x_0=0\}`$ and $`B_{m1}`$ is of type $`^0`$.
Let us distinguish in parts I.i) and I.ii) curves of even and odd degree.
I.i) -Harnack curves of even degree-
Let $`_{2k}`$ be a Harnack curve of degree $`2k`$ with polynomial $`B_{2k}(x_0,x_1,x_2)`$ of type $`^0`$ relatively to the line at infinity $`L:=x_0=0`$. We shall prove the following Proposition 2.2.4 of Chapter 2.2.4.
###### Proposition 2.2.4.
Let $`_{2k}`$ be a Harnack curve of degree $`2k`$ with polynomial $`B_{2k}(x_0,x_1,x_2)`$ of type $`^0`$ relatively to the line at infinity $`L:=x_0=0`$. Up to regular modification of $`B_{2k}`$, $`B_{2k1}(x_0,x_1,x_2)`$ is also of type $`^0`$ relatively to $`L`$.
proof:
The part I.i) until its end is devoted to the proof of Proposition 2.2.4 of Chapter 2.2.4.
The non-empty oval of $`_{2k}`$ is divided into $`2k`$ arcs which delimit a region which contains a segment of the line at infinity on its boundary. Denote $`𝒜`$ the set of the $`2k1`$ arcs which delimit a region which does not contain oval of $`𝐑_{2k}`$; $`k1`$ (resp, $`k`$) of them delimit regions $`\{x𝐑^2|b_{2k}(x)<0\}`$ (resp, $`\{x𝐑^2|b_{2k}(x)>0\}`$). The union of arcs of $`𝒜`$ is a connected orientable part of the non-empty oval which belongs to a tubular neighborhood of $`L`$. The connected surface (with boundary the non-empty oval) obtained from removing the interior of the inner ovals to the interior of the non-empty oval is also orientable. Thus, the orientation of the connected union of the $`2k1`$ arcs of $`𝒜`$ is induced by a half of $`𝐑L`$. We shall denote $`𝐑L^+`$ this half.
The tubular neighborhood of the line $`𝐑L`$ is a Möbius band $`=^{}^+`$ $`^\pm 𝐑L_1^\pm 𝐑L_2^\pm `$. The set $`𝒜`$ is the image of a smooth section of a tubular fibration of $`𝐑L^+`$ (-precisely, $`𝐑L^+`$ minus its extremities, these 2 points belong to the arc of $`𝐑_{2k}`$ bounding the non-empty region with a segment of the line at infinity on its boundary-)
It is easy to see that there exists an isotopy of $`𝐑P^2`$ which pushes the $`k`$ (oriented) arcs of $`𝒜𝐑^2`$ bounding positive regions $`\{x𝐑^2|b_{2k}(x)>0\}`$ to one positive line $`𝐑L_i^+`$, $`i\{1,2\}`$. Let it be $`𝐑L_1^+`$. In this way, there exists an isotopy of $`𝐑P^2`$ which pushes the $`k1`$ (oriented) arcs of $`𝒜𝐑^2`$ bounding negative regions $`\{x𝐑^2|b_{2k}(x)<0\}`$ to $`𝐑L_2^+`$. Let us prove that one can associate to the set of $`2k1`$ arcs $`𝒜`$ a set of critical points $`𝒫`$ of index $`1`$ of $`b_{2k}`$ and a set of arcs $`𝒜^{}`$ of $`𝐑_{2k}`$.
To this end, we shall study the pencil of curves $`x_0^{2k}.(b_{2k}(x_1/x_0,x_2/x_0)c)`$, $`c𝐑`$, over $`_{2k}`$.
Let us recall that (see Petrovskii’s Lemmas 2.1.2 of Chapter 2.1.2 2.1.3 of Chapter 2.1.3) as $`c`$ decreases from $`0`$ positive regions and positive ovals expand. As $`c`$ decreases from $`0`$ the number of regions $`G_i`$ of $`\{x𝐑^2|b_{2k}(x)>0\}`$ which contain a segment of the line at infinity of $`𝐑P^2`$ on their boundary decreases from $`k+1`$ to $`1`$.
As $`c`$ increases from $`0`$ the number of regions $`G_i`$ of $`\{x𝐑^2|b_{2k}(x)>0\}`$ which contain a segment of the line at infinity of $`𝐑P^2`$ on their boundary increases from $`k+1`$ to $`2k`$. Hence, to each arc of $`𝒜`$ is associated a critical of index $`1`$.
-Let us study the set $`𝒫`$ of critical point $`1`$ associated to the set $`𝒜`$ of arcs of $`𝐑_{2k}`$.
Recall that we consider on $`𝐑P^2`$ the Fubini-Study metric induced by the projection $`\pi _𝐑:S^2𝐑P^2`$.
Without loss of generality, one can assume that critical points of $`B_{2k}`$ and points of the intersection of $`_{2k}`$ with $`L`$ do not belong to the line $`x_2=0`$. On such assumption, consider the function $`b_{2k}(0,x_1/x_2,1)=b_{2k}(y)`$
$$\frac{}{x_1}b_{2k}(0,x_1/x_2,1)=\frac{}{y}b_{2k}(y)\frac{1}{x_2}$$
$$\frac{}{x_2}b_{2k}(0,x_1/x_2,1)=\frac{}{y}b_{2k}(y)\frac{x_1}{(x_2)^2}$$
Critical points of the function $`b_{2k}(0,x_1/x_2,1)`$ may be defined from critical points $`b_{2k}(y)`$. (Obviously, for any such critical point of $`b_{2k}(0,x_1/x_2,1)`$ the following equality is verified $`(0:x_1/x_2:1)=(0:x_1:x_2)=(0:x_1:x_2)`$.) The function $`b_{2k}(y)`$ has exactly $`2k`$ zeroes which coincide with intersection points of $`_{2k}`$ with the line $`L`$. Hence, by Rolle’s Theorem $`2k1`$ extrema at which $`b_{2k}^{}(y)`$ must change from positive to negative. By continuity of $`b_{2k}`$ and $`b_{2k}^{}`$, it follows the alternation of sign of $`b_{2k}^{}(x_0,x_1,x_2)`$ in an $`ϵ`$-tubular neighborhood $`_ϵ`$ of the line $`x_0=0`$ in $`𝐂P^2`$ and thereby the existence of a set of $`2k1`$ critical points of $`B_{2k}`$ (i.e critical points of the affine polynomial $`b_{2k}(x_1/x_0,x_2/x_0)`$) in $`_ϵ`$. According to Petrovskii’s theory, to describe critical points of $`b_{2k}(x_1/x_0,x_2/x_0)`$ by means of the pencil $`x_0^{2k}.(b_{2k}(x_1/x_0,x_2/x_0)c)`$ one can assume that $`B_{2k}`$ is a regular polynomial. (i.e none of its critical points belongs to $`x_0=0`$, any two of them have different critical values.) Up to slightly modify coefficients of $`b_{2k}`$, any critical point $`y=(0,x_1/x_2,1)`$ of $`b_{2k}(y)`$ gives rise to one critical point which may be chosen among two points $`(ϵ:\pm x_{0,1}:\pm x_{0,2})𝐑P^2`$.
Let us prove in Lemma 2.2.5 of Chapter 2.2.5 that without loss of generality one can assume that the $`2k1`$ points of $`𝒫`$ belong to a line $`L_{ϵ_1}`$ of $`_ϵ`$ the $`ϵ`$-tubular neighborhood of $`L:=x_0=0`$ where $`ϵ>0`$ is arbitrarily small.
###### Lemma 2.2.5.
Let $`L`$ be the real projective line at infinity. Denote by $`_ϵ`$ the $`ϵ`$-tubular neighborhood of $`L`$ in $`𝐑P^2`$.
Given $`B_{2k}`$ a Harnack polynomial of type $`^0`$ (relatively to $`L`$). There exists a regular modification of $`B_{2k}`$ such that:
1. the modified polynomial is of type $`^0`$ (relatively to $`L`$).
2. the set $`𝒫`$ is a set of $`(2k1)`$ points of a line $`L_{ϵ_1}_ϵ`$ and $`ϵ`$ is arbitrarily small.
proof:
Recall that any point $`p𝒫`$ is a critical point of index $`1`$ associated to one of the $`2k1`$ arcs of the set $`𝒜`$ (-the set of arcs of $`𝐑_{2k}`$ which delimit a region which contains a segment of the line at infinity on its boundary and does not contain an oval of $`𝐑_{2k}`$-)
-It is easy to get that we can regularly modify $`B_{2k}`$ such that any $`p𝒫`$ belongs to $`_ϵ`$ where $`ϵ>0`$ is arbitrarily small.
-Assume that any point $`p𝒫`$ belongs to $`_ϵ`$ the $`ϵ`$-tubular neighborhood of $`L`$ with $`ϵ`$ is arbitrarily small. The polynomial $`B_{2k}(x_0,x_1,x_2)`$ is of type $`^0`$ relatively to the line $`L`$.
Let us now prove that one can regularly modify $`B_{2k}`$ such that any $`p𝒫`$ belongs to a real projective line $`L_{ϵ_1}_ϵ`$.
Let $`L_{ϵ_1}`$ be a real projective line of $`_ϵ`$. Denote by $`x_0^{}=x_0+\alpha .x_1+\beta .x_2`$ the polynomial of $`L_{ϵ_1}`$. Then, consider the linear change of projective coordinates mapping $`(x_0:x_1:x_2)`$ to $`(x_0^{}:x_1:x_2)`$. Such transformation carries $`B_{2k}(x_0,x_1,x_2)=x_0.B_{2k1}(x_0,x_1,x_2)+C_{2k}(x_1,x_2)`$ to
$`B_{2k}^{}(x_0^{},x_1,x_2)=x_0^{}.B_{2k1}^{}(x_0^{},x_1,x_2)+C_{2k}^{}(x_1,x_2)`$ where
$`B_{2k1}^{}(x_0^{},x_1,x_2)=x_0^{}.B_{2k2}^{}(x_0^{},x_1,x_2)+C_{2k1}^{}(x_1,x_2)`$. We shall prove that, up to regular modification of $`B_{2k}^{}`$ and thus of $`B_{2k}`$, one can choose $`C_{2k1}^{}(x_1,x_2)`$ in such a way that the $`2k1`$ points of $`𝒫`$ belong to $`L_{ϵ_1}`$. Using the linear change of projective coordinates mapping $`(x_0^{}:x_1:x_2)`$ to $`(x_0:x_1:x_2)`$, we shall get the polynomial $`B_{2k}(x_0,x_1,x_2)`$.
Let us detail this construction.
Given $`B_{2k}(x_0,x_1,x_2)`$ of type $`^0`$, the polynomial $`C_{2k}(x_1,x_2)`$ such that $`B_{2k}(x_0,x_1,x_2)=x_0.B_{2k1}(x_0,x_1,x_2)+C_{2k}(x_1,x_2)`$ has $`2k`$ distinct roots on the line at infinity $`L:=x_0=0`$. Thus, $`C_{2k}(x_1,x_2)=\mathrm{\Sigma }_{i=0}^{2k}a_ix_1^{2ki}x_2^i`$ with $`a_i0`$.
It follows from the previous local study around $`p𝒫`$ that, up to regular modification, one can assume the polynomial $`B_{2k}`$ (of type $`^0`$ relatively to $`L`$) such that any point $`p`$ of $`𝒫`$ belongs to the $`ϵ`$-tubular neighborhood $`_ϵ`$ of $`L`$ where $`ϵ>0`$ is arbitrarily small. Denote by $`B_{2k,0}`$ such a polynomial and by $`𝒫_0`$ its respective set $`𝒫`$ of critical points.
Set $`x_0=ϵ>0`$.
$`B_{2k,0}(ϵ,x_1,x_2)=ϵ.B_{2k1,0}(ϵ,x_1,x_2)+C_{2k,0}(x_1,x_2)`$
$`=ϵ.(ϵ.B_{2k2,0}(ϵ,x_1,x_2)+C_{2k1,0}(x_1,x_2))+C_{2k,0}(x_1,x_2)`$ with $`B_i=x_0^i.b_i`$ where $`b_i`$ is the truncation of $`b_{2k,0}=B_{2k,0}(1,x_1,x_2)`$ monomials of $`x_1^a.x_2^b`$ of degree $`a+bi`$. Local coordinates defined in a neighborhood $`U(p)`$ of a point $`p`$ depend principally on the first derivative and the second derivative of the function $`B_{2k}`$ around $`p`$. Hence, to describe the curve in neighborhood of points $`𝒫`$, it is sufficient to consider the truncation $`D_{2k,0}`$ of $`B_{2k,0}`$ on monomials $`x_1^a.x_2^b`$ of degree $`2k2a+b2k`$. Therefore, to define a regular modification $`B_{2k,t}`$, $`t[0,1]`$, of $`B_{2k,0}`$ of type $`^0`$ with $`𝒫_0_ϵ`$ to $`B_{2k,1}`$ of type $`^0`$ with $`𝒫L_{ϵ_1}`$ it sufficient to consider the truncation $`D_{2k,0}`$ of $`B_{2k,0}`$.
Without loss of generality, one can assume that the $`2k1`$ points $`p𝒫_0`$ do not belong to $`x_2=0`$. In this way, for any $`p𝒫_0`$, $`\frac{B_{2k,0}}{x_0}(p)=0`$ and $`\frac{B_{2k,0}}{x_1}(p)=0`$. Moreover, $`B_{2k,0}`$ has $`2k`$ roots on the line $`x_0=0`$ and none of these points is a point of $`𝒫`$. Therefore, it follows that the set $`𝒫_0_ϵ`$ may be chosen such that the truncation $`D_{2k,0}(x_0,x_1,x_2)`$ of $`B_{2k,0}(x_0,x_1,x_2)`$ is of the form $`D_{2k,0}=\mathrm{\Sigma }_{2k2i+j2k}a_{i,j}x_0^{2k(i+j)}x_1^ix_2^j`$ with $`a_{i,j}0`$.
Set $`B_{2k}^{}=B_{2k,0}^{}=x_0^{}.B_{2k1}^{}+C_{2k}^{}`$. We shall prove that there exists a regular modification of $`B_{2k,0}^{}`$ (that is of $`B_{2k,0}`$) which pushes any point $`p𝒫_0`$ to a point of the set $`\stackrel{~}{𝒫}_0:=B_{2k1,1}^{}=0L_{ϵ_1}`$.
In such a way, any point $`p\stackrel{~}{𝒫}_0`$ verifies
(2.7)
$$\frac{B_{2k,1}^{}}{x_0^{}}(p)=0$$
Without loss of generality, on can assume that any point $`p\stackrel{~}{𝒫}_0`$ does not belong to $`x_2=0`$.
Thus, any point $`p\stackrel{~}{𝒫}_0`$ also verifies:
(2.8)
$$\frac{B_{2k,1}^{}}{x_1}(p)=0$$
Let us prove that the modification $`B_{2k,t}`$ $`t[0,1]`$, $`B_{2k,0}=B_{2k}`$ such that $`B_{2k,1}`$ verifies properties (2.7) and (2.8) is regular. We shall proceed by induction.
Denote $`p_i,1i2k1`$ the set of $`2k1`$ points of $`𝒫_0`$. Denote $`L_0^{i,j}`$ the unique line such that $`p_i,p_jL_0^{i,j}`$ Setting $`L_0^{i,j}:=x_0^{}=0`$, properties (2.7) and (2.8) are verified for $`p_i,p_j`$. Choose $`p_1,p_2`$ two points of $`𝒫_0`$ associated to respectively $`\gamma _1`$ $`\gamma _2`$ arcs of $`𝒜`$ which intersect each other in one point of the line at infinity $`L`$.
Let $`p_3𝒫_0`$ associated to $`\gamma _3𝒜`$, such that $`\gamma _3`$ and $`\gamma _2`$ intersect each other in one point.
Setting $`L_{ϵ_1}=L_0^{1,2}:=x_0^{}=0`$, $`L_0^{2,3}:=x\mathrm{"}_0=x_0^{}+\alpha ^{}.x_1+\beta ^{}.x_2=0`$. The linear change of projective coordinates mapping $`(x_0^{}:x_1:x_2)`$ to $`(x\mathrm{"}_0:x_1:x_2)`$ carries $`B_{2k}^{}(x_0^{},x_1,x_2)=x_0^{}.B_{2k1}(x_0^{},x_1,x_2)+C_{2k}(x_1,x_2)`$ to
$`B\mathrm{"}_{2k}(x\mathrm{"}_0,x_1,x_2)=x\mathrm{"}_0.B\mathrm{"}_{2k1}(x\mathrm{"}_0,x_1,x_2)+C\mathrm{"}_{2k}(x_1,x_2)`$
Moreover, $`\frac{B\mathrm{"}_{2k,0}}{x\mathrm{"}_0}(p_3)=0`$, $`\frac{B\mathrm{"}_{2k,0}}{x_1}(p_3)=0`$
One can modify regularly coefficients of $`B_{2k,0}`$ such that $`p_3L_{ϵ_1}=L_0^{1,2}`$. (i.e there exists regular modification on gradient trajectories such that $`p_3`$ is pushed to a point of the line $`L_{ϵ_1}=L_0^{1,2}`$) This modification may be expressed from the path $`x\mathrm{"}_{0,t}=x_0^{}+(\alpha ^{}t.\alpha ^{}).x_1+(\beta ^{}t.\beta ^{}).x_2`$, $`t[0,1]`$ from $`L_0^{2,3}`$ to $`L_{ϵ_1}`$ and polynomials $`B\mathrm{"}_{2k1,t}(x\mathrm{"}_{0,t},x_1,x_2)`$ and $`C\mathrm{"}_{2k,t}(x_1,x_2)`$ $`B\mathrm{"}_{2k,t}=x\mathrm{"}_{0,t}.B\mathrm{"}_{2k1,t}(x\mathrm{"}_{0,t},x_1,x_2)+C\mathrm{"}_{2k,t}(x_1,x_2)`$. The deformation moves $`p_{3,t}L_t`$ in such a way that $`\frac{B\mathrm{"}_{2k,t}}{x\mathrm{"}_{0,t}}(p_{3,t})`$ and $`\frac{B\mathrm{"}_{2k,t}}{x_1}(p_{3,t})`$. Note that $`B\mathrm{"}_{2k,0}=B\mathrm{"}_{2k}`$,$`p_{3,0}=p_3L_0^{2,3}`$, $`p_{3,t}L_t`$ where $`x\mathrm{"}_{0,t}`$ is the polynomial of $`L_t`$ and $`L_1=L_{ϵ_1}`$) ($`B\mathrm{"}_{2k,1}=x\mathrm{"}_{0,1}.B\mathrm{"}_{2k1,1}(x\mathrm{"}_{0,1},x_1,x_2)+C\mathrm{"}_{2k,1}(x_1,x_2)`$; $`B\mathrm{"}_{2k,1}=x_0^{}.B\mathrm{"}_{2k1,1}(x_0^{},x_1,x_2)+C_{2k,1}^{}(x_1,x_2)x_0^{}.B_{2k1}(x_0^{},x_1,x_2)+C_{2k}(x_1,x_2)`$.) Since $`p_1,p_2,p_3_ϵ`$, where $`ϵ>0`$ is arbitrarily small; $`|\alpha ^{}|`$ and $`|\beta ^{}|`$ are also arbitrarily small. In this way, (see Lemma 1 and Lemma 2 of ), the deformation $`B\mathrm{"}_{2k,t}=x\mathrm{"}_{0,t}.B\mathrm{"}_{2k1,t}(x\mathrm{"}_{0,t},x_1,x_2)+C\mathrm{"}_{2k,t}(x_1,x_2)`$ such that $`p_{3,t}`$ is a root of $`B\mathrm{"}_{2k1,t}(0,x_1,x_2)`$ (i.e of $`\frac{B\mathrm{"}_{2k,t}}{x\mathrm{"}_{0,t}}`$) and of $`\frac{B\mathrm{"}_{2k,t}}{x_1}`$, $`t[0,1]`$ is regular.
Iterating this argument for $`p_i`$, $`4i2k1`$, we get the announced regular deformation $`B_{2k,t}`$, $`t[0,1]`$, of $`B_{2k}`$. Any $`p_i𝒫_0`$ is pushed to a root of $`C_{2k1,1}^{}(x_1,1)`$ (that is a root of $`\frac{B_{2k,1}^{}}{x_0^{}}`$) which is also a root of $`\frac{C_{2k,1}^{}}{x_1}(x_1,1)`$ (that is a root of $`\frac{B_{2k,1}^{}}{x_1}`$) For any $`p_i`$, $`1i2k1`$, of $`\stackrel{~}{𝒫}_0:=B_{2k1,1}^{}=0L_{ϵ_1}`$ consider its representative $`(0:x_i:1)`$.
Let $`s_j`$, $`j\{1,\mathrm{},2k1\}`$ be the elementary roots symmetric polynomial of degree $`j`$ on the $`p_i`$.
The polynomial $`C_{2k,1}^{}(x_1,x_2)`$ of $`B_{2k,1}^{}(x_0^{},x_1,x_2)=x_0^{}.B_{2k1,1}^{}(x_0^{},x_1,x_2)+C_{2k,1}^{}(x_1,x_2)`$ is (up to multiplication by a constant) such that:
$$\frac{B_{2k,1}^{}}{x_1}(0,x_1,x_2)=x_1^{2k1}s_1x_1^{2k2}x_2+\mathrm{}(1)^{2k1}s_{2k1}x_2^{2k1}$$
Bringing together the resulting equalities, it follows the polynomial $`B_{2k,1}^{}(x_0^{},x_1,x_2)`$ and thus the definition of $`B_{2k,1}(x_0,x_1,x_2)`$ with $`𝒫=\stackrel{~}{𝒫}_0L_{ϵ_1}`$. Since $`L_{ϵ_1}_ϵ`$, with $`ϵ>0`$ arbitrarily small, the modification $`B_{2k,t}`$, $`t[0,1]`$, of $`B_{2k,0}`$ of type $`^0`$ with $`𝒫_0_ϵ`$ is such that $`B_{2k,1}`$ with $`𝒫_1L_{ϵ_1}`$ is also of type $`^0`$.
###### Remark 2.2.6.
Note that the rigid proposed above requires a balance of critical value. Hence, it does not push arcs $`𝒜`$ and $`𝒜^{}`$ to $`_ϵ`$. Indeed, consider the Fubini-Study metric on $`S^2`$ and in such a way consider $`𝐑_{2k}S^2`$. Any point $`p𝒫`$ has representative in $`S^2`$ $`(x_0,x_1,x_2)=(ϵ,p_1,p_2)`$ with $`ϵ<<1`$.
$$B_{2k}(x_0,x_1,x_2)=x_0^{2k2}B_{2k2}(x_0,x_1,x_2)+D_{2k}$$
Hence, since $`ϵ^{2k2}.B_{2k2}<<1`$ for $`ϵ<<1`$; In a neighborhood of $`p`$,
$$B_{2k}(ϵ,p_1,p_2)=ϵ^{2k2}B_{2k2}(ϵ,p_1,p_2)+D_{2k}(ϵ,p_1,p_2)D_{2k}(ϵ,p_1,p_2)$$
It follows that the ”depth of waves” $`𝒜𝒜^{}`$ depends principally on coefficients of monomials $`x_1^{2ki}x_2^i`$ of $`C_{2k}`$; $`B_{2k}=x_0.B_{2k1}+C_{2k}`$)
It concludes the proof of Lemma 2.2.5 of Chapter 2.2.5. Q.E.D
-Let us describe the set $`𝒜^{}`$ associated to $`𝒫`$ and $`𝒜`$. To this end, we shall study the pencil $`x_0^{2k}.(b_{2k}(x_0,x_1,x_2)c)`$ $`c𝐑`$.
Let us recall, for sake of clarity, some properties of this pencil (see Lemma 2.1.3 of Chapter 2.1.3). Let $`c_M`$ be the maximal critical value, let $`c_m`$ be the minimal critical value. Denote by $`M_c=\{x𝐑^2|b_{2k}>0\}`$ For $`c>c_M`$, $`M_c`$ has $`2k`$ connected components. For $`c<c_m`$, $`M_c`$ is connected.
Two components of the connected components can coalescence only when $`b_{2k}`$ passes a critical value $`1`$. Bringing together this statement with Lemma 2.1.2 of Chapter 2.1.2, and the fact that any curve of the pencil over $`_{2k}`$ intersects the line $`x_0=0`$ in $`2k`$ points, it follows easily that an arc of $`𝒜`$ bounding a positive region may not touch an arc of $`𝒜`$ bounding a negative region as $`c`$ varies from $`0`$.
Let us prove the following Lemma.
###### Lemma 2.2.7.
As $`c`$ varies from $`0`$ to $`+\mathrm{}`$ or $`\mathrm{}`$ any two arcs $`\gamma ,\gamma ^{}`$ of $`𝒜`$ may not be deformed in such a way that for $`c=c_0`$ they touch each other.
proof: Let $`_{2k}`$ be a curve of type $`^0`$ relatively to $`L`$ and $`B_{2k}`$ its polynomial. The proof of Lemma 2.2.7 of Chapter 2.2.7 is based on a study of $`𝐑_{2k}`$ on a tubular neighborhood $`_ϵ`$ of $`𝐑L`$ and in particular the two following facts.
According to Lemma 2.2.5 of Chapter 2.2.5, given a curve $`_{2k}`$ of type $`^0`$ relatively to $`L`$, one can push by a rigid isotopy the set $`𝒫`$ associated to $`𝒜`$ to a line of the tubular neighborhood $`_ϵ`$ of $`𝐑L`$ where $`ϵ`$ positive is arbitrarily small.
According to the Petrovskii’s theory, , up to slightly modify coefficients of $`B_{2k}`$ without changing topological structure of $`_{2k}`$ one can assume that none of its critical point belongs to the line at infinity. Choose $`_ϵ`$ the tubular neighborhood of $`𝐑L`$ such that any line of$`_ϵ`$ intersects the non-empty oval of $`_{2k}`$ in $`2k`$ points. According to proof of Lemma 2.2.5 of Chapter 2.2.5, one can assume that critical points of $`B_{2k}`$ of the set $`𝒫`$ belong to $`L_{ϵ_1}`$ where $`𝐑L_{ϵ_1}_ϵ`$ and $`ϵ>0`$ is arbitrarily small. In this way, as $`c`$ increases two arcs $`\gamma ,\gamma ^{}`$ of $`𝒜`$ bounding positive regions are deformed and for $`c=c_0`$ they touch each other in a point of the line $`𝐑L_{ϵ_1}_ϵ`$. The connected set $`𝒜`$ is such that these two arcs are separated from each other by an arc of $`𝒜`$ which bounds a negative region. As $`c`$ increases from $`c_0`$ to $`c_0+ϵ`$, a negative oval $`𝒪`$ appears. The oval $`𝒪`$ does not intersect $`𝐑L_{ϵ_1}`$, it intersects the line at infinity $`𝐑L`$ in two points. Up to slightly modify coefficients of $`B_{2k}`$ without changing topological structure of $`_{2k}`$ one can consider $`𝐑L_{ϵ_1}`$ as the line at infinity of $`𝐑P^2`$. The creation of this negative oval leads to contradiction with the number of critical points of $`B_{2k}`$ and also with Petrovskii’s Lemma 2.1.3 of Chapter 2.1.3. Indeed, as $`c`$ grows, this negative oval shrinks in point and then disappears; the number of intersection points of $`_{2k}`$ with the line $`L_{ϵ_1}`$ decreases from $`2k`$ to $`2k2`$.
An argumentation analogous to the previous one may be used to prove that as $`c`$ increases from $`0`$ to $`+\mathrm{}`$, two arcs bounding a negative region may not be deformed and touch each other.
Moreover, when consider the Petrovskii’s pencil over $`b_{2k}`$, it follows easily that one arc bounding a positive region may not touch an arc bounding a negative region as $`c`$ varies from $`0`$.
This proves Lemma 2.2.7 of Chapter 2.2.7.
Q.E.D
We may describe, according to Petrovskii’s Lemmas 2.1.2 of Chapter 2.1.2 and 2.1.3 of Chapter 2.1.3, critical points of index $`1`$ associated to arcs of $`𝒜`$ as follows. There exist $`2k3`$ critical points with critical value $`c_0>0`$ such that for $`c=c_0`$, an arc $`\gamma `$ of $`𝒜`$ touches an oval.-(Each of the $`(k1)`$ arcs bounding negative region touches a positive oval, the $`(k2)`$ arcs which bound positive region touch negative ovals)-
There exist $`2`$ arcs of $`𝒜`$ bounding a positive region $`\{x𝐑^2|b_{2k}(x)>0\}`$ for which there exists $`c_0<0`$ with the following property. As $`c`$ decreases from $`0`$ to $`c_0`$, the region expands and for $`c=c_0`$ the arc touches a positive oval. (Note that, at this time of our proof, any positive oval under consideration above may be the non-empty one). In such a way, according to Morse Lemma since any point $`p𝒫`$ is a critical point of index $`1`$ of $`B_{2k}(x_0,x_1,x_2)`$, any arc $`\gamma `$ of $`𝒜`$ is associated to an arc $`\gamma ^{}`$ of $`𝐑_{2k}`$. We shall denote by $`𝒜^{}`$ the set of arcs $`\gamma ^{}`$. Denote by $``$ a tubular neighborhood of the line at infinity. As already noticed, arcs $`𝒜𝐑^2`$ may be seen (up to isotopy) as arcs of $`𝐑L_1^+`$ and $`𝐑L_2^+`$. In the same way, arcs $`𝐑_{2k}`$ of $`𝒜^{}`$ may be pushed on the boundary and the inside of $``$. Any oriented arc of the $`k1`$ arcs of $`𝐑L_2^+`$ is associated (to a critical point with positive critical value) to an arc of $`𝐑L_1^+`$, any oriented arc of the $`k3`$ arcs of $`𝐑L_1^+`$ associated to a critical point with positive critical value is associated to an arc of $`𝐑L_2^+`$. (According to Petrovskii’s theory, as $`c`$ increases, these arcs recede from $`𝐑L`$. ) These arcs are associated to a critical point with positive critical value. The two other arcs of $`𝐑L_1^+`$ are associated to a critical point with negative critical value and to an arc lying in the inside of $`^{}`$. These two arcs contain an extremity of $`𝐑L^+`$. (Let $`(\gamma ,\gamma ^{})(𝒜,𝒜^{})`$ be such a pair of associated arcs. According to Petrovskii’s theory, as $`c`$ decreases, the arc $`\gamma `$ comes closer to $`𝐑L`$, the arc $`\gamma ^{}`$ comes closer to an extremity of $`𝐑L^+`$.)
-Let us extend this description of $`𝐑_{2k}`$ to a description in the whole $``$.
Let $`𝒜`$ be a curve and $`A`$ its polynomial. We shall say that the truncation $`D`$ of $`A`$ is sufficient for $`𝒞𝒜`$ in $`𝐑P^2`$ if the curve $`𝒟`$ with polynomial $`D`$ is such that : $`𝐑𝒟`$ is embedded in $`𝐑𝒜`$ and $`𝐑𝒟`$ is homeomorphic to a subset of $`𝐑𝒜`$ which contains $`𝒞`$.
Let us prove the following Lemma:
###### Lemma 2.2.8.
Let $`B_{2k}(x_0,x_1,x_2)`$ be a Harnack polynomial of degree $`2k`$ and type $`^0`$. Up to regular modification of $`B_{2k}(x_0,x_1,x_2)`$,
1. the truncation of $`b_{2k}^\mathrm{\Delta }(x,y)`$ on monomials of homogeneous degree $`2k2i2k`$ is sufficient for $`𝒜𝒜^{}`$.
2. there exists a truncation (on four monomials $`x^cy^d,x^{c+1}y^d,x^cy^{d+1},x^{c+1}y^{d+1}`$ with $`c+d=2k2`$) $`B_{2k}^S`$ of $`B_{2k}^\mathrm{\Delta }`$ for which $`B_{2k}^S`$ is sufficient for $`\gamma \gamma ^{}`$ where $`\gamma 𝒜`$ $`\gamma ^{}𝒜^{}`$.
Denote by $`_{2k}^\mathrm{\Delta }`$ the curve with polynomial $`B_{2k}^\mathrm{\Delta }`$. -Roughly speaking, it means that monomials which are not in $`b_{2k}^\mathrm{\Delta }`$ (resp, $`b_{2k}^S`$) have a small influence on $`𝒜𝒜^{}𝐑_{2k}`$ (resp, $`\gamma \gamma ^{}𝐑_{2k}`$)-
proof: The proof of Lemma 2.2.8 of Chapter 2.2.8 is based on Lemma 2.2.5 of Chapter 2.2.5.
According to the proof of Morse Lemma, local coordinates defined in a neighborhood $`U(p)`$ of a non-degenerate critical point $`p`$ of a function $`f`$ depend principally on the first derivative and the second derivative of the function around this point. Let $`L_{ϵ_1}`$ be a real projective line in the $`ϵ`$-tubular neighborhood of $`x_0=0`$, where $`ϵ>0`$ is arbitrarily small. On the assumption that any critical point $`p`$ belongs to the line $`L_{ϵ_1}`$, monomials which are not in $`b_{2k}^\mathrm{\Delta }`$ have a small influence on $`\{(x_0,x_1,x_2)𝐑P^2|B_{2k}(x_0,x_1,x_2)=0\}U(p)`$. Without loss of generality, we may assume that points at infinity of $`B_{2k}`$ and points of $`𝒫`$ do not belong $`x_2=0`$. In such a way, any point at infinity of an arc of $`𝒜`$ is a root of the polynomial in one variable $`b_{2k}(0,x_1/x_2,1)=b_{2k}(y)`$. Any arc of $`𝒜`$ intersects the line $`x_0=0`$ in two points. Up to regular modification the $`2k1`$ points of $`𝒫`$ belong to $`L_{ϵ_1}`$. According to Rolle’s Theorem, $`b_{2k}(y)`$ has exactly $`2k1`$ extrema at which $`b_{2k}^{}(y)`$ must change from positive to negative. Thus, $`b_{2k}\mathrm{"}(y)`$ has $`2k2`$ zeroes. By continuity, it gives the sign of the function $`b_{2k}(x_0,x_1,x_2)`$ in a neighborhood $`=\{(x_0:x_1:x_2)𝐑P^2||x_0|<ϵ\}`$ of $`x_0=0`$. In a neighborhood $`U(p)`$ of an extremum, $`\{x𝐑P^2|b_{2k}=0\}U(p)`$ is described as a desingularized crossing. From the previous study, it follows that up to regular modification of $`B_{2k}`$, properties 1. and 2. are verified. One can also notice that the alternation of sign of $`B_{2k}`$ in the inside of $``$ is equivalent to the modular property of the distribution of sign of $`b_{2k}^\mathrm{\Delta }`$ in the patchworking construction. Set $`𝒩`$ tubular neighborhood of $`L`$ in $`𝐂P^2`$ with real part $``$. It is easy to verify that replacing $`𝒩`$ in place of $``$ one get a general formulation of the Lemma 2.2.8 of Chapter 2.2.8 in $`𝐂P^2`$ This last remark concludes and completes our proof. Q.E.D
According to the Lemma 2.2.8 of Chapter 2.2.8 the previous description of $`𝐑_{2k}`$ in $`U(p)`$ enlarges to a description of $`𝐑_{2k}`$ in $`UU(p)`$ where
$$U=\{z=<u,p>=(u_0.p_0:u_1.p_1:u_2.p_2)𝐂P^2|u=(u_0:u_1:u_2)U_𝐂^3,p=(p_0:p_1:p_2)U(p)\}$$
Local coordinates $`y_1,y_2`$ in $`U(p)`$ extend to local coordinates in $`U`$ as follows. Given $`z=<u,p>U`$ with $`u=(1:u_1:u_2)U_𝐂^3`$, $`p=(p_0:p_1:p_2)U(p)\}`$, we may set $`y_1(z)=u_1.y_1(p)`$, $`y_2(z)=u_2.y_2(p)`$. In $`U(p)`$, the truncation $`b_{2k}^S`$ is $`ϵ`$-sufficient for $`b_{2k}`$.
$`b_{2k}^S(x,y)=l(x,y)+ϵ.k(x,y)`$
with $`l(x,y)=a_{c,d}x^cy^d+a_{c+1,d}x^{c+1}y^d`$, $`k(x,y)=a_{c,d+1}x^cy^{d+1}+a_{c+1,d+1}x^{c+1}y^{d+1}`$, (with $`a_{c,d}>0,a_{c+1,d}>0,a_{c,d+1}>0,a_{c+1,d+1}>0`$ and $`c+d=2k2`$ and $`ϵ>0`$)
Note that up to modify the coefficients $`a_{c,d},a_{c,d+1},a_{c+1,d},a_{c+1,d+1}`$ if necessary, the point $`p=(x_0,y_0)`$ is (up to homeomorphism) a critical point of the function $`\frac{l(x,y)}{k(x,y)}`$ with positive critical value. Hence, it follows from the equalities
$`l(x,y)=l(x,y)`$, $`k(x,y)=k(x,y)`$,
$`\frac{l}{x}(x,y)=\frac{l}{x}(x,y)`$, $`\frac{l}{y}(x,y)=\frac{l}{y}(x,y)`$,
$`\frac{k}{x}(x,y)=\frac{k}{x}(x,y)`$, $`\frac{k}{y}(x,y)=\frac{k}{y}(x,y)`$
that $`(x_0,y_0)`$ is also a critical point of the function $`\frac{l(x,y)}{k(x,y)}`$ with negative critical value.
The tubular neighborhood of $`𝐑L^+`$ is $`^+=_{p𝒫}U(p)`$. The transformation $`(x,y)(x,y)`$ defined locally inside any $`U(p)`$, $`p𝒫`$ maps the set of arcs $`𝒜𝒜^{}`$ to a set of arcs $``$ and the set of points $`𝒫`$ to a set of points $`𝒮`$. Any open $`U(p)`$ is mapped to an open $`U(s)`$ in such a way that the tubular neighborhood of $`L^{}`$ is $`^{}=_{s𝒮}U(s)`$ and $`L^+`$ is mapped to $`L^{}`$. Precisely, given $`x,y`$ local coordinates in a neighborhood $`U(p)`$, the pair $`(\gamma ,\gamma ^{})(𝒜,𝒜^{})`$ of $`\{(x_0:x_1:x_2)𝐑P^2|B_{2k}(x_0,x_1,x_2)=0\}U(p)`$ is up to homeomorphism the hyperbole $`x.y=1/2`$. By means of the transformation $`(x,y)(x,y)`$ one maps $`U(p)`$ onto $`U(s)`$ and the pair of arcs $`(\gamma ,\gamma ^{})(𝒜,𝒜^{})`$ defined up to homeomorphism by $`x.y=1/2`$ to the pair of arcs $`(\xi ,\xi ^{})`$ of $`\{(x_0:x_1:x_2)𝐑P^2|B_{2k}(x_0,x_1,x_2)=0\}U(p)`$ is up homeomorphism the hyperbole $`x.y=1/2`$. (In $`U(s)`$, one can say that arcs $`\xi `$,$`\xi ^{}`$ are associated and also associated to $`s`$.)
Let us in Lemma 2.2.9 of Chapter 2.2.9 describe the set $`𝒟=^{}𝐑_{2k}^\mathrm{\Delta }`$.
###### Lemma 2.2.9.
Let $`_{2k}`$ be a Harnack curve of type $`^0`$ with polynomial $`B_{2k}`$. Up to regular modification of the Harnack polynomial $`B_{2k}`$, $`𝐑_{2k}^\mathrm{\Delta }`$ intersects $`^{}`$ in $`2k2`$ positive ovals.
proof: The set $`𝒜`$ is a connected set of arcs which intersects $`𝐑L`$ the line at infinity of $`𝐑P^2`$ in $`2k`$ points which are intersection points of two arcs of $`𝒜`$. Using Morse Theory, one can consider these points as limit points of an half hyperbole. Up to isotopy of $`𝐑P^2`$, one can consider arcs of $`𝒜𝒜^{}`$ as arcs lying on $`𝐑L_1^+`$ and $`𝐑L_2^+`$.
Except the two arcs pushed by an isotopy of $`𝐑P^2`$ to a segment of $`𝐑L_1^+`$ which contains an extremity of $`𝐑L_1^+`$, each arc of $`𝒜`$ is associated to an arc of $`𝒜^{}`$ in such a way that the set $`𝒜𝒜^{}`$ is the union of two connected parts lying respectively on $`𝐑L_1^+`$ and $`𝐑L_2^+`$. Each of the two arcs which contains an extremity of $`𝐑L_1^+`$ is associated to an arc $`𝒜^{}`$ which belongs to the part $`^{}`$ of the Möbius band lying between $`𝐑L_1^{}`$ and $`𝐑L_2^{}`$. In our identification of arcs with connected real parts lying in the inside of $``$, one can assume that these two pairs remain the same under the action $`(x,y)(x,y)`$ which maps arcs $`𝒜𝒜^{}`$ to arcs of $``$. (Indeed, $`𝒫`$ intersects $`𝒮`$ in two points which are extremities of $`𝐑L_{ϵ_1}^+𝐑L`$.)
It follows that the union of $`(𝒜𝒜^{})`$ is mapped to $`2(k1)`$ ovals constituted by the union of arcs $``$ union two arcs of $`^+`$ which contains an extremity of $`𝐑L_1^+`$.
These ovals lie in the inside of the Möbius band $`^{}`$ and intersects $`𝐑L^{}`$ in two points. This last remark concludes our proof. Q.E.D
###### Remark 2.2.10.
It follows from the Lemma 2.2.9 of Chapter 2.2.9 a description, up to isotopy of $`𝐑P^2`$, of $`𝐑_{2k}`$ in $``$ which extends our previous description of $`𝒜𝒜^{}𝐑_{2k}`$ Note that this description is consistent with the study of the one variable function $`b_{2k}(y)=b_{2k}(0,x_1/x_2,1)`$. The function $`b_{2k}\mathrm{"}(y)`$ has $`2k2`$ zeroes. Any point at infinity of an arc of $`𝒜`$ is a root of the polynomial in one variable $`b_{2k}(0,x_1/x_2,1)=b_{2k}(y)`$. According to Rolle’s Theorem, $`b_{2k}(y)`$ has exactly $`2k1`$ extrema at which $`b_{2k}^{}(y)`$ must change from positive to negative. The zeroes of $`b_{2k}\mathrm{"}(y)`$ are simple zeroes. These zeroes are the intersection-points of the set $``$ with $`L`$.
The set $`𝒜`$ is the image of a smooth section of a tubular fibration of $`𝐑L^+`$ minus 2 points. In Lemma 2.2.11 of Chapter 2.2.11, we shall prove that there exists a subset $`𝒟^1`$ of $`𝒟=𝐑_{2k}^{}`$ with the property that the set $`𝒜𝒟^1`$ is the image of a smooth section of a tubular fibration of $`𝐑L`$ minus a finite number of points.
By cutting $`^{}`$ along $`L^{}`$, one get two surfaces. Denote by $`^{1,}`$ the one which contains $`L_1^{}`$. Consider the intersection $`𝐑_{2k}^{1,}`$
Any arc $`\xi `$ of $`𝒟=𝐑_{2k}`$ intersects the line at infinity $`L`$. Hence, it is divided into two halves with common point $`L\xi `$. Denote by $`\xi ^1`$ (resp, $`\xi ^2`$) the half of $`\xi `$ which belongs to the inside of Möbius delimited by $`𝐑L`$ and $`𝐑L_1`$ (resp, $`𝐑L`$ and $`𝐑L_2`$) minus its intersection with $`L`$. The set $`𝐑_{2k}^{1,}`$ is the set of arcs $`\xi ^1`$.
In Lemma 2.2.11 of Chapter 2.2.11 we prove that $`𝒟^1=𝐑_{2k}^{1,}`$
###### Lemma 2.2.11.
Let $`_{2k}`$ be a Harnack curve of type $`^0`$ with polynomial $`B_{2k}`$.
Up to regular modification of the Harnack polynomial $`B_{2k}`$, the set $`𝒜𝒟^1`$, where $`𝒟^1=𝐑_{2k}^\mathrm{\Delta }^{1,}`$, is the image of a smooth section of a tubular fibration of $`𝐑L`$ minus $`2k1`$ points.
proof: The set $`𝒜`$ is connected. There exists an isotopy of $`𝐑P^2`$ which pushes the union of the $`k`$ arcs of $`𝒜𝐑^2`$ bounding positive regions $`\{x𝐑^2|b_{2k}(x)>0\}`$ with $`k1`$ arcs of $`𝒜^{}`$ (associated to the $`k1`$ arcs of $`𝒜𝐑^2`$ bounding negative regions$`\{x𝐑^2|b_{2k}(x)>0\}`$) onto $`𝐑L_1^+`$. Let us denote by $`𝒜^1`$ this set of $`k`$ arcs of $`𝒜`$ with $`(k1)`$ arcs of $`𝒜^{}`$ arcs. The transformation $`(x,y)(x,y)`$ maps $`U(p)U(s)`$ and the pair of arcs $`(\gamma ,\gamma ^{})(𝒜,𝒜^{})`$ defined up to homeomorphism by $`x.y=1/2`$ to the pair of arcs $`(\xi ,\stackrel{~}{\xi })`$ of $``$ defined up to homeomorphism by $`x.y=1/2`$. Hence, identifying each arc of $`𝒜^1`$, with a segment of $`𝐑L_1^+U(p)`$, it follows the next descriptions in $`U(s)`$. Let us first consider an open $`U(s)`$ which does not contain an extremity of $`L`$. According to Morse Lemma in $`U(s)`$, it follows that the pair $`(\xi ^1,\stackrel{~}{\xi }^1)`$ of halves of associated arcs in $`U(s)`$ is the image of a smooth section of a tubular fibration of $`U(s)𝐑L_1^{}`$ minus one point.
If $`U(s)`$ contains an extremity of $`L`$, then the pair of arcs $`(\xi ^1,\stackrel{~}{\xi }^1)=(\xi ^1^{},\stackrel{~}{\xi }^1^+)`$ of halves of associated arcs in $`U(s)`$ is such that $`\xi ^1`$ is the image of a smooth section of a tubular fibration of $`U(s)𝐑L_1^{}`$ minus the extremity of $`𝐑L_1^{}U(s)`$.
By means of the transformation $`(x,y)(x,y)`$, inside any open $`U(p)`$, the tubular neighborhood $`^+=_{p𝒫}U(p)`$ of $`L^+`$ is mapped to the tubular neighborhood $`^{}=_{s𝒮}U(s)`$ of $`L^{}`$. It follows from Morse Lemma that the set $`𝒟^1`$ is the image of a smooth section of tubular fibration of $`𝐑L_1^{}`$ minus $`2k1`$ points.
Hence, the set $`𝒟^1`$ is also the image of a smooth section of tubular fibration of $`𝐑L^{}`$ minus $`2k1`$ points.
Since $`𝒫`$ intersects $`𝒮`$ in the extremities of $`𝐑L_{ϵ_1}^+`$ (see proof of Lemma 2.2.9 of Chapter 2.2.9), two of these $`2k1`$ points of $`𝐑L^{}`$ are extremities.
Hence, $`𝒟^1`$ is the image of a smooth section of tubular fibration of $`𝐑L^{}`$ minus $`𝒮`$. Therefore, since $`𝒜`$ is the image of a smooth section of tubular fibration of $`𝐑L_+`$ minus its extremities, $`𝒜𝒟^1`$ is the image of a smooth section of tubular fibration of $`𝐑L`$ minus $`2k1`$ points. Q.E.D
According to Proposition 2.1.4 of Chapter 2.1.4, given $`B_{2k}`$ and $`B_{2k1}`$ Harnack polynomials of type $`^0`$ and respective degree $`2k`$ and $`2k1`$
$$c(B_{2k})=(\frac{(k1)(k2)}{2},k.(2k1),\frac{3k.(k1)}{2}))$$
$$c(B_{2k1})=(\frac{(k1).(k2)}{2},(k1).(2k1),\frac{(k1).(3k4)}{2})$$
The Harnack curve $`_{2k}`$ of degree $`2k`$ has $`2(k2)`$ ( $`:=c_2(B_{2k})c_2(B_{2k1})`$) more positive ovals than the Harnack curve $`_{2k1}`$ of degree $`2k1`$. It follows from Lemma 2.2.8 of Chapter 2.2.8 and Lemma 2.2.9 of Chapter 2.2.9 that, up to regular to modification of the Harnack polynomial $`B_{2k}`$, these ovals are the ovals of $`𝐑_{2k}^\mathrm{\Delta }^{}`$. In such a way, it follows from Lemma 2.2.11 of Chapter 2.2.11 that, up to regular modification of $`B_{2k}`$, $`B_{2k}=x_0.B_{2k1}+C_{2k}`$ where $`B_{2k1}`$ is a Harnack polynomial of type $`^0`$ relatively to $`L:=x_0=0`$ the line at infinity. (Indeed, as we have noticed, the truncation $`B_{2k}^\mathrm{\Delta }`$ is sufficient for $`B_{2k}`$ and $`L_{ϵ_1}`$ may be chosen arbitrarily close to $`L`$.) It concludes our proof of Proposition 2.2.4 of Chapter 2.2.4. Q.E.D
I.ii)-Harnack curves of odd degree-
Let $`_{2k+1}`$ be a Harnack curve of degree $`2k+15`$ with polynomial $`B_{2k+1}(x_0,x_1,x_2)`$ of type $`^0`$ relatively to the line at infinity $`L:=x_0=0`$.
We shall prove the following Proposition 2.2.12 of Chapter 2.2.12.
###### Proposition 2.2.12.
Let $`_{2k+1}`$ be a Harnack curve of degree $`2k+1`$ with polynomial $`B_{2k+1}(x_0,x_1,x_2)`$ of type $`^0`$ relatively to the line at infinity $`L:=x_0=0`$. Up to regular modification of $`B_{2k+1}`$, $`B_{2k}(x_0,x_1,x_2)`$ is also of type $`^0`$ relatively to $`L`$.
proof:
The part I.ii) until its end is devoted to the proof of Proposition 2.2.12 of Chapter 2.2.12.
Our argumentation is a slightly modified version of the one given in the proof of Proposition 2.2.12 of Chapter 2.2.12. When it is possible, we refer to the first part I.i).
The odd component of $`_{2k+1}`$ is divided into $`2k+1`$ arcs which delimit a region bounding the line at infinity. One of this region contains ovals of $`_{2k+1}`$. Denote $`𝒜`$ the set of $`2k`$ arcs of $`𝐑_{2k+1}`$ which delimit a positive region $`\{x𝐑^2|b_{2k+1}>0\}`$ (i.e region which does not contain oval of $`_{2k+1}`$). The union of arcs of $`𝒜`$ is a connected orientable part of the odd component of $`𝐑_{2k+1}`$. The odd component of $`𝐑_{2k+1}`$ is homeomorphic to a projective line. Consider an orientation of the line at infinity $`𝐑L`$. The surface delimited by the $`2k`$ arcs of $`𝒜`$ and $`𝐑L`$ is orientable. Thus, one can assume that the orientation of the connected union of the $`2k`$ arcs of $`𝒜`$ is induced by a half of $`𝐑L`$. We shall denote by $`𝐑L^+`$ this half. The set $`𝒜`$ is the image of a smooth section of a tubular fibration of $`𝐑L^+`$ minus 2 points.
Let us prove that one can associate to the set $`𝒜`$ of the $`2k`$ arcs a set of critical points $`𝒫`$ of index $`1`$ of $`b_{2k+1}`$ and a set of arcs $`𝒜^{}`$ of $`𝐑_{2k+1}`$.
To this end, we shall study the Petrovskii’s pencil $`x_0^{2k+1}.(b_{2k+1}(x_1/x_0,x_2/x_0)c)`$, $`c𝐑`$, over $`_{2k+1}`$.
As $`c`$ decreases from $`0`$ positive regions and positive ovals expand. As $`c`$ decreases from $`0`$ the number of regions $`G_i`$ of $`\{x𝐑^2|b_{2k}(x)>0\}`$ which contain a segment of the line at infinity of $`𝐑P^2`$ on their boundary decreases from $`k+1`$ to $`1`$.
As $`c`$ increases from $`0`$ negative ovals disappear; the number of regions $`G_i`$ of $`\{x𝐑^2|b_{2k+1}(x)>0\}`$ which contain a segment of the line at infinity of $`𝐑P^2`$ on their boundary remains the same. Hence, to each arc of $`𝒜`$ is associated a critical of index $`1`$.
-Let us study the set $`𝒫`$ of critical point $`1`$ associated to the set $`𝒜`$ of arcs of $`𝐑_{2k}`$.
Let us prove the following Lemma:
###### Lemma 2.2.13.
There exist at most two arcs of the set $`𝒜`$ which may be deformed as $`c`$ varies from $`0`$ to $`\mathrm{}`$ or $`+\mathrm{}`$ in such a way that for $`c=c_0`$ they touch each other.
proof: The projective plane $`𝐑P^2`$ may be seen as a Möbius band $``$ and a disc $`D^2`$ glued along their common boundary where the core of the Möbius band is the line at infinity $`𝐑L`$. Assume that there exists $`c_0`$, such that as $`c`$ varies from $`0`$ to $`c_0`$ an arc $`\gamma `$ of $`𝒜`$ is deformed and for $`c=c_0`$ it touches an oval $`𝒪`$ of $`_{2k+1}`$. The arc $`\gamma `$ bounds positive region . Hence, if $`c_0`$ exists it is negative. The inside of the oval $`𝒪`$ is an orientable part of $`𝐑P^2`$ homeomorphic to $`D^2`$. The arc $`\gamma `$ may be pushed by an isotopy to the boundary $`𝐑L_1^+`$ of the Möbius band. In this way, for $`c=c_0`$, $`\gamma `$ may be also pushed to an arc $`\stackrel{~}{\gamma }`$ of the oval $`𝒪`$. Since $`c_0<0`$, according to Petrovskii’s Lemmas, for $`c=c_0+ϵ`$ one can trace a non-orientable branch (i.e an arc of $`𝐑L`$ which intersects $`𝐑L^+`$ and $`𝐑L^{}`$) of $`𝐑P^2`$ in the region $`b_{2k}<c_0ϵ`$ containing $`\gamma `$ and $`\stackrel{~}{\gamma }`$. The set $`𝒜`$ may be pushed by an isotopy onto $`𝐑L_1^+`$. In a neighborhood of $`𝐑L_1^+`$, a non-oriented branch of $`𝐑P^2`$ may be traced only near the extremities of $`𝐑L_1^+`$. Hence, the only two arcs of $`𝒜`$ which may glue with an oval are those which intersect an extremity of $`𝐑L_1^+`$. Q.E.D It follows from Petrovskii’s Lemmas 2.1.2 of Chapter 2.1.2, 2.1.3 of Chapter 2.1.3 and Lemma 2.2.13 of Chapter 2.2.13, that there exist $`2k2`$ critical points with critical value $`c_0<0`$ such that as $`c`$ varies from $`0`$ to $`c_0`$ two arcs are deformed in such a way that for $`c=c_0`$ they touch each other.
As already done for Harnack curves of even degree one can study critical points $`𝒫`$ of $`B_{2k+1}`$ using the function in one variable $`b_{2k+1}(y)=b_{2k+1}(0,x_1/x_2,1)`$. According to Rolle’s Theorem applied to $`b_{2k+1}(y)=b_{2k+1}(0,x_1/x_2,1)`$ for any pair ($`\gamma ,\gamma ^{})`$ of associated arcs of $`𝒜`$ there exists $`\gamma \mathrm{"}𝒜`$ such $`\gamma \gamma \mathrm{"}\mathrm{}`$ $`\gamma ^{}\gamma \mathrm{"}\mathrm{}`$.
According to Lemma 2.2.13 of Chapter 2.2.13 and its proof, for the $`2`$ arcs of $`𝒜`$ which intersect $`L`$ in one of its extremities, there exists $`c_0<0`$ such that as $`c`$ varies from $`0`$ to $`c_0`$, the arc is deformed and for $`c=c_0`$ it touches a positive oval.
In such a way, any arc $`\gamma `$ of $`𝒜`$ is associated to a critical point $`p`$ and to an other arc $`\gamma ^{}`$. We denote by $`𝒫`$ the set of points associated to $`𝒜`$ and by $`𝒜^{}`$ the set of arcs associated to $`𝒜`$. The sets $`𝒜^{}`$ and $`𝒜`$ consist of $`2k`$ arcs of $`_{2k+1}`$; $`2k2`$ of them belong to $`𝒜`$ and $`𝒜^{}`$.
-Let us give a description up to rigid isotopy of $`𝐑_{2k}^\mathrm{\Delta }`$ in the whole $``$.
Let $`B_{2k+1}`$ be a Harnack polynomial of degree $`2k+1`$ and type $`^0`$.
From a version of Lemma 2.2.5 of Chapter 2.2.5 for odd degree curves $`_{2k+1}`$, it follows that there exists a rigid isotopy $`B_{2k+1,t}`$ ,$`t[0,1]`$, $`B_{2k+1,0}=B_{2k+1}`$, such that the modified polynomial $`B_{2k+1,1}`$ has the following the property: the $`2k1`$ critical points $`𝒫`$ associated to $`𝒜`$ belong to a line $`L_{ϵ_1}`$ of $`_ϵ`$ the $`ϵ`$-tubular neighborhood of $`L`$ in $`𝐑P^2`$ with $`ϵ>0`$ arbitrarily small.
From a version of Lemma 2.2.8 of Chapter 2.2.8 for odd degree, the truncation $`b_{2k+1}^\mathrm{\Delta }(x,y)`$ of $`b_{2k+1}(x,y)`$ on monomials of homogeneous degree $`2k1i2k+1`$ is sufficient for $`𝒜𝒜^{}`$.
Denote by $`_{2k+1}^\mathrm{\Delta }`$ the curve with polynomial $`B_{2k+1}^\mathrm{\Delta }(x_0,x_1,x_2)=x_0^{2k+1}.b_{2k+1}^\mathrm{\Delta }(x,y)`$ the description of $`𝐑_{2k+1}^\mathrm{\Delta }`$ in any neighborhood $`U(p)`$ $`p𝒫`$ enlarges to a description of $`𝐑_{2k+1}^\mathrm{\Delta }`$ in $`UU(p)`$ where
$$U=\{z=<u,p>=(u_0.p_0:u_1.p_1:u_2.p_2)𝐂P^2|u=(u_0:u_1:u_2)U_𝐂^3,p=(p_0:p_1:p_2)U(p)\}$$
Consider $`^+`$ the tubular neighborhood of $`L^+`$, $`^+=_{p𝒫}U(p)`$ and $`^{}`$ the tubular neighborhood of $`L^{}`$.
From an argumentation analogous to the one given in the part I.i), according to Lemma 2.2.8 of Chapter 2.2.8 we define a set $``$ of arcs of $`𝐑_{2k+1}^\mathrm{\Delta }`$ in $`^{}`$.
Let $`(x,y)`$ be local coordinates in $`U(p)`$ , $`p𝒫`$. The transformation, defined locally in any $`U(p)`$ by $`(x,y)(x,y)`$ maps $`U(p)U(s)`$ and in this way $`^+=_{p𝒫}U(p)`$ to $`^{}=_{s𝒮}U(s)`$. It maps the set of arcs $`𝒜𝒜^{}`$ to a set of arcs $``$ and the set of points $`𝒫`$ to a set of points $`𝒮`$.
Let us in Lemma 2.2.14 of Chapter 2.2.14 describe the set $`𝒟=^{}𝐑_{2k+1}^\mathrm{\Delta }`$.
###### Lemma 2.2.14.
Up to regular modification of the Harnack polynomial $`B_{2k+1}`$, $`𝐑_{2k+1}^\mathrm{\Delta }`$ intersects $`^{}`$ in $`k1`$ negative ovals and $`k`$ positive ovals.
proof: Our argumentation is similar to the one given in the proof of Lemma 2.2.9 of Chapter 2.2.9.
The set $`𝒜`$ is a connected set of arcs which intersects $`L`$ the line at infinity of $`𝐑P^2`$ in $`2k`$ points. Using Morse Theory, one can consider intersection points of the set $`𝒜`$ with the line $`L`$ as limit points of an half hyperbole. Except the two arcs of $`𝒜`$ which contain an extremity of $`L^+`$, any arc of $`𝒜`$ is associated to another arc of $`𝒜`$. Given $`U(p)`$ with local coordinates $`(x,y)`$, the transformation $`(x,y)(x,y)`$ maps any pair of arcs of $`𝒜`$ to another pair of arcs $`𝐑_{2k+1}`$. It follows from the alternation of signs of $`b_{2k+1}`$ in any $`U(p)`$ of $`^+`$, the alternation of sign $`b_{2k}`$ in any $`U(s)`$ of $`^{}`$ and thereby around branches of $`𝐑_{2k+1}`$. It follows that any two arcs of $`𝒜`$ associated to each other in $`U(p)`$ are mapped to an arc of a positive oval and an arc of a negative oval in $`U(s)`$. Each of the two arcs which contains an extremity of $`𝐑L^+`$ is associated to an arc $`𝒜^{}`$ which belongs to $`^{}`$. One can assume that these two pairs remain the same under the action $`(x,y)(x,y)`$ which maps arcs $`𝒜𝒜^{}`$ to arcs of $``$. It follows that the union of $`(𝒜𝒜^{})`$ is mapped to the union $``$ of two arcs of the odd component of $`_{2k+1}`$ which contain extremities of $`𝐑L_1^+`$ (these arcs lie in $`^+`$) and $`k1`$ negative ovals and $`k`$ positive ovals of $`_{2k+1}`$ constituted by the union of arcs $``$. These ovals lie in the inside $`^{}`$ of the Möbius band and intersect $`𝐑L^{}`$ in two points. This last remark concludes our proof. Q.E.D
The set $`𝒜`$ is the image of a smooth section of a tubular fibration of $`𝐑L^+`$ minus 2 points. In Lemma 2.2.15 of Chapter 2.2.15, we shall prove that there exists a subset $`𝒟^1`$ of $`𝒟=𝐑_{2k+1}^\mathrm{\Delta }^{}`$ with the property that the set $`𝒜𝒟^1`$ is the image of a smooth section of a tubular fibration of $`𝐑L`$ minus a finite number of points.
By cutting $`^{}`$ along $`L^{}`$, one get two surfaces. Denote by $`^{1,}`$ the one which contains $`L_1^{}`$. Consider the intersection $`𝐑_{2k+1}^\mathrm{\Delta }^{1,}`$
Any arc $`\xi `$ of $`𝒟=𝐑_{2k+1}^\mathrm{\Delta }`$ intersects the line at infinity $`L`$. Hence, it is divided into two halves with common point $`L\xi `$. Denote by $`\xi ^1`$ (resp, $`\xi ^2`$) the half of $`\xi `$ which belongs to the inside of Möbius delimited by $`𝐑L`$ and $`𝐑L_1`$ (resp, $`𝐑L`$ and $`𝐑L_2`$) minus its intersection with $`L`$. The set $`𝐑_{2k}^\mathrm{\Delta }^{1,}`$ is the set of arcs $`\xi ^1`$.
In Lemma 2.2.11 of Chapter 2.2.11 we prove that $`𝒟^1=𝐑_{2k+1}^\mathrm{\Delta }^{1,}`$
###### Lemma 2.2.15.
Up to regular modification of the Harnack polynomial $`B_{2k+1}`$, the set $`𝒜𝒟^1`$ is the image of a smooth section of a tubular fibration of $`𝐑L`$ minus $`2k`$ points.
proof: Our proof is analogous to the one of Lemma 2.2.11 of Chapter 2.2.11. The set $`𝒜`$ is the image of a smooth section of a tubular fibration of $`𝐑L^+`$ minus 2 points The transformation $`(x,y)(x,y)`$ maps $`U(p)U(s)`$ and the pair of arcs $`(\gamma ,\gamma ^{})(𝒜,𝒜^{})`$ defined up to homeomorphism by $`x.y=1/2`$ to the pair of arcs $`(\xi ,\stackrel{~}{\xi })`$ of $``$ defined up to homeomorphism by $`x.y=1/2`$.
Hence, identifying each of arcs of $`𝒜`$, with a segment of $`𝐑L^+U(p)`$, it follows the next descriptions in $`U(s)`$.
Let us first consider open $`U(s)`$ which does not contain an extremity of $`L`$. According to Morse Lemma in $`U(s)`$, it follows that the pair $`(\xi ^1,\stackrel{~}{\xi }^1)`$ of halves of associated arcs in $`U(s)`$ is the image of a smooth section of a tubular fibration of $`U(s)𝐑L_1^{}`$ minus one point.
If $`U(s)`$ contains an extremity of $`L`$, then the pair of arcs $`(\xi ^1,\stackrel{~}{\xi }^1)=(\xi ^1^{},\stackrel{~}{\xi }^1^+)`$ of halves of associated arcs in $`U(s)`$ is such that $`\xi ^1`$ is the image of a smooth section of a tubular fibration of $`U(s)𝐑L_1^{}`$ minus the extremity of $`𝐑L_1^{}U(s)`$.
By means of the transformation $`(x,y)(x,y)`$, inside any open $`U(p)`$, the tubular neighborhood $`^+=_{p𝒫}U(p)`$ of $`L^+`$ is mapped to the tubular neighborhood $`^{}=_{s𝒮}U(s)`$ of $`L^{}`$. It follows from Morse Lemma that the set $`𝒟^1`$ is the image of a smooth section of tubular fibration of $`𝐑L^{}`$ minus $`2k`$ points.
Two of these $`2k`$ points of $`𝐑L^{}`$ are extremities of $`𝐑L^+`$
Hence, $`𝒟^1`$ is the image of a smooth section of tubular fibration of $`𝐑L_{}`$ minus $`𝒮`$. Therefore, $`𝒜𝒟^1`$ is the image of a smooth section of tubular fibration of $`𝐑L`$ minus the $`2k`$ points. Q.E.D
According to Proposition 2.1.4 of Chapter 2.1.4, given $`B_{2k+1}`$ and $`B_{2k}`$ Harnack polynomials of type $`^0`$ and respective degree $`2k`$ and $`2k1`$
$$c(B_{2k+1})=(\frac{(k).(k1)}{2},k.(2k+1),\frac{(k).(3k1)}{2})$$
$$c(B_{2k})=(\frac{(k1)(k2)}{2},k.(2k1),\frac{3k.(k1)}{2})$$
The Harnack curve $`_{2k+1}`$ of degree $`2k+1`$ has $`k1`$ negative ovals and $`k`$ positive ovals more than $`_{2k}`$.
It follows from the version of Lemma 2.2.8 of Chapter 2.2.8 for odd degree curves and Lemma 2.2.14 of Chapter 2.2.14 that up to regular to modification of the Harnack polynomial $`B_{2k+1}`$, these ovals are the ovals of $`𝐑_{2k+1}^\mathrm{\Delta }^{}`$. In such a way, it follows from Lemma 2.2.15 of Chapter 2.2.15 that, up to regular modification of $`B_{2k+1}`$, $`B_{2k+1}=x_0.B_{2k}+C_{2k+1}`$ where $`B_{2k}`$ is a Harnack polynomial of type $`^0`$ relatively to $`x_0=0`$ the line at infinity. It concludes our proof of Proposition 2.2.12 of Chapter 2.2.12. Q.E.D
II) For curves of degree $`i4`$ the Proposition 2.2.3 of Chapter 2.2.3 is immediate. It is sufficient to use the fact that Harnack curves of degree $`i4`$ are rigidly isotopy and are $`L`$-curves.
(When consider the curve $`_3`$, one can also give the following argument. Assume $`_4`$ is a curve of type $`^0`$ which results from deformation of $`_3L`$. The curve $`_3`$ has necessarily two real connected components. Noticing that none of the polynomials $`x_0^3+a_1.x_1^3+a_2.x_2^3`$ with $`a_1,a_2𝐑`$ is a polynomial of a curve of degree $`3`$ with two real connected components, we get the Proposition 2.2.3 of Chapter 2.2.3 for $`_3`$.) Q.E.D Conclusion: Proof of Theorem 2.2.2 of Chapter 2.2.2:
Bringing together Theorem 2.2.32 of Chapter 2.2.32 of the Appendix, Proposition 2.2.4 of Chapter 2.2.4 Proposition 2.2.12 of Chapter 2.2.12, and this last result for curves of degree $`4`$ we get by induction the Theorem 2.2.2 of Chapter 2.2.2. Q.E.D
The following remark may easily deduced from Theorem 2.2.2 of Chapter 2.2.2.
###### Remark 2.2.16.
Let $`L`$ be a real projective line and $`S`$ a set of $`m`$ real points lying on $`L`$. Denote by $`_m^{}`$ a curve of degree $`m`$ of type $`^0`$ relatively to a real projective line $`L^{}`$. There exists a rigid isotopy of $`𝐑P^2`$ which maps $`_m^{}`$ of type $`^0`$ to a curve $`\stackrel{~}{_m}`$ of type $`^0`$ relatively to $`L`$ such that $`\stackrel{~}{_m}L=S`$.
##### A generalization of Harnack’s Method
Harnack curves $`_i`$ of degree $`i4`$ are $`L`$-curves. Hence, up to rigid isotopy, they result from the successive classical deformation of the union $`_jL_{j+1}`$, $`1j4`$, where $`L_{j+1}`$ is a real projective line which intersects $`_j`$ in $`j`$ real points. In Theorem 2.2.30 of Chapter 2.2.30, we extend this property to Harnack curves $`_m`$ of arbitrary degree. Precisely, we prove that up to regular modification of its polynomial any Harnack curve results from the Harnack’s method (i.e $`L_{j+1}=L`$ for any $`j0`$). In this way, the Rigid Isotopy Classification Theorem 2.2.1 of Chapter 2.2.1. follows from the Theorem 2.2.30 of Chapter 2.2.30.
Given $`_m`$ a Harnack curve with polynomial $`B_m(x_0,x_1,x_2)`$, we denote by $`_i`$ the curve with polynomial $`B_i(x_0,x_1,x_2)=x_0^i.b_i(x_1/x_0,x_2/x_0)`$ where $`B_i(x_0,x_1,x_2)`$ is the homogeneous polynomial associated to the truncation $`b_i(x,y)`$ of $`b_m(x,y)`$ on the monomials $`x^\alpha .y^\beta `$ with $`0\alpha +\beta i`$.
Let us prove the following Proposition:
###### Proposition 2.2.17.
Up to rigid isotopy, one can always assume that $`_i`$, $`1i3`$, is of type $`^0`$ relatively to the line at infinity $`L:=x_0=0`$.
proof :
In the proof of the Rigid Isotopy Classification Theorem 2.2.1 of Chapter 2.2.1, we shall use only the fact that up to rigid isotopy, one can always assume that $`_i`$, $`1i3`$, is of type $`^0`$ relatively to the line at infinity $`L:=x_0=0`$. Our argumentation is based on the fact that up to rigid isotopy there exists only one $`M`$-curve of degree $`3`$.
The proof is based on two lemmas. I) Noticing that, up to regular modification of $`B_m(x_0,x_1,x_2)`$, the truncation $`B_3(x_0,x_1,x_2)`$ is the polynomial of a smooth curve, we shall prove that, up to regular modification of $`B_m(x_0,x_1,x_2)`$, $`B_3(x_0,x_1,x_2)`$ is of type $`^0`$. Our argumentation is based on the following statement:
A nonsingular cubic has exactly three real inflection points. These inflection points are collinear.
Denote by $`s_i`$, $`0i2`$, the real inflection points of $`B_3`$ and by $`L_i`$, $`0i2`$, the tangent to $`_3`$ at $`s_i`$.
Let us prove the following Lemma:
###### Lemma 2.2.18.
For suitable projective coordinates, and up to rigid isotopy of $`B_m(x_0,x_1,x_2)`$, $`_3`$ results from the classical deformation:
$`B_3(x_0,x_1,x_2)=x_0.B_2(x_0,x_1,x_2)+ϵ.L_0.L_1.L_2`$ where $`B_2(x_0,x_1,x_2)`$ is of type $`^0`$ relatively to $`x_0=0`$.
($`L_ix_0=0`$, $`B_3(0,1,1)=ϵ.L_0.L_1.L_2(0,1,1)`$)
proof: Consider the Taylor expansion for $`B_3(x_0,x_1,x_2)`$ at each inflection point $`s_i`$, $`0i2`$.
One can assume $`s_i\{(x_0:x_1:x_2)𝐑P^2|x_i=0\}`$ and consider $`b_{3,i}(x,y)=B_3(x_0,x_1,x_2)|_{x_i=1}`$. In a neighborhood $`U(s_i)`$ of $`s_i=(a_i,b_i)`$,
$`b_{3,i}=\mathrm{\Sigma }_{n=0}^3\frac{1}{n!}\mathrm{\Sigma }_{k=0}^n(\frac{^nB_3}{^kx^{nk}y}(a_i,b_i)(xa)^k(yb)^{nk})`$ Each tangent $`L_i`$ meets $`_3`$ with multiplicity $`3`$. Since $`\frac{^2B_3}{^kx^{2k}y}(a_i,b_i)=0`$ for any $`s_i`$, $`0i2`$, when consider the Taylor expansion of $`B_3`$, in any neighborhood $`U(s_i)𝐑P^2`$ of the inflection point $`s_i`$ $`B_2(s)=L_i(s)`$ for any $`sU(s_i)`$; and thus $`B_3(s)L_i(s)=C_3(s)`$ for any $`sU(s_i)`$.
Consider the linear change of coordinates which maps $`(x_0,x_1,x_2)`$ to $`(x_0^{}=L_0,x_1^{}=L_1,x_2^{}=L_2)`$. It maps $`B_3(x_0,x_1,x_2)`$ to $`B_3^{}(x_0^{},x_1^{},x_2^{})`$ . Let us prove that it also maps $`C_3`$ to $`ϵ.x_0^{}.x_1^{}.x_2^{}`$. It maps $`s_i`$ to $`s_i^{}`$ and one can assume $`s_i^{}\{(x_0^{}:x_1^{}:x_2^{})𝐑P^2|x_i^{}=0\}`$. Consider the Taylor expansion of $`b_{3,i}^{}`$ in a neighborhood $`U(s_i^{})=\{(x_0^{}:x_1^{}:x_2^{})𝐂P^2|x_i^{}0\}`$ of $`s_i^{}`$. For any $`i\{0,1,2\}`$, $`j,ki`$ $`j,k\{0,1,2\}`$, $`b_{3,i}^{}(s)x_j^{}(s).x_k^{}(s)=C_3^{}(s)`$ for any $`sU(s_i)`$. Hence, it follows that $`C_3^{}=ϵ.L_0.L_1.L_2`$. since $`L_0.L_1.L_2(s)ϵ.C_3^{}(s)=0`$ for any $`sU(s_i)`$. It concludes the proof of Lemma 2.2.18 of Chapter 2.2.18. Q.E.D II)
###### Definition 2.2.19.
Let $`𝒜`$ and $``$ be smooth algebraic curves of $`𝐂P^2`$ with respective order $`i`$ and $`j`$, $`ij`$. Let $`𝐑𝒜`$ ,resp. $`𝐑`$, be the real point set of $`𝒜`$, resp. $``$. We shall say $`𝒜`$ is $`\mathrm{𝑖𝑚𝑚𝑒𝑟𝑠𝑒𝑑}`$ in $``$ if, up to rigid isotopy of $`𝐑`$, $`𝐑𝒜`$ is embedded in $`𝐑`$.
###### Definition 2.2.20.
Given $`A(x_0,x_1,x_2)=x_0^i.a_(x_1/x_0,x_2/x_0)`$ the polynomial of a curve $`𝒜`$ immersed in $``$. Assume $`A`$ regular. Denote by $`𝒫_𝒜`$ the pencil of curves over $`𝒜`$. (i.e curves with polynomial $`x_0^i(a_(x_1/x_0,x_2/x_0)c)`$, $`c𝐑`$.) We shall say that $`𝒞𝒫_𝒜`$ is $`Mimmersed`$ (over $`𝒜`$) in $``$ if:
1. $`𝒞`$ is immersed in $``$
2. $`𝒞`$ has the maximal number of real components a curve of $`𝒫_𝒜`$ immersed in $``$ may have.
We shall prove in Lemma 2.2.21 of Chapter 2.2.21 that a curve $`_3`$ is immersed in $`_m`$.
###### Lemma 2.2.21.
Let $`_m`$ be the curve with polynomial $`B_m`$. Up to regular modification of $`B_m`$, the curve $`_3`$ is the curve $`_3`$ where $`_3`$ is immersed in $`_m`$ as the classical deformation of $`_2L`$ where $`L:=x_0=0`$.
The curve $`_3`$ is determined up to rigid isotopy by its real scheme. The two possible real schemes for $`_3`$ are $`J1`$ and $`J`$.
Hence, according to Lemma 2.2.18 of Chapter 2.2.18, up to rigid isotopy of $`B_m(x_0,x_1,x_2)`$, which is also a rigid isotopy of $`B_3(x_0,x_1,x_2)`$, one can assume that $`B_3(x_0,x_1,x_2)`$ results from the classical deformation of the union of a line $`L`$ with a curve $`_2`$ of degree $`2`$ with real scheme $`1`$; the deformation is directed to the union of the tangents at real inflection points of $`_3`$. Real inflection points of $`_3`$ are points of $`L`$. It is not hard to see that if none of these points belongs to the inner of $`_2`$, then $`_3`$ is of type $`^0`$ relatively $`L`$. Otherwise, if at least one of this point belongs to the inner of $`_2`$ ,then $`_3`$ has real scheme $`J`$.
Denote $`L`$ the line infinity $`x_0=0`$, and let $`b_3(x,y)=B_3(1,x_1,x_2)`$ be the affine polynomial associated to $`B_3(x_0,x_1,x_2)`$. Consider the pencil $`x_0^m.(b_3(x_1/x_0,x_2/x_0)c)`$. For any $`i4`$, let $`C_i(x_1,x_2)`$ be the polynomial of degree $`i`$ in the variables $`x_1,x_2`$ such that: $`B_m(x_0,x_1,x_2)=x_0^{m3}.B_3(x_0,x_1,x_2)+\mathrm{\Sigma }_{i=4}^mx_0^{mi}.C_i(x_1,x_2)`$.
Consider the pencil $`x_0^{m3}.(b_3(x_1/x_0,x_2/x_0)c)`$. Assume $`B_3`$ regular. For any critical point $`(1,x_{0,1},x_{0,2})`$ of $`B_3(1,x_1,x_2)=b_3(x_1,x_2)`$, choose its representative in $`S^2`$ $`\frac{1}{(1+x_{0,1}^2+x_{0,2})^{1/2}}(1,x_{0,1},x_{0,2})`$ In such a way, it is easy to see that a curve of degree $`3`$ is immersed in $`_m`$. Since $`_m`$ has the maximal number of real components a curve of degree $`m`$ may have, an $`M`$-curve of degree $`3`$ is $`M`$-immersed in $`_m`$. Indeed, if the curve $`_3`$ is not an $`M`$-curve, it is an easy consequence of the Petrovskii’s theory that as $`c`$ varies from $`[\mathrm{\Sigma }_{j=4}^mC_j,\mathrm{\Sigma }_{j=4}^mC_j]`$ the real point set of the curve $`_3`$ undergoes at least one Morse modification (i.e at least one critical point passes through its critical value.) According to Lemma 2.2.18 of Chapter 2.2.18, this is equivalent to say that any union of lines $`L_iL`$ $`1i3`$ is perturbed in such a way none of the resulting real branches intersects the oval $`_2`$. It follows that, up to regular deformation of $`B_m(x_0,x_1,x_2)`$, $`B_3(x_0,x_1,x_2)`$ is of type $`^0`$. Since curves of degree $`m3`$ are defined up to rigid isotopy by their real scheme, up to rigid isotopy of $`B_m(x_0,x_1,x_2)`$, $`B_2(x_0,x_1,x_2)`$, and in this way also $`B_1(x_0,x_1,x_2)`$ are of type $`^0`$. It concludes the proof of Lemma 2.2.21 of Chapter 2.2.21. Q.E.D The Proposition 2.2.17 of Chapter 2.2.17 is a straightforward consequence of the Lemma 2.2.21 of Chapter 2.2.21. Q.E.D
Denote $`\mathrm{\Omega }_j`$ the set of $`M`$-curves with real scheme:
-for even $`j=2k`$
$`1\alpha \beta `$
with $`\alpha +\beta =\frac{(j1)(j2)}{2}`$
-for odd $`j=2k+1`$
$`J\gamma `$
with $`\gamma =\frac{(j1)(j2)}{2}`$
###### Proposition 2.2.22.
Let $`_m`$ be a Harnack curve of degree $`m4`$. Denote by $`B_m`$ its polynomial. Up to regular modification of $`B_m`$, the curve $`_m`$ results from the deformation of the union $`𝒜_{mi}_i`$, $`i3`$, where $`𝒜_{mi}\mathrm{\Omega }_{mi}`$ and $`_i`$ is a Harnack curve of degree $`i3`$.
proof:
Set $`L:=x_0=0`$. It follows from Proposition 2.2.17 of Chapter 2.2.17 that one can set
$$B_m(x_0,x_1,x_2)=A_{mi}(x_0,x_1,x_2).B_i(x_0,x_1,x_2)+C_m(x_0,x_1,x_2)$$
where $`B_i(x_0,x_1,x_2)`$ is a Harnack of degree $`i3`$ and type $`^0`$ relatively to $`L`$ and $`A_{mi}(x_0,x_1,x_2)=x_0^{mi}`$.
Let us prove that there exists a regular modification of $`B_{m,t}(x_0,x_1,x_2)`$, $`t[0,1]`$,
$$B_{m,0}(x_0,x_1,x_2)=B_m(x_0,x_1,x_2)$$
$$B_{m,t}(x_0,x_1,x_2)=A_{mi,t}(x_0,x_1,x_2).B_{i,t}(x_0,x_1,x_2)+C_{m,t}(x_0,x_1,x_2)$$
such that:
1. $`B_{i,t}(x_0,x_1,x_2)`$ is a regular modification of $`B_i(x_0,x_1,x_2)`$ the Harnack polynomial of a Harnack curve $`_i`$,
2. $`A_{mi,1}`$ is the polynomial of a smooth curve of degree $`mi`$ ($`A_{mi,0}=x_0^{mi}`$)
3. The polynomial $`B_{m,1}`$ is the polynomial of a curve $`_m`$ which results from the classical deformation of the union of $`𝒜_{mi,1}_i`$ with polynomial $`A_{mi,1}B_{i,1}`$
(It is obvious that the regular modification $`B_{m,t}`$ of $`B_m`$ is not regular on $`A_{mi}`$) Our proof is based on the fact that curves of degree $`i4`$ are determined up to rigid isotopy by their real schemes and may be realized as $`L`$-curves (i.e may realized by classical small perturbation of $`i`$ lines in general position) . We shall consider immersion (see the definition 2.2.19 of Chapter 2.2.19 of the proof of Proposition 2.2.17 of Chapter 2.2.17), of such curves in the curve $`_m`$ with polynomial $`B_m`$.
We shall consider curves $`_m`$ of degree $`m7`$ and degree $`m8`$ separately.
1)
Let us consider curves $`_m`$ of degree $`m7`$. According to the Proposition 2.2.17 of Chapter 2.2.17, up to rigid isotopy, the polynomial $`B_m`$ of $`_m`$ is of the form
$$B_m(x_0,x_1,x_2)=x_0^{m3}.B_3(x_0,x_1,x_2)+C_m(x_0,x_1,x_2)$$
with $`m7`$ where $`B_3`$ is the polynomial of a Harnack curve of degree $`3`$.
- One can deform the polynomial $`A_{m3,0}=x_0^{m3}`$ into the polynomial $`\mathrm{\Pi }_{j=1}^{m3}(x_{0,j})`$ of $`m3`$ lines in general position in such a way that the path $`B_{m,t}`$, $`t[0,t_0]`$, is rigid isotopy from $`B_{m,0}=B_m`$ to $`B_{m,t_0}`$.
$$B_{m,t_0}(x_0,x_1,x_2)=\mathrm{\Pi }_{j=1}^{m3}(x_{0,j}).B_3(x_0,x_1,x_2)+C_m(x_0,x_1,x_2)$$
- According to the Lemma 2.2.21 of Chapter 2.2.21, one can assume that the curve $`_3`$ with polynomial $`B_3`$ is immersed in $`_m`$ as the classical deformation of $`_2L`$ where $`L:=x_0=0`$.
Therefore, one can consider a rigid isotopy $`B_{m,t}`$, $`t[t_0,t_1]`$,
$$B_{m,t_1}=B_3.B_{m3}+C_m$$
where $`B_{m3}`$ is a regular polynomial of an $`L`$-curve which results from the deformation of $`\mathrm{\Pi }_{j=1}^{m3}(x_{0,j})`$.
Let $`B_3`$ be the norm of $`B_3`$ and let $`b_{m3}.B_3`$ be the affine polynomial associated to $`B_{m3}.B_3`$. Consider the pencil of curves $`x_0^m(b_{m3}.||B_3||c)`$.
For any critical point $`(x_0,x_1,x_2)`$ of $`B_{m3}.B_3`$ choose its representative in $`S^2`$ $`\frac{1}{(1+x_{0,1}^2+x_{0,2}^2)^{1/2}}(1,x_{0,1},x_{0,2})`$. As $`c`$ varies from $`[C_m,C_m]`$, for any critical point of $`b_{m3}`$ which goes through its critical value; the real point set of a curve of the pencil undergoes a Morse modification. Since $`_m`$ has the maximal number of real connected components a curve of degree $`m`$ may have, it is not hard to see that an $`M`$-curve of degree $`m3`$ is $`M`$-immersed in $`_m`$. Up to rigid isotopy of $`B_m`$, it is realized as the $`L`$-curve which results from the perturbation of the union of $`m3`$ lines with polynomial $`\mathrm{\Pi }_{j=1}^{m3}(x_{r,j})`$.
- It follows that there exists a rigid isotopy $`B_{m,t}`$, $`t[0,1]`$ from $`B_{m,0}`$ to
$$B_{m,1}=B_3.B_{m3}+C_{m,1}$$
where $`B_{m3}`$ is the polynomial of $`M`$-curve of degree $`m37`$; and the curve the curve $`_m`$ results from classical small deformation of $`_{m3}_3`$ with polynomial $`B_{m3}.B_3`$. Since any curve $`_{m3}`$ of degree $`m34`$ belongs to $`\mathrm{\Omega }_{m3}`$, it proves the Proposition 2.2.22 of Chapter 2.2.22 for curves $`_m`$ of degree $`m7`$. 2)
Let us consider curves $`_m`$ of degree $`m8`$. According to the Proposition 2.2.17 of Chapter 2.2.17, up to rigid isotopy, the polynomial $`B_m`$ of $`_m`$ is of the form
$$B_m(x_0,x_1,x_2)=x_0^{m3}.B_3(x_0,x_1,x_2)+C_m(x_0,x_1,x_2)$$
where $`B_3`$ is the polynomial of a Harnack of degree $`3`$ and $`m3=4.r+s`$ with $`r,s𝐍,s3`$.
Set $`A_{m3,0}=x_0^{m3}=x_0^{4r+s}=(x_0^4)^r.x_0^s=(x_0^4+\mathrm{}x_0^4).x_0^s`$. For any of the $`n^{th}`$ ($`1nr`$) product $`x_0^4`$ of the polynomial $`(x_0^4)^r=x_0^4+\mathrm{}x_0^4)`$ consider the deformation of $`x_0^4`$ into the polynomial $`\mathrm{\Pi }_{j=1}^4(x_{n,j})`$ of the union of four lines in generic position. In the same way, consider the deformation of $`x_0^s`$ into the polynomial $`\mathrm{\Pi }_{j=1}^s(x,r+1,j)`$ of $`s`$ lines in generic position.
- For any of these union of $`4`$ (resp, $`s4`$), lines in generic position, consider the rigid isotopy $`B_{m,t}`$, $`t[0,t_0]`$ from $`B_{m,0}=B_m`$ to $`B_{m,t_0}`$.
For $`1nr`$,
$$B_{m,t_0}(x_0,x_1,x_2)=x_0^{4(r1)+s}.\mathrm{\Pi }_{j=1}^4(x_{n,j}).B_3(x_0,x_1,x_2)+C_m(x_0,x_1,x_2)$$
(resp,
$$B_{m,t_0}(x_0,x_1,x_2)=x_0^{4r}.\mathrm{\Pi }_{j=1}^s(x_{r+1,j}).B_3(x_0,x_1,x_2)+C_m(x_0,x_1,x_2))$$
From an argumentation analogous to the one given in $`\mathrm{𝟏})`$, it follows that, up to regular deformation of $`B_m`$, any of the union of $`4`$ (resp, $`s`$) lines in generic position may be deformed into an $`L`$-curve of degree $`4`$ (resp, $`s`$) which is an $`M`$-curve immersed in $`_m`$. Denote by $`_{4,n}`$, $`1nr`$, (resp, $`_s`$) the corresponding curves. Denote by $`B_{4,n}`$ (resp, $`B_s`$) the polynomial of $`_{4,n}`$ (resp, $`_s`$). - It follows that there exists a rigid isotopy $`B_{m,t}`$, $`t[0,t_1]`$, from $`B_{m,0}=B_m`$ to $`B_{m,t_1}`$.
$$B_{m,t_1}(x_0,x_1,x_2)=\mathrm{\Pi }_{n=1}^r(B_{4,n}).B_s.B_3+C_{m,t_1}(x_0,x_1,x_2)$$
###### Definition 2.2.23.
Call succesive small perturbation of a finite union of $`_{i=1}^n𝒜_i`$ of curve $`𝒜_i`$ the result of the following recursive classical small perturbation of the union of two curves:
1. the union $`𝒜_1𝒜_2`$ is a singular curve all of whose singular points are crossings; the classical small perturbation $`𝒜_1𝒜_2`$ leads to a curve $`_2`$
2. For $`2in1`$, the classical small perturbation $`_i𝒜_{i+1}`$ leads $`_{i+1}`$ such that the union $`_{i+1}𝒜_{i+2}`$ is a singular curve all of whose singular points are crossings.
Let us prove that a curve $`𝒜_{m3}\mathrm{\Omega }_{m3}`$ of degree $`m3`$ which results from a successive classical small deformation of the union $`_{n=1}^r_{4,r}_s`$ is $`M`$-immersed in $`_m`$.
Let us consider the union of curves $`_{n=1}^r_{4,n}_s`$ with polynomial $`\mathrm{\Pi }_{n=1}^rB_{4,n}.B_s`$. Up to rigid isotopy of $`B_m`$, one can assume $`_{4,n}=_{4,1}`$, $`2nr`$, where $`_{4,1}`$ is the $`L`$-curve of degree $`4`$ which results from the deformation of the union of four lines in generic position with polynomial $`\mathrm{\Pi }_{j=1}^4x_{1,j}`$. One can also assume that $`_s`$, $`s4`$,is the $`L`$-curve of degree $`s`$ which results from the classical deformation of the union of $`s`$ lines with polynomial $`\mathrm{\Pi }_{j=1}^sx_{1,j}`$.
Since $`_m`$ is an $`M`$-curve, for any $`i,1ir`$, there exists an $`M`$-curve of degree $`4i+s`$ immersed in $`_m`$ as the result of the successive classical perturbation of $`_{k=1}^i_{4,k}_s`$. Denote by $`_{s+4i}`$ the $`M`$-curve of degree $`4+si`$ immersed in $`_m`$ as the classical perturbation of $`_{k=1}^i_{4,k}_s`$.
The curve $`_{4+s}`$ immersed in $`_m`$ belongs to $`\mathrm{\Omega }_{4+s}`$. Indeed, none of the pair of ovals $`(𝒪,𝒪^{})(𝐑_{4,1},𝐑_s)`$ is injective. In the same way, since $`_{4,n}=_{4,1}`$ for $`2nr`$ none of the pair of ovals $`(𝒪,𝒪^{})(𝐑_{s+4i},𝐑_{4,i+1})`$ is an injective pair of ovals. Hence, at any step of the successive classical small deformation we get an $`M`$-curve $`_{s+4i}\mathrm{\Omega }_{s+4i}`$.
For clarity, let us verify that, up to rigid isotopy of $`B_m`$, the successive classical deformation of $`_{k=1}^r_{4,k}_s`$ as curve immersed in $`_m`$ does not depend on the order of the deformations.
Let us notice that:
###### Remark 2.2.24.
The classical small deformation of the union of two $`M`$-curves $`𝒜`$ leads to an $`M`$ curve only if common points of $`𝒜`$ and $``$ belong to one connected real connected of $`𝒜`$ and one connected real connected of $``$.
proof: Otherwise, it would lead to contradiction with the number of real connected components of $`_m`$. Q.E.D
According to remark 2.2.24 of Chapter 2.2.24, there exists $`𝒪𝐑_{4,1}`$ such that any two curves immersed in $`_m`$ as the successive classical deformation $`_{n=1}^l_{4,n}`$ and respectively $`_{n=1}^l^{}_{4,n}`$, $`l\mathit{}^{},2l,l^{}r`$ (where $`_{4,2}_{4,1}`$) intersects $`_{4,1}`$ in $`𝒪`$. In the same way, there exists $`𝒪𝐑_s`$ (if $`s`$ is odd, $`𝒪=J`$ the odd component of $`_s`$) such that any two curves immersed in $`_m`$ as the successive classical deformation $`_{n=1}^l_{4,n}`$ and $`_{n=1}^l^{}_{4,n}`$, $`ll^{},2l,l^{}r`$ (where $`_{4,2}_{4,1}`$) intersect $`_s`$ in $`𝒪`$.
Hence, the $`M`$-curve $`𝒜_{m3}`$ immersed in $`_m`$ as the result of the successive classical small deformation of $`_{k=1}^r_{4,k}_s`$ is well defined. Besides, it belongs to $`\mathrm{\Omega }_{m3}`$.
- It follows that there exists a rigid isotopy $`B_{m,t}`$, $`t[0,1]`$ from $`B_{m,0}`$ to
$$B_{m,1}=B_3.B_{m3}+C_{m,1}$$
where $`B_{m3}`$ is the polynomial of $`𝒜_{m3}\mathrm{\Omega }_{m3}`$ and $`_m`$ results from classical small deformation of the curve $`𝒜_{m3}_3`$ with polynomial $`B_{m3}.B_3`$.
It concludes the proof of Proposition 2.2.22 of Chapter 2.2.22. Q.E.D
In the subsection 2.2 of Chapter 2.2, we have proven that if $`_n`$ is of type $`^0`$, then, up to regular modification of its polynomial $`B_n`$, the polynomial $`B_{n1}`$ is also of type $`^0`$. Let us now in Proposition 2.2.25 of Chapter 2.2.25 prove the converse.
###### Proposition 2.2.25.
Let $`_n`$ be a Harnack curve of degree $`n`$, $`n4`$, and $`B_n(x_0,x_1,x_2)=x_0^n.b_n(x_1/x_0,x_2/x_0)`$ its polynomial. Denote by $`B_{n1}(x_0,x_1,x_2)=x_0^{n1}.b_{n1}(x_1/x_0,x_2/x_0)`$ the homogeneous polynomial associated to the truncation $`b_{n1}(x,y)`$ of $`b_n(x,y)`$ on the monomials $`x^\alpha .y^\beta `$ with $`0\alpha +\beta n1`$.
If $`B_{n1}(x_0,x_1,x_2)`$ is of type $`^0`$ relatively to the line $`L_n:=x_0=0`$, then up to regular modification of $`B_n`$, there exists a line $`L_{n+1}`$ with the property that $`_n`$ is of type $`^0`$ relatively to $`L_{n+1}`$.
proof :
Our proof uses an argumentation similar to the one given in the subsection 2.2 of Chapter 2.2 to prove that if $`_n`$ is of type $`^0`$ then up to regular modification of its polynomial $`B_n`$, the curve $`_{n1}`$ is also of type $`^0`$.
Let $`L_n:=x_0=0`$ and $``$ be a neighborhood of the line $`L_n`$. The curve $`_n`$ with polynomial $`B_n`$ results from deformation of $`_{n1}L_n`$
$$B_n=B_{n1}(x_0,x_1,x_2).x_0+C_n(x_0,x_1,x_2)$$
where $`C_n(x_0,x_1,x_2)`$ is homogeneous of degree $`n`$. Let $`U(s)𝐑P^2`$ be a neighborhood of a singular point of $`𝐑_{n1}𝐑L_n`$ such that $`𝐑_nU(s)\mathrm{}`$. According to the topology of $`_n`$, for each singular point $`s`$ there exists a homeomorphism $`h:U(s)D^1\times D^1`$ such that $`h(𝐑_{n1}L_n)=D^1\times 00\times D^1`$ and $`h(𝐑_nU(s))=\{(x,y)D^1\times D^1|xy=1/2\}`$. Since the polynomial of $`_n`$ is smooth, it follows by continuity, that real connected components of $`𝐑_n`$ define entirely and uniquely the homomorphism $`h`$.
We shall enlarge the previous description of $`\{x𝐑P^2|B_n(x)=0\}`$ in $`U(s)`$ to a description of in $`U_s`$ where
$$U_s=\{z=<u,p>=(u_0.p_0:u_1.p_1:u_2.p_2)𝐂P^2|u=(u_0:u_1:u_2)U_𝐂^3,p=(p_0:p_1:p_2)U(s)\}$$
To this end, let us prove the following Lemma 2.2.26 of Chapter 2.2.26.
###### Lemma 2.2.26.
Let $`B_n(x_0,x_1,x_2)=x_0^n.b_n(x_1/x_0,x_2/x_0)`$ be a Harnack polynomial of degree $`n`$. Denote by $`B_{n1}(x_0,x_1,x_2)=x_0^{n1}.b_{n1}(x_1/x_0,x_2/x_0)`$ the homogeneous polynomial associated to the truncation $`b_{n1}(x,y)`$ of $`b_n(x,y)`$ on the monomials $`x^\alpha .y^\beta `$ with $`0\alpha +\beta n1`$.
Assume that $`B_{n1}`$ is of type $`^0`$ relatively to $`L:=x_0=0`$.
1. there exists $``$ neighborhood of $`L`$ such that the polynomial $`B_n^\mathrm{\Delta }(x_0,x_1,x_2)=(x_0^n.b_n^\mathrm{\Delta }(x_1/x_0,x_2/x_0)`$ where $`b_n^\mathrm{\Delta }`$ is the truncation of $`b_n`$ on monomials of homogeneous degree $`n2in`$ is $`ϵ`$-sufficient for $`B_n(x_0,x_1,x_2)`$.
2. Up to regular modification of $`B_n`$, there exist $`_ϵ`$, $`_ϵ=_{s𝐑_{m1}L}U_s`$ $`ϵ`$-tubular neighborhood of $`L`$ and a polynomial $`\stackrel{~}{B}_n`$ with the following properties:
-$`\stackrel{~}{B}_n`$ is $`ϵ`$-sufficient for $`B_n`$ in $`_ϵ`$
-for any $`U_s`$, $`U_s_ϵ`$, the truncation $`\stackrel{~}{B}_n^S`$ of $`\stackrel{~}{B}_n^S`$ (on four monomials $`x^cy^dz^2,x^{c+1}y^dz,x^cy^{d+1}z,x^{c+1}y^{d+1}`$ with $`c+d+2=n`$) is $`ϵ`$-sufficient for $`\stackrel{~}{B}_n`$ (and thus for $`B_n`$) in $`U_s`$. (For any $`ss^{}`$, such that the truncation of $`\stackrel{~}{B}^n`$ on $`x^cy^dz^2,x^{c+1}y^dz,x^cy^{d+1}z,x^{c+1}y^{d+1}`$ (resp,on $`x^c^{}y^d^{}z^2,x^{c^{}+1}y^d^{}z,x^c^{}y^{d^{}+1}z,x^{c^{}+1}y^{d^{}+1}`$) is $`ϵ`$-sufficient for $`\stackrel{~}{B}_n`$ (and thus for $`B_n`$) in $`U_s`$ (resp, in $`U(s^{})`$, $`(c,d)(c^{},d^{})`$ )
###### Remark 2.2.27.
Note that if the curve $`_n`$ is not of type $`^0`$ (relatively to $`L`$), then $`_ϵ`$, $`_ϵ`$.
proof: Assume $`B_{n1}`$ of type $`^0`$ relatively to $`L_n:=x_0=0`$. According to the proof of Morse Lemma, local coordinates defined in a neighborhood $`U(s)`$ of a point $`s`$ depend principally on the first derivative and the second derivative of the function $`b_n`$ around this point. Hence, the truncation of $`b_n^\mathrm{\Delta }(x,y)`$ on monomials of homogeneous degree $`n2in`$ is $`ϵ`$-sufficient for $`b_n(x,y)`$.
According to Morse Lemma , in neighborhood $`U(s)`$ of any non-degenerate singular point $`s`$ of $`x_0.B_{n1}`$ one can choose local coordinates system $`z_1,z_2`$ with $`z_1(s)=0,z_2(s)=0`$ and $`x_0.B_{n1}=z_1.z_2`$. Real connected components of $`𝐑_n`$ define entirely and uniquely the homomorphism $`h`$ defined in $`U(s)`$ $`h:U(s)D^1\times D^1`$ such that $`h(𝐑_{n1}L_n)=D^1\times 00\times D^1`$ and $`h(𝐑_nU(s))=\{(x,y)D^1\times D^1|xy=1/2\}`$. In such a way, in a neighborhood $`U(s)`$ of a point $`s`$, one can choose local coordinates system $`y_1,y_2`$ with $`y_1(s)=0,y_2(s)=0`$; $`B_n=y_1.y_2+B_n(s)`$. In $`U(s)`$, any point $`s`$ is a local extremum of $`b_n`$. Any singular point $`s`$ of $`𝐑B_{n1}𝐑L_n`$ belongs to the line $`x_0=0`$, hence it is a local extremum of the function $`b_n(y)=b_n(0,x_1/x_2,1)`$. One can assume without loss of generality, that any point $`s`$ does not belong to the line $`x_2=0`$. On such assumption gradient trajectories of $`b_n(x_0,x_1,1)`$ in $`U(s)`$ give the direction of the perturbation of the crossing $`s`$. Thus, up to regular modification of $`B_n`$ and for sufficiently small $`ϵ>0`$, there exist $`U(s)_ϵ`$ and a polynomial $`\stackrel{~}{b}_n^S`$ on four monomials ($`x^cy^d,x^{c+1}y^d,x^cy^{d+1},x^{c+1}y^{d+1}`$ with $`c+d+2=n`$) which is $`ϵ`$-sufficient for $`B_n(1,x_1,x_2)`$ in $`U(s)_ϵ`$; $`\stackrel{~}{b}_n^S(x,y)=l(x,y)+k(x,y)`$ with $`l(x,y)=a_{c,d}x^cy^d+a_{c+1,d}x^{c+1}y^d`$
$`k(x,y)=a_{c,d+1}x^cy^{d+1}+a_{c+1,d+1}x^{c+1}y^{d+1}`$ with $`c+d+2=n`$. From this last observation, we deduce the definition of $`\stackrel{~}{b}_n`$. Q.E.D Let us consider the polynomial $`\stackrel{~}{b}_n`$ with truncation $`\stackrel{~}{b}_n^S`$ $`ϵ`$-sufficient for $`B_n(1,x_1,x_2)`$ in $`U_s`$.
$`\stackrel{~}{b}_n^S(x,y)=l(x,y)+k(x,y)`$ with $`l(x,y)=a_{c,d}x^cy^d+a_{c+1,d}x^{c+1}y^d`$
$`k(x,y)=a_{c,d+1}x^cy^{d+1}+a_{c+1,d+1}x^{c+1}y^{d+1}`$ with $`c+d+2=n`$.
Up to modify the coefficients $`a_{c,d},a_{c,d+1},a_{c+1,d},a_{c+1,d+1}`$ if necessary, the point $`s=(x_0,y_0)`$ is, up to homeomorphism, a critical point of the function $`\frac{l(x,y)}{k(x,y)}`$. Hence, it follows from the equalities $`l(x,y)=l(x,y)`$, $`k(x,y)=k(x,y)`$,
$`\frac{l}{x}(x,y)=\frac{l}{x}(x,y)`$, $`\frac{l}{y}(x,y)=\frac{l}{y}(x,y)`$,
$`\frac{k}{x}(x,y)=\frac{k}{x}(x,y)`$, $`\frac{k}{y}(x,y)=\frac{k}{y}(x,y)`$
that $`(x_0,y_0)`$ is also a critical point of the function $`\frac{l(x,y)}{k(x,y)}`$. In such a way, we define a map which maps the set $`𝒮`$ of singular points of $`_{n1}L_n`$ to a set $`𝒫`$ of $`n1`$ real points. In a neighborhood $`U(p)`$ of $`p𝒫`$, $`𝐑_m`$ is desingularized crossing i.e consists of two real arcs $`\gamma ,\gamma ^{}`$. The set $`𝒮`$ is a set of aligned points. Hence, up to regular modification, the set $`𝒫`$ is also a set of aligned points which belong to a line $`L`$. Let $`_ϵ=^+^{}`$ ,$`^\pm =L_1^\pm L_2^\pm `$, be a neighborhood of the line $`L`$; $`_{p𝒫}U(p)`$.
In the case $`n=2k`$, any pair of arcs of $`𝐑_{2k}U(p)`$, $`p𝒫`$, is such that one arc belongs to the non-empty oval of $`𝐑_{2k}`$. The non-empty oval is connected. Hence, it intersects a projective line $`L_1`$ in $`2k`$ points.
In the case $`n=2k+1`$, any pair of arcs of $`𝐑_{2k+1}U(p)`$, $`p𝒫`$, is such that at least one arc belongs to the odd component of $`𝐑_{2k+1}`$. The odd component is connected. Hence, it intersects a projective line $`L_1`$ in $`2k+1`$ points.
In such a way, we get that if $`_{n1}`$ is of type $`^0`$ relatively to $`L_n`$, there exists $`L_1:=L_{n+1}`$ such that the curve $`_n`$ is of type $`^0`$ relatively to $`L_1=L_{n+1}`$. This concludes our proof of Proposition 2.2.25 of Chapter 2.2.25. Q.E.D
Let us prove in Theorem 2.2.28 of Chapter 2.2.28 , that up to regular modification of the Harnack polynomial $`B_m(x_0,x_1,x_2)`$, there exists a sequence of lines $`L_{i+1}`$, $`1im`$, such that $`B_i`$ is of type $`^0`$ relatively to the line $`L_{i+1}`$ and describe relative position of these lines. We shall call rotation of $`𝐑P^2`$ the transformation of $`𝐑P^2`$ which extends continuously a usual rotation of $`𝐑^2`$. Relative position of lines $`L_i`$, $`1im+1`$ is described by means of rotations of $`𝐑P^2`$.
###### Theorem 2.2.28.
Let $`_m`$ be a Harnack curve of degree $`m`$. Up to rigid isotopy, there exists a sequence of lines $`_1=L_1`$, $`L_2`$, …, $`L_{m+1}`$ such that $`_i`$ is of type $`^0`$ relatively to $`L_{i+1}`$ and $`_{i+1}`$ results from the deformation of $`_iL_{i+1}`$.
The sequence of lines $`L_i`$ may be ordered in three different ways:
1. There exist $`pL_2L_j`$ , $`3jm+1`$, and a rotation of center $`p`$ and angle $`\theta `$ (where $`\theta `$ may be chosen arbitrarily small) which maps $`𝐑L_i`$ to $`𝐑L_{i+1}`$.
2. There exists $`pL_2L_j`$ , $`3jm+1`$, such that:
for $`i`$, $`2i2.[(m+1)/2]2`$ the rotation of center $`p`$ and angle $`\theta `$ maps $`𝐑L_i`$ to $`𝐑L_{i+2}`$. for $`i`$, $`1i2.[(m+1)/2]2`$ the rotation of center $`p`$ and angle $`\theta `$ maps $`𝐑L_{i+1}`$ to $`𝐑L_{i+3}`$ ($`\theta `$ may be chosen arbitrarily small)
\- \[m/2\] denotes the integer $`n`$, $`(m1)/2nm/2`$ -
3. There exists $`pL_2L_j`$ , $`3jm+1`$, such that:
for $`i`$, $`2il`$, arrangement of lines $`L_{i+1}`$ is given by 2
for $`i`$, $`lim+1`$ arrangement of lines $`L_{i+1}`$ is given by 1
proof: In a first part, we shall prove Theorem 2.2.28 of Chapter 2.2.28 for curves $`_m`$ of even degree $`m`$. Then, we shall extend our argumentation to Harnack curves of odd degree.
We shall use the rigid isotopy classification of Harnack curves of type $`^0`$ given in Theorem 2.2.2 of Chapter 2.2.2. The following remark may be easily deduced from the Harnack’s construction of curves. (see , , tables of isotopy types of Harnack $`M`$-curves)
###### Remark 2.2.29.
The Harnack curve $`_{n+1}`$ is the only $`M`$-curve of degree $`n+1`$ with:
-for even $`n+1`$ only one non-empty oval
\- for odd $`n+1`$ only empty ovals
which results from the classical deformation of the union $`_nL_{n+1}`$ of a Harnack $`_n`$ of type $`^0`$ relatively to a line $`L_{n+1}`$.
I) Harnack Curves of even degree $`_{2k}`$
Our proof is based on the construction of a sequence of curves $`𝒜_{2i}\mathrm{\Omega }_{2i}`$, $`1ik1`$ such that the deformation of $`𝒜_{2i}_{2k2i}`$ leads to the curve $`_{2k}`$.
According to Proposition 2.2.22 of Chapter 2.2.22, one can assume that $`_{2k}`$ results from the deformation of $`𝒜_{2k2}_2`$ where $`𝒜_{2k2}\mathrm{\Omega }_{2k2}`$. It follows from the Bezout’s theorem Morse Lemma and the real components of $`𝐑_{2k}`$ that the $`(2k2).2`$ common points of $`𝒜_{2k2}`$ and $`_2`$ belong to their non-empty ovals. We shall consider connected arcs $`I_j`$, $`3j2k`$ of the non-empty oval of $`𝒜_{2k2}`$; $`I_j𝐑_2𝒜_{2k2}`$ and identify these arcs $`I_j`$ with lines $`𝐑L_j`$. In this way, we shall construct a sequence of curves $`_j`$ of type $`^0`$ relatively to $`L_{j+1}`$ with the property that curves $`_{2k2i}`$ are such that $`𝒜_{2i}_{2k2i}`$ leads to the curve $`_{2k}`$.
We shall call extreme point a point of the intersection $`𝐑𝒜_{2k2}𝐑_2`$ the two points connected by an arc of $`𝐑𝒜_{2k2}`$ which does not contain other points $`𝐑𝒜_{2k2}𝐑_2`$. We shall call extreme arc an arc which contains an extreme point. Let us detail the construction of the sequences of curves $`𝒜_{2i}\mathrm{\Omega }_{2i}`$, $`1ik1`$. Any Harnack curves $`_2`$ are rigidly isotopic. Therefore, without changing relative position of real connected components of $`𝒜_{2k2}`$ and $`_2`$, we may assume $`_2`$ of type $`^0`$ relatively to an (extreme) arc $`I_3𝐑L_3`$ of $`𝒜_{2k2}`$. (Choose $`I_3`$ of small lenght. up to glue its extremities, we may identify $`I_3`$ with a projective line $`𝐑L_3`$.) In such a way, according to the remark 2.2.29 of Chapter 2.2.29, the deformation of $`I_3𝐑_2`$ leads necessarily to the Harnack curve $`_3`$. (Otherwise, it would lead to contradiction with the real connected components of $`_{2k}`$.) Moreover, it follows from Proposition 2.2.25 of Chapter 2.2.25, that the curve $`_3`$ is of type $`^0`$ relatively to one line. According to Theorem 2.2.2 of Chapter 2.2.2 and remark 2.2.16 of Chapter 2.2.16, we may, up to regular modification, identify (without changing relative position of real connected components of $`𝐑𝒜_{2k2}`$ and $`𝐑_2`$) this line with an other arc $`I_4`$ of the non-empty oval of the curve $`𝒜_{2k2}`$. The arc $`I_4`$ contains, besides a set of $`2`$ consecutive points which are up to rigid isotopy points of $`𝐑𝒜_{2k2}𝐑_2`$, an other point. Identifying arcs with lines $`I_3𝐑L_3`$, $`I_4𝐑L_4`$, this last point is provided by the intersection $`L_3L_4`$. From an argument analogous to the previous one, we get a curve $`_4`$ of type $`^0`$ relatively to an arc $`I_5𝐑L_5`$. By induction, it follows that $`_{2k}`$ results from the deformation of $`_4𝒜_{2k4}`$ where $`𝒜_{2k4}\mathrm{\Omega }_{2k4}`$.
Iterating our argumentation, we get a sequence of curves $`𝒜_{2i}\mathrm{\Omega }_{2i}`$, $`1ik1`$, such that the deformation of $`𝒜_{2i}_{2k2i}`$ leads to the curve $`_{2k}`$. We also define a set of arcs $`I_{j+1}`$,$`2j(2k1)`$, with the property that the curve $`_j`$ is of type $`^0`$ relatively to $`I_{j+1}𝐑L_{j+1}`$. Arcs $`I_j`$, $`3j2k`$, belong to the non-empty oval of $`𝒜_{2k2}`$; the set of arcs $`_jI_j`$, $`2k2i+1j2k`$ is mapped to a set of arcs of $`𝒜_{2i}`$, $`1ik1`$.
At the final step, we get $`_{2k}`$ from the deformation of $`𝒜_2_{2k2}`$.
In the construction of the sequence of curves $`𝒜_{2i}`$, we have to take into account the way $`_{2k2i}`$ intersects $`𝒜_{2i}`$. Curves $`𝒜_{2i}`$ may intersect curves $`_{2k2i}`$ in only two different ways: the non-empty oval of one curve contains all inner ovals of the other curve or it contains none of them. Let us first consider this last case. 1) Assume that the non-empty oval of $`𝒜_{2i}`$ does not contain ovals of the non-empty oval of $`_{2k2i}`$. The curve $`_{2k}`$ results from the deformation of the union $`_{2k2i}𝒜_{2i}`$. In this case, any arc $`I_j`$, $`2k2i+1j2k`$, may be identified with the same line $`𝐑L_{2k2i+1}`$. Thus, from step of the construction of $`_{2k2i+1}`$ the method proposed for the construction of $`_{2k}`$ is analogous to the Harnack’s construction. Since arcs $`I_j`$, $`2k2i+1j2k`$, are distinct, one can say equivalently -identifying $`I_j`$ with lines $`𝐑L_j`$\- that there exists $`pL_{2k2i+1}L_j`$ , $`2k2i+1jm+1`$, and a rotation of center $`p`$ and angle $`\theta `$ which maps $`𝐑L_i`$ to $`𝐑L_{i+1}`$. For simplicity, we choose the sequence of arcs $`I_j`$, $`2k2i+1j2k`$ as follows. The arc $`I_{j+1}`$ intersects the arcs $`I_j`$ and $`I_{j+2}`$; $`I_{2k}`$ is extreme. 2)Let us now assume that the non-empty oval of the curve $`_2`$ contains all inner ovals $`𝒜_{2k2}`$. In the case 1), we have chosen arcs $`I_j`$, $`3j2k`$ such that an arc intersects at most an other arc in one point. We need here one more assumption on the choice of arcs $`I_j`$ of the non-empty oval of $`𝒜_{2k2}`$.
Assume that for any $`i,1i<k`$, the non-empty oval of the curve $`𝒜_{2k2i}`$ does not contain inner ovals of $`_{2i}`$. Let $`I_3`$ and $`I_4`$ be the two extreme arcs of $`𝒜_{2k2}`$. Then, we get $`_4`$ and a curve $`𝒜_{2k4}\mathrm{\Omega }_{2k4}`$ with the property that $`_{2k}`$ results from the deformation of $`_4𝒜_{2k4}`$. We iterate this process. Consider the union $`I_j`$, $`3j2k`$, of arcs of $`𝒜_{2k2}`$ as follows: $`I_3`$ and $`I_4`$ are extreme arcs of $`𝒜_{2k2}`$ Up to rigid isotopy, arcs $`I_j`$, $`j>5`$, are arcs of the non-empty oval empty of $`𝒜_{2k4}`$. We choose $`I_5`$ and $`I_6`$ such that $`I_5`$ intersects $`I_3`$ in one point, $`I_6`$ intersects $`I_4`$ in one point. In such a way, they are mapped to extreme arcs of $`𝒜_{2k4}`$. By induction, we choose $`I_{3+2l}`$, $`I_{4+2l}`$, $`1l(k1)`$ such that $`I_{3+2l}`$ intersects $`I_{3+2(l1)}`$ in one point, $`I_{4+2l}`$ intersects $`I_{4+2(l1)}`$ in one point; and are extreme arcs of $`𝒜_{2k2(l+1)}`$. At the end of the process, we get $`_{2k}`$ from the deformation of $`𝒜_2_{2k2}`$ where $`𝒜_2=_2`$. (Identifying arcs $`I_j`$ with lines, $`𝐑L_{2k1}I_{2k1}`$, $`𝐑L_{2k}I_{2k}`$ the curve $`𝒜_2=_2`$, as $`L`$-curve, results from the perturbation of the union $`L_{2k1}L_{2k}`$.)
It is an easy property of the proposed construction that the oval of $`𝒜_2=_2`$ contains besides inner ovals of $`_{2k2}`$ inner ovals of $`_{2k}`$ in its inner component. Iterating this process, we get $`_{2k2i}`$, $`i<k`$, from $`_{2k2i2}𝒜_2`$ where the oval of $`𝒜_2=_2`$ contains, besides inner ovals of $`_{2k2i2}`$ inner ovals of $`_{2k2i}`$ in its inner component. (Identifying arcs $`I_j`$ with lines, the curve $`𝒜_2=_2`$, as $`L`$-curve, results from the perturbation of two lines $`𝐑L_{2k2i1}I_{2k2i1}`$, $`𝐑L_{2k2i}I_{2k2i}`$ in general position.)
If there exists $`jk`$ such that the non-empty oval of $`𝒜_{2k2j}`$ does not contain inner ovals of $`_{2j}`$, then from the step when the deformation $`_{2j}𝒜_{2k2j}`$ leads to $`_{2k}`$ to the final step \- the deformation of $`𝒜_2_{2k2}`$ leads to $`_{2k}`$ \- , we are in the case 1).
In the previous argumentation, we have assumed that the non-empty oval of $`𝐑𝒜_{2k2}`$ is the union of an arc which does not contain points of $`𝐑𝒜_{2k2}𝐑_2`$. with a set of arcs $`I_j`$, $`3j2k`$ which contains points of $`𝐑𝒜_{2k2}𝐑_2`$ From identification of these $`2k2`$ arcs $`I_j`$ with lines $`𝐑L_j`$, $`3j2k`$, we have described components of $`𝐑𝒜_{2k2}`$ as an arrangement of pieces of these $`2k2`$ lines. The degree of $`𝒜_{2k2}`$ is $`2k2`$. Hence, its real components may not result from more than $`2k2`$ lines. Moreover, since $`𝒜_{2k2}`$ is an $`M`$-curve, its real components may not result from less than $`2k2`$ lines. Our initial assumption is therefore always possible. Moreover, one can assume that $`_{2k}`$ results from the deformation of $`𝒜_{2k2}_2`$ where $`𝒜_{2k2}\mathrm{\Omega }_{2k2}`$. and $`𝐑𝒜_{2k2}`$ is entirely covered by the set of arcs $`_{j=3}^{2k}I_j`$. Hence, since $`_2`$ is, up to rigid isotopy, an $`L`$-curve; we get Theorem 2.2.28 of Chapter 2.2.28 for Harnack curves $`_{2k}`$ of even degree. II) Harnack Curves of odd degree $`_{2k+1}`$ According to Proposition 2.2.17 of Chapter 2.2.17, up to regular modification of its Harnack polynomial $`B_{2k+1}`$, one can assume that $`_{2k+1}`$ results from the deformation of the union $`𝒜_{2k2}_3`$, where $`𝒜_{2k2}\mathrm{\Omega }_{2k2}`$ and $`_3`$ is the Harnack curve of degree $`3`$. Curves $`_3`$ are rigidly isotopic. One can consider $`_3`$ as an $`L`$-curve i.e as the result of the perturbation of three real lines $`L_0`$, $`L_1`$,$`L_2`$ in general position. According to remark 2.2.16 of Chapter 2.2.16, we shall assume lines $`L_0`$,$`L_1`$, $`L_2`$ chosen as follows. Outside a neighborhood $``$ of singular points $`L_0L_1L_2`$, there exists a rigid isotopy $`j_t`$ of $`𝐑P^2`$ which pushes $`_{i=0}^2𝐑L_i\backslash `$ to $`𝐑_3\backslash `$ such that $`𝐑_3`$ intersects $`𝐑𝒜_{2k2}`$ only its part $`j_1(_{i=0}^2𝐑L_i\backslash )\backslash j_1(𝐑L_0)`$. Consider the deformation $`_2𝒜_{2k2}`$ where $`_2`$ is an $`L`$-curve which results from the perturbation of the two real lines $`L_1`$,$`L_2`$ in general position. Combining this construction of curves $`_{2k+1}`$ and our previous study of curves of even degree, we obtain the Theorem 2.2.28 of Chapter 2.2.28 for curves $`_{2k+1}`$. Q.E.D To prove the rigid isotopy Theorem 2.2.1 of Chapter 2.2.1 it is sufficient to note that according to Proposition 2.2.25 of Chapter 2.2.25 and the Theorem 2.2.28 of Chapter 2.2.28 and its proof, any curve $`_m`$ is of type $`^0`$. Thus, the rigid isotopy Theorem 2.2.1 of Chapter 2.2.1 is a straightforward consequence of the Theorem 2.2.2 of Chapter 2.2.2.
However, for clarity, we propose to prove the Theorem 2.2.30 of Chapter 2.2.30.
###### Theorem 2.2.30.
Up to rigid isotopy, any curve $`_m`$, $`m1`$, results from the successive classical small deformation of the union $`L_{i+1}_i`$ , $`1in1`$, where $`_i`$ is a Harnack curve of degree $`i`$, $`L_{i+1}`$ is a projective line and the curve $`_i`$ is of type $`^0`$ relatively to $`L_{i+1}`$.
Let us give in Proposition 2.2.31 of Chapter 2.2.31 an other formulation of the Theorem 2.2.30 of Chapter 2.2.30.
###### Proposition 2.2.31.
Let $`_m`$ be a Harnack curve of degree $`m`$ and $`B_m(x_0,x_1,x_2)`$ its polynomial.
Denote by $`b_m(x,y)`$ the affine polynomial associated to $`B_m(x_0,x_1,x_2)`$,
$$B_m(x_0,x_1,x_2)=x_0^m.b_m(x_1/x_0,x_2/x_0)$$
Denote by $`b_i(x,y)`$ the truncation of $`b_m(x,y)`$ on the monomials $`x^\alpha .y^\beta `$ with $`0\alpha +\beta i`$ and by $`B_i(x_0,x_1,x_2)=x_0^i.b_i(x_1/x_0,x_2/x_0)`$ the homogeneous polynomial associated to $`b_i`$. Then, up to linear change of projective coordinates, and up to slightly perturb coefficients of $`B_m(x_0,x_1,x_2)`$ without changing either order or topological structure of $`_m`$, we may assume that any polynomial $`B_i`$ is of type $`^0`$ relatively to $`x_0=0`$.
proof: According to Theorem 2.2.28 of Chapter 2.2.28, up to rigid isotopy, there exists a sequence of lines $`_1=L_1`$, $`L_2`$, …, $`L_{m+1}`$ such that $`_i`$ is of type $`^0`$ relatively to $`L_{i+1}`$ and $`_{i+1}`$ results from the deformation of $`_iL_{i+1}`$. The sequence of lines $`L_i`$ may be ordered in three different ways described in (1), (2) and (3). To prove to Theorem 2.2.30, it is sufficient to prove that the description (1) is rigidly isotopic to the description (2).
Without loss of generality, one can assume that points $`_iL_{i+1}`$ do not belong to the line at infinity of $`𝐑P^2`$. Any affine line $`l_j=𝐑^2L_j`$ is divided into two halves $`l_j^+`$ $`l_j^{}`$ with common point $`p`$. (The rotation of center $`p`$ maps a positive (resp, negative) half of one line to a positive (resp, negative) half of an other line.) Consider the half plane $`𝐑^+`$ (resp, $`𝐑^{}`$) which contains positive (resp, negative) halves lines. Let $`H`$ be the common line of $`𝐑^+`$ and $`𝐑^{}`$ (where $`HL_j=p`$, for any $`L_j`$, $`1jm`$). The symmetry $`s_p`$ of center $`p`$ exchanges halves of lines. Let us denote $`s_H`$ the orthogonal reflection with respect to $`H`$. Assume, as described in (2), that there exists $`pL_2L_j`$ , $`3jm+1`$, such that:
\- for $`i`$, $`2i2.[(m+1)/2]2`$ the rotation of center $`p`$ and angle $`\theta `$ maps $`𝐑L_i`$ to $`𝐑L_{i+2}`$.
\- for $`i`$, $`1i2.[(m+1)/2]2`$ the rotation of center $`p`$ and angle $`\theta `$ maps $`𝐑L_{i+1}`$ to $`𝐑L_{i+3}`$ It follows from the remark 2.2.16 that without loss of generality, one can assume that $`_i`$ intersects $`L_{i+1}`$ in its half $`l_{i+1}^+`$. Choose the lines $`L_j`$ such that the symmetry $`s_H`$ exchanges lines $`L_{2l1}`$ and $`L_{2l}`$, $`s_H(L_{2l1})=L_{2l}`$,$`s_H(L_{2l})=L_{2l1}`$ According to remark 2.2.16 of Chapter 2.2.16, the rigid isotopy is trivially preserved under the action of $`s_H`$.
Let us now choose lines $`L_j`$ such that the rotation of center $`p`$ and angle $`\pm \frac{\theta }{2}`$ maps $`s_H(L_{2l})`$ to $`L_{2l\pm 1}`$. Obviously, the description of $`_m`$ which results from this arrangement of lines $`L_j`$ and the one where $`s_H(L_{2l1})=L_{2l}`$,$`s_H(L_{2l})=L_{2l1}`$ are rigidly isotopic. The rigid isotopy is preserved by the symmetry $`s_p`$ which maps the intersection of $`_i`$ with $`L_{i+1}`$ (namely, with $`l_{i+1}^+`$) to a set of the half $`l_{i+1}^{}`$. The rigid isotopy is also preserved when consider $`s_H(L_{2j})=L_{2j}^{}`$ and $`L_{2j1}`$ i.e $`l_{2j}^{}_{}{}^{}+`$ and $`l_{2j1}^+`$ instead of halves lines $`l_{2j1}^{}`$, $`l_{2j}^{}`$ In this way, we get that the description (2) is rigidly isotopic to the description (1). This proves the Theorem 2.2.30 of Chapter 2.2.30. Q.E.D
Harnack curves $`_m`$ constructed from the Harnack’s method are rigidly isotopic. It follows that the Theorem 2.2.30 of Chapter 2.2.30 implies the Rigid Isotopy Classification Theorem 2.2.1 of Chapter 2.2.1.
Appendix
Curves $`_m`$ of type $`^0`$ may result from the Harnack’s inductive construction (, ) as follows.
Given a projective line $`L`$, each curve $`_{i+1}`$, $`1im1`$, results from the classical small perturbation of the union $`_iL`$ where $`_i`$ is of type $`^0`$ relatively to $`L`$; $`_1`$ is a projective line which intersects $`L`$ in one point.
We take the auxiliary curve $`𝒞_{i+1}`$ which perturb the union $`_iL`$ to be a union of $`i+1`$ lines. Consider $`L`$ as the line at infinity of $`𝐑P^2`$. From the Harnack’s method initiated with $`L`$, one can obtain isotopy of curves different from $`_m`$ and curves $`_m`$ which are not of type $`^0`$ relatively to $`L`$. The isotopy type of curves (projective and affine) which result from the Harnack’s method depends on the choice of auxiliary curves.
Let us prove in Theorem 2.2.32 of Chapter 2.2.32 that isotopy implies also rigid isotopy for Harnack curves of type $`^0`$ obtained from this method. For curves of degree $`m6`$, the Theorem is obvious.
###### Theorem 2.2.32.
Any two Harnack curves $`_m`$ of type $`^0`$ constructed from the Harnack’s method are rigidly isotopic.
proof: We shall proceed by ascending induction on the degree.
The Theorem 2.2.32 of Chapter 2.2.32 is trivial for Harnack curves of degree $`_i`$, $`i6`$. Let us start our induction with Harnack curves of degree $`1`$. Let $`L_1`$ and $`L_2`$ be two projective lines, and $`_1^1`$, $`_1^2`$ be Harnack curves of degree $`1`$. Consider the classical small perturbation of $`_1^1L_1`$ and $`_1^2L_2`$. Denote by $`_2^1`$,(resp, $`_2^2`$), the Harnack curve of degree $`2`$ which results from the classical deformation of $`_1^1L_1`$ (resp, $`_1^2L_2`$).
Any two real projective lines are rigidly isotopic. Let $`h_t`$ be the rigid isotopy $`[0,1]𝐑𝒞_1\backslash 𝐑𝒟_1`$ ,$`h_0(𝐑L_1)=(𝐑L_1)`$, $`h_1(𝐑L_1)=𝐑L_2`$. Denote by $`𝐑L_{t+1}`$ the line $`h_t(𝐑L_1)`$. Any two projective lines intersect each other in one point. Hence, one can assume, without loss of generality that $`_1=_2=L_0`$. Along the rigid isotopy $`h_t`$, any line $`𝐑L_{t+1}=h_t(𝐑L_1)`$ intersects $`𝐑_0`$ in one point. Noticing that the classical perturbation of the union of two lines which intersect each other in one point leads to the Harnack curve $`_2`$, the path $`h_tid`$: $`[0,1]𝐑𝒟_2`$ extends to a path $`g_t:`$ $`[0,1]𝐑𝒞_2\backslash 𝐑𝒟_2`$, $`g(0)=𝐑_2^1`$, $`g(1)=𝐑_2^2`$ which is a rigid isotopy.
Given $`_m^1`$, (resp, $`_m^2`$), $`i2`$ the Harnack curve $`m`$ deduced from the Harnack’s method initiated with the line $`L_1`$, (resp, $`L_2`$). We shall prove that curves $`_m^1`$ and $`_m^2`$ are rigidly isotopic for arbitrary degree $`m`$.
Let us start with the following remark.
###### Lemma 2.2.33.
Let $`L`$ be a real projective line. Consider the Harnack’s inductive construction of curves initiated with the line $`L`$. Any two Harnack curves $`_m`$ of type $`^0`$ obtained from the Harnack’s method initiated with the line $`L`$ are rigidly isotopic.
proof: Indeed, consider the Harnack’s construction of the curve $`_m`$ initiated with the line $`L`$. In this construction, any curve $`_i`$, $`1im1`$, is of type $`^0`$ relatively to $`L`$. Consider $`L`$ as the line at infinity of $`𝐑P^2`$. Given a Harnack curve $`_i`$, denote by $`h_i`$ the affine corresponding curve. Denote by $`𝒞_i`$ the curve which deforms the union of $`L`$ with the curve of degree $`i1`$, $`im`$. Let the auxiliary curve of degree $`j`$ $`𝒞_j`$ be the union of $`j`$ lines of a pencil through a point. Noticing that the isotopy type of the affine curve $`h_i`$ ($`1im`$) depends exclusively on the position of $`𝒞_jL`$ ,$`1ji`$, we shall prove the Lemma 2.2.33 of Chapter 2.2.33. Let us detail our argumentation.
Let us prove that any two curves $`_{i+1}`$ and $`_{i+1}^{}`$ resulting from the classical perturbation of $`_iL`$ directed to a union of $`i+1`$ lines $`𝒞_{i+1}`$, resp $`𝒞_{i+1}^{}`$ are rigidly isotopic.
Denote by $`x_0`$, $`B_i`$, $`C_{i+1}`$,(resp, $`C_{i+1}^{}`$), $`B_{i+1}`$ ,(resp,$`B_{i+1}^{}`$) the polynomial of $`L`$, $`_i`$, and $`𝒞_{i+1}`$, (resp $`𝒞_{i+1}^{}`$), $`_{i+1}`$, (resp, $`_{i+1}^{}`$).
In the construction of $`_{i+1}`$ of type $`^0`$ from $`_iL`$, the curve $`𝒞_{i+1}`$ intersects $`L`$ in $`i+1`$ distinct points lying
-for even $`i+1`$ (i.e construction of $`_{2k}`$ from $`_{2k1}L`$) in the component $`S_i`$ of $`𝐑L`$ which is a boundary of the unique non-empty positive region $`\{x𝐑^2|b_i(x)>0\}`$ the non-empty region of $`𝐑^2`$ with boundary a part of the odd component of $`𝐑_i`$).
-for odd $`i+1`$, (i.e construction of $`_{2k+1}`$ from $`_{2k}L`$) in the component $`S_i`$ of $`𝐑L\backslash 𝐑_i`$ containing $`𝐑L𝐑_{i1}`$.
(It is also necessary that $`𝒞_{i+1}`$ does not intersect $`_iL`$ in its singular points).
i.) If $`𝒞_{i+1}L=𝒞_{i+1}^{}L`$, one can vary continuously the direction of lines of $`𝒞_{i+1}`$ and in this way define a one parameter family $`𝒞_{i+1,t;t[0,1]}`$, $`𝒞_{i+1,0}=𝒞_{i+1}`$, $`𝒞_{i+1,1}=𝒞_{i+1}^{}`$, and thus a rigid isotopy $`x_0.B_i+C_{i+1,t}`$ from $`_{i+1}`$ to $`_{i+1}^{}`$.
ii.) Otherwise $`𝒞_{i+1}L𝒞_{i+1}^{}L`$. In the construction of $`_{i+1}`$ of type $`^0`$ from classical deformation of $`_iL`$, any auxiliary curve which deforms the union $`_iL`$ intersects $`𝐑L`$ in $`i+1`$ points lying in a connected part of $`S_i`$, $`S_i𝐑L`$. Given $`𝒞_{i+1}`$, denote $`𝒞_{i+1}L`$ by $`I_{i+1}`$ and let $`I_{i+1}^{}`$ be a set of $`i+1`$ real points of $`S_i`$. Let us distinguish the cases $`i+1`$ even and odd.
-In case $`i+1=2k`$ even, (i.e construction of $`_{2k}`$ from $`_{2k1}L`$) $`S_i=S_{2k1}`$ is connected. According to the previous study, one can assume that lines of $`𝒞_{2k}`$ and $`𝒞_{2k}^{}`$ have the same direction. It is not hard to see that there exists a rigid isotopy $`B_{2k,t}`$ with $`t[0,1]`$ of $`B_{2k,0}=x_0.B_{2k1}+C_{2k}`$, $`B_{2k,1}=x_0.B_{2k1}+C_{2k}^{}`$ such that $`𝒞_{2k}^{}(x_0,x_1,x_2)L=I_{2k}^{}`$.
-The case $`i+1=2k+1`$ odd (i.e construction of $`_{2k+1}`$ from $`_{2k}L`$) differs from the case $`i`$ odd in the sense that the set $`S_i=S_{2k}`$ has two connected parts. These two connected parts may be defined using a projection to $`L`$ in a direction perpendicular to $`L`$ as follows. The non-empty oval of $`𝐑_{2k}`$ intersects the line $`L`$ in $`2k`$ points. They belong to a segment of $`L`$ with extremities $`2`$ points of $`𝐑_{2k}L`$. Using a pencil of lines in a direction (for example in a direction perpendicular) to $`L`$, project the non-empty oval of $`𝐑_{2k}`$ to a segment $`S`$ of $`L`$. The set $`S_{2k}`$ consists of the two segments $`S_{2k}^1=[a_1,b_1[`$ and $`S_{2k}^2=[a_2,b_2[`$ with extremities an extremity of $`S`$ and a point of $`𝐑_{2k}L`$ (this last one is the open extremity of the segment). Changing the direction of the projection, one defines two other segments. (The idea in our proof is to find a way to join this two segments. It may be done by changing the direction of the projection.)
Let $`_{2k+1}`$, such that $`_{2k+1}S_{2k}^2=\mathrm{}`$ and $`_{2k+1}^{}S_{2k}^2=\mathrm{\varnothing ̸}`$. To prove that the curves $`_{2k+1}`$ and $`_{2k+1}^{}`$ are rigidly isotopic, it is sufficient to prove that there exists a continuous path $`_{2k+1,t}`$ $`t[0,1]`$, $`_{2k+1,0}=_{2k+1}`$, $`_{2k+1,1}`$ with the following properties:
-for any $`t[0,1]`$ $`_{2k+1,t}`$ is a Harnack curve of degree $`2k+1`$
-as $`t`$ varies from $`0`$ to $`1`$ a line of $`_{2k+1}`$ (i.e of $`𝒞_{2k}`$) intersecting $`S_{2k+1}^1`$ is deformed to a line of ($`_{2k+1,1}^{}`$) intersecting $`S_{2k}^2`$.
(Recall that if $`𝒞_{2k+1}L=𝒞_{2k+1}^{}L`$ there exists a rigid of isotopy $`x_0.B_{2k}+C_{2k+1,t}`$, $`t[0,1]`$, from $`_{2k+1}`$ to $`_{2k+1}^{}`$. Hence, we shall consider $`𝒞_{2k+1}`$ and $`𝒞_{2k+1}^{}`$ up to such rigid isotopy.) Let $`\gamma `$ be the non-empty arc (of the non-empty oval) of $`_{2k}`$ which intersects $`L`$ in two points $`a_1S_{2k}^1=[a_1,b_1[`$ $`a_2S_{2k}^2=[a_2,b_2[`$. Let $`p\gamma `$, consider a path $`\gamma _t`$, $`\gamma _0=p`$ which moves $`p`$ along $`\gamma `$.
Let $`p_1S_{2k}^1`$, and $`L_1=(pp_1)`$ be a line through $`p`$. Consider the pencil of lines through the point $`p`$ which intersect $`\gamma `$. Move continuously the point $`p`$ along the path $`\gamma _t\gamma `$ and in this way lines $`(p_1\gamma _t)`$ of the pencil of lines through $`p_1`$. The projection of $`p`$ (in a direction perpendicular to $`L`$) to $`L`$ belongs to $`S_{2k}^1`$. It is not hard to see that one can choose the path $`\gamma _t`$ such that, as $`p`$ moves along $`\gamma _t`$, the projection $`\gamma _t`$ to $`L`$ moves from $`S_{2k}^1`$ to $`p_2S_{2k}^2`$.
In this way, considering together pencil of lines through any point $`p_t:=\gamma _t`$ and pencils through the point $`p_1S_{2k}^1`$ and pencils through the point $`p_2S_{2k}^2`$ it follows that there exists a rigid isotopy of $`_{2k+1}`$ from $`_{2k+1}^{}`$ which maps a line of $`𝒞_{2k+1}`$ which intersects $`S_{2k}^1`$ to a line of a curve $`𝒞_{2k+1}^{}`$ which intersects $`S_{2k}^2`$.
It concludes the proof of the Lemma 2.2.33 of Chapter 2.2.33. Q.E.D
It follows from the Lemma 2.2.33 of Chapter 2.2.33 that when one considers Harnack’s construction of curves $`_m`$ up to rigid isotopy, we may restrict our study to classical deformation of $`_iL`$, $`1i(m1)`$, directed to a union of $`i+1`$ lines of a pencil through a point $`p`$ chosen outside $`L`$ and $`_i`$.
Denote $`𝒞_{i+1}`$ the union of lines which deform the union of $`L`$ with the curve of degree $`i`$, $`im`$. The isotopy type of curves of degree $`m`$ obtained from the Harnack’s recursive method depends on the relative position of intersection points of $`𝒞_{i+1}`$ with the line $`L`$.
Consider the Harnack’s construction of curves $`_i`$, $`1im`$. Denote by $`B_i`$ the polynomial of $`_i`$. Consider $`L`$ as the line at infinity, and denote by $`b_i`$ the affine polynomial associated to $`B_i`$. Curves $`_i`$ constructed from the Harnack’s method are of type $`^0`$ relatively to $`L`$.
Let $`_5`$ be the Harnack curve of degree $`5`$. There exists a unique positive region $`\{x𝐑^2|b_5(x)>0\}`$ with a segment $`S_5`$ of the line $`𝐑L`$ on its boundary which contains an oval of $`_5`$. To get $`_6`$ from $`_5L`$, it is necessary that the intersection of $`𝒞_6`$ with $`L`$ consists of $`6`$ points lying in $`S_5`$. The non-empty oval of $`_6`$ intersects $`L`$ in $`6`$ points. Denote $`S_6`$ the connected part of the line $`𝐑L`$ which contains the intersection points of $`_6`$ with $`L`$. To get $`_7`$ from $`_6L`$, it is necessary that the intersection of $`𝒞_7`$ with $`L`$ consists of $`7`$ points lying in $`L\backslash S_6`$ and that it intersects $`_6`$ in its non-empty real component. Iterating this process, (set $`2k+1`$ (resp, $`2k`$) instead of $`5`$ (resp, $`6`$)) we define a sequence $`S_i`$ of connected components of $`𝐑L`$ such that $`S_i𝒞_{i+1}L`$.
In other words, we define relative position of lines $`𝒞_{i+1}`$ as follows. We take $`𝒞_m`$ to be a union of $`m`$ lines which intersect $`L`$ in $`m`$ distinct points lying, for even $`m`$ in the component of $`𝐑L`$ which is a boundary of the unique non-empty positive region $`\{x𝐑^2|b_{m1}(x)>0\}`$, for odd $`m`$ in the component of $`𝐑L\backslash 𝐑_{m1}`$ containing $`𝐑L𝐑_{m2}`$. (It is also necessary to choose $`𝒞_m`$ such that it does not intersect $`_{m1}L`$ in its singular points).
Denote by $`B_i^1`$ (resp, $`B_i^2`$) the polynomial of $`_i^1`$ (resp, $`_i^2`$) and by $`C_{i+1}^1`$ (resp, $`C_{i+1}^2`$) the union of $`i+1`$ lines which deform $`_i^1L_1`$ (resp, $`_i^2L_1`$) to $`_{i+1}^1`$ (resp, $`_{i+1}^2`$).
To prove Theorem 2.2.32 of Chapter 2.2.32, is sufficient to construct for $`i,1i(m1)`$ a continuous one parameter family with parameter $`t[0,1]`$ of curves $`𝒞_{i+1}^{1+t}`$ such that as $`t`$ varies from $`0`$ to $`1`$ the relative position of lines $`𝒞_{i+1}^{1+t}`$ remains the same.
The rigid isotopy $`h_t`$ is continuous. Therefore, the function which associates to any line $`h_t(𝐑L_1)`$ its normal vector $`\stackrel{}{n}_t`$ and its tangent vector $`\stackrel{}{v}_t`$ is also continuous. It follows from the continuity of $`h_t`$, that one can choose lines $`𝒞_{i+1}^{1+t}`$ such that the relative position of intersections $`𝒞_{i+1}^{1+t}L^{1+t}`$, and also the direction of lines of $`𝒞_{i+1}^{1+t}`$ relatively to $`(\stackrel{}{v}_t,\stackrel{}{n}_t)`$ remains the same as $`t`$ varies from $`0`$ to $`1`$. In such a way, for any $`i(m2)`$ we construct by induction curves $`_{i+2}^t`$, $`t[0,1]`$, of degree $`i+2`$ rigidly isotopic. Hence, for any $`i(m2)`$, curves $`_{i+2}^1`$ and $`_{i+2}^2`$ are rigidly isotopic. In particular, curves $`_m^1`$ and $`_m^2`$ are rigidly isotopic. Q.E.D
#### 2.3. Harnack Curves from a Complex viewpoint
In this section, we shall first construct particular deformations of Harnack polynomials. Then, we shall deduce from the properties of these deformations a characterization of the complex set point of Harnack curves in $`𝐂P^2`$.
We shall divide this section into two subsections. Denote $`_m`$ the Harnack curve of degree $`m`$ defined up to isotopy of real points set. In the first section, we define (in proposition 2.3.6 of Chapter 2.3.6) deformation of any Harnack curve $`_m`$ to a singular curve of which singular points are critical points index $`1`$ of the Harnack curve $`_m`$. At first, we characterize in proposition 2.3.1 of Chapter 2.3.1 such deformation for Harnack curves obtained via the pathchworking method and then generalize it in proposition 2.3.6 of Chapter 2.3.6 to any Harnack curve. Considering patchworking method seems at first glance of less interest at this time of our proof of the Rokhlin’s Conjecture. Nonetheless, it will be useful in the second part (Perestroika Theory on Harnack curves). In the next section, we deduce from the proposition 2.3.6 of Chapter 2.3.6 a description of $`(𝐂P^2,𝐂_m)`$ up to conj-equivariant isotopy and give the main result (Theorem 2.3.9 of Chapter 2.3.9) of this chapter.
##### Deformation of Harnack Curves
Given a non-singular curve $`𝒫_m`$ in $`𝐂P^2`$, call simple deformation $`𝒫_{m;t}`$ of the curve $`𝒫_m`$ a path
$`[0,1]𝐑𝒞_m`$
$`t𝒫_{m;t}`$
with the following properties:
1. For any $`0t<1`$, $`𝒫_{m;t}`$ is a smooth curve. If $`𝒫_{m;1}`$ is singular, then it is irreducible and any of its singular points is a real crossing which is critical point with positive critical value of the affine Harnack polynomial. Denote $`S`$ the set of singular points of$`𝒫_{m;1}`$.
2. For $`ϵ>0`$, let $`𝒟_ϵ=_{aS}D_(a,ϵ)`$ be the union of disc $`D(a,ϵ)`$ (in the Fubini-Study metric) of center $`a`$ and radius $`ϵ`$ taken over crossings $`a`$ of $`𝒫_{m;1}`$.
There exists $`ϵ_0>0`$, such that: for any curve $`𝒫_{m,t}`$, $`t]0,1[`$, $`𝐂P_{m,t}`$ lies in an $`ϵ_0`$-tubular $`N`$ neighborhood of $`𝐂P_{m,1}\backslash \{aS\}`$.
$`𝐂P_{m,t}`$ can be extented to the image of a smooth section of the tubular fibration $`N𝐂P_{m,1}\backslash \{aS\}`$.
Under two distinct simple deformations $`𝒫_{m;t}`$ and $`\stackrel{~}{𝒫_{m;t}}`$ of $`𝒫_m`$, the number of real singular points of $`𝒫_{m;1}`$ and $`\stackrel{~}{𝒫_{m;1}}`$ may be distinct.
Among simple deformations we shall distinguish deformations which lead to a curve $`𝒫_{m;1}`$ with the maximal number of real singular points a curve $`𝒫_{m;1}`$ may have. Denote by $`\alpha `$ this number. We shall call maximal simple deformation $`𝒫_{m;t}`$, $`t[0,1]`$, of the curve $`𝒫_m`$ a simple deformation $`𝒫_{m;t}`$ such that the curve $`𝒫_{m;1}`$ has $`\alpha `$ real singular points.
In proposition 2.3.6 of Chapter 2.3.6, we define some deformation of $`𝐑_m`$,
$`[0,1]𝐑𝒞_m`$
$`t𝐑_{m;t}`$
to the real points set $`𝐑𝒜`$ of a singular curve $`𝒜`$ of degree $`m`$ of which singular points are critical points of index $`1`$ of the Harnack curve $`_m`$. At first, we characterize in proposition 2.3.1 of Chapter 2.3.1 such deformation for Harnack curves obtained via the pathchworking method and then generalize it in proposition 2.3.6 of Chapter 2.3.6 to any Harnack curve. Recall that given a Harnack polynomial $`B_{2k}`$ of degree $`2k`$ type $`^0`$ and associated affine polynomial $`b_{2k}`$, $`S_{2k}^{}`$ denotes the subset of critical points $`(x_0,y_0)`$ of positive critical value $`c_0`$ with the property that as $`c`$ increases from $`c_0ϵ`$ to $`c_0+ϵ`$ the number of real connected components of $`M_c=\{((x,y)𝐑^2|b_{2k}>c\}`$ bounding the line at infinity increases.
###### Maximal Deformation of $`T`$-Harnack curves
Let us in Proposition 2.3.1 of Chapter 2.3.1 characterize maximal simple deformation of $`T`$-Harnack curves. It is obious that the Proposition 2.3.1 of Chapter 2.3.1 may be generalized to any Harnack curve. The general statement 2.3.6 of Chapter 2.3.6 may be proven independently of Proposition 2.3.6 of Chapter 2.3.6. For sake of clarity, we present here deformation of $`T`$-Harnack curves.
###### Proposition 2.3.1.
1. Along any simple deformation of a $`T`$-Harnack curve of odd degree $`2k+1`$ curves are smooth.
2. Given $`_{2k}`$ a T-Harnack curve of even degree $`2k`$, any maximal simple deformation of $`_{2k}`$
$`[0,1]𝐂𝒞_m`$
$`t_{m;t}`$
is such that $`_{2k;1}`$ has $`2k3`$ crossings.
Let $`S`$ be the set of singular points of $`_{2k;1}`$. Then, there exists a Harnack polynomial $`B_{2k}`$ of degree $`2k`$ and type $`^0`$ of which the set of critical points contains the set $`S`$.
These points are critical points with positive critical value. Besides, $`k1`$ of these points are points of $`S_{2k}^{}`$.
As explained in section 0.1 of Chapter 0.1, one can give via the Patchworking method an inductive construction of Harnack curves called $`T`$-inductive contruction.
Recall that in the $`T`$-inductive construction of Harnack polynomials $`\stackrel{}{a}_m,t_m`$ denotes the pair constituted by a vector $`\stackrel{}{a}_m`$ and a real $`t_m`$ such that for any $`t]0,t_m[`$, $`\stackrel{~}{X}_{m;\stackrel{}{a}_m,t}=_{(i,j)verticesofT_m}ϵ_{i,j}a_{i,j}x_1^ix_2^jx_0^{mij}t^{\nu (i,j)}`$ is a Harnack polynomial of degree $`m`$.
Given $`\stackrel{}{a}_m,t_m`$, $`\stackrel{}{a}_{m+1}=(\stackrel{}{a}_m,\stackrel{}{c}_{m+1})`$.
###### Remark 2.3.2.
1. There exists $`\stackrel{}{c}_{2k+1}`$ such that $`t_{2k+1}=t_{2k}`$.
2. There exists $`\stackrel{}{c}_{2k}`$ and $`\tau <t_{2k1}`$ such that for $`t=\tau `$ the curve given by $`\stackrel{~}{X}_{2k;\stackrel{}{a}_{2k},\tau }`$ has $`2k3`$ singular points which are collinear crossings.
proof: Consider a $`T`$-Harnack curve of degree $`m`$.
For $`m=1`$ and $`m=2`$, the proposition 2.3.1 of Chapter 2.3.1 is trivially verified. For $`m>2`$, it is based on the previous proposition 2.1.4 of Chapter 2.1.4. We shall use the terminology introduced in the chapter 1 and denote $`\stackrel{~}{X}_{m;t}`$, $`t]0,t_m[`$ a T-Harnack polynomial of degree $`m`$. It is obvious that any $`T`$-Harnack polynomial is of type $`^0`$. In what follows, we shall deal with one-parameter polynomials $`\stackrel{~}{X}_{m;t}`$, $`t[t_m,t_{m1}[`$.
When no confusion is possible, we shall denote $`B_m`$ a regular Harnack polynomial of degree $`m`$ of type $`^0`$ and $`S_m^{}`$, respectively $`S_m^+`$, the set of its critical points of negative, respectively positive, critical value.
Let us recall a result of the patchworking theory.
###### Lemma 2.3.3.
(, p.190) Let $`a`$ be a polynomial such that $`a=0`$ admits an $`ϵ`$-tubular neighborhood. If a set $`U𝐂P^2`$ is compact and contains no singular points of $`a=0`$, then for any $`ϵ>0`$ and any polyhedron $`\mathrm{\Delta }\mathrm{\Delta }(a)`$ there exists $`\alpha >0`$ such that for any polynomial $`b`$ with $`\mathrm{\Delta }(b)\mathrm{\Delta }`$, $`ba<\alpha `$ and $`b^{\mathrm{\Delta }(a)}=a`$ the truncation $`b^{\mathrm{\Delta }(a)}`$ is $`ϵ`$-sufficient in $`U`$.
We shall consider polynomials $`\stackrel{~}{X}_{m+1;\tau }`$ with the following properties. For any elementary triangle $`\mathrm{\Delta }`$ of the triangulation of $`T_{m+1}`$:
1. the truncation $`\stackrel{~}{X}_{m+1;\tau }^\mathrm{\Delta }`$ is completely non-degenerate.
2. $`\stackrel{~}{X}_{m+1;\tau }^\mathrm{\Delta }`$ is $`ϵ`$-sufficient for $`\stackrel{~}{X}_{m+1;\tau }`$ in $`\rho ^{m+1}(𝐑_+\mathrm{\Delta }^0\times U_{𝐂^2})`$
We shall say that such polynomial $`\stackrel{~}{X}_{m+1;\tau }`$ satisfies the ”good truncation properties”. It is easy to see that any polynomial $`\stackrel{~}{X}_{m+1;\tau }`$ with good truncation properties is such that for any elementary triangle $`\mathrm{\Delta }`$, the truncation $`𝐑\stackrel{~}{X}_{m+1;\tau }^\mathrm{\Delta }`$ is isotopic to $`𝐑_{m+1}`$ in the open $`\rho ^m(𝐑_+\mathrm{\Delta }^0\times U_{𝐂^2})(𝐂^{})^2`$. (Obviously, any Harnack polynomial has good truncation properties.) Let us give the main ideas which motivate our study and explain the method of our proof. Motivation
The $`T`$-inductive construction of Harnack polynomials can be considered as a slightly modified version of Harnack’s initial one.
Fix $`t_0]0,t_m[`$, and denote by $`B_m:=\stackrel{~}{X}_{m;t_0}=_{(i,j)verticesofT_m}ϵ_{i,j}a_{i,j}x_1^ix_2^jx_0^{mij}t_0^{\nu (i,j)}`$ the corresponding T-Harnack polynomial.
The T-Harnack polynomial $`\stackrel{~}{X}_{m+1}=_{(i,j)verticesofT_{m+1}}ϵ_{i,j}a_{i,j}x_1^ix_2^jx_0^{m+1ij}t_0^{\nu (i,j)}`$ of degree $`m+1`$ may be deduced by the formula: $`\stackrel{~}{X}_{m+1;t}=x_0.B_m+t_0.C_{m+1}`$ where $`x_0`$ is a line, the curve given by $`C_{m+1}`$ is the union of $`m+1`$ parallels lines which intersect $`x_0=0`$ and do not pass through the singular points of the curve $`𝒜_{m+1}`$ of degree $`m+1`$ with polynomial $`x_0.B_m`$. For $`t`$ sufficiently small, in particular for $`t_0<t_{m+1}`$, $`\stackrel{~}{X}_{m+1;t}=x_0.B_m+t.C_{m+1}`$ is a Harnack polynomial $`B_{m+1}`$. Therefore, outside $`𝐑C_{m+1}`$ the curve $`𝐑\stackrel{~}{X}_{m+1;t_0}`$ is a level curve of the function $`\frac{x_0.B_m}{C_{m+1}}`$. On $`𝐑𝒜_{m+1}\backslash 𝐑C_{m+1}`$, this level curve has critical points only at the singular points of $`𝐑𝒜_{m+1}`$. These points are non-degenerate singular points. Hence, the behavior of $`𝐑\stackrel{~}{X}_{m+1;t_0}`$ outside $`𝐑C_{m+1}`$ is described by the implicit function theorem and Morse Lemma. In particular, we have the following description:
1. Let $`\{a_1,\mathrm{},a_m\}`$ be the set of crossings of $`x_0.B_m=0`$. Denote $`D(a_i,ϵ)`$ a small neighborhood of radius $`ϵ`$ around $`a_i`$ in $`𝐑P^2`$.
Denote $`𝒟_ϵ=_{i=1}^mD(a_i,ϵ)`$ the neighborhood of the set of singular points of $`𝐑𝒜_{m+1}`$ in $`𝐑P^2`$. Let $`N`$ be a tubular neighborhood of $`𝐑𝒜_{m+1}\backslash 𝒟_ϵ`$ in $`𝐑P^2\backslash 𝒟_ϵ`$. Then, there exists $`ϵ_0`$ such that for any $`0<ϵϵ_0`$ and any $`D(a_i,ϵ)`$ of $`𝒟_ϵ`$ there exists a homeomorphism $`h:D(a_i,ϵ)D^1\times D^1`$ (where $`D^1`$ is a one-disc of $`𝐑`$) such that
$$h(𝐑_{m+1}D(a_i,ϵ))=\{(x,y)D^1\times D^1|x.y=\frac{1}{2}\}$$
2. Moreover, for any $`0<ϵϵ_0`$, $`𝐑_{m+1}\backslash 𝒟_ϵ`$ is a section of the tubular fibration
$`N𝐑𝒜_{m+1}\backslash 𝒟_ϵ`$.
Method
Assume that there exists a singular polynomial $`\stackrel{~}{X}_{m+1;\tau }`$ with good truncation properties and denote by $`S`$ the set of its singular points. On this assumption, the Harnack curve $`_{m+1}`$ is the image of a smooth section of the tubular fibration $`N\{(x_0:x_1:x_2)𝐂P^2|\stackrel{~}{X}_{m+1;\tau }=0\}\backslash S`$.
Moreover, outside $`x_0.B_m=0`$ and outside $`C_{m+1}=0`$, curves $`\stackrel{~}{X}_{m+1;\tau }=0`$ and $`\stackrel{~}{X}_{m+1;t}=0`$, $`0<t<t_{m+1}`$, are isotopic. Thus, $`_{m+1}`$ and $`\stackrel{~}{X}_{m+1;\tau }=0`$ may be not isotopic only in $`ϵ`$-neighborhood $`U(p)`$ of $`_{m+1}`$ defined from points of faces $`\mathrm{\Gamma }l_mT_{m+1}`$ (see definition 0.1.8 of Chapter 0.1.8).
Denote $`U(p_i)`$, $`1im`$, the $`ϵ`$-neighborhood defined from points of a face $`\mathrm{\Gamma }_i=\{x(mi),yi,x+y=m\}`$, $`\mathrm{\Gamma }_il_mT_{m+1}`$. The union $`=_{i=1}^mU(p_i)`$ is a neighborhood of the set of singular points of $`𝐂𝒜_{m+1}`$ in $`𝐂P^2`$. Let $`N`$ be the $`ϵ`$-tubular neighborhood of $`𝐂𝒜_{m+1}\backslash `$ in $`𝐂P^2\backslash `$. It follows immediately from the Lemma 2.3.3 of Chapter 2.3.3 that given $`U𝐂P^2`$ compact which contains no singular points of $`𝐂𝒜_{m+1}`$, any polynomial $`x_0.\stackrel{~}{X}_{m,t}`$ with $`t]0,t_m[`$, is $`ϵ`$-sufficient for $`\stackrel{~}{X}_{m+1,t}`$ in $`U`$. In other words, for $`t]0,t_{m+1}[`$, the intersection $`U𝐂_{m+1}`$ is contained in $`N`$ and can be extended to the image of a smooth section of a tubular fibration $`N𝐂𝒜_{m+1}\backslash `$.
According to the corollary 1.2.4 of Chapter 1.2.4 of chapter 1 and its proof, in the patchworking scheme crossings of $`_mL`$ are in bijective correspondence with faces $`\mathrm{\Gamma }_i=\{x(mi),yi,x+y=m\}`$ of the triangulation of $`T_{m+1}`$. Hence, we may assume that any crossing $`a_i_mL`$ belongs to the $`ϵ`$-neighborhood $`U(p_i)`$ of $`_{m+1}`$ defined from points $`\mathrm{\Gamma }_i^0`$, and consider the truncation of the affine restriction of a Harnack polynomial $`b_{m+1}`$ on the monomials $`x^{mi}y^{i1},x^{mi}y^i,x^{mi+1}y^{i1},x^{mi+1}y^i`$ which is $`ϵ`$-sufficient for $`b_{m+1}`$ in $`U(p_i)`$.
Using these truncations, and the fact that for any neighborhood $`B(p)`$ of a singular point $`p`$ of $`\stackrel{~}{X}_{m+1;\tau }=0`$, $`𝐂_{m+1}B(p)`$ is the image of a smooth section of the tubular fibration
$`NB(p)\{(x_0:x_1:x_2)𝐂P^2|\stackrel{~}{X}_{m+1;\tau }=0\}B(p)\backslash p`$ we shall characterize singular points of curves $`\stackrel{~}{X}_{m+1;\tau }=0`$. Then, according to the definition of the vector $`\stackrel{}{a}_{m+1}`$, we shall construct singular polynomials $`\stackrel{~}{X}_{m+1;\tau }`$ with singular points on $``$ and thus characterize maximal simple deformation of $`T`$-Harnack curves. The exposition of the main ideas of our proof is now finished and we shall now proceed to precise arguments.
In order to construct singular polynomials $`\stackrel{~}{X}_{m;t}`$, $`t[t_m,t_{m1}[`$ with good truncation properties, we shall distinguish the cases of even and odd $`m`$.
We shall work with notations and definitions related to the patchworking theory given in the chapter 1.
i) Consider the case $`m=2k`$.
i.1) Let us prove the following Lemma.
###### Lemma 2.3.4.
Any polynomial $`\stackrel{~}{X}_{2k;\stackrel{}{a}_{2k},t}`$ which satisfies good truncation properties has at most $`2k3`$ singular points. These points are crossings.
proof: The proof is based on the Patchworking theory. From a local study on the $`e`$-neighborhoods $`U(p_i)`$ of $`_{2k}`$ defined from points $`\mathrm{\Gamma }_i^0`$ of faces $`\mathrm{\Gamma }_i=\{x(mi),yi,x+y=m\}`$, we shall deduce where singular points of a polynomial $`\stackrel{~}{X}_{2k;t}`$ may appear. i.a) Let $`S`$ be a square defined by vertices
$`(c,d),(c+1,d),(c,d+1),(c+1,d+1)`$ with $`c,d`$ odd and
$`c+d=2k2`$. It is contained in the interior of the triangle $`T_{2k}`$. (Obviously, there are $`(k1)`$ such squares contained in the interior of $`T_{2k}`$.)
Consider polynomials $`\stackrel{~}{X}_{2k;t}`$, $`t]0,t_{2k1}[`$
$`\stackrel{~}{X}_{2k;t}=x_0.\stackrel{~}{X}_{2k1;\stackrel{}{a}_{2k1},t}+C_{2k;\stackrel{}{c}_{2k},t}`$ satisfying good truncation properties. (Given the vector $`\stackrel{}{a}_{2k1}`$, the vector $`\stackrel{}{c}_{2k}((𝐑^{})^+)^{2k+1}`$ is the vector which may be chosen.)
Denote $`x_{2k;t}`$ the affine restriction of a polynomial $`\stackrel{~}{X}_{2k;t}`$, $`t]0,t_{2k1}[`$.
Denote $`x_{2k;t}^S(x,y)`$ its truncation on the monomials
$`x^cy^d,x^cy^{d+1},x^{c+1}y^d,x^{c+1}y^{d+1}`$. Namely,
$$x_{2k;t}^S(x,y)=a_{c,d}t^{\nu (c,d)}x^cy^d+a_{c,d+1}t^{\nu (c,d+1)}x^cy^{d+1}+a_{c+1,d}t^{\nu (c+1,d)}x^{c+1}y^d$$
$$a_{c+1,d+1}t^{\nu (c+1,d+1)}x^{c+1}y^{d+1}$$
with $`a_{c,d}>0,a_{c+1,d}>0,a_{c,d+1}>0,a_{c+1,d+1}>0`$.
Let $`\mathrm{\Gamma }`$ be the face in the triangulation of $`T_{2k}`$ given by coordinates $`(c,d+1),(c+1,d)`$. According to the patchworking construction, it follows that for $`t]0,t_{2k}[`$ the truncation $`x_{2k;t}^S`$ is $`ϵ`$-sufficient for $`x_{2k;t}`$ in an $`ϵ`$-neighborhood $`U(p)`$ of $`x_{2k;t}^S=0`$ defined from points of $`\mathrm{\Gamma }^0`$.)
Fix $`t]0,t_{2k1}[`$, and let $`a_t`$ be the crossing of $`_{2k1}L`$ contained in $`U(p)`$.
According to the patchworking construction, there exists an homeomorphism $`\stackrel{~}{h}:𝐂_{2k}U(p)\{(x,y)(𝐂^{})^2|x_{2k;t}^S=0\}U(p)`$ such that $`\stackrel{~}{h}(a_t)=(x_1,y_1)_t`$ is a critical point of $`x_{2k;t}^S(x,y)`$. Obviously, (see corollary 1.2.4 of Chapter 1.2.4), $`(x_1,y_1)_t`$ is such that $`x_1.y_1<0`$.
Set $`x_{2k;t}^S(x,y)=l_t(x,y)+t.k_t(x,y)`$ with
$`l_t(x,y)=a_{c,d}t^{\nu (c,d)}x^cy^d+a_{c+1,d}t^{\nu (c,d+1)}x^{c+1}y^d`$
$`k_t(x,y)=a_{c,d+1}t^{\nu (c+1,d)1}x^cy^{d+1}+a_{c+1,d+1}t^{\nu (c+1,d+1)1}x^{c+1}y^{d+1}`$. Note that up to modify the coefficients $`a_{c,d},a_{c,d+1},a_{c+1,d},a_{c+1,d+1}`$ if necessary, the point $`\stackrel{~}{h}(a_t)=(x_1,y_1)_t`$ is a critical point of the function $`\frac{l_t(x,y)}{k_t(x,y)}`$ with positive critical value $`t<t_{2k1}`$.
On this assumption, it follows from the equalities
$`l_t(x,y)=l_t(x,y)`$, $`k_t(x,y)=k_t(x,y)`$,
$`\frac{l_t}{x}(x,y)=\frac{l_t}{x}(x,y)`$, $`\frac{l_t}{y}(x,y)=\frac{l_t}{y}(x,y)`$,
$`\frac{k_t}{x}(x,y)=\frac{k_t}{x}(x,y)`$, $`\frac{k_t}{y}(x,y)=\frac{k_t}{y}(x,y)`$
that $`(x_1,y_1)_t`$ is a critical point of the function $`\frac{l_t(x,y)}{k_t(x,y)}`$ with negative critical value $`t>t_{2k1}`$.
Therefore, for fixed $`t[t_{2k},t_{2k1}[`$, one can choose $`C_{2k;t}`$ and thus $`\stackrel{~}{X}_{2k;t}`$, in such a way that $`\stackrel{~}{h}^1(x_1,y_1)_t`$ with $`x_1.y_1<0`$ is a singular point of $`\stackrel{~}{X}_{2k;t}=0`$, and one-parameter curves $`\stackrel{~}{X}_{2k;t}`$ are Harnack curves for $`t>0`$ sufficiently small. Moreover, the singular situation, we shall denote by $`𝒮^{}`$, is as follows:
$`𝒮^{}`$: one outer oval of the part $`j_1(𝐑_{2k1}\backslash 𝒟_ϵ)`$ of $`𝐑_{2k}`$ touches the non-empty outer oval in the part $`𝐑_{2k}\backslash j_1(𝐑_{2k}\backslash 𝒟_ϵ)`$.
(Namely, one branch of a outer oval contained in the patchworking scheme in the subset $`\rho ^{2k}(D_{2k1,2k2}\times U_𝐑^2)𝐑P^2`$ touches a branch of the non-empty outer oval.) Thus, according to the Proposition 2.1.4 of Chapter 2.1.4 and Petrovskii’s theory (Petrovskii’s Lemmas (Lemma.2 and Lemma.3)), it follows that the point $`\stackrel{~}{h}^1(x_1,y_1)_t`$ is a critical point of a Harnack polynomial $`B_{2k}`$ of type $`^0`$ and belongs to $`S_{2k}^{}`$. i.b) Likewise, consider squares $`S`$ defined by vertices
$`(c+1,d),(c,d),(c,d+1),(c+1,d+1)`$ with $`c+d=2k2`$ and
$`c>0,d>0`$ even. (Obviously, there are $`(k2)`$ such squares contained into $`T_{2k}`$.) Denote $`x_{2k;t}^S(x,y)`$ the truncation of a polynomial $`x_{2k;t}`$ on the monomials $`x^cy^d,x^cy^{d+1},x^{c+1}y^d,x^{c+1}y^{d+1}`$. Let $`\mathrm{\Gamma }`$ be the face in the triangulation of $`T_{2k}`$ given by coordinates $`(c,d+1),(c+1,d)`$.
(From the patchworking construction, it follows that for $`t]0,t_{2k}[`$ the truncation $`x_{2k;t}^S`$ is $`ϵ`$-sufficient for $`x_{2k;t}`$ in an $`ϵ`$-neighborhood $`U(p)`$ of $`x_{2k;t}^S=0`$ defined from points of $`\mathrm{\Gamma }^0`$.)
Fix $`t]0,t_{2k1}[`$, and let $`a_t`$ be a crossing of $`_{2k1}L`$ contained in $`U(p)`$.
According to the patchworking construction, there exists an homeomorphism $`\stackrel{~}{h}:𝐂_{2k}U(p)\{(x,y)(𝐂^{})^2|x_{2k;t}^S=0\}U(p)`$ such that $`\stackrel{~}{h}(a_t)=(x_1,y_1)_t`$ is a critical point of $`x_{2k;t}^S(x,y)`$. Obviously, (see corollary 1.2.4 of Chapter 1.2.4), $`(x_1,y_1)_t`$ is such that $`x_1.y_1<0`$.
From an argumentation similar to the one above, keeping the same notations, it follows that, up to modify the coefficients $`a_{c,d},a_{c,d+1},a_{c+1,d},a_{c+1,d+1}`$ if necessary, one can get a critical point $`\stackrel{~}{h}(a_t)`$ $`(x_1,y_1)_t`$, $`x_1.y_1<0`$ of the function $`\frac{l_t(x,y)}{k_t(x,y)}`$ with positive critical value $`t<t_{2k1}`$. Therefore, for fixed $`t]0,t_{2k1}[`$, one can choose $`C_{2k;t}`$ and thus $`\stackrel{~}{X}_{2k;t}`$, in such a way that $`\stackrel{~}{h}^1(x_1,y_1)_t`$ with $`x_1.y_1<0`$ is a singular point of $`\stackrel{~}{X}_{2k;t}=0`$, and one-parameter curves $`\stackrel{~}{X}_{2k;t}`$ are Harnack curves for $`t`$ sufficiently small. Moreover, the singular situation, we shall denote by $`𝒮\mathrm{"}`$, is as follows:
$`𝒮\mathrm{"}`$: one branch of a inner oval of the part $`j_1(𝐑_{2k1}\backslash 𝒟_ϵ)`$ of $`𝐑_{2k}`$ not isotopic to a part of $`𝐑_{2k2}`$
(namely, in the patchworking scheme, one inner oval contained in the subset $`\rho ^{2k}(D_{2k1,2k2}\times U_𝐑^2)𝐑P^2`$) touches a branch of the non-empty outer oval. Thus, according to the Proposition 2.1.4 of Chapter 2.1.4 and the Petrovskii’s theory, it follows that the point $`\stackrel{~}{h}^1((x_1,y_1)_t)`$ is a critical point of positive critical value of a Harnack polynomial $`B_{2k}`$. (Obviously, $`\stackrel{~}{h}^1((x_1,y_1)_tS_{2k}^+\backslash S_{2k}^{}`$.) From an argumentation similar to the previous one, it is easy to deduce that there are no other points which may be singular points of a polynomial $`\stackrel{~}{X}_{2k;t}`$. Otherwise, it would contradict the strict positiveness of the coefficients of the vector $`\stackrel{}{a}_{2k}`$. (Indeed, consider the squares $`S`$ given by vertices $`(c+1,d),(c,d),(c,d+1),(c+1,d+1)`$ with $`c+d=2k2`$ and $`c=0`$ or $`d=0`$. Keeping the previous notations, from a crossing $`a_t`$ of $`_{2k1}L`$, one can get a critical point $`\stackrel{~}{h}(a_t)=(x_1,y_1)_t`$ of the function $`\frac{l_t(x,y)}{k_t(x,y)}`$ with negative critical value.)
Hence, we have obtained the Lemma 2.3.4 of Chapter 2.3.4.
i.2)
Let us now construct a polynomial $`\stackrel{~}{X}_{2k;\tau }`$, $`\tau [t_{2k},t_{2k1}[`$ with $`2k3`$ singular points: $`k1`$ points described locally by the singular situation $`𝒮^{}`$, and $`k2`$ points described locally by the singular situation $`𝒮\mathrm{"}`$.
Let us fix $`\tau <t_{2k1}`$. Crossings of the union $`_{2k1}L`$ belong to the line $`L`$.
Let $`p_i=\stackrel{~}{h}^1(x,y)_\tau `$ with $`(x,y)_\tau `$ defined up to homeomorphism from a crossing of the curve $`_{2k1}L`$ with polynomial $`x_0.\stackrel{~}{X}_{2k1;\stackrel{}{a}_{2k1},\tau }`$. We shall prove that given the $`2k3`$ points $`p_1`$,…,$`p_{2k3}`$ the polynomial $`\stackrel{~}{X}_{2k;\tau }`$ may be entirely defined.
According to the previous description of points $`p_1,\mathrm{},p_{2k3}`$, we shall assume that the $`2k3`$ singular points $`\stackrel{~}{X}_{2k,\tau }`$ belong to the line $`x_0^{}=0`$ with $`x_0^{}=x_0+\alpha .x_1+\beta .x_2`$. Then, consider the linear change of complex projective coordinates mapping $`(x_0:x_1:x_2)`$ to $`(x_0^{}:x_1:x_2)`$. Such transformation carries $`x_0.\stackrel{~}{X}_{2k1,\tau }(x_0,x_1,x_2)+C_{2k,\tau }`$ to $`\stackrel{~}{X}_{2k,\tau }^{}=x_0^{}.\stackrel{~}{X}_{2k1,\tau }^{}(x_0^{},x_1,x_2)+C_{2k,\tau }^{}`$. We shall prove that one can choose $`C_{2k}^{}(x_1,x_2)=0`$ a union of $`m`$ parallel lines in such a way that $`C_{2k}^{}(x_1,x_2)=0`$, $`\stackrel{~}{X}_{2k1}^{}(x_0^{},x_1,x_2)=0`$, and $`x_0^{}=0`$ have $`2k3`$ common points.
In such a way, using the linear change of projective coordinates mapping $`(x_0^{}:x_1:x_2)`$ to $`(x_0:x_1:x_2)`$, we shall get the polynomial $`\stackrel{~}{X}_{2k,\tau }(x_0,x_1,x_2)=x_0.\stackrel{~}{X}_{2k1,\tau }(x_0,x_1,x_2)+C_{2k}(x_1,x_2)`$ from the polynomial $`\stackrel{~}{X}_{2k,\tau }^{}(x_0,x_1,x_2)=x_0^{}.\stackrel{~}{X}_{2k1,\tau }^{}(x_0^{},x_1,x_2)+C_{2k}^{}(x_1,x_2)`$ and in that way the vector $`\stackrel{}{c}_{2k}`$. Let us detail this construction. Assume singular points $`p=(0:p_1:p_2)`$ of $`\stackrel{~}{X^{}}_{2k}=0`$ such that $`\frac{\stackrel{~}{X^{}}_{2k}}{x_0^{}}(p)=0`$ (i.e belong to the curve $`C_{2k}^{}=0`$.)
Assume $`(\frac{\stackrel{~}{X^{}}_{2k1}}{_{x_0^{}}}=0)(x_0^{}=0)(\stackrel{~}{X}_{2k1}^{}=0)=\mathrm{}`$,
Set
$$D_{2k}(x)=x_1.\frac{\stackrel{~}{X}_{2k}^{}}{x_1}+x_2.\frac{\stackrel{~}{X}_{2k}^{}}{x_2}$$
$$D_{2k1}(x)=x_1.\frac{^2\stackrel{~}{X}_{2k}^{}}{x_0^{}x_1}+x_2.\frac{^2\stackrel{~}{X}_{2k}^{}}{x_0^{}x_2}$$
where $`D_{2k1}(0,x_1,x_2)`$ is an homogeneous polynomial of degree $`(2k1)`$ in the variables $`x_1,x_2`$. Since $`(\frac{\stackrel{~}{X^{}}_{2k1}}{_{x_0^{}}}=0)(x_0^{}=0)(\stackrel{~}{X^{}}_{2k1}=0)=\mathrm{}`$, $`\frac{^2\stackrel{~}{X^{}}_{2k}}{^2x_0^{}}(p)0`$ for any singular point $`p`$ of $`\stackrel{~}{X^{}}_{2k}`$.
Thus, for $`2k>2`$, $`D_{2k}(x_0^{},x_1,x_2)`$ is of degree at least $`1`$ in $`x_0^{}`$ and any point $`p`$ belongs to $`\frac{D_{2k}}{_{x_0^{}}}=D_{2k1}=0`$.
According to the Euler formula, for any singular point $`p`$ of $`\stackrel{~}{X}_{2k}^{}`$ the following equalities are verified:
(2.9) $`D_{2k}(p)=p_1.{\displaystyle \frac{\stackrel{~}{X}_{2k}^{}}{x_1}}(p)+p_2.{\displaystyle \frac{\stackrel{~}{X}_{2k}^{}}{x_2}}(p)=0`$
(2.10) $`D_{2k1}(p)=p_1.{\displaystyle \frac{^2\stackrel{~}{X}_{2k}^{}}{x_0^{}x_1}}(p)+p_2.{\displaystyle \frac{^2\stackrel{~}{X}_{2k}^{}}{x_0^{}x_2}}(p)=0`$
(2.11) $`{\displaystyle \frac{\stackrel{~}{X}_{2k}^{}}{x_1}}(p)=0`$
$`or`$
$`{\displaystyle \frac{\stackrel{~}{X}_{2k}^{}}{x_2}}(p)=0`$
It is easy to deduce from the Newton’s binomial formula (applied to $`(x_0^{})^i=(x_0+\alpha .x_1+\beta .x_2)^i`$. that coefficients of the polynomial $`\stackrel{~}{X}_{2k}^{}=x_0^{}.\stackrel{~}{X}_{2k1}^{}+C_{2k}^{}`$ which depend on the vector $`\stackrel{}{c}_{2k}`$ (and do not depend on $`\stackrel{}{c}_{2k1}`$) are coefficients of monomials $`x_1^{2ki}x_2^i`$, (namely coefficients of $`C_{2k}^{}`$) and coefficients of monomials $`x_0^{}.x_1^{2k1i}x_2^i`$ and in that way coefficients of $`D_{2k1}(0,x_1,x_2)`$.
Therefore, any singular point of $`\stackrel{~}{X}_{2k}^{}=0`$ belongs to the intersection of $`x_0^{}=0`$ with $`C_{2k}^{}=0`$ and $`\stackrel{~}{X}_{2k1}^{}=0`$. Let us detail the construction of $`\stackrel{~}{X}_{2k}^{}`$. Since any singular point $`p`$ of $`\stackrel{~}{X}_{2k}^{}`$ belongs to $`O_1=\{(0:x_1:x_2)𝐂P^2|x_10\}`$ or $`O_2=\{(0:x_1:x_2)𝐂P^2|x_20\}`$, it is sufficient to consider singular points of $`\stackrel{~}{X}_{2k}^{}`$ in the affine charts associated to $`O_1`$ and $`O_2`$.
It follows from the form of $`\stackrel{~}{X}_{2k}^{}(x_0,x_1,x_2)`$ that points $`(0:0:1)`$ and $`(0:1:0)`$ are not singular points of the polynomial $`\stackrel{~}{X}_{2k}^{}(x_0,x_1,x_2)`$. Consequently, any singular point $`p`$ of $`\stackrel{~}{X}_{2k}^{}`$ is such that
$$\frac{\stackrel{~}{X}_{2k}^{}}{x_1}(p)=0$$
or
$$\frac{\stackrel{~}{X}_{2k}^{}}{x_2}(p)=0$$
Up to constant factor, the following equalities are verified:
(2.12)
$$\frac{\stackrel{~}{X}_{2k}^{}}{x_2}(0,x_1,x_2)=x_2^{2k1}c_1x_2^{2k2}x_1+\mathrm{}(1)^{2k1}c_{2k1}x_1^{2k1}$$
where $`c_i`$, $`i\{1,\mathrm{},2k1\}`$ is the elementary roots symmetric polynomial of degree $`i`$ of $`\frac{\stackrel{~}{X}_{2k}^{}}{x_2}(0,1,x_2)`$
(2.13)
$$\frac{\stackrel{~}{X}_{2k}^{}}{x_1}(0,x_1,x_2)=x_1^{2k1}s_1x_1^{2k2}x_2+\mathrm{}(1)^{2k1}s_{2k1}x_2^{2k1}$$
where $`s_i`$, $`i\{1,\mathrm{},2k1\}`$ is the elementary roots symmetric polynomial of degree $`i`$ of $`\frac{\stackrel{~}{X}_{2k}^{}}{x_1}(0,x_1,1)`$.
Set $`x_1=1`$ (respectively, $`x_2=1`$) in the equality (2.12) (respectively, (2.13)). Bringing together the resulting equalities, it follows that $`\stackrel{~}{X}_{2k}^{}(0,x_1,1)`$ has at most $`[\frac{2k1}{2}]`$ singular points (Indeed, any singular point $`(0,p_1,1)`$ is of multiplicity $`1`$ of $`D_{2k1}(0,x_1,1)`$ and of multiplicity at least $`2`$ of $`\stackrel{~}{X}_{2k,\tau }^{}(0,x_1,1)`$). Besides, $`\stackrel{~}{X}_{2k}^{}(0:1:x_2)`$ has at most $`[\frac{2k1}{2}]`$ singular points (since any singular point $`(0,1,p_2)`$ is of multiplicity $`1`$ of $`D_{2k1}(0,1,x_2)`$ and of multiplicity at least $`2`$ of $`\stackrel{~}{X}_{2k,\tau }^{}(0,1,x_2)`$)).
Since $`\stackrel{~}{X}_{2k}^{}=0`$ is singular, it follows immediately from the equalities (2.12), (2.13) and (2.10), (2.11) that at least one point $`p=(0:1:1)=(0:1:1)`$ is a singular point of $`\stackrel{~}{X}_{2k}^{}`$.
Therefore, the construction of the polynomial $`\stackrel{~}{X}_{2k;\tau }^{}(x_0,x_1,x_2)`$ with $`2k3`$ singular points $`p_1`$,…,$`p_{2k3}`$ on the intersection of $`x_0^{}=0`$ with $`C_{2k}^{}=0`$ and $`\stackrel{~}{X}_{2k1}^{}=0`$ follows immediately.
Moreover, according to the Lemma 2.3.3 of Chapter 2.3.3, one can perturb coefficients of one monomials $`x_1^{2k}`$ or $`x_2^{2k}`$ of the polynomial $`\stackrel{~}{X}_{2k}^{}`$ in such a way that the modified polynomial is a Harnack polynomial of degree $`2k`$ with critical points $`p_1`$,..,$`p_{2k3}`$. This concludes the part (2) of proposition 2.3.1 of Chapter 2.3.1.
ii) Consider the case $`m=2k+1`$. From an argumentation similar to the previous one, it is easy to deduce that any polynomial $`\stackrel{~}{X}_{2k+1;t}`$, $`t]0,t_{2k}[`$ which satisfies good truncation properties is a Harnack polynomial. Besides, it is also obvious that the existence of a singular polynomial $`\stackrel{~}{X}_{2k+1;t}(x_0,x_1,x_2).x_0+C_{2k+1;t}(x_1,x_2)`$ with good truncation properties is not compatible with the Harnack’s distribution of signs (i.e the strict positiveness of the coefficients of the vector $`\stackrel{}{a}_{2k+1}`$.)
Indeed, consider the squares $`S`$ given by vertices
$`(c+1,d),(c,d),(c,d+1),(c+1,d+1)`$ with $`c+d=2k1`$, $`c0`$ and $`d0`$. Keeping the notations introduced in the case $`m=2k`$, from a crossing $`a_t`$ of $`_{2k}L`$, one can get a critical point $`\stackrel{~}{h}(a_t)=(x_1,y_1)_t`$ of the function $`\frac{l_t(x,y)}{k_t(x,y)}`$ with negative critical value. Therefore, from an argument similar to the one given in case $`m=2k`$, it follows easily that any polynomial $`\stackrel{~}{X}_{2k+1;t}(x_0,x_1,x_2).x_0+C_{2k+1;t}(x_1,x_2)`$, $`t]0,t_{2k}[`$, with good truncation properties is smooth.
Q.E.D
###### Maximal Deformation of any Harnack curves
We shall now in Proposition 2.3.6 of Chapter 2.3.6 describe maximal simple deformation of any Harnack curve.
Let us start with Corollary 2.3.5 of Chapter 2.3.5 of Theorem 2.2.1 of Chapter 2.2.1.
###### Corollary 2.3.5.
Let $`_m`$ be a Harnack curve of degree $`m`$. Then, one can assume, without changing its topological properties, that any Harnack curve $`_m`$ is obtained via the $`T`$-inductive construction of Harnack curves.
proof: It is straightforward consequence of the rigid isotopy Theorem 2.2.1 of Chapter 2.2.1. Q.E.D.
The generalization of Proposition 2.3.1 of Chapter 2.3.1 to any Harnack curve is a straightforward consequence of the Corollary 2.3.5 of Chapter 2.3.5.
###### Proposition 2.3.6.
Denote $`_m`$ the Harnack curve of degree $`m`$ defined up to isotopy of real points set.
1. Along any simple deformation of a Harnack curve of odd degree $`2k+1`$ curves are smooth.
2. Given $`_{2k}`$ a Harnack curve of even degree $`2k`$, any maximal simple deformation of $`_{2k}`$
$`[0,1]𝐂𝒞_m`$
$`t_{m;t}`$
is such that $`_{2k;1}`$ has $`2k3`$ crossings.
Let $`S`$ be the set of singular points of the singular curve $`_{2k;1}`$. Then, there exists a Harnack polynomial $`B_{2k}`$ of degree $`2k`$ and type $`^0`$ of which the set of critical points contains the set $`S`$.
These points are critical points with positive critical value. Besides, $`k1`$ of these points are points of $`S_{2k}^{}`$.
proof: The generalization of Proposition 2.3.1 of Chapter 2.3.1 is a straightforward consequence of the Corollary 2.3.5 of Chapter 2.3.5.
We propose here a generalization of the argumentation given in the proof of proposition 2.3.6 of Chapter 2.3.6. Consider $`B_m(x_0,x_1,x_2)`$ a Harnack polynomial of degree $`m`$.
According to Proposition 2.2.31 of Chapter 2.2.31, it is possible to choose projective coordinates such that $`B_m(x_0,x_1,x_2)=x_0.B_{m1}(x_0,x_1,x_2)+C_m(x_1,x_2)`$ where $`B_{m1}(x_0,x_1,x_2)`$ is the Harnack polynomial of degree $`m1`$.
According to Morse Lemma , around any non-degenerate singular point $`s`$ of $`B_{m1}.x_0`$ one can choose a local coordinates system $`y_1,y_2`$ with $`z_1(s)=0,z_2(s)=0`$ and $`x_0.B_{m1}(x_0,x_1,x_2)=y_1.y_2`$.
Analogously, in a neighborhood of any critical point $`p`$ of $`B_m`$, one can choose a local coordinates system $`y_1,y_2`$ with $`y_1(p)=0,y_2(p)=0`$ and $`B_m(x_0,x_1,x_2)=B_m(p)+y_1.y_2`$. Therefore, given a singular crossing $`s`$ of $`x_0.B_{m1}`$, since $`s`$ is critical point of the curve $`B_m=0`$ considered as level curve of the function $`\frac{x_0.B_{m1}}{C_m}`$ the crossing $`s`$ is perturbed in such a way that around $`s`$, one can choose local coordinates in an open $`U(s)`$ with $`y_1(s)0,y_2(s)0`$ and $`B_m(x_0,x_1,x_2)=\frac{y_1}{y_2}+t`$.
Such description enlarges to a description of $`B_m(x_0,x_1,x_2)=0`$ inside $`UU(s)`$ $`U=\{z=<u,p>=(u_0.p_0:u_1.p_1:u_2.p_2)𝐂P^2|u=(u_0:u_1:u_2)U_𝐂^3,p=(p_0:p_1:p_2)U(s)\}`$ Local coordinates $`y_1,y_2`$ inside the open $`U(s)`$ extend to local coordinates inside $`U`$ as follows. Given $`z=<u,p>U`$ with $`u=(1:u_1:u_2)U_𝐂^3`$, $`p=(p_0:p_1:p_2)U(s)\}`$, we may set $`y_1(z)=u_1.y_1(p)`$, $`y_2(z)=u_2.y_2(p)`$.
In such a way, given a singular crossing $`s`$ of $`x_0.B_{m1}`$, any point $`pU`$ with local coordinates $`y_1(p),y_2(p)`$ such that $`|y_1(p)|=|y_1(s)|`$, $`|y_2(p)|=|y_2(s)|`$ and $`y_1(p).y_2(p)=y_1(s).y_2(s)`$, is a critical point of $`B_m(x_0,x_1,x_2)`$. Moreover, according to Rolle’s Theorem, around any singular point $`s`$ of $`x_0.B_{m1}`$, in a neighborhood of the line $`x_0=0`$, the sign of $`B_m(0,x_1/x_2,1)`$ alternates. (One can refer also to the marking method (see ) where a distribution of sign of $`C_m`$ and $`x_0.B_{m1}`$ is defined around any point of intersection of $`C_m`$ with $`x_0.B_{m1}`$.) Therefore, when consider the critical points of $`B_m`$ deduced locally from singular points of $`x_0.B_{m1}`$, the sign of $`B_m(0,x_1/x_2,1)`$ alternates as well around these points. By continuity of $`B_m`$, it follows the alternation of sign of $`B_m(x_0,x_1,x_2)`$ in an $`ϵ`$-tubular neighborhood of the line $`x_0=0`$ in $`𝐂P^2`$. It is obvious that such alternation of sign is equivalent to the modular property of the distribution of sign in the patchworking construction of Harnack’s curves. Hence, our proposition is a straightforward consequence of the version of proposition 2.3.6 of Chapter 2.3.6 for $`T`$-Harnack curves. For sake of clarity, we refer to proposition 2.3.1 of Chapter 2.3.1 of this chapter where a proof of the version of proposition 2.3.6 of Chapter 2.3.6 for $`T`$-Harnack curves is given.
Hence, it follows from proposition 2.3.1 of Chapter 2.3.1, that along any simple deformation of a Harnack curve of odd degree $`2k+1`$ curves are smooth. From an argumentation similar to the one given in the second part of proof of it follows that in case $`m=2k`$ one can modify coefficients of the polynomial in such a way that it has at most $`2k3`$ singular points. Q.E.D
##### Description of $`(𝐂P^2,𝐂_m)`$ up to conj-equivariant isotopy
Let us recall in introduction some properties of the real point set of Harnack curves which may be easily deduced from the Harnack curves construction and, according to proposition 2.3.6 of Chapter 2.3.6, may be generalized to any Harnack curve
Given $`𝒜_{m+1}=_mL`$. Let $`𝒟_ϵ=_{a_i_mL}D(a_i,ϵ)𝐑P^2`$ be a neighborhood of the set of singular points of $`𝒜_{m+1}`$ in $`𝐑P^2`$ (where $`D(a_i,ϵ)`$ denotes a 2-disc of radius $`ϵ`$ around $`a_i`$ in the Fubini-Study metric).
Let $`N`$ be a tubular neighborhood of $`𝐑𝒜_{m+1}\backslash 𝒟_ϵ`$ in $`𝐑P^2\backslash 𝒟_ϵ`$. Then, the Harnack curve $`_{m+1}`$ is a non-singular curve of degree $`m+1`$ such that: $`𝐑_{m+1}\backslash 𝒟_ϵ`$ is a section of the tubular fibration $`N𝐑𝒜_{m+1}\backslash 𝒟_ϵ`$.
Therefore, there exists $`ϵ>0`$ and $`\stackrel{~}{j}_t`$, with $`t[0,1]`$, an isotopy of $`𝐑P^2`$ which pushes $`𝐑_m\backslash 𝒟_{ϵ_0}`$ onto a subset of $`𝐑_{m+1}`$.
Besides, there is the following biunivoque correspondence between critical points of index $`0`$ and $`2`$ of Harnack polynomials $`B_m`$ and ovals of curves:
1. ovals of the part $`\stackrel{~}{j}_1(𝐑_m\backslash 𝒟_ϵ)`$ of $`𝐑_{m+1}`$ with ovals of $`_m`$ and thus with critical points of index $`0`$ and $`2`$ of $`B_{m+1}`$.
2. $`m`$ ovals of $`𝐑_{m+1}\backslash \stackrel{~}{j}_1(𝐑_m\backslash 𝒟_ϵ)`$ with critical points of index $`0`$ and $`2`$ of $`B_{m+1}`$.
Any oval of these $`m`$ ovals intersects a $`2`$-disc $`D(a_i)𝐑P^2`$ defined around a point $`a_i_mL`$.
We shall in Proposition 2.3.7 of Chapter 2.3.7 define the minimal (minimal in the sense of the inclusion of subsets of $`𝐂P^2`$), subset $`_ϵ`$ of $`𝐂P^2`$ which consists of union of $`4`$-ball (in the Fubini-Study metric) of $`𝐂P^2`$ of radius $`ϵ`$, $`𝒟_ϵ_ϵ`$, outside of which the isotopy $`\stackrel{~}{j}_t`$, $`t[0,1]`$, extends to a conj-equivariant isotopy $`𝐂P^2`$ which pushes $`𝐂_m\backslash `$ onto a subset of $`𝐂_{m+1}`$.
In such a way, we shall enlarge the above description of Harnack curves to the complex domain and then in theorem 2.3.9 of Chapter 2.3.9 describe the pair $`(𝐂_m,𝐂P^2)`$ up to conj-equivariant isotopy of $`𝐂P^2`$.
In way of preparation, let us make some remarks and recall generalities connected with Harnack curves and complex topological characteristics of real curves. A curve $`𝒜_m`$ of type $`I`$ is such its real set of points $`𝐑𝒜_m`$ divides its complex set of points into two connected pieces called halves $`𝐂𝒜_m^+`$, $`𝐂𝒜_m^{}`$. These two halves are exchanged by complex conjugation $`conj(𝐂𝒜_m^+)=𝐂𝒜_m^{}`$. The natural orientations of these two halves determine two opposite orientations on $`𝐑𝒜_m`$ (as their common boundary), called the complex orientations of $`𝒜_m`$.
It is easy to see that the deformation turning $`_mL`$, into $`_{m+1}`$ brings the complex orientations of $`_m`$ and $`L`$ to the orientations of the corresponding pieces of $`𝐑_{m+1}`$ induced by a single orientation of the whole $`𝐑_{m+1}`$.
In such a way, given an orientation of $`𝐑_m`$, there exists only one orientation of $`𝐑L`$ which induces the orientation of $`𝐑_{m+1}`$. Moreover, it is easy to deduce from the recursive construction of Harnack curves that orientations of $`𝐑L`$ alternate. Namely, given an orientation of $`𝐑_{m1}`$, $`m>1`$, if $`𝐂L^+`$ is the half of $`𝐂L`$ which induces the orientation on $`𝐑L`$ carried on $`𝐑_m`$, then $`𝐂L^{}=conj(𝐂L^+)`$ induces the orientation on $`𝐑L`$ carried on $`𝐑_{m+1}`$.
Let us recall the different ways to perturb one crossing $`p`$ of a singular curve $`𝒜`$. In a neighborhood invariant by complex conjugation of the crossing $`p`$ , the curve $`𝒜`$ can be considered up to conj-equivariant isotopy as the intersection of two lines in the point $`p`$. Thus, let $`L_1`$ and $`L_2`$ be lines embedded in $`𝐂P^2`$ which intersect each other in a single point $`p`$. Denote by $`C`$ the result of the perturbation of the union $`L_1L_2`$. From the complex viewpoint, there are essentially two ways to perturb the singular union $`L_1L_2`$. On the other hand, there are two ways to connect the halves of their complexifications. Indeed, the curve $`C`$ is a curve of type $`I`$ since it is a non-empty conic. The halves of $`𝐂L_i`$ $`i\{1,2\}`$ connected each other after perturbation correspond to the complex orientation of $`𝐑L_i`$ which agrees with some perturbation of $`𝐑C`$. One distinguishes two ways to connect the halves with each other. We shall define them as perturbation of type 1 and type 2 of the crossing (see figure $`1.1`$ and figure $`1.2`$). Let $`B`$ be a complex $`4`$-ball globally invariant by complex conjugation around the crossing. After a perturbation of type $`1`$, $`𝐑CB`$ does not divide $`𝐂CB`$ into two connected pieces. After a perturbation of type $`2`$, $`𝐑CB`$ divides $`𝐂CB`$ into two connected pieces. According to Theorem 2.2.30 of Chapter 2.2.30, one can assume that Harnack polynomials are deduced by induction on the degree as follows:
i) The Harnack polynomial $`B_1`$ is the polynomial of a real projective line.
Given $`B_{2k1}`$ a Harnack polynomial of degree $`2k13`$. Then, (up to linear change of projective coordinates and up to slightly modify its coefficients), $`B_{2k1}`$ is a Harnack polynomial of degree $`2k1`$, and type $`^0`$ deduced from classical small deformation $`B_{2k1}=x_0.B_{2k2}+ϵ_{2k}C_{2k1}`$ of $`x_0.B_{2k2}`$ where $`B_{2k2}`$ is Harnack polynomial of degree $`2k2`$ and type $`^0`$. We shall denote $`𝒜_{2k1}=_{2k2}L`$ the curve given by $`x_0.B_{2k2}`$.
For $`k1`$, we shall denote $`A_{2k1}`$ and call set of points perturbed in a maximal simple deformation of $`_{2k1}`$ the set of points $`A_{2k1}=\{a_1,\mathrm{},a_{2k}\}`$ where $`a_1,\mathrm{}a_{2k}`$ are the crossings of $`_{2k2}L`$.
ii) The Harnack polynomial $`B_2`$ is the polynomial of a curve of degree 2 of which real part consists of an oval. We shall denote $`A_2`$ the set $`A_2=\{a_1\}`$ where $`a_1`$ is the crossing of $`_1L`$.
Given $`B_{2k}`$ a Harnack polynomial of degree $`2k4`$. Then, (up to linear change of projective coordinates and up to slightly modify its coefficients), $`B_{2k}`$ is a Harnack polynomial of degree $`2k`$ and type $`^0`$ deduced from classical small deformation $`B_{2k}=x_0.B_{2k1}+C_{2k}`$ (with $`C_{2k}`$ arbitrarily small) of $`B_{2k1}.x_0`$, where $`B_{2k1}`$ is a Harnack polynomial of degree $`2k1`$ and type $`^0`$. We shall denote $`𝒜_{2k}=_{2k1}L`$ the curve given by $`x_0.B_{2k1}`$.
Let $`_{2k;t}`$, $`t[0,1]`$, be a maximal simple deformation of the Harnack curve $`_{2k}`$ given by the polynomial $`B_{2k}`$. For $`k2`$, we shall denote $`A_{2k}`$ and call set of points perturbed in a maximal simple deformation of $`_{2k}`$ the ordered set of points $`A_{2k}=\{a_1,\mathrm{},a_{2k1},\mathrm{},a_{4k2}\}`$ where $`a_1,\mathrm{}a_{2k1}`$ are the crossings of $`_{2k1}L`$ and $`a_{2k},\mathrm{},a_{4k2}`$ are the crossings of $`_{2k;1}`$.
One can easily notice that any point of $`A_m`$ is perturbed by a perturbation of type 1 to give $`_m`$. (i.e locally around any crossing of $`_mL`$, the deformation turning $`_mL`$ into $`_{m+1}`$ is a perturbation of type 1 of the crossing. In case of even $`m=2k`$, given a maximal simple deformation $`_{2k,t}`$, $`0t1`$, of $`_{2k}`$; locally around any crossing $`a`$ of $`_{2k;1}`$, the deformation turning $`_{2k;1}=0`$ into $`_{2k}=0`$ is a perturbation of type 1 of $`a`$.)
Let $`𝐂_m^+`$ (resp, $`𝐂L^+`$) be the half of $`𝐂_m`$ (resp, $`𝐂L`$) which induces orientation on $`𝐑_m`$ (resp,$`𝐑L`$). Denote $`𝐂𝒜_m^+`$ the union $`𝐂_m^+𝐂L^+`$.
In proposition 2.3.7 of Chapter 2.3.7, we construct a subset $`_ϵ`$ of $`𝐂P^2`$, (minimal in the sense of the inclusion of subsets of $`𝐂P^2`$), $`𝒟_ϵ_ϵ`$, with the property that outside $`_ϵ`$, the isotopy $`\stackrel{~}{j}_t`$ which pushes $`𝐑_m\backslash 𝒟_ϵ`$ onto a subset of $`𝐑_{m+1}`$ extends to a a conj-equivariant isotopy of $`𝐂P^2`$ which maps $`𝐂_{m+1}^+\backslash _ϵ`$ to $`𝐂_m^+𝐂L^+\backslash _ϵ`$. The construction of $`_ϵ`$ is based on the definition of a maximal simple deformation of $`_m`$ and the way points of the set $`A_m`$ are perturbed to give $`_m`$.
###### Proposition 2.3.7.
1. Let $`_ϵ=_{aA_{2k+1}}B_(a,ϵ)𝐂P^2`$, the union of conj-equivarant 4-ball $`B_(a,ϵ)`$ of radius $`ϵ`$ taken over crossings $`a`$ of $`𝒜_{2k+1}=_{2k}L`$, be a neighborhood of the set of singular points of $`𝒜_{2k+1}=_{2k}L`$ in $`𝐂P^2`$. Let $`N`$ be an oriented tubular neighborhood of $`𝐂𝒜_{2k+1}^+\backslash _ϵ`$ in $`𝐂P^2\backslash _ϵ`$.
Then, there exists $`ϵ_0`$ such that, $`𝐂_{2k+1}^+\backslash _{ϵ_0}`$ is the image of a non-zero section of the oriented tubular fibration $`N𝐂𝒜_{2k+1}^+\backslash _{ϵ_0}`$ and there exists $`j_t`$ with $`t[0,1]`$ a conj-equivariant isotopy of $`𝐂P^2`$ which maps $`𝐂_{2k+1}^+\backslash _{ϵ_0}`$ to $`𝐂_{2k}^+𝐂L^+\backslash _{ϵ_0}`$.
2. Let $`_ϵ=_{aA_{2k}}B(a,ϵ)𝐂P^2`$ the union of conj-equivarant 4-ball $`B(a,ϵ)`$ of radius $`ϵ`$ taken over points $`a`$ of $`A_{2k}`$,be a neighborhood in $`𝐂P^2`$ of the set of points perturbed in a maximal simple deformation of $`_{2k}`$. ($`_ϵ`$ contains besides neighborhoods of singular point of $`𝒜_{2k}=_{2k1}L`$, the union of neighborhood of $`2k3`$ crossings which appear in a maximal simple deformation of $`_{2k}`$.) Let $`N`$ be an oriented tubular neighborhood of $`𝐂𝒜_{2k}^+\backslash _ϵ`$ in $`𝐂P^2\backslash _ϵ`$.
Then, there exists $`ϵ_0`$ such that $`𝐂_{2k}^+\backslash _{ϵ_0}`$ is the image of a non-zero section of the oriented tubular fibration $`N𝐂𝒜_{2k}^+\backslash _{ϵ_0}`$ and there exists $`j_t`$, with $`t[0,1]`$, a conj-equivariant isotopy of $`𝐂P^2`$ which maps $`𝐂_{2k}^+\backslash _{ϵ_0}`$ to $`𝐂_{2k1}^+𝐂L^+\backslash _{ϵ_0}`$.
proof: Our proof is based on results of proposition 2.3.1 of Chapter 2.3.1 and the fact that orientation of $`𝐑L`$ alternates in the recursive construction of Harnack curves. Let us explain briefly the method of our proof. We shall consider the usual handlebody decomposition of $`𝐂P^2=B_0B_1B`$ where $`B_0`$,$`B_1`$,$`B`$ are respectively 0,2 and 4 handles. In such a way, the canonical $`𝐑P^2`$ can be seen as the union of a Möbius band $``$ and the disc $`D^2B`$ glued along their boundary. The Möbius band $``$ lies in $`B_0B_1S^1\times D^2`$. The complex conjugation switches $`B_0`$ and $`B_1`$ and $``$ belongs to the set of its fixed points.
Let $`𝒟_ϵ`$ be a neighborhood of the set of singular points of $`𝒜_{m+1}`$ in $`𝐂P^2`$ which is the union of $`2`$-disc of radius $`ϵ`$ around point of the set $`𝒜_m`$.
We shall assume that there exists an isotopy $`\stackrel{~}{j}_t`$ of $`𝐑P^2`$ which maps $`𝐑_{m+1}\backslash 𝒟`$ on the boundary of the Möbius band and study whether or not $`\stackrel{~}{j}_t`$ extends to $`j_t`$ a conj-equivariant isotopy of $`𝐂P^2`$ which maps $`𝐂_{m+1}^+\backslash 𝒟_ϵ`$ to $`𝐂_m^+𝐂L^+\backslash 𝒟_ϵ`$.
We shall notice that given a $`4`$-ball $`B(p)(B_0B_1)𝐂P^2`$ globally invariant by complex conjugation and such that $`B(p)^{}=B(p)B_1`$ and $`B(p)^+=B(p)B_0`$, if $`𝐂_{m+1}^+B(p)\mathrm{}`$ and $`𝐂_{m+1}^+B(p)𝐂_{m+1}^+B(p)^\pm `$ then there is no conj-equivariant isotopy $`j_t`$ of $`𝐂P^2`$ which maps $`𝐂_{m+1}^+\backslash 𝒟_ϵB(p)`$ to $`((𝐂_m^+𝐂L^+)\backslash 𝒟_ϵ)B(p)`$.
Thereby, since orientation of $`𝐑L`$ alternates in the recursive construction of Harnack curves an obstruction to construct a conj-equivariant isotopy in a $`4`$-ball $`B(p)`$ will be provided by a crossing $`aB(p)`$ of a curve $`_{m+1;1}`$ where $`_{m+1;t}`$ $`t[0,1]`$ is a simple deformation of the Harnack curve $`_{m+1}`$. We shall now proceed to precise arguments.
Assume $`m+1=2k4`$. First, we shall prove that the existence of a conj-equivariant isotopy $`j_t`$ which maps $`𝐂_{2k}^+\backslash 𝒟_ϵ`$ to $`𝐂_{2k1}^+𝐂L^+\backslash 𝒟_ϵ`$ is not compatible with elementary topological properties of $`𝐂_{2k}`$. Then, using proposition 2.3.1 of Chapter 2.3.1, we shall define a subset $`_ϵ𝐂P^2`$, $`_ϵ𝒟_ϵ`$, such that there exists a conj-equivariant isotopy $`j_t`$ which maps $`𝐂_{2k}^+\backslash _ϵ`$ to $`𝐂_{2k1}^+𝐂L^+\backslash _ϵ`$.
Since $`_{2k}`$ is of even degree, its real point set $`𝐑_{2k}`$ divides the projective plane $`𝐑P^2`$ in two components respectively orientable and non-orientable whose common boundary is $`𝐑_{2k}`$.
Without loss of generality, one can assume that $`D^2𝐑P^2`$ is a small real disc with boundary the line at infinity $`L`$ of $`𝐑P^2`$ with $`𝒟_ϵ𝐑P^2D^2`$. The tubular neighborhood of the line $`L`$ is the Möbius band $``$ of $`𝐑P^2`$. Consider $`\stackrel{~}{j}_t`$, with $`t[0,1]`$ an isotopy of $`𝐑P^2`$ which lets fix the interior of the disc $`D^2`$ of $`𝐑P^2`$ and pushes the rest of $`𝐑_{2k}\backslash 𝒟_ϵ`$ to the boundary of the Möbius band.
Assume that the isotopy $`\stackrel{~}{j}_t`$ extends to a conj-equivariant isotopy $`j_t`$ of $`𝐂P^2`$. Then, one can set $`j_1(𝐂_{2k}^+\backslash B)B_0`$. Given $`N`$ a tubular neighborhood of $`j_1(𝐂_{2k})`$ in $`𝐂P^2`$, $`j_1(𝐂_{2k})`$ is the image of a non-zero section of the oriented tubular fibration $`Nj_1(𝐂_{2k})`$.
Since $`j_t`$ is conj-equivariant, there exists $`N^+B_0`$ a tubular neighborhood of $`𝐂_{2k}^+\backslash B`$ in $`𝐂P^2\backslash 𝒟_ϵ`$ such that
$`N^+𝐂_{2k}^+\backslash B`$ is an oriented tubular fibration.
Hence, since $`_{2k}`$ is two-sided and has orientable real part, there exists $`N`$ an oriented tubular neighborhood of $`\stackrel{~}{j}_1(𝐑_{2k}\backslash 𝒟_ϵ)`$ in $`𝐑P^2\backslash 𝒟_ϵ`$ such that $`\stackrel{~}{j}_1(𝐑_{2k}\backslash 𝒟_ϵ`$ is the image of a non-zero smooth section of the tubular fibration $`N\stackrel{~}{j}_1(𝐑_{2k}\backslash 𝒟_ϵ)`$.
It leads to contradiction since the tubular neighborhood of $`j_1(𝐑_{2k}\backslash (D^2)^0)`$ in $`𝐑P^2\backslash (D^2)^0`$ is homeomorphic to the Möbius band and therefore has non-orientable normal fibration.
Therefore, on $`𝐑_{2k}\backslash 𝒟_ϵ`$, the isotopy $`\stackrel{~}{j}_t`$ does not extend to a conj-equivariant isotopy $`j_t`$ of $`𝐂P^2`$.
Nonetheless, using the proposition 2.3.6 of Chapter 2.3.6, we shall define a set $`_ϵ𝒟_ϵ`$ for which on $`𝐑_{2k}\backslash `$ the isotopy $`\stackrel{~}{j_t}`$ extends to a conj-equivariant isotopy $`j_t`$ which maps $`𝐂_{2k}^+\backslash _ϵ`$ to $`𝐂_{2k1}^+𝐂L^+\backslash _ϵ`$.
Given $`_{2k}`$ the Harnack curve of degree $`2k`$, consider $`_{2k;t}`$, $`t[0,1]`$, a simple deformation of the Harnack curve $`_{2k}`$. Let us prove the following Lemma.
###### Lemma 2.3.8.
There exists $`B(p)B_0B_1`$ such that
$`𝐂_{2k}^+B(p)𝐂_{2k}^+B(p)^\pm `$ if and only if there exists a simple deformation $`_{2k;t}`$ $`t[0,1]`$ of the Harnack curve $`_{2k}`$ such that a crossing of $`_{2k;1}`$ belongs to $`B(p)`$.
proof:
For $`ϵ>0`$ small, the sets of complex points of curves $`_{2k;t}`$ $`t[1ϵ,1]`$ belong to a tubular neighborhood of the set of complex points of the Harnack curves $`_{2k;1}`$.
Let $`B(p)𝐂P^2`$ be a $`4`$-ball centered in a singular point of $`_{2k;1}`$ such that $`𝐂_{2k;t,t[1ϵ,1]}B(p)\mathrm{}`$. Inside $`B(p)`$, consider the gradient trajectories of the deformation turning $`_{2k;1}`$ into $`_{2k;1ϵ}`$. Up to isotopy $`\stackrel{~}{j}_t`$ of $`𝐑P^2`$, one can assume $`B(p)B_0B_1`$ and set $`B(p)^{}=B(p)B_1`$ and $`B(p)^+=B(p)B_0`$.
Since orientation of the line $`L`$ alternates in the recursive construction of Harnack curves, the complex conjugation acts on gradient trajectories as symmetry of center $`p`$. Therefore, inside $`B(p)`$ we have $`𝐂_{2k}^+B(p)𝐂_{2k}^+B(p)^\pm `$. (i.e around $`B(p)`$, the deformation turning $`_{2k;1}`$ into $`_{2k;1ϵ}`$ is a deformation of type 1.)
Reciprocally, it is easy to see that if there exists $`B(p)B_0B_1`$ such that
$`𝐂_{2k}^+B(p)𝐂_{2k}^+B(p)^\pm `$ then there exists a simple deformation $`_{2k;t}`$ $`t[0,1]`$ such that a crossing of $`_{2k;1}`$ belongs to $`B(p)`$.
Q.E.D
The Lemma 2.3.8 of Chapter 2.3.8 implies that outside a neighborhood $`_ϵ`$ of the set of singular points of $`_{2k;1}`$, $`\stackrel{~}{j}_t`$ with $`t[0,1]`$ extends to a conj-equivariant isotopy of $`𝐂P^2`$ which maps $`𝐂_{2k}^+\backslash _ϵ`$ to $`𝐂_{2k1}^+𝐂L^+\backslash _ϵ`$.
Assume now $`m+1=2k+13`$, Without loss of generality one can assume that $`D^2`$ is a small real disc such that $`D^2_{2k+1}L=𝒟`$. Consider $`\stackrel{~}{j}_t`$, with $`t[0,1]`$, an isotopy of $`𝐑P^2`$ which lets fix the disc $`D^2`$ of $`𝐑P^2`$ and pushes $`𝐑_{2k+1}\backslash 𝒟_ϵ`$ to the boundary of the Möbius band. As any curve of odd degree, $`_{2k+1}`$ is one-sided. Therefore, we may not use the argument given in case of curve of even degree to refute the existence of a conj-equivariant isotopy $`j_t`$ of $`𝐂P^2`$ which maps $`𝐂_{2k+1}^+\backslash 𝒟_ϵ`$ to $`𝐂_{2k}^+𝐂L^+\backslash 𝒟_ϵ`$.
Besides, from an argument analogous to the one given in even degree, using the odd version of the previous Lemma 2.3.8 of Chapter 2.3.8, it follows that on $`𝐑_{2k}\backslash 𝒟_ϵ`$, $`\stackrel{~}{j}_t`$ extends to a conj-equivariant isotopy $`j_t`$ of $`𝐂P^2`$.
Otherwise, there would exist a simple deformation $`_{2k+1;t}`$ $`t[0,1]`$ of the Harnack curve $`_{2k+1}`$ of degree $`2k+1`$ which intersects the discriminant hypersurface $`𝐑𝒟_{2k+1}`$. According to proposition 2.3.6 of Chapter 2.3.6, it is impossible. Q.E.D
The following theorem is a straighforward consequence of the Proposition 2.3.7 of Chapter 2.3.7.
###### Theorem 2.3.9.
Let $`_m`$ be the Harnack curve of degree $`m`$. There exists a finite number $`I`$ ($`I=1+2\mathrm{}m+\mathrm{\Sigma }_{k=2}^{k=[m/2]}2k3`$) of disjoint $`4`$-balls $`B(a_i)`$ invariant by complex conjugation centered in points $`a_i`$ of $`𝐑P^2`$ such that up to conj-equivariant isotopy of $`𝐂P^2`$,
1. $`_m\backslash _{iI}B(a_i)=_{i=1}^mL_i\backslash _{i=1}^IB(a_i)`$ where $`L_1,\mathrm{},L_m`$ are $`m`$ distinct projective lines with
$$L_i\backslash _{i=1}^IB(a_i)L_j\backslash _{i=1}^IB(a_i)=\mathrm{}$$
for any $`ij`$, $`1i,jm`$.
2. situation inside any $`4`$-ball $`B(a_i)`$ is described by the perturbation of type $`1`$ of the crossing $`a_i`$.
proof:
For $`m=1`$, the theorem is trivially verified: $`_1`$ is a projective line. For $`m>1`$, it may be deduced by induction on $`m`$. The induction is based on the proposition 2.3.7 of Chapter 2.3.7 and the inductive construction of Harnack curves.
As related above, for each integer $`m`$ ($`m0`$), on can assume without loss of generality, that the curve $`_{m+1}`$ results from classical deformation of $`_mL`$.
According to the proposition 2.3.6 of Chapter 2.3.6,
1. let $`A_{2k}=\{a_1,\mathrm{},a_{2k1},\mathrm{},a_{4k2}\}`$ be the set of points perturbed in a maximal simple deformation of $`_{2k}`$
2. denote $`A_{2k+1}=\{a_1,\mathrm{},a_{2k}\}`$ the set of points perturbed in a maximal simple deformation of $`_{2k+1}`$
For any $`a_iA_{m+1}`$, we may always choose a small open conj-symmetric $`4`$-ball $`B(a_i)`$ of $`a_i`$, in such a way that $`B(a_i)\rho ^{m+1}(D(p_j)\times U_𝐂^2)`$ with $`D(p_j)`$ is a small open real disc around $`p_j\mathrm{\Gamma }_j`$.
1. Given $`_{2k1}`$ the Harnack curve of degree $`2k1`$, from the proposition 2.3.7 of Chapter 2.3.7, it results the following characterization of $`_{2k}`$:
There exists $`j_t`$, $`t[0,1]`$, a conj-equivariant isotopy of $`𝐂P^2`$, which maps $`_{2k}\backslash _{i=1}^{4k2}B(a_i)`$ onto
1. $`_{2k1}\backslash _{i=1}^{4k2}B(a_i)`$ (in the patchworking scheme, such part is contained in the restriction $`\rho ^{2k}(T_{2k1}\times U_𝐂^2)`$ of $`𝐂P^2𝐂T_{2k}`$.)
2. union $`L\backslash _{i=1}^{4k2}B(a_i)`$ (in the patchworking scheme, such part is contained in the restriction $`\rho ^{2k}(D_{2k,2k1}\times U_𝐂^2)`$ of $`𝐂P^2𝐂T_{2k}`$.)
2. Given $`_{2k}`$ the Harnack curve of degree $`2k`$. From the proposition 2.3.7 of Chapter 2.3.7, it results the following characterization of $`_{2k+1}`$:
There exists $`j_t`$, $`t[0,1]`$, a conj-equivariant isotopy of $`𝐂P^2`$, which maps $`_{2k+1}\backslash _{i=1}^{2k}B(a_i)`$ onto
1. $`_{2k}\backslash _{i=1}^{2k}B(a_i)`$ (in the patchworking scheme, such part is contained in the restriction $`\rho ^{2k+1}(T_{2k}\times U_𝐂^2)`$ of $`𝐂P^2𝐂T_{2k+1}`$.)
2. union $`L\backslash _{i=1}^{2k}B(a_i)`$ (in the patchworking scheme, such part is contained in the restriction $`\rho ^{2k+1}(D_{2k+1,2k}\times U_𝐂^2)`$ of $`𝐂P^2𝐂T_{2k+1}`$.)
Therefore, it follows easily by induction on $`m`$ that outside a finite number of $`4`$-balls $`B(a_i)`$ the curve $`_m`$ is, up to conj-equivariant isotopy, the union of $`m`$ (non-intersecting) projective lines minus their intersections with $`4`$-balls $`B(a_i)`$.
The $`4`$-balls $`B(a_i)`$ are centered in points $`a_i`$. Inside any $`4`$-ball $`B(a_i)`$, one can get the whole description of $`_{m+1}`$ from the real set of points $`𝐑_{m+1}B(a_i)`$ and its orientation (i.e the type of deformation).
Situations inside $`4`$-balls centered in crossings $`a_i`$ of $`_mL`$ are easily deduced from the construction of $`_{m+1}`$: the deformation turning $`(_{2k}L)B(a_i)`$ into $`_{2k+1}B(a_i)`$ is of type $`1`$.
In case $`m+1=2k`$, one has to consider $`4`$-balls centered in crossings $`B_{2k;1}`$. The situation inside $`4`$-balls centered in crossings $`B_{2k;1}`$ has been explicitly described inside $`D(a_i)`$, and respectively inside $`B(a_i)`$, in the proof of proposition 2.3.1 of Chapter 2.3.1, and respectively in the proof of proposition 2.3.7 of Chapter 2.3.7. The deformation turning $`_{2k;1}B(a_i)`$ into $`_{2k}B(a_i)`$ is of type $`1`$.
Q.E.D
### Chapter 3 Arnold surfaces of Harnack curves
In this chapter, we valid in Theorem 3.0.1 of Chapter 3.0.1 the Rokhlin Conjecture for the Harnack curves of even degree. In other words, we prove in Theorem 3.0.1 of Chapter 3.0.1 that Arnold surfaces of Harnack curves split into connected sum of $`𝐑P^2`$ and $`\overline{𝐑P^2}`$. This result is a straightforward consequence of Theorem 2.3.9 of Chapter 2.3.9.
###### Theorem 3.0.1.
Arnold surfaces of Harnack curves of even degree are standard. Moreover, let $`_{2k}`$ be the Harnack curve of degree $`2k`$ and $`𝔄_+`$, $`𝔄_{}`$ its Arnold surfaces then, up to isotopy of $`S^4`$: for $`k=1`$
$$𝔄_+D^2$$
$$𝔄_{}\overline{𝐑P^2}$$
for $`k2`$
$$𝔄_+p𝐑P^2\text{}q_+\overline{𝐑P^2}$$
$$𝔄_{}p𝐑P^2\text{}q_{}\overline{𝐑P^2}$$
where
$`p=q_+=\frac{(k1)(k2)}{2};q_{}=\frac{k(5k3)}{2}`$
proof:
The proof, is based on the theorem 2.3.9 of Chapter 2.3.9 stated previously, and on the next Livingston’s statement: Theorem : Let $`FD^4S^4`$ be a closed surface which lies in $`S^3=D^4`$, except several two-dimensional discs standardly embedded inside $`D^4`$. Then $`F`$ is a standard surface in $`S^4`$. Let us recall in details the handles decomposition of $`𝐂P^2`$ and
$`𝐂P^2/conjS^4`$. Choose $`b`$ a point in $`𝐑P^2`$ and $`B`$ a small neighborhood of $`b`$ in $`𝐂P^2`$. Setting $`b=(b_0:b_1:1)𝐑P^2`$, $`B`$ is defined, up to homeomorphism, in $`U_2=\{(x_0:x_1:x_2)𝐂P^2|x_20\}𝐂P^2`$ $`U_2𝐂^2`$ as a usual $`4`$-ball of $`𝐂^2`$: $`B\{p=(p_0:p_1:1)𝐂^2||(p_0b_0,p_1b_1)||ϵ\}`$.
The $`4`$-ball $`B`$ is globally invariant by complex conjugation. The real set of points of the ball $`B`$ is a two-dimensional disc $`D^2`$ which remains fix by complex conjugation. Then, consider the central projection
$$p:𝐂P^2\backslash B𝐂L$$
from $`b`$ to some set of complex points of a real projective line $`L`$, $`b𝐂L`$. Denote $`𝐂L_0`$ and $`𝐂L_1`$ the two connected parts of $`𝐂L`$ such that :
1. $`𝐂L=𝐂L_0𝐂L_1`$
2. $`𝐂L_0𝐂L_1=𝐑L`$
Since $`𝐂L_0D^2`$,$`𝐂L_1D^2`$, $`𝐑LS^1`$, the fibrations:
$$p^1(𝐂L_0)𝐂L_0D^2$$
$$p^1(𝐂L_1)𝐂L_1D^2$$
$$p^1(𝐑L)𝐑LS^1$$
are, up to homeomorphism, trivial fibrations.
Therefore, setting $`p^1(𝐂L_0)=B_0`$ and $`p^1(𝐂L_1)=B_1`$ it follows :
$$B_0D^2\times D^2,B_1D^2\times D^2,B_0B_1=p^1(𝐑L)S^1\times D^2$$
Thus, we shall take the handlebody decomposition of $`𝐂P^2`$ $`𝐂P^2=B_0B_1B`$ where $`B_0`$,$`B_1`$,$`B`$ are respectively 0,2 and 4 handles. Moreover, from the action of complex conjugation on
$`𝐂P^2=B_0B_1B`$ it is easy to deduce a decomposition of $`𝐂P^2/conjS^4`$. The canonical $`𝐑P^2`$ can be seen as the union of a Möbius band $``$ and the disc $`D^2`$ glued along their boundary. In such a way, $``$ lies in $`B_0B_1S^1\times D^2`$. The complex conjugation switches $`B_0`$ and $`B_1`$ and $``$ belongs to the set of its fixed points. The quotient $`(B_0B_1)/conj`$ is a $`4`$-ball as well which contains the quotient $`(B_0B_1)/conj`$.
Consider now $`_{2k}`$ the Harnack curve of degree $`2k`$ and $`𝔄_+`$ and $`𝔄_{}`$ its Arnold surfaces.
Acording to Theorem 2.3.9 of Chapter 2.3.9, there exists a conj-equivariant isotopy of $`𝐂P^2`$, $`h_t`$, $`t[0,1]`$, $`h(0)=_{2k}`$ and a finite number $`I`$ of disjoint $`4`$-balls $`B(a_i)`$ invariant by complex conjugation centered in points $`a_i`$ of $`𝐑P^2`$ such that:
1. $`h_1(_{2k}\backslash _{i=1}^IB(a_i))=_{i=1}^{2k}L_i\backslash _{i=1}^IB(a_i)`$ where $`L_1,\mathrm{},L_{2k}`$ are $`2k`$ distinct projective lines with
$$(L_i\backslash _{i=1}^IB(a_i))(L_j\backslash _{i=1}^IB(a_i))=\mathrm{}$$
for any $`i,j`$, $`1ij2k`$,
2. inside any $`4`$-ball $`B(a_i)`$, $`h_1(_mB(a_i))`$ is the perturbation of a crossing as described in theorem 2.3.9 of Chapter 2.3.9.
Choose the $`4`$-ball $`B`$ of the handlebody decomposition of $`𝐂P^2`$ centered in $`b`$ in such a way that:
1. $`b𝐑P_\pm ^2`$ (For example, choose $`b=(0,0,1)`$ inside $`𝐑P_{}^2`$ and
$`b=(ϵ,ϵ,1)`$, $`0<ϵ<1`$, inside $`𝐑P_+^2`$).
2. $`B`$ does not intersect with $`_{iI}B(a_i)/conj`$. and lines $`L_1,\mathrm{},L_{2k}`$ provided by the isotopy $`h_t`$.
Then, delete a small tubular neighborhood $`U`$ of $`(B_0B_1)/conj`$ from $`(B_0B_1)/conj`$ and denote by $`B^{}`$ the resulting 4-ball.
From the decomposition of $`𝐂P^2/conjS^4`$ and the conj-equivariant isotopy $`h_t`$, it follows an isotopy $`h_t/conj`$ of $`S^4`$ which pushes $`𝐑P_\pm ^2\backslash _{i=1}^IB(a_i)/conj`$ on the boundary of $`B^{}`$ and $`(_{2k}/conj\backslash _{i=1}^IB(a_i)/conj)`$ to $`2k`$ standards discs of $`B^{}`$. (such discs are defined by quotient of lines $`L_1/conj`$, …,$`L_{2k}/conj`$ minus their intersection with $`_{iI}B(a_i)/conj`$.)
Thus, since $`𝔄_\pm \backslash _{iI}B(a_i)/conjIntU=\mathrm{}`$, the problem of construction of an isotopy which pushes $`𝔄_\pm `$ to the boundary of a $`4`$-disc and to several $`2`$-discs standardly embedded is reduced to the problem of construction of such an isotopy inside $`_{iI}B(a_i)/conj`$.
Moreover, since the $`4`$-balls $`B(a_i)`$ are disjoints, it is sufficient to study local questions inside $`2`$-discs $`B(a_i)/conj`$.
Inside any $`4`$-ball $`B(a_i)`$, $`h_1(_{2k}B(a_i))`$ is described by the perturbation of the singular crossing $`a_i`$.
Thus, inside any $`2`$-disc $`B(a_i)/conj`$,
$`h_1/conj(_{2k}/conj𝐑P_\pm ^2B(a_i)/conj)`$ is determined by the relative positions of the point $`a_i`$, the point $`b`$ and $`𝐑P_\pm ^2`$. Using an argumentation similar to the one given in \[, see p.6\], in any case it is always possible to push by an isotopy the part $`𝔄_\pm `$, contained inside $`B(a_i)/conj`$ in $`B(a_i)/conj`$ leaving a $`2`$-disc inside $`B(a_i)/conj`$ if necessary.
Hence, the part $`𝔄_\pm `$, contained inside $`U`$, can be pushed by an isotopy (which coincides with $`h_t/conj`$ inside $`S^4\backslash _{i=1}^IB(a_i)/conj`$ and inside any $`B(a_i)/conj`$ satisfies the requirement above) into a $`4`$-ball $`B\mathrm{"}`$ obtained by taking union of $`B^{}`$ with disc $`B(a_i)/conj`$ or excising $`B(a_i)/conj`$. (The choice of union or excision of the disc $`B(a_i)/conj`$ depends on relative position of $`b`$ and $`𝐑P_\pm ^2B(a_i)/conj`$ \[see \]) After such isotopy, $`𝔄_\pm `$ lies in the boundary of $`B\mathrm{"}`$, except $`2k`$ discs left from $`_{2k}/conj\backslash _{i=1}^IB(a_i)/conj`$ and several discs which lie all inside $`B\mathrm{"}`$ and are unknotted. Thus from Livingston’s Theorem it follows that :
Arnold surfaces $`𝔄_\pm `$ of Harnack curve $`_{2k}`$ are standard surfaces in $`S^4`$. The decomposition announced in theorem 3.0.1 of Chapter 3.0.1 is immediate when considers the two double coverings branched respectively over $`𝔄_{}`$ and $`𝔄_+`$ (see ), However, we propose to recover this decomposition from the study of local situations inside $`4`$-balls $`B(a_i)`$ centered in crossings points $`a_i`$ perturbed in maximal simple deformation of the Harnack curves.
###### Lemma 3.0.2.
Let $`𝒞_0`$ be a singular curve of degree $`m`$ with one singular crossing. Consider a variation $`𝒞_t`$, $`t[1,1]`$ which crosses transversally the discriminant $`\mathrm{\Delta }_m`$ at $`𝒞_0`$ Let $`p`$ be the singular crossing of $`𝒞_0`$. Let $`𝔄_t`$ denote either the positive or negative Arnold surface of $`𝒞_t`$. Assume that inside a small neighborhood of $`p`$:
1. For $`t0`$, $`𝔄_t`$ is the union of $`𝒞_t/conj`$ and a disc component.
2. For $`t>0`$, $`𝔄_t`$ is the union of $`𝒞_t/conj`$ and a Möbius component. Then the perturbation $`𝒞_0`$ to $`𝒞_1`$ corresponds on Arnold surface to $`\text{}\overline{𝐑P^2}`$ (The direction of twisting of the Möbius band is obviously standard.)
proof: It can be easily deduced from an algebraic model. In a neighborhood of $`p`$, one can choose coordinates such that $`𝒞_t`$ is given by $`x_1^2x_2^2=t`$, the projection $`q:𝐂P^2𝐂P^2\backslash conj`$ is given by the map $`𝐂^2𝐂^2`$ $`(x_1,x_2)(x_1^2,x_2)`$. Q.E.D
###### Lemma 3.0.3.
Let $`𝒞_0`$ be a singular curve of degree $`2k`$ with two crossings $`p_1`$ and $`p_2`$ such that each $`p_i`$ $`i\{1,2\}`$ belongs to a real branch and an oval (with a given orientation);the two real branches which contain $`p_i`$ have opposite orientation. (see figure 3.1) Consider the deformation $`𝒞_t`$, $`t[0,1]`$ which crosses transversally the discriminant $`𝐑𝒟_m`$ at $`𝒞_0`$. Then the deformation $`𝒞_0`$ to $`𝒞_1`$ implies $`𝔄_1=𝔄_0\text{}𝐑P^2\text{}\overline{𝐑P^2}`$.
proof:
In case of Arnold surface $`𝔄_{}`$, it may be deduced from the preceding Lemma 3.0.2 of Chapter 3.0.2. It is easy to see that the directions of twisting of the Möbius band given by $`p_1`$ and $`p_2`$ are opposite.
In case of Arnold surface $`𝔄_+`$, consider the Morse-function $`f:𝒞_0\times ]1,+1[𝒞_{t,t]1,1[}]1,+1[`$, $`f(𝒞_t)=t`$.
Obviously, up to isotopy, one can identify $`𝒞_1`$ with any curve $`𝒞_{1ϵ}`$ $`ϵ>0`$ small. Let us fix such $`ϵ>0`$.
Consider the descending one-manifolds $`D_t`$ of $`p_1`$ and $`D_t^{}`$ of $`p_2`$. Each of them reaches the boundary of the quotient curve $`𝒞_{1ϵ}\backslash conj`$ and defines a normal bundle $`N`$ of the real set point $`𝐑𝒞_{1ϵ}`$ of $`𝒞_{1ϵ}`$. Choose a smooth tangent vector field $`V`$ on $`𝐑P^2`$ such that on $`𝐑𝒞_{1ϵ}`$ it is to tangent to $`𝐑𝒞_{1ϵ}`$ and directed according to the orientation of $`𝐑𝒞_{1ϵ}`$. In such a way, for $`x𝐑𝒞_{1ϵ}`$, $`N(x)=iV(x)`$ is directed inside the half of $`𝐂𝒞_{1ϵ}`$ which induces orientation on $`𝐑𝒞_{1ϵ}`$. Recall that multiplication by $`i`$ makes a real vector normal to the real plane and leaves any vector tangent to $`𝐑𝒞_{1ϵ}`$ tangent to $`𝐂𝒞_{1ϵ}`$. Extend the tangent vector field $`V`$ to a tangent vector field of $`𝒞_{1ϵ+s}`$,$`0<s<ϵ,s+e<1`$, in a neighborhood of $`𝐑𝒞_{1ϵ}𝐑P_+^2`$ (where $`𝐑P_+^2=\{x𝐑P^2|C_{1ϵ}(x)0\}`$ with $`C_{1ϵ}`$ is a polynomial giving $`𝒞_{1ϵ}`$.)
Since complex conjugation acts on gradient-trajectories of $`D_t`$ (resp $`D_t^{}`$) as symmetry of center $`p_1`$ (resp, $`p_2`$), the perturbation of $`𝒞_0`$ to $`𝒞_{1ϵ}`$ implies on Arnold surface $`𝔄_+`$ one connected sum with the fibre bundle of the fibration with base space $`𝐑𝒞_{1ϵ}`$ and fiber defined for any point $`x𝐑𝒞_{1ϵ}`$ as the gradient-trajectory through $`x`$ quotiented by the action of complex conjugation.
Namely, it adds $`\text{}𝐑P^2\text{}\overline{𝐑P^2}`$ on Arnold surface $`𝔄_+`$. Obviously, an argument analogous to this last one may be also used to describe Arnold surfaces $`𝔄_{}`$.
Q.E.D
From the Lemmas 3.0.2 of Chapter 3.0.2 and 3.0.3 of Chapter 3.0.3, we shall get topological effect of perturbations of crossings in the recursive construction of Harnack curves on Arnold surfaces and thus the description given in Theorem 3.0.1 of Chapter 3.0.1. We shall use the results of proposition 2.3.6 of Chapter 2.3.6. and (according to the notations introduced in proposition 2.3.6 of Chapter 2.3.6) we shall denote by $`S\mathrm{"}_{2k}`$ the complementary subset of $`S_{2k}^{}`$ inside $`S`$. It consists of the $`(k2)`$ points of $`S_{2k}^+\backslash S_{2k}^{}`$.
Let $`A_{2k}=\{a_1,\mathrm{},a_{2k1},\mathrm{},a_{4k2}\}`$ be the set of points perturbed in a maximal simple deformation of $`_{2k}`$ (i.e the ordered set of points $`a_i`$ where $`a_1,\mathrm{}a_{2k1}`$ are the crossings of $`_{2k1}L`$ and $`a_{2k},\mathrm{},a_{4k2}`$ are the crossings of $`B_{2k;1}`$). Let $`A_{2k+1}=\{a_1,\mathrm{},a_{2k}\}`$ be the set of points perturbed in a maximal simple deformation of $`_{2k+1}`$ (i.e the set of crossings $`a_i`$ of $`_{2k}L`$).
For any $`k1`$, denote $`𝔄_{}^{2k}`$ and $`𝔄_+^{2k}`$ the Arnold surfaces. of the Harnack curve $`_{2k}`$ of degree $`2k`$.
For $`k=1`$, it follows easily from the Lemma 3.0.2 of Chapter 3.0.2 that: $`𝔄_{}^2\overline{𝐑P^2}`$ and $`𝔄_+^2D^2`$
For $`k>1`$, pairs $`(S^4,𝔄_{}^{2k})`$, $`(S^4,𝔄_+^{2k})`$ may be deduced by induction on $`k`$.
Let $`_{2k2}`$ be the Harnack curve of even degree $`2k21`$ and $`𝔄_{}^{2k2}`$ and $`𝔄_+^{2k2}`$ be its Arnold surfaces.
Let $`_{2k}`$ be the Harnack curve of even degree $`2k`$. and $`𝔄_{}^{2k}`$ and $`𝔄_+^{2k}`$ be its Arnold surfaces.
Then, one can deduce the pair $`(S^4,𝔄_{}^{2k})`$, (resp,$`(S^4,𝔄_+^{2k})`$) from the pair $`(S^4,𝔄_{}^{2k2})`$ (resp,$`(S^4,𝔄_+^{2k2})`$).
The Harnack curve $`_{2k}`$ is obtained from $`_{2k2}`$ by intermediate construction of $`_{2k1}`$ (see proposition 2.3.1 of Chapter 2.3.1 and theorem 2.3.9 of Chapter 2.3.9 of this chapter). From Lemma 3.0.2 of Chapter 3.0.2, perturbations of singular crossings of $`_{2k2}L`$ and then of $`_{2k1}L`$ imply $`(2k2)+(2k1)`$ $`\text{}\overline{𝐑P^2}`$ on $`𝔄_{}^{2k}`$, and imply $`(2k2)+(2k1)`$ $`\text{}D^2`$ on $`𝔄_+^{2k}`$.
From the Lemma 3.0.2 of Chapter 3.0.2 and Lemma 3.0.3 of Chapter 3.0.3, keeping the notations of the proposition 2.3.1 of Chapter 2.3.1, it follows that perturbations of singular crossings of $`S_{2k}^{}S_{2k}^+`$ imply $`(k1)`$ $`\text{}\overline{𝐑P^2}`$ on $`𝔄_{}^{2k}`$, and imply $`(k1)`$ $`\text{}D^2`$ on $`𝔄_+^{2k}`$. Furthermore, the perturbations of the $`(k2)`$ singular crossings of $`S\mathrm{"}_{2k}S_{2k}^+`$ imply
(joined with perturbations of singular crossings of $`S_{2k2}^{}S_{2k2}^+`$)
$`(k2)`$ $`\text{}𝐑P^2`$ on $`𝔄_{}^{2k}`$ From the Lemma 3.0.3 of Chapter 3.0.3, it follows that perturbations of singular crossings of
$`S\mathrm{"}_{2k}S_{2k}^+`$ joined with perturbations of singular crossings of $`S_{2k2}^{}S_{2k2}^+`$ imply $`(k2)`$ $`\text{}𝐑P^2\text{}\overline{𝐑P^2}`$ on $`𝔄_+^{2k}`$.
Hence $`𝔄_{}^{2k}𝔄_{}^{2k2}\text{}(k2)\text{}𝐑P^2\text{}(5k4)\overline{𝐑P^2}`$ and $`𝔄_+^{2k}𝔄_+^{2k2}\text{}(k2)\text{}𝐑P^2\text{}(k2)\overline{𝐑P^2}`$.
Q.E.D
## Part III Perestroika Theory on Harnack curves
As already noticed in the last section, the classification problem of pairs $`(S^4,𝒜_\pm )`$ amounts to the classification problem of real algebraic curves up to conjugate equivariant isotopy of $`𝐂P^2`$. Thus, there is an obvious connection with classification of Arnold surfaces up to isotopy of $`S^4`$ and Hilbert’s Sixteen Problem on the arrangements of ovals of real algebraic curves.
First of all we shall detail the construction of Harnack curves. Then, we shall deduce from this detailed construction, a construction of curves $`𝒜_m`$ of degree $`m`$ which provides a description of the pair $`(𝐂P^2,𝐂𝒜_m)`$ up to conj-equivariant isotopy of $`𝐂P^2`$ (see Theorem 5.2.9 of Chapter 5.2.9 and Theorem 5.3.7 of Chapter 5.3.7). Moreover, this method gives all possible arrangements of real connected components of curves and therefore is an advancement in the Hilbert’s Sixteen Problem. After that, we shall deal with Arnold surfaces defined on any curve of even degree with non-empty real part.
### Chapter 4 Recursive Morse-Petrovskii’s theory of recursive Harnack curves
In this section, we shall go on the study of critical points of Harnack’s polynomial initiated in chapter 2.
Recall that given a Harnack polynomial $`B_m(x_0,x_1,x_2)=x_0^m.b_m(x_1/x_0,x_2/x_0)`$ we call critical point of $`B_m(x_0,x_1,x_2)`$ any point $`(x_0,y_0)`$ such that $`b_x(x_0,y_0)=0`$ and $`b_y(x_0,y_0)=0`$. We shall consider only Harnack curves given by regular polynomials. According to the classification of Harnack curves $`_m`$ up to rigid isotopy (Theorem 2.2.1), one can assume that any Harnack curve $`_m`$ results of the recursive construction of Harnack curves $`_i`$, $`1im`$, where the curve $`_{i+1}`$ is deduced from classical small perturbation of the union $`_iL`$ of the curve $`_i`$ with a line $`L`$. This section until its end is devoted to settle how critical points of a Harnack polynomial of degree $`m`$ may be associated with crossings of curves $`_{mi}L`$, $`1i(m1)`$, and how a critical point associated with a crossing of $`_{mi}L`$, $`1i(m1)`$, varies in the recursive construction of Harnack curves $`_m`$.
Given a regular Harnack polynomial $`B_m(x_0,x_1,x_2)`$ with affine associated polynomial $`b_m(x_1/x_0,x_2/x_0)`$, we shall study the pencil of curves given by polynomial $`x_0^m.(b_m(x_1/x_0,x_2/x_0)c)`$, $`c𝐑`$. Since Harnack curves realize the maximal number of real components curves of the pencil may have, to each critical point of index 1 of a regular Harnack polynomial corresponds a ”gluing” of real components in the pencil. These gluings are the subject of this section.
The main result of this section is gathered in the Lemma 4.2.6 of Chapter 4.2.6, Lemma 4.2.7 of Chapter 4.2.7 and Lemma 4.2.9 of Chapter 4.2.9 where we describe critical points of index $`1`$ of Harnack polynomials.
#### 4.1. Preliminaries
Given a regular polynomial $`R(x_0,x_1,x_2)=x_0^m.r(x_1/x_0,x_2/x_0)`$, for any real finite critical point $`(x_0,y_0)`$ of $`r(x,y)`$ , $`r(x_0,y_0)=c_0`$, there exists $`ϵ>0`$ sufficiently small such that $`r^1[c_0ϵ,c_0+ϵ]`$ contains no critical point other than $`(x_0,y_0)`$.
We shall call topological meaning of $`(x_0,y_0)`$ the homeomorphism type of the following triad of spaces $`(W(x_0,y_0);_0W(x_0,y_0),_1W(x_0,y_0))`$ with
$$W(x_0,y_0)=_{c[c_0ϵ,c_0+ϵ]}\{(x_0:x_1:x_2)𝐂P^2|x_0^m(r(x_1/x_0,x_2/x_0)c))=0\}$$
$$_0W(x_0,y_0)=\{(x_0:x_1:x_2)𝐂P^2|x_0^m(r(x_1/x_0,x_2/x_0)(c_0ϵ))=0\}$$
$$_1W(x_0,y_0)=\{(x_0:x_1:x_2)𝐂P^2|x_0^m(r(x_1/x_0,x_2/x_0)(c_0+ϵ))=0\}$$
Let $`B(x_0,y_0)𝐂P^2`$ be a small 4-ball around $`(x_0,y_0)`$ globally invariant by complex conjugation such that $`_0W(x_0,y_0)B(x_0,y_0)\mathrm{}`$ and $`_1W(x_0,y_0)B(x_0,y_0)\mathrm{}`$
We shall call local topological meaning of $`(x_0,y_0)`$ the homeomorphism type of the following triad of spaces
$`(W(x_0,y_0)B(x_0,y_0);_0W(x_0,y_0)B(x_0,y_0),_1W(x_0,y_0)B(x_0,y_0))`$
In what follows, we shall associate to each critical point $`p`$ of index 1 of a Harnack polynomial of degree $`m`$, two real branches of the Harnack curve of degree $`m`$ involved in the local topological meaning of $`p`$.
Terminology Recall that we distinguish two ways to perturb a crossing $`p`$ of a singular curve $`𝒜`$ defined in the section 2.3 of Chapter 2.3 as perturbation of type 1 and perturbation of type 2 of $`p`$.
Thus, local topological meaning of a crossing $`(x_0,y_0)`$ can be deduced from one topological space $`_1W(x_0,y_0)B(x_0,y_0)`$ or $`_0W(x_0,y_0)B(x_0,y_0)`$ and the type $`1`$ or $`2`$ of the perturbation of the point $`(x_0,y_0)`$. Let us introduce the notations: $`(_1W(x_0,y_0)B(x_0,y_0))_i`$ and $`(_0W(x_0,y_0)B(x_0,y_0))_i`$ where the subscript $`i`$ stands for the type of the perturbation of the crossing. Although any perturbation of the union of two projective lines leads to a conic with orientable real part, the perturbation of a singular curve of degree $`m>2`$ of which singular points are crossings may lead to a curve with non-orientable part. Nonetheless, in case of a Harnack curve $`_m`$, any deformation of a crossing which appears in a maximal simple deformation of $`_m`$ (see Proposition 2.3.6 of Chapter 2.3.6 and Proposition 2.3.7 of Chapter 2.3.7) agrees with a complex orientation of its real part $`𝐑_m`$. Besides, it is easy to deduce from relative orientation and location of real branches of $`𝐑_m`$ that any such crossing is deformed by a perturbation of type 1.
#### 4.2. Critical points and recursive construction of Harnack curves
##### Introduction
Let $`_m`$ be the Harnack curve of degree $`m`$. In what follows, for any $`m>0`$, we shall consider only Harnack polynomials $`B_m(x_0,x_1,x_2)=x_0^m.b_m(x_1/x_0,x_2/x_0)`$ of degree $`m`$ and type $`^0`$. Since no confusion is possible, we shall call such polynomials Harnack polynomials. As already introduced in the preceding section (see proof of propositions 2.3.6 of Chapter 2.3.6 and 2.3.1 of Chapter 2.3.1), one can associate critical points of index $`1`$ of $`B_m`$ with crossings of curves $`_{m1}L`$. Formally, we shall say that:
###### Definition 4.2.1.
Let $`a`$ be a crossing of the curve $`_{mi}L`$. Given a disc $`D(a,ϵ)`$ (in the Fubini-Study metric) of radius $`ϵ`$ around of $`a`$ in $`𝐑P^2`$.
Let $`U(a,ϵ)`$ be the neighborhood of $`_{mi+1}`$ in $`𝐂P^2`$ deduced from the action of $`U_𝐂^3`$ on $`U(a)`$:
$$U(a,ϵ)=\{z=<u,p>=(u_0.p_0:u_1.p_1:u_2.p_2)𝐂P^2u=(u_0,u_1,u_2)U_𝐂^3,pU(a)\}$$
.
A critical point $`(x_0,y_0)`$ of index 1 of $`B_m`$ is associated with a crossing $`a`$ of the curve $`_{mi}L`$ if there exists $`ϵ_0`$ such that:
1. $`a,(x_0,y_0)U(a,ϵ_0)`$
2. the perturbation on the real part of $`_m`$ involved in the local topological meaning of $`(x_0,y_0)`$ is a deformation on $`𝐂_mU(a,ϵ_0)`$
As already noticed corollary 2.3.5 of Chapter 2.3.5 of the previous section 2.3 of Chapter 2.3, without loss of generality, one can consider that Harnack curves are obtained via the Patchworking method.
Hence, for sake of simplicity, we shall consider the $`T`$-Harnack curves introduced in the Chapter 1. For the patchworking construction of Harnack we can give the following of version definition 4.2.1.
Recall that $`\rho ^m:𝐑_+T_m\times U_𝐂^2𝐂T_m𝐂P^2`$ and its restriction $`\rho ^m|_{𝐑_+T_m\times U_𝐑^2}:𝐑_+T_m\times U_𝐑^2𝐑T_m𝐑P^2`$ denote the natural surjections.
Recall (see definition 0.1.8 of Chapter 0.1.8) that given $`\mathrm{\Gamma }`$ a face of the triangulation of $`T_m`$ and $`D(p,ϵ)𝐑_+T_m^0`$ an (euclidian) open 2-disc such that the moment map $`\mu :𝐂T_mT_m`$ maps $`D(p,ϵ)`$ to a two-disc $`\mu (D(p,ϵ)`$ which contains $`\mathrm{\Gamma }`$ and intersects only the face $`\mathrm{\Gamma }`$ of the triangulation of $`T_m`$; we call $`U(p)=\rho ^m(D(p,ϵ)\times U_𝐂^2)`$ the $`ϵ`$-neighborhood of $`𝐂_m`$ in $`𝐂T_m𝐂P^2`$ defined from $`\mathrm{\Gamma }^0`$; we call $`ϵ`$-tubular neighborhood of $`𝐂_m`$ defined from $`\mathrm{\Gamma }^0`$ the $`ϵ`$-tubular neighborhood of $`𝐂_m`$ in $`U(p)`$.
###### Definition 4.2.2.
Let $`\mathrm{\Gamma }`$ be a face of the triangulation of $`T_m`$ and $`U(p)`$ be the $`ϵ`$-neighborhood of $`𝐂_m`$ defined from $`\mathrm{\Gamma }^0`$. A critical point $`(x_0,y_0)`$ of index 1 of $`B_m`$ is associated with a crossing $`a`$ of the curve $`_{mi}L`$ if :
1. $`a,(x_0,y_0)U(p)`$
2. the perturbation on the real part of $`_m`$ involved in the local topological meaning of $`(x_0,y_0)`$ is a deformation on $`𝐂_mU(p)`$
###### Definition 4.2.3.
Let $`p_m`$ (resp $`p_{m+j}`$, $`j>0`$) be a critical point of index 1 of a Harnack polynomial of degree $`m`$ (resp $`m+j`$). We say that $`p_{m+j}`$ is equivalent to $`p_m`$ if there exist small conj-equivariant open $`4`$-balls $`B(p_{m+j})`$ and $`B(p_m)`$ around $`p_{m+j}`$ and $`p_m`$ with the following properties:
1. $`B(p_{m+j})B(p_m)`$
2. The triad of spaces
$$(W(p_{m+j}B(p_{m+j});_0W(p_{m+j})B(p_{m+j}),_1W(p_{m+j})B(p_{m+j}))$$
is the local topological meaning of $`p_{m+j}`$. The triad of spaces
$$(W(p_mB(p_m);_0W(p_m)B(p_m),_1W(p_m)B(p_m))$$
is the local topological meaning of $`p_{m+j}`$.
3. Local topological meanings of $`p_{m+j}`$ and $`p_m`$ are homeomorphic.
###### Remark 4.2.4.
Let $`B_m(x_0,x_1,x_2)`$ be a $`T`$-Harnack polynomial. It is an easy consequence of the $`T`$-construction of Harnack curves, that one can define $`4`$-ball of $`𝐂P^2`$ around any critical point of $`B_m(x_0,x_1,x_2)`$ as usual $`4`$-ball of $`(𝐂^{})^2`$. Indeed, none of the critical points of $`B_m(x_0,x_1,x_2)`$ belongs to the coordinates axes.
We shall denote $`p_{m+j}`$ is equivalent to $`p_m`$ by $`p_{m+j}p_m`$.
###### Definition 4.2.5.
Given $`a`$ a crossing of $`_{m1}L`$, denote by $`p_l`$ the critical point of index 1 of a Harnack polynomial of degree $`lm`$ associated with $`a`$.
1. We say that $`p_m`$ is the simple point of $`_m`$ associated with $`a`$.
2. If for any $`j1`$, $`p_{m+j}`$ is equivalent to $`p_m`$, we say that a simple point of $`_{m+j}`$ with $`j0`$ is associated with $`a`$.
3. If for any integers $`l`$,$`l^{}`$ ,$`0<ll^{}k`$, $`p_{m+l}`$ is not equivalent to $`p_{m+l^{}}`$; for any $`lk`$, we say that a $`l`$-point $`(p_m,p_{m+1},\mathrm{},p_{m+l})`$ of $`_{m+l}`$ is associated with $`a`$.
4. If $`k`$ is the smallest integer such that: for any integers $`l,l^{}`$,
$`0<ll^{}k`$, $`p_{m+l}`$ is not equivalent to $`p_{m+l^{}}`$; for any $`lk`$, we say that a $`k`$-point $`(p_m,p_{m+1},\mathrm{},p_{m+k})`$ of $`_{m+l}`$ is associated with $`a`$.
##### Critical points of index $`1`$ of Harnack polynomials
We shall now proceed to the study of critical points of index $`1`$ of Harnack polynomials. We shall work with the notations introduced in the chapter 2. Given $`B_m(x_0,x_1,x_2)=x_0^m.b_m(x_1/x_0,x_2/x_0)`$ a Harnack polynomial of type $`^0`$, we denote by $`S_m`$ the set of critical points of index 1 of $`b_m(x,y)`$, and by $`S_m^{}`$ (resp, $`S_m^+`$) the subset of $`S_m`$ consisting respectively of critical points of index 1 with negative (resp, positive) critical value.
In case $`m=2k`$, we distinguish the two subsets $`S_{2k}^{}`$, $`S\mathrm{"}_{2k}`$ of $`S_{2k}^+`$ with the properties $`S_{2k}^{}S\mathrm{"}_{2k}=S_{2k}^+`$ $`S_{2k}^{}S\mathrm{"}_{2k}=\mathrm{}`$. Given $`B_{2k}(x_0,x_1,x_2)`$ a Harnack polynomial of type $`^0`$ and $`b_{2k}(x_1/x_0,x_2/x_0)`$ its affine associated polynomial. The set $`S_{2k}^{}`$ is constituted by the $`c_1^{}(B_{2k})`$ critical points $`(x_0,y_0)`$ of index 1 $`b_{2k}(x_0,y_0)=c_0`$, $`c_0>0`$ with the property that as $`c`$ increases from $`c_0ϵ`$ to $`c_0+ϵ`$ the number of real connected components of $`𝒜_c=\{(x_0:x_1:x_2)𝐑P^2|x_0^{2k}.(b_{2k}c)=0\}`$ intersecting the line at infinity increases by 1. The subset $`S\mathrm{"}_{2k}`$ denotes the complementary set of $`S_{2k}^{}`$ inside $`S_{2k}^+`$. These sets were already under consideration in propositions 2.1.4 of Chapter 2.1.4, 2.3.6 of Chapter 2.3.6 and theorem 3.0.1 of Chapter 3.0.1. Let us start by the Lemma 4.2.6 of Chapter 4.2.6 in which we prove that critical points of negative critical value of Harnack polynomials of type $`^0`$ are simple points in the recursive construction of Harnack curves.
###### Lemma 4.2.6.
Let $`_m`$ be a Harnack curve of degree $`m`$ obtained via the patchworking method and given by a Harnack polynomial of type $`^0`$.
1. In case of odd $`m`$, let $`\mathrm{\Gamma }`$ be a face of the triangulation of $`T_m`$ which belongs to $`l_{m1}`$.
2. In case of even $`m`$, let $`\mathrm{\Gamma }`$ given by vertices $`(c,d+1)(c+1,d)`$ with $`c=0`$ or $`d=0`$.
Let $`U(p)`$ be the $`ϵ`$-neighborhood of $`𝐂_m`$ in $`𝐂P^2`$ defined from $`\mathrm{\Gamma }^0`$. Denote $`a`$ the unique crossing of $`_{m1}L`$ which belongs to $`U(p)`$. Denote $`p_m`$ the critical point of $`_m`$ associated with $`a`$, and let $`B(p_m)U(p)`$ a conj-equivariant $`4`$-ball around $`p_m`$. Then,
1. the critical point $`p_m`$ of $`_m`$ belongs to $`S_m^{}`$
Moreover, $`_mB(p_m)(_1W(p_m)B(p_m))_1`$
2. a simple-point of curves $`_{m+j}`$, $`j0`$, is associated with $`a`$
proof: The proof is based on the Petrovskii’s theory and propositions 2.1.4 of Chapter 2.1.4, 2.3.6 of Chapter 2.3.6 and 2.3.1 of Chapter 2.3.1. Our argumentation is similar to the one used in the proof of proposition 2.3.1 of Chapter 2.3.1.
Assume $`\mathrm{\Gamma }`$ given by vertices $`(c,d+1)(c+1,d)`$. For any $`j0`$, given
$`B_{m+j}=x_0^{m+j}b_{m+j}(x_1/x_0,x_2/x_0)`$ the Harnack polynomial of degree $`m+j`$, we shall denote $`b_{m+j}^S`$ the truncation of $`b_{m+j}`$ on the monomials $`x^cy^d`$,$`x^cy^{d+1}`$ $`x^{c+1}y^d`$, $`x^{c+1}y^{d+1}`$.
We shall use local description of Harnack curves, $`_{m+j}`$, $`j0`$, inside $`U(p)`$ provided by the patchworking theory.
Recall that for any $`j0`$, there exists an homeomorphism
$`\stackrel{~}{h}:𝐂_{m+j}U(p)\{(x,y)(𝐂^{})^2|b_{m+j}^S=0\}U(p)`$ such that
$`\stackrel{~}{h}(a)=(x_0,y_0)`$ is a critical point of $`b_{m+j}^S`$.
We shall divide our proof in two parts. In the first part, we shall consider the curve $`_m`$ and study local topological meaning of crossings of $`_{m1}L`$. Then, we shall consider curves $`_{m+j}`$, $`j1`$. (1) Let us distinguish the cases $`m`$ even and $`m`$ odd.
(1).a In case of even $`m=2k`$ Let $`\mathrm{\Gamma }l_{2k1}`$ given by vertices $`(c,d+1)(c+1,d)`$ with
$`c=0`$ or $`d=0`$. (Namely, $`\mathrm{\Gamma }`$ has vertices $`(0,2k1),(1,2k2)`$ or $`(2k1,0),(2k2,1)`$.) From the patchworking theory, (see corollary 1.2.4 of Chapter 1.2.4 of the chapter 1) $`\stackrel{~}{h}(a)=(x_0,y_0)`$ $`x_0>0,y_0<0`$ is a critical point of $`b_{2k}^S(x,y)`$.
Keeping the notations introduced in the chapter 2, set $`b_{2k}=x_{2k;t}`$ with $`t]0,t_{2k}[`$. Then, for fixed $`t]0,t_{2k}[`$, let $`a_t`$ be the crossing of $`_{2k1}L`$ contained in $`U(p)`$.
Therefore, $`\stackrel{~}{h}(a_t)=(x_0,y_0)_t`$ $`x_0>y_0<0`$ is a critical point of $`x_{2k,t}^S(x,y)`$.
Set $`x_{2k;t}^S(x,y)=l_t(x,y)+t.k_t(x,y)`$ with
$`l_t(x,y)=a_{c,d}t^{\nu (c,d)}x^cy^d+a_{c+1,d}t^{\nu (c,d+1)}x^{c+1}y^d`$;
$`k_t(x,y)=a_{c,d+1}t^{\nu (c+1,d)1}x^cy^{d+1}+a_{c+1,d+1}t^{\nu (c+1,d+1)1}x^{c+1}y^{d+1}`$. Then, it is easy to see that, up to modify the coefficients
$`a_{c,d},a_{c,d+1},a_{c+1,d},a_{c+1,d}`$ if necessary, the point $`\stackrel{~}{h}(a_t)=(x_1,y_1)_t`$ is a critical point of the function $`\frac{l_t(x,y)}{k_t(x,y)}`$ with negative critical value $`t>t_{2k1}`$.
On this assumption, it follows from the equalities
$`l_t(x,y)=l_t(x,y)`$, $`k_t(x,y)=k_t(x,y)`$,
$`\frac{l_t}{x}(x,y)=\frac{l_t}{x}(x,y)`$, $`\frac{l_t}{y}(x,y)=\frac{l_t}{y}(x,y)`$,
$`\frac{k_t}{x}(x,y)=\frac{k_t}{x}(x,y)`$, $`\frac{k_t}{y}(x,y)=\frac{k_t}{y}(x,y)`$
that $`(x_1,y_1)_t`$ is a critical point of the function $`\frac{l_t(x,y)}{k_t(x,y)}`$ with positive critical value $`t<t_{2k1}`$. Therefore, up to modify coefficients of $`b_{2k}^S`$, $`(x_0,y_0)`$ is a critical point with negative critical value of $`b_{2k}^S`$. Obviously, (see ), such modification does not affect the order and the topological structure of $`b_{2k}^S`$. Therefore, it does not change local topological meaning of critical points of the Harnack polynomial $`B_{2k}`$.
We shall now proceed to the study of the Petrovskii’s pencil of curves over $`_{2k}`$.
Let $`d_{2k}`$ be the unique polynomial such that $`b_{2k}=b_{2k}^S+d_{2k}`$. Consider curves of the Petrovskii’s pencil over $`_{2k}`$ outside $`d_{2k}=0`$, as level curves of the function $`\frac{b_{2k}^Sc}{d_{2k}}`$.
Inside $`b_{2k}^Sc=0\backslash d_{2k}=0`$, these curves have critical points the singular points of $`b_{2k}^Sc=0`$.
Bringing together Petrovskii’s theory and the implicit function theorem applied to one-parameter polynomial
$`x_0^{2k}.(b_{2k}c)=x_0^{2k}.(b_{2k}^S+d_{2k}c)`$ with parameter $`c`$, it follows that a point $`p_{2k}`$ with $`b_{2k}(p_{2k})=c_0<0`$ is associated with $`a`$.
When $`c=c_0`$, one positive oval touches an other positive oval.
(1).b In case of odd $`m=2k+1`$, one has to distinguish whether $`\mathrm{\Gamma }`$ intersects or not the coordinates axes.
(1).b.1
Let $`\mathrm{\Gamma }l_{2k}`$ given by vertices $`(c,d+1)(c+1,d)`$ with
$`c=0`$ or $`d=0`$. (Namely, $`\mathrm{\Gamma }`$ has vertices $`(0,2k),(1,2k1)`$ or $`(2k,0),(2k1,1)`$.)
According to the patchworking theory, (see corollary 1.2.4 of Chapter 1.2.4 of the first section) $`\stackrel{~}{h}(a)=(x_0,y_0)`$ with $`x_0.y_0>0`$ is a critical point of $`b_{2k+1}^S(x,y)`$. Set $`\stackrel{~}{h}(a)=(x_0,y_0)x_0>0,y_0>0`$.
On this assumption, following the previous argumentation, up to modify coefficients of $`b_{2k+1}^S`$ if necessary (without changing the order and the topological structure of $`b_{2k+1}`$) $`(x_0,y_0)`$ is a critical point with negative critical value of $`b_{2k+1}^S`$.
Therefore, according to Petrovskii’s theory, from an argumentation similar to the one given in even case, a point $`p_{2k+1}`$ with $`b_{2k+1}(p_{2k+1})=c_0<0`$ is associated with $`a`$.
When $`c=c_0`$, a positive oval touches the one-side component of the curve.
(1).b.2
We shall now consider faces which do not intersect the coordinates axes. Our argumentation is a slightly modified version of the previous one. We shall study the situation inside $`ϵ`$-neighborhood defined from two faces together. Let $`\mathrm{\Gamma }l_{2k}`$ given by vertices $`(c,d+2),(c+1,d+1)`$ with $`c`$ odd, $`c0`$ and $`d0`$ and $`\mathrm{\Gamma }^{}l_{2k}`$ given by vertices $`(c+2,d),(c+1,d+1)`$. Consider the convex polygon $`KT_{2k+1}`$ with vertices
$`(c+1,d),(c+2,d),(c+2,d+1),(c+1,d+2),(c,d+2),(c,d+1)`$. It is contained in the triangle $`T_{2k+1}`$ and triangulated by the triangulation $`\tau `$ of $`T_{2k+1}`$. Denote $`U(p)`$ the $`ϵ`$-neighborhood of $`𝐂_{2k+1}`$ defined from $`\mathrm{\Gamma }^0`$. Denote $`U(p^{})`$ the $`ϵ`$-neighborhood of $`𝐂_{2k+1}`$ defined from $`\mathrm{\Gamma }_{}^{}{}_{}{}^{0}`$. Let $`a`$, (resp, $`a^{}`$) be the crossing of $`_{2k}L`$ which belongs to $`U(p)`$, (resp, $`U(p^{})`$). Denote $`b_{2k+1}^S`$ (resp, $`b_{2k+1}^S^{}`$) the truncation of $`b_{2k+1}`$ on the monomials $`x^cy^{d+1}`$,$`x^cy^{d+2}`$ $`x^{c+1}y^{d+2}`$, $`x^{c+1}y^{d+1}`$ (resp, $`x^{c+1}y^d`$,$`x^{c+1}y^{d+1}`$ $`x^{c+2}y^{d+1}`$, $`x^{c+2}y^d`$.)
As previously, consider homeomorphisms
$`\stackrel{~}{h}:𝐂_{2k+1}U(p)\{(x,y)(𝐂^{})^2|b_{2k+1}^S=0\}U(p)`$ such that $`\stackrel{~}{h}(a)=(x_0,y_0)`$ $`x_0>0,y_0>0`$ is a critical point of $`b_{2k+1}^S`$ and $`\stackrel{~}{h}^{}:𝐂_{2k+1}U(p^{})\{(x,y)(𝐂^{})^2|b_{2k+1}^S^{}=0\}U(p^{})`$ such that $`\stackrel{~}{h}^{}(a^{})=(x_0^{},y_0^{})`$ $`x_0>0,y_0>0`$ is a critical point of $`b_{2k+1}^S^{}`$. Keeping the notations introduced in the chapter 2, set $`b_{2k+1}=x_{2k;t}`$, $`t]0,t_{2k+1}[`$.
For fixed $`t]0,t_{2k+1}[`$, let $`a_t`$ (resp $`a_t^{}`$) be the crossing of $`_{2k}L`$ contained in $`U(p)`$ (resp $`U(p^{})`$). Set $`x_{2k+1;t}^K(x,y)=l_t(x,y)+t.k_t(x,y)`$ with
$`l_t(x,y)=a_{c+1,d}t^{\nu (c,d)}x^{c+1}y^d+a_{c+2,d}t^{\nu (c+2,d)}x^{c+2}y^d+a_{c+1,d+2}t^{\nu (c,d)}x^{c+1}y^{d+2}+a_{c+2,d+2}t^{\nu (c+2,d)}x^{c+2}y^{d+2}`$;
$`k_t(x,y)=a_{c,d+1}t^{\nu (c+1,d)1}x^cy^{d+1}a_{c+1,d+1}t^{\nu (c+1,d+1)1}x^{c+1}y^{d+1}+a_{c+2,d+1}t^{\nu (c+1,d)1}x^{c+2}y^{d+1}`$ Then, it is easy to see that, up to modify the coefficients $`a_{i,j}`$ if necessary, points $`\stackrel{~}{h}(a_t)=(x_0,y_0)_t`$ and $`\stackrel{~}{h^{}}(a_t^{})=(x_0^{},y_0^{})_t`$ are critical points of the function $`\frac{l_t(x,y)}{k_t(x,y)}`$ with negative critical value $`t>t_{2k1}`$. On this assumption, it follows from the equalities
$`l_t(x,y)=l_t(x,y)`$, $`k_t(x,y)=k_t(x,y)`$,
$`\frac{l_t}{x}(x,y)=\frac{l_t}{x}(x,y)`$, $`\frac{l_t}{y}(x,y)=\frac{l_t}{y}(x,y)`$,
$`\frac{k_t}{x}(x,y)=\frac{k_t}{x}(x,y)`$, $`\frac{k_t}{y}(x,y)=\frac{k_t}{y}(x,y)`$ that points $`(x_0,y_0)_t`$ and $`(x_0^{},y_0^{})_t`$ are critical point of the function $`\frac{l_t(x,y)}{k_t(x,y)}`$ with positive critical value $`t<t_{2k1}`$. Therefore, up to modify the coefficients of $`b_{2k+1}^K`$, $`(x_0,y_0)`$ and $`(x_0^{},y_0^{})`$ are critical points with negative critical value of $`b_{2k+1}^K`$. Obviously, (see ), such modification does not modify the order and the topological structure of $`b_{2k+1}^K`$. Therefore, it does not change local topological meanings of critical points of the Harnack polynomial $`B_{2k+1}`$.
We shall now proceed to the study of the Petrovskii’s pencil of curves over $`_{2k+1}`$. Let $`U(K^0)`$ be the subset $`\rho ^{2k+1}(𝐑_+K^0\times U_𝐂^2)`$ of $`𝐂P^2`$.
Obviously, $`U(p)U(p^{})U(K^0)`$. According to patchworking theory, the truncation $`b_{2k+1}^K`$ of $`b_{2k+1}`$ is $`ϵ`$-sufficient for $`b_{2k+1}`$ in $`U(K^0)`$. Therefore,
$`\stackrel{~}{h}:𝐂_{2k+1}U(p)\{(x,y)(𝐂^{})^2|b_{2k+1}^S=0\}U(p)`$
and $`\stackrel{~}{h}:𝐂_{2k+1}U(p^{})\{(x,y)(𝐂^{})^2|b_{2k+1}^S^{}=0\}U(p^{})`$ extend to the homeomorphism $`\stackrel{~}{h}:𝐂_{2k+1}U(K^0)\{(x,y)(𝐂^{})^2|b_{2k+1}^K=0\}U(K^0)`$ such that $`\stackrel{~}{h}(a)=(x_0,y_0)`$ $`x_0.y_0>0`$ $`\stackrel{~}{h}(a^{})=(x_0^{},y_0^{})`$ $`x_0^{}.y_0^{}>0`$. Let $`d_{2k+1}`$ be the unique polynomial such that $`b_{2k+1}=b_{2k+1}^K+d_{2k+1}`$. Consider curves of the Petrovskii’s pencil over $`_{2k+1}`$ outside $`d_{2k+1}=0`$, as level curves of the function $`\frac{b_{2k+1}^Kc}{d_{2k+j}}`$.
From an argumentation similar to the previous one, according to the Petrovskii’s theory, it follows that: a point $`p_{2k+1}`$ with $`b_{2k+1}(p_{2k+1})=c_0<0`$ is associated with $`a`$ and a point $`p_{2k+1}^{}`$, $`b_{2k+1}(p_{2k+1}^{})=c_0^{}<0`$ is associated with $`a^{}`$. When $`c=c_0`$, the one-side component of the curve touches itself. When $`c=c_0^{}`$, the one-side component of the curve touches itself.
(2) Let $`a`$ be one singular point of the curve $`_{m1}L`$ under consideration in the part 1 of the proof and $`p_m`$ be the critical point of the Harnack polynomial of degree $`m`$ associated with $`a`$. We shall now prove that for any such crossing $`a`$ and any $`j1`$, there exists a critical point $`p_{m+j}`$ of a Harnack polynomial $`B_{m+j}`$ associated with $`a`$. The point $`p_{m+j}`$ is equivalent to the critical point $`p_m`$ of $`B_m`$.
Let $`j1`$ and $`B_{m+j}`$ be a Harnack polynomial of degree $`m+j`$
As previously, consider the homeomorphism
$`\stackrel{~}{h}:𝐂_{m+j}U(p)\{(x,y)(𝐂^{})^2|b_{m+j}^S=0\}U(p)`$. (2).a
Assume $`\mathrm{\Gamma }`$ intersects the coordinates axes ($`\mathrm{\Gamma }l_{2k1}`$ or $`\mathrm{\Gamma }l_{2k}`$) Let $`d_{m+j}`$ be the unique polynomial such that $`b_{m+j}=b_{m+j}^S+d_{m+j}`$. Consider curves of the Petrovskii’s pencil over $`_{m+j}`$ outside $`d_{m+j}=0`$, as level curves of the function $`\frac{b_{m+j}^Sc}{d_{m+j}}`$. Inside $`b_{m+j}^Sc=0\backslash d_{m+j}=0`$, these curves have critical points the singular points of $`b_{m+j}^Sc=0`$.
From an argumentation similar to the previous one, according to the Petrovskii’s theory (Petrovskii’s pencil over $`_{m+j}`$), it follows easily that for any $`j1`$, a point $`p_{m+j}`$ is associated with $`a`$ and is equivalent to the point $`p_m`$.
(2).b Assume $`\mathrm{\Gamma }`$ does not intersect the coordinates axes ($`\mathrm{\Gamma }l_{2k}`$). Consider as previously the convex polygon $`K`$ with vertices
$`(c+1,d),(c+2,d),(c+2,d+1),(c+1,d+2),(c,d+2),(c,d+1)`$. Let $`d_{m+j}`$ be the unique polynomial such that $`b_{m+j}=b_{m+j}^K+d_{m+j}`$.
Consider curves of the Petrovskii’s pencil over $`_{m+j}`$ outside $`d_{m+j}=0`$, as level curves of the function $`\frac{b_{m+j}^Kc}{d_{m+j}}`$. From an argumentation similar to the previous one, according to the Petrovskii’s theory, it follows that for any $`j1`$, a point $`p_{m+j}`$ is associated with $`a`$ and is equivalent to the point $`p_m`$. Q.E.D. Then, consider the case when a critical point $`p_m`$ of $`S_m^+`$ is associated to a crossing of $`_{m1}L`$. Such critical point appears the first time as critical point of a Harnack polynomial of degree $`4`$. Recall (see the Proposition 2.1.4 of Chapter 2.1.4 ) that for a Harnack polynomial $`B_{2k}(x_0,x_1,x_2)=x_0^{2k}.b_{2k}(x_1/x_0,x_2/x_0)`$ of type $`^0`$ we denote $`S_{2k}^{}`$ the set constituted by the $`c_1^{}(B_{2k})`$ critical points $`(x_0,y_0)`$ of index 1 $`b_{2k}(x_0,y_0)=c_0`$, $`c_0>0`$ with the property that as $`c`$ increases from $`c_0ϵ`$ to $`c_0+ϵ`$ the number of real connected components intersecting the line at infinity increases by $`1`$. In the Lemma 4.2.7 of Chapter 4.2.7, we shall study how critical points $`S_{2k}^{}`$ of a Harnack polynomial of type $`^0`$ and degree $`2k`$ vary in the recursive construction of Harnack curves.
###### Lemma 4.2.7.
Let $`_{2k}`$ be the Harnack curve of degree $`2k`$ obtained via the patchworking method and given by a Harnack polynomial of type $`^0`$. Let $`\mathrm{\Gamma }`$ be a face of the triangulation of $`T_{2k}`$ contained into the line $`l_{2k1}`$ and given by vertices $`(c+1,d),(c,d+1)`$ with $`c`$ and $`d`$ odd. Let $`U(p)`$ be the $`ϵ`$-neighborhood of $`𝐂_m`$ defined from $`\mathrm{\Gamma }^0`$. Denote $`a`$ the unique crossing of $`_{2k1}L`$ such that $`aU(p)`$. Then, a $`3`$-point $`(p_{2k},p_{2k+1},p_{2k+2})`$ of Harnack curves $`_{2k+j}`$ ,$`j3`$, is associated with the point $`a`$.
1. The point $`p_{2k}S_{2k}^{}S_{2k}^+`$. Besides, there exists a positive oval $`𝒪`$ such that the point $`p_{2k}S_{2k}^+`$ has the property:
$$𝒪_1W(p_{2k}),𝒪_0W(p_{2k})$$
Moreover, given $`B(p_{2k})`$ a conj-equivariant $`4`$-ball around $`p_{2k}`$,
$$_{2k}B(p_{2k})(_0W(p_{2k})B(p_{2k}))_1$$
2. There exists a negative oval $`𝒪`$ such that the point
$`p_{2k+1}S_{2k+1}^+`$ has the property:
$$𝒪_0W(p_{2k+1}),𝒪_1W(p_{2k+1})$$
Moreover, given $`B(p_{2k+1})`$ a conj-equivariant $`4`$-ball around $`p_{2k+1}`$,
$$_{2k+1}B(p_{2k+1})(_0W(p_{2k+1})B(p_{2k+1}))_1$$
3. There exists a positive oval $`𝒪`$ of $`_{2k+2}`$ such that the point
$`p_{2k+2}S_{2k+2}^{}`$ has the property:
$$𝒪_1W(p_{2k+2}),𝒪_0W(p_{2k+2})$$
Moreover, given $`B(p_{2k+2})`$ a conj-equivariant $`4`$-ball around $`p_{2k+2}`$,
$$_{2k+2}B(p_{2k+2})(_1W(p_{2k+2})B(p_{2k+2}))_1$$
###### Definition 4.2.8.
We shall call $`3`$-point of the first kind a $`3`$-point verifying all assumptions of Lemma 4.2.7 of Chapter 4.2.7.
proof:
As in the proof of Lemma 4.2.6 of Chapter 4.2.6, we shall use local description of Harnack curves $`_{m+j}`$, $`j0`$, inside $`U(p)`$ provided by the patchworking theory.
Recall that for any $`j0`$, there exists an homeomorphism
$`\stackrel{~}{h}:𝐂_{m+j}U(p)\{(x,y)(𝐂^{})^2|b_{m+j}^S=0\}U(p)`$ such that $`\stackrel{~}{h}(a)=(x_0,y_0)`$ $`x_0.y_0<0`$ is a critical point of $`b_{m+j}^S`$.
Our proof uses the results of chapter 2 and in particular proposition 2.1.4 of Chapter 2.1.4, and proposition 2.3.6 of Chapter 2.3.6. We shall work with notations introduced in the chapter 2. (1) Let $`_{2k}`$ the Harnack curve of degree $`2k`$ $`(k>2)`$ and $`B_{2k}`$ be a Harnack polynomial of degree $`2k`$. According to proposition 2.3.6 of Chapter 2.3.6, a point $`p_{2k}S_{2k}^{}`$ is associated with $`a`$.
Set $`b_{2k}(p_{2k})=c_0>0`$, then when $`c=c_0`$, the non-empty positive oval touches a positive outer oval. (2) Let $`_{2k+1}`$ be the Harnack curve of degree $`2k+1`$ and $`B_{2k+1}`$ be a Harnack polynomial of degree $`2k+1`$.
Consider the convex polygon $`KT_{2k+1}`$ with vertices $`(c,d),(c,d+2),(c+1,d+2),(c+2,d),(c+2,d+1)`$. It is contained in the triangle $`T_{2k+1}`$ and triangulated by the triangulation $`\tau `$ of $`T_{2k+1}`$.
Let $`U(K^0)`$ be the subset $`\rho ^{2k+1}(𝐑_+K^0\times U_𝐂^2)`$ of $`𝐂P^2`$. Obviously, $`U(p)U(K^0)`$.
Denote $`b_{2k+1}^K`$ the truncation of $`b_{2k+1}`$ on the monomials $`x^cy^d`$, $`x^cy^{d+1}`$, $`x^cy^{d+2}`$, $`x^{c+1}y^d`$, $`x^{c+1}y^{d+1}`$, $`x^{c+1}y^{d+2}`$, $`x^{c+2}y^d`$, $`x^{c+2}y^{d+1}`$. (Namely, $`b_{2k+1}^K=b_{2k+1}^S+a_{c,d+2}x^cy^{d+2}+a_{c+1,d+2}x^{c+1}y^{d+2}+a_{c+2,d}x^{c+2}y^d+a_{c+2,d+1}x^{c+2}y^{d+1}`$ with $`a_{c,d+2}>0,a_{c+1,d+2}>0,a_{c+2,d}>0,a_{c+2,d+1}>0`$.
According to the patchworking theory, the truncation $`b_{2k+1}^K`$ of $`b_{2k+1}`$ is $`ϵ`$-sufficient for $`b_{2k+1}`$ in $`U(K^0)`$.
Therefore,
$`\stackrel{~}{h}:𝐂_{2k+1}U(p)\{(x,y)(𝐂^{})^2|b_{2k+1}^S=0\}U(p)`$ extends to the homeomorphism $`\stackrel{~}{h}:𝐂_{2k+1}U(K^0)\{(x,y)(𝐂^{})^2|b_{2k+1}^K=0\}U(K^0)`$ such that $`\stackrel{~}{h}(a)=(x_0,y_0)`$ $`x_0.y_0<0`$.
According to proposition 2.3.1 of Chapter 2.3.1, one can assume that $`(x_0,y_0)`$, $`x_0.y_0>0`$, $`x_0<0,y_0<0`$, is a critical point of $`b_{2k+1}^S=b_{2k}^S`$ with positive critical value.
Thus, up to modify coefficients of the polynomials $`b_{2k+1}^K`$ (without changing the order and the topological structure of $`b_{2k+1}`$), the point $`(x_0,y_0)`$ is a critical point of $`b_{2k+1}^K`$ with $`b_{2k+1}(x_0,y_0)>0`$.
Hence, we shall proceed to the study of the Petrovskii’s pencil of curves over $`_{2k+1}`$. Let $`d_{2k+1}`$ be the unique polynomial such that $`b_{2k+1}=b_{2k+1}^K+d_{2k+1}`$. Consider curves of the Petrovskii’s pencil over $`_{2k+1}`$ outside $`d_{2k+1}=0`$, as level curves of the function $`\frac{b_{2k+1}^Kc}{d_{2k+1}}`$.
Inside $`b_{2k+1}^Kc=0\backslash d_{2k+1}=0`$, these curves have critical points the singular points of $`b_{2k+1}c=0`$.
Bringing together Petrovskii’s theory and the implicit function theorem applied to one-parameter polynomial $`x_0^{2k+1}.(b_{2k+1}c)`$ with parameter $`c`$, it follows that a point $`p_{2k+1}`$ with $`b_{2k+1}(p_{2k+1})=c_0>0`$ is associated with $`a`$.
When $`c=c_0`$, the one-side component of the curve $`x_0^{2k+1}.(b_{2k+1}(c_0)`$ touches itself. Besides, according to Petrovskii’s Lemma 2 (3.a) and Lemma 3 as $`c`$ increases from $`c_0`$ to $`c_0+ϵ`$ one negative oval $`𝒪^{}`$ (of the curve $`x_0^{2k+1}.(b_{2k+1}(c_0ϵ)=0`$) disappears. (In the pencil of curves, $`𝒪^{}`$ as oval of $`_{2k}`$ has been created in the construction of $`_{2k+1}`$ from $`_{2k}L`$) (3) Let $`_{2k+2}`$ be the Harnack curve of degree $`2k+2`$ and $`B_{2k+2}`$ be a Harnack polynomial of degree $`2k+2`$.
Consider $`J`$ the square with vertices $`(c,d),(c,d+2),(c+2,d),(c+2,d+2)`$. It is contained in the triangle $`T_{2k+2}`$ and triangulated by the triangulation $`\tau `$ of $`T_{2k+2}`$. Let $`U(J^0)`$ be the subset $`\rho ^{2k+2}(𝐑_+J^0\times U_𝐂^2)`$ of $`𝐂P^2`$. Obviously, $`U(p)U(K^0)U(J^0)`$.
Denote $`b_{2k+2}^J`$ the truncation of $`b_{2k+2}`$ on the monomials $`x^cy^d`$, $`x^cy^{d+1}`$, $`x^cy^{d+2}`$, $`x^{c+1}y^d`$, $`x^{c+1}y^{d+1}`$, $`x^{c+1}y^{d+2}`$, $`x^{c+2}y^d`$, $`x^{c+2}y^{d+1}`$,$`x^{c+2}y^{d+2}`$. (Namely, $`b_{2k+1}^J=b_{2k+1}^K++a_{c+2,d+2}x^{c+2}y^{d+2}`$ with $`a_{c+2,d+2}>0`$, $`b_{2k+1}^J=a_{c,d}x^{c,d}+a_{c+1,d}x^{c+1}y^d++a_{c,d+1}x^{c,d+1}a_{c+1,d}x^{c+1}y^d++a_{c,d+2}x^cy^{d+2}+a_{c+1,d+2}x^{c+1}y^{d+2}+a_{c+2,d}x^{c+2}y^d+a_{c+2,d+1}x^{c+2}y^{d+1}`$ with $`a_{c,d+2}>0,a_{c+1,d+2}>0,a_{c+2,d}>0,a_{c+2,d+1}>0`$.
According to the patchworking theory, the truncation $`b_{2k+2}^J`$ of $`b_{2k+2}`$ is $`ϵ`$-sufficient for $`b_{2k+2}`$ in $`U(J^0)`$. Therefore, $`\stackrel{~}{h}:𝐂_{2k+2}U(p)\{(x,y)(𝐂^{})^2|b_{2k+2}^S=0\}U(p)`$ extends to the homeomorphism $`\stackrel{~}{h}:𝐂_{2k+2}U(J^0)\{(x,y)(𝐂^{})^2|b_{2k+2}^J=0\}U(J^0)`$ such that $`\stackrel{~}{h}(a)=(x_0,y_0)`$ $`x_0.y_0<0`$.
It is easy to see that, up to modify coefficients of the polynomials $`b_{2k+2}^J`$, (without changing the order and the topological structure of $`b_{2k+2}`$), the point $`(x_0,y_0)`$, $`x_0<0,y_0<0`$ is a critical point of $`b_{2k+2}^J`$ with $`b_{2k+2}(x_0,y_0)<0`$.
Hence, we shall proceed to the study of the Petrovskii’s pencil of curves over $`_{2k+2}`$.
Denote $`d_{2k+2}`$ the unique polynomial such that $`b_{2k+2}=b_{2k+2}^J+d_{2k+2}`$.
As previously, consider curves of the Petrovskii’s pencil over $`_{2k+2}`$ outside $`d_{2k+2}=0`$, as level curves of the function $`\frac{b_{2k+2}^Jc}{d_{2k+2}}`$.
Inside $`(b_{2k+2}^Jc=0)\backslash (d_{2k+2}=0)`$, these curves have critical points the singular points of $`b_{2k+2}^Jc=0`$. Bringing together Petrovskii’s theory, the implicit function theorem applied to one-parameter polynomial $`x_0^{2k+2}.(b_{2k+2}c)=x_0^{2k}.(b_{2k+2}^J+d_{2k+2}c)`$ with parameter $`c`$, it follows that a point $`p_{2k+2}`$ with $`b_{2k+2}(p_{2k+2})=c_0<0`$ is associated with $`a`$.
Let $`b_{2k+2}(p_{2k+2})=c_0`$. When $`c=c_0`$, one positive oval touches another positive oval. Then, as $`c`$ decreases from $`c_0`$ to $`c_0ϵ`$, a positive oval disappears. Let $`_{2k+j}`$ be the Harnack curve of degree $`2k+j`$, $`j3`$ and let $`B_{2k+j}`$ be a Harnack polynomial of degree $`m+j`$
As previously, consider $`b_{2k+j}^J`$ the truncation of $`b_{2k+j}`$ on the monomials $`x^cy^d`$, $`x^cy^{d+1}`$, $`x^{c+1}y^d`$, $`x^{c+1}y^{d+1}`$.
According to patchworking theory, the truncation $`b_{2k+j}^K`$ of $`b_{2k+j}`$ is $`ϵ`$-sufficient for $`b_{2k+j}`$ in $`U(J^0)`$.
Therefore,
$`\stackrel{~}{h}:𝐂_{2k+j}U(p)\{(x,y)(𝐂^{})^2|b_{2k+j}^S=0\}U(p)`$ extends to the homeomorphism
$`\stackrel{~}{h}:𝐂_{2k+j}U(K^0)\{(x,y)(𝐂^{})^2|b_{2k+1}^K=0\}U(K^0)`$ such that $`\stackrel{~}{h}(a)=(x_0,y_0)`$ $`x_0.y_0>0`$.
Let $`d_{2k+j}`$ be the unique polynomial such that $`b_{2k+j}=b_{2k+j}^J+d_{2k+j}`$. Consider curves of the Petrovskii’s pencil over $`_{2k+j}`$ outside $`d_{2k+j}=0`$, as level curves of the function $`\frac{b_{2k+j}^Jc}{d_{2k+j}}`$. From an argumentation similar to the previous one, it follows that a point $`p_{2k+j}`$ of $`S_{2k+j}^{}`$ equivalent to $`p_{2k+2}`$ is associated with $`a`$.
Therefore, a 3-point of curves $`_{2k+j}`$, $`j3`$, is associated with $`a`$.
Q.E.D
Recall (see propositiony 2.3.6 of Chapter 2.3.6) that given $`B_{2k}=x_0^{2k}.b_{2k}(x_1/x_0,x_2/x_0)`$ a Harnack polynomial of degree $`2k`$ and type $`^0`$, $`S\mathrm{"}_{2k}`$ denotes $`S_{2k}^+\backslash S_{2k}^{}`$ the complementary set of $`S_{2k}^{}`$ inside $`S_{2k}^+`$. In the Lemma 4.2.9 of Chapter 4.2.9, we shall study how critical points $`S\mathrm{"}_{2k}`$ of a Harnack polynomial of type $`^0`$ and degree $`2k`$ vary in the recursive construction of Harnack curves.
###### Lemma 4.2.9.
Let $`_{2k}`$ be the Harnack curve of degree $`2k`$ obtained via the patchworking method and given by a Harnack polynomial of type $`^0`$. Let $`\mathrm{\Gamma }`$ be a face of the triangulation of $`T_{2k}`$ contained into the line $`l_{2k1}`$ and given by vertices $`(c+1,d),(c,d+1)`$ with $`c`$ and $`d`$ even and strictly positive. Let $`U(p)`$ be the $`ϵ`$-neighborhood of $`𝐂_m`$ defined from $`\mathrm{\Gamma }^0`$. Denote $`a`$ the unique crossing of $`_{2k1}L`$ such that $`aU(p)`$. Then, a $`3`$-point $`(p_{2k},p_{2k+1},p_{2k+2})`$ of Harnack curves $`_{2k+j}`$, $`j3`$, is associated with $`a`$.
1. There exists a negative oval $`𝒪`$ such that the point
$`p_{2k}S\mathrm{"}_{2k}S_{2k}^+`$ has the property: $`𝒪_0W(p_{2k})`$ $`𝒪_1W(p_{2k})`$.
Moreover, given $`B(p_{2k})`$ a conj-equivariant $`4`$-ball around $`p_{2k}`$,
$$_{2k}B(p_{2k})(_0W(p_{2k})B(p_{2k}))_1$$
2. There exists a negative oval $`𝒪`$ such that the point
$`p_{2k+1}S_{2k+1}^+`$ has the property: $`𝒪_0W(p_{2k+1})`$ $`𝒪_1W(p_{2k+1})`$.
Moreover, given $`B(p_{2k+1})`$ a conj-equivariant $`4`$-ball around $`p_{2k+1}`$,
$$_{2k+1}B(p_{2k+1})(_0W(p_{2k+1})B(p_{2k+1}))_1$$
3. There exists a negative oval $`𝒪`$ such that the point
$`p_{2k+2}S_{2k+2}^+`$ has the property: $`𝒪_0W(p_{2k+2})`$ $`𝒪_1W(p_{2k+2})`$.
Moreover, given $`B(p_{2k+2})`$ a conj-equivariant $`4`$-ball around $`p_{2k+2}`$,
$$_{2k+2}B(p_{2k+2})(_0W(p_{2k+2})B(p_{2k+2}))_1$$
###### Definition 4.2.10.
We shall call $`3`$-point of the second kind a $`3`$-point verifying all assumptions of Lemma 4.2.9 of Chapter 4.2.9.
proof Our proof is based on arguments similar to the one given in the proof of Lemma 4.2.7 of Chapter 4.2.7.
We shall use local description of Harnack curves $`_{m+j}`$ $`j0`$ inside $`U(p)`$ provided by the patchworking theory.
Recall that for any $`j0`$, there exists an homeomorphism
$`\stackrel{~}{h}:𝐂_{m+j}U(p)\{(x,y)(𝐂^{})^2|b_{m+j}^S=0\}U(p)`$ such that $`\stackrel{~}{h}(a)=(x_0,y_0)`$ is a critical point of $`b_{m+j}^S`$. (1) Let $`_{2k}`$ be the Harnack curve of degree $`2k`$, $`(k>3)`$ and $`B_{2k}`$ be Harnack polynomial of degree $`2k`$.
From proposition 2.3.6 of Chapter 2.3.6, a point $`p_{2k}`$ of $`S_{2k}^+`$ is associated with $`a`$. Set $`b_{2k}(p_{2k})=c_0`$. Then, when $`c=c_0`$, the non-empty positive oval touches a negative (inner) oval; as $`c`$ increases from $`c_0`$ to $`c_0+ϵ`$, one negative oval $`𝒪^{}`$ disappears. (In the pencil of curves, $`𝒪^{}`$ as oval of $`_{2k1}`$ has been created in the construction of $`_{2k1}`$ from $`_{2k2}L`$).
(2)
Let $`_{2k+1}`$ be the Harnack curve of degree $`2k+1`$ and $`B_{2k+1}`$ be a Harnack polynomial of degree $`2k+1`$.
Consider the convex polygon $`KT_{2k+1}`$ with vertices $`(c,d),(c,d+2),(c+1,d+2),(c+2,d),(c+2,d+1)`$. It is contained in triangle $`T_{2k+1}`$ and triangulated by the triangulation $`\tau `$ of $`T_{2k+1}`$.
Let $`U(K^0)`$ be the subset $`\rho ^{2k+1}(𝐑_+K^0\times U_𝐂^2)`$ of $`𝐂P^2`$. Obviously, $`U(p)U(K^0)`$.
Denote $`b_{2k+1}^K`$ the truncation of $`b_{2k+1}`$ on the monomials $`x^cy^d`$, $`x^cy^{d+1}`$, $`x^cy^{d+2}`$, $`x^{c+1}y^d`$, $`x^{c+1}y^{d+1}`$, $`x^{c+1}y^{d+2}`$, $`x^{c+2}y^d`$, $`x^{c+2}y^{d+1}`$.
According to the patchworking theory, the truncation $`b_{2k+1}^K`$ of $`b_{2k+1}`$ is $`ϵ`$-sufficient for $`b_{2k+1}`$ in $`U(K^0)`$.
Therefore, $`\stackrel{~}{h}:𝐂_{2k+1}U(p)\{(x,y)(𝐂^{})^2|b_{2k+1}^S=0\}U(p)`$ extends to the homeomorphism $`\stackrel{~}{h}:𝐂_{2k+1}U(K^0)\{(x,y)(𝐂^{})^2|b_{2k+1}^K=0\}U(K^0)`$ such that $`\stackrel{~}{h}(a)=(x_0,y_0)`$ $`x_0.y_0<0`$. According to proposition 2.3.1 of Chapter 2.3.1, one can assume that $`(x_0,y_0)`$ is a critical point of $`b_{2k+1}^S=b_{2k}^S`$. It is easy to see that, up to modify coefficients of the polynomials $`b_{2k+1}^K`$ (without changing the order and the topological structure of $`b_{2k+1}`$), the point $`(x_0,y_0)`$ $`x_0>0,y_0>0`$ is a critical point of $`b_{2k+1}^K`$ with $`b_{2k+1}(x_0,y_0)>0`$.
Let $`d_{2k+1}`$ be the unique polynomial such that $`b_{2k+1}=b_{2k+1}^K+d_{2k+1}`$. Consider curves of the Petrovskii’s pencil over $`_{2k+1}`$ outside $`d_{2k+1}=0`$, as level curves of the function $`\frac{b_{2k+1}^Kc}{d_{2k+1}}`$.
Inside $`(b_{2k+1}^Kc=0)\backslash (d_{2k+1}=0)`$, these curves have critical points the singular points of $`b_{2k+1}c=0`$.
Bringing together Petrovskii’s theory and the implicit function theorem applied to one-parameter polynomial $`x_0^{2k+1}.(b_{2k+1}c)`$ with parameter $`c`$, it follows that a point $`p_{2k+1}`$ with $`b_{2k+1}(p_{2k+1})=c_0>0`$ is associated with $`a`$.
When $`c=c_0`$, the one-side component of the curve touches itself. Then, as $`c`$ increases from $`c_0`$ to $`c_0+ϵ`$ one negative oval $`𝒪^{}`$ disappears. (In the pencil of curves, $`𝒪^{}`$ as oval of $`_{2k1}`$ has been created in the construction of $`_{2k1}`$ from $`_{2k2}L`$).
(3) Let $`_{2k+2}`$ be the Harnack curve of degree $`2k+2`$ and $`B_{2k+2}`$ be the Harnack polynomial of degree $`2k+2`$
Consider $`J`$ the square with vertices $`(c,d),(c,d+2),(c+2,d),(c+2,d+2)`$. It is contained in the triangle $`T_{2k+2}`$ and triangulated by the triangulation $`\tau `$ of $`T_{2k+2}`$.
Let $`U(J^0)`$ be the subset $`\rho ^{2k+2}(𝐑_+J^0\times U_𝐂^2)`$ of $`𝐂P^2`$. Obviously, $`U(p)U(K^0)U(J^0)`$.
Denote $`b_{2k+2}^J`$ the truncation of $`b_{2k+2}`$ on the monomials $`x^cy^d`$, $`x^cy^{d+1}`$, $`x^cy^{d+2}`$, $`x^{c+1}y^d`$, $`x^{c+1}y^{d+1}`$, $`x^{c+1}y^{d+2}`$, $`x^{c+2}y^d`$, $`x^{c+2}y^{d+1}`$,$`x^{c+2}y^{d+2}`$.
According to patchworking theory, the truncation $`b_{2k+2}^J`$ of $`b_{2k+2}`$ is $`ϵ`$-sufficient for $`b_{2k+2}`$ in $`U(J^0)`$.
Therefore, $`\stackrel{~}{h}:𝐂_{2k+2}U(p)\{(x,y)(𝐂^{})^2|b_{2k+2}^S=0\}U(p)`$ extends to the homeomorphism $`\stackrel{~}{h}:𝐂_{2k+2}U(J^0)\{(x,y)(𝐂^{})^2|b_{2k+2}^J=0\}U(J^0)`$ such that $`\stackrel{~}{h}(a)=(x_0,y_0)`$ $`x_0.y_0<0`$.
It is easy to see that up to modify coefficients of the polynomials $`b_{2k+2}^J`$ (without changing the order and the topological structure of $`b_{2k+2}`$), the point $`(x_0,y_0)`$ $`x_0>0,y_0<0`$ is a critical point of $`b_{2k+2}^J`$ with $`b_{2k+2}^J(x_0,y_0)>0`$.
According to the Petrovskii’s theory, from an argumentation similar to the previous one, it follows that a point $`p_{2k+2}`$ with $`b_{2k+2}(p_{2k+2})=c_0>0`$ is associated with $`a`$.
When $`c=c_0`$, the negative oval $`𝒪^{}`$ touches another negative oval. (In the pencil of curves, this last negative oval as oval of $`_{2k+1}`$ is, in the patchworking scheme, situated in front of $`𝒪^{}`$ and has been created in the construction of $`_{2k+1}`$ from $`_{2k}L`$).
Then, as $`c`$ increases from $`c_0`$ to $`c_0+ϵ`$, one negative oval disappears.
Let $`j3`$ be an integer and $`_{2k+j}`$ be the Harnack curve of degree $`2k+j`$.
Let $`B_{2k+j}`$ be a Harnack polynomial of degree $`m+j`$. As previously, consider $`b_{2k+j}^J`$ the truncation of $`b_{2k+j}`$ on the monomials $`x^cy^d`$, $`x^cy^{d+1}`$, $`x^{c+1}y^d`$, $`x^{c+1}y^{d+1}`$. According to the patchworking theory, the truncation $`b_{2k+j}^K`$ of $`b_{2k+j}`$ is $`ϵ`$-sufficient for $`b_{2k+j}`$ in $`U(J^0)`$. Therefore, $`\stackrel{~}{h}:𝐂_{2k+j}U(p)\{(x,y)(𝐂^{})^2|b_{2k+j}^S=0\}U(p)`$. extends to the homeomorphism $`\stackrel{~}{h}:𝐂_{2k+j}U(K^0)\{(x,y)(𝐂^{})^2|b_{2k+1}^K=0\}U(K^0)`$ such that $`\stackrel{~}{h}(a)=(x_0,y_0)`$ $`x_0.y_0<0`$.
From an argumentation similar to the previous one, it follows that a point $`p_{2k+j}`$ of $`S_{2k+j}^+`$ equivalent to $`p_{2k+2}`$ is associated with $`a`$.
Therefore, a 3-point is associated with $`a`$. Q.E.D
###### Remark 4.2.11.
Given $`_{2k}`$ the Harnack of degree $`2k`$. It follows immediately from the Lemmas 4.2.6 of Chapter 4.2.6, 4.2.7 of Chapter 4.2.7, 4.2.9 of Chapter 4.2.9 above that $`(k2)^2`$ $`3`$-points are associated to crossing of curves $`_{2k3j}L`$ $`(L𝐂l_{2k3j})`$ $`(0j(2k7))`$. Likewise, given $`_{2k+1}`$ the Harnack curve of odd degree $`2k+1`$, $`(k2)^2`$ $`3`$-points are associated to crossing of curves $`_{2k3j}L`$ $`(L𝐂l_{2k3j})`$ $`(0j(2k7))`$, and $`(2k3)`$ $`2`$-points are associated with crossing points of $`_{2k1}L`$. Hence, for any Harnack curve $`_m`$ of degree $`m`$ with polynomial $`B_m`$, one can recover from the Lemmas above local topological meaning of critical points of index 1 of $`B_m`$.
### Chapter 5 Perestroika theory on Harnack curves
-Construction of any curve of degree $`m`$ with non-empty real part-
In this chapter, we shall give a method of construction of any curve of a given degree $`m`$ which provides description of pairs $`(𝐂P^2,𝐂𝒜_m)`$ up to conj-equivariant isotopy. In a few words, our method of construction of curves of degree $`m`$ is based on the definition of a chain of modification on the real connected components of the Harnack curve $`_m`$.
We shall divide this section into three sections. Our method of construction of algebraic curve uses invariants of real algebraic curves similar to the Arnold’s invariants of generic immersion of the circle into the plane. These invariants and their analogous for algebraic curves were respectively introduced in and . In the first section, we state the problem of construction of algebraic curves from the pair $`(𝐂P^2,𝐂_m)`$ and recall definitions of these invariants. Classification of pairs $`(𝐂P^2,𝐂𝒜_m)`$ leads us to distinguish curves with orientable set of real points usually called curves of type $`I`$, and curves with non-orientable set of real points usually called curves of type $`II`$. In the second section, we state a method of construction of curves of type $`I`$. Then, in the third section, we enlarge the method to curves of type $`II`$. The main result of this section is gathered in Theorem 5.2.9 of Chapter 5.2.9 and Theorem 5.3.7 of Chapter 5.3.7. This statement (Theorem 5.2.9 of Chapter 5.2.9 and 5.3.7 of Chapter 5.3.7) is the counterpart of Theorem 2.3.9 of Chapter 2.3.9 established for Harnack curves. Namely, we present any curve $`𝒜_m`$ of degree $`m`$ with non-empty real part up to conj-equivariant isotopy of $`𝐂P^2`$ as follows: Outside a finite number of $`4`$-balls $`B(a_i)`$ globally invariant by complex conjugation, $`𝒜_m`$ is the union of $`m`$ non-intersecting projective lines; inside any $`4`$-ball $`B(a_i)`$ it is the perturbation of a crossing.
#### 5.1. Introduction-Classical Topological Facts-
In any cases, (curves $`𝒜_m`$ with an arbitrary number of real components and arbitrary type), we shall bring the problem of description of $`(𝐂P^2,𝐂𝒜_m)`$ to the problem of description of all possible moves of the real components of $`_m`$. Let us recall in introduction general facts (see for example) which will state the space of moves providing an arbitrary curve of degree $`m`$ from the Harnack curve $`_m`$.
##### Introduction
Passing from polynomials to real set of points, it follows easily that real algebraic curves form a real projective space of dimension $`\frac{m(m+3)}{2}`$. We shall denote this space by the symbol $`𝐑𝒞_m`$. It easy to see that one can trace the real part of a curve of degree $`m`$ through any $`\frac{m(m+3)}{2}`$ real points, and that it is uniquely defined for points in general position. Moreover, real algebraic curves can be considered as complex curves of special kind. Similarly, passing from polynomials to complex set of points, it follows easily that complex algebraic curves of degree $`m`$ form a complex projective space of dimension $`\frac{m(m+3)}{2}`$. We shall denote this space by $`𝐂𝒞_m`$. Obviously, $`𝐑𝒞_m`$ coincides with the real part of $`𝐂𝒞_m`$. It is easy to see that through any $`\frac{m(m+3)}{2}`$ points one can always trace a curve of degree $`m`$. Moreover, the set of $`\frac{m(m+3)}{2}`$-points for which such a curve is unique is open and dense in the space of all $`\frac{m(m+3)}{2}`$-points subset of $`𝐂P^2`$.
Let $`𝐑𝒟_m`$ denote the subset of $`𝐑𝒞_m`$ corresponding to real singular curves. Let $`𝐂𝒟_m`$ denote the subset of $`𝐂𝒞_m`$ corresponding to singular curves.
We call a path in the complement $`𝐑𝒞_m\backslash 𝐑𝒟_m`$ of the discriminant hypersurface in $`𝐑𝒞_m`$ a rigid isotopy of real points set of nonsingular curves of degree $`m`$. We call a smooth path in the complement $`𝐂𝒞_m\backslash 𝐂𝒟_m`$ of the discriminant hypersurface in $`𝐂𝒞_m`$ a rigid isotopy of complex point set of nonsingular curves of degree $`m`$. These definitions give rise naturally to the classification problem of non-singular curves of degree $`m`$ up to rigid isotopy. From the complex viewpoint, the rigid isotopy classification problem has trivial solution: the complex point sets of any two nonsingular curves of degree $`m`$ are rigidly isotopic. Moreover, they are diffeotopic in $`𝐂P^2`$. From the real viewpoint, even if we consider only rigid isotopy, this property can not be extended.
##### State of the problem
From now on, one can state some properties of the sets of points $`𝐂𝒜_m`$, $`𝐑𝒜_m`$ of a curve $`𝒜_m`$. Let $`_m`$ be the Harnack curve of degree $`m`$ and $`𝒜_m`$ be a curve of degree $`m`$ such that the pair $`(𝐑P^2,𝐑𝒜_m)`$ is non-homeomorphic to the pair $`(𝐑P^2,𝐑_m)`$. From the elementary topological facts recalled previously, there exists a smooth path $`h:[0,1]𝐂𝒞_m\backslash 𝐂𝒟_m`$ with $`h(0)=𝐂_m,h(1)=𝐂𝒜_m`$ such that the family $`𝐂h(t)`$ is a diffeotopy of submanifolds of $`𝐂P^2`$. Thus, if there exists $`B𝐂P^2`$ globally invariant by complex conjugation such that $`h_t(𝐂_m)B)B`$, $`t[0,1]`$, topological spaces $`(B_m)`$ and $`(B𝒜_m)`$ are isotopic in $`B`$.
According to corollary 2.3.5 of Chapter 2.3.5, assume $`_m`$ obtained via $`T`$-inductive construction of Harnack curves. From the Lemma 4.2.6 of Chapter 4.2.6, Lemma 4.2.7 of Chapter 4.2.7 and Lemma 4.2.9 of Chapter 4.2.9 of the previous section, any critical point $`(x_0,y_0)`$ of index 1 of $`B_m`$ is associated with a crossing $`a`$ of one curve $`_{mi}L`$. Namely, there exists $`U(p)=\rho ^m(D(p,ϵ)\times U_𝐂^2)(𝐂^{})^2`$ $`ϵ`$-neighborhood of $`𝐂_m`$ defined the interior $`\mathrm{\Gamma }^0`$ of a face $`\mathrm{\Gamma }`$ of the triangulation of $`T_m`$ such that :
1. $`a,(x_0,y_0)U(p)`$
2. the perturbation on the real part of $`_m`$ involved in the local topological meaning of $`(x_0,y_0)`$ is a deformation on $`𝐂_mU(p)`$
Let us denote by $`𝒫_m`$ the set of points $`p`$ of $`T_m`$ with the property that for any $`p𝒫_m`$, the $`ϵ`$-neighborhood $`U(p)`$ verifies the two properties above. Note that for any two distinct points $`p,p^{}`$ the subsets $`U(p)`$ and $`U(p^{})`$ have empty-intersection. It easy to see that the set $`𝒫_m`$ contains $`\frac{m(m1)}{2}`$ points of $`T_m`$.
Moreover, for any $`p𝒫_m`$, $`U(p)=\rho ^m(D(p,ϵ)\times U_𝐂^2)`$ contains conj-equivariant $`4`$-balls $`B(x_0,y_0)`$ and $`B(a)`$ around $`(x_0,y_0)`$ and $`a`$ such that
$`𝐂_mB((x_0,y_0))\mathrm{}`$ ,$`𝐂_mB(a)\mathrm{}`$. Thus, according to the density of $`\frac{m(m+3)}{2}`$-uple of points which define uniquely a curve in the space of $`\frac{m(m+3)}{2}`$-uple of points of $`𝐂P^2`$, if there exists a diffeotopy of submanifolds of $`𝐂P^2`$ $`h:[0,1]𝐂𝒞_m\backslash 𝐂𝒟_m`$ with $`h(0)=𝐂_m,h(1)=𝐂𝒜_m`$ such that $`h_t(𝐂_mU(p))U(p)`$, $`t[0,1]`$, the curve $`𝒜_m`$ is entirely defined from its restriction on the union $`_{\{p𝒫_m\}}U(p)`$ of $`ϵ`$-neighborhood $`U(p)`$ taken over the $`\frac{m(m1)}{2}`$ points of $`𝒫_m`$.
In such a way, we bring the problem of construction of a curve $`𝒜_m`$ to the definition of a path $`h:[0,1]𝐂𝒞_m\backslash 𝐂𝒟_m`$ with $`h(0)=𝐂_m,h(1)=𝐂𝒜_m`$.
##### Varieties of irreducible curves of type $`I`$, degree $`m`$, and genus $`g`$
As already introduced for smooth curves, curves with orientable real set of points are called curves of type $`I`$.
In what follows, we shall consider smooth curves of type $`I`$ in the broader set of irreducible curves of type $`I`$. Any irreducible curve of type $`I`$ is such that the real part of its normalization divides the set of complex points of its normalization into two halves. We shall call the images of the halves of the normalization in the set of points of the curve, the halves of the curve. Each of the halves is oriented and induces an orientation on the real part as on its boundary. We shall consider irreducible curves of degree $`m`$, genus $`g`$ and type $`I`$ with a distinguished complex orientation. Curves of this kind constitute a finite-dimensional stratified real algebraic variety we shall denote $`𝒞_{I;g}`$.
We shall set $`𝒞_{I,m}=_{0g\frac{(m1)(m2)}{2}}𝒞_{I;g}`$ the subset of $`𝒞_m`$ constituted by all irreducible curves of degree $`m`$, genus $`g`$, $`0g\frac{(m1)(m2)}{2}`$, and type $`I`$.
Passing from polynomials to real (resp, complex) set of points, it follows easily that varieties $`𝒞_{I;g}`$ define subspaces $`𝐑𝒞_{I;g}`$ (resp,$`𝐂C_{I;g}`$) of the space $`𝐑𝒞_m`$ (resp, $`𝐂𝒞_m`$.)
We shall set $`𝐑𝒞_{I,m}=_{0g\frac{(m1)(m2)}{2}}𝐑𝒞_{I;g}`$ the subset of $`𝐑𝒞_m`$ constituted by real set of points of irreducible curves of degree $`m`$ and type $`I`$, $`𝐂𝒞_{I,m}=_{0g\frac{(m1)(m2)}{2}}𝐂𝒞_{I;g}`$ the subset of $`𝐂𝒞_m`$ constituted by complex set of points of irreducible curves of degree $`m`$ and type $`I`$. Obviously, $`𝐑𝒞_{I,m}`$ coincides with the real part of $`𝐂𝒞_{I,m}`$. Let us study the set of singular algebraic curves in the variety $`𝒞_{I;g}`$ of curves of degree $`m`$, genus $`g`$ and type $`I`$. Generic curves Following usual terminology, we call ordinary double point singularity of a real irreducible algebraic curve $`𝒜_m`$ a non-degenerate singular point of $`𝒜_m`$. All ordinary double points singularity are equivalent from the complex viewpoint. From the real viewpoint, one distinguishes several types of such points.
1. real double point of intersection of two real branches called crossing
2. real point of intersection of two complex branches conjugated to each other called solitary double point
3. imaginary double point of intersection of the different halves of $`𝐂𝒜_m\backslash 𝐑𝒜_m`$ denoted $`\alpha `$-point.
4. imaginary double point of self-intersection of one of the halves of $`𝐂𝒜_m\backslash 𝐑𝒜_m`$ denoted $`\beta `$-point.
Define a generic curve as a real irreducible algebraic curve with only ordinary double singularities. As it is well known, generic curves in the variety $`𝒞_{I;g}`$ constitute a Zarisky open set in the variety $`𝒞_{I;g}`$.
For a generic curve $`𝒜_m`$ of degree $`m`$ and genus $`g`$ ,$`0<g<\frac{(m1)(m2)}{2}`$, by smoothing of its real point set $`𝐑𝒜_m`$ we shall understand a smooth oriented $`1`$-dimensional submanifold of $`𝐑P^2`$ obtained from $`𝐑𝒜_m`$ by the modification at each double point determined by the complex orientation (see figure 5.1).
By smoothing of its complex point set we shall understand a smooth oriented $`1`$-dimensional complex submanifold of $`𝐂P^2`$ obtained from $`𝐂𝒜_m`$ by the modification at each real double point as above and by the modification at each complex double point. The two complex branches which merge in $`\alpha `$-point (resp $`\beta `$-point) become after smoothing two branches of different halves of $`𝐂𝒜_m\backslash 𝐑𝒜_m`$ (resp, two branches of one of the halves of $`𝐂𝒜_m\backslash 𝐑𝒜_m`$) We shall call smoothing of a generic curve $`𝒜_m`$, the smooth $`1`$-dimensional complex submanifold of $`𝐂P^2`$ deduced from $`𝒜_m`$ by smoothing of $`𝐑𝒜_m`$ and $`𝐂𝒜_m`$ as described above. We shall say that the singular points of $`𝒜_m`$ are smoothed. Discriminant Hypersurface Consider now the complement of the set of generic curves of all real algebraic curves in the variety $`𝒞_{I;g}`$. It can be considered as a discriminant hypersurface of which strata consist of curves with only one singular point which is not an ordinary double point. One can distinguish six main strata defined by the type of the singularity of the curves it contains:
1. real cusp,
2. real point of direct ordinary tangency,
3. real point of inverse ordinary tangency,
4. real point of ordinary tangency of two imaginary branches,
5. real ordinary triple point of intersection of three real branches,
6. real ordinary triple point of intersection of a real branch and two conjugate imaginary branches.
Perestroika A generic path in the variety $`𝒞_{I;g}`$ intersects the discriminant hypersurface in a finite number of points, and these points belong to the main strata. We call perestroika a change experienced by a generic path in the space $`𝒞_{I;g}`$ when it goes through the main strata. The notion of perestroika was initially introduced in the context of generic immersion of the circle into the plane (see for example , ).
Given a perestroika, we call smoothed perestroika the change obtained by smoothing the fragments involved in the perestroika. Moreover, to define how a generic path crosses the strata of the discriminant hypersurface one has to specify a coorientation of the strata. Let us define the positive direction of a perestroika. The opposite direction is naturally called negative direction.
1. In the case of cusp, define the positive direction the direction from curve with one more crossing point to curve with one more solitary double point.
2. In case of real point of tangency, there is a natural positive direction from curve with less real double points to curve with more double points.
3. In case of triple-point, the coorientation of the stratum can be defined as follows. The crossing of a triple point by a path gives rise to the vanishing triangle: the dying triangle which existed just before the crossing and the new born existed just after it. A cyclic order on the sides of the triangle is given by the order the sides are visited. In the case of triple-point, the coorientation rule assigns signs to triangle. Let $`q`$ be the number of the sides of the vanishing triangle whose directions coincide with that given by the cyclical order. The sign of a triangle is $`(1)^q`$. The crossing of the triple point is positive (resp, negative) if the newborn triangle is positive (resp, negative) (and hence the dying one is negative (resp, positive)).
Let us study perestroikas and smoothed perestroikas in the positive direction. Properties of perestroikas in the negative direction can be naturally deduced from properties in the positive one. We shall first consider perestroikas in which no imaginary double point is involved; namely, cusp perestroika and triple-point perestroika. It is obvious that the corresponding smoothed perestroikas do not change the complex part. Therefore, we shall consider only real part.
1. Cusp Perestroika
It is easy to see that smoothed cusp perestroika does not change real part.
2. Triple-point perestroika. Relatives orientation of the three real branches of a real triple-point gives rise to distinguish two kinds of real triple-point. Consider vectors at the triple point tangent to the branches and directed according to their orientations. If one of the vectors can be presented as linear combination of the two others, then the triple-point is said weak (see figure 5.2); otherwise it is said strong( see figure 5.3) Smoothed weak triple-point perestroika does not change real part.
Smoothed strong triple-point perestroika changes real part as shown in figure 5.(3.b).
Then, consider perestroika in which imaginary singularities are involved.
Introduce a function $`J`$ counting the difference between the number of (resp, smoothed) $`\alpha `$-points and $`\beta `$-points under (resp, smoothed) perestroika. Obviously, if $`J(\alpha )`$ and $`J(\beta )`$ are the values of $`J`$ under one perestroika $`\pi `$ in the positive direction the function $`J`$ takes values $`J(\alpha )`$, $`J(\beta )`$ under the perestroika $`\pi `$ in the negative direction.
1. As already known, relative orientation of the real branches of a real self-tangency point gives rise to distinguish real point of direct ordinary tangency and real point of inverse ordinary tangency.
1. It is easy to see that smoothed real direct ordinary tangency perestroika does not change smoothed real part. Nonetheless, the complex part is changed in such a way that:
$`J(\alpha )=0,J(\beta )=2`$ under real direct ordinary tangency perestroika.
2. Smoothed real inverse tangency perestroika changes the real part (see figure 5.4 ). Besides, the complex part is changed in such a way that : $`J(\alpha )=2,J(\beta )=0`$ under real inverse tangency perestroika.
2. Perestroika of solitary self tangency is as follows: two solitary double points come from the world of imaginary to form a solitary self-tangency point; then arise two solitary double points with opposite orientation. Besides, the complex part is changed in such a way that: $`J(\alpha )=2,J(\beta )=0`$ under a solitary inverse tangency perestroika.
3. Perestroika of a triple-point with imaginary branches changes the real part as shown in figure 5.5. Besides, it changes the complex part in such a way that $`J(\alpha )=+2,J(\beta )=2`$.
#### 5.2. Construction of Curves of type $`I`$
In this section, we shall consider curves $`𝒜_m`$ of degree $`m`$ and type $`I`$ with pair $`(𝐑P^2,𝐑𝒜_m)`$ non-homeomorphic to the pair $`(𝐑P^2,𝐑_m)`$. As already noticed, for any curve $`𝒜_m`$ of type $`I`$, the pair $`(𝐑P^2,𝐑𝒜_m)`$ with orientation of the real point set $`𝐑𝒜_m`$ provides the pair $`(𝐂P^2,𝐂𝒜_m)`$, up to conj-equivariant isotopy of $`𝐂P^2`$. Bringing together the recursive Morse-Petrovskii’s theory for Harnack’s curves $`_m`$ of chapter 4 and properties of generic paths in the varieties $`𝒞_{I;g}`$, we define in Proposition 5.2.4 of Chapter 5.2.4 a path $`S:[0,1]𝐑𝒞_m`$ with $`S(0)=𝐑_m`$, $`S(1)=𝐑𝒜_m`$. Such path is obtained from lifting a generic path in the space of generic immersion of the circle into $`𝐑^2`$ and is described as a set of moves defined on the connected real components of $`𝐑_m`$.
In this way, the classification of $`M`$-curves of a prescribed degree $`m`$ amounts to the description of all the possible locations of the real components of the Harnack curve $`_m`$.
The first two parts of this section are devoted to the definition of the lifting. Along a generic path in the space of generic immersion of the circle into $`𝐑^2`$, three type of events ”perestroikas” (namely, the instateneous triple crossings and the inverse and direct self-tangency) may happen. By means of perestroikas and their counterparts for algebraic curves, we bring the problem of the definition of the path $`S`$ to the one of lifting perestroikas in the space $`𝐑𝒞_m`$.
In the third part, we give in Theorem 5.2.9 of Chapter 5.2.9 a description of pairs $`(𝐂P^2,𝐂𝒜_m)`$ up to conj-equivariant isotopy which extends the properties of the Harnack curve $`_m`$ stated in Theorem 2.3.9 of Chapter 2.3.9 to any smooth curve of type $`I`$.
##### Smoothing generic immersion of the circle into $`𝐑^2`$
Recall that by a generic immersion of the circle $`S^1`$ into the plane $`𝐑^2`$, one means an immersion with only ordinary double points of transversal intersection, namely without triple points and without points of self-tangency. The space of all immersions is an infinite-dimensional manifold which consists of an infinite countable set of irreducible components described by H.Whitney in 1937.
Whitney’s Theorem (1937) The space of the immersions of a circle into the plane with the same Whitney index is pathwise connected The Whitney index of an immersion of an oriented curve into the plane is the rotation number of the tangent vector (the degree of the Gauss map).
Arbitrary differentiable immersion of the circle into the plane does not admit complexification. Nonetheless, we shall generalize Whitney’s Theorem to the case of real algebraic curves of type $`I`$. Let us recall in introduction the following property of the genus. Genus Property If a collection of disjoint circles embedded into a closed orientable surface of genus $`g`$ does not divide the surface, the number of circles is at most $`g`$. In particular $`g+1`$ disjoint circles always divide the surface.
The set of complex points of any smooth curve of degree $`m`$ is an orientable surface of genus $`g_m=\frac{(m1)(m2)}{2}`$. Therefore, one can easily conjecture the following one-to-one correspondence between the $`g_m+1`$ (resp, $`lg_m+1`$) real connected components of an $`M`$-curve (resp, an $`Mi`$-curve) $`𝒜_m`$ and the $`g_m`$ $`1`$-handles and the sphere $`S^2`$: the sphere $`S^2`$ and each $`1`$-handle contains one real connected component of $`𝐑𝒜_m`$ which divides it into two halves (resp, the sphere $`S^2`$ and each $`1`$-handle contains a part of a real connected component of $`𝐑𝒜_m`$; such part divides it into two halves.) Call regular curve, the smooth oriented submanifold of $`𝐑P^2`$ deduced from the generic immersion of the circle by modification at each real double point which is either the Morse modification in $`𝐑^2`$ in the direction coherent to a complex orientation or the Morse modification at infinity in $`𝐑P^2`$ (in the direction coherent to a complex orientation) of the double point which associates to the double point of $`𝐑^2`$ two points of the line at infinity of $`𝐑P^2`$.
In Lemma 5.2.1 of Chapter 5.2.1, we give a generalization of Whitney’s theorem to the case of real algebraic curves of type $`I`$.
###### Lemma 5.2.1.
Let $`𝒜_m`$ be a smooth curve of degree $`m`$ and type $`I`$, then its real point set $`𝐑𝒜_m𝐑P^2`$ is a regular curve deduced from a generic immersion of the circle $`S^1`$ into the plane $`𝐑^2`$ with $`\frac{(m1)(m2)}{2}n\frac{(m1)(m2)}{2}+[\frac{m}{2}]`$, double points and Whitney index $`\frac{(m1)(m2)}{2}+1`$. The smoothing is such that $`\frac{(m1)(m2)}{2}`$ double points are smoothed in $`𝐑^2`$ by Morse modification in the direction coherent to a complex orientation. The $`n\frac{(m1)(m2)}{2}`$ others double points are smoothed by Morse modification at infinity in $`𝐑P^2`$.
###### Remark 5.2.2.
One can associate to any smooth curve of type $`I`$ a three-dimensional rooted tree.
Recall that a generic immersion of the circle into the plane is a tree-like curve if any of its double points subdivides it into two disjoints loops. On the assumption of Lemma 5.2.1 of Chapter 5.2.1, let $`𝒜_m`$ be a smooth curve of degree $`m`$ and type $`I`$, and $`\varphi `$ the corresponding generic immersion of the circle $`S^1`$ into the plane $`𝐑^2`$ with $`\frac{(m1)(m2)}{2}n\frac{(m1)(m2)}{2}+[\frac{m}{2}]`$. Smooth the double-points of $`\varphi `$ which result from points of $`𝐑𝒜_m`$ at infinity in such a way that it results a curve in $`𝐑^3`$ with the property that any of its double points subdivides it into two disjoints loops; it follows the rooted-tree associated to $`𝒜_m`$.
In case of Harnack curves $`_m`$, an immediate definition of the root of its rooted-tree is the following: assume $`_m`$ obtained via patchworking process;
for odd $`m`$, the root consists of the part of the odd component of the curve which intersects $`T_m`$;
for even $`m`$, the root consists of the part of the positive oval (non-empty for $`m4`$) which intersects $`T_m`$.
proof:
Our proof makes use of properties of the complex point set $`𝐂𝒜_m`$ embedded in $`𝐂P^2`$.
Consider the usual handlebody decomposition of $`𝐂P^2=B_0B_1B`$ where $`B_0`$,$`B_1`$,$`B`$ are respectively 0, 2 and 4 handles.
The balls $`B_0`$ and $`B_1`$ meet along an unknotted solid torus $`S^1\times B^2`$. The gluing diffeomorphism $`S^1\times B^2S^1\times B^2`$ is given by the $`+1`$ framing map. In such a way, the canonical $`𝐑P^2`$ can be seen as the union of a Möbius band $``$ and the disc $`D^2B`$ glued along their boundary. The Möbius band $``$ lies in $`B_0B_1S^1\times D^2`$ with $``$ as the $`(2,1)`$ torus knot, and $`D^2B`$ as the properly imbedded unknotted disc. The complex conjugation switches $`B_0`$ and $`B_1`$ and lets fix $``$, it rotates $`B`$ around $`D^2`$.
The set of complex points of $`𝒜_m`$ is an orientable surface of genus $`g_m=\frac{(m1)(m2)}{2}`$, i.e it is diffeomorphic to a sphere $`S^2`$ with $`g_m`$ $`1`$-handles $`S_{g_m}^2`$. We shall denote $`h`$ the diffeomorphism $`h:𝐂𝒜_mS_{g_m}^2𝐑^3`$.
Recall that surfaces $`S_{g_m}^2`$ are constructed as follows. From the sphere $`S^2`$, $`(m1)`$ pairwise non-intersecting open discs are removed, and then the resulting holes are closed up by $`g_m`$ orientable cylinders $`S^1\times D^1`$ connecting the boundary circles of the discs removed.
Assume $`S^2`$ provided with a complex conjugation with fixed point set a circle $`S^1`$ which divides the sphere $`S^2`$ into two halves. Without loss of generality, one can assume that any disc removed from $`S^2`$ intersects $`S^1`$ and the two halves of $`S^2`$. In such a way, $`𝐑𝒜_m𝐑P^2`$ intersects each $`1`$-handle.
Fix $`D^2`$ the two disc of $`𝐑P^2=D^2`$ in such a way that the boundary circle of $`D^2`$ is $`S^1`$ and therefore each one handle belongs to the solid torus $`S^1\times D^2`$.
Let $`D_ϵ^2D^2`$ be the disc $`D^2`$ thickened. Since the interior $`(D_ϵ^2)^0`$ of $`D_ϵ^2`$ is homeomorphic to $`𝐑^2`$, one can project $`𝐑𝒜_m`$ (up to homeomorphism $`𝐑^2(D_ϵ^2)^0`$) in a direction perpendicular to $`𝐑^2`$ onto $`𝐑^2`$. We may suppose that the direction of the projection is generic i.e all points of self-intersection of the image on $`𝐑^2`$ are double and the angles of intersection are non-zero.
Let $`r:𝐑^3𝐑^2`$ be the projection which maps $`h(𝐑𝒜_m)𝐑^3`$ to $`𝐑^2`$.
Consider an oriented tubular fibration $`N𝐂𝒜_m`$. Since $`𝐂𝒜_m`$ is diffeomorphic to a sphere $`S^2`$ with $`g_m`$ $`1`$-handles, one can consider the restriction of $`𝐂𝒜_m`$ diffeomorphic to each torus $`T^2`$ given by the sphere $`S^2`$ with one of the $`g_m`$ $`1`$-handles. The oriented tubular neighborhood of $`h(𝐂𝒜_m)T^2`$, as oriented tubular neighborhood of the torus $`T^2`$, intersects the solid torus $`S^1\times D^2`$ with $`MS^1\times S^1`$ as the $`(2,1)`$ torus knot. Hence, since in $`𝐂P^2`$ each real line is split by its real part into two halves lines conjugate to each others, and $`2`$ disjoint circles always divide the torus; the real part of the restriction of $`h(𝐂𝒜_m)T^2`$ belongs to the boundary of the Möbius band in such a way that its projection to $`𝐑^2`$ gives one crossing. Besides these doubles, some double-points of $`r(h(𝐑𝒜_m))`$ may result either from two points of $`𝐑𝒜_m`$ which belong to two different handles or from two points which belong respectively to a $`1`$-handle and to $`S^2`$. It is easy to deduce from the relative location of the sphere $`S^2`$ and the $`1`$-handles that the projection of $`𝐑𝒜_m`$ leads to an even number of these double-points of $`r(h(𝐑𝒜_m))`$. Therefore, the projection $`r(h(𝐑𝒜_m))`$ contains $`g_m(mod2)`$ crossings and $`𝐑𝒜_m`$ can be seen as a regular curve.
By use of the Whitney expression of the index $`ind`$ of a plane curve with $`n`$ double points - $`ind=ϵ_i\pm 1`$ where the summation is over the set of $`n`$ double points, $`ϵ_i`$ is a sign associated to each double point, and the term $`\pm `$ depends on the orientations of the curve and of the plane - we shall deduce the Whitney index of $`r(h(𝐑𝒜_m))`$.
According to the definition of $`𝐑P^2`$, each real part of $`h(𝐑𝒜_m)`$ contained in a $`1`$-handle has the same projection to $`𝐑^2`$ as the $`(2,1)`$ torus knot on the torus defined from the sphere $`S^2`$ with the $`1`$-handle.
Therefore, according to the gluing diffeomorphism
$`S^1\times B^2S^1\times B^2`$, given an orientation on the circle $`S^1`$ induced by one of the halves of $`S^2`$, the projection of each real part of $`h(𝐑𝒜_m)`$ contained in a $`1`$-handle gives a crossing with sign $`+1`$ of the plane $`𝐑^2`$. Besides, any $`2`$ crossings which result from the projection of points of two different handles or of points of a $`1`$-handle and the sphere $`S^2`$ have opposite Whitney index. Hence, $`r(h(𝐑𝒜_m))`$ has Whitney index $`g_m+1`$.
Moreover, since $`𝒜_m`$ is of type $`I`$, the real part $`𝐑𝒜_m`$ is deduced from $`r(h(𝐑𝒜_m))`$ by smoothings which are either standard Morse modifications in $`𝐑^2`$ at the double point determined by the complex orientation, or smoothings at infinity in $`𝐑P^2\backslash 𝐑^2`$ which associate to a crossing its two pre-image under $`r^1`$.
Two real branches of $`r(h(𝐑𝒜_m))`$ which intersect in a crossing become after smoothing two oriented non-intersecting real branches which belong either to two different connected components of $`𝐑𝒜_m`$ or to one branch which lies as a part of the $`(2,1)`$ torus knot on the boundary Möbius $`𝐑P^2`$. According to the Rokhlin’s orientation formula, for any $`m>2`$, $`𝐑𝒜_m`$ has at least two connected components.
Therefore, for $`m>2`$, $`g_m`$ crossings are smoothed in $`𝐑^2`$ by Morse Modification in the direction coherent to the complex orientation and become two oriented non-intersecting real branches of distinct connected components of $`𝐑𝒜_m`$ in $`𝐑P^2`$. The other double-points are smoothed in such a way that it results two points of $`𝐑𝒜_m`$ in $`𝐑P^2\backslash 𝐑^2`$.
According to the Bezout’s theorem, since each of these $`ng_m`$ double-points results from two points which belong to $`𝐑P^2\backslash 𝐑^2`$ and the number of points in the intersection of $`𝐑𝒜_m`$ with the real line $`𝐑P^2\backslash 𝐑^2`$ is less or equal to $`m`$ and congruent to $`m(mod2)`$, the difference $`ng_m`$ is less or equal to $`[\frac{m}{2}]`$. It is obvious that one can assume, without loss of generality, that these $`[\frac{m}{2}]`$ double points results from points of $`𝐑𝒜_m`$ which belong to two different handles or from two points which belong respectively to a $`1`$-handle and to $`S^2`$.
Q.E.D. Given a generic immersion of the circle $`S^1`$ into the plane $`𝐑^2`$ and $`𝒮`$ the set of its singular points. We shall call partially regular curve the smooth oriented submanifold of $`𝐑P^2`$ deduced from a generic immersion of the circle by modification of a set $`𝒦𝒮`$, $`𝒦𝒮`$, of real double points where modification at each real double point of $`𝒦`$ is either the Morse modification in $`𝐑^2`$ (in the direction coherent to a complex orientation) or the Morse modification at infinity of the double point which associates to the double point of $`𝐑^2`$ two points of the real line $`𝐑P^2\backslash 𝐑^2`$.
The next Lemma enlarges the preceding statement to generic curves of degree $`m`$, type $`I`$ and genus $`g`$, and more generally to singular curves of degree $`m`$ and type $`I`$ with non-degenerate singular points.
###### Lemma 5.2.3.
Let $`𝒜_m`$ be a singular curve of degree $`m`$ and type $`I`$ with non-degenerate singular points, then its real point set $`𝐑𝒜_m𝐑P^2`$ is a partially regular curve deduced from a generic immersion of the circle $`S^1`$ into the plane $`𝐑^2`$ with $`\frac{(m1)(m2)}{2}n\frac{(m1)(m2)}{2}+[\frac{m}{2}]`$, double points and Whitney index $`\frac{(m1)(m2)}{2}+1`$ if and only if its set of singular points consists of at most $`\frac{(m1)(m2)}{2}`$ crossings.
proof: It follows from an argument similar to the one of the proof of Lemma 5.2.1 of Chapter 5.2.1. Q.E.D
##### Lifting of a generic path in the space of generic immersion of the circle into the plane to $`𝐑𝒞_m`$
According to the Lemma 5.2.1 of Chapter 5.2.1, one can consider $`𝐑_m`$ and $`𝐑𝒜_m`$ as smoothed immersions of the circle into the plane with the same Whitney index. Let us denote $`𝐑\stackrel{~}{}_m`$, $`𝐑\stackrel{~}{𝒜}_m`$ the corresponding immersions. From the Whitney’s Theorem, there exists a path $`\stackrel{~}{h}`$ which connects $`𝐑\stackrel{~}{}_m`$ and $`𝐑\stackrel{~}{𝒜}_m`$. In this section, we shall in Proposition 5.2.4 of Chapter 5.2.4 and Theorem 5.2.7 of Chapter 5.2.7 define a path $`S`$ in $`𝐑𝒞_m`$ $`S(0)=𝐑_m`$,
$`S(1)=𝐑𝒜_m`$ from lifting a path $`\stackrel{~}{h}`$ which connects $`𝐑\stackrel{~}{}_m`$ and $`𝐑\stackrel{~}{𝒜}_m`$ in the space of immersions of the circle into $`𝐑^2`$.
By means of the Lemmas 5.2.1 of Chapter 5.2.1 and 5.2.3 of Chapter 5.2.3, we shall develop properties provided by the genus property and interpret the path $`\stackrel{~}{h}:𝐑\stackrel{~}{}_m𝐑\stackrel{~}{𝒜}_m`$ as a change of relative position of the $`1`$-handles of $`(𝐂_m)S_{g_m}^2`$, $`(𝐂_m)𝐑_m`$ which gives $`(𝐂𝒜_m)S_{g_m}^2`$, $`(𝐂𝒜_m)𝐑𝒜_m`$. The path $`S`$ appears as the track on $`𝐑P^2`$ of a diffeotopy $`h_t`$ of $`𝐂P^2`$, $`t[0,1]`$, $`h(0)=𝐂_m`$, $`h(1)=𝐂𝒜_m`$. ( 5.2 of Chapter 5.2).A Lifting Perestroikas in the space $`𝐑𝒞_m`$
Recall that given a perestroika along a path in the space $`𝒞_{I,g}`$ we call smoothed perestroika the change obtained by smoothing the fragments involved in the perestroika. We shall call diffeotopic perestroika of perestroika defined along a path $`h`$ in the space $`𝒞_{I,g})`$ of irreducible curves of type $`I`$ and genus $`g`$ and degree $`1nm`$, the change experienced by a curve in $`𝒞_{I,g}`$ with the following property.
For any point $`a`$ which participated in a perestroika (along the path $`h`$) let $`U(a)`$ be a small $`ϵ`$-neighborhood around $`a`$ of the complex point set of the initial curve embedded in $`𝐂P^2`$, (one can choose also $`U(a)`$ as follows: $`U(a)=\{z=<u,p>=(u_0.p_0:u_1.p_1:u_2.p_2)𝐂P^2u=(u_0,u_1,u_2)U_𝐂^3,p=(p_0:p_1:p_2)D(a,ϵ)\}`$ where $`D(a,ϵ)`$ denotes a small disc (in the Fubini-Study metric) of radius $`ϵ`$ around $`a`$ in $`𝐑P^2`$. Restrictions to $`U(a)`$ of the complex point set of the initial curve and of the complex point set which results from the diffeotopic perestroika are diffeotopic. (As will become clear later, in most cases, given a curve $`𝒜_m`$ the change experienced by the curve $`𝒜_m`$ after just one diffeotopic perestroika does not lead to a curve.)
Given a diffeotopic perestroika, we call smoothed diffeotopic perestroika the change obtained by smoothing the fragments involved in the diffeotopic perestroika (i.e fragments inside $`U(a)`$ for any $`a`$ which participated in the perestroika).
( 5.2 of Chapter 5.2).B Chain of Diffeotopic Perestroika
Let us recall that up to regular deformation of its polynomial (see Theorem 2.2.30 of Chapter 2.2.30), curve $`_m`$ results from the recursive construction of Harnack curves $`_i`$, $`1im`$, where $`_{i+1}`$ is deduced from classical small perturbation of the union $`_iL`$ of the curve $`_i`$ with a line $`L`$.
Hence, when searching to interpret the path $`\stackrel{~}{h}:𝐑\stackrel{~}{}_m𝐑\stackrel{~}{𝒜}_m`$ as a change of relative position of the $`1`$-handles of $`(𝐂_m)S_{g_m}^2`$, $`(𝐂_m)𝐑_m`$, one can proceed by induction and search for $`1j(m1)`$ to move the union of the $`g_j`$ 1-handles of $`𝐂_jS_{g_j}^2`$ with the $`(j1)`$ more 1-handles of $`𝐂_{j+1}S_{g_j+(j1)}^2`$ in such a way that at the end of the process it results $`(𝐂𝒜_m)S_{g_m}^2`$, $`(𝐂𝒜_m)𝐑𝒜_m`$.
Using the recursive Morse-Petroskii’s theory of chapter 4 as a tool to describe the lifting of $`\stackrel{~}{h}`$ as a track of a diffeotopy $`h_t`$ of $`𝐂P^2`$, $`t[0,1]`$, $`h(0)=𝐂_m`$, $`h(1)=𝐂𝒜_m`$, we shall need a global order on the set of crossings of curves $`_iL`$ and the associated $`k`$-uple.
a) global order
According to Theorem 2.3.9 of Chapter 2.3.9, for any Harnack curve $`_m`$ of degree $`m`$, there exists a finite number $`I`$ ($`I=\frac{m(m+1)}{2}+\mathrm{\Sigma }_{k=2}^{k=[m/2]}(2k3)`$) of disjoint $`4`$-balls $`B(a_i)`$ of radius $`ϵ`$ (in the Fubini-Study metric) invariant by complex conjugation centered in points $`a_i`$ of $`𝐑P^2`$ such that up to conj-equivariant isotopy of $`𝐂P^2`$, $`_m\backslash _{iI}B(a_i)=_{i=1}^mL_i\backslash _{i=1}^IB(a_i)`$ where $`L_1,\mathrm{},L_m`$ are $`m`$ distinct projective lines with
$$L_i\backslash _{i=1}^IB(a_i)L_j\backslash _{i=1}^IB(a_i)=\mathrm{}$$
for any $`ij`$, $`1i,jm`$.
The proof of Theorem 2.3.9 of Chapter 2.3.9 is based on an induction on the degree $`m`$ of the curve $`_m`$ which provides a natural order on the lines $`L_j`$, $`1jm`$. Since any $`4`$-ball $`B(a_i)`$) intersects exactly $`2`$ lines $`L_{j1}`$, $`L_j`$ the natural order on the set of lines $`L_j`$ extends to an order on the set of points $`a_i`$. We shall say that $`a_i`$ has order $`j`$, $`2jm`$, if and only if $`B(a_i)`$ intersects the lines $`L_{j1}`$ and $`L_j`$. (In other words, in the inductive construction, any crossing of $`_jL`$ has order $`j`$; and according to Definition 4.2.1 of Chapter 4.2.1, the first point of a $`k`$-uple has the order of the crossing to which it is associated.)
Such order extends on $`3`$-uple (resp the $`2`$-uple) as follows. Given $`a_j`$ of global order $`j`$ the first point of a $`3`$-point $`(a_j,a_{j+1},a_{j+2})`$, ($`j`$ even) associated to a $`a_{j+1}`$ has global order $`j+1`$, $`a_{j+2}`$ has global order $`j+2`$.
Consider the set $`a_j`$ $`1jJ=\frac{m(m+1)}{2}`$ of crossings of curves $`_iL`$ , $`1i(m1)`$ in the inductive construction of $`_m`$. It may be easily easily extracted from the proof of Theorem 2.3.9 of Chapter 2.3.9 (see also Proposition 2.3.7 of Chapter 2.3.7 and Definition 4.2.1 of Chapter 4.2.1) that $`_{i=1}^IB(a_i)_{j=1}^{\frac{m(m+1)}{2}}U(a_j)`$ where $`B(a_i)`$ is the $`4`$-ball around $`a_i`$ of radius $`ϵ`$ and $`U(a_j)`$ is the $`ϵ`$-neighborhood around $`a_j`$. (According to terminology of section 2.3 of Chapter 2.3 (see proposition 2.3.7 of Chapter 2.3.7) the set of points $`_{iI}a_i=_{n=1}^mA_n`$ where $`A_n`$ denotes the set of points perturbed in a maximal simple of deformation of $`_n`$.)
For sake of simplicity, according to corollary 2.3.5 of Chapter 2.3.5, we shall consider the $`T`$-inductive construction of Harnack curves and $`e`$-tubular neighborhoods $`U(p)`$ (see definition 4.2.2 of Chapter 4.2.2) defined from faces $`\mathrm{\Gamma }l_j`$ , $`l_j=\{(x,y)(𝐑^+)^2|x+y=j\}`$. As already introduced in subsection 5.1 of Chapter 5.1, $`𝒫_m`$ denotes the set of points $`p`$ with the property: $`_{p𝒫_m}U(p)=_{a_iI}U(a_i)`$. (In such a way, we have the following equivalent definition of the global order for the patchworking construction. Any point $`a`$ of the set $`a_iI`$ has order $`j+1`$, $`2j+1m`$, if and only if $`aU(p)`$ where $`U(p)`$ is the $`ϵ`$-tubular neighborhood defined from the interior $`\mathrm{\Gamma }^0`$ of a face $`\mathrm{\Gamma }l_j`$ , $`l_j=\{(x,y)(𝐑^+)^2|x+y=j\}`$. This global order extends on $`k`$-uple as previously explained.)
Let us extend the global order on the set points $`a_i,iI`$ to the set $`_{a_iI}U(a_i)`$. We shall say that $`U(a_i)`$ has order $`j`$ if $`a_i`$ has global order $`j`$. Equivalently, in the pathchworking scheme, $`U(p)=U(a_i)`$ has order $`j`$ if it is defined from the interior $`\mathrm{\Gamma }^0`$ of a face $`\mathrm{\Gamma }l_{j1}`$.
b) Covering of $`𝐂P^2`$
The union $`_{iI}U(a_i)`$ does not cover $`𝐂P^2`$. We shall consider consider the family $`\stackrel{~}{U}(a_j)`$, $`_{j=1}^J\stackrel{~}{U}(a_j)_{i=1}^IB(a_i)`$, of neighborhoods of points $`a_j`$, $`jJ`$, verifying the following properties:
\- any $`\stackrel{~}{U}(a_l)`$ of order $`j`$ intersects two neighborhoods $`\stackrel{~}{U}(a_k)`$ of order $`j1`$,
$`_{j=2}^m_{a_ioforderj}\stackrel{~}{U}(a_i)`$ $`=𝐂P^2`$
$`forj=2:`$
$`_{a_ioforder2}\stackrel{~}{U}(a_i)`$ $`𝐂L_2\backslash _{a_ioforder3}\stackrel{~}{U}(a_i)`$
$`_{a_ioforder2}\stackrel{~}{U}(a_i)`$ $`𝐂L_1`$
$`for3j<m:`$
$`_{a_ioforderj}\stackrel{~}{U}(a_i)`$ $`𝐂L_j\backslash _{a_ioforderj+1}\stackrel{~}{U}(a_i)`$
$`_{a_ioforderj}\stackrel{~}{U}(a_i)`$ $`𝐂L_{j1}\backslash _{a_ioforderj1}\stackrel{~}{U}(a_i)`$
$`forj=m:`$
$`_{a_ioforderm}\stackrel{~}{U}(a_i)`$ $`𝐂L_m`$
$`_{a_ioforderm}\stackrel{~}{U}(a_i)`$ $`𝐂L_{m1}\backslash _{a_moforderm1}\stackrel{~}{U}(a_i)`$
Given $`_m`$ the Harnack curve of degree $`m`$ and $`𝒜_m`$ a smooth curve of type $`I`$, we shall in the Proposition 5.2.4 of Chapter 5.2.4, define a path
$$S:[0,1]𝐑𝒞_m$$
$`S(0)=𝐑_m`$, $`S(1)=𝐑𝒜_m`$ described locally in the opens $`\stackrel{~}{U}(a_i)`$ and sequentially by induction on the global order on the set points $`a_i`$,$`iI`$ and the associated $`k`$-uple by diffeotopic perestroikas of perestroikas in the spaces $`𝒞_{I,g}`$ of irreducible curves of degree $`1nm`$, type $`I`$ and genus $`g`$. Besides, any singular point which participated in diffeotopic perestroika belongs to $`_{i=1}^IB(a_i)`$.
###### Proposition 5.2.4.
Let $`_m`$ be the Harnack curve of degree $`m`$ and $`𝒜_m`$ be a smooth curve of type $`I`$. There exists a path
$$S:[0,1]𝐑𝒞_m,S(0)=𝐑_m,S(1)=𝐑𝒜_m$$
which crosses the discriminant hypersurface $`𝐑𝒟_m`$. Up to conj-equivariant isotopy of $`𝐂P^2`$, the path may be sequentially defined in the open $`\stackrel{~}{U}(a_i)`$ by induction on the global order by smoothed diffeotopic perestroikas of diffeotopic perestroikas on the spaces $`𝒞_{I,g}`$ of curves of type $`I`$, degree $`1nm`$, and genus $`g`$ as follows:
The sequence of diffeotopic perestroika is defined as follows:
Assume $`_m`$ obtained via the patchworking method.
1. one can assume that only points $`_{iI}a_i`$ are double-points involved in a diffeotopic perestroika.
2. any diffeotopic perestroika of which real ordinary double-point are points of order at most $`j`$ is a diffeotopic perestroika in the space $`𝒞_{I,g}`$ of curves of degree $`j`$, type $`I`$ and genus $`g`$.
3. a diffeotopic perestroika defines double-points involved in next as follows. As branches involved in the topological meaning of the critical point associated to a real ordinary double-point of order $`j`$ involved in the diffeotopic perestroika move under the diffeotopic perestroika they define real ordinary double-point of order $`j+1`$ involved in a next diffeotopic perestroika. in the space $`𝒞_{I,g}`$ of curves of degree $`j+1`$, type $`I`$ and genus $`g`$.
Any imaginary point which participated to such sequence of diffeotopic perestroika belongs to the intersection of a neighborhood $`\stackrel{~}{U}(a_j)`$ of a point order $`j`$ and a neighborhood of a point $`\stackrel{~}{U}(a_{j+1})`$ of order $`j+1`$. Besides, imaginary points participated in such a way that if one $`\alpha `$-point (or $`\beta `$-point) appears (resp, disappears) after a diffeotopic perestroika on real points of order $`j`$, then it disappears (resp, appears) after a diffeotopic perestroika on real points of order $`j+1`$.
proof: Let us explain the method of our proof. We shall prove that up to slightly modify the coefficients of the polynomial giving the curve $`𝒜_m`$, one can always assume that there exists a diffeotopy $`h_t`$ of $`𝐂P^2`$ $`h(0)=𝐂_m`$, $`h(1)=𝐂𝒜_m`$ with the property $`h_t(𝐂𝒜_mU(p))U(p)`$ for any $`p𝒫_m`$.
Our argumentation is based on the Lemma 5.2.1 of Chapter 5.2.1 and its proof. According to the Lemma 5.2.1 of Chapter 5.2.1, one can consider $`𝐑_m`$ and $`𝐑𝒜_m`$ as smoothed immersions of the circle into the plane with the same Whitney index. We shall denote $`𝐑\stackrel{~}{}_m`$, $`𝐑\stackrel{~}{𝒜}_m`$ the corresponding immersions. From the Whitney’s theorem, there exists a path $`\stackrel{~}{h}`$ which connects $`𝐑\stackrel{~}{}_m`$ and $`𝐑\stackrel{~}{𝒜}_m`$.
Therefore, the definition of the path $`S`$ in $`𝐑𝒞_m`$ $`S(0)=𝐑_m`$,
$`S(1)=𝐑𝒜_m`$ may be deduced from a lifting of the path $`\stackrel{~}{h}`$ which connects $`𝐑\stackrel{~}{}_m`$ and $`𝐑\stackrel{~}{𝒜}_m`$ in the space of immersions of the circle into $`𝐑^2`$.
Since any smooth curve is irreducible, and any reducible polynomial is the product of a finite number of reducible polynomials, we shall, according to Lemma 5.2.1 of Chapter 5.2.1 and Lemma 5.2.3 of Chapter 5.2.3, lift the path $`\stackrel{~}{h}:𝐑\stackrel{~}{_m}𝐑\stackrel{~}{𝒜}_m`$ in the space of immersion of the circle into the plane from smoothed diffeotopic perestroikas of diffeotopic perestroikas in the spaces $`𝒞_{I;g}`$ of curves of degree $`nm`$, type $`I`$ and genus $`0g\frac{(n1).(n2)}{2}`$.
For any diffeotopic perestroika $`\pi `$ experienced along a generic path in $`𝒞_{I;g}`$, locally a description of the real components of the resulting curve is deduced from double points involved in the diffeotopic perestroika. Thus, any smoothed diffeotopic perestroika $`\pi `$ experienced along $`S`$, may be defined from its smoothed double points.
From the detailed construction of $`_m`$ given in chapter 4, we shall get smoothed points of smoothed diffeotopic perestroika experienced along the path $`S`$.
As already introduced in the proof of Lemma 5.2.1 of Chapter 5.2.1, we shall consider $`h`$ the diffeomorphism $`h:𝐂_mS_{g_m}^2𝐑^3`$ and $`r:𝐑^3𝐑^2`$ the projection which maps $`h(𝐑_m)`$ to $`r(h(𝐑_m)=𝐑\stackrel{~}{}_m𝐑^2`$. We shall prove that only points of $`𝐑\stackrel{~}{}_m`$ smoothed in $`𝐑^2`$ to give $`𝐑_m`$ may lift to smoothed double points of an irreducible curve of degree $`n`$, $`1nm`$ type $`I`$ and genus $`g`$ involved in smoothed diffeotopic perestroika. Thus, double-points of $`𝐑\stackrel{~}{}_m`$ smoothed at infinity are involved along the path $`\stackrel{~}{h}:𝐑\stackrel{~}{}_m𝐑\stackrel{~}{𝒜}_m`$ as the relative position of the $`1`$-handles of $`(𝐂_m)S_{g_m}^2`$, $`(𝐂_m)𝐑_m`$ changes to give $`(𝐂𝒜_m)S_{g_m}^2`$, $`(𝐂𝒜_m)𝐑𝒜_m`$.
In such a way, using the path provided by Lemma 5.2.1 of Chapter 5.2.1 and Whitney’s theorem, according to the Lemma 5.2.3 of Chapter 5.2.3 and the fact that $`𝐂_m`$ and $`𝐂𝒜_m`$ are diffeotopic, we shall characterize the track $`S`$ on $`𝐑P^2`$ of a diffeotopy $`h_t`$ of $`𝐂P^2`$, $`t[0,1]`$, $`h(0)=𝐂_m`$, $`h(1)=𝐂𝒜_m`$.
We shall divide our proof into three parts:
1. In the first part, we define double-points of curves along $`S`$.
2. The second part deals with smoothed diffeotopic perestroikas along $`S`$.
3. In the third part, we give a method which provides any curve $`𝒜_m`$ of degree $`m`$ and type $`I`$ from $`_m`$, that is $`𝐑𝒜_m`$ from $`𝐑_m`$.
The first two parts make use of the idea that locally a real branch does not differ from another real branch. In these parts, we consider the situation more locally than globally and fix definitions we shall need in the third part. In the third part, we give a method of organizing these local moves in such a way that globally it results a real algebraic curve in $`𝐂P^2`$. The exposition of the main ideas of our proof is now finished, and we shall proceed to precise arguments.
1. Double points of curves along $`S`$.
Let us start our definition of double points of curves along the path $`S`$ we have planed to construct by the following remark. Assume $`_m`$ constructed via the patchworking method. Given $`_m`$, and its projection $`r(h(𝐑_m))=𝐑\stackrel{~}{_m}`$, it is easy to deduce from the Lemma 5.2.1 of Chapter 5.2.1 and its proof and Lemma 5.2.3 of Chapter 5.2.3, that for $`m>2`$ only the $`g_m`$ points of $`𝐑\stackrel{~}{}_m`$ smoothed in $`𝐑^2`$ to give $`𝐑_m`$ may lift to real double points of curves along $`S`$.
Denote $`\rho ^m:(T_m\times U_𝐂^2)𝐂T_m𝐂P^2`$ the natural surjection. According to the chapter 2 (see Proposition 2.3.7 of Chapter 2.3.7 and Theorem 2.3.9 of Chapter 2.3.9), the Harnack curve $`_m`$ of degree $`m`$ is obtained from recursive perturbation of curves $`_{mi}L`$, $`0<i<m`$, ( with $`L𝐂l_{mi}=\rho ^m(l_{mi}\times U_𝐂^2)`$) in such way that outside a finite number $`I`$ of $`4`$-ball $`B(a_i)𝐂P^2`$ there exists a conj-equivariant isotopy $`j_t`$ of subset of $`𝐂P^2`$ which maps:
1. $`_{mi}\backslash _{i=1}^IB(a_i)`$ onto a part of the curve $`_{mi+1}`$.
(in the patchworking scheme, such part is contained in the restriction $`\rho ^m(T_{mi}\times U_𝐂^2)`$ of $`𝐂P^2𝐂T_m`$.)
2. $`(_{mi}\backslash j_1(_{mi+1}\backslash _{i=1}^IB(a_i))`$ onto a part of the projective line $`L`$.
(in the patchworking scheme, such part is contained in the restriction $`\rho ^m(D_{mi+1,mi}\times U_𝐂^2)`$ of $`𝐂P^2𝐂T_m`$.)
Consequently, (see Theorem 2.3.9 of Chapter 2.3.9), outside a finite number of $`4`$-balls $`B(a_i)`$, the curve $`_m`$, is up to conj-equivariant isotopy of $`𝐂P^2`$, the union of $`m`$ projective lines minus their intersections with $`4`$-balls $`B(a_i)`$. Since projective lines are diffeotopic to circles, outside a finite number of $`4`$-balls $`B(a_i)`$, components of $`𝐑_m`$ can be considered as components of $`m`$ concentric circles. Thus, smoothed real points involved in diffeotopic perestroikas along $`S`$ belong to the union of real components of $`_m`$ contained in $`\rho ^m(D_{mi+1,mi}D_{mi,mi1}\times U_𝐂^2)`$ where $`D_{mi+1,mi}`$ and $`D_{mi,mi1}`$ are bands of $`T_m`$ (for $`j\{mi1,mi\}`$, $`D_{j+1,j}=\{x0,y0|jx+yj+1\}`$) with intersection $`D_{mi+1,mi}D_{mi,mi1}=l_{mi}`$.
Therefore, locally any smoothed real point can be considered as the smoothed intersection of two real branches of $`𝐑_m\rho ^m(D_{mi+1,mi}D_{mi,mi1}\times U_𝐑^2)`$ of which intersection point belongs to a neighborhood of $`\rho ^m(l_{mi}\times U_𝐑^2)`$.
Consequently, it is easy to deduce from relative orientation and relative location of real branches of $`𝐑_m\rho ^m(D_{mi+1,mi}D_{mi,mi1}\times U_𝐂^2)`$ that smoothed real points involved in diffeotopic perestroika along $`S`$ belong to $`\rho ^m(D(p)\times U_𝐑^2)=U(p)𝐑^2`$ with $`p𝒫_m`$ and that any $`U(p)`$ contains one such smoothed double-point.
In such a way, according to Lemma 5.2.1 of Chapter 5.2.1 and Lemma 5.2.3 of Chapter 5.2.3, smoothed real points along $`S`$ are smoothed crossings which belong to the set
$`_{p𝒫_m}\rho ^m(D(p)\times U_𝐑^2)=_{p𝒫_m}U(p)𝐑^2`$.
2.Smoothed diffeotopic perestroika on smoothed double points of curves along $`S`$.
Recall that from the previous Morse-Petrovskii’s study of Harnack curves (see chapter 4) for any $`p𝒫_m`$, there exists a crossing $`a_{mi}L`$ with $`aU(p)=\rho ^m(D(p)\times U_𝐂^2)`$ such that the crossing $`a𝒫_m`$ is associated to:
1. either a simple-point of $`_m`$
2. either a $`3`$-point $`(p_{mi+1},p_{mi+2},p_{mi+3})`$ (where $`p_{mi+3}p_m`$ of $`_m`$
3. or in case $`m`$ odd a $`2`$-point $`(p_{m1},p_m)`$ of $`_m`$
As recalled previously, cusp perestroika and weak triple point perestroika do not change real and complex smoothed parts. Therefore, there is not need to consider them. Besides, since smoothed real points of generic curves along $`S`$ are smoothed crossings, no smoothed imaginary self-tangency perestroika is undergone along $`S`$. (This restriction is not in contradiction with the fact that $`𝒜_m`$ is an arbitrary curve of degree $`m`$ and type $`I`$ since $`_m`$ is an $`M`$-curve.)
Obviously, in $`U(p)`$ any change implied by a smoothed diffeotopic perestroika is a perturbation in which, besides real branches of the crossing $`aU(p)`$ of $`_jL`$, real branches involved in the local topological meaning of simple point and of any point of a $`3`$-point (or $`2`$-point in case $`m`$ odd) participated.
We shall study the situations inside the four $`4`$-balls with non-empty real part and globally invariant by complex conjugation of $`U(p)(𝐂^{})^2`$ simultaneously.
We shall distinguish perestroikas in which only simple points are involved and perestroikas in which $`k`$-points are involved.
2.a Smoothed diffeotopic perestroikas defined on simple points.
Let us consider smoothed diffeotopic perestroikas defined on simple points.
Recall that given $`a`$ a crossing of $`_jL`$, we say that $`p_l`$ the critical point of $`_l`$ $`j+1lm`$, associated to $`a`$ is simple point if:
1. $`lj+1`$,
2. any critical point $`p_k`$, $`j+1kl`$ is equivalent to $`p_{j+1}`$.
Given $`a`$ a crossing of $`_jL`$ ($`L𝐂l_j`$) $`1jm`$ denote $`(x_0,y_0)`$ the simple point of $`_m`$, associated the crossing $`a`$. Let $`p𝒫_m`$, such that $`a,(x_0,y_0)\rho ^m(D(p)\times U_𝐑^2)=U(p)(𝐑^{})^2`$ Real branches in $`U(p)`$ are the real branches involved in the topological meaning of $`(x_0,y_0)`$ and the real branches which result from the perturbation of the crossing $`a`$. Let us remark that if the relative position of real branches which result from the perturbation of the crossing $`a`$ is preserved, then relative position of the real real branches involved in the topological meaning of $`(x_0,y_0)`$ is also preserved. Otherwise, according to Lemmas 4.2.6 of Chapter 4.2.6, 4.2.7 of Chapter 4.2.7, 4.2.9 of Chapter 4.2.9 and Lemma 5.2.1 of Chapter 5.2.1, it would lead with contradiction with the fact that we consider only curves of type $`I`$ that is only Morse modification in the direction coherent with complex orientation.
As already noticed, in most perestroikas imaginary points are involved. Considering together imaginary and real points involved in these perestroikas, we shall define smoothed diffeotopic perestroikas of perestroika along a path $`h`$ in the space $`𝒞_{I,g})`$ (see section 5.1 of Chapter 5.1) and get relative position of real branches in $`U(p)`$ after such diffeotopic perestroikas.
###### Lemma 5.2.5.
Let $`\pi `$ be either a direct or an inverse self tangency smoothed diffeotopic perestroika or a triple point with imaginary branches smoothed diffeotopic perestroika along the path
$$S:[0,1]𝐑𝒞_m$$
$`S(0)=𝐑_m`$, $`S(1)=𝐑𝒜_m`$. Let $`c\rho ^m(D(p)\times U_𝐑^2)=U(p)(𝐑^2)`$, with $`p𝒫_m`$ (which may be $`a`$ or $`(x_0,y_0)`$) a real point involved in $`\pi `$.
Then, under $`\pi `$ relative position of real branches in $`U(p)`$ changes as follows.
One can always assume that the crossing $`a_jL`$ is the real (smoothed) point involved in the perestroika, the relative position of real branches involved in the local topological meaning of $`(x_0,y_0)`$ is changed as follows. Draw a dotted (unknotted) line between the two real branches involved in the local topological meaning of the critical point $`(x_0,y_0)`$ The diffeotopic perestroika moves the real component around $`a`$ according to the perestroika and leaves the two real branches linked by the unknotted dotted line.
proof:
(a)
Let $`\pi `$ be a direct self tangency perestroika or inverse self tangency perestroika in the negative direction.
Denote $`(x_0,y_0)`$ the critical point of $`B_{j+1}`$ associated to the crossing $`a`$ such that
$`a,(x_0,y_0)\rho ^m(D(p)\times U_𝐑^2)=U(p)(𝐑^{})^2`$.
Denote $`c,c^{}`$ the two real points involved in $`\pi `$.
Let $`p,p^{}𝒫_m`$ such that $`c\rho ^m(D(p)\times U_𝐑^2)=U(p)𝐑^2`$,
$`c^{}\rho ^m(D(p^{})\times U_𝐑^2)=U(p^{})𝐑^2`$ Then, imaginary points $`i,i^{}`$ involved in $`\pi `$ are such that: $`i\rho ^m(D(p)\times U_𝐂^2)=U(p)`$,
$`i^{}\rho ^m(D(p^{})\times U_𝐂^2)=U(p^{})`$.
In an $`ϵ`$-neighborhood $`U(p)`$, $`p𝒫_m`$, the complex point set of the Harnack curve $`_{j+1}`$ can be considered as the image of a smooth section of a tubular fibration $`N(𝐂V\backslash \{a\}))`$ where $`V=(_{j+1}L)U(p)`$ (i.e, from Morse Lemma, in a neighborhood of $`a`$, $`V`$ looks like the intersection of two real lines in the point $`a`$) and $`N`$ is an $`ϵ`$-tubular neighborhood of $`V`$.
Assume $`c=a`$, (resp $`c=(x_0,y_0)`$), relative location of the real branches in $`U(p)`$ is changed after diffeotopic perestroika in such a way that after diffeotopic perestroika the complex point set is the image of a smooth section of a tubular fibration $`N(𝐂V\backslash \{a\}))`$. Therefore, since $`(x_0,y_0)`$ (resp, $`a`$) is not a real point involved in the perestroika, the complex point $`i`$ lies in complex 4-ball around $`(x_0,y_0)`$ which contains real branches involved in the topological meaning of $`(x_0,y_0)`$ (resp, which results from the perturbation of the singular crossing $`a`$). Hence, it follows the relative position of the real branches in $`U(p)`$. Obviously, the relative position of the real branches inside $`U(p^{})`$ may be deduced from the same argument. (b)
Consider now $`\pi `$ a triple point perestroika with imaginary branches. Obviously, triple point with imaginary branches perestroika in the positive (resp, negative) direction may appear only after direct (resp, inverse) self tangency perestroika. Such perestroika provides two imaginary smoothed points, each of which belongs to an $`ϵ`$-neighborhood $`U(p)`$ $`p𝒫_m`$.
Assume an oval $`𝒪`$ is involved in a smoothed triple point perestroika with imaginary branches. Then, $`𝒪`$ intersects $`_{p𝒫_m}U(p)`$ in four of its neighborhood $`U(p)`$. Each of these four neighborhoods $`U(p)`$ contains one imaginary point involved in the perestroika: two of them contain one of the two imaginary points required before the perestroika, and the two others contain one of the two imaginary points which appear after the perestroika.
Q.E.D
2.b Smoothed diffeotopic perestroikas defined on 2-points and 3-points.
We shall now consider crossings associated to $`3`$-point of $`_m`$ and $`2`$-point in case of odd $`m`$ in the construction of $`_m`$. Let us prove the following Lemma.
###### Lemma 5.2.6.
A smoothed strong real triple point diffeotopic perestroika $`\pi `$ may be experienced along the path
$$S:[0,1]𝐑𝒞_m$$
$`S(0)=𝐑_m`$, $`S(1)=𝐑𝒜_m`$ only in a neighborhood which contains a crossing $`a_jL`$ associated to a triple-point or a double-point of the curve $`_m`$.
proof:
Recall that given $`a`$ a crossing of one curve $`_{j1}L`$ (with $`j`$ even) associated to a $`3`$-uple $`(p_j,p_{j+1},p_{j+2})`$ of $`_{j+2}`$, $`j+2m`$, we consider the following natural order on the $`3`$-uple $`(p_j,p_{j+1},p_{j+2})`$: $`p_j`$ has order $`j`$, $`p_{j+1}`$ has order $`j+1`$, $`p_{j+2}`$ has order $`j+2`$.
1Let us first consider the $`2`$-uple $`(p_j,p_{j+1})`$.
1.a Let $`b_j(p_j)=c_0`$ with $`(p_j,p_{j+1})`$ a $`2`$-uple of the first or the second kind. It follows from Lemma 4.2.7 of Chapter 4.2.7 of chapter 4 that in case of $`2`$-uple of the first kind: when $`c=c_0`$, the non-empty positive oval touches a positive outer oval; and from Lemma 4.2.9 of Chapter 4.2.9 of chapter 4 that in case of $`2`$-uple of the second kind: when $`c=c_0`$, the non-empty positive oval touches a negative outer oval.
1.b Let $`b_{j+1}(p_{j+1})=c_0`$ with $`(p_j,p_{j+1})`$ a $`2`$-uple of the first or the second kind. When $`c=c_0`$, the one-side component of the curve $`_{j+1}`$ touches itself. Then, as $`c`$ increases from $`c_0`$ to $`c_0+ϵ`$, one negative oval $`𝒪^{}`$ disappears.
Besides, when consider together local topological meaning of $`p_j`$ and $`p_{j+1}`$ three real branches are involved which belong respectively to a part of $`_{j1}`$ isotopic to the real line $`L_{j1}𝐑l_{j2}`$, a part of $`_j`$ isotopic to the real line $`L_j𝐑l_{j1}`$, a part of $`_{j+1}`$ isotopic to the real line $`L_{j+1}𝐑l_j`$. Besides, the branches which are respectively isotopic to a part of the line $`L_j𝐑l_{j1}`$ and isotopic to a part of the line $`L_{j+1}𝐑l_j`$ belong to an oval $`𝒪`$ (negative in case of $`2`$-uple of the first kind and positive in case of $`2`$-uple of the second kind) which results from the perturbation of $`_jL`$, $`L_{j+1}𝐂l_j`$, around two crossings.
Consider $`𝐂_{j+1}`$ inside $`\rho ^m(D_{j,j1}D_{j+1,j}\times U_𝐑^2)𝐑P^2`$ where
$`D_{j+1,j}=\{x0,y0|jx+yj+1\}`$,
$`D_{j,j1}=\{x0,y0|j1x+yj+1\}`$.
It is an easy consequence of the Lemma 4.2.7 of Chapter 4.2.7 and 4.2.9 of Chapter 4.2.9 of chapter 4 that one can consider the change from $`p_j`$ to $`p_{j+1}`$ as a ”jump” of the branch of $`_{j1}`$ isotopic to the real line $`L𝐑l_{j2}`$ from $`\rho ^m(D_{j,j1}\times U_𝐑^2)`$ to $`\rho ^m(D_{j+1,j}\times U_𝐑^2)`$, namely as a smoothed triple point perestroika.
2 Consider now $`3`$-uple $`(p_j,p_{j+1},p_{j+2})`$ of the first and second kind.
Let $`b_{j+2}(p_{j+2})=c_0`$,
2.a
in case of $`3`$-uple of the first kind, it follows from Lemma 4.2.7 of Chapter 4.2.7 of chapter 4, that when $`c=c_0`$ one positive outer oval of the curve $`_{j+2}`$ touches another outer positive oval;
2.b
in case of $`3`$-uple of the first kind, it follows from Lemma 4.2.9 of Chapter 4.2.9 of chapter 4, that when $`c=c_0`$ one negative inner oval of the curve $`_{j+2}`$ touches another negative oval.
It is obvious that the previous description of the strong triple point perestroika defined on the $`2`$-uple $`(p_j,p_{j+1})`$ extends to the definition of a strong triple point perestroika defined on the $`3`$-uple $`(p_j,p_{j+1},p_{j+2})`$ in such a way that relative position and orientation of the two real branches involved in the topological meaning of $`p_{j+2}`$ does not change under the perestroika.
Therefore, the Lemma is straightforward consequence of the previous study of the $`2`$-uple and $`3`$-uple. The varieties $`𝐂_m`$ and $`𝐂𝒜_m`$ are diffeotopic; under the perestroika the diffeotopy is preserved locally in any neighborhood $`U(p)`$, $`p𝒫_m`$. Q.E.D
As explained in section 5.1 of Chapter 5.1, any smoothed perestroika $`\pi `$ required a given location and orientation of real branches. In particular, smoothed triple-point perestroika is possible only in a neighborhood of an oval and of three real branches close to this oval with relative location and orientation as described in section 5.1 of Chapter 5.1. Namely, the real branches involved in the local topological meaning of $`p_j`$ may be involved in a triple point diffeotopic perestroika only if they have been reproached by previous perestroikas. It is easy to deduce from the relative location of branches involved in the topological meaning of critical point that in case of $`2`$-uple of the first kind, the singular situation $`b_{j+1}(p_{j+1})=c_0`$ is possible if and only if two positive outer ovals coalesce with the one-side component of the curve. In case of $`2`$-uple of the second kind, the singular situation $`b_{j+1}(p_{j+1})=c_0`$ is possible if and only if two negative inner ovals coalesce with the one-side component of the curve. We shall precise this remark in the third part.
3. Method which provides $`(𝐂P^2,𝐂𝒜_m)`$ from $`(𝐂P^2,𝐂_m)`$
Denote $`𝒮_{I,m}`$ the set $`𝒞_{I,m}\backslash 𝒟_m`$ of smooth curves of degree $`m`$ and type $`I`$.
We shall define a family of moves $`\varphi :_m𝒮_{I,m}`$ with the property that in these moves the set $`_{i=1}^Ia_i=_{n=1}^mA_n`$ of crossing-points perturbed in the construction of $`_m`$ participate as follows. Moves of real components of $`_m`$ correspond to an other choice of perturbation on these crossings. Briefly the method of our proof is the following:
We shall apply successively Lemma 5.2.5 of Chapter 5.2.5 and 5.2.6 of Chapter 5.2.6 inside the sets:
$$\rho ^m((D_{1,0}D_{2,1})\times U_𝐂^2)$$
$$\rho ^m((D_{2,,1}D_{3,2})\times U_𝐂^2)$$
$$\rho ^m((D_{3,2}D_{2,1})\times U_𝐂^2)$$
$$\mathrm{}$$
$$\rho ^m((D_{mi+1,mi}D_{mi+2,mi+1})\times U_𝐂^2)$$
$$\rho ^m((D_{mi+3,mi+2}D_{mmi+2,mi+2})\times U_𝐂^2)$$
$$\mathrm{}$$
$$\rho ^m((D_{m1,m2}D_{m2,m3})\times U_𝐂^2)$$
$$\rho ^m((D_{m,m1}D_{m1,m2})\times U_𝐂^2)$$
to describe the track $`S`$ on $`𝐑P^2`$ of a diffeotopy $`h_t`$ of $`𝐂P^2`$, $`t[0,1]`$, $`h(0)=𝐂_m`$, $`h(1)=𝐂𝒜_m`$ with $`𝒜_m`$ a curve of type $`I`$ and degree $`m`$.
We shall now proceed to precise argument.
3.a Preliminaries Recall that, according to chapter 4, for any point $`a`$ perturbed in the construction of $`_m`$, there exists $`p𝒫_m`$ (where $`\mathrm{}𝒫_m=\frac{m(m1)}{2}`$) such that $`aU(p)`$ where $`U(p)`$ is the $`ϵ`$-tubular neighborhood defined from the interior $`\mathrm{\Gamma }^0`$ of a face $`\mathrm{\Gamma }`$ of the triangulation of $`T_m`$. According to Lemma 5.2.1 of Chapter 5.2.1 $`𝐑_m`$ (resp, $`𝐑𝒜_m`$) is a smoothed generic immersion of the circle $`S^1`$ into the plane $`𝐑^2`$, $`𝐑\stackrel{~}{}_m`$ (resp, $`𝐑\stackrel{~}{𝒜}_m`$) with $`n\frac{(m1)(m2)}{2}`$ double points and Whitney index $`\frac{(m1)(m2)}{2}+1`$. The smoothing is such that $`\frac{(m1)(m2)}{2}`$ double points are smoothed in $`𝐑^2`$ by Morse modification in the direction coherent to the complex orientation. We shall denote by $`𝒢_m`$ the set of the $`g_m=\frac{(m1)(m2)}{2}`$ double points of $`𝐑\stackrel{~}{}_m`$ smoothed in $`𝐑^2`$ by Morse modification in the direction coherent to the complex orientation.
According to the part 1, we shall assume that the set $`𝒢_m`$ of double-points of $`\stackrel{~}{𝐑_m}`$ has the following properties:
1. any two distinct points of $`𝒢_m`$ belong to two distinct neighborhoods $`U(p)`$ with $`p𝒫_m`$
2. there exists a bijective correspondence between points of $`𝒢_m`$ and the $`g_m`$ ovals $`𝒪`$ of $`_m`$ which do not intersect the line at infinity.
Properties (1) and (2) are straightforward consequences of the patchworking construction (see for example Corollary 1.2.4 of Chapter 1.2.4 and its proof) that any oval $`𝒪`$ of $`_m`$ which does not intersect the line at infinity, intersects two neighborhoods $`U(p)`$ defined from a face $`\mathrm{\Gamma }l_j`$ ,$`1jm1`$, $`l_j=\{(x,y)(𝐑^+)^2|x+y=j\}`$, namely two neighborhoods $`U(p)`$ of global order $`j+1`$.
Besides, an oval $`𝒪`$ may intersect:
1. either $`4`$ neighborhoods $`U(p)`$ of the union $`_{p𝒫_m}U(p)`$ such that two neighborhoods $`U(p)`$ are of global order $`j`$, $`3jm1`$, the two others are of global order $`j1`$, $`j+1`$
2. or $`3`$ neighborhoods $`U(p)`$ of the union $`_{p𝒫_m}U(p)`$ such that two neighborhoods $`U(p)`$ are of global order $`m`$, and the other is of global order $`m1`$.
Therefore, one can choose the set $`𝒢_m`$ such that any point of $`𝒢_m`$ belongs to a neighborhood $`U(p)`$, $`p𝒫_m`$ of global order $`j`$, $`2jm1`$. Consider the family $`\stackrel{~}{U}(p)U(p)`$, $`p𝒫_m`$ $`_{p𝒫_m}\stackrel{~}{U}(p)=𝐂P^2`$ verifying the properties (5.2), (5.2), (5.2), (5.2). Using the patchworking construction, according to the moment map $`\mu :𝐂T_mT_m`$ and the natural surjection $`\rho ^m:𝐑^+T_m\times U_𝐂^2𝐂T_m`$ (see the preliminary section ), one can choose for example:
1. for any $`U(p)`$ of global order $`2jm1`$,
$`\stackrel{~}{U}(p)=\rho ^m(S\times U_𝐂^2)`$ where $`S`$ is the unique square $`S`$ with vertices $`(c,d),(c+1,d),(c,d+1),(c+1,d+1)`$, $`c+d+1=j`$ such that $`U(p)\stackrel{~}{U(p)}`$.
2. for any $`U(p)`$ of global order $`m`$, $`\stackrel{~}{U}(p)=\rho ^m(K\times U_𝐂^2)`$ where $`K`$ is the polygon with vertices $`(c,d),(c+2,d),(c,d+2)`$, $`c+d+1=m1`$ such that $`U(p)\stackrel{~}{U}(p)`$.
3.b Statement of the Method
We shall now deduce a lifting of the path $`\stackrel{~}{h}`$ in the space of immersions of the circle into $`𝐑^2`$ with index $`g_m=\frac{(m1)(m2)}{2}`$ in the space $`𝐑𝒞_m`$ and describe the topological pair $`(𝐂P^2,𝐂𝒜_m)`$. Such lifting is defined by smoothed diffeotopic perestroikas on $`𝐑_m`$ inside $`_{p𝒫_m}\stackrel{~}{U}(p)`$.
Recall that, according to parts 1 and 2 of our proof, we shall consider only direct or inverse self tangency perestroika, triple point with imaginary branches perestroika, strong real triple point perestroika. Obviously, in any smoothed diffeotopic perestroika $`\pi `$ undergone along the lifting $`S`$ of $`\stackrel{~}{h}:𝐑\stackrel{~}{}_m𝐑\stackrel{~}{𝒜}_m`$, at least one oval $`𝒪`$ of $`𝐑_m`$ which does not intersect the line at infinity is involved. As already noticed, in most perestroikas imaginary points are involved. We shall consider the function $`J`$ counting increment and decrement of $`\alpha `$-point (double point of different halves) and $`\beta `$-point (double point of one of the halves) along the path $`S:[0,1]𝐑𝒞_m`$. Let $`𝒪`$ be an oval of $`_m`$. It intersects $`_{p𝒫_m}\stackrel{~}{U}(p)`$ in four neighborhood $`\stackrel{~}{U}(p)`$, $`p𝒫_m`$: $`U(p_1)`$ of global order $`i`$ with $`2im1`$, , $`U(p_2)`$ and $`U(p_3)`$ of global order $`i+1`$, $`U(p_4)`$ of global order $`i+2`$.
We shall, from the Lemma 5.2.5 of Chapter 5.2.5 and 5.2.6 of Chapter 5.2.6, get moves of the real branches contained in $`U(p_i)`$, $`1i4`$, under any smoothed diffeotopic perestroika $`\pi `$.
We shall prove that smoothed diffeotopic perestroikas on the real components of $`_m`$ of which smoothed points are points of global order $`2jm1`$ ( in the union of neighborhoods $`\stackrel{~}{U}(p)`$, $`p𝒫_m`$ of global order $`j`$, $`2jm1`$), induce smoothed diffeotopic perestroikas on the real component of $`_m`$ in neighborhoods $`\stackrel{~}{U}(p)`$ $`p𝒫_m`$ of global order $`j+1`$.
We shall call this process the step $`j`$ of the lifting, and define the lifting of $`\stackrel{~}{h}:𝐑\stackrel{~}{}_m𝐑\stackrel{~}{𝒜}_m`$ by induction on the global order on the set of points $`a_i,iI`$ and the associated $`k`$-uple.
1) Let us first consider direct and inverse self-tangency and triple point with imaginary branches diffeotopic perestroika. According to Lemma 5.2.5 of Chapter 5.2.5 one can assume, without loss of generality that the smoothed real points involved in triple point with imaginary branches diffeotopic perestroika are the two crossings which belong to the neighborhoods $`U(p_2)`$, $`U(p_3)`$ of global order $`im1`$. In such a way, according to the Lemma 5.2.5 of Chapter 5.2.5, we get description of relative position of real branches in $`U(p_2)`$ and $`U(p_3)`$ and the smoothed points of global order $`j+1`$ which may be involved in a next diffeotopic perestroika is deduced
Besides, along the path $`S`$, as the curve $`_m`$ experiences various diffeotopic perestroikas $`J(\alpha )`$ and $`J(\beta )`$ increases or decreases. It follows from the Lemma 5.2.5 of Chapter 5.2.5 that imaginary points involved in a perestroika $`\pi `$ belong to neighborhoods $`U(p_i)`$, $`1i4`$, in such a way that diffeotopic perestroikas in neighborhoods of global order $`j`$ provide or smooth the imaginary points involved in diffeotopic perestroika in neighborhoods $`\stackrel{~}{U}(p)`$ of order $`j+1`$. 2)Consider now $`\pi `$ a strong triple diffeotopic perestroika.
From the Lemma 5.2.6 of Chapter 5.2.6, strong triple point smoothed diffeotopic perestroika $`\pi `$ requires a neighborhood $`\stackrel{~}{U}(p)`$ which contains a crossing $`a`$ associated to a triple-point or a double-point of the curve $`_m`$. It is an easy consequence of the Lemma 5.2.6 of Chapter 5.2.6, that, according to our notations, the $`k`$-uple ($`k\{2,3\}`$) is of the form $`(p_i,p_{i+1},p_{i+2})`$, $`i`$ even (resp, $`(p_i,p_{i+1})`$, $`i`$ even) with $`p_iU(p_1)`$ belongs to the neighborhood $`\stackrel{~}{U}(p_1)`$ of global order $`i`$ with $`i+1jm1`$. As already noticed in the proof of Lemma 5.2.6 of Chapter 5.2.6, the restriction of $`3`$-uple $`(p_i,p_{i+1},p_{i+2})`$ to $`(p_i,p_{i+1})`$ is sufficient to define $`\pi `$; relative position of real branches involved in the topological meaning of $`p_{i+2}`$ under $`\pi `$ follows by induction.
Besides,the diffeotopic perestroika $`\pi `$ is possible only if orientation and relative location of the real branches in the neighborhoods $`U(p_1)`$, $`U(p_2)`$, $`U(p_3)`$, $`U(p_4)`$ is one of the two required (see section 5.1 of Chapter 5.1). It follows from the Lemma 5.2.6 of Chapter 5.2.6 except around real branches involved in $`\pi `$, the situation remains the same inside $`_{1i4}\stackrel{~}{U}(p)`$ under diffeotopic perestroika.
In such a way, at the end of the step $`j`$ of the lifting, we get an irreducible curve of degree $`j`$ and type $`I`$ inside $`\rho ^j(T_j\times U_𝐂^2)`$. This irreducible curve of degree $`j`$ has no real double-point singularity but may have an even number of imaginary double-point singularity. The step $`j+1`$ of the lifting can be interpreted as the perturbation of the union of this curve of degree $`j`$ with a line -namely, according to the Lemma 5.2.1 of Chapter 5.2.1 and its proof as the way to move the union of the $`g_j`$ 1-handles of $`𝐂_jS_{g_j}^2`$ with the $`(j1)`$ more 1-handles of $`𝐂_{j+1}S_{g_j+(j1)}^2`$-.
We have considered only neighborhoods $`\stackrel{~}{U}(p)`$ of global order $`j`$, $`2jm1`$, and $`\mathrm{}𝒫_mg_m=m1`$. Nonetheless, the lifting of $`\stackrel{~}{h}:𝐑\stackrel{~}{}_m𝐑\stackrel{~}{𝒜}_m`$ is entirely defined by induction on the global order, and we get the topological pair $`(𝐂P^2,𝐂𝒜_m)`$ from the method proposed. Indeed, consider the $`m2`$ ovals $`𝒪`$ which intersect two neighborhoods $`\stackrel{~}{U}(p)`$ of global order $`m`$, and one neighborhood $`\stackrel{~}{U}(p)`$ of global order $`m1`$. (The union of these $`m2`$ ovals intersects any neighborhood $`\stackrel{~}{U}(p)`$, $`p𝒫_m`$ of global order $`m`$.) Any of these m-2 ovals $`𝒪`$ is involved in a diffeotopic perestroika implied (in the chain of diffeotopic perestroika) by a diffeotopic perestroika on neighborhood of order $`m1`$. Hence, the lifting of $`\stackrel{~}{h}:𝐑\stackrel{~}{}_m𝐑\stackrel{~}{𝒜}_m`$ follows by induction.
Consequently, we get description of $`𝐂𝒜_m`$ in any neighborhood $`\stackrel{~}{U}(p)`$, $`p𝒫_m`$ and thus in $`_{p𝒫_m}\stackrel{~}{U}(p)=𝐂P^2`$.
According to Lemma 5.2.1 of Chapter 5.2.1 and Lemma 5.2.3 of Chapter 5.2.3, we have defined a lifting of the path $`\stackrel{~}{h}:𝐑\stackrel{~}{}_m𝐑\stackrel{~}{𝒜}_m`$ in the space of generic immersion of the space of immersions of a circle into the plane to $`S:𝐑_m𝐑𝒜_m`$ which is the track on $`𝐑P^2`$ of a diffeotopy $`h_t`$ of $`𝐂P^2`$, $`t[0,1]`$, $`h(0)=𝐂_m`$, $`h(1)=𝐂𝒜_m`$. Q.E.D
##### Description of pairs $`(𝐂P^2,𝐂𝒜_m)`$ up to conj-equivariant isotopy
In this section, we shall first give a combinatorial method to describe pairs $`(𝐂P^2,𝐂𝒜_m)`$, up to conj-equivariant isotopy, where $`𝒜_m`$ is a smooth curve of degree $`m`$ and type $`I`$. Then, we extend in Theorem 5.2.9 of Chapter 5.2.9 the description of the pair $`(𝐂P^2,𝐂_m)`$ stated in Theorem 2.3.9 of Chapter 2.3.9 for Harnack curve to any smooth curve of type $`I`$.
Combinatorial Description of pairs $`(𝐂P^2,𝐂𝒜_m)`$ In this part, we shall detail the method given in the third part of the proof of the Proposition 5.2.4 of Chapter 5.2.4 and describe the family of moves $`\varphi :_m𝒮_{I,m}`$ where $`𝒮_{I,m}=𝒞_{I,m}\backslash 𝒟_m`$. For any curve $`𝒜_m`$ of degree $`m`$ and type $`I`$, we have defined a lifting of the path $`\stackrel{~}{h}:𝐑\stackrel{~}{}_m𝐑\stackrel{~}{𝒜}_m`$ which is the track $`S`$ on $`𝐑P^2`$ of a diffeotopy $`h_t`$ of $`𝐂P^2`$, $`t[0,1]`$, $`h(0)=𝐂_m`$, $`h(1)=𝐂𝒜_m`$.
(Obviously, the path is not necessarily lifted continuously in the space of real algebraic curves in the sense that we can not affirm that any immersion of the circle into $`𝐑^2`$ with index $`g_m`$ along an arbitrary path from $`𝐑\stackrel{~}{}_m`$ to $`𝐑\stackrel{~}{𝒜}_m`$ lift to a curve of degree $`m`$ and type $`I`$.)
In this part, we shall describe the family of moves $`\varphi :_m𝒮_{I,m}`$
We shall work with the terminology introduced in section 2.3 of Chapter 2.3, (see Proposition 2.3.7 of Chapter 2.3.7, and also Theorem 2.3.9 of Chapter 2.3.9 and consider the set of points $`_{iI}a_i=_{n=1}^mA_n`$ where $`A_n`$ denotes the set of points perturbed in a maximal simple of deformation of $`_n`$. Recall that $`A_{2k}=\{a_1,\mathrm{},a_{2k1},\mathrm{},a_{4k2}\}`$ where $`a_1,\mathrm{}a_{2k1}`$ are the crossings of $`_{2k1}L`$ and $`a_{2k},\mathrm{},a_{4k2}`$ are the crossings of $`_{2k;1}`$ where $`_{2k;t}`$, $`t[0,1]`$ is a maximal simple deformation of $`_{2k}`$; $`A_{2k+1}=\{a_1,\mathrm{},a_{2k}\}`$ where $`a_1,\mathrm{}a_{2k}`$ are the crossings of $`_{2k}L`$. In such a way, we assign to the set $`A2k=\{a_1,\mathrm{},a_{2k1},\mathrm{},a_{4k2}\}`$ the global order $`2k`$; and to the set $`A_{2k+1}=\{a_1,\mathrm{},a_{2k}\}`$ the global order $`2k+1`$.
Fold the real point set of $`_m`$ in such a way for any $`1<j<m`$, the two real points of $`_{j+1}`$ resulting from the perturbation of one crossing of $`_jL`$ are glued each other in the initial crossing. Then, fold the curve $`_m`$ in any other points perturbed in its construction; namely, glue the two branches involved in the local topological meaning of any critical point $`a_i`$, $`2ki4k2`$ of Harnack polynomials $`B_{2k}`$, $`2km`$.
The curve $`𝒜_m`$ results from desingularization of all these points with multiplicity $`2`$ compatible with the way $`S`$ goes through the discriminant hypersurface. It is equivalent to unfold in a way different from the one given $`_m`$. Restriction on the unfolding is given by the set of smoothed diffeotopic perestroika experienced along the path $`S`$ from $`𝐑_m`$ to $`𝐑𝒜_m`$. Along the path only triple-point with imaginary branches, strong triple point, inverse and direct self-tangency diffeotopic perestroika are encountered.
Moreover, according to the third part of the proof of the Proposition 5.2.4 of Chapter 5.2.4,
1. desingularization (unfolding) are done in increasing global order on the set of points perturbed in the construction of $`_m`$. Locally, any choice of a desingularization is compatible with one diffeotopic perestroika or lets fix the real components of $`_m`$.
2. Imaginary points are involved in one diffeotopic perestroika $`\pi `$ along the path $`S`$ in such a way that if $`J_\pi (\gamma )=2`$, $`\gamma \{\alpha ,\beta \}`$, then it provides imaginary points involved in -at least one and at most two- next diffeotopic perestroikas $`\pi ^{}`$ along $`S`$ such that $`J_\pi ^{}(\gamma )=2`$.
From the Lemma 5.2.5 of Chapter 5.2.5, any real or imaginary point involved in a diffeotopic perestroika belongs to a neighborhood $`\stackrel{~}{U}(p)U(p)=\rho ^m(D(p)\times U_𝐂^2)`$, $`p𝒫_m`$, and the next perestroikas are provided by move of real branches in $`U(p)`$. Consequently, in the chain of moves, imaginary singular points are smoothed in such a way a the end of the procedure one gets a curve without singular points.
Let us summarize properties of the path $`S:[0,1]𝐑𝒞_m`$ in the following Theorem 5.2.7 of Chapter 5.2.7.
Recall that given the triangle $`T_m`$, the real projective space $`𝐑P^2`$ may be deduced from the square $`T_m^{}`$ made of $`T_m`$ and it symmetric copies $`T_{m,x}=s_x(T_m)`$, $`T_{m,y}=s_y(T_m)`$ $`T_{m,xy}=s(T_m)`$ where $`s_x,s_y,s=s_xs_y`$ are reflections with respect to the coordinates axes. Given a triangle $`T`$ of the set $`T_m`$,$`T_{m,x}`$, $`T_{m,y}`$, $`T_{m,xy}`$ $`s(T)`$ denotes the symmetric copy of $`T`$.
###### Theorem 5.2.7.
Let $`𝒜_m`$ be a curve of degree $`m`$ and type $`I`$. Let $`_m`$ be the Harnack curve of degree $`m`$. Then, up to conj-equivariant isotopy of $`𝐂P^2`$, there exists a path $`S:[0,1]𝐑𝒞_m`$ with $`S(0)=𝐑_m`$, $`S(1)=𝐑𝒜_m`$ which crosses the discriminant hypersurface $`𝐑𝒟_m`$ in curves with crossings, real points of intersection of a real branch and two conjugate imaginary branches, strong triple-points, direct and inverse self-tangency points.
Curves along $`S`$ may be deduced one from another from smoothed diffeoopic perestroikas with the property that any smoothed diffeotopic perestroika along $`S`$ defines double-points involved in next smoothed diffeotopic perestroikas along $`S`$. Furthermore,
1. any real ordinary double point of curves along $`S`$ is a crossing which belongs to the set of points perturbed in the construction of $`_m`$
2. any double-point involved in a diffeotopic perestroika is a crossing which belongs to the set of points perturbed in the whole construction of $`_m`$.
Besides, smoothed diffeotopic perestroikas may be described in the patchworking scheme as follows: Let $`T\{T_m,T_{m,x},T_{m,y},T_{m,xy}\}`$;
1. inverse self tangency diffeotopic perestroika in the positive direction acts as follows: an oval $`𝒪`$ of $`_m`$ contained in a triangle $`T`$ becomes an oval $`s(𝒪)`$ of the symmetric triangle $`s(T)`$, in the negative direction, according to the definition of the self tangency diffeotopic perestroika $`s(𝒪)`$ disappears in two real branches with opposite orientation.
Moreover, in the positive (resp negative) direction, it requires (resp, smooths) imaginary points provided by a triple point with imaginary diffeotopic perestroika in the negative direction.
2. direct self tangency diffeotopic perestroika acts as follows: it does not change the real part but in the negative direction it provides the imaginary points of a triple-point diffeotopic perestroika with imaginary conjugate branches in the positive direction; in the positive direction it smooths the imaginary point of a triple-point diffeotopic perestroika with imaginary conjugate branches in the negative direction.
3. triple-point with imaginary conjugate branches diffeotopic perestroika is as follows: an oval $`𝒪`$ of $`_m`$ contained in a triangle $`T`$ becomes an oval $`s(𝒪)`$ of the symmetric triangle $`s(T)`$.
4. strong triple point perestroika may be experienced only in an open which contains a $`k`$-point (in case $`m`$ even : $`k=3`$; in case $`m`$ odd: $`k=3`$ or $`k=2`$) of $`_m`$.
###### Remark 5.2.8.
1. One can easily deduce the maximal nest curve of degree $`m=2k`$ (which are $`L`$-curves and thus have standard Arnold surfaces) from the Harnack curve $`_{2k}`$.
Indeed, choose an orientation on the recursive line $`𝐑l_{2kj}`$ opposite to the one chosen in the construction of $`_m`$. A slight perturbation of the resulting curves given by standard modification on real double point singularity compatible with the orientation of the resulting curve leads to the maximal nest curve.
Besides, according to the Theorem 5.2.7 of Chapter 5.2.7, there exists a path from $`_{2k}`$ to the maximal nest curves $`<1>^k`$. Such path may be described by smoothed diffeotopic perestroika, some of which (for $`k>2`$) are smoothed strong triple point diffeotopic perestroikas.
2. Obviously, from the method summarized in theorem 5.2.7 of Chapter 5.2.7, one can deduce whether there exist curves of a degree $`m`$ and type $`I`$ with a given real scheme.
For example, the non-existence of the curve of degree $`7`$ and real scheme $`J<1<14>>`$ is easy to deduce. Indeed, it can be easily proved that the existence of such curve is in contradiction with the possible situations inside the three neighborhoods $`U(p)`$, $`p𝒫_7`$ of global order $`6`$ which contain crossings associated to $`2`$-uple. (More generally, one can notice that the $`2k3`$ $`2`$-uples of $`_{2k+1}`$ define possible intersection of the one-side component of $`_{2k+1}`$ with the boundary of the Möbius band of $`𝐑P^2`$ embedded in $`𝐂P^2`$. )
Description of pairs $`(𝐂P^2,𝐂𝒜_m)`$ up to conj-equivariant isotopy We shall now in Theorem 5.2.9 of Chapter 5.2.9 extend the results of Theorem 2.3.9 of Chapter 2.3.9 to any smooth curve of type $`I`$.
According to Theorem 2.3.9 of Chapter 2.3.9 there exists a finite number $`I`$ of $`4`$-balls $`B(a_i)`$ globally invariant by complex conjugation; such that, up to conj-equivariant isotopy of $`𝐂P^2`$ $`_m\backslash _{iI}B(a_i)=_{i=1}^mL_i\backslash _{i=1}^IB(a_i)`$; inside any $`B(a_i)`$ $`_m`$ is the perturbation of crossing centered the in point $`a_i`$. The following statement is a consequence of the definition of $`S:[0,1]𝐑𝒞_m`$ given in Proposition 5.2.4 of Chapter 5.2.4 and the fact that any real singular point which participated in diffeotopic perestroika along $`S`$ belongs to $`_{i=1}^Ia_i`$.
###### Theorem 5.2.9.
Let $`𝒜_m`$ be a curve of degree $`m`$ and type $`I`$. There exists a finite number $`I`$ ($`I=1+2\mathrm{}m+\mathrm{\Sigma }_{k=2}^{k=[m/2]}2k3`$) of disjoints $`4`$-balls $`B(a_i)`$ invariant by complex conjugation and centered in points $`a_i`$ of $`𝐑P^2`$ such that, up to conj-equivariant isotopy of $`𝐂P^2`$ :
1. $`𝒜_m\backslash _{iI}B(a_i)=_{i=1}^mL_i\backslash _{iI}B(a_i)`$ where $`L_1`$,…,$`L_m`$ are $`m`$ projective lines with $`L_i\backslash _{iI}B(a_i)L_j\backslash _{iI}B(a_i)=\mathrm{}`$ for any $`ij`$,$`1i,jm`$.
2. situations inside $`4`$-balls $`B(a_i)`$ are covered by perturbation of type $`1`$ or type $`2`$ of the crossing $`a_i`$.
proof:
Our proof is based on the proof of Proposition 5.2.4 of Chapter 5.2.4, Theorem 5.2.7 of Chapter 5.2.7 and Theorem 2.3.9 of Chapter 2.3.9.
According to Theorem 2.3.9 of Chapter 2.3.9, there exists a finite number $`I`$ of disjoints $`4`$-balls $`B(a_i)`$ invariant by complex conjugation and centered in points $`a_i`$ of $`𝐑P^2`$ such that , up to conj-equivariant isotopy of $`𝐂P^2`$ :
1. $`_m\backslash _{iI}B(a_i)=_{i=1}^mL_i\backslash _{iI}B(a_i)`$ where $`L_1`$,…,$`L_m`$ are $`m`$ projective lines with $`L_i\backslash _{iI}B(a_i)L_j\backslash _{iI}B(a_i)=\mathrm{}`$ for any $`ij`$,$`1i,jm`$.
2. situations inside $`4`$-balls $`B(a_i)`$ are perturbations of type 1 of crossing.
Let $`𝐂_m^+`$ be the half of $`𝐑_m`$ which induces orientation on the real part of $`𝐑_m`$. The conj-equivariant isotopy brings $`𝐂_m^+\backslash _{i=1}^IB(a_i)`$ to halves $`𝐂L_i^+\backslash _{i=1}^IB(a_i)`$ of lines $`L_i\backslash _{i=1}^IB(a_i)`$, $`1im`$, which induce an orientation on the real part $`𝐑L_i`$.
The path $`S:[0,1]𝐑𝒞_{m,I}`$ $`S(0)=𝐂_m`$, $`S(1)=𝐂𝒜_m`$ can be seen as the track on $`𝐑P^2`$ of a diffeotopy $`h_t`$ of $`𝐂P^2`$, $`t[0,1]`$, $`h(0)=𝐂_m`$, $`h(1)=𝐂𝒜_m`$ with $`h_t(𝐂_m)U(p)U(p)`$ and may be described as a family moves defined by smoothed diffeotopic perestroikas.
According to the proof of proposition 5.2.4 of Chapter 5.2.4, any real double point involved in a diffeotopic perestroika belongs to $`_{iI}B(a_i)`$.
Besides, nonetheless relative location of any part of $`_m\backslash _{iI}B(a_i)`$ conj-equivariant isotopic to a part $`L_i\backslash _{iI}B(a_i)`$ is changed under a diffeotopic perestroika, its orientation remains the same before and after diffeotopic perestroika.
Therefore, for any line $`L_i`$, the half $`𝐂L_i^+\backslash _{iI}B(a_i)`$ of $`L_i\backslash _{iI}B(a_i)`$ which induces orientation on $`𝐑L_i\backslash _{iI}B(a_i)`$ is up to conj-equivariant isotopy a part of the half $`𝐂𝒜_m^+`$ of $`𝐂𝒜_m`$ which induces orientation on $`𝐑𝒜_m`$.
Hence, outside $`_{i=1}^IB(a_i)`$, the curve $`𝒜_m`$ is union of $`m`$ lines minus their intersection with $`_{i=1}^IB(a_i)`$.
Moreover, situations inside $`4`$-balls $`B(a_i)`$ are covered by two cases which are locally perturbations of type $`1`$ or type $`2`$ of a crossing.
Q.E.D
#### 5.3. Construction of Curves of type $`II`$
As already known, if an $`M`$-curve is constructed curves with fewer components are easily constructed. Besides, $`M`$-curves, as curves of type $`I`$, have been studied in the section 5.2 of Chapter 5.2. We shall now introduce a method inspired from the preceding one, which will give all $`Mi`$-curves of degree $`m`$ and therefore any curves of type $`II`$ from the Harnack curve $`_m`$. In the previous section, it was essential that all curves have orientable real point set. In order to consider curves with non-orientable real set of points we shall modify the argument introduced in the preceding part. We shall call perturbation in the non-coherent direction of a crossing of a curve the desingularization which induces non-orientability of the real part of the resulting curve. For a generic curve $`𝒜_m`$ of degree $`m`$ and genre $`g`$, $`0<g<\frac{(m1)(m2)}{2}`$, by smoothing in the coherent direction of one double point of its real point set $`𝐑𝒜_m`$, we shall understand the Morse modification in $`𝐑^2`$ in the direction coherent to a complex orientation of $`𝒜_m`$, by smoothing in the non-coherent direction of one double point of its real point set $`𝐑𝒜_m`$, we shall understand the Morse modification in $`𝐑^2`$ at the double point in the direction non-coherent to a complex orientation of $`𝒜_m`$, i.e not compatible with an orientation of the real point set of the resulting curve. Let us start our study of curves of type $`II`$ by a Lemma analogous to the Lemma 5.2.1 of Chapter 5.2.1 stated for curves of type $`I`$.
Call smoothed generic immersion of the circle $`S^1`$ into the plane $`𝐑^2`$, the smooth submanifold of $`𝐑P^2`$ deduced from the generic immersion of the circle by modification at each real double point which is either a Morse modification in $`𝐑^2`$ coherent or non-coherent to the complex orientation, or the Morse modification in $`𝐑P^2`$ which associates to the double point of $`𝐑^2`$ two points of the line at infinity of $`𝐑P^2`$.
A generalization of Whitney’s theorem to the case of real algebraic curves will be provided by the following Lemma:
###### Lemma 5.3.1.
Let $`𝒜_m`$ be a smooth curve of degree $`m`$ and type $`II`$ with non-empty real points set $`𝐑𝒜_m`$ then $`𝐑𝒜_m𝐑P^2`$ is a smoothed generic immersion of the circle $`S^1`$ into the plane $`𝐑^2`$ with
$`\frac{(m1)(m2)}{2}n\frac{(m1)(m2)}{2}+[\frac{m}{2}]`$, double points and Whitney index $`\frac{(m1)(m2)}{2}+1`$. The smoothing is such that $`\frac{(m1)(m2)}{2}`$ double points are smoothed in $`𝐑^2`$ by Morse modification and at least one Morse modification is in the non coherent direction.
###### Remark 5.3.2.
One can associate to any smooth curve of type $`II`$ a three-dimensional rooted tree.
proof: Our proof makes use of properties of the complex point set $`𝐂𝒜_m`$ embedded in $`𝐂P^2`$. It is deduced from an argument similar to the one given for curves of type $`I`$.
Consider the usual handlebody decomposition of $`𝐂P^2=B_0B_1B`$ where $`B_0`$,$`B_1`$,$`B`$ are respectively 0,2 and 4 handles.
The ball $`B_0`$ and $`B_1`$ meet along an unknotted solid torus $`S^1\times B^2`$. The gluing diffeomorphism $`S^1\times B^2S^1\times B^2`$ is given by the $`+1`$ framing map. In such a way, the canonical $`𝐑P^2`$ can be seen as the union of a Möbius band $``$ and the disc $`D^2B`$ glued along their boundary. The Möbius band $``$ lies in $`B_0B_1S^1\times D^2`$ with $``$ as the $`(2,1)`$ torus knot, and $`D^2B`$ as the properly imbedded unknotted disc. The complex conjugation switches $`B_0`$ and $`B_1`$ and lets fix $``$, it rotates $`B`$ around $`D^2`$.
The set of complex points of $`𝒜_m`$ is an orientable surface of genus $`g_m=\frac{(m1)(m2)}{2}`$, i.e it is diffeomorphic to a sphere $`S^2`$ with $`g_m`$ $`1`$-handles $`S_{g_m}^2`$. We shall denote $`h`$ the diffeomorphism $`h:𝐂𝒜_mS_{g_m}^2𝐑^3`$.
Assume $`S^2`$ provided with a complex conjugation with fixed point set a circle $`S^1`$ which divides the sphere $`S^2`$ into two halves.
Without loss of generality, one can assume that any disc removed from $`S^2`$ and then closed up to give $`S_{g_m}^2`$ intersects $`S^1`$ and the two halves of $`S^2`$. In such a way, $`𝐑𝒜_m𝐑P^2`$ intersects each $`1`$-handle.
Fix $`D^2`$ the two disc of $`𝐑P^2=D^2`$ in such a way that the boundary circle of $`D^2`$ is $`S^1`$ and therefore each one handle belongs to the solid torus $`S^1\times D^2`$.
Since the interior $`(D^2)^0`$ of $`D^2`$ is homeomorphic to $`𝐑^2`$, one can project $`𝐑𝒜_m`$ (up to homeomorphism $`𝐑^2(D^2)^0`$) in a direction perpendicular to $`𝐑^2`$ onto $`𝐑^2`$. We may suppose that the direction of the projection is generic i.e all points of self-intersection of the image on $`𝐑^2`$ are double and the angles of intersection are non-zero.
Let $`r:𝐑^3𝐑^2`$ be the projection which maps $`h(𝐑𝒜_m)𝐑^3`$ to $`𝐑^2`$. Consider an oriented tubular fibration $`N𝐂𝒜_m`$. Since $`𝐂𝒜_m`$ is diffeomorphic to a sphere $`S^2`$ with $`g_m`$ $`1`$-handles, one can consider the restriction of $`𝐂𝒜_m`$ diffeomorphic to the torus $`T^2`$ given by the sphere $`S^2`$ with one of $`g_m`$ $`1`$-handle.
The oriented tubular neighborhood of $`h(𝐂𝒜_m)T^2`$, as oriented tubular neighborhood of the torus $`T^2`$, intersects the solid torus
$`S^1\times D^2`$ with $`MS^1\times S^1`$ as the $`(2,1)`$ torus knot.
Hence, since in $`𝐂P^2`$, each real line is split by its real part into two halves lines conjugate to each others, and $`2`$ disjoint circles always divide the torus, one can assume, without loss of generality, that the real part of the restriction of $`h(𝐂𝒜_m)T^2`$ belongs to the boundary of the Möbius band in such a way that its projection to $`𝐑^2`$ gives one crossing.
Besides, some double-points of $`r(h(𝐑_m))`$ may result either from two points of $`𝐑𝒜_m`$ which belong to two different handles or one point which belongs to a $`1`$-handle and the other point belongs to $`S^2`$. It easy to see that, in both cases, the projection leads to an even number of such double-points.
Hence, from an argumentation similar to the one given in case of curves of type $`I`$, it follows that $`𝐑𝒜_m`$ is the smoothed immersion of a generic immersion of the circle $`S^1`$ into the plane $`𝐑^2`$ with $`n`$, $`\frac{(m1)(m2)}{2}n\frac{(m1)(m2)}{2}+[\frac{m}{2}]`$, double points and Whitney index $`\frac{(m1)(m2)}{2}+1`$.
Besides, it is obvious, since $`𝒜_m`$ is of type $`II`$, that at least one-double point of $`r(h(𝐑\stackrel{~}{𝒜}_m))`$ is smoothed in the non-coherent direction (in other words, at least two circles of the $`g_m+1`$ which divide the surface $`S_{g_m}^2`$ glue and disappear in one circle.) Q.E.D.
Given a generic immersion of the circle $`S^1`$ into the plane $`𝐑^2`$ and $`𝒮`$ the set of its singular points. We shall call partially smoothed generic immersion of the circle $`S^1`$ into the plane $`𝐑^2`$ the singular real part deduced from the generic immersion of the circle by modification of a set $`𝒦𝒮`$, $`𝒦𝒮`$ of real double points. The following Lemma, analogue to the Lemma 5.2.3 of Chapter 5.2.3 of the previous part, enlarges the preceding statement to singular curves.
###### Lemma 5.3.3.
Let $`𝒜_m`$ be a singular curve of degree $`m`$ and type $`II`$ with non-empty real points set $`𝐑𝒜_m`$ and non-degenerate singular points, then its real point set $`𝐑𝒜_m𝐑P^2`$ is a partially smoothed generic immersion of the circle $`S^1`$ into the plane $`𝐑^2`$ with $`n\frac{(m1)(m2)}{2}`$ double points and Whitney index $`\frac{(m1)(m2)}{2}+1`$ if and only if its set of singular points consists of at most $`\frac{(m1)(m2)}{2}`$ crossings and at least one crossing is smoothed in $`𝐑^2`$ in the non-coherent direction.
proof: It follows from an argument similar to the one used in the proof of Lemma 5.3.1 of Chapter 5.3.1. Q.E.D Although the stratification of the set singular curves in the variety $`𝒞_{I;g}`$ may not be extended to curves of type $`II`$, we shall lift the path provided by Whitney’s theorem and Lemma 5.3.1 of Chapter 5.3.1 to algebraic curves of type $`II`$.
As previously, we shall consider the space $`𝐑𝒞_m`$ of real algebraic curves of degree $`m`$ and its subset $`𝐑𝒟_m`$ constituted by real singular algebraic curves of degree $`m`$. The set $`𝐑𝒟_m`$ has an open every dense subset which consist of curves with only one singular point which is a non-degenerate double point (i.e a solitary real double-point, or a crossing). This subset is called principal part of the set $`𝐑𝒟_m`$. A generic path in $`𝐑𝒞_m`$ intersects $`𝐑𝒟_m`$ only in its principal part and only transversally. When a singular curve occurs in a generic one-parameter family curves, the moving curve is passing through a Morse modification. Here are simple properties of Morse modifications which motivate our study:
1. Under a Morse modification a curve of type $`I`$ can turn only to a curve of type $`II`$ and the number of its real components decreases.
2. A complex orientation of a non-singular curve turns into an orientation of the singular curve which appears at the moment of the modification: the orientations defined by complex orientation of the curve on the two arcs which approach each other and merge should be coherent with an orientation of the singular curve.
We shall generalize the previous method giving any smooth curve $`𝒜_m`$ of type $`I`$ to smooth curves of type $`II`$ and deduce curves of type $`II`$ from smoothing (including smoothing in a direction non-coherent to a complex orientation) generic curves of type $`I`$, genus $`g`$ and type $`I`$.
Assume $`_m`$ the Harnack of degree $`m`$ obtained via the patchworking method. For any point $`p𝒫_m`$, the subset $`U(p)=\rho ^m(D(p,ϵ)\times U_𝐂^2)(𝐂^{})^2`$ intersects $`(𝐑^{})^2`$ in four discs. We shall denote by $`_m`$ the set constituted by the union of the four $`2`$-discs of $`U(p)(𝐑^{})^2`$ taken over the points $`p𝒫_m`$.
###### Proposition 5.3.4.
Let $`_m`$ be the Harnack curve of degree $`m`$ and $`𝒜_m`$ be a smooth curve of type $`II`$. Then, up to conj-equivariant isotopy of $`𝐂P^2`$, there exists a path
$$S:[0,1]𝐑𝒞_m$$
$`S(0)=𝐑_m`$, $`S(1)=𝐑𝒜_m`$ which intersects the discriminant hypersurface $`𝐑𝒟_m`$ in such a way that:
1. any real ordinary double point of curves along $`S`$ is a crossing which belongs to a real $`2`$-disc of $`_m`$.
2. the curve $`𝒜_m`$ is deduced from smoothing the real part of a generic curve of degree $`m`$ of which double points are crossings. At least one point is smoothed in a direction non-coherent with a complex orientation of the real part of the curve.
proof: As already done in case of curves of type $`I`$, using the path provided by Lemma 5.3.1 of Chapter 5.3.1 and Whitney’s theorem, according to the Lemma 5.3.3 of Chapter 5.3.3, and the fact that $`𝐂_m`$ and $`𝐂𝒜_m`$ are diffeotopic, we shall characterize the track $`S`$ on $`𝐑P^2`$ of a diffeotopy $`h_t`$ of $`𝐂P^2`$, $`t[0,1]`$, $`h(0)=𝐂_m`$, $`h(1)=𝐂𝒜_m`$.
Method which provides $`(𝐂P^2,𝐂𝒜_m)`$ from $`(𝐂P^2,𝐂_m)`$.
Our method is a slightly modified version of the preceding one where any curve $`𝒜_m`$ of type $`I`$ is obtained from $`_m`$. We shall prove that up to modify the coefficients of the polynomial giving the curve $`𝒜_m`$, one can always assume that there exists a diffeotopy $`h_t`$ of $`𝐂P^2`$ $`h(0)=𝐂_m`$, $`h(1)=𝐂𝒜_m`$ with the property $`h_t(𝐂𝒜_mU(p))U(p)`$ for any $`p𝒫_m`$. We shall refer to the method of section 5.2 of Chapter 5.2 and stress only the modification. According to the Lemma 5.3.1 of Chapter 5.3.1, consider $`𝐑_m`$ and $`𝐑𝒜_m`$ as smoothed immersions of the circle into the plane with the same Whitney index. Denote $`𝐑\stackrel{~}{}_m`$, $`𝐑\stackrel{~}{𝒜}_m`$ the corresponding immersions. From the Whitney’s theorem (), there exists a path $`\stackrel{~}{h}`$ which connects $`𝐑\stackrel{~}{}_m`$ and $`𝐑\stackrel{~}{𝒜}_m`$. Thus, the definition of the path $`S`$ is reduced to the definition of a lifting of the path which connects $`𝐑\stackrel{~}{}_m`$ and $`𝐑\stackrel{~}{𝒜}_m`$ in the space of immersions of the circle into $`𝐑^2`$ to a path in the space of curve of degree $`m`$ and type $`II`$.
Define now the smoothed perestroika of a perestroika, the change obtained by smoothing the fragments involved in the perestroika where smoothing before perestroika are taken only in a direction coherent and smoothing after perestroika may be taken in a direction non-coherent.
Then, as in the proof of proposition 5.2.4 of Chapter 5.2.4 of section 5.2 of Chapter 5.2 where curves of type $`I`$ were under consideration, since any smooth curve is irreducible, and any reducible polynomial is the product of a finite reducible one, according to Lemma 5.3.1 of Chapter 5.3.1 and Lemma 5.3.3 of Chapter 5.3.3, we may lift the path $`\stackrel{~}{h}`$ in the space of immersion of the circle into the plane from smoothed perestroikas in the space $`𝒞_{I;g}`$ of curve of degree $`nm`$, type $`I`$ and genus $`0g\frac{(n1).(n2)}{2}`$.
In such a way, using the path provided by Lemma 5.3.1 of Chapter 5.3.1 and Whitney’s theorem, according to the Lemma 5.3.3 of Chapter 5.3.3 and the fact that $`𝐂_m`$ and $`𝐂𝒜_m`$ are diffeotopic, we characterize the track $`S`$ on $`𝐑P^2`$ of a diffeotopy $`h_t`$ of $`𝐂P^2`$, $`t[0,1]`$, $`h(0)=𝐂_m`$, $`h(1)=𝐂𝒜_m`$. Q.E.D
###### Remark 5.3.5.
As already known, if two $`M2`$-curves are obtained one from the other by a deformation through a double nondegenerate point, one of them is of type $`I`$ and the other is of type $`II`$. Since, any $`M1`$-curve is of type $`II`$, according to our method it remains to say that two deformations in the non-coherent direction may turn a curve of type $`I`$ to a curve of type $`I`$.
As a curve of degree $`m`$ moves as a point of $`𝐑𝒞_m`$ along an arc which has a transversal intersection with the principal part, then the set of real points of this curve undergoes either a Morse modification of index $`0`$ or $`2`$ (the curves acquires a solitary singular point, which then becomes a new oval, or else one of the ovals contracts to point) or Morse modification of index $`0`$, an oval contracts to a point or a Morse modification of index $`1`$ ( two real arcs of the curve approach one other and merge, then diverge in their modified form.)
As already introduced, we denote $`A_{m+1}`$ the set of crossings perturbed in the construction of $`_{m+1}`$ from $`_m`$ introduced in proposition 2.3.1 of Chapter 2.3.1.
In case $`m+1=2k`$, $`A_{2k}=\{a_1,\mathrm{},a_{2k1},\mathrm{},a_{4k2}\}`$ where $`a_1,\mathrm{}a_{2k1}`$ are crossings of $`_{2k1}L`$ and $`a_{2k},\mathrm{},a_{4k2}`$ are crossings of $`\stackrel{~}{X}_{2k;\tau }`$.
In case $`m+1=2k+1`$, $`A_{2k+1}=\{a_1,\mathrm{},a_{2k}\}`$ where $`a_1,\mathrm{}a_{2k}`$ are crossings of $`_{2k}L`$.
We may characterize Morse modifications on a curves of degree $`m`$ as follows:
###### Corollary 5.3.6.
Any Morse modification on a curve of degree $`m`$ may be described, up to conj-equivariant isotopy, by perturbations of crossings of a subset of $`_{n=1}^mA_m`$
proof: It is a straightforward consequence of proposition 5.3.4 of Chapter 5.3.4. Q.E.D.
Description of pairs $`(𝐂P^2,𝐂𝒜_m)`$ up to conj-equivariant isotopy
###### Theorem 5.3.7.
Let $`𝒜_m`$ be a curve of degree $`m`$ and type $`II`$ with non-empty real part. There exists a finite number $`I`$ of disjoints $`4`$-balls $`B(a_i)`$ invariant by complex conjugation and centered in points $`a_i`$ of $`𝐑P^2`$ such that, up to conj-equivariant isotopy of $`𝐂P^2`$ :
1. $`𝒜_m\backslash _{iI}B(a_i)=_{i=1}^mL_i\backslash _{iI}B(a_i)`$ where $`L_1`$,…,$`L_m`$ are $`m`$ projective lines with $`L_i\backslash _{iI}B(a_i)L_j\backslash _{iI}B(a_i)=\mathrm{}`$ for any $`ij`$,$`1i,jm`$.
2. situations inside $`4`$-balls $`B(a_i)`$ are covered by perturbation of type $`1`$ or type $`2`$ of the crossing $`a_i`$.
proof:
Our proof is based on the previous method and the theorem 2.3.9 of Chapter 2.3.9. The argument is analogous to the one of the proof of theorem 5.2.9 of Chapter 5.2.9 where curves of type $`I`$ were under consideration. Q.E.D
### Chapter 6 Arnold surfaces of curves of even degree <br>-Proof of the Rokhlin’s conjecture-
In this section, we shall prove the Rokhlin’s conjecture: Theorem (Rokhlin’s conjecture)
Arnold surfaces $`𝔄`$ are standard for all curves of even degree with non-empty real part.
proof: The proof is based on the Livingston’s theorem and the following statement (deduced from theorem 5.2.9 of Chapter 5.2.9 and theorem 5.3.7 of Chapter 5.3.7 )
Given $`𝒜_m`$ a curve of degree $`m`$ and type $`I`$ or type $`II`$ with non-empty real part. There exists a finite number $`I`$ of disjoints $`4`$-balls $`B(a_i)`$ invariant by complex conjugation and centered in points $`a_i`$ of $`𝐑P^2`$ such that up to conj-equivariant isotopy of $`𝐂P^2`$:
1. $`𝒜_m\backslash _{iI}B(a_i)=_{i=1}^mL_i\backslash _{iI}B(a_i)`$ where $`L_1`$,…,$`L_m`$ are $`m`$ projective lines with $`L_i\backslash _{iI}B(a_i)L_j\backslash _{iI}B(a_i)=\mathrm{}`$ for any $`ij`$,$`1i,jm`$.
2. situations inside any $`4`$-balls $`B(a_i)`$ are covered by perturbation of type $`1`$ or type $`2`$ of the crossing $`a_i`$.
It follows from an argumentation similar to the one given in Chapter 3 (see theorem 3.0.1 of Chapter 3.0.1 where it is stated that Arnold surfaces of Harnack curves $`_{2k}`$ are standard surfaces in $`S^4`$) that any Arnold surface of a curve of even degree $`2k`$ with non-empty real part is standard.
Q.E.D
|
warning/0003/astro-ph0003484.html
|
ar5iv
|
text
|
# Propagation of ultra-high energy protons in the nearby universe
## I Introduction
The world statistics of ultra high energy cosmic ray (UHECR) events of energy above 10<sup>20</sup> eV has now grown to 20 events . It is very difficult to accelerate particles to such high energies in astrophysical shocks, the process thought to be responsible for the majority of the galactic cosmic rays . This has led to a large number of production models, many of them based on exotic particle physics scenarios . The gyroradii of 10<sup>20</sup> eV protons are significantly larger than our own Galaxy and this suggests an extragalactic origin for any astrophysical scenario ($`r_g=100\mathrm{k}\mathrm{p}\mathrm{c}\times (E/10^{20}\mathrm{eV})\times (1\mu \mathrm{G}/B)`$ with $`E`$ and $`B`$ being the proton energy and the magnetic field strength, respectively). The large distances between potential UHECR sources and Earth leads to another set of problems first pointed out independently by Greisen and by Zatsepin & Kuzmin, now widely known as the GZK effect . UHECR protons interact with photons of the microwave background radiation and lose their energy relatively rapidly during propagation over distances of tens of megaparsecs. This should result in a cutoff in the cosmic ray spectrum at an energy just below 10<sup>20</sup> eV.
Many different calculations , performed using various techniques, of the modification of the cosmic ray spectrum due to propagation have been published since the original suggestion. As a result, the general features of the cosmic ray spectrum after propagation are well established. Differences between the various approaches are, however, significant and the accuracy achieved is not sufficient for the interpretation of the existing experimental data, and more accurate calculations are needed for the expected significant increase of the experimental statistics .
Previous calculations can be divided into two classes dealing mainly with: (a) the energy loss processes , and (b) the deflection and scattering of protons in the extragalactic magnetic field . The first group of calculations shows that small differences in the realization of the proton energy loss processes generate observable differences in the predicted spectra at Earth. Such calculations, however, cannot establish an accurate relation between the distance of a potential source and the modification of the proton spectrum emitted by this source because the influence of the extragalactic magnetic field is neglected. Among the calculations of the second kind, Refs. do not consider the proton energy losses in a satisfactory way, and Refs. mostly discuss their results in a specific context. Only Achterberg et al. give a detailed discussion of the fundamental aspects of UHECR propagation in extragalactic magnetic fields, which we are interested in here.
We present here calculations performed with the photoproduction event generator SOPHIA , which is proven to reproduce well the cross section and final state composition in nucleon-photon interactions for energies from the particle production threshold up to hundreds of GeV in the center-of-mass system. We also account for all other energy loss processes of UHECR nucleons, and calculate the proton deflection in the extragalactic magnetic field in three dimensions.
We restrict ourselves to proton injection energies up to 10<sup>22</sup> eV, and consider (with few exceptions) proton propagation for source distances less than 200 Mpc. The calculations are carried out using a Monte Carlo technique, and we propagate individual protons injected as either a mono-energetic beam, or with energies sampled from a fixed source energy spectrum. This approach has the advantage of representing fluctuations in the proton energy losses very well, thereby giving us a good handle on the correlations between energy loss, time of flight and angular deviation of the flight direction. As we will show, these important UHECR characteristics are deeply interconnected. For a given source distance, there is a strong correlation between the amount of energy lost, the time delay, and the scattering angle.
Our calculations are thus mainly relevant to scenarios of UHECR acceleration at astrophysical shocks, for which 10<sup>22</sup> eV is a very generous upper energy limit. With this paper we wish to establish limits for the distance of potential UHECR proton sources as a function of proton energy and the average strength of the extragalactic magnetic field. We also study the angular distribution of UHECR with respect to the source direction (arrival angle) and the time delays after propagation over different distances. In addition, the neutrino fluxes produced during the propagation are presented.
The article is organized as follows. We describe the propagation method, including the relevant features of the event generator SOPHIA, in Section 2. Section 3 gives some interesting results on the propagation of mono-energetic proton beams, and compares our results with other work. Section 4 analyzes the formation and development of the primary and secondary particle spectra for protons injected with a power law spectrum. In section 5 we discuss the results, present our conclusions, and make suggestions for future work.
## II Cosmic Ray Propagation
This section provides a description of our simulation code for propagating protons in intergalactic space. We treat energy losses due to hadronic and electromagnetic interactions of the nucleons with photons of the cosmic microwave background radiation as well as the deflection of particles by the intergalactic magnetic field. Although we present here only results on nucleon propagation in random magnetic fields, our approach also allows us to follow the particles in complicated magnetic field topologies. Because of the time-consuming detailed simulation of each nucleon propagation path by Monte Carlo, the propagation method described below is not suitable for calculations involving large cosmological distances.
### A Interactions and energy loss processes
Particles of energy $`E>10^{18}`$ eV interact with photons of the cosmic microwave background radiation giving rise to secondary particle production and nucleon energy loss. The most important processes are:
* photoproduction of hadrons, and
* Bethe-Heitler (BH) production of $`e^+e^{}`$ pairs by protons.
We also account for the adiabatic losses due to cosmological expansion of the Universe, and for the decay of neutrons produced in hadronic production process. Since we restrict our calculation to models of UHECR acceleration in astrophysical shocks, and energies below 10<sup>22</sup> eV, we consider only interactions with cosmic microwave background photons. The calculation of nucleon propagation at higher energies would require the use of models of the radio background (see e.g. Ref. ). Since we are not presenting results on the development of electromagnetic cascades initiated by secondary particles produced in proton-photon interactions, we can safely neglect interactions on the universal optical/infrared background as well. We keep track, however, of the individual energies of all secondaries of photoproduction interactions and are thus able to show the spectra of neutrinos generated by primary protons after propagation over different distances.
Hadron production and energy loss in nucleon-photon interactions is simulated with the event generator SOPHIA . This event generator samples collisions of nucleons with photons from isotropic thermal or power law energy distributions, using standard Monte Carlo techniques. In this paper the code has been used with a blackbody spectrum with $`T=2.726`$ K to represent the cosmic microwave background. According to the respective partial cross sections, which have been parametrized using all available accelerator data, it invokes an interaction either via baryon resonance excitation, one-particle $`t`$-channel exchange (direct one-particle production), diffractive particle production and (non-diffractive) multiparticle production using string fragmentation. The distribution and momenta of the final state particles are calculated from their branching ratios and interaction kinematics in the center-of-mass frame, and the particle energies and angles in the lab. frame are calculated by Lorentz transformations. The decay of all unstable particles except for neutrons is treated subsequently using standard Monte Carlo methods of particle decay according to the available phase space. The neutron decay is implemented separately into the present propagation code. The SOPHIA event generator has been tested and shown to be in good agreement with available accelerator data. A detailed description of the code including the sampling methods, the interaction physics used, and the performed tests can be found in Ref. .
The Monte Carlo treatment of photoproduction is very important, because nucleons lose a large fraction of their energy in each interaction. As early as 1985 Hill & Schramm pointed out that the use of a continuous energy loss approximation for this process neglects the intrinsic spread of arrival energies due to the variation of the energy loss $`\mathrm{\Delta }E`$ per interaction, and the Poissonian distribution in the number of pion production interactions during propagation. This results in a certain “survival probability” of cosmic rays arriving at Earth with energies above the GZK-cutoff, as estimated in the assumption of continuous energy loss.
Fig. 1a shows the energy dependence of all parameters relevant to the average proton energy loss in the microwave background (T=2.726 K) for redshift $`z`$ = 0. The photoproduction interaction length $`\lambda _{\mathrm{ph}}`$ for protons is shown as a dashed line. Denoting the proton-photon center-of-mass energy by $`\sqrt{s}`$, the interaction length can be written as
$`{\displaystyle \frac{1}{\lambda _{\mathrm{ph}}(E)}}`$ $`=`$ (2)
$`{\displaystyle \frac{1}{8E^2\beta }}{\displaystyle _{ϵ_{\mathrm{th}}}^{\mathrm{}}}𝑑ϵ{\displaystyle \frac{n(ϵ)}{ϵ^2}}{\displaystyle _{s_{\mathrm{min}}}^{s_{\mathrm{max}}(ϵ,E)}}𝑑s(sm_p^2c^4)\sigma _{p\gamma }(s)`$
with
$`s_{\mathrm{min}}`$ $`=`$ $`(m_pc^2+m_{\pi ^0}c^2)^2`$ (3)
$`s_{\mathrm{max}}(ϵ,E)`$ $`=`$ $`m_p^2c^4+2Eϵ(1+\beta )`$ (4)
$`ϵ_{\mathrm{th}}`$ $`=`$ $`{\displaystyle \frac{s_{\mathrm{min}}m_p^2c^4}{2E(1+\beta )}},\beta ^2=1{\displaystyle \frac{m_p^2c^4}{E^2}}.`$ (5)
Here $`E`$ ($`ϵ`$) is the proton (photon) energy and the proton and neutral pion masses are $`m_p`$ and $`m_{\pi ^0}`$, respectively. The CMB photon density is given by $`n(ϵ)`$ in units of cm<sup>-3</sup> eV<sup>-1</sup> and the photoproduction cross section, $`\sigma _{p\gamma }(s)`$, is taken from the parametrization implemented in SOPHIA.
The mean energy loss distance $`x_{\mathrm{loss}}(E)`$, shown in Fig. 1a as triple-dot-dashed curve, is calculated as
$$x_{\mathrm{loss}}(E)=\frac{E}{dE/dx}=\frac{\lambda (E)}{\kappa (E)}$$
(6)
with $`\kappa (E)`$ being the mean inelasticity
$$\kappa (E)=\frac{\mathrm{\Delta }E}{E}.$$
(7)
The mean energy loss of the nucleon due to the hadron production, $`\mathrm{\Delta }E`$, has been calculated by simulating $`10^4`$ interactions for each given proton energy, resulting in a statistical error of the order of 1%. For $`E>10^{20}`$ eV losses through photomeson production dominate with a loss distance of about $`15`$ Mpc at $`E8\times 10^{20}`$ eV. Below this energy, Bethe-Heitler pair production and adiabatic losses due to the cosmological expansion in the Hubble flow determine the proton energy losses.
Both the photoproduction interaction and the pair production are characterized by strongly energy dependent cross sections and threshold effects. Fig. 1a shows $`\lambda _{\mathrm{ph}}`$ decreasing by more than three orders of magnitude for a proton energy increasing by a factor of three. After the minimum $`\lambda _{\mathrm{ph}}`$ is reached, the proton energy loss distance is approximately constant. It is worth noting that the threshold region of $`\lambda _{\mathrm{ph}}`$ is very important for the shape of the propagated proton spectrum. As pointed out by Berezinsky & Grigoreva , a pile–up of protons will be formed at the intersection of the photoproduction and pair production energy loss distances. Another, smaller pile–up will develop at the intersection of the pair production and adiabatic loss functions.
In the current calculation we treat pair production as a continuous loss process which is justified considering its small inelasticity of $`2m_e/m_p10^3`$ (with $`m_e`$,$`m_p`$ being the electron and proton masses, respectively) compared to pion-photoproduction ($`\kappa 0.20.5`$). We use the analytical fit functions given by Chodorowsky et al. to calculate the mean energy loss distance for Bethe-Heitler pair production. This result is in excellent agreement with results obtained by simulating this process via Monte Carlo as done by Protheroe & Johnson .
The turning point from pion production loss dominance to pair production loss dominance lies at $`E6\times 10^{19}`$ eV, with a mean energy loss distance of $`1`$ Gpc. The minimum of the pair production loss length is reached at $`E(24)\times 10^{19}`$ eV. For $`E(23)\times 10^{18}`$ eV continuous losses due to the expansion of the universe dominate. For an Einstein-de Sitter (flat, matter-dominated) universe as considered here, the cosmological energy loss distance scales with redshift $`z`$ as
$$x_{\mathrm{loss},\mathrm{ad}}(E,z)=\frac{c}{H_0}(1+z)^{3/2}4000\mathrm{Mpc}(1+z)^{3/2},$$
(8)
for a Hubble constant of $`H_0=75`$ km/s/Mpc, which we use throughout this paper. All other energy loss distances, $`x_{\mathrm{loss},\mathrm{BH}}`$ for Bethe-Heitler pair production and $`x_{\mathrm{loss},\mathrm{ph}}`$ for photomeson production, scale as
$$x_{\mathrm{loss}}(E,z)=(1+z)^3x_{\mathrm{loss}}[(1+z)E,z=0].$$
(9)
We also show the mean decay distance of $`9\times 10^9\gamma _n`$ kpc for neutrons, where $`\gamma _n`$ is the Lorentz factor of the neutron. Obviously, neutrons of energy below 10<sup>21</sup> eV tend to decay, whereas at higher energies neutrons tend to interact.
Since the details of the proton energy loss directly affect the proton spectra after propagation, we present the ratio of the loss distance in previous calculations to that of our work on a linear scale in Fig. 1b. Generally all values of the energy loss distance are in a good qualitative agreement. Rachen & Biermann treat both Bethe-Heitler and pion production losses very similarly to our work except for the threshold region of pion production. In the pair production region our work is also in perfect agreement with Protheroe & Johnson . An overestimate of the loss distance due to pion production of $`1020\%`$ in Ref. , however, will result in a small shift of the GZK cutoff to higher energies in comparison to the present calculations. Berezinsky & Grigoreva used a very good approximation for the pion production losses, but underestimate the energy loss in pair production interactions by at least 30-40%. The largest deviation of the combined loss distance from our model appears in the calculations of Yoshida & Teshima . As already pointed out in Ref. the largest difference occurs at $`5\times 10^{19}`$ eV where Ref. underestimates pair production losses and uses $`x_{\mathrm{loss}}`$ values larger by about 60%, while photoproduction losses are overestimated by up to 50%. In the work of Lee pion as well as pair production losses are treated in fair agreement with our work, with differences up to 40% in the threshold region of pion production, and 10-20% otherwise. The energy loss code of Lee was also used by Sigl and collaborators . The simple analytical estimate of photoproduction losses in the recent work of Achterberg et al. underestimates the photoproduction loss distance by $`1040\%`$, while $`x_{\mathrm{loss}}`$ due to pair production losses is overestimated by about $`20\%`$.
### B Method of particle propagation
UHECR propagation involves two main distance scales: (a) the hadronic interaction length $`\lambda _{\mathrm{ph}}`$ of typically 3 to 7 Mpc, and (b) the much smaller length scale $`\mathrm{}_{\mathrm{mag}}`$ of typically 10 kpc needed for a precise numerical integration of the equations of motion in a random magnetic field. A straightforward Monte Carlo treatment of the propagation using a step size of $`\mathrm{}_{\mathrm{mag}}`$ for both hadronic interactions and the equations of motion leads to severe efficiency problems for total propagation distances of hundreds of Mpc. Hence, the Monte Carlo simulation is done in the following way. First the path length $`X_{\mathrm{dist}}`$ from the current particle position to the next possible hadronic interaction is determined from
$$X_{\mathrm{dist}}=\lambda _{\mathrm{ph},\mathrm{min}}\mathrm{ln}(\xi ),$$
(10)
where $`\lambda _{\mathrm{ph},\mathrm{min}}`$ is the minimum interaction length for hadronic interactions (at maximum redshift possible for a given total propagation distance) and $`\xi `$ is a random number uniformly distributed in $`(0,1]`$. The nucleon is then propagated over the path length $`X_{\mathrm{dist}}`$ in steps of $`\mathrm{}_{\mathrm{mag}}`$, and for charged particles Bethe–Heitler losses are taken into account and the deflection angle is calculated. A hadronic interaction is then simulated with the probability $`\lambda _{\mathrm{ph},\mathrm{min}}/\lambda _{\mathrm{ph}}(E,z)`$, $`\lambda _{\mathrm{ph}}(E,z)`$ being the interaction length for the energy $`E`$ and redshift $`z`$. It is shown in Appendix A that this method corresponds exactly to a propagation simulation using Eq. (10) with $`\lambda _{\mathrm{ph}}(E,z)`$ for the calculation of the interaction distance at each step with the length $`\mathrm{}_{\mathrm{mag}}`$.
The reduction of the proton energy due to BH pair production and of all nucleons due to adiabatic expansion is calculated at every propagation step, whereas the corresponding loss lengths are updated after a simulated path length of $`\lambda _{\mathrm{ph},\mathrm{min}}`$ and every photoproduction interaction. In the case of neutrons the decay path length is sampled using Eq. (10) with the neutron decay length. The smaller of both the hadronic interaction and the decay lengths determines then the larger scale of the simulation.
If a photoproduction interaction has occurred, the new energy of the proton (neutron) is substituted for the old one, and the energies and particle types of the secondary particles are recorded. The event generator SOPHIA generates the full set of secondary particles, including nucleon–antinucleon pairs. Thus the total flux of nucleons after propagation is slightly higher than the injected proton flux. Although this is not essential for the main results of this paper, it may occasionally affect the normalization of the proton arrival spectra.
The propagation is completed when the distance between the injection point and the particle location exceeds the predefined source distance. To obtain precise results for the time delay (e.g. total nucleon path length compared to the path length of a light ray), the last integration step is adjusted to end exactly at the desired distance.
Particles are injected at a point in space with a randomly chosen small angular deviation from the $`z`$-axis which defines the main propagation direction. The space along the $`z`$-axis is subdivided into $`32\times 32\times 512`$ cubes of side 250 kpc, each filled with a random magnetic field of average strength $`B`$ = 10<sup>-9</sup> Gauss (1 $`n`$G) satisfying a Kolmogorov spectrum with three logarithmic scales. In practice three field vectors of random orientation are sampled at scales $`\mathrm{}`$ = 1000, 500, and 250 kpc with amplitudes proportional to $`\mathrm{}^{\mathrm{\hspace{0.33em}1}/3}`$ (see Appendix B). The final magnetic field in each of the 250 kpc cubes is the vectorial sum of these three vectors. Cyclic boundary conditions are imposed in case a particle leaves the space of pre-calculated magnetic fields. This means that the magnetic field experienced by a particle at location $`𝐱`$ is the same as the field calculated at $`𝐱^{}`$,
$$x_i^{}=x_iN_iR_i,i=x,y,z$$
(11)
with $`R_i`$ being the size of the pre-calculated magnetic field region in direction $`i`$. $`N_i`$ is the largest integer number satisfying $`x_iN_iR_i0`$. The magnetic field values are refreshed after the calculation of 100 propagations to exclude systematic effects by our choice of field vectors. We have verified numerically that the magnetic field constructed in this way obeys approximately div($`𝐁`$) = 0 and that recalculations of the field at smaller intervals do not change the final result. We assume that the magnetic field strength does not scale with redshift. More information about the implementation of the random magnetic field is given in Appendix B.
The value chosen for $`\mathrm{}_{\mathrm{mag}}`$, in principle, depends strongly on the average magnetic field and nucleon energy, and is a compromise between the precision of the calculation and computing time limits. We have chosen $`\mathrm{}_{\mathrm{mag}}`$ = 10 kpc for $`B`$ = 1 $`n`$G, with an inverse linear scaling for other $`B`$ values. A step size of 1 kpc has been used for short distance propagations to ensure accurate results for arrival angle and time delay distributions.
Finally, it should be mentioned that the calculation of the redshift at a given distance can be done only approximately. The reason is the unknown total travel time of a particle from the source to Earth at injection time. The actual travel time (path length) can be significantly larger than the light travel time along a geodesic and is, in general, different for each simulated particle trajectory. In the following we use the proper distance-redshift relation to define the redshift of the source and along the travel path at observation time. This approximation does not strongly affect our results since we consider here mainly distances with redshifts smaller than 0.06 and weak magnetic fields. However, it should be noted that, in the case of a strong magnetic field, cosmological evolution might become important already at relatively short distances.
## III Results and comparison with previous work
In this section, we present results from the simulation of proton propagation. We start with mono-energetic proton fluxes for which we can compare our results with previous work, and which reflect more directly the different treatments of the energy loss processes. We then compare results for the propagation of protons injected with a power law spectrum.
One can divide previous calculations into two general groups: Monte Carlo based methods, like our own one, and analytical/numerical calculations. Protheroe & Johnson have used a matrix technique to follow the particles over cosmological distances and calculate the $`\gamma `$-ray, neutrino and nucleon spectra arriving at Earth. The energy loss matrices for all particles are calculated with Monte Carlo event generators. We have compared our SOPHIA event generator with the one of Ref. by propagating with the same method an $`E^2`$ proton spectrum with different exponential cutoffs (see Eq. (12)). For this purpose we have used SOPHIA and the event generator of Ref. to calculate the corresponding photoproduction matrices and have applied the two matrices to propagation over the same set of distances. A comparison of the resulting secondary particle spectra yields excellent agreement, pointing to a similar treatment of the particle production process in the different codes. We have also compared the matrix method with our Monte Carlo approach by propagating an exponentially modified power law injection spectrum over 200 Mpc. Again good agreement is found for the resulting $`\nu _\mu `$-spectra, while the $`\overline{\nu }_e`$\- and neutron spectra are at variance with our calculations, which we attribute to a different treatment of the neutron decay. Also, our Monte Carlo method results in more losses due to pair production for distances $`200`$ Mpc and a sharp spike at the injection energy for very short distance propagation, a consequence of the Poisson nature of photon-proton encounters. This feature is discussed in detail in Sect. III.A.
The approach used by Berezinsky & Grigoreva and Rachen & Biermann is to solve the transport equation quasi-analytically by approximating the collisional terms as continuous energy loss terms. This does not take into account the Poissonian nature of the pion production process as pointed out above, and introduces artifacts into the resulting nucleon spectra in form of sharp pile-ups. Lee used a numerical technique to solve the transport equation for particle propagation without using the continuous loss approximation.
The common assumption in all this work is to consider the spatial propagation as strictly along a null-geodesic, with the consequence of not being able to gain knowledge about time delays and arrival angles of the cosmic rays with respect to light and neutrino propagation.
A hybrid model, combining a Monte Carlo particle transport code with analytical techniques was presented by Achterberg et al. . Besides simplifying the properties of the energy losses by analytical estimates (see Fig. 1b), this code also describes the scattering in the magnetic field as a diffusion process employing stochastic differential equations. This approach has the advantage to allow large propagation steps, and is thus computationally very fast, but has a disadvantage at small propagation distances which we discuss further below. Our approach is to use the Monte Carlo technique for simulating particle production and to follow closely cosmic ray orbits in 3D-magnetic field configurations while traveling through the nearby Universe to Earth. This concept, while being the most accurate one, limits our propagation calculation to small source distances.
### A Propagation of mono-energetic protons
In this section we present distributions of arrival energy, arrival direction and time delay of the nucleons, as well as neutrino spectra, for mono-energetic injection of protons at distances of 2, 8, 32, 128 and 512 Mpc from Earth. Protons are injected with energy $`10^{21.5}`$ eV. At this energy, propagated protons can easily suffer several photoproduction interactions, and this tends to emphasize the pion production features.
Fig. 2 shows the distribution of arrival energy of protons and neutrons. Clearly visible is the effect of the statistical nature of photon–proton encounters, also found qualitatively in Ref. . At a distance of 2 Mpc, roughly 60% of all injected particles do not interact, and this generates a sharp spike at the injection energy. This effect due to Poisson statistics remains visible for distances up to $`30`$ Mpc, showing up as a high–energy spike in the cosmic ray spectrum. At larger distances, essentially all injected particles undergo interactions, and therefore, the high-energy spike vanishes. The arrival energy distributions then become much narrower, and in propagation over larger distances would scale simply with the energy loss distance for pair production and adiabatic losses, modified by the increasing scattering in the magnetic field.
Fig. 3 shows the distribution of the average time delay of the cosmic rays arriving at Earth with respect to propagation along a geodesic with the speed of light. This delay is caused by scattering of the charged particles by the intergalactic magnetic field, leading to an increase of the particle’s effective path length. Thus, the average time delay increases with propagation distance, as visible in Fig. 3. Like the arrival energy distributions, the distributions of the time delay also show signs of Poisson statistics, visible especially when propagating over short distances.
The time delay effectively reflects the arrival energy distribution $`t_{\mathrm{del}}1/E_{\mathrm{arr}}^2`$ as a result of the random walk process . This also emphasizes the importance of an accurate treatment of energy losses. For example, a direct comparison with the propagation code of Achterberg et al. for (almost) the same propagation parameters has shown differences in the time delay up to one order of magnitude for $`D=32`$ Mpc. For the same propagation distance, the code by Achterberg et al. produces a peak in the arrival energy distribution about a factor of 2 lower than found in the present work, due to its $`20\%`$ overestimation of energy losses in the photoproduction regime. Together with a difference in the magnetic field sampling, which leads to an effective correlation length $`\mathrm{}_{\mathrm{corr}}390`$ kpc for the Kolmogorov spectrum used in the present work (see Appendix B) compared to $`\mathrm{}_{\mathrm{corr}}=1`$ Mpc for the homogeneous cell approach used in Ref. , the observed differences can then be fully understood by the relation $`t_{\mathrm{del}}\mathrm{}_{\mathrm{corr}}/E_{\mathrm{arr}}^2`$, as derived in Ref. .
Protons with injection energy $`10^{19}`$ eV suffer mainly continuous BH pair production and adiabatic losses that are proportional to their path length. The substantial deflection in the random magnetic field at such energies results in a significant increase of the path length. For protons injected at a sufficiently large distance this can also lead to excessive time delays. For example, cosmic rays with energy of about $`10^{19}`$ eV, injected at distances greater than $`500`$ Mpc in a 1 $`n`$G magnetic field, show a time delay exceeding the Hubble time. This gives a strict constraint on the cosmic ray horizon.
The diffusion coefficient for an effective description of the scattering process in the magnetic field is strongly energy dependent, and so is the time delay, $`t_{\mathrm{del}}`$. To emphasize this correlation, and demonstrate the advantages of the Monte Carlo approach, we show in Fig. 4 the scatter plot of proton energy versus delay after propagation over 32 Mpc. There is a strong correlation suggesting that energy changes of one and a half orders of magnitude lead to differences in delay times of more than three orders of magnitude, i.e. we find an energy dependence similar to $`t_{\mathrm{del}}(BD/E)^2`$ as derived by Achterberg et al. in the small scattering-angle approximation and the quasi-linear approximation of wave-particle interactions. The correlation becomes less pronounced when propagating over significantly larger distances simply because the arrival energy distributions become much narrower and the statistical nature of the energy loss is smoothed by the prevailing pair production and adiabatic losses. This correlation, however, would have very important implications for specific models of UHECR production, where the duration of an active phase of the source competes with the time delay of the protons during propagation. The extreme case would be the acceleration of UHECR in gamma ray bursts. The particles with the highest energies are expected to arrive first, followed by a dissipating widening halo of lower energy protons, as emphasized by Waxman & Miralda–Escudé .
For large propagation distances, even protons injected with 10<sup>21.5</sup> eV show time delays that are a considerable fraction of the light propagation time ($`510\%`$ for $`512`$ Mpc). This would lead to a limiting proton horizon for a large set of source distances and magnetic field values . 512 Mpc is already a limiting horizon for protons injected with 10<sup>19</sup> eV in 1 $`n`$G fields, as noted above.
The scattering that leads to time delay also causes angular deviations from the direction to the source, as shown in Fig. 5 for the injection of mono-energetic protons at the same set of distances. Note that in our propagation code the ‘observer’ sits on a sphere surrounding the injection point. The angle shown is the angle between the particle’s arrival direction and direction to the injection point. This ‘arrival angle’ is somewhat different from the angle between particle’s arrival direction and the injection direction. This method may lead to an underestimate of the scattering angle and the time delay when the particle fluxes become nearly isotropic and many particles have a high probability to scatter back through the ‘observer’s sphere’. It will not, however, affect strongly the results presented in this paper, because, as Fig. 5 demonstrates, we do not reach the limit of isotropic 3D diffusion.
The features of the angular distribution closely follow the time delay distributions already shown. For large propagation distances, the cosmic ray arrival directions are distributed uniformly up to a maximum deflection angle, which increases with propagation distance to reach more than 20 at 512 Mpc. At propagation distances smaller than $``$30 Mpc, thus a few times the proton interaction length $`\lambda _{\mathrm{ph}}`$, a peak at small deflection angles occurs due to the effect of Poisson statistics for proton–photon interactions.
Finally Fig. 6 shows the electron and muon neutrino spectra generated by the injection of 10<sup>21.5</sup> eV protons at the same set of distances. The muon neutrino spectra develop as a function of the proton arrival energy spectra folded with the photoproduction cross section. The fluxes grow with propagation distance, and the maximum neutrino energy shifts to lower energy reflecting the decreasing proton energy. The growth rate with distance decreases for very large distances, where the average proton energy significantly decreases and $`\lambda _{\mathrm{ph}}`$ is correspondingly significantly longer.
Electron neutrino spectra show another, very interesting feature, that develops with distance. At a minimum distance of 2 Mpc the $`\nu _e`$-flux reaches its maximum of 1/2 of the $`\nu _\mu `$ spectrum and shows a somewhat wider energy spectrum, enhanced at low energy. At larger distances an additional $`\nu _e`$ component develops at significantly lower energy. As already noted in Ref. , these are $`\overline{\nu }_e`$’s from neutron decay. The resulting protons from the decay process carry most of the energy, leaving for the $`\overline{\nu }_e`$’s an average energy of only $`5\times 10^4`$ of the original neutron energy, and the $`\nu _e`$-peak is placed at about two orders of magnitude to lower energy with respect to the $`\nu _\mu `$-peak. The strength of this component increases with distance relative to the direct $`\nu _e`$ component from $`\mu ^\pm `$ decay.
### B Cosmological modification of the cosmic ray source spectrum
Berezinsky & Grigoreva introduced the modification factor $`M(E,z)`$ to represent the cosmological evolution of the UHECR spectra. $`M(E,z)`$ gives the ratio of propagated to injected protons at the same energy $`E`$, for a fixed injection spectrum, as a function of the redshift of the injection distance compensating for the proton adiabatic losses. $`M(E,z)`$ is thus exactly unity for proton energies below the $`p\gamma `$-particle production energy threshold.
At the highest injection energies the modification factor shows the GZK cutoff, followed by a pile–up at the crossover of photoproduction and pair production energy loss. This pile–up is a direct consequence of the resonance nature of photoproduction and the hadronic particle production threshold. The next feature at still lower energy is a shallow dip corresponding to the pair production loss, followed by a small pile–up below it. The magnitude of the pile–ups and dips depend not only on the distance and the mean loss distance at the photoproduction/pair production crossover, but also on the shape of the proton injection spectrum. Flatter spectra create bigger pile–ups, because of the increased number of higher energy protons that have interacted to lose energy. The pile–up energy is linked to the energy where losses due to pair production take over from pion production losses, and is therefore strongly dependent on the details of the loss processes in the simulations. Fig. 7a shows $`M(E,z)`$ for propagation without magnetic field for the sole reason of comparison with previous work. An $`E^2`$ proton spectrum with a sharp cutoff at $`E_c=3\times 10^{20}`$ eV is injected, and we propagate over a distance of 256 Mpc in our calculation (solid line) compared to Refs. (dotted line, $`D`$=240 Mpc), (dashed-dotted line, D=256 Mpc), (dashed line, $`D`$=228 Mpc, $`E_c=10^{20}`$ eV) and (dashed-dot-dot-dot line, $`D`$=256 Mpc). There is excellent agreement at all energies with the work of Protheroe & Johnson . The sharp photoproduction peak of Rachen & Biermann is an artifact coming from their continuous loss approximation for pion photoproduction. As noted previously, Yoshida & Teshima used a loss curve which shows a significant deviation from that used in the present paper, and hence their corresponding pile–up height is also larger than in our work. We agree with the position of the pile–up of Lee . However, due to an overestimate of the loss rate at this energy, the magnitude of the pile–up in this paper is smaller than in our model. The dip just below the pile-up is in reasonable agreement with all other works.
Fig. 7b illustrates the effect of scattering in the magnetic field by comparing the resulting corresponding modification factors. The ‘no scattering’ curve (dashed line, as in Fig. 7a) is much higher than the more realistic ‘scattering curve’ in the energy range between 10<sup>18</sup> and 10<sup>19</sup> eV. The reason is that particles in this energy range have considerable time delays and correspondingly much higher total energy loss in pair production interactions. Another consequence of the increased proton travel time due to scattering is the development of a higher pile–up at about 10<sup>18</sup> eV, corresponding to the large number of particles moved to lower energies from the region of that dip. Note that simulation of 10<sup>18</sup> eV particles in a 1 $`n`$G field is at the threshold of our direct Monte Carlo approach, and the calculation is not carried to lower energy where it might show an additional pile–up content. Fig. 7b thus demonstrates the importance of the proton scattering in the extragalactic magnetic fields for the shape of the final spectrum on arrival at Earth.
It is important to note that the curves shown in Fig. 7b are calculated for a source with unlimited lifetime. In addition, by construction, energy loss due to cosmological evolution does not enter the modification factor $`M(E,z)`$. Imposing a constraint on the source lifetime will change the modification factor considerably for low energies because, for a given distance, the time delay due to the scattering in the turbulent magnetic field might become comparable to or even exceed the source lifetime.
## IV Formation of the primary and secondary particle spectra during propagation
To study the development of the primary and secondary particle spectra we followed the propagation of protons injected with a $`E^2`$ power law spectrum with an exponential cutoff at 10<sup>21.5</sup> eV, i.e.
$$\frac{dN}{dE}=AE^2\mathrm{exp}[E/(10^{21.5}\mathrm{eV})].$$
(12)
We recorded the spectra after propagation over 10 Mpc intervals up to a source distance of 200 Mpc. The results of this calculation are relevant for models of UHECR acceleration at astrophysical shock fronts, although the cutoff energy adopted in this calculation is fairly high. 10,000 protons were injected with a power law spectrum (integral spectral index $`\gamma `$ = 1) in each of 30 energy bins covering energies from 10<sup>19</sup> to 10<sup>22</sup> eV, i.e. 10 bins per decade of energy. We did not simulate the propagation of lower energy particles, which do not experience photoproduction interactions, but followed the secondaries down to arbitrary low energies.
Fig. 8a shows the evolution of the particles injected in the highest energy bin 10<sup>21.9</sup> to 10<sup>22</sup> eV. The size of each rectangle is proportional to the fractional energy distribution after propagation over 10, 20, etc., Mpc. The rate of energy degradation is dramatic. After only 10 Mpc the spectrum of protons injected in a 0.1 logarithmic bin have spread over one and a half orders of magnitude. The width of the energy distribution increases with the propagation distance up to $``$30 Mpc and then decreases. Qualitatively this behavior is very similar to the calculation of Aharonian & Cronin , although the direct comparison is difficult because of the different approach to the calculation. The average behavior of all protons injected with energy above about $`3\times 10^{20}`$ eV is similar, although the magnitude of the spread decreases — particles of energy below 10<sup>20</sup> eV suffer much smaller losses. After propagation over about 100 Mpc the spectrum shown in Fig. 8 is already final - it is concentrated within roughly 1/2 order of magnitude around $`8\times 10^{19}`$ eV. This energy slowly decreases because of pair production and adiabatic losses during propagation over larger distances, but without change in the shape of the distribution.
The lower panel of Fig. 8 shows the fractional energy carried by different particles after propagation in terms of the total energy of the protons injected with energy spectrum described by Eq. 12. The proton curve, which also includes neutrons, always dominates. The energy content in protons, however, is only about 50% of that injected for distances above 120 Mpc. The rest of the injected energy is distributed between the electromagnetic component and neutrinos. Note the difference between the photon (and electron) components from photoproduction (long dashed line), and from pair production (short dashed line). While the photoproduction component rises very quickly and changes very little after 100 Mpc, the pair production component is almost proportional to the distance, as most of the injected protons, despite the high threshold of 10<sup>19</sup> eV, have similar pair production losses. At distances of 100 (200) Mpc 51% (43%) of the injected power is carried by nucleons, 31% (37%) by the electromagnetic component and 18% (20%) by neutrinos. The neutrino fluxes will remain at the same level during propagation over larger distances, and the established energy balance will only slightly change as nucleons yield some of their power to the electromagnetic component through pair production. Adiabatic losses will, of course, affect all components in the same way.
In addition to distributing a fraction of the energy of the injected protons to secondary particles, the propagation changes the energy spectrum of protons. The most energetic nucleons lose energy fast and are downgraded after a short propagation distance. The number of nucleon with energy above $`10^{21}`$ eV decreases by 10%, 50% and 90% from the injected number of protons after only 1, 6, and 20 Mpc. The corresponding distances for nucleons of energy above $`10^{20}`$ eV are 10, 40 and 85 Mpc. The magnitude of these changes emphasizes the importance of detection of very high energy particles: for particles of energy above 3$`\times `$10<sup>20</sup> eV (same as the highest energy event detected by the Fly’s Eye ) these distances are 1, 10 and 30 Mpc. The rapid absorption of the highest energy cosmic rays implies that the horizon of the highest energy protons is very small, and increases the energetics requirements for potential UHECR sources.
## V Discussion, conclusions and outlook
The Monte Carlo propagation of ultra high energy protons in a random extragalactic magnetic field has obvious advantages over other approaches to calculations of proton propagation in the cosmologically nearby Universe. To start with, this approach takes fully into account fluctuations in the positions of proton interactions, and thus also in the proton energy losses and production of secondary particle fluxes. It also naturally generates the correlations between the proton’s arrival energy, its time delay, and its angular deviation from the source direction. We have also shown that mathematical approaches which use a diffusion description of magnetic scattering, although superior in computational speed, can lead to significant systematic errors for propagation distances smaller then $`100`$ Mpc.
These features of the calculation become extremely valuable when applied to specific models of UHECR acceleration, especially models that involve a relatively short (compared to light travel time and proton time delay) active phase of the source. An extreme example for such a model is the GRB model for UHECR acceleration. However, other models involving interacting galaxies or radio galaxies of specific morphology could also be affected, especially if embedded in regions of high (random) magnetic field.
At energies that allow protons to photoproduce, namely above 10<sup>20</sup> eV, the energy degradation is extremely rapid. This is not very surprising because of the very short photoproduction interaction length at energies corresponding to the maximum cross section – i.e. $`\lambda _{\mathrm{ph}}`$ below 4 Mpc for energies between 4$`\times `$10<sup>20</sup> eV and 10<sup>21</sup> eV. This energy range is very relevant, as it is just above the highest energy particles detected by the Fly’s Eye and AGASA arrays . A large part of this rapid energy dissipation in our calculation is due to the correct implementation of the fluctuations in photoproduction interactions in SOPHIA. A good example for the size of the fluctuations is the proton energy distribution after propagation over 10 Mpc shown in Fig. 8, which covers more than one and a half orders of magnitude. This is an extreme case. However, every particle injected with an energy well above the photoproduction threshold would very rapidly result in a distribution extending down to the threshold, within the first 10 Mpc.
This rapid energy dissipation creates additional problems for models of cosmic ray acceleration at astrophysical shocks. Apart from the difficult question of the maximum acceleration energy, such models require that a significant fraction (0.01 to 0.1) of their source luminosity contributes to the UHECR flux. The rapid energy dissipation increases the energy requirements in terms of total luminosity and severely limits the source distance. Because of magnetic scattering, such limits could also be set for particles injected with energy below the photoproduction threshold.
Fig. 9 shows the 50% horizon for UHECR sources as a function of source particle energy for $`B`$ values of 0.1, 1 and 10 $`n`$G. The 50% horizon $`R_{50}`$ is defined here as the light propagation distance to the source at which 1/$`e`$ of all injected protons have retained 50% or more of their energy, i.e. $`R_{50}`$ is achieved when
$$_{\frac{E_0}{2}}^{E_0}\frac{dN}{dE}𝑑E=N_0\mathrm{exp}(1),$$
(13)
where $`N_0`$ is the number of particles injected with energy $`E_0`$.
To start with, $`R_{50}`$ is small at any energy, and demonstrates the resonant nature of the photoproduction cross section. At E = 10<sup>20</sup> eV $`R_{50}`$ is about 100 Mpc, while at 2$`\times `$10<sup>20</sup> eV it decreases to 20 Mpc and becomes smaller than 10 Mpc for energies above 3$`\times `$10<sup>20</sup> eV. For injection energies above 10<sup>20</sup> eV the horizon energy dependence is similar to that of the energy loss distance shown in Fig. 1. These protons are not affected much by the magnetic field since their scattering angles are small, but suffer mainly from energy degradation due to $`p\gamma `$ encounters. Below 10<sup>20</sup> eV the picture changes. The scattering in the magnetic field increases the propagation time and thus causes additional energy loss and an increase of the ratio $`x_{\mathrm{loss}}/R_{50}`$.
Stronger magnetic fields create delays, that could be longer than the light propagation time from the source and reverse the trend – the horizon starts decreasing below $`6\times 10^{19}`$ eV and is restricted to 75 Mpc at 10<sup>19</sup> eV. Since the average time delay is inversely proportional to $`E^2`$, the decrease of $`R_{50}`$ is expected to become more drastic at lower energy. One consequence of the strong energy dependence of $`R_{50}`$ is, for example, that our attempts to correlate the arrival directions of UHECR with different types of astrophysical objects should use only objects within the particle horizon depending on the magnetic fields strength in different regions of the Universe. Independently of the magnetic field value, however, the horizon defined above is much smaller than the conventional numbers of 50 or 100 Mpc for the highest energy cosmic ray events.
There are many relevant astrophysical problems which can be studied with the approach described in this paper. We plan to use the code for proton propagation in regular magnetic fields associated with large scale structures (local supercluster, supergalactic plane). The regular fields, especially if they reach the observationally allowed limits of 0.03 $`\mu `$G and even 0.1 $`\mu `$G, could change the propagation patterns for 10<sup>19</sup> eV cosmic ray protons and alter the horizon values shown in Fig. 9. We also plan to set limits on models of slow UHECR acceleration on shocks of very large dimensions and to look for possibilities of ultra-high energy $`\gamma `$–ray halos around the sources and along the tracks of the UHECR protons.
## Acknowledgments
The authors are indebted to A. Achterberg and the authors of Ref. for sharing their corrected results prior to publication, and to P.P. Kronberg for careful reading of the manuscript and valuable discussions. The research of TS is supported in part by NASA Grant NAG5-7009. RE is supported in part by the US Department of Energy contract DE-FG02 91ER 40626. The work of RJP and AM was supported in part by a grant from the Australian Research Council. JPR is supported through the TMR network Astro–Plasma Physics, funded by the EU under contract FMRX-CT98-0168. AM thanks BRI for its hospitality during her visit. These calculations are performed on DEC Alpha and Beowulf clusters funded by NSF grant PHY–9601834.
## A Monte Carlo sampling of interaction points
In the following we discuss the application of the veto algorithm to the sampling of interaction points along a nucleon propagation path. The probability of having no hadronic interaction with a photon of the CMB within a path length interval $`(s_1,s_2)`$ reads
$$P_{\mathrm{no}}(s_1,s_2)=\mathrm{exp}\left\{_{s_1}^{s_2}\frac{ds}{\lambda _{\mathrm{ph}}(E(s))}\right\}.$$
(A1)
The interaction length itself depends only on the nucleon energy. However, because of the treatment of Bethe-Heitler losses as continuous process, this energy depends on the path length $`s`$. Correspondingly, the probability for one interaction in the interval $`(s,s+ds)`$ is given by
$$P_{\mathrm{int}}(s)ds=P_{\mathrm{no}}(0,s)\frac{ds}{\lambda _{\mathrm{ph}}(E(s))},$$
(A2)
where $`P_{\mathrm{no}}(0,s)`$ is the probability that no interaction has occurred before. In our approach we replace $`\lambda _{\mathrm{ph}}(E(s))`$ by the constant $`\lambda _{\mathrm{ph},\mathrm{min}}`$ and use (A2) to sample the path length distance from the current location ($`s=0`$) to the next interaction. This interaction point is then accepted with the probability $`\lambda _{\mathrm{ph},\mathrm{min}}/\lambda _{\mathrm{ph}}(E(s))`$. Hence the interaction probability can be written as
$`P_{\mathrm{int}}(s)ds`$ $`=`$ $`[\stackrel{~}{P}_{\mathrm{no}}(0,s)`$ (A7)
$`+{\displaystyle _0^s}{\displaystyle \frac{ds_1}{\lambda _{\mathrm{ph},\mathrm{min}}}}\stackrel{~}{P}_{\mathrm{no}}(0,s_1)\left(1{\displaystyle \frac{\lambda _{\mathrm{ph},\mathrm{min}}}{\lambda _{\mathrm{ph}}(E(s_1))}}\right)\stackrel{~}{P}_{\mathrm{no}}(s_1,s)`$
$`+{\displaystyle _0^s}{\displaystyle \frac{ds_1}{\lambda _{\mathrm{ph},\mathrm{min}}}}\stackrel{~}{P}_{\mathrm{no}}(0,s_1)\left(1{\displaystyle \frac{\lambda _{\mathrm{ph},\mathrm{min}}}{\lambda _{\mathrm{ph}}(E(s_1))}}\right)`$
$`\times {\displaystyle _{s_1}^s}{\displaystyle \frac{ds_2}{\lambda _{\mathrm{ph},\mathrm{min}}}}\stackrel{~}{P}_{\mathrm{no}}(s_1,s_2)(1{\displaystyle \frac{\lambda _{\mathrm{ph},\mathrm{min}}}{\lambda _{\mathrm{ph}}(E(s_2))}})\stackrel{~}{P}_{\mathrm{no}}(s_2,s)`$
$`+\mathrm{}\left]\right({\displaystyle \frac{\lambda _{\mathrm{ph},\mathrm{min}}}{\lambda _{\mathrm{ph}}(E(s))}}){\displaystyle \frac{ds}{\lambda _{\mathrm{ph},\mathrm{min}}}},`$
where we have used
$$\stackrel{~}{P}_{\mathrm{no}}(s_2,s_1)=\mathrm{exp}\left\{\frac{s_1s_2}{\lambda _{\mathrm{ph},\mathrm{min}}}\right\}.$$
(A8)
The first term in square brackets corresponds to the probability that no interaction was sampled in the interval $`(0,s)`$. The second term is the contribution which comes from an interaction point sampled at $`s_1`$ but rejected with the probability $`1\lambda _{\mathrm{ph},\mathrm{min}}/\lambda _{\mathrm{ph}}`$.
The integration limits in (A7) ensure the ordering of the interaction points according to the simulation method, $`0<s_1<s_2<\mathrm{}<s`$. Symmetrizing the integration limits yields
$`P_{\mathrm{int}}(s)ds`$ $`=`$ $`{\displaystyle \frac{ds}{\lambda _{\mathrm{ph}}(E(s))}}\mathrm{exp}\left\{{\displaystyle \frac{s}{\lambda _{\mathrm{ph},\mathrm{min}}}}\right\}`$ (A10)
$`\times {\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n!}}\left[{\displaystyle _0^s}ds^{}({\displaystyle \frac{1}{\lambda _{\mathrm{ph},\mathrm{min}}}}{\displaystyle \frac{1}{\lambda _{\mathrm{ph}}(E(s^{}))}})\right]^n`$
$`=`$ $`\mathrm{exp}\left\{{\displaystyle _0^s}{\displaystyle \frac{ds^{}}{\lambda _{\mathrm{ph}}(E(s^{}))}}\right\}{\displaystyle \frac{ds}{\lambda _{\mathrm{ph}}(E(s))}},`$ (A11)
which is identical to (A2) and shows that the described simulation method reproduces the correct, energy-dependent interaction length.
## B Implementation of the magnetic field
A turbulent magnetic field which is frozen into a fluid with fully developed hydrodynamic turbulence would follow a Kolmogorov spectrum, which is defined by
$$I(k)=I_0(k/k_0)^{5/3}$$
(B1)
where $`k`$ is the wavenumber . $`I(k)`$ is the energy density per unit wave number, $`k_0`$ the smallest wavenumber of the turbulence, the inverse $`k_0^1`$ is sometimes called the “cell size” of the turbulence. Hence we have for the total energy density
$$U_{\mathrm{tot}}=\frac{B_{\mathrm{rms}}^2}{8\pi }=dkI(k).$$
(B2)
In the propagation program we consider 3 discrete wave numbers. Thus we have to rewrite this integral in terms of a discrete spectrum in $`k`$, starting with $`k_0`$ and continuing with $`k_i=2k_{i1}`$, $`i=1,2`$. These are equally spaced apart in $`\mathrm{log}_2k`$, with $`\mathrm{\Delta }(\mathrm{log}_2k)`$=1. Hence the energy density we should ascribe to each of the three wavenumbers is approximately
$`U_i`$ $``$ $`{\displaystyle \frac{I(k)\mathrm{d}k}{\mathrm{d}(\mathrm{log}_2k)}}|_{k_i}\mathrm{\Delta }(\mathrm{log}_2k)`$ (B3)
$`=`$ $`I_0k_0\mathrm{ln}2\left({\displaystyle \frac{k_i}{k_0}}\right)^{2/3}.`$ (B4)
The total energy density is then a simple sum,
$$\frac{B^2}{8\pi }=U_0+U_1+U_2.$$
(B5)
We normalize the field to a total energy density corresponding to $`|B|=1\mathrm{nG}`$, i.e. $`U_{\mathrm{tot}}4\times 10^{20}\mathrm{erg}\mathrm{cm}^3`$.
The technical implementation of the magnetic field into our propagation code is as follows. We divide the propagation volume into cubes of $`1\mathrm{Mpc}`$ side length, and attach to each of them a homogeneous field $`𝐁_\mathrm{𝟎}`$ with magnitude $`B_0`$ and random direction. Each of these cubes is divided into 8 cubes of $`0.5\mathrm{Mpc}`$ side length, to which a field $`𝐁_\mathrm{𝟏}`$ of magnitude $`B_1`$ and random direction is vectorially added to the field $`𝐁_0`$. The procedure is repeated once more, so that our field is eventually realized on elementary cubes of $`0.25\mathrm{Mpc}`$ side length, each of which carries a magnetic field $`𝐁_\mathrm{𝟎}+𝐁_\mathrm{𝟏}+𝐁_\mathrm{𝟐}`$. We check that $`\mathrm{div}𝐁0`$ by approximating the surface integral with the sum of the outward normal component of $`𝐁`$ over the surface of the 8$`\times `$8$`\times `$128 Mpc<sup>3</sup> volume $`V`$. The volume averaged value of div($`𝐁`$) is calculated as
$$B=\frac{1}{V}𝐁_{}ds.$$
(B6)
The r.m.s. value of $`B`$ for 10,000 field realizations is $`B_{\mathrm{rms}}`$ = $`3.7\times 10^6`$ $`n`$G/kpc.
We also calculate the effective correlation length $`\mathrm{}_{\mathrm{corr}}`$ by equating
$$𝐁(𝐱)𝐁(𝐱+\xi )=𝐁^2(𝐱)\mathrm{exp}\left(\frac{|\xi |}{\mathrm{}_{\mathrm{corr}}}\right).$$
(B7)
The best fit value of $`\mathrm{}_{\mathrm{corr}}`$ is 390 kpc.
|
warning/0003/hep-ph0003251.html
|
ar5iv
|
text
|
# Construction of a Minimal Higgs 𝑆𝑂(10) SUSY GUT Model
## I INTRODUCTION
In several recent papers \[1-4\] we have developed a highly predictive model of quark and lepton masses based on the grand unified group $`SO(10)`$. This model grew out of our attempt to construct a realistic grand unified theory (GUT) in which $`SO(10)`$ was broken down to the standard model in the simplest possible, or “minimal” way . In this model there emerged a new mechanism based on certain well-known features of $`SU(5)`$ for explaining the large mixing between the mu and tau neutrinos that is seen at SuperKamiokande. In we gave the structure of the quark and lepton mass matrices for the second and third families, treating the first family as massless. In , it was shown how to extend the model to include the first family, which leads to several interesting predictions. In , it was observed that the mixing of the electron neutrino very naturally falls either within the range $`0.004\mathrm{sin}^22\theta _{e\mu }0.008`$, corresponding to the small angle MSW solution of the solar neutrino problem, or very near to the value $`\mathrm{sin}^2\theta _{e\mu }=1`$, corresponding to what is called “bimaximal mixing”.
In this paper we present the model in fuller detail, especially in regard to neutrino phenomenology, and to the structure of the Higgs sector, Yukawa interactions, and flavor symmetries that underlie the quark and lepton mass matrices.
The paper is organized as follows. In Section II we discuss in general terms, that is apart from a particular model, our mechanism for explaining the large mixing of $`\nu _\mu `$ and $`\nu _\tau `$. In Section III, we explain what we mean by a “minimal” $`SO(10)`$ breaking scheme, and show how such minimal breaking and the requirements of simplicity lead one naturally to a certain form for the mass terms of the heavier two families of quarks and leptons. We then observe that this form realizes the general mechanism for large $`\nu _\mu \nu _\tau `$ mixing described in the previous Section. It is important to emphasize that this mechanism emerged not from an attempt to explain neutrino phenomenology, but from other considerations entirely, in particular the attempt to simplify the Higgs structure of $`SO(10)`$. It is most interesting that the same mechanism has also independently been found by other groups attempting to make sense of neutrino phenomenology. In Section IV, it will be explained how this model is best extended to the first family of quarks and leptons, and how this gives rise to several distinctive predictions. Accurate analytic expressions for the predictions at the GUT scale will be presented. In Section V, the neutrino sector will be examined in detail. It will be seen how either the small-angle MSW solution of the solar neutrino problem or bimaximal mixing can result with equal simplicity. Finally, in Section VI, a concrete model, including all the details of flavor symmetries and of the Higgs and Yukawa superpotentials, will be presented, showing that the basic scheme is technically natural.
## II MECHANISM FOR LARGE $`\nu _\mu \nu _\tau `$ MIXING
Before explaining our mechanism, it will be helpful to explain why the observed large mixing of $`\nu _\mu `$, presumably with $`\nu _\tau `$, has been a theoretical puzzle. The basic reason is simple: the mixing that is seen between the quarks of the second and third families is described by a small mixing angle, namely $`V_{cb}0.04`$, and therefore it was expected that the mixing between the second and third family of leptons would also be small.
The grounds for this expectation were twofold. First, there is the empirical fact that the masses of the quarks and leptons exhibit roughly similar “hierarchical” patterns, and therefore it was natural to assume that their mixing angles would be similar also. Second, the most promising theoretical approaches to understanding the pattern of quark and lepton masses, namely grand unification and flavor symmetry, tend to treat quarks and leptons in similar ways. For instance, small quark mixing angles might suggest an underlying fundamental “family symmetry” that is weakly broken, in which case the lepton mixings might be expected also to be small. And in grand unification based on $`SO(10)`$ there is a close connection between the quark and lepton mass matrices.
There are actually two puzzles associated with the mixing of the second and third families: First, why is the lepton mixing $`|U_{\mu 3}|0.7`$ so large? And, second, why is the quark mixing $`V_{cb}0.04`$ so small? What we mean by saying that these are distinct puzzles is that they are both unexpected within the most commonly assumed framework for explaining quark and lepton masses, the Weinberg-Wilczek-Zee-Fritzsch (WWZF) idea .
The WWZF idea was that the Cabibbo angle could be understood if the mass matrices of the first and second families of quarks had the following form:
$$L_{mass}=(\overline{u_R},\overline{c_R})\left(\begin{array}{cc}0& b\\ b& a\end{array}\right)\left(\begin{array}{c}u_L\\ c_L\end{array}\right)+(\overline{d_R},\overline{s_R})\left(\begin{array}{cc}0& b^{}\\ b^{}& a^{}\end{array}\right)\left(\begin{array}{c}d_L\\ s_L\end{array}\right).$$
(1)
This gives $`\left|m_d/m_s\right|\left|b^{}/a^{}\right|^2`$, $`\left|m_u/m_c\right|\left|b/a\right|^2`$, and $`V_{us}b^{}/a^{}b/a`$, and thus the famous relation
$$V_{us}\sqrt{m_d/m_s}e^{i\alpha }\sqrt{m_u/m_c}.$$
(2)
Since $`\left|V_{us}\right|0.22`$, $`\sqrt{m_d/m_s}0.22`$, and $`\sqrt{m_u/m_c}0.07`$, this relation is satisfied for $`\alpha \pm \pi /2`$. If we apply the same idea to the leptons of the first two families we get
$$U_{e2}\sqrt{m_e/m_\mu }e^{i\beta }\sqrt{m_{\nu _1}/m_{\nu _2}}.$$
(3)
The second term on the right is not known, but if it is assumed to be small one has the rough prediction that $`U_{e2}\sqrt{m_e/m_\mu }0.07`$. This could be consistent with the small angle MSW solution of the solar neutrino problem, which requires that $`U_{e2}0.04`$. Thus the WWZF idea appears to work well where it was originally applied, namely to the first and second families.
Fritzsch later extended this idea to explain the mixing of the third family. If a WWZF form is assumed to hold for the second and third family, i.e., if one takes $`(u,c)(c,t)`$ and $`(d,s)(s,b)`$ in Eq. (1), one obtains
$$V_{cb}\sqrt{m_s/m_b}e^{i\gamma }\sqrt{m_c/m_t}$$
(4)
and
$$U_{\mu 3}\sqrt{m_\mu /m_\tau }e^{i\delta }\sqrt{m_{\nu _2}/m_{\nu _3}}.$$
(5)
Since $`\sqrt{m_s/m_b}0.14`$, and $`\sqrt{m_c/m_t}0.04`$, one sees that the observed value of $`V_{cb}0.04`$ is too small by a factor of three or so. Assuming that the neutrino mass ratio in Eq. (5) is small, and given that $`\sqrt{m_\mu /m_\tau }0.24`$, one sees that the nearly maximal value of $`U_{\mu 3}1/\sqrt{2}0.7`$ that is observed is too large by a factor of three or so.
These Eqs. (2-5) are based on the assumption of a hierarchical and symmetric form for the mass matrices. A key feature in our mechanism for understanding the large mixing of the tau neutrino is that it involves highly asymmetric mass matrices. As we shall see, the assumption of asymmetric mass matrices naturally explains why $`U_{\mu 3}`$ is larger than the Fritzsch value and $`V_{cb}`$ is smaller than the Fritzsch value by approximately the same factor.
Consider, a toy model with $`SU(5)`$ symmetry, which has a set of Yukawa terms of the following form: $`\lambda _{33}(\overline{\mathrm{𝟓}}_3\mathrm{𝟏𝟎}_3)\overline{\mathrm{𝟓}}_H+\lambda _{23}(\overline{\mathrm{𝟓}}_2\mathrm{𝟏𝟎}_3)\overline{\mathrm{𝟓}}_H+\lambda _{32}(\overline{\mathrm{𝟓}}_3\mathrm{𝟏𝟎}_2)\overline{\mathrm{𝟓}}_H`$, with $`\lambda _{32}\lambda _{23}\lambda _{33}`$ and the subscript $`H`$ denoting a Higgs representation. These terms yield the following mass matrices for the second and third families of down quarks and charged leptons:
$$(\overline{d_{2R}},\overline{d_{3R}})\left(\begin{array}{cc}0& \sigma \\ ϵ& 1\end{array}\right)\left(\begin{array}{c}d_{2L}\\ d_{3L}\end{array}\right)M_D+(\overline{l_{2R}},\overline{l_{3R}})\left(\begin{array}{cc}0& ϵ\\ \sigma & 1\end{array}\right)\left(\begin{array}{c}l_{2L}\\ l_{3L}\end{array}\right)M_D,$$
(6)
with $`ϵ\sigma 1`$. Here we have labelled the fermions with a family index, instead of the names $`s`$, $`b`$, $`\mu `$, and $`\tau `$, since the mass matrices in this case are far from diagonal. A crucial point to notice is that the matrix for the leptons, which we will denote by $`L`$, is the transpose of the matrix for the down quarks, which we will denote by $`D`$. This is a feature of minimal $`SU(5)`$. It arises from the fact that the $`\overline{\mathrm{𝟓}}`$ representation of fermions contains the left-handed leptons, $`l_L`$, and the charge conjugate of the right-handed down-quarks, $`d_R`$, while the $`\mathrm{𝟏𝟎}`$ representation of fermions contains the charge conjugate of the right-handed leptons, $`l_R`$, and the left-handed down-quarks, $`d_L`$. Thus, $`SU(5)`$ relates $`D`$ to $`L`$, but only up to a left-right transposition: $`D=L^T`$.
The transposition feature of $`SU(5)`$ unification appearing in Eq. (6) results in the large element, $`\sigma `$, of $`L`$ producing an O(1) mixing of $`l_{2L}`$ with $`l_{3L}`$ for the leptons, while in $`D`$ for the quarks it produces a large mixing of the right-handed fields $`d_{2R}`$ and $`d_{3R}`$. The mismatch between the large $`l_{2L}l_{3L}`$ mixing and the $`\nu _{2L}\nu _{3L}`$ mixing, which is small (as will soon be seen), leads to a large $`U_{\mu 3}`$ mixing element. But the right-handed mixings of the quarks are not observable through standard model physics. What matters is the left-handed mixing of $`d_{2L}`$ with $`d_{3L}`$, which contributes to $`V_{cb}`$, and is controlled by the small parameter $`ϵ`$.
The common statement that grand unification relates quark and lepton mixing angles, and thus $`V_{cb}`$ to $`U_{\mu 3}`$, is very misleading. What is really true in general is that grand unification relates the mixing of quarks of one handedness to the mixing of leptons of the other handedness. Thus $`V_{cb}`$ and $`U_{\mu 3}`$ need not be directly related to each other. Of course, if the mass matrices are symmetric, as has almost always been assumed, the left-handed and right-handed mixings are the same, and hence $`V_{cb}`$ is directly related to $`U_{\mu 3}`$. The most natural interpretation, then, of the experimental discovery that $`|U_{\mu 3}||V_{cb}|`$ is that the mass matrices are highly asymmetric. This is the essential point first made in .
Not only does a highly asymmetric, or, as we will call it, “lopsided,” form of the mass matrices explain the difference between the size of $`U_{\mu 3}`$ and $`V_{cb}`$, but it also explains the fact, noted above, that $`U_{\mu 3}`$ is larger than the Fritzsch value and $`V_{cb}`$ is smaller than the Fritzsch value by about the same factor. The point is that the product of the two off-diagonal elements, $`ϵ`$ and $`\sigma `$, is controlled by the fermion mass ratio. As is evident from Eq. (6), $`m_s/m_b\frac{ϵ\sigma }{1+\sigma ^2}ϵ\sigma `$. That means that the Fritzsch prediction for the mixing of $`d_L`$ and $`s_L`$, which is $`\sqrt{m_s/m_b}`$, goes approximately as $`\sqrt{ϵ\sigma }`$. That shows that the Fritzsch prediction for the mixing angles is roughly the geometric mean between the true value of $`U_{\mu 3}\sigma `$ and the true value of $`V_{cb}`$ $`ϵ`$. In other words, in our hypothesis of lopsided mass matrices, the surprising largeness of $`U_{\mu 3}`$ and the surprising smallness of $`V_{cb}`$ are two sides of the same coin.
Another important feature of this mechanism should be emphasized. Almost all published explanations of the largeness of the $`\nu _\mu \nu _\tau `$ mixing trace it to some special feature or form of the neutrino mass matrix. Perhaps this is due to the purely linguistic fact that we talk about “neutrino mixing angles”. But they could just as well be called the “charged-lepton mixing angles”. They are really the angles expressing the mismatch between the neutrino mass eigenstates and the charged lepton mass eigenstates, just as the CKM angles are the mismatch between the up and down quark eigenstates. In our mechanism, the large value of $`U_{\mu 3}`$ is traceable to a peculiarity of the charged lepton mass matrix $`L`$, namely, having a large off-diagonal entry $`\sigma `$. As we shall see in the next Section, having such a large entry helps to explain several other features of the quark and lepton mass spectrum.
To sum up, the mechanism for explaining large $`\nu _\mu \nu _\tau `$ mixing proposed in has three salient features: (1) the largeness of this mixing is due to the charged lepton mass matrix, which is (2) highly asymmetric, and which is (3) related to the transpose of the down quark mass matrix by $`SU(5)`$.
In the next Section we will see how a model with precisely these features arises very naturally in $`SO(10)`$ from very different considerations.
## III FERMION MASS MATRICES IN MINIMAL SCHEMES OF $`SO(10)`$ BREAKING
The model that we shall examine in this paper emerged originally from our attempt to construct a realistic model based on $`SO(10)`$ in which $`SO(10)`$ is broken to the standard model group, $`G_{SM}=SU(3)_c\times SU(2)_L\times U(1)_Y`$ in the simplest possible way. We shall therefore start by explaining what we mean by minimal $`SO(10)`$ breaking.
Since $`SO(10)`$ is a rank 5 group, it requires for its breakdown to $`G_{SM}`$ at least two Higgs fields. One Higgs field is needed to break the rank of the group to 4, but this generally leaves an unbroken $`SU(5)`$. The second Higgs field is needed to break $`SU(5)`$ down to $`G_{SM}`$. The two breakings can occur in either sequence depending upon which Higgs field has the larger VEV and effects the first breaking.
Whatever Higgs field gives superlarge mass to the right-handed neutrinos, as required for the standard seesaw explanation of the lightness of the left-handed neutrinos, will also break $`SO(10)`$ down to $`SU(5)`$, and thus the rank to 4. There are two simple choices for this Higgs field: either an antisymmetric five-index tensor $`\overline{\mathrm{𝟏𝟐𝟔}}`$ or a spinor $`\overline{\mathrm{𝟏𝟔}}`$. In either case, one also expects a Higgs field in the conjugate representation, $`\mathrm{𝟏𝟐𝟔}`$ or $`\mathrm{𝟏𝟔}`$, to go along with it. A nice feature of the $`\overline{\mathrm{𝟏𝟐𝟔}}`$ is that this tensorial representation leaves unbroken a $`Z_2`$ subgroup of the center of $`SO(10)`$ that acts as an automatic matter parity, whereas if a spinor Higgs is introduced, then matter parity is not automatic. On the other hand, to introduce $`\mathrm{𝟏𝟐𝟔}+\overline{\mathrm{𝟏𝟐𝟔}}`$ is to introduce quite large representations that tend to make the unified gauge coupling go non-perturbative below the Planck scale, and that may be hard to obtain from superstring theory. In any event, it would seem that the use of a spinor-antispinor pair, $`\mathrm{𝟏𝟔}+\overline{\mathrm{𝟏𝟔}}`$, is more economical. Thus we assume that the rank of $`SO(10)`$ is broken at the unification scale and the right-handed neutrinos get mass from one such spinor-antispinor pair of Higgs fields.
To break the group the rest of the way to $`G_{SM}`$ requires the existence of Higgs fields in the adjoint representation $`\mathrm{𝟒𝟓}`$ and/or in the symmetric two-index tensor representation $`\mathrm{𝟓𝟒}`$. Most published realistic $`SO(10)`$ models have several of both kinds of multiplets. However, it has been shown that it is possible to break $`SO(10)`$ to $`G_{SM}`$ with only a single adjoint Higgs and no larger representations .
This, then, is what we call the “minimal breaking scheme for $`SO(10)`$”: The breaking of $`SO(10)`$ to $`G_{SM}`$ is accomplished by the expectation values of a set of Higgs fields consisting of $`\mathrm{𝟒𝟓}_H+\mathrm{𝟏𝟔}_H+\overline{\mathrm{𝟏𝟔}}_H`$, with the model containing no multiplets larger than the single $`\mathrm{𝟒𝟓}`$. There, of course, have to be other Higgs fields to break the $`SU(2)_L\times U(1)_Y`$ group of the electroweak interactions.
This minimality assumption is restrictive enough that it is possible to say in which direction the expectation values of these fields point. This can be done by considering the problem of doublet-triplet splitting, whereby the colored partners of the weak-doublet Higgs fields of the standard model become superheavy while the weak-scale masses of the doublets themselves are preserved. In $`SO(10)`$ the only known way of doing this in a technically natural manner is the Dimopoulos-Wilczek or “missing VEV” mechanism . The idea is that if an adjoint Higgs field that has an expectation value proportional to the $`SO(10)`$ generator $`BL`$ couples to Higgs fields in the vector representation, it will make their color-triplet components heavy (since they have $`BL=\pm 2/3`$) while leaving their weak-doublet components massless (since they have $`BL=0`$). The needed coupling is simply of the form $`\mathrm{𝟏𝟎}_{1H}\mathrm{𝟒𝟓}_H\mathrm{𝟏𝟎}_{2H}`$. Of course, the expectation value of the adjoint, by virtue of the definition of the adjoint representation, is necessarily a linear combination of generators of the group. $`BL`$ is one of the $`SO(10)`$ generators that is picked out by simple forms of the Higgs superpotential for the adjoint multiplet. It should be noted that there is another version of the missing VEV mechanism that works in $`SO(10)`$, in which the VEV of the adjoint is proportional to the generator $`I_{3R}`$ of the $`SU(2)_R`$ subgroup of $`SO(10)`$ . However, that version is significantly more complicated. Therefore, simplicity dictates the choice that
$$\mathrm{𝟒𝟓}_HBL.$$
(7)
Since the assumption of a minimal $`SO(10)`$ breaking scheme included the supposition that only one adjoint exists in the model, no adjoint exists except the one that points in the $`BL`$ direction. As we shall see, this puts an important limitation on the possibilities for constructing realistic mass matrices for the quarks and leptons. The assumption of a minimal $`SO(10)`$ breaking scheme thus acts as an important guide in searching for good models.
The simplest possible terms that would give mass to quarks and leptons in $`SO(10)`$ would be $`\lambda _{ij}\mathrm{𝟏𝟔}_i\mathrm{𝟏𝟔}_j\mathrm{𝟏𝟎}_H`$, where the subscripts $`i`$ and $`j`$ are family indices. This would lead to four proportional Dirac mass matrices for the up-quarks ($`U`$), down quarks ($`D`$), charged leptons ($`L`$), and neutrinos ($`N`$). In fact one would have $`D=LU=N`$. Moreover, all these matrices would be symmetric, which is why one can write $`D=L`$ instead of $`D=L^T`$ as in minimal $`SU(5)`$. Some of the predictions that follow from these relations are good, notably the famous prediction $`m_b^0=m_\tau ^0`$, where the superscript zero stands for quantities evaluated at the unification scale $`M_G`$. However, $`D=L`$ also predicts that $`m_s^0=m_\mu ^0`$ and $`m_d^0=m_e^0`$. Empirically, one finds instead that $`m_s^0\frac{1}{3}m_\mu ^0`$ and $`m_d^03m_e^0`$. These factors of three are called the Georgi-Jarlskog factors . The simplest possible $`SO(10)`$ Yukawa terms also predict that all the CKM angles vanish, since $`UD`$. While not exactly true, this is at least a good zeroth order relation, since the CKM angles are all small compared to unity. By contrast, in $`SU(5)`$ the matrices $`D`$ and $`U`$ are not related by the unified symmetry and so the CKM angles are unconstrained. The smallness of the CKM angles can be regarded, therefore, as evidence for $`SO(10)`$. On the other hand, the proportionality of $`D`$ and $`U`$ in $`SO(10)`$ also predicts that $`m_c^0/m_t^0=m_s^0/m_b^0`$, which fails badly by over an order of magnitude.
What one can conclude is that a way of going beyond the simplest possible $`SO(10)`$ Yukawa scheme must be found which preserves some of its predictions while breaking others. One way to do this involves using larger representations to break the electroweak interactions. For instance, in the original Georgi-Jarlskog model, a $`\overline{\mathrm{𝟒𝟓}}`$ multiplet of $`SU(5)`$ (not to be confused with the adjoint of $`SO(10)`$) participates in breaking $`SU(2)_L\times U(1)_Y`$. In the context of $`SO(10)`$, this $`\overline{\mathrm{𝟒𝟓}}`$ is contained in a $`\overline{\mathrm{𝟏𝟐𝟔}}`$, which is inconsistent with our minimality assumptions. More economical is to assume that the Higgs fields that break $`SO(10)`$ at the unification scale, i.e., the $`\mathrm{𝟒𝟓}_H+\mathrm{𝟏𝟔}_H+\overline{\mathrm{𝟏𝟔}}_H`$, couple to quarks and leptons and thus introduce the effects of that $`SO(10)`$ breaking into the quark and lepton mass relations. This is the assumption we make.
To describe the third family it is simplest to assume the minimal Yukawa term $`\mathrm{𝟏𝟔}_3\mathrm{𝟏𝟔}_3\mathrm{𝟏𝟎}_H`$ as pictured in Fig. 1(a). By itself, this would make all the mass matrices have the form
$$\left(\begin{array}{ccc}0& 0& 0\\ 0& 0& 0\\ 0& 0& 1\end{array}\right).$$
(8)
That would give the following predictions, all of which are at least good zeroth approximations to reality: $`m_b^0=m_\tau ^0`$, $`V_{cb}=0`$, and $`m_1/m_3=m_2/m_3=0`$, where $`m_i`$ is a mass of a fermion of the $`i^{th}`$ family. Note that these are just the “good” $`SO(10)`$ predictions mentioned above.
The second family presents more of a challenge. The main issue is how to get the Georgi-Jarlskog factor of 3 between $`m_\mu ^0`$ and $`m_s^0`$. Breaking of $`SU(5)`$ must be involved, since the bad relation $`m_s^0=m_\mu ^0`$ arises already at the $`SU(5)`$ level. The only field that breaks $`SU(5)`$ in the framework of minimal $`SO(10)`$ breaking is the adjoint, $`\mathrm{𝟒𝟓}_H`$. Since $`\mathrm{𝟒𝟓}_HBL`$, and the $`BL`$ of leptons is $`3`$ times that of quarks, this field has the possibility of giving the needed Georgi-Jarlskog factor. Thus one must seek an effective Yukawa term that involves the $`\mathrm{𝟒𝟓}_H`$. The simplest such term , in the sense of the term of lowest dimension, is of the form $`(\mathrm{𝟏𝟔}_i\mathrm{𝟏𝟔}_j)\mathrm{𝟏𝟎}_H\mathrm{𝟒𝟓}_H/M_G`$. Moreover, this term can arise in a simple way by the integration out of a $`\mathrm{𝟏𝟔}+\overline{\mathrm{𝟏𝟔}}`$ family plus antifamily at the unification scale, as shown in Fig. 1(b).
There are actually two ways to contract the $`SO(10)`$ indices of such a term: the product $`(\mathrm{𝟏𝟔}_i\mathrm{𝟏𝟔}_j)`$ can be contracted symmetrically or antisymmetrically. It is easy to show that if $`\mathrm{𝟒𝟓}BL`$, only the antisymmetric contraction contributes to the quark and lepton mass matrices. (The reason is simple. If the VEV of the adjoint is proportional to a generator $`Q`$, then the symmetric/antisymmetric contractions give contributions to fermion masses that go as $`Q(f)\pm Q(\overline{f})`$. Since $`BL`$ of an antifermion is minus that of the fermion, the contribution cancels for the symmetric contraction.) Thus, one need only consider the flavor-antisymmetric term, which means only $`ij=23`$ and not $`ij=22`$, or $`33`$. Consequently, the only operator of interest is $`(\mathrm{𝟏𝟔}_2\mathrm{𝟏𝟔}_3)\mathrm{𝟏𝟎}_H\mathrm{𝟒𝟓}_H`$ which, together with the operator $`\mathrm{𝟏𝟔}_3\mathrm{𝟏𝟔}_3\mathrm{𝟏𝟎}_H`$, gives
$$\begin{array}{cc}U=\left(\begin{array}{ccc}0& 0& 0\\ 0& 0& ϵ/3\\ 0& ϵ/3& 1\end{array}\right)M_U,\hfill & D=\left(\begin{array}{ccc}0& 0& 0\\ 0& 0& ϵ/3\\ 0& ϵ/3& 1\end{array}\right)M_D,\hfill \\ & \\ N=\left(\begin{array}{ccc}0& 0& 0\\ 0& 0& ϵ\\ 0& ϵ& 1\end{array}\right)M_U,\hfill & L=\left(\begin{array}{ccc}0& 0& 0\\ 0& 0& ϵ\\ 0& ϵ& 1\end{array}\right)M_D.\hfill \end{array}$$
(9)
The desired factor of 3 has been achieved between leptons and quarks, due to the generator $`BL`$ to which the adjoint VEV is proportional. One also can see that the $`ϵ`$ entries are flavor antisymmetric for reasons already explained. As they stand, these forms of the matrices are inadequate to explain even the features of the second and third families of fermions. There are three inadequacies. (1) The factor of 3 comes in squared between the mass of the leptons and quarks of the second family. The reason is that, for $`ϵ`$ small due to the mass hierarchy between families, the second eigenvalue of $`L`$ is given by the seesaw formula $`m_\mu ^0ϵ^2M_D`$, while the second eigenvalue of $`D`$ is given by $`m_s^0(ϵ/3)^2M_D`$. (2) The matrices $`D`$ and $`U`$ are still exactly proportional. This is a consequence of the fact that the generator $`BL`$ does not distinguish up and down quarks. Therefore, the CKM angle $`V_{cb}`$ still exactly vanishes. (3) Because $`D`$ and $`U`$ are exactly proportional, one still has the bad prediction $`m_c^0/m_t^0=m_s^0/m_b^0`$.
It is clear that the breaking of $`SO(10)`$ due to the adjoint cannot cure all of these problems, since $`BL`$ does not distinguish $`D`$ from $`U`$. Thus the breaking of $`SO(10)`$ done by $`\mathrm{𝟏𝟔}_H+\overline{\mathrm{𝟏𝟔}}_H`$ must come into play. As we shall now show, a single simple operator exists, which involves one of these spinor Higgs and cures at one stroke all three of the problems we have identified.
The lowest dimension effective Yukawa operators that involve the spinor Higgs fields are quartic in spinors. Consider, therefore, operators of the form $`\mathrm{𝟏𝟔}_i\mathrm{𝟏𝟔}_j\mathrm{𝟏𝟔}_{}^{}{}_{H}{}^{}\mathrm{𝟏𝟔}_H/M_G`$. The $`\mathrm{𝟏𝟔}_H`$ is the spinor Higgs field that breaks $`SO(10)`$ at $`M_G`$ down to $`SU(5)`$. The $`\mathrm{𝟏𝟔}_{}^{}{}_{H}{}^{}`$ is a spinor Higgs that has a weak-scale VEV that breaks $`SU(2)_L\times U(1)_Y`$. In principle, these two spinors could be the same field. However, if they were, it would mean that they had to be contracted symmetrically by Bose statistics, which in turn would mean that $`\mathrm{𝟏𝟔}_i`$ and $`\mathrm{𝟏𝟔}_j`$ would also have to be contracted symmetrically. A careful examination shows that the resulting flavor-symmetric contributions to the mass matrices do not lead to realistic forms, though it is possible to achieve realistic mass matrices by adding yet another Yukawa operator, as in the interesting model of Babu, Pati, and Wilczek . Therefore $`\mathrm{𝟏𝟔}_𝐇^{}`$ must be a distinct field. As will be seen later, introducing this $`\mathrm{𝟏𝟔}_H^{}`$ involves no loss of economy, since it allows a very elegant explanation of the largeness of the ratio $`m_t/m_b`$ without making $`\mathrm{tan}\beta `$ large.
There are still several operators of this type to be considered: the family indices can take the values $`ij=33,22,23`$, or $`32`$, and there are three ways to contract the four spinors to make an $`SO(10)`$ singlet. Here again, one must examine the various cases to see which gives the most realistic mass matrices. As it turns out, there is one operator that is much superior to the others, in the sense that it much more cleanly and simply fits the data. It is of the form $`[\mathrm{𝟏𝟔}_2\mathrm{𝟏𝟔}_H][\mathrm{𝟏𝟔}_3\mathrm{𝟏𝟔}_{}^{}{}_{H}{}^{}]`$, where $`[\mathrm{}]`$ means that the spinors inside are contracted into a $`\mathrm{𝟏𝟎}`$. This can arise very simply by integrating out a $`\mathrm{𝟏𝟎}`$ of fermions, as shown in Fig. 1(c).
Let us write the resulting mass operator in $`SU(5)`$ language. Denote by $`𝐩(𝐪)`$ a $`𝐩`$ multiplet of $`SU(5)`$ that is contained in a $`𝐪`$ multiplet of $`SO(10)`$. The VEV of $`\mathrm{𝟏𝟔}_H`$ lies, of course, in the $`\mathrm{𝟏}(\mathrm{𝟏𝟔})`$ direction, while the VEV of $`\mathrm{𝟏𝟔}_H^{}`$ that breaks the weak interactions lies in the $`\overline{\mathrm{𝟓}}(\mathrm{𝟏𝟔})`$ direction. Therefore, the resulting mass term is of the form $`[\overline{\mathrm{𝟓}}(\mathrm{𝟏𝟔}_2)\mathrm{𝟏}(\mathrm{𝟏𝟔}_H)][\mathrm{𝟏𝟎}(\mathrm{𝟏𝟔}_3)\overline{\mathrm{𝟓}}(\mathrm{𝟏𝟔}_H^{})]`$, which in $`SU(5)`$ terms gives effectively the operator $`(\overline{\mathrm{𝟓}}_2\mathrm{𝟏𝟎}_3)\overline{\mathrm{𝟓}}_H`$. Note that this has the same form as the $`SU(5)`$ operator discussed in the last Section, which gave the $`\sigma `$ entries in Eq. (6). The result, then, of including this operator is to make the mass matrices take the form:
$$\begin{array}{cc}U=\left(\begin{array}{ccc}0& 0& 0\\ 0& 0& ϵ/3\\ 0& ϵ/3& 1\end{array}\right)M_U,\hfill & D=\left(\begin{array}{ccc}0& 0& 0\\ 0& 0& \sigma +ϵ/3\\ 0& ϵ/3& 1\end{array}\right)M_D,\hfill \\ & \\ N=\left(\begin{array}{ccc}0& 0& 0\\ 0& 0& ϵ\\ 0& ϵ& 1\end{array}\right)M_U,\hfill & L=\left(\begin{array}{ccc}0& 0& 0\\ 0& 0& ϵ\\ 0& \sigma +ϵ& 1\end{array}\right)M_D\hfill \end{array}.$$
(10)
The new term has given the entries we call $`\sigma `$. Note that these lopsided entries appear only in $`D`$ and $`L`$. The reason is simply that the $`\mathrm{𝟏𝟔}_H^{}`$ contains a $`\overline{\mathrm{𝟓}}`$ of $`SU(5)`$ but no $`\mathrm{𝟓}`$.
It is easy to see that the new term with $`\sigma ϵ`$ cures at once all three of the problems we identified with the forms given in Eq. (9): (1) Instead of $`m_\mu ^0ϵ^2M_D`$ and $`m_s^0(ϵ/3)^2M_D`$, one has approximately that $`m_\mu ^0(ϵ)(\sigma +ϵ)ϵ\sigma `$ and $`m_s^0(ϵ/3)(\sigma +ϵ/3)ϵ\sigma /3`$. More exact expressions will be given later. Thus the desired Georgi-Jarlskog factor of 1/3 is obtained, instead of 1/9. The $`\sigma `$ entry has dominated over one of the factors of $`ϵ/3`$ and thus prevented the factor of 1/3 from coming in squared.
(2) The $`\sigma `$ entry comes into $`D`$ but not $`U`$, and thus breaks the proportionality of the two matrices. As a result, $`V_{cb}`$ no longer vanishes, but is given approximately by $`(ϵ/3)(\frac{\sigma ^2}{\sigma ^2+1})`$. Note that this is of the same order in $`ϵ`$ as $`m_s/m_b(ϵ/3)(\frac{\sigma }{\sigma ^2+1})`$, rather than $`\sqrt{m_s/m_b}`$ as is the case with Fritzsch forms, and accords much better with the actual experimental values.
(3) The fact that $`\sigma `$ breaks the proportionality of $`U`$ and $`D`$ also means that the bad relation $`m_c^0/m_t^0=m_s^0/m_b^0`$ is broken. Specifically, $`m_s^0/m_b^0`$ is of order $`ϵ`$, while $`m_c^0/m_t^0`$ is still of order $`ϵ^2`$ and therefore much smaller. This also accords well with the experimental numbers. In fact, as we shall see, if one uses $`V_{cb}`$ and $`m_\mu /m_\tau `$ to fix the parameters $`\sigma `$ and $`ϵ`$, one finds that $`m_c`$(1 GeV) is predicted to be in agreement with the experimentally determined value of $`1.27\pm 0.1`$ GeV. It should also be noted that the prediction $`m_b^0=m_\tau ^0`$ is only very slightly affected by the addition of the $`\sigma `$ term, both $`m_b^0`$ and $`m_\tau ^0`$ being given to leading order in $`ϵ`$ by $`\sqrt{\sigma ^2+1}M_D`$.
The economy of the above mass matrix forms is seen in the fact that five quantities ($`V_{cb}`$, $`m_\mu /m_\tau `$, $`m_s/m_b`$, $`m_c/m_t`$, and $`m_\tau /m_b`$) are successfully fit with only the two parameters $`\sigma `$ and $`ϵ`$. No other published form succeeds in accurately reproducing the masses and mixing of the the heavier two families with so few parameters. The predictions and fits will be discussed in detail in Section IV.
We see that the matrices in Eq. (10) were arrived at by a process of reasoning that had nothing to do with the question of neutrino mass but rather with an attempt to get realistic masses and mixings for the quarks and charged leptons using as simple a Higgs sector as possible in $`SO(10)`$. But what has emerged is a structure with precisely the three critical features identified in the last Section as giving a simple explanation of the large mixing of $`\nu _\mu `$ and $`\nu _\tau `$. In fact, a fit of $`V_{cb}`$ and $`m_\mu /m_\tau `$ gives $`\sigma 1.8`$ and $`ϵ0.14`$. Consequently, as can be seen directly from Eq. (10), the angle $`\theta _{\mu \tau }=\mathrm{sin}^1U_{\mu 3}=\mathrm{tan}^1\sigma O(ϵ)=60^{}O(8^{})`$. This is quite consistent with what is observed. We will look more carefully at these predictions later.
To summarize, with two parameters, $`ϵ`$ and $`\sigma `$, four mass ratios and two mixing angles are satisfactorily accounted for, if we include $`U_{\mu 3}`$. No greater economy could be hoped for in explaining the spectrum of the heavy two families. Moreover, as we shall see in the next Section, the forms in Eq. (10) can be extended to include the first family with equal economy: the introduction of two new parameters (one of which is complex) will nicely account for seven quantities pertaining to the first family.
Before explaining how the model is extended to the first family we will expand on a couple of points made earlier. First, we said that the introduction of the $`\mathrm{𝟏𝟔}_H^{}`$ allows a simple explanation of why $`m_tm_b`$ that does not require a $`\mathrm{tan}\beta 1`$. The point is that the Higgs doublet of the MSSM that is often called $`H_U`$ is purely contained in the $`\mathrm{𝟏𝟎}_H`$ that couples to $`\mathrm{𝟏𝟔}_3\mathrm{𝟏𝟔}_3`$ and gives rise to the “1” entry in the mass matrices of Eq. (10). However, the Higgs doublet of the MSSM that is called $`H_D`$ does not come purely from $`\mathrm{𝟏𝟎}_H`$. Rather it is a mixture of doublets in $`\mathrm{𝟏𝟎}_H`$ and $`\mathrm{𝟏𝟔}_H^{}`$, since they both contain $`\overline{\mathrm{𝟓}}`$’s of $`SU(5)`$. Thus we may write
$$\begin{array}{c}H_U=H(\mathrm{𝟏𝟎}_H),\hfill \\ \\ H_D=\overline{H}(\mathrm{𝟏𝟎}_H)\mathrm{cos}\gamma +\overline{H}(\mathrm{𝟏𝟔}_H^{})\mathrm{sin}\gamma ,\hfill \end{array}$$
(11)
where $`\gamma `$ is some mixing angle that depends on the parameters of the Higgs sector. Since the 33 elements of the mass matrices all arise purely from the coupling of the $`\mathrm{𝟏𝟎}_H`$, the parameters we called $`M_U`$ and $`M_D`$ in Eq. (10) are given by
$$\begin{array}{c}M_U=\lambda _{33}H_U,\hfill \\ \\ M_D=\lambda _{33}H_D\mathrm{cos}\gamma .\hfill \end{array}$$
(12)
leading to the ratio $`M_U/M_D=H_U/(H_D\mathrm{cos}\gamma )=\mathrm{tan}\beta /\mathrm{cos}\gamma `$, where $`\mathrm{tan}\beta `$ is defined to be the ratio of the Higgs VEV $`H_U`$ giving mass to the top quark to the Higgs VEV $`H_D`$ giving mass to the bottom quark. Hence
$$m_t^0/m_b^0(\sigma ^2+1)^{1/2}(\mathrm{tan}\beta /\mathrm{cos}\gamma )$$
(13)
The point is simply that the large ratio of the top to bottom masses could be the result of $`\mathrm{cos}\gamma `$ being small rather than $`\mathrm{tan}\beta `$ being large. In fact, since we do not know anything a priori about the angle $`\gamma `$, we cannot say whether $`\mathrm{tan}\beta `$ is large or small. It should be noted that if one assumes that $`H_D`$ lies mostly in the $`\mathrm{𝟏𝟔}_H^{}`$ (so that $`\mathrm{cos}\gamma 1`$), it would explain why the parameter $`\sigma `$ is large since it comes from a coupling to $`\mathrm{𝟏𝟔}_H^{}`$, and also explain why $`\mathrm{tan}\beta `$ might be small.
A second point we wish to underline here has to do with the reasonableness of asymmetric mass matrices. In many models it is assumed that all the mass matrices are symmetric. However, this is not something that is called for by the group theory of grand unification. It is true that with the minimal Yukawa terms $`SU(5)`$ gives a symmetric $`U`$. But even with minimal Yukawa terms $`SU(5)`$ does not predict any symmetry of the $`D`$ and $`L`$ matrices. And in $`SO(10)`$, as we have seen, once one introduces the effects of $`SO(10)`$ breaking into the Yukawa sector, as one virtually must, one easily obtains effective Yukawa terms that are asymmetric. Fig. 1(c) shows that very simple diagrams can give terms that are lopsided, in the sense that they contribute only above or below the diagonal. From the point of view of the fundamental grand unified theory, then, lopsided terms are as natural as symmetric ones. The preference for symmetric terms has been the result not of examining what kinds of terms are obtained in a simple way in unification, but rather from the desire to reduce the number of parameters at the level of mass matrices with the aim of making models which are highly predictive. However, putting oneself in the straightjacket of symmetric matrices makes it hard to get a good fit to all the quark and lepton masses and mixings. It turns out, as we have seen, and will see further below, that allowing asymmetric matrices makes possible a model which gives both a very good fit to the data and is actually much more predictive than most models which assume symmetric matrices.
## IV EXTENSION to the FIRST FAMILY
In arriving at the form of the mass matrices for the heavy two families we were limited in the choices that were possible by the assumption we made about the simplicity of the $`SO(10)`$-breaking sector. In extending to the first family we are not quite so limited. Nevertheless, the number of simple possibilities is not very large. There are several discrete choices: Should the contributions to the first row and column of the mass matrices be flavor symmetric like the 1’s in Eq. 10, antisymmetric like the $`ϵ`$’s, or lopsided like the $`\sigma `$’s? Should they contribute to all the matrices equally like the 1’s, to all the matrices but with non-trivial Clebsch factors like the $`ϵ`$’s or only to $`D`$ and $`L`$ like the $`\sigma `$’s? It is fairly easy to run through the various cases and see what kinds of relations among masses and mixings result. As it turns out, one of the simplest possibilities gives a remarkably good fit to the data. This uniquely simple choice is the following:
$$\begin{array}{cc}U=\left(\begin{array}{ccc}\eta & 0& 0\\ 0& 0& ϵ/3\\ 0& ϵ/3& 1\end{array}\right)M_U,\hfill & D=\left(\begin{array}{ccc}\eta & \delta & \delta ^{}\\ \delta & 0& \sigma +ϵ/3\\ \delta ^{}& ϵ/3& 1\end{array}\right)M_D,\hfill \\ & \\ N=\left(\begin{array}{ccc}\eta & 0& 0\\ 0& 0& ϵ\\ 0& ϵ& 1\end{array}\right)M_U,\hfill & L=\left(\begin{array}{ccc}\eta & \delta & \delta ^{}\\ \delta & 0& ϵ\\ \delta ^{}& \sigma +ϵ& 1\end{array}\right)M_D.\hfill \end{array}$$
(14)
We have already mentioned that fits give $`\sigma 1.8`$ and $`ϵ0.14`$. The new parameters $`\delta `$ and $`\delta ^{}`$ both have magnitude of about $`0.008`$. The parameter $`\eta `$ is by far the smallest, being about $`8\times 10^6`$. The only role that $`\eta `$ plays in the sector of quarks and charged leptons is in giving the up quark a mass, for it makes negligible contributions to the down quark and electron masses as determined from $`D`$ and $`L`$, respectively. In Fig. 2(a) we have illustrated a higher-order diagram that can contribute to the parameter $`\eta `$. Since it is not excluded that the up quark is exactly massless, it is possible to set $`\eta `$ to zero. In any event, one can see that $`\eta m_u^0/m_t^0`$, which is by orders of magnitude smaller than any other interfamily ratio of masses in the standard model. It will, however, be of some significance for neutrino masses. If $`\eta `$ vanishes, this model gives only the small-angle MSW solution to the solar neutrino problem. But even if $`\eta `$ is as small as $`8\times 10^6`$, it allows either the small-angle MSW solution or bimaximal neutrino mixing to arise in a simple way.
Turning to the parameters $`\delta `$ and $`\delta ^{}`$, we see that they appear symmetrically and only in $`D`$ and $`L`$. Such terms are easily obtained in $`SO(10)`$ from simple diagrams such as that shown in Figs. 2(b) and 2(c). The effective operators arising from these diagrams are of the form $`[\mathrm{𝟏𝟔}_1\mathrm{𝟏𝟔}_j][\mathrm{𝟏𝟔}_H\mathrm{𝟏𝟔}_H^{}]`$ with $`j=2,3`$, where again the spinors in brackets are contracted symmetrically into a $`\mathrm{𝟏𝟎}`$ of $`SO(10)`$ which is integrated out. Note, however, that the symmetric contributions $`\delta `$ and $`\delta ^{}`$ from the two Higgs contraction $`[\mathrm{𝟏}(\mathrm{𝟏𝟔}_H)\overline{\mathrm{𝟓}}(\mathrm{𝟏𝟔}_H^{})]`$ contributes only to $`D`$ and $`L`$ by virtue of their $`SU(5)`$ structure. Contrast these effective operators for $`j=2,3`$ with that occurring previously for the term $`\sigma `$ arising from the diagram shown in Fig. 1(c).
The three new parameters we have introduced are, as we shall see, sufficient to account for everything about the first family. Before proceeding, however, we must be careful about complex phases. It is easy to show that if we allow all parameters of the model to be complex all but two phase angles can be rotated away from the mass matrices $`U`$, $`D`$, $`L`$ and $`N`$, provided we now neglect the negligible $`\eta `$ contributions to $`D`$ and $`L`$. We will call these physical phases $`\alpha `$ and $`\varphi `$ which appear as follows,
$$\begin{array}{cc}U=\left(\begin{array}{ccc}\eta & 0& 0\\ 0& 0& ϵ/3\\ 0& ϵ/3& 1\end{array}\right)M_U,\hfill & D=\left(\begin{array}{ccc}0& \delta & \delta ^{}e^{i(\varphi +\alpha )}\\ \delta & 0& \sigma +ϵe^{i\alpha }/3\\ \delta ^{}e^{i\varphi }& ϵ/3& 1\end{array}\right)M_D,\hfill \\ & \\ N=\left(\begin{array}{ccc}\eta & 0& 0\\ 0& 0& ϵ\\ 0& ϵe^{i\alpha }& 1\end{array}\right)M_U,\hfill & L=\left(\begin{array}{ccc}0& \delta & \delta ^{}e^{i\varphi }\\ \delta & 0& ϵ\\ \delta ^{}e^{i(\varphi +\alpha )}& \sigma +ϵe^{i\alpha }& 1\end{array}\right)M_D,\hfill \end{array}$$
(15)
where in these matrices and henceforth $`ϵ`$, $`\sigma `$, $`\delta `$, $`\delta ^{}`$ and $`\eta `$ denote the magnitudes of these parameters and the phases are written explicitly. The phase $`\alpha `$ only comes into the fits of masses at higher order in the small quantity $`ϵ/\sigma 0.08`$. Numerically, its effect is only a few percent; moreover, the fits (especially to $`m_c`$) prefer a value near zero. Therefore, we can ignore $`\alpha `$ and will do so from now on. That leaves only the phase $`\varphi `$. Its only significant effect, but a very important one, is to give the CP-violating phase angle $`\delta _{CP}`$.
Instead of using the parameters $`\delta `$, $`\delta ^{}`$ and $`e^{i\varphi }`$, it will be somewhat more convenient to use the parameters $`t_L`$, $`t_R`$, and $`e^{i\theta }`$, which are defined in terms of them as follows:
$$t_Le^{i\theta }\frac{\delta \sigma \delta ^{}e^{i\varphi }}{\sigma ϵ/3},$$
(16)
and
$$t_R\frac{\delta \sqrt{\sigma ^2+1}}{\sigma ϵ/3}.$$
(17)
The significance of these parameters is that they are essentially the left-handed and right-handed Cabbibo angles. This can be seen by taking the forms for $`D`$ and $`L`$ given in Eq. (15) and diagonalizing the 2-3 block. When this is done the 1-2 blocks of these matrices take the form
$$D^{[12]}\left(\begin{array}{cc}0& t_R\\ t_L& 1\end{array}\right),L^{[12]}\left(\begin{array}{cc}0& t_L\\ t_R& 3\end{array}\right).$$
(18)
In terms of the five dimensionless parameters $`ϵ`$, $`\sigma `$, $`t_L`$, $`t_R`$, and $`e^{i\theta }`$ with $`\eta `$ set equal to zero, we now write down expressions for fourteen observable quantities: seven ratios of quark and lepton masses, three CKM angles and one phase, and three lepton mixing angles.
$$\begin{array}{ccc}m_b^0/m_\tau ^0\hfill & & 1\frac{2}{3}\frac{\sigma }{\sigma ^2+1}ϵ,\hfill \\ m_c^0/m_t^0\hfill & & \frac{1}{9}ϵ^2[1\frac{2}{9}ϵ^2],\hfill \\ m_\mu ^0/m_\tau ^0\hfill & & ϵ\frac{\sigma }{\sigma ^2+1}[1+ϵ\frac{1\sigma ^2\sigma ϵ}{\sigma (\sigma ^2+1)}+\frac{1}{18}(t_L^2+t_R^2)],\hfill \\ m_s^0/m_b^0\hfill & & \frac{1}{3}ϵ\frac{\sigma }{\sigma ^2+1}[1+\frac{1}{3}ϵ\frac{1\sigma ^2\sigma ϵ/3}{\sigma (\sigma ^2+1)}+\frac{1}{2}(t_L^2+t_R^2)],\hfill \\ m_u^0/m_t^0\hfill & =& 0,\hfill \\ m_e^0/m_\mu ^0\hfill & & \frac{1}{9}t_Lt_R[1ϵ\frac{\sigma ^2+2}{\sigma (\sigma ^2+1)}+ϵ^2\frac{\sigma ^4+9\sigma ^2/2+3}{\sigma ^2(\sigma ^2+1)^2}\frac{1}{9}(t_L^2+t_R^2)],\hfill \\ m_d^0/m_s^0\hfill & & t_Lt_R[1\frac{1}{3}ϵ\frac{\sigma ^2+2}{\sigma (\sigma ^2+1)}(t_L^2+t_R^2)+(t_L^4+t_L^2t_R^2+t_R^4)],\hfill \\ V_{cb}\hfill & & \frac{1}{3}ϵ\frac{\sigma ^2}{\sigma ^2+1}[1+\frac{2}{3}ϵ\frac{1}{\sigma (\sigma ^2+1)}],\hfill \\ V_{us}\hfill & & t_L[1\frac{1}{2}t_L^2t_R^2+t_R^4+\frac{5}{2}t_L^2t_R^2+\frac{3}{8}t_L^4\frac{ϵ}{3\sigma \sqrt{\sigma ^2+1}}\frac{t_R}{t_L}e^{i\theta }],\hfill \\ V_{ub}\hfill & & \frac{1}{3}t_Lϵ\frac{1}{\sigma ^2+1}[\sqrt{\sigma ^2+1}\frac{t_R}{t_L}e^{i\theta }(1\frac{1}{3}ϵ\frac{\sigma }{\sigma ^2+1})(1\frac{2}{3}ϵ\frac{\sigma }{\sigma ^2+1})],\hfill \\ U_{\mu 3}^0\hfill & & \mathrm{sin}\theta _{\mu \tau }\frac{\sigma }{\sqrt{\sigma ^2+1}}+O(ϵ),\hfill \\ U_{e2}^0\hfill & & \mathrm{cos}\theta _{\mu \tau }\left(\frac{1}{3}t_R\right)[1+ϵ(\frac{\mathrm{tan}\theta _{\mu \tau }}{\sigma ^2+1}\frac{t_L}{t_R}e^{i\theta }\frac{(1+\sigma \mathrm{tan}\theta _{\mu \tau })}{\sigma \sqrt{\sigma ^2+1}})\frac{1}{18}t_R^2\frac{1}{9}t_L^2],\hfill \\ U_{e3}\hfill & & \mathrm{tan}\theta _{\mu \tau }U_{e2}[1+ϵ\frac{2}{\mathrm{sin}2\theta _{\mu \tau }}\left(\frac{t_L}{t_R}e^{i\theta }\frac{1}{\sqrt{\sigma ^2+1}}\frac{1}{\sigma ^2+1}\right)],\hfill \end{array}$$
(19)
These expansions have been carried to sufficiently high order in small quantities to be accurate to within 0.2% and are useful in doing the fits to the data. However, the leading terms in these expansions have much simpler forms and thus allow one to see more readily the relationships among various quantities in this model. We therefore write these simpler expressions for purposes of discussion.
$$\begin{array}{ccc}m_b^0/m_\tau ^0\hfill & & 1,\hfill \\ m_c^0/m_t^0\hfill & & \frac{1}{9}ϵ^2,\hfill \\ m_\mu ^0/m_\tau ^0\hfill & & ϵ\frac{\sigma }{\sigma ^2+1},\hfill \\ m_s^0/m_b^0\hfill & & \frac{1}{3}ϵ\frac{\sigma }{\sigma ^2+1},\hfill \\ m_u^0/m_t^0\hfill & =& 0,\hfill \\ m_e^0/m_\mu ^0\hfill & & \frac{1}{9}t_Lt_R,\hfill \\ m_d^0/m_s^0\hfill & & t_Lt_R,\hfill \\ V_{cb}^0\hfill & & \frac{1}{3}ϵ\frac{\sigma ^2}{\sigma ^2+1},\hfill \\ V_{us}^0\hfill & & t_L,\hfill \\ V_{ub}^0\hfill & & \frac{1}{3}t_Lϵ\frac{1}{\sigma ^2+1}(\sqrt{\sigma ^2+1}\frac{t_R}{t_L}e^{i\theta }1),\hfill \\ U_{\mu 3}^0\hfill & & \mathrm{sin}\theta _{\mu \tau }\frac{\sigma }{\sqrt{\sigma ^2+1}}+O(ϵ)0.7,\hfill \\ U_{e2}^0\hfill & & \mathrm{cos}\theta _{\mu \tau }\left(\frac{1}{3}t_R\right),\hfill \\ U_{e3}\hfill & & \mathrm{sin}\theta _{\mu \tau }\left(\frac{1}{3}t_R\right).\hfill \end{array}$$
(20)
It might at first seem surprising that without any information about the Majorana mass matrix $`M_R`$ of the right-handed neutrinos we are able to write down predictions for the three neutrino mixing angles. However, if $`\eta =0`$, as we are assuming at present, then the Dirac mass matrix of the neutrinos ($`N`$) has vanishing first row and column, and therefore, obviously, the same will be true of the mass matrix of the light neutrinos, which is given by the well-known “see-saw” formula $`M_\nu =N^TM_R^1N`$. This means that the two mixing elements of the electron neutrino, $`U_{e2}`$ and $`U_{e3}`$, get no contribution from diagonalizing $`M_\nu `$, but come entirely from diagonalizing $`L`$. Since $`L`$ is a known matrix in our model, these two mixing elements are predicted. In the case of the mixing of the mu and tau neutrinos, $`U_{\mu 3}`$ does receive a contribution from diagonalizing $`M_\nu `$. However, as can be seen from the form of $`N`$ this is an effect of $`O(ϵ)`$. The contribution to $`U_{\mu 3}`$ coming from diagonalizing $`L`$, on the other hand, is of order unity, since it arises from the large parameter $`\sigma `$. Thus $`U_{\mu 3}`$ is predicted, although not precisely.
Since we have written fourteen quantities in terms of five parameters, there are altogether nine predictions of the model. Which quantities one takes as “predicted” depends on which quantities are used to determine the values of the parameters. We will use the lepton mass ratios and the angles $`V_{cb}`$ and $`V_{us}`$ for this purpose as they are the best measured. As one can see from the third and eighth of Eqs. (20), one can get the value $`\sigma `$ from the ratio $`3V_{cb}^0/(m_\mu ^0/m_\tau ^0)`$. One finds (of course, taking the renormalization effects into account as was done in ) that numerically $`\sigma \sqrt{3}`$. Substituting this into the expression for $`m_\mu ^0/m_\tau ^0`$, one obtains that $`ϵ0.14`$. This is the small parameter of the model that is responsible for the hierarchy between the second and third families, and is small enough that the expressions in Eqs. (20) are fairly accurate. One can use $`m_e/m_\mu `$ and $`V_{us}`$ to determine $`t_L`$ and $`t_R`$. A careful fit, described later, gives $`t_L=0.236`$ and $`t_R=0.205`$. That $`t_Lt_R`$ is easily understood from Eq. (18) and the well-known Weinberg-Wilczek-Zee-Fritzsch result that the Cabbibo angle is well acounted for by symmetric mass matrices for the first two families; cf. Eq. (1). The near equality of $`t_L`$ and $`t_R`$ is also apparant from the seventh and ninth relations of Eqs. (20) and the fact that numerically $`V_{us}\sqrt{m_d/m_s}`$. The phase factor $`e^{i\theta }`$ will be determined from the CP-violating phase $`\delta _{CP}`$.
The nine predictions, then, are the following. To begin with, there are the three famous predictions, (1) $`m_b^0/m_\tau ^01`$, (2) $`m_s^0\frac{1}{3}m_\mu ^0`$, and (3) $`m_d^03m_e^0`$. The first is the “good” prediction of minimal $`SU(5)`$ unification, and the latter two are the Georgi-Jarlskog relations. These predictions are manifest from the first, third, fourth, sixth, and seventh of Eqs. (20). It is hardly surprising that the model gives these relations, since we were guided by them in constructing the model. The fourth prediction is (4) $`m_u^0/m_t^0=0`$. Even if the $`u`$ quark is not exactly massless this relation is a very good approximation to reality. If one takes the favored value of $`m_u4`$ MeV, then, with reasonable assumptions about thresholds in doing the running up to the GUT scale, one obtains $`\eta m_u^0/m_t^08\times 10^6`$. This is far smaller than any other interfamily ratio of masses. For instance, the comparable ratio for down quarks is $`m_d^0/m_b^010^3`$, and for leptons is $`m_e^0/m_\tau ^03\times 10^4`$. Like the previous three relations, $`m_u0`$ is a reflection of basic group-theoretical aspects of the model. It comes from the fact, explained above, that the $`\delta `$ and $`\delta ^{}`$ entries only appear in $`D`$ and $`L`$.
The remaining five predictions are not simple group-theoretical relations like the foregoing, but are non-trivial quantitative predictions. They are predictions for (5) $`m_c`$, (6) $`V_{ub}`$, (7) $`U_{\mu 3}`$, (8) $`U_{e2}`$, and (9) $`U_{e3}`$.
The prediction for $`m_c`$ is particularly interesting. We see immediately that, for reasons having to do with the group-theoretic structure of the model, the ratio $`m_c^0/m_t^0`$ is much less than the corresponding ratio $`m_s^0/m_b^0`$ for the down quarks because it is of higher order in the small parameter $`ϵ`$. This is a highly significant success, because the minimal Yukawa terms of $`SO(10)`$ notoriously give these ratios to be equal. Moreover, the success is not merely a qualitative one. When $`ϵ`$ and $`\sigma `$ are fit (using $`V_{cb}`$ and $`m_\mu /m_\tau `$) and the renormalization effects are later taken into account, it is found that $`m_c`$ comes out within about 5% of the experimentally preferred value, which is quite remarkable given the various experimental and theoretical uncertainties. This success is non-trivial, because the reasoning that led to the forms of the mass matrices did not depend upon the value of $`m_c`$, and hence it could have been expected that $`m_c`$ would come out wrong by a large factor.
Another non-trivial quantitative hurdle for the model is the prediction for $`V_{ub}`$. The eighth, ninth, and tenth relations of Eqs. (20) give $`V_{ub}^0V_{us}^0V_{cb}^0\frac{1}{\sigma ^2}(\sqrt{\sigma ^2+1}\frac{t_R}{t_L}e^{i\theta }1)`$. If we use the facts that $`\sigma \sqrt{3}`$ and $`t_Lt_R`$, this gives $`V_{ub}V_{us}V_{cb}(\frac{2}{3}e^{i\theta }\frac{1}{3})`$. A careful fit gives
$$V_{ub}=V_{us}V_{cb}(0.558e^{i\theta }0.315).$$
(21)
In other words the model predicts that $`V_{ub}`$ should lie on a certain circle in the complex plane. As can be seen from Fig. 3, the circle for $`V_{ud}V_{ub}^{}`$ slices neatly through the middle of the presently allowed region. Again, this is a very significant success, since the reasoning that led to the forms in Eq. (15) was not based on the value of $`V_{ub}`$.
The prediction for the mixing of $`\nu _\mu `$ and $`\nu _\tau `$ has already been discussed. It is one of the key successes of this model that this mixing turns out to be nearly maximal. The fact that $`\sigma \sqrt{3}`$ tells us that the first term in the expression for $`U_{\mu 3}^0`$ in Eq. (20) corresponds to an angle near $`\pi /3`$. As we shall see in the next Section, the $`O(ϵ)`$ corrections easily bring this down close to the maximal mixing value of $`\pi /4`$.
The prediction of this model for the mixing of $`\nu _e`$ and $`\nu _\mu `$ with $`\eta =0`$ is quite interesting. From the sixth relation of Eqs. (20) and the fact that $`t_Lt_R`$, one sees that $`\frac{1}{3}t_R\sqrt{m_e/m_\mu }`$. Thus the model predicts that $`U_{e2}\mathrm{cos}\theta _{\mu \tau }\sqrt{m_e/m_\mu }`$. The factor of $`\mathrm{cos}\theta _{\mu \tau }`$ is crucial since without it one would have $`\mathrm{sin}^22\theta _{\mathrm{solar}}=4|U_{e1}|^2|U_{e2}|^24(m_e/m_\mu )2\times 10^2`$, which is about twice the value needed for the small-angle MSW solution to the solar neutrino problem. Since atmospheric neutrino data tells us that $`\mathrm{cos}\theta _{\mu \tau }1/\sqrt{2}`$, the model gives just the correct value for the small-angle MSW solution.
In the future both $`V_{ub}`$ and $`\mathrm{sin}^22\theta _{\mathrm{solar}}`$ will be known better and will provide a sharp test of the model. The theoretical uncertainties in the predictions for $`V_{ub}`$ and $`U_{e2}`$ are estimated to be only a few percent.
In discussing the $`\nu _e\nu _\mu `$ mixing above, we have assumed that $`\eta =0`$. If $`\eta `$ does not vanish, but is around $`8\times 10^6`$, corresponding to $`m_u4.5`$ MeV, then it turns out that both the small-angle MSW solution we just discussed and bimaximal mixing are possible. This will be discussed in detail in the next Section.
Finally, there is the prediction of $`\nu _e\nu _\tau `$ mixing. One sees from Eq. (20) that there is a prediction that $`U_{e3}\mathrm{tan}\theta _{\mu \tau }U_{e2}0.05`$. It is interesting that even for the bimaximal mixing case that will be discussed in the next Section, the numerical value of $`U_{e3}`$ is virtually unaffected. Thus this prediction is a “robust” one of this model.
## V NEUTRINO MIXING
### A Mixing of $`\nu _\mu \nu _\tau `$
In analyzing the predictions of this model for $`\nu _\mu \nu _\tau `$ mixing, we may make the approximation that $`\eta =0`$. This means that $`N`$ has vanishing first row and column. Therefore, in computing $`M_\nu =N^TM_R^1N`$ the first row and column of $`M_R^1`$ are irrelevant. Thus we may write $`M_R^1`$ as
$$M_R^1=\left(\begin{array}{ccc}& & \\ & X& Y\\ & Y& Z\end{array}\right),$$
(22)
where $`X`$, $`Y`$, and $`Z`$ are in general complex. There are consequently five real parameters (the over all phase does not matter) that come into the masses and mixing of $`\nu _\mu `$ and $`\nu _\tau `$ from $`M_R`$. As observed earlier, this does not prevent us from making a qualitative prediction for the mixing parameter $`U_{\mu 3}`$, since the contribution to it from diagonalizing the mass matrix $`M_\nu `$ is only of order $`ϵ`$ and $`U_{\mu 3}`$ comes predominantly from diagonalizing the known matrix $`L`$. However, in order to see if more precise predictions can be obtained, we shall look at two simple special cases:
$$(I)M_R=\left(\begin{array}{ccc}& 0& 0\\ 0& 0& Be^{i\beta }ϵ\\ 0& Be^{i\beta }ϵ& 1\end{array}\right)\mathrm{\Lambda }_Re^{i\gamma },$$
(23)
and
$$(II)M_R=\left(\begin{array}{ccc}& 0& 0\\ 0& Be^{i\beta }ϵ^2& 0\\ 0& 0& 1\end{array}\right)\mathrm{\Lambda }_Re^{i\gamma }.$$
(24)
In these cases only three parameters in $`M_R`$, namely $`\mathrm{\Lambda }_R`$, $`B`$ and $`e^{i\beta }`$, contribute to the neutrino observables of the second and third families, since $`ϵ`$ has appeared previously and is used as a natural scaling parameter. In the first case,
$$M_\nu ^I=\left(\begin{array}{ccc}0& 0& 0\\ 0& 0& ϵ\\ 0& ϵ& 2+B^1e^{i\beta }\end{array}\right)\frac{M_U^2}{B\mathrm{\Lambda }_R}e^{i(\beta +\gamma )}.$$
(25)
The neutrino mixing matrix $`U`$, now known as the MNS mixing matrix , is given by $`U=U_L^{}U_\nu `$, where $`U_L`$ is the unitary matrix that diagonalizes $`L^{}L`$, and $`U_\nu `$ is the unitary matrix that diagonalizes $`M_\nu ^{}M_\nu `$. For case I, $`U_\nu `$ is given by
$$U_\nu =\left(\begin{array}{ccc}1& 0& 0\\ 0& \mathrm{cos}\theta _{23}^\nu & \mathrm{sin}\theta _{23}^\nu \\ 0& \mathrm{sin}\theta _{23}^\nu & \mathrm{cos}\theta _{23}^\nu \end{array}\right),$$
(26)
where $`\mathrm{tan}2\theta _{23}^\nu =2ϵ/K`$, and $`K2+B^1e^{i\beta }`$. The ratio of eigenvalues of $`M_\nu `$ gives $`m_{\nu _2}/m_{\nu _3}(ϵ^2/\left|K\right|^2)(1ϵ^2/\left|K\right|^2+\mathrm{})`$. One can choose $`\mathrm{cos}\theta _{23}^\nu `$ to be real, and one can write
$$\mathrm{sin}\theta _{23}^\nu \sqrt{m_{\nu _2}/m_{\nu _3}}e^{i\xi }\left[1+(\frac{1}{2}e^{2i\xi })\frac{m_{\nu _2}}{/m_{\nu _3}}\right].$$
(27)
where $`e^{i\xi }`$ is the phase of $`K`$. One readily sees from the form of the charged-lepton mass matrix $`L`$ in Eq. (15) that $`\mathrm{sin}\theta _{23}^L(U_L)_{23}`$ is given by $`\mathrm{tan}2\theta _{23}^L=\frac{2(\sigma +ϵ)}{\sigma ^21+2\sigma ϵ}=\frac{2\sigma }{\sigma ^21}+O(ϵ)`$. Since $`\sigma \sqrt{3}`$ it is evident that $`\theta _{23}^L60^{}`$. Using the best fit values of $`\sigma `$ and $`ϵ`$ one finds, more precisely, that $`\theta _{23}^L63^{}`$.
Altogether, then, the mixing parameter of $`\nu _\mu `$ and $`\nu _\tau `$ is given by
$$\begin{array}{ccc}U_{\mu 3}\hfill & & \mathrm{sin}\theta _{\mu \tau }\hfill \\ & & \\ & =& \mathrm{sin}\theta _{23}^L\mathrm{cos}\theta _{23}^\nu +\mathrm{cos}\theta _{23}^L\mathrm{sin}\theta _{23}^\nu \hfill \\ & & \\ & & 0.898(1m_{\nu _3}/m_{\nu _2})+0.441\sqrt{m_{\nu _3}/m_{\nu _2}}e^{i\xi }\hfill \end{array}$$
(28)
If neutrino masses are hierarchical, and atmospheric neutrino oscillations are $`\nu _\mu \nu _\tau `$ oscillations, then $`m_{\nu _3}0.06`$ eV. And if one further assumes the small-angle MSW solution to the solar neutrino problem, then $`m_{\nu _2}0.003`$ eV. Thus, $`m_{\nu _2}/m_{\nu _3}0.05`$, within a factor of two or so. Taking it to have the value 0.05, and the phase $`\xi `$ to vanish, Eq. (28) gives $`U_{\mu 3}0.756`$, and $`\mathrm{sin}^22\theta _{\mu \tau }0.984`$. With the same value of the neutrino mass ratio and $`\xi `$ taken to be $`\pi /4`$, $`\mathrm{sin}^22\theta _{\mu \tau }0.943`$. We see that there is excellent agreement with the experimental limits from SuperKamiokande if the complex phase is not too large. But if $`\xi =\pi /2`$, with the same mass ratio, $`\mathrm{sin}^22\theta _{\mu \tau }0.77`$.
The value of $`m_{\nu _2}/m_{\nu _3}=0.05`$ corresponds to $`|K|=0.63`$. Since $`B^1e^{i\beta }=|K|e^{i\xi }2`$, for $`\xi =0`$ this gives $`B=0.73`$, or $`Bϵ=0.1`$. In other words, no very great hierarchy is required in $`M_R`$.
Turning now to case II, we have that
$$M_\nu ^{II}=\left(\begin{array}{ccc}0& 0& 0\\ 0& ϵ^2& ϵ\\ 0& ϵ& 1+B^1e^{i\beta }\end{array}\right)\frac{M_U^2}{\mathrm{\Lambda }_R}e^{i\gamma }.$$
(29)
Consequently, for this case
$$\mathrm{tan}2\theta _{23}^\nu 2ϵ/K^{},$$
(30)
where $`K^{}1+B^1e^{i\beta }`$. The ratio of eigenvalues of $`M_\nu `$ gives $`m_{\nu _2}/m_{\nu _3}ϵ^2\sqrt{1|K^{}1|^2}/|K^{}|^2`$. If we take, $`m_{\nu _2}/m_{\nu _3}=0.05`$, as before, and assume that $`K^{}`$ is real, we have that $`K^{}0.6`$. This gives $`\theta _{23}^\nu 12.3^{}`$, $`\theta _{\mu \tau }50.5^{}`$, and $`\mathrm{sin}^22\theta _{\mu \tau }0.96`$. Again, there is good agreement with the SuperKamiokande results. Moreover, since this value of $`K^{}`$ corresponds to $`B=2.5`$, or $`Bϵ^20.05`$, we see that in this case also no great hierarchy is needed in $`M_R`$.
The foregoing discussion is all based on the assumption that the mixing with the first family is small, so that one has the small-angle MSW solution to the solar neutrino problem. This will certainly be the case if $`\eta =0`$. As we will now see, if instead $`\eta 8\times 10^6`$, as needed to have $`m_u4.5`$ MeV, either the small mixing of $`\nu _e`$ that we have been considering or large mixing of $`\nu _e`$ is possible, depending on the form of $`M_R`$.
### B Mixing of the First Family
In the previous discussion, we set $`\eta =0`$ in which case, no matter what the form of $`M_R`$, the matrix $`M_\nu =N^TM_R^1N`$ has vanishing first row and column, and the matrix $`U_\nu `$ that diagonalizes $`M_\nu ^{}M_\nu `$ has the form of Eq. (26). It is easy to show that the matrix $`U_L`$ which diagonalizes $`L^{}L`$ has the form
$$U_L\left(\begin{array}{ccc}\mathrm{cos}\theta _{12}^L& \mathrm{sin}\theta _{12}^L& 0\\ \mathrm{sin}\theta _{12}^L& \mathrm{cos}\theta _{12}^L& 0\\ 0& 0& 1\end{array}\right)\left(\begin{array}{ccc}1& 0& 0\\ 0& \mathrm{cos}\theta _{23}^L& \mathrm{sin}\theta _{23}^L\\ 0& \mathrm{sin}\theta _{23}^L& \mathrm{cos}\theta _{23}^L\end{array}\right),$$
(31)
where $`\mathrm{sin}\theta _{12}^L\frac{1}{3}t_R`$, $`t_R`$ is defined in Eq. (17), and $`\theta _{23}^L`$ is given after Eq. (27). Putting these together, one has that the total mixing matrix of the neutrinos, $`U=U_L^{}U_\nu `$, is
$$U=\left(\begin{array}{ccc}\mathrm{cos}\theta _{12}^L& \mathrm{sin}\theta _{12}^L\mathrm{cos}\theta _{\mu \tau }& \mathrm{sin}\theta _{12}^L\mathrm{sin}\theta _{\mu \tau }\\ \mathrm{sin}\theta _{12}^L& \mathrm{cos}\theta _{12}^L\mathrm{cos}\theta _{\mu \tau }& \mathrm{cos}\theta _{12}^L\mathrm{sin}\theta _{\mu \tau }\\ 0& \mathrm{sin}\theta _{\mu \tau }& \mathrm{cos}\theta _{\mu \tau }\end{array}\right),$$
(32)
where $`\theta _{\mu \tau }=\theta _{23}^L\theta _{23}^\nu `$. This yields the results, already given in Eq. (20), for $`U_{e2}`$ and $`U_{e3}`$.
Now we will consider what happens under what is presumably the more realistic assumption that $`\eta 8\times 10^6`$.
With $`\eta 0`$, there are two basic possibilities to consider. One possibility is that $`M_R`$ has a form in which its 12, 21, 13, and 31 elements all vanish or are negligibly small. If such is the case, then the previous analysis applies, and the mixing of $`\nu _e`$ is due entirely to the matrix $`L`$. The only effect of the parameter $`\eta `$ in the lepton sector is then to give $`\nu _1`$ a mass of about $`4\times 10^7`$ eV. The second possibility is that $`M_R`$ does have significant 12, 21 and/or 13, 31 elements. If this is the case then a strikingly different situation can arise , namely “bimaximal” mixing , .
We will first illustrate what happens with a simple example. Consider the following form for $`M_R`$:
$$M_R=\left(\begin{array}{ccc}0& Aϵ^3& 0\\ Aϵ^3& Bϵ^2& 0\\ 0& 0& 1\end{array}\right)\mathrm{\Lambda }_R.$$
(33)
We normalize $`A`$ and $`B`$ by powers of $`ϵ`$ simply for later convenience. The mass matrix of light neutrinos resulting from this form is
$$M_\nu =N^TM_R^1N=\left(\begin{array}{ccc}\frac{\eta ^2}{ϵ^4}\frac{B}{A^2}& 0& \frac{\eta }{ϵ^2}\frac{1}{A}\\ 0& ϵ^2& ϵ\\ \frac{\eta }{ϵ^2}\frac{1}{A}& ϵ& 1\end{array}\right)\frac{M_U^2}{\mathrm{\Lambda }_R}.$$
(34)
One sees that the 2-3 block has vanishing determinant, so that a rotation in the 2-3 plane by an angle $`\theta _{23}^\nu ϵ`$ brings $`M_\nu `$ to the form
$$M_\nu ^{}\left(\begin{array}{ccc}\frac{\eta ^2}{ϵ^4}\frac{B}{A^2}& \frac{\eta }{ϵ}\frac{1}{A}& \frac{\eta }{ϵ^2}\frac{1}{A}\\ \frac{\eta }{ϵ}\frac{1}{A}& 0& 0\\ \frac{\eta }{ϵ^2}\frac{1}{A}& 0& 1\end{array}\right)\frac{M_U^2}{\mathrm{\Lambda }_R}.$$
(35)
This can be put in a more transparent form by a rotation in the 1-3 plane by an angle $`\theta _{13}^\nu \eta /(ϵ^2A)`$. This angle is less than or of order $`10^4`$ and thus negligible, practically speaking, so
$$M_\nu ^{\prime \prime }\left(\begin{array}{ccc}\frac{\eta ^2}{ϵ^4}\frac{(B1)}{A^2}& \frac{\eta }{ϵ}\frac{1}{A}& 0\\ \frac{\eta }{ϵ}\frac{1}{A}& 0& 0\\ 0& 0& 1\end{array}\right)\frac{M_U^2}{\mathrm{\Lambda }_R}.$$
(36)
It is clear that the 11 element, being higher order in $`\eta `$, is likely to be much smaller than the 12 and 21 elements. The condition for this to be the case is that $`A/(B1)>\eta /ϵ^32\times 10^3`$. If this very weak condition is satisfied, then the form of the matrix manifestly corresponds to the situation in which the $`\nu _e`$ and $`\nu _\mu `$ together form a pseudo-Dirac pair. That in turn would mean that the mixing of these two neutrinos is very close to maximal.
One sees from Eq. (36) that $`m_{\nu _3}=M_U^2/\mathrm{\Lambda }_R`$, and that the splitting between $`m_{\nu _1}`$ and $`m_{\nu _2}`$ is given by $`\mathrm{\Delta }m_{21}^22(\eta ^3/ϵ^5)((B1)/A^3)(M_U^2/\mathrm{\Lambda }_R)^2`$. For the vacuum solution to the solar neutrino problem, one has $`\mathrm{\Delta }m_{21}^210^{10}`$ eV<sup>2</sup>, so that $`\mathrm{\Delta }m_{21}^2/m_{\nu _3}^23\times 10^82(\eta ^3/ϵ^5)(B1)/A^3`$. This gives $`A(B1)^{1/3}0.06`$. Thus no great hierarchy is required in $`M_R`$ to get the vacuum oscillation solution. The reason for this is that in this scheme the smallness of $`\mathrm{\Delta }m_{21}^2`$ is due to the extreme smallness of the parameter $`\eta `$, which is equal to the ratio $`m_u/m_t`$.
It is easy to see from what has already been said that the matrix $`U_\nu `$ needed to diagonalize $`M_\nu ^{}M_\nu `$ is of the form
$$U_\nu \left(\begin{array}{ccc}1& 0& 0\\ 0& \mathrm{cos}\theta _{23}^\nu & \mathrm{sin}\theta _{23}^\nu \\ 0& \mathrm{sin}\theta _{23}^\nu & \mathrm{cos}\theta _{23}^\nu \end{array}\right)\left(\begin{array}{ccc}1/\sqrt{2}& 1/\sqrt{2}& 0\\ 1/\sqrt{2}& 1/\sqrt{2}& 0\\ 0& 0& 1\end{array}\right).$$
(37)
where we have neglected the tiny rotation $`\theta _{13}^\nu `$. The matrix $`U_L`$ is already given in Eq. (31), so that the full neutrino mixing matrix can be written
$$UU_{SMA}\left(\begin{array}{ccc}1/\sqrt{2}& 1/\sqrt{2}& 0\\ 1/\sqrt{2}& \sqrt{2}& 0\\ 0& 0& 1\end{array}\right),$$
(38)
where $`U_{SMA}`$ is given in Eq. (32), and is just the form that results in the small mixing angle (SMA) MSW case of this model. In other words, the net result of the large mxing of the first family produced by the $`A`$ entries in Eq. (34) is simply to multiply the (SMA) MSW form of $`U`$ on the right by a rotation of $`\pi /4`$ in the 1-2 plane. Consequently, the predictions for $`U_{\mu 3}`$ and $`U_{e3}`$ are essentially unaffected. However, $`U_{e2}`$ becomes $`1/\sqrt{2}`$ instead of the value given in Eq. (20).
The interesting lesson is that “bimaximal mixing” is easy to achieve if the large mixing of $`\nu _\mu `$ and $`\nu _\tau `$ comes from the charged lepton sector, i.e., from diagonalizing $`L`$, while the large mixing of $`\nu _e`$ comes from diagonalizing $`M_\nu `$.
The simple form given in Eq. (33) gives $`\theta _{23}^\nu ϵ8^{}`$, and thus $`\theta _{\mu \tau }=\theta _{23}^L\theta _{23}^\nu 55^{}`$, corresponding to $`\mathrm{sin}^22\theta _{\mu \tau }=0.88`$. Somewhat larger values of $`\mathrm{sin}^22\theta _{\mu \tau }`$ can arise from a more general form $`M_R`$. Consider, for example,
$$M_R=\left(\begin{array}{ccc}0& Aϵ^3& Cϵ^2\\ Aϵ^3& Bϵ^2& 0\\ Cϵ^2& 0& 1\end{array}\right)\mathrm{\Lambda }_R.$$
(39)
Then
$$M_\nu =\left(\begin{array}{ccc}\frac{\eta ^2}{ϵ^4}B& \frac{\eta }{ϵ}BC& \frac{\eta }{ϵ}(BCA)\\ \frac{\eta }{ϵ}BC& ϵ^2A^2& ϵA(A+C)\\ \frac{\eta }{ϵ}(BCA)& ϵA(A+C)& (A+C)^2\end{array}\right)(A^2+BC^2)^1\frac{M_U^2}{\mathrm{\Lambda }_R}$$
(40)
Here, as in Eq. (34), the 2-3 block has vanishing determinant. The crucial difference is that the diagonalization of this matrix involves a rotation in the 2-3 plane by an angle $`\theta _{23}^\nu ϵ\frac{A}{A+C}`$. With $`C=\frac{1}{2}A`$, for instance, $`\theta _{\mu \tau }`$ comes out very close to $`45^{}`$. Otherwise, this case is quite similar to that of Eq. (33).
## VI DETAILS OF A SPECIFIC MODEL
In the previous Sections we have presented the construction of our $`SO(10)`$ minimal Higgs model in the framework of effective $`SO(10)`$ and $`SU(5)`$ operators. We now show that one can construct a technically-natural realization of this scheme by introducing sets of Higgs and matter superfields with a well-defined family symmetry. We first address the Higgs sector.
### A Higgs Sector with $`U(1)\times Z_2\times Z_2`$ Family Symmetry
The doublet-triplet splitting problem in $`SU(5)`$, and therefore $`SO(10)`$, arises because the colored Higgs in the $`\mathrm{𝟓}\overline{\mathrm{𝟓}}`$ pairs of each $`\mathrm{𝟏𝟎}_H`$ must be made superheavy at the GUT scale, while just one pair of Higgs doublets should remain massless there and be free to develop VEV’s at the electroweak scale. This problem has been addressed and solved in by the introduction of just one $`\mathrm{𝟒𝟓}`$ adjoint Higgs field with its VEV pointing in the $`BL`$ direction, together with two pairs of $`\mathrm{𝟏𝟔}+\overline{\mathrm{𝟏𝟔}}`$ spinor Higgs fields, two $`\mathrm{𝟏𝟎}`$ Higgs in the vector representation plus several Higgs singlets. We shall briefly summarize the solution, but first we note that it is necessary to introduce several more Higgs fields in the vector and singlet representations in order to generate the Yukawa structure for the fermion mass matrices presented earlier.
The authors of found that the Higgs superpotential required to solve the doublet-triplet splitting problem could be neatly obtained from their list of Higgs fields by introducing a global family symmetry group of the type $`U(1)\times Z_2\times Z_2`$, which can arise in a natural fashion from string theory. With this in mind, we now list in Table I all the Higgs fields to be considered together with their family charge assignments.
| Higgs Fields Needed to Solve the 2-3 Problem: | |
| --- | --- |
| $`\mathrm{𝟒𝟓}_{BL}`$: | $`A(0)^+`$ |
| $`\mathrm{𝟏𝟔}`$: | $`C(\frac{3}{2})^+,C^{}(\frac{3}{2}p)^{++}`$ |
| $`\overline{\mathrm{𝟏𝟔}}`$: | $`\overline{C}(\frac{3}{2})^{++},\overline{C}^{}(\frac{3}{2}p)^+`$ |
| $`\mathrm{𝟏𝟎}`$: | $`T_1(1)^{++},T_2(1)^+`$ |
| $`\mathrm{𝟏}`$: | $`X(0)^{++},P(p)^+,Z_1(p)^{++},Z_2(p)^{++}`$ |
| Additional Higgs Fields for the Mass Matrices: | |
| $`\mathrm{𝟏𝟎}`$: | $`T_0(1+p)^+,T_o^{}(1+2p)^+`$, |
| | $`\overline{T}_o(3+p)^+,\overline{T}_o^{}(13p)^+`$ |
| $`\mathrm{𝟏}`$: | $`Y(2)^+,Y^{}(2)^{++},S(22p)^{},S^{}(23p)^{}`$, |
| | $`V_M(4+2p)^{++}`$ |
| Table I. Higgs superfields in the proposed model. | |
| --- | --- |
As noted in the table, in order to complete the construction of the Dirac mass matrices, four more vector Higgs fields and four additional Higgs singlets are needed, while one Higgs singlet is introduced to generate the right-handed Majorana neutrino mass matrix.
It is then possible to write down explicitly the full Higgs superpotential from the Higgs $`SO(10)`$ and family assignments, where we have written it as the sum of five terms:
$$\begin{array}{ccc}\hfill W_{\mathrm{Higgs}}& =& W_A+W_{CA}+W_{2/3}+W_{H_D}+W_R\hfill \\ \hfill W_A& =& trA^4/M+M_AtrA^2\hfill \\ \hfill W_{CA}& =& X(\overline{C}C)^2/M_C^2+F(X)\hfill \\ & & +\overline{C}^{}(PA/M_1+Z_1)C+\overline{C}(PA/M_2+Z_2)C^{}\hfill \\ \hfill W_{2/3}& =& T_1AT_2+Y^{}T_2^2\hfill \\ \hfill W_{H_D}& =& T_1\overline{C}\overline{C}Y^{}/M+\overline{T}_0CC^{}+\overline{T}_0(T_0S+T_0^{}S^{})\hfill \\ \hfill W_R& =& \overline{T}_0\overline{T}_0^{}V_M\hfill \end{array}$$
(41)
The Higgs singlets are all assumed to develop VEV’s at the GUT scale. We can then determine the fate of the other Higgs fields from the F-flat and D-flat conditions. In particular, $`W_A`$ fixes $`A`$ through the $`F_A=0`$ condition where one solution is $`ABL`$, the Dimopoulos-Wilczek solution . $`W_{CA}`$ gives a GUT-scale VEV to $`\overline{C}`$ and $`C`$ by the $`F_X=0`$ condition and also couples the adjoint $`A`$ to the spinors $`C,\overline{C},C^{}`$ and $`\overline{C}^{}`$ without destabilizing the Dimopoulos-Wilczek solution or giving Goldstone modes, as shown in . $`W_{2/3}`$ gives the doublet-triplet splitting by the Dimopoulos-Wilczek mechanism , . $`W_{H_D}`$ mixes the $`(1,2,1/2)`$ doublet in $`T_1`$ with those in $`C^{}`$ (by $`F_{\overline{C}}=0`$), and in $`T_0`$ and $`T_0^{}`$ (by $`F_{\overline{T}_0}=0`$).
### B Yukawa Sector
We now turn to the Yukawa sector and specify the matter fields and their $`U(1)\times Z_2\times Z_2`$ charge assignments which will complete the realization of the specific model in question. For this purpose, we require three spinor fields $`\mathrm{𝟏𝟔}_i`$, one for each light family, two vector-like pairs of $`\mathrm{𝟏𝟔}\overline{\mathrm{𝟏𝟔}}`$ spinors which can get superheavy, a pair of superheavy $`\mathrm{𝟏𝟎}`$ fields in the vector representation, and three pairs of superheavy $`\mathrm{𝟏}\mathrm{𝟏}^c`$ fermion singlets. The complete listing is given in Table II.
| $`\mathrm{𝟏𝟔}_1(\frac{1}{2}2p)^+`$ | $`\mathrm{𝟏𝟔}_2(\frac{1}{2}+p)^{++}`$ | $`\mathrm{𝟏𝟔}_3(\frac{1}{2})^{++}`$ |
| --- | --- | --- |
| $`\mathrm{𝟏𝟔}(\frac{1}{2}p)^+`$ | $`\mathrm{𝟏𝟔}^{}(\frac{1}{2})^+`$ | |
| $`\overline{\mathrm{𝟏𝟔}}(\frac{1}{2})^+`$ | $`\overline{\mathrm{𝟏𝟔}}^{}(\frac{3}{2}+2p)^+`$ | |
| $`\mathrm{𝟏𝟎}_1(1p)^+`$ | $`\mathrm{𝟏𝟎}_2(1+p)^{++}`$ | |
| $`\mathrm{𝟏}_1(2+2p)^+`$ | $`\mathrm{𝟏}_2(2p)^{++}`$ | $`\mathrm{𝟏}_3(2)^{++}`$ |
| $`\mathrm{𝟏}_1^c(22p)^+`$ | $`\mathrm{𝟏}_2^c(2)^+`$ | $`\mathrm{𝟏}_3^c(2p)^{++}`$ |
| Table II. Matter superfields in the proposed model. | | |
In terms of these fermion fields and the Higgs fields previously introduced, one can then spell out all the terms in the Yukawa superpotential which follow from their $`SO(10)`$ and $`U(1)\times Z_2\times Z_2`$ assignments:
$$\begin{array}{cc}\hfill W_{Yukawa}& =\mathrm{𝟏𝟔}_3\mathrm{𝟏𝟔}_3T_1+\mathrm{𝟏𝟔}_2\mathrm{𝟏𝟔}T_1+\mathrm{𝟏𝟔}^{}\mathrm{𝟏𝟔}^{}T_1\hfill \\ & +\mathrm{𝟏𝟔}_3\mathrm{𝟏𝟔}_1T_0^{}+\mathrm{𝟏𝟔}_2\mathrm{𝟏𝟔}_1T_0+\mathrm{𝟏𝟔}_3\overline{\mathrm{𝟏𝟔}}A\hfill \\ & +\mathrm{𝟏𝟔}_1\overline{\mathrm{𝟏𝟔}}^{}Y^{}+\mathrm{𝟏𝟔}\overline{\mathrm{𝟏𝟔}}P+\mathrm{𝟏𝟔}^{}\overline{\mathrm{𝟏𝟔}}^{}S\hfill \\ & +\mathrm{𝟏𝟔}_3\mathrm{𝟏𝟎}_2C^{}+\mathrm{𝟏𝟔}_2\mathrm{𝟏𝟎}_1C+\mathrm{𝟏𝟎}_1\mathrm{𝟏𝟎}_2Y\hfill \\ & +\mathrm{𝟏𝟔}_3\mathrm{𝟏}_3\overline{C}+\mathrm{𝟏𝟔}_2\mathrm{𝟏}_2\overline{C}+\mathrm{𝟏𝟔}_1\mathrm{𝟏}_1\overline{C}\hfill \\ & +\mathrm{𝟏}_3\mathrm{𝟏}_3^cZ+\mathrm{𝟏}_2\mathrm{𝟏}_2^cP+\mathrm{𝟏}_1\mathrm{𝟏}_1^cX\hfill \\ & +\mathrm{𝟏}_3^c\mathrm{𝟏}_3^cV_M+\mathrm{𝟏}_1^c\mathrm{𝟏}_2^cV_M\hfill \end{array}$$
(42)
where the coupling parameters have been suppressed. To obtain the GUT scale structure for the fermion mass matrix elements, all but the three chiral spinor fields in the first line of Table II. will be integrated out to yield Froggatt-Nielsen diagrams of the type pictured earlier. Note that the right-handed Majorana matrix elements will all be generated through the Majorana couplings of the $`V_M`$ Higgs field with conjugate singlet fermions given in the last two terms of Eq. (42).
In order to present a clearer description of how the GUT scale mass matrices are determined from the Yukawa and Higgs superpotentials, we shall illustrate the procedure for the up quark mass matrix $`U`$. The three massless color-triplet quark states each with charge $`2/3`$ are obtained as linear combinations of all such color and charge states within the fermion supermultiplets given in (41). In particular, the basis for the left-handed states $`u_L`$ and left-handed conjugate states $`u_L^c`$ can be ordered as follows:
$$\begin{array}{cc}\hfill _{u_L}=& \{|3,2,\frac{1}{6}>_{10(16_1)},|3,2,\frac{1}{6}>_{10(16_2)},|3,2,\frac{1}{6}>_{10(16_3)},|3,2,\frac{1}{6}>_{10(16)},\hfill \\ & |3,2,\frac{1}{6}>_{10(16^{})},|3,1,\frac{2}{3}>_{\overline{10}(\overline{16})},|3,1,\frac{2}{3}>_{\overline{10}(\overline{16}^{})}\}\hfill \end{array}$$
(43)
$$\begin{array}{cc}\hfill _{u_L^c}=& \{|\overline{3},1,\frac{2}{3}>_{10(16_1)},|\overline{3},1,\frac{2}{3}>_{10(16_2)},|\overline{3},1,\frac{2}{3}>_{10(16_3)},|\overline{3},1,\frac{2}{3}>_{10(16)},\hfill \\ & |\overline{3},1,\frac{2}{3}>_{10(16^{})},|\overline{3},2,\frac{1}{6}>_{\overline{10}(\overline{16})},|\overline{3},2,\frac{1}{6}>_{\overline{10}(\overline{16}^{})}\}\hfill \end{array}$$
(44)
where the states are labeled by their representations and hypercharge according to $`|SU(3)_c,SU(2)_L,Y_{SU(5)(SO(10))}`$.
We then form the Yukawa contribution $`u_L^{cT}C^1D_uu_L`$ by using the above bases and the superpotentials to obtain the matrix
$$D_u=\left(\begin{array}{ccccccc}0& 0& 0& 0& 0& 0& y^{}\\ 0& 0& 0& t_2& 0& 0& 0\\ 0& 0& t_3& 0& 0& a& 0\\ 0& t_2& 0& 0& 0& p& 0\\ 0& 0& 0& 0& t^{}& 0& s^{\prime \prime }\\ 0& 0& a& p& 0& 0& 0\\ y^{}& 0& 0& 0& s^{\prime \prime }& 0& 0\end{array}\right)$$
where we have introduced the following shorthand notation:
$$\begin{array}{cccc}\hfill t_3& =\lambda _{16_316_3T_1}T_1,\hfill & \hfill t_2=& \lambda _{16_216T_1}T_1,t^{}=\lambda _{16^{}16^{}T_1}T_1,\hfill \\ \hfill a& =\lambda _{16_3\overline{16}A}A,\hfill & \hfill p=& \lambda _{16\overline{16}P}P,\hfill \\ \hfill s^{\prime \prime }& =\lambda _{16^{}\overline{16}^{}S}S,\hfill & \hfill y^{}=& \lambda _{16_1\overline{16}^{}Y^{}}Y^{}\hfill \end{array}.$$
(45)
We can then determine from this matrix the three pairs of zero-mass eigenstates at the GUT scale where the electroweak VEV of $`T_1`$ vanishes:
$$\begin{array}{cc}\hfill |u_{1L}=& \left[|10(16_1)\frac{y^{}}{s^{\prime \prime }}|10(16^{})\right]/\sqrt{1+y^2/s^{\prime \prime 2}}\hfill \\ \hfill |u_{2L}=& |10(16_2)\hfill \\ \hfill |u_{3L}=& \left[|10(16_3)\frac{a}{p}|10(16)\right]/\sqrt{1+a^2/p^2}\hfill \\ \hfill |u_{1L}^c=& \left[|10(16_1)\frac{y^{}}{s^{\prime \prime }}|10(16^{})\right]/\sqrt{1+y^2/s^{\prime \prime 2}}\hfill \\ \hfill |u_{2L}^c=& |10(16_2)\hfill \\ \hfill |u_{3L}^c=& \left[|10(16_3)\frac{a}{p}|10(16)\right]/\sqrt{1+a^2/p^2}\hfill \end{array}$$
(46)
and where the states are now simply labeled by their $`SU(5)`$ and $`SO(10)`$ representations.
Finally, the Dirac matrix $`U`$ for the three light quark states $`u,c,t`$ is obtained by bracketing the electroweak contributions by the appropriate $`u_{iL}^c`$ state on the left and the $`u_{jL}`$ state on the right. The result obtained for $`U`$ has exactly the form found earlier in Eq. (15) from the previous effective operator approach, with the identifications:
$$\begin{array}{cc}\hfill M_U=& (t_3)_{5(10)}\hfill \\ \hfill ϵM_U=& |3(a_q/p)(t_2)_{5(10)}|\hfill \\ \hfill \eta M_U=& (y^{}/s^{\prime \prime })^2(t^{})_{5(10)}\hfill \end{array}$$
(47)
Here the subscript on $`a_q`$ signifies a factor of $`1/3`$ arising from the $`BL`$ VEV of the $`A`$ in the adjoint representation, while the subscripts on the $`t`$ terms specify the appropriate doublet VEV in the $`\mathrm{𝟏𝟎}`$ for $`T_1`$. We have neglected the state normalization factors in (47) but will later argue that they can all be taken to be approximately unity.
The Dirac matrices, $`D,N,L`$ are constructed in a similar fashion. In the case of $`D`$ and $`L`$, the bases corresponding to Eqs. (43) and (44) are enlarged by two states lying in the $`\mathrm{𝟏𝟎}_1`$ and $`\mathrm{𝟏𝟎}_2`$ representations of $`SO(10)`$. For $`N`$, on the other hand, in addition to the two above states, one must add the singlet fermion contributions from the representations $`\mathrm{𝟏}_k,k=1,2,3`$ for $`_{\nu _L}`$ and $`\mathrm{𝟏}_k^c,k=1,2,3`$ for $`_{\nu _L^c}`$. We list the zero-mass state vectors for $`D`$, $`N`$ and $`L`$ in analogy to Eqs. (46):
$$\begin{array}{cc}\hfill |d_{1L}=& \left[|10(16_1)\frac{y^{}}{s^{\prime \prime }}|10(16^{})\right]/\sqrt{1+y^2/s^{\prime \prime 2}}\hfill \\ \hfill |d_{2L}=& |10(16_2)\hfill \\ \hfill |d_{3L}=& \left[|10(16_3)\frac{a}{p}|10(16)\right]/\sqrt{1+a^2/p^2}\hfill \\ \hfill |d_{1L}^c=& \left[|\overline{5}(16_1)\frac{y^{}}{s^{\prime \prime }}|\overline{5}(16^{})\right]/\sqrt{1+y^2/s^{\prime \prime 2}}\hfill \\ \hfill |d_{2L}^c=& \left[|\overline{5}(16_2)\frac{c}{y}|\overline{5}(10_2)\right]/\sqrt{1+c^2/y^2}\hfill \\ \hfill |d_{3L}^c=& \left[|\overline{5}(16_3)\frac{a}{p}|\overline{5}(16)\right]/\sqrt{1+a^2/p^2}\hfill \end{array}$$
(48)
$$\begin{array}{cc}\hfill |n_{1L}=& \left[|\overline{5}(16_1)\frac{y^{}}{s^{\prime \prime }}|\overline{5}(16^{})\right]/\sqrt{1+y^2/s^{\prime \prime 2}}\hfill \\ \hfill |n_{2L}=& \left[|\overline{5}(16_2)\frac{c}{y}|\overline{5}(10_2)\right]/\sqrt{1+c^2/y^2}\hfill \\ \hfill |n_{3L}=& \left[|\overline{5}(16_3)\frac{a}{p}|\overline{5}(16)\right]/\sqrt{1+a^2/p^2}\hfill \\ \hfill |n_{1L}^c=& \left[|1(16_1)\frac{y^{}}{s^{\prime \prime }}|1(16^{})\frac{\overline{c}_1}{x}|1_1^c\right]/\sqrt{1+y^2/s^{\prime \prime 2}+\overline{c}_1^2/x^2}\hfill \\ \hfill |n_{2L}^c=& \left[|1(16_2)\frac{\overline{c}_2}{p_{22}}|1_2^c\right]/\sqrt{1+\overline{c}_2^2/p_{22}^2}\hfill \\ \hfill |n_{3L}^c=& \left[|1(16_3)\frac{a}{p}|1(16)\frac{\overline{c}_3}{z}|1_3^c\right]/\sqrt{1+a^2/p^2+\overline{c}_3^2/z^2}\hfill \end{array}$$
(49)
$$\begin{array}{cc}\hfill |\mathrm{}_{1L}=& \left[|\overline{5}(16_1)\frac{y^{}}{s^{\prime \prime }}|\overline{5}(16^{})\right]/\sqrt{1+y^2/s^{\prime \prime 2}}\hfill \\ \hfill |\mathrm{}_{2L}=& \left[|\overline{5}(16_2)\frac{c}{y}|\overline{5}(10_2)\right]/\sqrt{1+c^2/y^2}\hfill \\ \hfill |\mathrm{}_{3L}=& \left[|\overline{5}(16_3)\frac{a}{p}|\overline{5}(16)\right]/\sqrt{1+a^2/p^2}\hfill \\ \hfill |\mathrm{}_{1L}^c=& \left[|10(16_1)\frac{y^{}}{s^{\prime \prime }}|10(16^{})\right]/\sqrt{1+y^2/s^{\prime \prime 2}}\hfill \\ \hfill |\mathrm{}_{2L}^c=& |10(16_2)\hfill \\ \hfill |\mathrm{}_{3L}^c=& \left[|10(16_3)\frac{a}{p}|10(16)\right]/\sqrt{1+a^2/p^2}\hfill \end{array}$$
(50)
In the above we have introduced, in analogy with Eqs. (45), the additional shorthand notation:
$$\begin{array}{cccc}\hfill c=& \lambda _{16_210_1C}C,\hfill & \hfill \overline{c}_i=& \lambda _{16_i1_i\overline{C}}\overline{C},i=1,2,3,\hfill \\ \hfill x=& \lambda _{1_11_1^cX}X,\hfill & \hfill y=& \lambda _{10_110_2Y}Y,\hfill \\ \hfill z=& \lambda _{1_31_3^cZ}Z,\hfill & \hfill p_{22}=& \lambda _{1_21_2^cP}P\hfill \end{array}$$
(51)
The Dirac matrices $`D,N`$ and $`L`$ are found by forming matrix elements of the electroweak symmetry breaking VEV’s with the appropriate basis vectors. Again these matrices have exactly the structures given in Eqs. (15), provided the state normalization factors are approximated by unity, i.e., we assume that the zero-mass states have their large components in the chiral representations $`\mathrm{𝟏𝟔}_1,\mathbf{16}_2`$ and $`\mathrm{𝟏𝟔}_3`$ and that all the other components are small. We shall return to this point in the next Section. In the meantime we make the identifications
$$\begin{array}{cc}\hfill M_D=& (t_3)_{\overline{5}(10)}\hfill \\ \hfill ϵM_D=& |3(a_q/p)(t_2)_{\overline{5}(10)}|\hfill \\ \hfill \sigma M_D=& (c/y)(c^{})_{\overline{5}(16)}\hfill \\ \hfill \mathrm{\Delta }M_D=& t_0\overline{t}_0/s\hfill \\ \hfill \delta ^{}e^{i\varphi }M_D=& t_0^{}\overline{t}_0/s^{}\hfill \end{array}$$
(52)
in terms of the notation given in Eqs. (45) and (51) and the following:
$$\begin{array}{cccc}\hfill t_0=& \lambda _{16_116_2T_0},\hfill & \hfill t_0^{}=& \lambda _{16_116_3T_0^{}}\hfill \\ \hfill \overline{t}_0=& \lambda _{CC^{}\overline{T}_0}CC^{},\hfill & \hfill c^{}=& \lambda _{16_310_2C^{}}C^{},\hfill \\ \hfill s=& \lambda _{T_0\overline{T}_0S}S,\hfill & \hfill s^{}=& \lambda _{T_0^{}\overline{T}_0S^{}}S^{}\hfill \end{array}$$
(53)
The phase $`\varphi `$ appearing in the $`\delta ^{}`$ term can be understood to arise from a phase in the VEV for $`S^{}`$. The structures of the Dirac matrix elements given in Eqs. (15), (47) and (52) can be understood in terms of the simple Froggatt-Nielsen diagrams of Fig. 1 and 2, with the Higgs fields labeled as in Table I.
Turning to the right-handed Majorana mass matrix, we use the zero mass left-handed conjugate states that were obtained implicitly above for the Dirac matrix $`N`$ to form the basis for $`M_R`$. The right-handed Majorana matrix is then obtained by bracketing the Majorana Higgs $`V_M`$ with the appropriate zero mass conjugate neutrino states in (49). We obtain
$$M_R=\left(\begin{array}{ccc}0& Aϵ^3& 0\\ Aϵ^3& 0& 0\\ 0& 0& 1\end{array}\right)\mathrm{\Lambda }_R$$
(54)
where
$$\begin{array}{cc}\hfill M_3=\mathrm{\Lambda }_R=& \lambda _{1_3^c1_3^cV_M}V_M(\overline{c}_3/z)^2,\hfill \\ \hfill M_2=M_1=Aϵ^3\mathrm{\Lambda }_R=& \lambda _{1_1^c1_2^cV_M}V_M(\overline{c}_1/x)(\overline{c}_2/p_{22})\hfill \end{array}$$
(55)
The lighter two right-handed Majorana masses are degenerate and have opposite CP-parity. Note that the whole right-handed Majorana matrix has been generated in this simple model by one Majorana VEV coupling superheavy conjugate fermion singlets as shown in Fig. 4.
We conclude this Section with a summary of the GUT scale predictions derived from the Dirac and Majorana mass matrices with the particular parameters appropriate for the model in question. For convenience we give the whole set equations which are the counterpart of Eqs. (19).
$$\begin{array}{cc}\hfill m_t^0/m_b^0& (\sigma ^2+1)^{1/2}M_U/M_D,m_u^0/m_t^0\eta ,\hfill \\ \hfill m_c^0/m_t^0& \frac{1}{9}ϵ^2[1\frac{2}{9}ϵ^2],m_b^0/m_\tau ^01\frac{2}{3}\frac{\sigma }{\sigma ^2+1}ϵ,\hfill \\ \hfill m_s^0/m_b^0& \frac{1}{3}ϵ\frac{\sigma }{\sigma ^2+1}[1+\frac{1}{3}ϵ\frac{1\sigma ^2\sigma ϵ/3}{\sigma (\sigma ^2+1)}+\frac{1}{2}(t_L^2+t_R^2)],\hfill \\ \hfill m_d^0/m_s^0& t_Lt_R[1\frac{1}{3}ϵ\frac{\sigma ^2+2}{\sigma (\sigma ^2+1)}(t_L^2+t_R^2)\hfill \\ & +(t_L^4+t_L^2t_R^2+t_R^4)],\hfill \\ \hfill m_\mu ^0/m_\tau ^0& ϵ\frac{\sigma }{\sigma ^2+1}[1+ϵ\frac{1\sigma ^2\sigma ϵ}{\sigma (\sigma ^2+1)}+\frac{1}{18}(t_L^2+t_R^2)],\hfill \\ \hfill m_e^0/m_\mu ^0& \frac{1}{9}t_Lt_R[1ϵ\frac{\sigma ^2+2}{\sigma (\sigma ^2+1)}+ϵ^2\frac{\sigma ^4+9\sigma ^2/2+3}{\sigma ^2(\sigma ^2+1)^2}\hfill \\ & \frac{1}{9}(t_L^2+t_R^2)],\hfill \\ \hfill V_{cb}^0& \frac{1}{3}ϵ\frac{\sigma ^2}{\sigma ^2+1}[1+\frac{2}{3}ϵ\frac{1}{\sigma (\sigma ^2+1)}],\hfill \\ \hfill V_{us}^0& t_L[1\frac{1}{2}t_L^2t_R^2+t_R^4+\frac{5}{2}t_L^2t_R^2+\frac{3}{8}t_L^4\hfill \\ & \frac{ϵ}{3\sigma \sqrt{\sigma ^2+1}}\frac{t_R}{t_L}e^{i\theta }],\hfill \\ \hfill V_{ub}^0& \frac{1}{3}t_Lϵ\frac{1}{\sigma ^2+1}[\sqrt{\sigma ^2+1}\frac{t_R}{t_L}e^{i\theta }(1\frac{1}{3}ϵ\frac{\sigma }{\sigma ^2+1})\hfill \\ & (1\frac{2}{3}ϵ\frac{\sigma }{\sigma ^2+1})],\hfill \\ \hfill m_2^0/m_3^0& \left(\frac{\eta }{Aϵ\sqrt{1+ϵ^2}}\right)\left[1+\frac{\eta }{Aϵ^3\sqrt{1+ϵ^2}}\right],\hfill \\ \hfill m_1^0/m_3^0& \left(\frac{\eta }{Aϵ\sqrt{1+ϵ^2}}\right)\left[1\frac{\eta }{2Aϵ^3\sqrt{1+ϵ^2}}\right],\hfill \\ \hfill U_{\mu 3}^0& \frac{1}{\sqrt{\sigma ^2+1}}(\sigma ϵ\frac{\sigma ^2}{\sigma ^2+1}),\hfill \\ \hfill U_{e2}^0& \frac{1}{\sqrt{2}}[1\frac{ϵ}{3\sigma }t_Le^{i\theta }\hfill \\ & +\frac{1}{3\sqrt{\sigma ^2+1}}(1+ϵ\sigma )t_R],\hfill \\ \hfill U_{e3}^0& \frac{1}{3\sqrt{\sigma ^2+1}}(\sigma ϵ)t_R\frac{\eta }{Aϵ^2}\hfill \end{array}$$
(56)
To the quark equations we have added the ratio $`m_t^0/m_b^0`$ which involves the ratio of $`\mathrm{𝟓}(T_1)`$ to $`\overline{\mathrm{𝟓}}(T_1)`$, i.e., $`M_U/M_D`$, as well as giving the leptonic mass ratios and mixings specific to the model in question.
### C Numerical Evaluation of Matrix Parameters
We have elaborated above how the simple explicit model proposed gives precisely the structure for the Dirac mass matrices that was obtained from the effective operator approach. We now show that the entries are also numerically in the range to fit the quark and lepton mass and mixing data.
In order to compare the GUT scale predictions in Eq. (56) with the low scale data, the GUT scale values are first evolved from $`\mathrm{\Lambda }_G=2\times 10^{16}`$ GeV down to the SUSY scale which is taken to be $`\mathrm{\Lambda }_{\mathrm{SUSY}}=m_t(m_t)=165`$ GeV by use of the 2-loop MSSM beta functions. For this purpose, the mass ratios at the two scales are related by the $`\eta _{i/j}`$ running factors, while the quark mixing elements are scaled by the $`\eta _{ij}`$ factor according to
$$\begin{array}{cc}\hfill \left(\frac{m_i}{m_j}\right)_{\mathrm{SUSY}}=& \left(\frac{m_i^0}{m_j^0}\right)/\eta _{i/j},\hfill \\ \hfill (V_{ij})_{\mathrm{SUSY}}=& V_{ij}^0/\eta _{ij},(ij)=(ub),(cb),(td),(ts)\hfill \end{array}$$
(57)
The remaining evolutions to the bottom and charm quark or tau lepton running mass scales, or to the 1 GeV scale in the case of the light quarks and leptons, is carried out with the 3-loop QCD and 1-loop QED renormalization group equations. Here the running factors are $`\eta _i`$ with the mass ratios scaled according to
$$m_i(m_i)=(m_i)_{\mathrm{SUSY}}/\eta _i(m_t)$$
(58)
or similarly, with the running mass scale $`m_i`$ replaced by 1 GeV. With $`\mathrm{tan}\beta =5`$ used for the numerical evaluations for reasons that will become apparently shortly, $`\alpha _s(M_Z)=0.118,\alpha (M_Z)=1/127.9`$ and $`\mathrm{sin}^2\theta _W=0.2315`$, the running factors are given by
$$\begin{array}{cccc}\hfill \eta _{u/t}=\eta _{c/t}=& 0.6927,\hfill & \hfill \eta _{d/b}=& \eta _{s/b}=0.8844\hfill \\ \hfill \eta _{\mu /\tau }=& 0.9988,\hfill & \hfill \eta _{b/t}=& 0.5094\hfill \\ & \multicolumn{3}{c}{\eta _{ub}=\eta _{cb}=\eta _{td}=\eta _{ts}=0.8853}\\ \hfill \eta _u(m_t)=0.4235,& \eta _c(m_t)=0.4733,\hfill & \hfill \eta _t(m_t)=1.0000& \\ \hfill \eta _d(m_t)=0.4262,& \eta _s(m_t)=0.4262,\hfill & \hfill \eta _t(m_t)=0.6540& \\ \hfill \eta _e(m_t)=0.9816,& \eta _\mu (m_t)=0.4816,\hfill & \hfill \eta _\tau (m_t)=0.9836& \end{array}$$
(59)
Finally, finite corrections must be applied to $`m_s`$, $`m_b`$ and the evolved quark mixing matrix elements which arise from gluino and chargino loops. The correction factors are conventionally written as $`(1+\mathrm{\Delta }_s),(1+\mathrm{\Delta }_b)`$ and $`(1+\mathrm{\Delta }_{cb})`$ where we have used
$$\mathrm{\Delta }_s=0.20,\mathrm{\Delta }_b=0.15,\mathrm{\Delta }_{cb}=0.05$$
(60)
as explained below.
Using the quantities $`m_t(m_t)=165\mathrm{GeV},m_\tau =1.777\mathrm{GeV},m_\mu =105.7\mathrm{MeV},m_e=0.511\mathrm{MeV},m_u=4.5\mathrm{MeV},V_{us}=0.220,V_{cb}=0.0395`$, and $`\delta _{CP}=64^o`$ to determine the input parameters, one obtains for them $`M_U113\mathrm{GeV},M_D1\mathrm{GeV},\sigma =1.780,ϵ=0.145,t_L=0.236,t_R=0.205,\theta =34^o(\mathrm{corresponding}\mathrm{to}\delta =0.0086,\delta ^{}=0.0079,\varphi =54^o)`$, and $`\eta =8\times 10^6`$. With these inputs the remaining quark masses and mixings are obtained, to be compared with the experimental values in parentheses:
$$\begin{array}{ccc}\hfill m_c(m_c)=& 1.23\mathrm{GeV}\hfill & (1.27\pm 0.1\mathrm{GeV})\hfill \\ \hfill m_b(m_b)=& 4.25\mathrm{GeV}\hfill & (4.26\pm 0.11\mathrm{GeV})\hfill \\ \hfill m_s(1\mathrm{GeV})=& 148\mathrm{MeV}\hfill & (175\pm 50\mathrm{MeV})\hfill \\ \hfill m_d(1\mathrm{GeV})=& 7.9\mathrm{MeV}\hfill & (8.9\pm 2.6\mathrm{MeV})\hfill \\ \hfill |V_{ub}/V_{cb}|=& 0.080\hfill & (0.090\pm 0.008)\hfill \end{array}$$
(61)
where the finite SUSY loop corrections for $`m_b,m_s`$ and $`V_{cb}`$ have been rescaled to give $`m_b(m_b)4.25`$ GeV for $`\mathrm{tan}\beta =5.`$ Had we chosen $`\delta _{CP}=70^o`$ as input, on the other hand, we would find instead $`|V_{ub}/V_{cb}|=0.085`$. With the numerical values in (61) we find for $`\overline{\rho },\overline{\eta }`$ and the $`\alpha ,\beta `$ and $`\gamma `$ angles of the unitarity triangle pictured in Fig. 3
$$\overline{\rho }=0.143,\overline{\eta }=0.305,\alpha =96^o,\beta =20^o,\gamma =64^o$$
(62)
The upper vertex of the triangle appears to be circled precisely in the allowed experimental region.
Additional predictions follow for the neutrino sector. The effective light neutrino mass matrix of Eq. (34) or (36) with $`B=0`$ leads to bimaximal mixing with a large angle solution for atmospheric neutrino oscillations and the “just-so” vacuum solution involving two pseudo-Dirac neutrinos, if we set $`\mathrm{\Lambda }_R=2.4\times 10^{14}`$ GeV and $`A=0.05`$. We then find
$$\begin{array}{cc}\multicolumn{2}{c}{m_3=54.3\mathrm{meV},m_2=59.6\mu \mathrm{eV},m_1=56.5\mu \mathrm{eV}}\\ \multicolumn{2}{c}{M_3=2.4\times 10^{14}\mathrm{GeV},M_2=M_1=3.66\times 10^{10}\mathrm{GeV}}\\ \multicolumn{2}{c}{U_{e2}=0.733,U_{e3}=0.047,U_{\mu 3}=0.818,\delta _{CP}^{}=0.2^o}\\ \mathrm{\Delta }m_{23}^2=3.0\times 10^3\mathrm{eV}^2,\hfill & \mathrm{sin}^22\theta _{\mathrm{atm}}=4|U_{\mu 3}|^2|U_{\tau 3}|^2=0.89\hfill \\ \mathrm{\Delta }m_{12}^2=3.6\times 10^{10}\mathrm{eV}^2,\hfill & \mathrm{sin}^22\theta _{\mathrm{solar}}=4|U_{e1}|^2|U_{e2}|^2=0.99\hfill \\ \mathrm{\Delta }m_{13}^2=3.0\times 10^3\mathrm{eV}^2,\hfill & \mathrm{sin}^22\theta _{\mathrm{reac}}=4|U_{e3}|^2(1|U_{e3}|)^2=0.009\hfill \end{array}$$
(63)
The effective scale of the right-handed Majorana mass contribution occurs two orders of magnitude lower than the SUSY GUT scale of $`\mathrm{\Lambda }_G=1.2\times 10^{16}`$ GeV. The effective two-component reactor mixing angle given above should be observable at a future neutrino factory, whereas the present limit from the CHOOZ experiment is approximately 0.1 for the above $`\mathrm{\Delta }m_{23}^2`$. In principle, the parameter $`A`$ appearing in $`M_R`$ can also be complex, but we find that in no case does the leptonic CP-violating phase, $`\delta _{CP}^{}`$ exceed $`10^o`$ in magnitude. Hence the model predicts leptonic CP-violation will be unobservable.
The vacuum solar solution depends critically on the appearance of the parameter $`\eta `$ in the matrix $`N`$, corresponding to the non-zero $`\eta `$ entry in $`U`$ which gives the up quark a mass at the GUT scale. Should we set $`\eta =0`$, only the small-angle MSW solution would be obtained for the solar neutrino oscillations. The large angle MSW solution is disfavored by the larger hierarchy, i.e., very small $`A`$ value, required in $`M_R`$.
Finally we address the issue that the state normalization factors were all replaced by unity in Eqs. (47) and (52) for the various matrix parameters. This is a good approximation provided the three fermion spinor states $`|16_1,|16_2,|16_3`$ provide the dominant contributions to the zero-mass quark and lepton states at the GUT scale. In particular, the following ratios must be much less than unity:
$$(a/p)^2,(y^{}/s^{\prime \prime })^2,(c/y)^2,(\overline{c}_1/x)^2,(\overline{c}_2/p_{22})^2,(\overline{c}_3/z)^21$$
(64)
Let us assume for simplicity that the electroweak couplings of $`T_1`$ in $`t_3,t_2`$ and $`t^{}`$ in Eq. (45) and of $`C^{}`$ in $`c^{}`$ of (53) are identical. Then with $`ϵ=|3(a_q/p)|=|a_{\mathrm{}}/p|=0.14`$, we find the condition $`(a/p)^20.021`$ holds. To obtain an up quark mass $`m_u(1GeV)4.5`$ MeV, we need $`\eta (y^{}/s^{\prime \prime })^28\times 10^6`$ at the GUT scale, which easily satisfies (64).
Requiring that $`(c/y)^21`$ and with the result from Eqs. (52) that
$$\sigma \left|\frac{c}{y}\frac{\overline{5}(C^{})}{\overline{5}(T_1)}\right|1.8$$
(65)
leads us to the results that
$$\begin{array}{cc}\hfill \mathrm{tan}\gamma & \frac{\overline{5}(C^{})}{\overline{5}(T_1)}\sigma \hfill \\ \hfill \mathrm{tan}\beta & \sqrt{\sigma ^2+1}m_t^0(cos\gamma )/m_b^0m_t^0/m_b^0\hfill \end{array}$$
(66)
in terms of the $`T_1C^{}`$ mixing angle, $`\gamma `$, in Eq. (11). With $`c/y0.1`$, for example, $`\mathrm{tan}\gamma 18`$ which implies $`\mathrm{tan}\beta 6`$, a very reasonable mid-range value allowed by experiment. For this reason, we have chosen to illustrate the numerical results above with $`\mathrm{tan}\beta =5`$.
The remaining ratios in Eq. (64) can also be satisfied. For comparable $`\overline{c}_i`$’s, $`Aϵ^31.4\times 10^4`$ obtained from Eq. (55) requires that $`Z/\sqrt{XP}0.01`$. This ratio is consistent with the VEV’s needed in the Higgs superpotential of Eq. (41) in order to solve the doublet-triplet splitting problem.
Turning now to the parameters $`\delta `$ and $`\delta ^{}`$, we note that the near equality of their magnitudes leads to the ratio $`\delta /|\delta ^{}|s/|s^{}|1`$. Moreover, if we assume $`yy^{}`$, we obtain the estimate $`\delta cy^{}/(ys)\mathrm{tan}\gamma 5\times 10^3`$ with the numbers obtained earlier, whereas the actual value needed is $`\delta 0.008`$.
Thus we have found that not only are the desired forms of the Dirac (and Majorana) matrices generated by the model of this Section, but that the numerical values required for the matrix parameters are also quite reasonable.
## VII SUMMARY
Both the largeness of the atmospheric neutrino mixing $`U_{\mu 3}`$ and the smallness of the quark mixing $`V_{cb}`$ can be elegantly accounted for by the idea that the charged lepton mass matrix $`L`$ is highly asymmetric or “lopsided” and that the down-quark mass matrix $`D`$ is related to the transpose of $`L`$ by an $`SU(5)`$ symmetry. This idea was discovered independently by several groups and has since been used in numerous models of fermion masses. Remarkably, exactly such mass matrices emerged in our work from quite other considerations than neutrino masses and mixings, specifically from an attempt to construct the simplest possible realistic $`SO(10)`$ model.
Advances have been made in recent years in simplifying the Higgs structure of SUSY $`SO(10)`$ models. If one assumes the minimal set of Higgs fields that can break $`SO(10)`$ down to the standard model group, the possibilities for Yukawa terms for the quarks and leptons become significantly restricted. It turns out that there is what seems to be a uniquely simple set of $`SO(10)`$ Yukawa terms that gives realistic masses and mixings. This set consists of only six effective Yukawa terms (five if $`m_u=0`$) which satisfactorily fits all nine masses of the quarks and charged leptons as well as the four CKM parameters. In addition, large $`\nu _\mu \nu _\tau `$ mixing emerges automatically. Moreover, in this uniquely simple model, the simplest possibilities for the Majorana mass matrix $`M_R`$ of the right-handed neutrinos lead either to small angle MSW values for the solar neutrino mixing or to vacuum oscillation values. In this paper we have studied in detail the consequences of different forms of $`M_R`$ for the neutrino mixing angles and mass ratios.
In the published literature no more predictive and economical a model for quark and lepton masses than the one discussed here exists that is also consistent with present knowledge. It is striking that in this model a single term and a single parameter (which we call $`\sigma `$) accounts for no less than four puzzling aspects of the light fermion spectrum: the largeness of $`U_{\mu 3}`$, the smallness of $`V_{cb}`$, the smallness of $`m_c/m_t`$ compared to $`m_s/m_b`$, and the Georgi-Jarlskog factor of three between $`m_\mu `$ and $`m_s`$ at the GUT scale. It should be emphasized that, while many satisfactory neutrino mixing ideas and also many interesting ideas for explaining the pattern of quark and charged lepton masses have been proposed, very few models exist which not only give a satisfactory account of neutrino phenomenology but are at the same time highly predictive.
We have shown that the model defined by the existence of these five (or six) effective Yukawa terms can be realized in a complete and specific renormalizable SUSY $`SO(10)`$ model that is technically natural. We have presented the details of such a model, including all the Higgs and quark and lepton superfields, the abelian flavor symmetries, and the transformation properties of the fields under these symmetries. Finally, we have done a quantitative comparison of the predictions of the model to experiment.
In the future this model will be rigorously testable in several ways. The most important are (1) a relation between the real and imaginary parts of $`V_{ub}`$ including a precise test of the angles of the unitarity triangle; (2) a prediction for $`U_{e2}`$, which in the small angle MSW case gives a sharp relation between the solar and atmospheric angles; and (3) a definite prediction for $`U_{e3}`$.
The research of SMB was supported in part by Department of Energy Grant Number DE FG02 91 ER 40626 A007. One of us (CHA) thanks the Fermilab Theoretical Physics Department for its kind hospitality where much of his work was carried out. Fermilab is operated by Universities Research Association Inc. under contract with the Department of Energy.
|
warning/0003/quant-ph0003040.html
|
ar5iv
|
text
|
# Unified approach to quantum capacities: towards quantum noisy coding theorem
## Abstract
Basing on unified approach to all kinds of quantum capacities we show that the rate of quantum information transmission is bounded by the maximal attainable rate of coherent information. Moreover, we show that, if for any bipartite state the one-way distillable entanglement is no less than coherent information, then one obtains Shannon-like formulas for all the capacities. The inequality also implies that the decrease of distillable entanglement due to mixing process does not exceed of corresponding loss information about a system.
The challenge for the present quantum information theory is to determine the quantum capacity of noisy channels . The problem is difficult mainly for two reasons. First, according to the present knowledge, unlike in the classical case, there are at least five different types of quantum capacities . This is because quantum information channel can be supplemented by one- or two-way classical channel . Moreover, there are teleportation channels , for which a bipartite state is a resource. The second reason is that quantum capacity exhibits a kind of nonadditivity that makes them extremely hard to deal with.
As one knows, the key success of classical information theory is the famous Shannon noisy coding theorem, giving the formula for capacity of noisy channel
$$C=\underset{X}{sup}I(X;Y)$$
(1)
where the supremum is taken over all sources $`X`$; $`I(X;Y)=H(X)+H(Y)H(X;Y)`$ is the Shannon mutual information (with $`H`$ being the Shannon entropy); $`Y`$ is the random variable resulting from action of the noise to $`X`$. In quantum information theory the candidate for the counterpart of mutual information has been found . It is the so-called coherent information (CI). To define it, consider two coherent informations of the bipartite state $`\varrho `$ with reductions $`\varrho _A`$ and $`\varrho _B`$:
$$I^X(\varrho )=S(\varrho _X)S(\varrho ),X=A,B,$$
(2)
for $`S(\varrho ^X)S(\varrho )0`$ and $`I^X=0`$ otherwise. Here $`S(\varrho )=\text{Tr}\varrho \mathrm{log}_2\varrho `$ is the von Neumann entropy. Then one defines the coherent information for a channel $`\mathrm{\Lambda }`$ and a source state $`\sigma `$ as
$$I(\sigma ,\mathrm{\Lambda })=I^B((I\mathrm{\Lambda })(|\psi \psi |)),$$
(3)
where $`\psi `$ is a pure state with reduction $`\sigma `$ (the quantity does not depend on the choice of $`\psi `$). Now, the following connection between coherent information and a quantum capacity is known
$$Q_ø\underset{n\mathrm{}}{lim}\frac{1}{n}\underset{\varrho _n}{sup}I(\varrho _n,\mathrm{\Lambda }^n)I_ø^B(\mathrm{\Lambda })$$
(4)
where $`Q_ø`$ is the maximal number of qubits that can be reliably sent down the channel without any supplementary classical channel. Note that if, instead of inequality, there were equality, then we would have analogue of Shannon formula.
Unfortunately, despite a huge effort devoted to the problem , the equality has not been proven so far. Moreover, remarkably, the similar inequality is not known for other capacities than $`Q_ø`$, i.e. the ones attainable at the support of backword ($`Q_{}`$) or two-way ($`Q_{}`$) classical communication . There are also capacities of quantum teleportation channels, where the resource is bipartite state rather than channel. In the latter case, transmission requires prior manipulations (called distillation ) over the shared pairs, transforming them into pairs in pure maximally entangled states. Then the quantum information can be transmitted by using teleportation . The manipulations include any local actions, and one- or two-way classical communication. Correspondingly, we have two kinds of one-way distillable entanglement of a state $`\varrho `$, $`D_{}`$ or $`D_{}`$ (since $`\varrho `$ need not be symmetric, one distinguishes directions of classical communication) and two-way distillable entanglement $`D_{}`$. In fact, $`Q_{}`$ also necessarily involves the distillation and teleportation process.
The latter processes are so exotic from the point of view of classical information theory, that no analogue of theory of error correcting codes has been worked out for them so far. In contrast, in the case of the capacity $`Q_ø`$, there exists a huge theory of quantum codes, being a generalization of classical error correcting codes theory .
In the above context the basic question arise: Is there possible a consistent approach for all of the capacities? In particular: Is there a single counterpart of the Shannon mutual information? In this paper, we provide a unified framework for all capacities. We show that one and the same coherent information, although in different contexts, is a basic quantity in each case. More specifically, we show that the inequality (4) is, in a sense, universal. Given any type of supplementary resources, the maximal rate of quantum communication (quantum capacity) is bounded by the maximal rate of coherent information attainable via these resources (CI capacity).
Now, there remains the fundamental question: Are the quantum capacities equal to corresponding CI capacities? We will show that the following hypothetical inequality (call it hashing inequality )
$$D_{}(\varrho )I^B(\varrho ),$$
(5)
if satisfied for all bipartite states $`\varrho `$, implies that the capacities are equal to one another in all cases (see Theorem 2). In other words, the hashing inequality implies the Shannon-like formulas for quantum capacities, providing the quantum noisy coding theorem. Consequently, we argue that to prove (or disprove) this inequality is one of fundamental tasks of the present information theory. In particular, if the inequality holds, then to evaluate $`Q_ø`$ one would need to consider only the maximization problem of the right-hand-side of the inequality (4). Finally, we show that the power of the above inequality is even more surprising. Namely, it also implies the relation between loss of classical information, and loss of distillable entanglement .
Surprisingly, the reasoning leading to our results on capacities is extremely simple. Namely, the capacity of a channel or bipartite state at given supplementary resource is the optimal rate of reliable transmission of qubits. However, this is equivalent to optimal rate of reliable sharing two-qubit pairs in maximally entangled (in short, singlet) states . Thus, we can imagine that Alice and Bob, started with large number $`n`$ of pairs in initial state $`\varrho ^n`$ (or disposed $`n`$ uses of quantum channel $`\mathrm{\Lambda }`$), aim to share the maximal attainable number $`k`$ of singlet pairs. The capacity is just the optimal rate $`k/n`$. Then, what is the coherent information $`I_{out}^X`$ of the output of such protocol attaining capacity? Since the coherent information is additive, and for singlet state $`I^X=1`$, then $`I_{out}^X`$ equals to the number $`k`$ of final singlet pairs. Thus the obtained rate of the coherent information $`I_{out}^X/n`$ is equal to capacity. But the maximal attainable rate of coherent information (i.e. CI capacity) is no less than the one acheievable in some particular protocol, so that it is no less than the capacity.
Assume now, that for any bipartite state, its capacity (i.e. distillation rate) is no less than its coherent information. Then we consider the following protocol. Alice and Bob start from $`k`$ groups of $`n`$ pairs (or divide $`kn`$ uses of channel into $`k`$ groups), with $`k,n`$ being large, and for any group obtain the final state $`\varrho `$ of maximal attainable coherent information. Then they distill the latter state, and, by assumption, obtain the final number of singlet pair no less than the coherent information of $`\varrho `$. Since the latter was maximal attainable one, we conclude that the capacity of the input state $`\varrho `$ or channel $`\mathrm{\Lambda }`$ is no less than CI capacity. Therefore, due to previous paragraph, the quantities must be equal.
The presented argumentation is very intuitive. It is similar to the approach of Ref. , that has already appeared to be fruitful in a different context . The rigourous version of the above heuristic approach is more or less immediate. Indeed, the main simplification we made was the assumption that exact singlets are produced. In fact, they are always impure (the impurity vanishes in the asymptotic or “thermodynamic” limit of infinite number of input pairs). However such simplification does not lead to wrong conclusions, if only the involved functions exhibit suitable continuity. In the rigorous proofs below we will use continuity of coherent information.
Let us now pass to the rigorous part of the paper. As mentioned, we will be concerned with four supplementary resources $`C\{,,,ø\}`$. (The last one symbolizes no supplementary resource). If Alice and Bob dispose one use of a channel $`\mathrm{\Lambda }`$ (directed, by convention, from Alice to Bob) and the supplemetary resources symbolized by $`C`$, then they can share a bipartite state $`\varrho `$. An operation that produced in this way the state $`\varrho `$ from $`\mathrm{\Lambda }`$ will be denoted by $`_C`$ so that
$$\varrho =_C(\mathrm{\Lambda }).$$
(6)
If Alice and Bob share initially a bipartite state $`\varrho _{in}`$, then we will use notation $`𝒟_C`$
$$\varrho _{out}=𝒟_C(\varrho _{in})$$
(7)
(The letters used in our notation follows from the common associations: usual channel capacity – error correction, teleportation channels – distillation). Now, the CI capacitites are defined by
$$I_C^X(\mathrm{\Lambda })=\underset{n\mathrm{}}{lim}\frac{1}{n}\underset{_C}{sup}I^X(_C(\mathrm{\Lambda }^n))$$
(8)
for channels, and
$$I_C^X(\varrho )=\underset{n\mathrm{}}{lim}\frac{1}{n}\underset{𝒟_C}{sup}I^X(𝒟_C(\varrho ^n))$$
(9)
for bipartite states . Throughout the paper, the symbol $`X`$ stands for $`A`$ or $`B`$.
We also define quantum capacities as follows . Define maximally entangled state on the space $``$ by
$$P_+()=|\psi _+()\psi _+()|,\psi _+()=\frac{1}{\sqrt{d}}\underset{i=1}{\overset{d}{}}|ii,$$
(10)
where $`|i`$ are basis vectors in $``$, while $`d=dim`$. Given a state $`\varrho `$, consider sequence of operations $`\{𝒟_C^n\}`$ (called protocol) transforming the input state $`\varrho ^n`$ into the state $`\sigma _n`$ acting on the Hilbert space $`_n_n`$ and with $`dim_n=d_n`$, satisfying
$$F_n\psi _+(_n)|\sigma _n|\psi _+(_n)1.$$
(11)
The asymptotic ratio attainable via given protocol is then given by
$$D_{\{𝒟_C^n\}}(\varrho )=\underset{n\mathrm{}}{lim}\frac{\mathrm{log}_2dim_n}{n}$$
(12)
Then the capacity $`D_C(\varrho )`$ (call it $`C`$-distillable entanglement) is defined by maximum over all possible protocols
$$D_C(\varrho )=supD_{\{𝒟_C^n\}}(\varrho ).$$
(13)
The usual channel capacities can be defined in the same way. We only need make the following substitutions: $`DQ`$, $`\varrho \mathrm{\Lambda }`$ and $`𝒟`$. The protocols $`\{_C^n\}`$ and $`\{𝒟_C^n\}`$ that achieve the considered suprema will be called optimal error correction and optimal distillation protocol, respectively. The quantity $`D_ø`$ is a bit pathological, but certainly interesting quantity. We will not be concerned with it here. However, it is likely, that $`D_ø`$ is the amount of pure entanglement that can be drawn from the state reversibly.
We will need a lemma, stating that coherent information $`I^X`$ is continuous on isotropic state. The latter is defined on $``$ (cf. )
$$\varrho (F,d)=pP_+()+(1p)\frac{1}{d^2}I,0p1,$$
(14)
with $`\text{Tr}\left[\varrho (F,d)P_+()\right]=F`$, $`d=dim`$.
Lemma. For a sequence of isotropic states $`\varrho (F_n,d_n)`$, such that $`F_n1`$ and $`d_n\mathrm{}`$ we have
$$\underset{n\mathrm{}}{lim}\frac{1}{\mathrm{log}_2d_n}I^X(\varrho (F_n,d_n))1.$$
(15)
Proof. This can be checked by direct calculation.
We note an important property of the isotropic state . Namely, any state $`\sigma `$ acting on $``$, if subjected to $`UU^{}`$ twirling (cf. ), i.e. random unitary transformations of the form $`UU^{}`$, becomes isotropic state $`\varrho (d,F)`$ with $`F=\text{Tr}\left[\sigma P_+()\right]`$, $`d=dim`$.
Now we can state the theorems being the main results of this paper.
Theorem 1. Quantum capacities are bounded from above by CI capacities:
$`Q_C(\mathrm{\Lambda })I_C^X(\mathrm{\Lambda }),`$ (16)
$`D_C(\varrho )I_C^X(\varrho )`$ (17)
for any $`\mathrm{\Lambda }`$, $`\varrho `$ and $`C\{,,,ø\}`$. Theorem 2. If the hashing inequality
$$D_{}(\varrho )I^B(\varrho )$$
(18)
holds for any bipartite state $`\varrho `$ then the qantum capacities are equal to corresponding CI capacities
$`Q_{}(\mathrm{\Lambda })=I_{}^X(\mathrm{\Lambda }),Q_{\genfrac{}{}{0pt}{}{}{()}}(\mathrm{\Lambda })=I_{\genfrac{}{}{0pt}{}{}{()}}^{B(A)}(\mathrm{\Lambda }),`$ (19)
$`Q_ø(\mathrm{\Lambda })=I_ø^B(\mathrm{\Lambda }),`$ (20)
$`D_{}(\varrho )=I_{}^X(\varrho ),D_{\genfrac{}{}{0pt}{}{}{()}}(\mathrm{\Lambda })=I_{\genfrac{}{}{0pt}{}{}{()}}^{B(A)}(\varrho ).`$ (21)
Remarks. (i) If we assumed “dual” hashing inequality $`D_{}I^B`$, we would get the same results modulo change $`AB`$. Our choice of $`D_{}I^B`$ is motivated by investigations of Refs. . (ii) It follows that the hashing inequality implies $`I_{}^A=I_{}^B`$. (iii) Our results apply to other kind of supplementary resurces such as e.g. public bound entanglement.
Proof of Theorem 1. We will prove the “Q” part of the theorem. The proof for “D” part is similar. Let $`\{_C^n\}`$ be the optimal error correction protocol for $`\mathrm{\Lambda }`$. Then we have the following estimates
$$I_C^X(\mathrm{\Lambda })\frac{1}{n}I^X(_C^n(\mathrm{\Lambda }^n))\frac{1}{n}I^X(\varrho (d_n,F_n))Q_C(\mathrm{\Lambda }),$$
(22)
where $`\varrho (d_n,F_n)`$ is the twirled state $`\sigma _n=_C^n(\mathrm{\Lambda }^n)`$. The first inequality comes from the very definition of $`I_C^X`$. The second one follows from convexity of $`I^X`$ , and its invariance under product unitary tranformations (as the twirled state is a mixture of product unitary tranformations of the initial one). Finally, since in optimal error correction protocol we have $`F_n1`$, and $`\mathrm{log}d_n/nQ_C`$, we obtain the right-hand-side limit by applying continuity of $`I_X`$ (see lemma).
Proof of Theorem 2. We will also prove only the “Q” part. For $`C\{,,\}`$, consider the following particular error correcting protocol for the channel $`\mathrm{\Gamma }=\mathrm{\Lambda }^n`$. One applies to $`\mathrm{\Gamma }`$ the operation $`_C`$ that produces the state $`\sigma =_C(\mathrm{\Lambda }^n)`$ of maximal attainable coherent information. Subsequently, one performs optimal distillation protocol for the state $`\sigma `$. Then we find
$`Q_C(\mathrm{\Lambda })={\displaystyle \frac{1}{n}}Q_C(\mathrm{\Lambda }^n){\displaystyle \frac{1}{n}}D_C(_C(\mathrm{\Lambda }^n))`$ (23)
$`{\displaystyle \frac{1}{n}}I^X(_C(\mathrm{\Lambda }^n))I_C^X(\mathrm{\Lambda })`$ (24)
where $`C\{,,\}`$; $`X=A,B`$ for $`C=`$, and $`X=A(B)`$ for $`C=()`$. The equality follows from the very definition of $`Q_C`$. The first inequality comes from the fact, that $`Q_C`$ is supremum over all error correction protocols, so it is no less from the rate obtained in the protocol above. The second inequality follows from the hashing inequality, and from the obvious inequality $`D_{}D_{\genfrac{}{}{0pt}{}{}{()}}`$. Finally, the limit is due to the definition of $`I_C^X(\mathrm{\Lambda })`$. The above estimate together with Theorem 1 gives all the desired equalities apart from the one involving $`C=ø`$. That the latter one is also implied by the hashing inequality, it follows immediately from the facts: (a) trivially $`I_ø^BI_{}^B`$; (b) $`Q_{}=Q_ø`$ ; (c) as just proved, the hashing inequality implies $`Q_{}=I_{}^B`$.
Let us now prove yet another important implication of the hashing inequality. Namely, consider the process of discarding information
$$\{p_i,\varrho _i\}\varrho =\underset{i}{}p_i\varrho _i.$$
(25)
In Ref. it was shown that, for a class of ensembles $`\{p_i,\varrho _i\}`$, the amount of information lost in the process is no less than the loss of distillable entanglement $`D_{}`$, and it was conjectured to hold in general. The loss of information is quantified by average increase of entropy, so that the problem is whether the folowing inequality holds
$$\underset{i}{}p_iD_C(\varrho _i)D_C(\varrho )S(\varrho )\underset{i}{}p_iS(\varrho _i).$$
(26)
Note, that for pure state $`\psi `$, $`D_C(\psi )=S(\varrho ^X)`$, where $`\varrho ^X`$ is either of the reductions of $`\psi `$ . Therefore, for pure states $`\varrho _i`$, the inequality reads
$$D_C(\varrho )\underset{i}{}p_iS(\varrho _i^X)S(\varrho )$$
(27)
Applying convexity of entropy we see that the hashing inequality implies the above one. It is interesting, that it does not seem to imply the inequality for impure $`\varrho _i`$’s.
Let us list that recent results concerning entanglement distillation, implying that it is reasonable to conjecture that the hashing inequality (5) holds. (i) In all cases where one has sufficiently tight lower bounds for $`D_{}`$, the inequality is known to be satisfied. For pure states, and other ones with entanglement of formation equal to entanglement of distillation we have $`D_{}=I^B`$. For mixtures of two-qubit Bell states we have $`D_{}I^B`$ by hashing protocol . In particular, for some of them there is equality , while for other ones one has $`D_{}>I^B`$ . (ii) If the hashing inequality is true, then any upper bound for $`D_{}`$ should be no less than $`I^X`$. This was shown for entanglement of formation and, quite recently, for relative entropy of entanglement . We do not know yet, if the inequality holds for the new bound for $`D`$ derived in . (iii) If a state is bound entangled , then we should have $`I^X=0`$. It is indeed the case. According to Ref. , the bound entangled states must satisfy the so called reduction criterion of separability. This implies that the entropic inequality $`S(\varrho )S(\varrho _X)`$ is also satisfied, hence $`I^X=0`$. Thus we see that there is a strong evidence that the inequality is true. We believe that the present results will stimulate to prove (or disprove) it.
The work is supported by Polish Committee for Scientific Research, contract No. 2 P03B 103 16, and by the IST project EQUIP, contract No. IST-1999-11053.
|
warning/0003/math0003168.html
|
ar5iv
|
text
|
# The connectedness of the moduli space of maps to homogeneous spaces
## 0. Introduction
Let $`X`$ be a compact algebraic homogeneous space: $`X=𝐆/𝐏`$ where $`𝐆`$ is a connected complex semisimple algebraic group and $`𝐏`$ is a parabolic subgroup. Let $`\beta H_2(X,)`$. The (coarse) moduli space $`\overline{M}_{g,n}(X,\beta )`$ of $`n`$-pointed genus $`g`$ stable maps parameterizes the data
$$[\mu :CX,p_1,\mathrm{},p_n]$$
satisfying:
1. $`C`$ is a complex, projective, connected, reduced, (at worst) nodal curve of arithmetic genus $`g`$.
2. The points $`p_iC`$ are distinct and lie in the nonsingular locus.
3. $`\mu _{}[C]=\beta `$.
4. The pointed map $`\mu `$ has no infinitesimal automorphisms.
Since $`X`$ is convex, the genus 0 moduli space $`\overline{M}_{0,n}(X,\beta )`$ is of pure dimension
$$\text{dim}(X)+_\beta c_1(T_X)+n3.$$
Moreover, $`\overline{M}_{0,n}(X,\beta )`$ is locally the quotient of a nonsingular variety by a finite group. For general $`g`$, the space $`\overline{M}_{g,n}(X,\beta )`$ may have singular components of different dimensions. Stable maps in algebraic geometry were first defined in \[Ko\]. Basic properties of the moduli space $`\overline{M}_{g,n}(X,\beta )`$ can be found in \[BM\], \[FP\], and \[KoM\]. The following connectedness result is proven here.
###### Theorem 1.
$`\overline{M}_{g,n}(𝐆/𝐏,\beta )`$ is a connected variety.
This result may be viewed as analogous to the connectedness of the Hilbert scheme of projective space proven by Hartshorne. As in \[Har\], connectedness is obtained via maximal degenerations.
Since $`\overline{M}_{0,n}(X,\beta )`$ has quotient singularities, connectedness is equivalent to irreducibility.
###### Corollary 1.
$`\overline{M}_{0,n}(𝐆/𝐏,\beta )`$ is an irreducible variety.
Corollary 1 is easy to verify in case $`X`$ is a projective space. When $`X`$ is a Grassmannian, the irreducibility follows from Strømme’s Quot scheme analysis \[S\]. A proof of Corollary 1 can be found in case $`𝐆=\mathrm{𝐒𝐋}`$ in \[MM\]. For the variety of partial flags in $`^n`$, a proof of irreducibility using flag-Quot schemes is established in \[Ki\]. Results of Harder closely related to Corollary 1 appear in \[Ha\]. There is an independent proof by J. Thomsen for the irreducibility of $`\overline{M}_{0,n}(𝐆/𝐏,\beta )`$ in \[T\].
The moduli space $`\overline{M}_{g,n}(X,\beta )`$ has a natural locally closed decomposition indexed by stable, pointed, modular graphs $`\tau `$ (see \[BM\]). The strata correspond to maps with domain curves of a fixed topological type and a fixed distribution $`\beta _\tau `$ of $`\beta `$. The graph $`\tau `$ determines a complete moduli space of stable maps
$$\overline{M}_{\tau ,n}(X,\beta _\tau )$$
together with a canonical morphism:
(1)
$$\pi _\tau :\overline{M}_{\tau ,n}(X,\beta _\tau )\overline{M}_{g,n}(X,\beta ).$$
A closed decomposition is determined by the images of these morphisms (1). Theorem 1 is a special case of the following result.
###### Theorem 2.
$`\overline{M}_{\tau ,n}(𝐆/𝐏,\beta _\tau )`$ is a connected variety.
Since $`\overline{M}_{\tau ,n}(X,\beta _\tau )`$ is normal in the genus 0 case, we obtain the corresponding corollary.
###### Corollary 2.
Let $`g=0`$. $`\overline{M}_{\tau ,n}(𝐆/𝐏,\beta _\tau )`$ is an irreducible variety.
In particular, all the boundary divisors of $`\overline{M}_{0,n}(X,\beta )`$ are irreducible.
Theorem 2 is proven by studying the maximal torus action on $`X`$. The method is to degenerate a general $`𝐆`$-translate of a map $`\mu :CX`$ onto a canonical $`1`$-dimensional configuration of $`^1`$’s in $`X`$ determined by the maximal torus and the Bialynicki-Birula stratification of $`X`$.
In the genus 0 case, we study the Bialynicki-Birula stratification of $`\overline{M}_{0,n}(X,\beta )`$. The following result is then deduced from the rationality of torus fixed components.
###### Theorem 3.
$`\overline{M}_{0,n}(𝐆/𝐏,\beta )`$ is rational.
The fixed component rationality is equivalent to a rationality result for certain quotients of $`\mathrm{𝐒𝐋}_2`$-representations proven by Katsylo and Bogomolov \[Ka\], \[Bog\]. It should be noted that the fixed components will in general be contained in the boundary of the moduli space of maps – the compactifaction by stable maps therefore plays an important role in the proof.
The rationality of the Hilbert schemes of rational curves in projective space (birational to $`\overline{M}_{0,0}(^r,d)`$) is a consequence of Katsylo’s results \[Ka\] and was also studied by Hirschowitz in \[Hi\].
The main part of this paper was written in 1996 at the Mittag-Leffler Institute where the authors benefitted from discussions with many members. Thanks are especially due to I. Ciocan-Fontanine, B. Fantechi, W. Fulton, T. Graber, and B. Totaro. Conversations with F. Bogomolov were also helpful. B. K. was partially supported by KOSEF grant 1999-2-102-003-5 and POSTECH grant 1999. R. P. was partially supported by NSF grant DMS-9801574 and an A. P. Sloan foundation fellowship.
## 1. The torus action on $`𝐆/𝐏`$
Let $`𝐆`$ be a connected complex semisimple algebraic group. Let $`𝐏`$ be a parabolic subgroup. Select a maximal algebraic torus $`𝐓`$ and Borel subgroup $`𝐁`$ of $`𝐆`$ satisfying:
$$𝐓𝐁𝐏𝐆.$$
Let $`(𝐆/𝐏)^𝐓`$ denote the fixed point set of the left $`𝐓`$-action on $`𝐆/𝐏`$. Three special properties of this $`𝐓`$-action will be needed:
1. The $`𝐓`$-action has isolated fixed points.
2. For every point $`p(𝐆/𝐏)^𝐓`$, there exits a $`𝐓`$-invariant open set $`U_p`$ containing $`p`$ which is $`𝐓`$-equivalent to a vector space representation of $`𝐓`$.
3. Let $`^{}𝐓`$ correspond to an interior point of a Weyl chamber. Then, $`(𝐆/𝐏)^{^{}}=(𝐆/𝐏)^𝐓`$, and the Bialynicki-Birula decomposition obtained from the $`^{}`$-action is an affine stratification of $`𝐆/𝐏`$.
A stratification is a decomposition such that the closures of the strata are unions of strata. In general, the Bialynicki-Birula decomposition obtained from a $`^{}`$-action on a nonsingular variety need not be a stratification.
The claims (i)-(iii) are well known. Only a brief summary of the arguments will be presented here. Let $`𝐖`$ be the Weyl group of $`𝐆`$ relative to $`𝐓`$.
###### Lemma 1.
$`|(𝐆/𝐁)^𝐓|=|𝐖|`$, and $`𝐖`$ acts transitively on $`(𝐆/𝐁)^𝐓`$.
Proof. See, for example, \[Bor\]. ∎
In particular, $`(𝐆/𝐁)^𝐓`$ is a finite set.
###### Lemma 2.
The natural map $`(𝐆/𝐁)^𝐓(𝐆/𝐏)^𝐓`$ is surjective.
Proof. Let $`p(𝐆/𝐏)^𝐓`$. The invariant fiber (isomorphic to $`𝐏/𝐁`$) over the fixed point $`p`$ is a nonsingular projective variety, and hence contains a $`𝐓`$-fixed point by the Borel fixed point theorem (or, alternatively, this is a Hamiltonian action on a compact manifold). ∎
Therefore, $`𝐖`$ acts transitively on the finite set $`(𝐆/𝐏)^𝐓`$.
A representation $`\psi :𝐓\mathrm{𝐆𝐋}(V)`$ is fully definite if there exists a $`^{}`$-basis of $`𝐓`$ for which all the weights of the representation are positive integers. Equivalently, a fully definite representation can be written
$$\psi (t_1,\mathrm{},t_r)v_j=\underset{i=1}{\overset{r}{}}t_i^{\lambda _{ij}}v_j$$
where $`\lambda _{ij}>0`$ for some choice of $`^{}`$-basis of $`𝐓`$ and $``$-basis $`\{v_j\}`$ of $`V`$.
The point $`1𝐆/𝐁`$ corresponding to the identity element of $`𝐆`$ is a $`𝐓`$-fixed point. The $`𝐓`$-action induces a representation
$$\varphi :𝐓\mathrm{𝐆𝐋}(\text{Tan}_1𝐆/𝐁).$$
###### Lemma 3.
The representation $`\varphi `$ is fully definite.
Proof. The natural quotient map $`q:𝐆𝐆/𝐁`$ is $`𝐓`$-equivariant for the conjugation action on $`𝐆`$ and the left action on $`𝐆/𝐁`$. The differential of $`q`$ yields an isomorphism from the Adjoint representation of $`𝐓`$ on $`\mathrm{Lie}(𝐆)/\mathrm{Lie}(𝐁)`$ to $`\varphi `$. $`\mathrm{Lie}(𝐆)/\mathrm{Lie}(𝐁)`$ is the space of positive roots. This $`𝐓`$-representation space has $`n`$ simple roots (where $`n`$ is the rank of $`𝐆`$). All the 1-dimensional representations in $`\mathrm{Lie}(𝐆)/\mathrm{Lie}(𝐁)`$ are non-negative tensor products of these simple roots. Moreover, the $`n`$ weight vectors of these simple roots are independent in the lattice of 1-dimensional representations of the torus $`𝐓`$. Lemma 3 now follows from Lemma 4 below. ∎
###### Lemma 4.
Let $`\psi :𝐓\mathrm{𝐆𝐋}(^n)`$ be an $`n`$ dimensional representation of a rank $`n`$ torus $`𝐓`$. If the $`n\times n`$ matrix of weights is nonsingular, then the representation is fully definite.
Proof. See \[Bi\]. ∎
###### Lemma 5.
The $`𝐓`$-representation $`\text{Tan}_1𝐆/𝐏`$ is fully definite.
Proof. There is a surjection of $`𝐓`$-modules given by the differential $`\mathrm{Tan}_1𝐆/𝐁\mathrm{Tan}_1𝐆/𝐏`$.
###### Proposition 1.
For every $`p(𝐆/𝐏)^𝐓`$, there exists a $`𝐓`$-invariant Zariski open set $`U_p𝐆/𝐏`$ of $`p`$ which is $`𝐓`$-equivalent to a vector space representation of $`𝐓`$.
Proof. By a theorem of Bialynicki-Birula \[Bi\], it suffices to show the tangent representation of $`𝐓`$ is fully definite at $`p`$. This is a consequence of Lemma 5 and the transitivity of the $`𝐖`$-action on $`(𝐆/𝐏)^𝐓`$. (In fact, only definiteness of the tangent representation is needed in \[Bi\].) ∎
Let $`^{}𝐓`$ correspond to an interior point of a Weyl chamber. By the analysis of the tangent representation $`\varphi `$, every point of $`(𝐆/𝐁)^𝐓`$ is an isolated fixed point of $`^{}`$. The equality $`(𝐆/𝐁)^{^{}}=(𝐆/𝐁)^𝐓`$ follows. Since the map $`(𝐆/𝐁)^{^{}}(𝐆/𝐏)^{^{}}`$ is surjective, $`(𝐆/𝐏)^{^{}}=(𝐆/𝐏)^𝐓`$.
For each $`p(𝐆/𝐏)^𝐓`$, let $`A_p`$ be the set of points $`x𝐆/𝐏`$ such that
$$\underset{t0}{lim}tx=p.$$
By Proposition 1, $`A_p`$ is isomorphic to the affine space $`^{r_p}`$ where $`r_p`$ is the number of positive weights in the $`^{}`$-representation $`\text{Tan}_p𝐆/𝐏`$. The set $`\{A_p\}`$ is the Bialynicki-Birula affine decomposition of $`𝐆/𝐏`$. In fact, $`\{A_p\}`$ coincides (up to the Weyl group action) with the (open) Schubert cell stratification of $`𝐆/𝐏`$. This is essentially proven in \[Bor\] for the case $`𝐆/𝐁`$. The general case $`𝐆/𝐏`$ is proven in \[A\]. Therefore, $`\{A_p\}`$ is a stratification.
## 2. The $`^{}`$-flow
Let $`^{}𝐓`$ correspond to an interior point of a Weyl chamber. Let $`s,x_1,\mathrm{},x_l(𝐆/𝐏)^𝐓`$ be the fixed points corresponding to the unique maximal dimensional stratum $`A_s`$ and the complete set of codimension 1 strata, $`A_1,\mathrm{},A_l`$, respectively. The points of $`A_s`$ flow $`(t0)`$ to $`s`$, and the points of $`A_i`$ flow $`(t0)`$ to $`x_i`$. Let $`U=A_sA_1\mathrm{}A_l`$. Since the Bialynicki-Birula decomposition $`\{A_p\}`$ is a stratification, $`U`$ is a Zariski open set with complement of codimension at least 2.
The inverse action of $`^{}`$ on $`𝐆/𝐏`$ is also a torus action on $`𝐆/𝐏`$ with the same fixed point set. Let $`A_s^{},A_1^{},\mathrm{},A_l^{}`$ be the affine strata for the inverse action corresponding to the fixed points $`s,x_1,\mathrm{},x_l`$. Let $`\text{dim}(𝐆/𝐏)=m`$. Since,
$$\text{dim}(A_p)+\text{dim}(A_p^{})=m,$$
$`A_1^{},\mathrm{},A_l^{}`$ are the complete set of $`1`$-dimensional strata for the inverse action. Moreover, the closure $`P_i=\overline{A^{}}_i`$ can contain only the unique $`0`$-dimensional stratum $`A_s^{}=s`$. We have shown the closures $`P_i`$ are contained in $`U`$. Each $`P_i`$ is isomorphic to $`^1`$ (Chevelley \[C\] proves the closed Schubert cells have singularities in codimension at least 2). The intersection pairing
$$P_i\overline{A}_j=\delta (ij)$$
follows from the above analysis. Since the closed strata of the inverse action freely generate the integral homology, the classes
$$[P_1],\mathrm{},[P_l]H_2(𝐆/𝐏,)$$
span an integral basis of $`H_2(𝐆/𝐏,)`$.
Let $`f:C𝐆/𝐏`$ be a non-constant stable map satisfying the following properties:
1. The image $`f(C)`$ lies in $`U`$.
2. $`C`$ intersects (via $`f`$) the divisors $`A_i`$ transversely at nonsingular points of $`C`$.
3. All the markings of $`C`$ have image in $`A_s`$.
If $`[f]`$ represents the class
$$\beta =\underset{i=1}{\overset{l}{}}a_i[P_i]H_2(𝐆/𝐏,),$$
then let $`C`$ meet $`A_i`$ at the $`a_i`$ distinct points
$$\{x_{i,1},\mathrm{},x_{i,a_i}\}.$$
We will study the induced $`^{}`$-action on $`\overline{M}_{g,n}(𝐆/𝐏,\beta )`$ by translation of maps. Let $`F:C_0𝐆/𝐏`$ be the limit in the space of stable maps,
$$F=\underset{t0}{lim}tf$$
where $`t^{}`$.
Define a map $`\stackrel{~}{F}:\stackrel{~}{C}𝐆/𝐏`$ as follows. Let the domain $`\stackrel{~}{C}`$ be:
$$\stackrel{~}{C}=C\underset{i=1}{\overset{l}{}}(_{j=1}^{a_i}𝐏_{i,j}^1)$$
where $`𝐏_{i,j}^1`$ is a projective line attached to $`C`$ at the point $`x_{i,j}`$. Let the markings of $`\stackrel{~}{C}`$ coincide with the markings of $`C`$ (note the markings of $`C`$ are disjoint from the set $`\{x_{i,j}\}`$ by condition (ii)). Define $`\stackrel{~}{F}`$ by $`\stackrel{~}{F}(C\stackrel{~}{C})=s`$ and
$$\stackrel{~}{F}|_{𝐏_{i,j}^1}:𝐏_{i,j}^1\stackrel{}{=}P_i$$
for each $`i`$ and $`j`$.
###### Proposition 2.
If $`f`$ satisfies conditions (i-iii), then the $`t0`$ limit $`F`$ equals the stabilization of $`\stackrel{~}{F}`$.
Proof. Let $`^{}`$ be the punctured holomorphic disk at the origin. Let
$$h:C\times ^{}𝐆/𝐏$$
be the map defined by $`h(c,t)=tf(c)`$. The $`^{}`$-action on $`A_s`$ extends to a map
$$\times A_sA_s$$
since the $`^{}`$-action on $`A_s`$ is a vector space representation with positive weights. The map $`h`$ thus extends to a map
$$h:C\times \{x_{i,j}\times 0\}𝐆/𝐏$$
since the $`f`$-image of $`C\{x_{i,j}\}`$ lies in $`A_s`$. Note,
(2)
$$h(C\{x_{i,j}\},0)=s.$$
After a suitable blow-up
$$\gamma :SC\times $$
supported along the isolated nonsingular points $`\{x_{i,j}\times 0\}`$ of $`C\times `$, there is a morphism $`h^{}:S𝐆/𝐏`$.
The limit as $`t0`$ of $`tf(x_{i,j})`$ equals $`x_i`$. Hence, the exceptional divisor $`C_{i,j}`$ of $`\gamma `$ over $`x_{i,j}`$ connects the points $`x_i`$ to $`s`$ under the map $`h^{}`$. The image $`h^{}(C_{i,j})`$ thus represents an effective curve class containing the class $`[P_i]`$. By degree considerations over all the exceptional divisors $`C_{i,j}`$, we conclude $`h^{}(C_{i,j})`$ is of curve class exactly $`[P_i]`$. As $`P_i`$ is the unique $`^{}`$-fixed curve of class $`[P_i]`$ connecting the points $`x_i`$ and $`s`$,
$$h^{}(C_{i,j})=P_i.$$
We may assume $`S`$ to be nonsingular (away from the original nodes of $`C`$) and each $`C_{i,j}`$ to be a normal crossings divisor – possibly after further blow-ups and base changes altering only the special fiber over $`0`$. We then conclude each $`C_{i,j}`$ has a single component which is mapped to $`P_i`$ isomorphically (and the other components of $`C_{i,j}`$ are contracted).
After blowing-down the $`h^{}`$-contracted components of each $`C_{i,j}`$, we obtain a map $`h^{\prime \prime }:S^{\prime \prime }𝐆/𝐏`$ which is a family of nodal maps over $``$. The fiber of $`S^{\prime \prime }`$ over $`t=0`$ is isomorphic to $`\stackrel{~}{C}`$. Moreover, the condition $`\stackrel{~}{F}(C\stackrel{~}{C})=s`$ follows directly from (2).
The limit stable map $`F`$ is then simply obtained by stabilizing the map $`\stackrel{~}{F}`$. We have carried out the stable reduction of the family of maps $`tf`$ (see \[FP\]). ∎
## 3. Connectedness
Let $`[\mu ]`$ denote the point $`[\mu :CX,p_1,\mathrm{},p_n]\overline{M}_{g,n}(X,\beta ).`$ The stable, pointed, modular graph $`\tau `$ with $`H_2(X,)`$-structure canonically associated to $`[\mu ]`$ consists of the following data:
1. The pointed dual graph of $`C`$:
1. The vertices $`V_\tau `$ correspond to the irreducible components of the curve $`C`$.
2. The edges correspond to the nodes.
3. The markings correspond to the marked points $`p_i`$.
2. The genus function, $`g_\tau :V_\tau ^0`$, where $`g_\tau (v)`$ is the geometric genus of the corresponding component of $`C`$.
3. The $`H_2(X,)`$-structure, $`\beta _\tau :V_\tau H_2(X,\beta )`$, where $`\beta _\tau (v)`$ equals the $`\mu `$ push-forward of the fundamental class of the corresponding component of $`C`$.
Following \[BM\], define $`M_{\tau ,n}(X,\beta _\tau )`$ to be the moduli space of maps $`\mu `$ together with an isomorphism of $`\tau _\mu `$ with a fixed stable graph $`\tau `$. The space $`\overline{M}_{\tau ,n}(X,\beta _\tau )`$ is the compactification via stable maps where the vertices of $`V_\tau `$ may correspond to nodal curves. Note $`M_{\tau ,n}(X,\beta _\tau )`$ may not be dense in $`\overline{M}_{\tau ,n}(X,\beta _\tau )`$.
There is a canonical morphism
$$\pi _\tau :\overline{M}_{\tau ,n}(X,\beta _\tau )\overline{M}_{g,n}(X,\beta ).$$
As $`\tau `$ varies over possible graphs, the images of $`\pi _\tau `$ determine a (closed) decomposition of the moduli space of maps.
Let $`\tau `$ be a stable, pointed, modular graph with $`H_2(𝐆/𝐏,)`$-structure. The connectedness of $`\overline{M}_{\tau ,n}(𝐆/𝐏,\beta _\tau )`$ will now be established.
Proof of Theorem 2. If $`\beta _\tau =0`$, the irreducibility of $`\overline{M}_{\tau ,n}(𝐆/𝐏,\beta _\tau )`$ is a direct consequence of the irreducibility of the corresponding stratum in $`\overline{M}_{g,n}`$ and the irreducibility of $`𝐆/𝐏`$. We may thus assume $`\beta _\tau 0`$.
Fix the $`^{}`$-action on $`𝐆/𝐏`$ as studied in Section 2. Consider an arbitrary point
$$[\mu ]\overline{M}_{\tau ,n}(𝐆/𝐏,\beta _\tau ).$$
By the Kleiman-Bertini Theorem, a general $`𝐆`$-translate $`f`$ of $`\mu `$ satisfies conditions (i-iii) of Section 2. As $`𝐆`$ is connected, $`[\mu ]`$ is connected to its general $`𝐆`$-translate $`[f]`$.
The point $`[f]`$ is connected to the limit:
$$[F]=\text{lim}_{t0}[tf].$$
To prove the connectedness of $`\overline{M}_{\tau ,n}(𝐆/𝐏,\beta _\tau )`$, it suffices to prove the set of limits $`F`$ lies in a connected locus of the moduli space. We will first construct the required connected locus of $`\overline{M}_{\tau ,n}(𝐆/𝐏,\beta _\tau )`$.
The pair $`(\tau ,\beta _\tau )`$ canonically determines a family of maps $`\gamma _b`$ with nodal domains over a base $`bB`$. For $`vV_\tau `$, let $`\beta _\tau (v)=_ia_i^v[P_i]`$. Define the base space $`B`$ as follows:
$$B=\underset{vV_\tau }{}\overline{M}_{g(v),\text{val}(v)+_ia_i^v},$$
where $`\text{val}(v)`$ is the valence of $`v`$ in $`\tau `$ (including nodes and markings). The extra $`_ia_i^v`$ markings each correspond to a basis homology element – with $`a_j^v`$ of these markings corresponding to $`[P_j]`$. The degenerate cases $`\overline{M}_{0,1}`$ and $`\overline{M}_{0,2}`$ in the product $`B`$ are taken to be points. $`B`$ is irreducible and hence connected.
For $`b=_v[b_v]B`$, let
$$\gamma _b:D_b𝐆/𝐏$$
be defined as follows:
1. $`D_b`$ is obtained by attaching the curves $`b_v`$ by connecting nodes as specified by $`\tau `$ and further attaching $`^1`$’s to each of the extra points $`_ia_i^v`$.
2. For each subcurve $`b_vD_b`$, $`\gamma _b(b_v)=s`$.
3. For each $`^1`$ corresponding to $`[P_j]`$, $`\gamma _b(^1)\stackrel{}{=}P_j`$.
The family of maps $`\gamma _b`$ over $`B`$ then defines a morphism (via stabilization):
$$ϵ:B\overline{M}_{\tau ,n}(𝐆/𝐏,\beta _\tau ).$$
Certainly the image variety $`ϵ(B)`$ is connected.
By Proposition 2, the limit $`F`$ is simply the stabilization of $`[\stackrel{~}{F}]`$. Since $`\stackrel{~}{F}=\gamma _b`$ for some $`b`$, the set of limits $`F`$ lies in a connected locus of $`\overline{M}_{\tau ,n}(𝐆/𝐏,\beta _\tau )`$. This concludes the proof of Theorem 2. ∎
Theorem 1 is a special case of Theorem 2 (where $`\tau `$ has a single vertex). Corollary 2 is a simple consequences of Theorem 2.
Proof of Corollary 2. In the genus 0 case, $`\tau `$ is a tree with genus function identically zero. The moduli stack
(3)
$$\overline{}_{\tau ,n}(𝐆/𝐏,\beta _\tau )$$
is constructed as a fiber product over the evaluation maps obtained from the edges of $`\tau `$. We will prove $`\overline{}_{\tau ,n}(𝐆/𝐏,\beta _\tau )`$ is a nonsingular Deligne-Mumford stack by induction on the number of vertices of $`\tau `$.
First, suppose $`\tau `$ has only 1 vertex $`v`$. Then, the moduli stack (3) is $`\overline{}_{0,\text{val}(v)}(𝐆/𝐏,\beta _\tau (v))`$ – a nonsingular moduli stack by the convexity of $`𝐆/𝐏`$.
Next, let $`\tau `$ have $`m`$ vertices and let $`v`$ be an extremal vertex ($`v`$ is incident to exactly 1 edge). Let $`p𝐆/𝐏`$ be a point. By the Kleiman-Bertini Theorem,
(4)
$$\text{ev}_1^1(p)\overline{}_{0,\text{val}(v)}(𝐆/𝐏,\beta _\tau (v))$$
is a nonsingular Deligne-Mumford stack for the general point $`p`$ (and hence every point $`p`$). Let $`\tau ^{}`$ be the graph obtained by removing $`v`$ from $`\tau `$ and adding an extra marking corresponding to the broken node. The moduli stack (3) is fibered over
(5)
$$\overline{}_{\tau ^{},n^{}+1}(𝐆/𝐏,\beta _\tau ^{})$$
with fiber (4). As (5) is nonsingular by induction, the stack (3) is thus nonsingular. This completes the induction step.
Finally, since $`\overline{}_{\tau ,n}(𝐆/𝐏,\beta _\tau )`$ is a nonsingular and connected Deligne-Mumford stack, it is irreducible. ∎
## 4. Rationality
We first review a basic rationality result proven in a sequence papers by Katsylo and Bogomolov \[Ka\], \[Bog\]. Let $`V=^2`$ be a vector space. Let $`a_1,a_2,\mathrm{},a_n`$ be a sequence of positive integers with $`_ia_i3`$. Then, the quotient
(6)
$$(\text{Sym}^{a_1}V^{})\times \mathrm{}\times (\text{Sym}^{a_n}V^{})//\mathrm{𝐏𝐆𝐋}(V)$$
is a rational variety – we may take any non-empty invariant theory quotient. Geometrically, the quotient (6) is birational to the moduli space quotient
(7)
$$M_{0,_ia_i}/\mathrm{\Sigma }_{a_1}\times \mathrm{\Sigma }_{a_2}\times \mathrm{}\times \mathrm{\Sigma }_{a_n}$$
where $`\mathrm{\Sigma }`$ is the symmetric group. Essentially, the rationality of (6) is deduced from rationality in case $`n=1`$ \[Ka\]. Proofs in the $`n=1`$ case may be found in \[Ka\], \[Bog\].
We will also need the following simple Lemma.
###### Lemma 6.
Let $`W`$ be any finite dimensional linear representation of $`𝐀`$ where $`𝐀=\mathrm{\Sigma }_2`$ or $`𝐀=\mathrm{\Sigma }_3`$. Then, $`W/𝐀`$ is rational.
Proof. By the complete reducibility of representations and the fact that a $`\mathrm{𝐆𝐋}`$-bundle is locally trivial in the Zariski topology, it suffices to prove the Lemma in case $`W`$ is an irreducible representation. It is then easily checked by hand the two irreducible representation of $`\mathrm{\Sigma }_2`$ and the three irreducible representations of $`\mathrm{\Sigma }_3`$ have rational quotients. ∎
Proof of Theorem 3. Fix the $`^{}`$-action on $`𝐆/𝐏`$ as studied in Section 2. We first consider the moduli space
$$\overline{M}=\overline{M}_{0,n}(𝐆/𝐏,\beta =\underset{i}{}a_i[P_i])$$
where the property
(8)
$$n+\underset{i}{}a_i4$$
is satisfied.
Let $`\tau `$ be the graph with a single vertex $`v`$ with $`n`$ markings, and let $`\beta _\tau (v)=_ia_i[P_i]`$. Let $`\gamma _b`$ over $`B`$ be the family of maps constructed canonically from $`(\tau ,\beta _\tau )`$ in the proof of Theorem 2. The base $`B`$ is simply:
(9)
$$B=\overline{M}_{0,n+_ia_i}.$$
The map $`\gamma _b`$ over a general point $`bB`$ has no map automorphisms (as $`n+_ia_i4`$). Hence, the image $`ϵ(B)`$ in $`\overline{M}`$ intersects the nonsingular (automorphism-free) locus of the moduli space $`\overline{M}^0\overline{M}`$. Let
$$ϵ(B)^0=ϵ(B)\overline{M}^0,$$
and let $`B^0=ϵ^1(ϵ(B)^0).`$ The map
$$B^0ϵ(B)^0$$
is simply a quotient of $`B^0`$ by the natural $`\mathrm{\Sigma }_{a_1}\times \mathrm{}\times \mathrm{\Sigma }_{a_n}`$ action on (9). By the rationality result (6), $`ϵ(B)^0`$ is rational.
Consider now the $`^{}`$-action on $`\overline{M}^0`$ by translation. As $`\overline{M}^0`$ is a nonsingular, irreducible, quasi-projective variety, we may study the Bialynicki-Birula stratification of $`\overline{M}^0`$. By the proof of Theorem 2, $`ϵ(B)^0`$ is a $`^{}`$-fixed locus which contains the limit,
$$\text{lim}_{t0}t[f],$$
of the general point $`[f]\overline{M}^0`$. By \[Bi\], $`\overline{M}^0`$ is birational to an affine bundle over $`ϵ(B)^0`$. Therefore, $`\overline{M}`$ is rational. The proof of Theorem 3 is complete in case (8) is satisfied.
Next, we will consider the case where the sum (8) is at most 3. In this case, the base $`B`$ is a point. If $`ϵ(B)`$ lies in the automorphism-free locus, the previous argument proving the rationality of $`\overline{M}_{0,n}(𝐆/𝐏,\beta )`$ is still valid. There are exactly four cases in which the point $`ϵ(B)`$ corresponds to a map with nontrivial automorphisms:
1. $`n=0`$, $`\beta =3[P_i]`$.
2. $`n=0`$, $`\beta =2[P_i]+[P_j]`$, $`ij`$.
3. $`n=0`$, $`\beta =2[P_i]`$.
4. $`n=1`$, $`\beta =2[P_i]`$.
Here, the Deligne-Mumford stack structure of these moduli spaces is important. The automorphism group in case (i) is $`\mathrm{\Sigma }_3`$ and in cases (ii-iv) is $`\mathrm{\Sigma }_2`$. In each case, we will show the coarse moduli space $`\overline{M}_{0,n}(𝐆/𝐏,\beta )`$ is birational to a quotient of a linear representation of the corresponding automorphism group.
Consider first the case (i): $`n=0`$, $`\beta =3[P_i]`$. Let $`ϵ(B)=[\gamma ]`$. Let $`[\mu ]`$ denote the unique 3-pointed stable map obtained from $`\gamma `$ by marking each $`^1\stackrel{}{=}P_i`$ by a point lying over $`x_i`$. Certainly, $`[\mu ]\overline{M}_{0,3}^0(𝐆/𝐏,\beta )`$
We will study:
$$N\overline{M}_{0,3}^0(𝐆/𝐏,\beta )$$
where $`N`$ is the component of the locus of transverse intersection of the three divisors $`\text{ev}_1^1(\overline{A}_i)`$, $`\text{ev}_2^1(\overline{A}_i)`$, and $`\text{ev}_3^1(\overline{A}_i)`$ containing $`[\mu ]`$. The torus $`^{}`$ acts on $`N`$ by translation. By an argument exactly parallel to the flow result of Proposition 2, we deduce
$$\text{lim}_{t0}t[f]=[\mu ]$$
for a general element $`[f]N`$. As $`N`$ is a nonsingular, quasi-projective scheme, Theorem 2.5 of \[Bi\] implies that $`N`$ is $`^{}`$-equivariantly birational to the tangent $`^{}`$-representation at $`[\mu ]`$.
There is a $`\mathrm{\Sigma }_3`$-action on $`N`$ by permutation of the markings. The $`^{}`$ and $`\mathrm{\Sigma }_3`$ actions commute. A slightly refined version of Theorem 2.5 of \[Bi\] shows $`N`$ is $`^{}\times \mathrm{\Sigma }_3`$-equivariantly birational to the tangent $`^{}\times \mathrm{\Sigma }_3`$-representation at $`[\mu ]`$. Lemma 7 below explains the refinements of the results of \[Bi\] needed here. $`N/\mathrm{\Sigma }_3`$ is birational to $`\overline{M}_{0,0}(𝐆/𝐏,\beta )`$. Hence, by Lemma 6, Theorem 3 is proven in case (i).
A similar strategy is used in cases (ii-iv). In each of these cases, let $`ϵ(B)=[\gamma ]`$ and let $`[\mu ]`$ denote the rigidification by adding 2 new markings $`,^{}`$ which lie over $`x_i`$. The locus $`N`$ is chosen as the corresponding transverse intersection locus of $`\text{ev}_{}^1(\overline{A}_i)`$ and $`\text{ev}_{^{}}^2(\overline{A}_i)`$ in the maps space with the new markings. $`N`$ is then $`^{}\times \mathrm{\Sigma }_2`$-equivariantly birational to the tangent $`^{}\times \mathrm{\Sigma }_2`$-representation of $`N`$ at $`[\mu ]`$ by the refined Lemma 7. Theorem 3 is then a consequence of Lemma 6 since $`N/\mathrm{\Sigma }_2`$ is birational to the moduli space of maps considered in the case. ∎
###### Lemma 7.
Let $`𝐀`$ be a finite group. Let $`S`$ be a nonsingular, irreducible, quasi-projective scheme with a $`^{}\times 𝐀`$-action and a $`^{}\times 𝐀`$-fixed point $`sS`$. Let $`T_s`$ denote the $`^{}\times 𝐀`$-representation on the tangent space at $`s`$. Suppose the $`^{}`$-action is fully definite at $`s`$. Then, there is $`^{}\times 𝐀`$-equivariant isomorphism between an open set of $`(S,s)`$ and $`(T_s,0)`$.
Proof. We note $`^{}\times 𝐀`$ is a linearly reductive group. By Theorem 2.4 of \[Bi\] for linearly reductive group actions, we may find a third nonsingular irreducible pointed space $`(Z,z)`$ with a $`^{}\times 𝐀`$-action and equivariant, étale, morphisms:
$$\pi _1:(Z,z)(S,s),$$
$$\pi _2:(Z,z)(T_s,s).$$
In the proof of Theorem 2.5 of \[Bi\], such morphisms $`\pi _1`$ and $`\pi _2`$ are proven to be open immersions by a study of only the $`^{}`$-action. Hence, the morphisms $`\pi _1`$ and $`\pi _2`$ are open immersions in our case. By the full definiteness of the $`^{}`$-representation on $`T_s`$, the morphism $`\pi _2`$ is then an isomorphism. ∎
|
warning/0003/nlin0003061.html
|
ar5iv
|
text
|
# 1 Introduction
## 1 Introduction
There is a remarkable difference between the quantum properties of the classically integrable and fully chaotic (ergodic) systems. While in the integrable case the wavefunctions possess an ordered structure, the eigenstates of classically ergodic systems appear random (Berry 1977, Voros 1979) in the semiclassical limit $`\mathrm{}0`$. They can be locally well represented as a superposition of plane waves with equal wavevector magnitude but random phases, leading to the Gaussian distribution of the wavefunction amplitude. The assumption of random phases may, however, break down on dynamical grounds when $`\mathrm{}`$ is not sufficiently small. This leads to the phenomenon of scars, which are the regions of amplified wavefunction amplitude close to the short and weakly unstable classical periodic orbits (Heller 1984).
In the classically mixed systems the energy (hyper)surface is split into both chaotic regions, within which the motion is ergodic, and the regular regions where motion is, as in the fully integrable case, quasiperiodic and confined to invariant tori within the energy surface. We call such systems generic since this is the most general type of Hamiltonian dynamics. The principle of uniform semiclassical condensation (PUSC, see Robnik 1988,1998) states that the quantal phase space distribution of any eigenstate, given by its Wigner function, should be uniformly distributed on a classically invariant object in the phase space when $`\mathrm{}0`$. This object can be either an invariant torus in one of the regular regions in the phase space or a whole chaotic component.
While the PUSC is valid for fully integrable and ergodic systems, its full potential is shown when applied to the generic case. Here it predicts that the eigenstates are separated into regular and irregular ones, depending on whether the classical object onto which their Wigner function condenses is a regular torus or a chaotic component, respectively. This has far reaching consequences leading, for example, to the picture of Berry and Robnik (1984) for the statistics of energy levels, where the regular and chaotic states are assumed to contribute independent level sequences to the total spectrum. This was confirmed by many numerical computations (Prosen and Robnik 1993b, 1994 and 1999, Prosen 1995, 1996 and 1998, Robnik 1998).
In this work we are interested in geometrical and statistical properties of high lying eigenfunctions. It is an extension of previous work done on mixed type systems by Prosen and Robnik (1993c,1994) and by Li and Robnik (1995a,b), but also relates to and further develops the paper by Li and Robnik (1994) concerning the statistical properties of chaotic states, which drastically differ from the properties of regular states in classically integrable systems (Robnik and Veble 1998). The main step forward in the numerical direction was the use of the so-called scaling method, introduced by Vergini and Saraceno (1995), with which we obtained the states with consecutive indices around $`2.510^6`$ and thus it enabled us to go much farther into the semiclassical region. On the other hand, we were able to semiclassically reconstruct the numerically obtained regular states by employing the Einstein-Brillouin-Keller (EBK) torus quantization in section 3.1, where we offer a new approach to this problem, especially on the numerical side. While we were not able to semiclassically reconstruct the chaotic states, we obtained a good prediction of their statistical properties by strictly employing the PUSC (see Robnik 1998) in section 3.2.
## 2 Our catalogues of states
We dealt with a model billiard system obtained by conformally mapping the unit circle with the complex quadratic polynomial, as introduced by Robnik (1983, 1984),
$$zw(z)=z+\lambda z^2,w(z)=x+iy.$$
(1)
The range of $`x`$ at $`y=0`$ inside the billiard is $`x[1+\lambda ,+1+\lambda ]`$. We used the value of $`\lambda =0.15`$, where the classical phase space is roughly equally divided into components of regular and chaotic motion. We chose the Poincaré surface of section (SOS) to lie on the symmetry axis $`y=0`$ with coordinate $`x`$ and the conjugate momentum $`p_x`$ as the parameters of the surface. The intersection of the main chaotic component of our billiard with the SOS is shown in the figure 1. The coordinate $`x`$ is taken relative to the center of the billiard (i.e. it is shifted by $`\lambda `$ w.r.t. $`x=0`$ of equation (1), so that now the range of $`x`$ is $`x[1,+1]`$), while $`p_x`$ is the $`x`$-component of the unit momentum vector.
The quantum mechanics of billiards is described by the Helmholtz equation
$$(\mathrm{\Delta }+k^2)\psi =0,$$
(2)
with the Dirichlet boundary conditions, where $`k^2=2mE/\mathrm{}^2`$. We limited ourselves to the states with even parity with respect to reflection across the symmetry line $`y=0`$.
For each state we calculated the smoothed projection of the Wigner function. The Wigner function of a state $`\psi (𝐪)`$ in the general case of $`N`$ degrees of freedom is defined in the full phase space $`(𝐪,𝐩)`$ as (Berry 1983)
$$W(𝐪,𝐩)=\frac{1}{(2\pi \mathrm{})^N}d^N𝐗\mathrm{exp}(i𝐩𝐗/\mathrm{})\psi ^{}(𝐪𝐗/2)\psi (𝐪+𝐗/2).$$
(3)
In our case the eigenfunctions $`\psi (x,y)`$ generate their Wigner transforms $`W(x,y,p_x,p_y)`$ through (3), where $`N=2`$. In order to compare the Wigner function of a state of our system with the classical SOS plot we took its value on the symmetry line ($`y=0`$) and integrated it over $`p_y`$,
$$\rho _{SOS}(x,p_x)=dp_yW(x,y=0,p_x,p_y).$$
(4)
The result is
$$\rho _{SOS}(x,p_x)=\frac{1}{2\pi \mathrm{}}𝑑X\mathrm{exp}(ip_xX/\mathrm{})\psi ^{}(xX/2,y=0)\psi (x+X/2,y=0).$$
(5)
Here we see the reason for considering the even parity states only, because $`\psi (x,y=0)`$ is exactly zero for odd states, and therefore a different approach must be used to analyze them.
As is well known, the Wigner function is not positive definite but exhibits small oscillations that can blur the overall picture. We chose to smooth the projection of the Wigner function by a suitable Gaussian. It was chosen narrower than the minimum uncertainty Gaussian (in which case the Wigner function becomes the positive definite Husimi distribution) in order not to smooth out too many features, but still wide enough to reduce the oscillations.
The first catalogue of eigenstates and the corresponding smoothed Wigner function projections comprises the first 1000 even states. They were obtained by the conformal mapping diagonalization technique (as described in Robnik 1984). We stress that by this method no levels and eigenstates were lost, as can happen with other approaches such as Heller’s plane wave decomposition method (Heller 1984) and/or the boundary integral method (see e.g. Berry and Wilkinson 1984) etc.
Such a complete catalogue gives us a good picture of the overall behaviour of the system. Many of the states in this low energy region can be associated with the shortest classical periodic orbits. In a mixed type system such as ours, there are both stable orbits that are found within the islands of stability, or unstable ones that lie within one of the chaotic components.
The shortest periodic orbits of our system are shown in figure 2. These are the stable (labeled by 1) and unstable (5) periodic orbit with two bounces, the stable (2) and unstable (6) three bounce periodic orbit and the stable periodic orbit with four bounces (4). As our billiard is convex, there also exists an infinity of stable periodic orbits skipping along the boundary (3) of the billiard, which are associated with and support the whispering gallery modes (Lazutkin 1981,1991, Li and Robnik 1995a).
Examples of the states that correspond to the stable periodic orbits with the indices (1-4) in the figure 2 are shown in the corresponding rows of the figure 3 with the smoothed projections of their Wigner functions shown in the figure 5. Each row shows the states of the same type with increasing energy. The states corresponding to the stable periodic orbits appear consistently and systematically across the catalogue of states. As we will see later, they can be attributed to the quantized tori in the regular regions, and can be found for those tori whose actions satisfy the Einstein-Brillouin-Keller quantization condition (see Robnik 1998). If we compare the Wigner plots of these states with the SOS plot in figure 1, we notice that the areas of greatest intensities of Wigner functions are found within the corresponding islands of stability in the SOS plot.
The states that correspond to the unstable periodic orbits with indices (5) and (6) are shown in the rows of the figure 4 with the corresponding smoothed projections of their Wigner functions given in the figure 6, with their energy again increasing along the rows. The states corresponding to the unstable periodic orbits emerge with varying intensities with respect to the background, and are less frequent with increasing energy. The approximate position of their emergence in the spectrum can be determined by the condition that the classical action along the periodic orbit should be a multiple of Planck’s constant plus the Maslov’s phase corrections due to the caustics formed by nearby trajectories (Robnik 1989). Therefore, in 2-dim billiards, the intervals (either in energy $`E`$ or the cumulative number of states $`𝒩`$) of consecutive re-appearance of the eigenstates of the same type grow as the square root of $`E`$ or $`𝒩`$. Such states with increased intensities of amplitude close to the classically unstable periodic orbits are called scars (Heller 1984, 1986, Bogomolny 1988).
With increasing energy the longer periodic orbits start to manifest themselves in the structure of eigenstates, however, many of the states are becoming increasingly difficult to associate with simple periodic orbits. Two examples of such states are shown in figure 7. The left state is a regular state spanned by a torus in the neighbourhood of the stable periodic orbit with three bounces, characterised by two quantum numbers. It is interesting to note that while the quantum number along the direction of the periodic orbit is quite large (about $`140`$), the transversal quantum number is equal to $`1`$, giving rise to a single nodal line along the direction of the periodic orbit. On the other hand the state on the right can not be associated only with a single unstable periodic orbit, but with a large portion of the chaotic component, this being more prominent in the smoothed projections of the corresponding Wigner function shown in the lower row of the figure 7.
The separation of states into regular and irregular ones becomes fully explicit in our second catalogue of states. The catalogue consists of $`100`$ consecutive states starting at the consecutive index of about $`2.510^6`$. These states were obtained by the scaling method first introduced by Vergini and Saraceno (1995), that enables us to find a few states in the neighbourhood of a chosen wavenumber $`k`$. As this is a diagonalizational method no levels were missed. Almost each level in our small catalogue can be clearly identified as regular or irregular, an idea proposed already by Percival (1973). The only exception in the catalogue is a pair of states lying close together with respect to the mean level spacing, shown in figure 8, where both of the states are superpositions of a regular ($`|\psi _r`$) and irregular ($`|\psi _i`$) state. These two states are close to a degeneracy of two energetically equal but structurally different quantized classical objects.
This mixing is the consequence of the fact that the regular and irregular state are not exact solutions of the Hamiltonian which leads to the matrix element
$$H_{ri}=\psi _r|\widehat{H}|\psi _i$$
(6)
not being equal to $`0`$. When two such states, or, more precisely, their adiabatically corresponding states, are brought close together on the energy scale by varying a parameter of the system (e.g. $`\lambda `$ in our billiard), their eigenenergies do not cross but show a phenomenon of level repulsion (avoided crossing). In the cases when the effect of other levels can be neglected, the smallest energy spacing between the two levels reached is twice the value of the matrix element (6). At this point the eigenstates are exactly the symmetric and antisymmetric superpositions of the well-separated states. If we vary the parameter of the system further, the adiabatic equivalents of the original states will exchange identities.
While the avoided level crossings are typical for irregular states, the regular states do not exhibit the phenomenon of level repulsion amongst themselves (Berry 1983), except possibly on an exponentially small scale (due to the tunneling effects). Furthermore, according to the PUSC, there should be no level repulsion between the regular and irregular states as $`\mathrm{}0`$. As this repulsion is directly connected to the mixing of states, the relative number of mixed states, such as the pair shown in figure 8, is expected to tend to $`0`$ when effective Planck’s constant of the system tends to $`0`$.
More states of both the regular and irregular type from this catalogue with their appropriate analysis will be shown in the next section.
## 3 Analysis of states
The main purpose of our work was to understand to what extent it is possible to describe the structure of individual eigenstates by semiclassical methods. For completely integrable systems there exists the Einstein-Brillouin-Keller method of torus quantization (see for example Berry 1983, Robnik 1998), where each semiclassical eigenstate is spanned by an invariant torus in classical phase space, for which the classical actions are integer multiples of $`\mathrm{}`$ with corrections due to the singularities of projection of the torus onto the configuration space. For fully chaotic systems, and general mixed systems (generic systems), the Gutzwiller periodic orbit theory (Gutzwiller 1990) can in principle be employed to obtain the density of states (the spectrum) and the wavefunctions, to the leading semiclassical approximation.
The semiclassical methods cannot, however, predict individual energy levels within a vanishing fraction of the mean level spacing even in the limit $`\mathrm{}0`$ (Prosen and Robnik 1993a, Robnik and Salasnich 1997). This limits their use in the analysis of statistical properties of spectra. Furthermore, in the chaotic case the levels are also very sensitive to perturbations of the system (Percival 1973). The Gutzwiller approach and method is very useful in the qualitative analysis, and in certain context also quantitative (in describing the collective and statistical properties), but it is still just the leading term in a certain semiclassical expansion, not good enough to resolve the fine structures with sufficient accuracy to make the analysis of individual states and energy levels reliable. Therefore, because of the sensitivity of chaotic eigenstates and the approximating nature of the theory, to some extent the questions about the fine structure of individual chaotic states are irrelevant. It is therefore more appropriate to discuss the statistical properties of chaotic states, which are, however, less sensitive to perturbations.
In a mixed system the phase space is divided into chaotic and regular components. Our work was guided by the principle of uniform semiclassical condensation (PUSC, see Robnik 1998), stating that when $`\mathrm{}`$ tends to $`0`$ the Wigner function of any eigenstate uniformly condenses on an invariant object in phase space. This can be either a torus in the regular region or a whole chaotic component. Each state could thus be labeled as either regular or irregular (chaotic) in the semiclassical limit. By looking at the catalogue of states at high (and to some extent even at low) energies, one can see that this can indeed be done, though there is still the localization phenomenon present due to the still insufficiently low value of the effective Planck’s constant.
### 3.1 Regular states
We start the analysis by considering the regular states. These are the states that can be attributed to quantized tori within the regular regions. For these states we tried to employ the EBK torus quantization. We construct a wavefunction on the torus as a sum of contributions
$$\psi _j(𝐪)=A_j(𝐪)\mathrm{exp}(\frac{i}{\mathrm{}}S_j^{cl}(𝐪)+\varphi _j)$$
(7)
of different projections $`j`$ of the torus onto configuration space. $`S_j`$ is the classical action with respect to some point on the torus and $`A_j^2`$ the classical density of trajectories on this projection. The phase of the wavefunction must change by an integer multiple of $`2\pi `$ when going around any closed contour of the torus. This gives us the quantization conditions
$$I_i=\frac{1}{2\pi }_{\gamma _i}𝐩𝑑𝐪=\mathrm{}(n_i+\beta _i/4)$$
(8)
where $`\gamma _i`$ are the irreducible closed contours on the torus and $`n_i`$ the torus quantum numbers. The integers $`\beta _i`$ are Maslov’s corrections and arise due to the changes of phase $`\varphi _j`$ at the singularities of projection of the torus onto configuration space. At each caustic encountered along the contour $`\gamma _i`$ the wavefunction acquires a negative phase shift of $`\pi /2`$, and shifts by $`\pi `$ when reflected from a hard wall<sup>2</sup><sup>2</sup>2If the contour passes the singularity in the contrary direction to that of the Hamiltonian flow on the torus, the phase shifts are of the opposite sign. From this consideration it follows that $`\beta _i`$ counts the number of caustics plus twice the number of hard walls encountered along the contour.
The main problem arises since, unlike in many completely integrable cases, the transformation to action-angle variables for regular components of mixed type systems is usually not known. The object we are dealing with is only the numerically calculated trajectory, so we must make the best of it.
The task of finding the semiclassical EBK wavefunctions can be divided into two parts. The first one is finding the torus with the desired quantum numbers $`n_i`$, the second one being the construction of its appropriate wavefunction in configuration space.
In our billiard system all the regular tori except the whispering gallery ones wind around a stable periodic orbit with a finite number of bounces $`l_b`$. We choose $`\theta _1`$ to represent the movement of the trajectories along the corresponding periodic orbit while $`\theta _2`$ represents the winding of the trajectories around it. The Maslov index for the contour along $`\theta _2`$ is $`\beta _2=2`$ since there are two singularities of projection (caustics) encountered, while for the contour along $`\theta _1`$ this index is equal to $`\beta _1=2l_b`$. In the case of whispering gallery modes one similarly obtains $`\beta _1=0`$ and $`\beta _2=3`$, if $`\theta _1`$ is taken along the boundary of the billiard and $`\theta _2`$ ’perpendicular’ to it.
We may find the appropriate torus by iteration. We start with a trajectory $`\eta `$ in a regular island on the SOS and follow it sufficiently long until it returns to the desired neighbourhood of the initial point, thus approximately completing $`N_1`$ integer number of cycles in $`\theta _1`$ and $`N_2`$ in $`\theta _2`$ in the time $`T`$. One can obtain the numbers $`N_1`$ and $`N_2`$ by knowing the total number of bounces $`N_b`$ and $`N_c`$ of the caustics encountered during the process through the topological properties of the torus. For a torus winding around a periodic orbit, $`N_1=N_b/l_b`$ and $`N_2=N_c/2`$. The winding frequencies on the torus are then given by
$$\omega _i=2\pi N_i/T.$$
(9)
Finding the number of bounces $`N_b`$ along the trajectory is straightforward. The number of caustics $`N_c`$ may be obtained by starting a trajectory $`\eta ^{}`$ near the original trajectory $`\eta `$ on the same torus (more are needed in case of more than two dimensions). Whenever the two trajectories cross there is a singularity of the density of trajectories and hence a caustic.
One could of course use the monodromy matrix to find the caustics. There are two reasons for not doing so. The first one is the numerical simplicity of our approach. The more important reason is, however, that if the trajectories $`\eta `$ and $`\eta ^{}`$ are taken too close together (infinitesimally close together in the monodromy matrix approach) we may observe caustics due to the possible graining of the desired torus into smaller islands of stability, which quantum mechanics is at the given value of effective Planck’s constant still unable to resolve. A good criterion is that the trajectories should be separated by $`1/n_i`$ for each $`\theta _i`$
We still have to calculate the action integrals on the torus. The integral along $`\theta _2`$ can be calculated by the integral
$$I_2=\frac{1}{2\pi }_{\gamma _2}𝐩𝑑𝐪$$
(10)
along the curve $`\gamma _2`$ formed by the crossing points of the trajectory with the SOS. This integral is just $`1/2\pi `$ of the area of the intersection of the torus with the SOS. Another action integral we can calculate is the action along the chosen orbit $`\eta `$,
$$I=\frac{1}{2\pi }_\eta 𝐩𝑑𝐪.$$
(11)
This integral is the sum of $`I=N_1I_1+N_2I_2`$, so the integral along $`\theta _2`$ is equal to
$$I_1=(IN_2I_2)/N_1.$$
(12)
We must iterate this procedure by choosing different starting points until the proper torus fulfilling the conditions (8) has been found. In general there are as many conditions as there are degrees of freedom. For two dimensional billiard systems whose dynamics is independent of energy, all action integrals can be written in the form
$$I_m=G_m\sqrt{E},$$
(13)
where $`G_m`$ is an integral dependent purely upon the geometry of the classical object in question. In order to obtain the correct quantized torus, instead of solving two separate equations, due to this scaling property we need only to find the appropriate ratio of the (geometric) actions
$$\frac{I_1}{I_2}=\frac{n_1+\beta _1/4}{n_2+\beta _2/4}.$$
(14)
We used the robust bisection method to fulfill this condition since the dependence of actions upon initial conditions is not smooth due to the previously mentioned graining of the tori. Note that this procedure yields only the correct geometry of the quantized torus, the energy of which must still be determined by the quantization conditions 8.
The classical action as a function of time can then be written as
$$_{𝐪_0}^{𝐪(t)}L(𝐪,\dot{𝐪})𝑑t=S(t)=(I_1\omega _1+I_2\omega _2)t,$$
(15)
so one immediately obtains the semiclassical energy as
$$E=\frac{S}{t}=I_1\omega _1+I_2\omega _2.$$
(16)
In our case we took the quantized values for $`I_1`$ and $`I_2`$ and taken into account that for billiard systems the angular frequencies are of the form
$$\omega _m=\lambda _m\sqrt{E},$$
(17)
where $`\lambda _m`$ are frequencies dependent again only upon the geometry of the chosen torus.
In principle one could obtain the semiclassical energy without considering the angular frequencies $`\omega _i`$ by simply finding the appropriate torus through the quantization conditions 8 and reading its energy. In a KAM system such as ours, however, these quantization conditions may be only approximately fullfiled due to the fine structure of the phase space. In contrast to the actions $`I_i`$ the frequencies $`\omega _i`$ do not depend as strongly on the initial conditions (close to the main periodic orbit of an island of stability the transversal frequency typically does not vanish but is characteristic of the periodic orbit). So by taking the quantized values for actions $`I_i`$ and the numerical values for $`\omega _i`$ in equation (16) the best estimate of the semiclassical energy is obtained.
Once the proper torus has been found, we proceed to the second step. We need to span the semiclassical wavefunction (7) on this torus. Again we deal only with a trajectory, so we must find means of representing the wavefunction with it. We may write the wavefunction in the form
$$\psi _j(𝐪)=\underset{T\mathrm{}}{lim}1/T_0^Tdt\delta (𝐪_{cl}(t)𝐪)D_j(𝐪_{cl}(t))\mathrm{exp}(\frac{i}{\mathrm{}}S_j(𝐪_{cl}(t)+\varphi _j),$$
(18)
where $`𝐪_{cl}(t)`$ is the classical trajectory in configuration space. This definition is appropriate only in the sense of integrals of the wavefunction over configuration space. We can hence determine the function $`D_j`$ by integrating the wavefunction over a small volume $`V^{}`$,
$$_V^{}𝑑V\psi _j=\underset{T\mathrm{}}{lim}D_j(𝐪)\mathrm{exp}\left(\frac{i}{\mathrm{}}S_j(𝐪)+\varphi _j\right)\frac{T^{}}{T},$$
(19)
where $`T^{}`$ is the time the trajectory spends in the volume $`V^{}`$. This time is proportional to the density of trajectories $`A_j^2`$ in the volume $`V^{}`$,
$$T^{}=TA_j^2(𝐪)V^{}.$$
(20)
If we want the expression (19) to be consistent with (7),
$$D_j=\frac{1}{A_j}$$
(21)
must hold.
We still have to account for the phase shifts when the trajectory traverses from one projection of the torus onto another. This can be done hand in hand with the estimation of the density of trajectories $`A_j^2`$. Again we start a trajectory $`𝐪_{cl}^{}`$ close to the original one (more trajectories are needed in more than two degrees of freedom). We could again use the monodromy matrix approach, but the same criticism applies here as in the case of finding the appropriate tori. Let us imagine a bundle of trajectories inside a small parallelogram spanned by the three points $`𝐪_{cl}`$, $`𝐪_{cl}^{}`$ and $`𝐪_{cl}+\dot{𝐪}_{cl}\delta t`$, where all three points lie on the given torus. The area of the parallelogram is given by the absolute value of
$$𝐏=\dot{𝐪}_{cl}\times (𝐪_{cl}^{}𝐪_{cl})\delta t.$$
(22)
As we are interested only in the relative sizes of this parallelogram, we may set $`\delta t=1`$. The reciprocal value of this area is proportional to the density of trajectories,
$$A_j^2=\alpha /|P|.$$
(23)
The value of $`\alpha `$ does not change along the trajectory and can be subsequently determined by the normalization of the semiclassical wavefunction. Whenever the value of $`P`$ changes its sign, the trajectory has encountered a caustic and has passed from one projection $`j`$ of the torus to another one $`k`$. At this point the wavefunction acquires a phase shift
$$\varphi _k=\varphi _j\pi /2.$$
(24)
The phase shift is equal to $`\pi `$ if a hard wall is encountered.
The construction (18) is, in its original form, of course unsuitable for numerical computation. We must substitute the $`\delta `$ function in the integral by a function of finite width. The other important reason for doing so is to again smooth the classical behaviour which can have a far more detailed structure than the quantum mechanics can yet resolve. This width may not necessarily need be isotropic. A good estimate for it is again to be of the order of $`1/n_i`$ in coordinates $`\theta _i`$ projected onto configuration space.
We used a Gaussian for the wide delta function. If we adjusted the width and the amplitude of the Gaussian so that it followed the classical density of trajectories, a remarkable similarity of our method with that of Heller’s wavepacket approach was observed (Heller 1991). There are, however, two important differences. The wavepacket approach starts with a wavepacket with the expected value of energy equal to that of the exact wavefunction and then tries to construct its semiclassical approximation. Our approach is independent of the exact eigenstate. All the information we have to supply is the quantum numbers and the geometrical properties of the (projection of the) torus, obtaining from them both the semiclassical energy (16) and the wavefunction. The other important difference is that the wavepacket approach relies on the monodromy matrix of the trajectory, the use of which can be questionable to obtain semiclassical wavefunctions in mixed type systems due to the fine structure of the classical phase space, as was already pointed out before.
We show three examples of the regular states in the figures 9 to 11. For each of the states we present the exact numerical quantum probability density(top left), the probability density of its semiclassical approximation (top right), the classical density of trajectories on the appropriate torus (bottom left) and the smoothed projection of the exact Wigner function (bottom right). The semiclassical wavefunctions shown are remarkable as they possess all of the features of their exact counterparts that are larger than the appropriate wavelength. Note that for each torus there are two characteristic wavelengths since there are two quantum numbers associated with it. As it happens in our case, the two wavelengths can be of different orders of magnitude.
There were, however, some ’regular’ states that we were unable to semiclassically reproduce. The first class of these states can be described as localized chaotic states since their Wigner function clearly shows that they lie in the chaotic region, yet very close to an island of stability, which gives them a regular appearance. The other class are the states whose Wigner transforms lie in the regular regions, but where the primary tori have already been destroyed by the perturbation and now form secondary tori interwoven by small regions of chaotic motion, and even the smoothing of classical dynamics, as described above, fails.
The accuracy of the semiclassical energies that we were indeed able to reproduce may seem remarkable, since the error is approximately $`5`$ units of energy at the energies around $`210^7`$. Such accuracy, however, is still insufficient to perform short range spectral statistics since the mean level spacing in our system is approximately $`8`$ units of energy. This experience is of course in agreement with the proposition and conclusion that the semiclassical methods (to the leading order) cannot resolve the energy spectra within the vanishing fraction of the mean level spacing, and also not the structures of the wavefunctions smaller than de Broglie wavelength (Prosen and Robnik 1993a, Robnik and Salasnich 1997).
### 3.2 Irregular states
While for the regular states it was quite straightforward to find their semiclassical approximations, the nature of irregular states is very much different. The chaotic component of a system does not possess any obvious structure. While Gutzwiller’s approach can yield the properties of a quantum system by a summation over all periodic orbits of its classical counterpart, the relevance of examining individual chaotic states becomes questionable. These states are very sensitive to small perturbations of the system, so in any physical system the individual features of the states are lost when the effective Planck’s constant tends to $`0`$. The features that are insensitive to small perturbations are, however, the statistical properties of spectra and eigenstates.
One measure of statistical properties of the wavefunctions is the wavefunction autocorrelation function,
$$C(𝐪,𝐱)=\frac{\psi ^{}(𝐪^{}𝐱/2)\psi (𝐪^{}+𝐱/2)_{𝐪^{}ϵ(𝐪)}}{\psi ^{}(𝐪^{})\psi (𝐪^{})_{𝐪^{}ϵ(𝐪)}}.$$
(25)
The area of averaging $`ϵ(𝐪)`$ close to the point $`𝐪`$ should be taken such that its linear size is many wavelengths across, however small enough that the local properties of classical mechanics within it are largely uniform.
If one takes the Fourier transform of the Wigner function (3), it is easy to show that
$$W(𝐪,𝐩)\mathrm{exp}(i𝐩𝐱/\mathrm{})d^N𝐩=\psi ^{}(𝐪𝐱/2)\psi (𝐪+𝐱/2).$$
(26)
By knowing the Wigner function of an eigenstate, it is then possible to use this result to calculate its autocorrelation function.
According to the principle of uniform semiclassical condensation, the Wigner function of any chaotic state should uniformly condense on the whole chaotic component when the effective $`\mathrm{}`$ tends to $`0`$. Let us limit ourselves only to the cases of the Hamiltonians with an isotropic dependence upon $`𝐩`$. We can write the semiclassical Wigner function in the form of a conditional delta function
$$W_{𝒟_i}(𝐪,𝐩)=\alpha \delta (\{𝐪,𝐩\}𝒟_i;EH(𝐪,p))$$
(27)
where $`𝒟_i`$ denotes a chaotic component and $`\alpha `$ is the normalization constant. The Fourier transform (26) of this Wigner function is
$$W_{𝒟_i}(𝐪,𝐩)p^{N1}\mathrm{exp}(i𝐩𝐱/\mathrm{})𝑑p𝑑\mathrm{\Omega }_p=\alpha \frac{p(𝐪)^{N1}}{\frac{H}{p}(𝐪,p(𝐪))}_{\mathrm{\Omega }_p𝒟_i(𝐪)}𝑑\mathrm{\Omega }_p\mathrm{exp}(i𝐩𝐱/\mathrm{}),$$
(28)
where $`p(𝐪)`$ denotes the absolute value of momentum at the point $`𝐪`$. The integration over the spatial angle $`\mathrm{\Omega }_p`$ is performed along all the directions of momentum that constitute the chaotic component $`𝒟_i`$ at the point $`𝐪`$. The autocorrelation function is then equal to
$$C_{𝒟_i}(𝐪,𝐱)=\frac{_{\mathrm{\Omega }_p𝒟_i(𝐪^{})}𝑑\mathrm{\Omega }_p\mathrm{exp}(i𝐩𝐱/\mathrm{})_{𝐪^{}ϵ(𝐪)}}{_{\mathrm{\Omega }_p𝒟_i(𝐪^{})}𝑑\mathrm{\Omega }_p_{𝐪^{}ϵ(𝐪)}}.$$
(29)
The averaging area should again stretch across many wavelengths.
If the chaotic component is equal to the whole energy surface, as is the case in completely ergodic systems, in the case of two degrees of freedom one obtains the well known Berry’s result (Berry 1977)
$$C_{\mathrm{e}rgodic}(𝐱)=J_0(p(𝐪)r/\mathrm{}),r=|𝐱|.$$
(30)
However, when the system is of the mixed type the autocorrelation function ceases to be isotropic and acquires contributions of higher order Bessel functions. We can obtain these contributions by rewriting integrals $`_{\varphi _p𝒟_i}f(\varphi _p)𝑑\varphi _p`$ by integrals of the characteristic function $`\chi _{𝒟_i}(\varphi _p)f(\varphi _p)𝑑\varphi _p`$, where in two degrees of freedom the spatial angle is replaced by a simple angle $`\varphi _p`$ and $`f(\varphi _p)`$ is an arbitrary function of $`\varphi _p`$. If we write the characteristic function as a Fourier series,
$$\chi _{𝒟_i}(𝐪;\varphi _p)=\underset{m=\mathrm{}}{\overset{\mathrm{}}{}}\kappa _m^{𝒟_i}(𝐪)\mathrm{exp}(im\varphi _p),$$
(31)
it is quite straightforward to show by using the integral representations of the Bessel functions that
$$C_{𝒟_i}(𝐪,𝐱)=\frac{_{m=\mathrm{}}^{\mathrm{}}\kappa _m^{𝒟_i}(𝐪^{})i^mJ_m(p(𝐪^{})r/\mathrm{})\mathrm{exp}(im\varphi _x)_{𝐪^{}ϵ(𝐪)}}{\kappa _0^{𝒟_i}(𝐪^{})_{𝐪^{}ϵ(𝐪)}},$$
(32)
where $`\varphi _x`$ is the polar angle of the vector $`𝐱`$.
As in the case of the regular states, numerically we cannot deal with components of phase space but with trajectories. So we start a trajectory within the chaotic component $`𝒟_i`$ and not the direction of its momentum $`\varphi _p^j(𝐪)`$ at each passage $`j`$ through the neighbourhood of the point $`𝐪`$. The averaged characteristic function for this neighbourhood can then be represented as
$$\chi _{𝒟_i}(𝐪;\varphi _p)/\kappa _0^{𝒟_i}(𝐪)=\frac{_{j=1}^nd_j\delta (\varphi _p\varphi _p^j)}{_{j=1}^nd_j},$$
(33)
where $`d_j`$ are the lengths covered by the particle in the averaging neighbourhood at the passage $`j`$, and $`n`$ is the number of passages. One can check that this method tends to the proper angular distribution $`\chi _{𝒟_i}`$ as $`n\mathrm{}`$ for any shape of the averaging neighbourhood, if one assumes the homogeneity of trajectories at a given angle (which is true for a small enough neighbourhood, over which the classical phase space picture does not vary). For any incident angle $`\varphi `$ the conditional expected value of a contribution to the equation (33) is proportional to $`d_j𝑑x`$, where $`x`$ is the homogeneously distributed impact parameter (the direction perpendicular to the incident angle). This integral gives just the area of the averaging neighbourhood and is clearly the same for all incident angles.
If the area of averaging is large enough so that the variations in the classical phase space picture become important, the above procedure is still valid as long as the value of the momentum does not change appreciably (as is the case in our billiard system, where between the bounces the momentum remains constant). One can then imagine the large averaging area as being cut into smaller ones within which the above assumptions still hold true.
We compared the autocorrelation functions for a few chaotic states with the semiclassical prediction in figures 12 to 14. The averaging area $`ϵ(𝐪)`$ was taken as a circle of radius $`0.2`$ around the point $`(x,y)=(0.65,0)`$ (the coordinates are as defined in equation (1)). It was taken the same for both the semiclassical prediction and for the numerical results. The averaging radius was taken quite large in order to reduce the localization properties of the wavefunctions, which are still apparent at the values of effective $`\mathrm{}`$ that we were able to obtain. But this radius still has to be taken small enough in order not to completely smooth out the classical dynamics. The agreement with the semiclassical prediction is quite good particularly in figure 14. In all cases it clearly deviates from the Berry’s prediction for fully ergodic systems (30), as it must for mixed systems, and tends towards our semiclassical result. Although in some cases there are amplitude deviations from our prediction, in all of the plots the phase of the numerical correlation function matches the phase of its semiclassical prediction and is significantly different from the phase predicted for the fully ergodic case.
## 4 Discussion and conclusion
As already presented by Prosen and Robnik (1993c), the classification of states into regular and irregular ones is well founded when the effective Planck’s constant tends to $`0`$. Its theoretical foundation is the principle of uniform semiclassical condensation of Wigner functions of eigenstates (Robnik 1988, 1998). This separation is not strictly a semiclassical phenomenon since even some of the lowest levels in our catalogue of low lying states can be classified as either regular or irregular. In the high energy catalogue of states each state can easily be classified as either chaotic or regular, with only one notable exception, where two close lying states are a superposition of a regular and an irregular state. These exceptions are expected to disappear with higher energies. While the states can be separated with respect to classical dynamics, the chaotic states in our high energy catalogue still exhibit the phenomenon of dynamical quantal localization. Their Wigner functions are not uniformly extended over the whole chaotic component, but are significant only on a part of it. This localization is expected to disappear at sufficiently small effective $`\mathrm{}`$, when the quantum mechanical break time $`t_{break}=\mathrm{}/\mathrm{\Delta }E`$, where $`\mathrm{\Delta }E`$ is the mean level spacing, becomes longer than the time for a typical trajectory to explore the whole chaotic component (diffusion time).
The most important part of this work is of course the semiclassical analysis of states. We were able to reconstruct both the semiclassical wavefunction and the semiclassical energy of the regular states by using the EBK quantization. Since our system is of a KAM type, most of the resonant tori are destroyed, forming smaller islands of stability interwoven with chaotic components. The classical mechanics thus shows a rich structure that quantum mechanics at a fixed value of $`\mathrm{}`$ is still unable to resolve. In order to obtain the regular semiclassical wavefunctions, we had to appropriately smooth out this fine classical behaviour as explained in the section 3.1.
We were unable to predict the individual properties of chaotic states. We made a step forward, however, in describing their semiclassical statistical properties. We obtained a semiclassical prediction for the autocorrelation function of their wavefunction, which differs from the one for fully ergodic systems as it is not isotropic. The numerical results confirm this prediction, although there are still localization phenomena at the currently attainable effective $`\mathrm{}`$ that cause deviations from it.
One aspect that needs to be investigated further is the localization properties of the chaotic states. As Casati and Prosen (1998) show in the example of a fully chaotic stadium billiard, in the diffusive regime ($`ϵ`$-stadium, having very large ergodic time), the quantum diffusion is stopped by the cantori in phase space leading to localization. How these and similar ideas translate to the case of a mixed type system remains so far an open question.
## Acknowledgments
We thank Dr. Tomaž Prosen for assistance and advise with some computer programs. This work was supported by the Ministry of Science and Technology of the Republic of Slovenia and by the Rector’s Fund of the University of Maribor.
## References
Berry M V 1977, J. Phys. A: Math. Gen. 10 2083
Berry M V 1983 Chaotic Behaviour of Deterministic Systems (Proc. NATO ASI Les Houches Summer School ed G Iooss, R H G Helleman, R Stora (Amsterdam: Elsevier) p 171
Berry M V and Robnik M 1984 J. Phys. A: Math. Gen. 17 2413
Berry M V and Wilkinson M 1984 Proc. Roy. Soc. London A 392 15
Bogomolny E 1988 Physica 31D 169
Casati G and Prosen T 1999 Phys. Rev. E 59 2516
Gutzwiller M 1990, Chaos in Classical and Quantum Mechanics, Springer New York
Heller E J 1984 Phys. Rev. Lett. 53 1515
Heller E J 1986 Lecture Notes in Physics 263 162
Heller E J 1991 Chaos and Quantum Physics (Proc. NATO ASI Les Houches Summer School ed M-J Giannoni, A Voros and J Zinn-Justin (Amsterdam: Elsevier) p 547
Lazutkin V F 1981 The Convex Billiard and the Eigenfunctions of the Laplace Operator (Leningrad: University Press) (in Russian)
Lazutkin V F 1991 KAM Theory and Semiclassical Approximations to Eigenfunctions (Heidelberg: Springer)
Li Baowen and Robnik M 1994 J. Phys. A: Math. Gen. 27 5509
Li Baowen and Robnik M 1995a J. Phys. A: Math. Gen. 28 2799
Li Baowen and Robnik M 1995b J. Phys. A: Math. Gen. 28 4483
Percival I C 1973 J. Phys. B: At. Mol. Phys. 6 L229
Prosen T 1995 J. Phys. A: Math. Gen. 28 L349
Prosen T 1996 Physica D 91 244
Prosen T 1998 J. Phys. A: Math. Gen. 34 7023
Prosen T and Robnik M 1993a J. Phys. A: Math. Gen. 26 L37
Prosen T and Robnik M 1993b J. Phys. A: Math. Gen. 26 2371
Prosen T and Robnik M 1993c J. Phys. A: Math. Gen. 26 5365
Prosen T and Robnik M 1994 J. Phys. A: Math. Gen. 27 8059
Prosen T and Robnik M 1999 J. Phys. A: Math. Gen. 32 1863
Robnik M 1983 J. Phys. A: Math. Gen. 16 3971
Robnik M 1984 J. Phys. A: Math. Gen. 17 1049
Robnik M 1988 in Atomic Spectra and Collisions in External Fields Eds. K.T. Taylor, M.H. Nayfeh and C.W. Clark (New York: Plenum Press) pp251-274
Robnik M 1989 Bound-State Eigenfunctions of Classically Ergodic Hamilton Systems: A Theory of Scars Preprint Institute for Theoretical Physics, University of California Santa Barbara, unpublished
Robnik M 1998 Nonlinear Phenomena in Complex Systems 1 1
Robnik M and Salasnich L 1997 J. Phys. A: Math. Gen. 30 1711
Robnik M and Veble G 1998 J. Phys. A: Math. Gen. 31 4669
Vergini E and Saraceno M 1995 Phys. Rev. E 52 2204
Voros A 1979 Lecture Notes in Physics 93 326
|
warning/0003/astro-ph0003462.html
|
ar5iv
|
text
|
# A Central Engine for Cosmic Gamma-Ray Burst Sources
## 1 Introduction
Gamma-ray bursts (GRBs) are observed daily from sources at distances extending out to those of the oldest galaxies in our Universe. To account for details of these bursts, “central engines” (CEs) of the GRB sources should have the following properties (see Kluźniak & Ruderman 1998, hereafter KR, for details and references).
(a) Energy. Some CEs must store and release of order $`10^{53}`$ ergs (assuming modest beaming of the energy outflow).
(b) Fluctuations. There are often large temporal variations in the CE power output. A CE should be capable of attaining peak power within tens of milliseconds and exhibiting large fluctuations thereafter. The main power emission is often in sub-bursts between which the CE is relatively dormant, typically for about $`10`$ seconds, but sometimes for as long as several $`10^2`$ seconds or as short as $`10^1`$ seconds.
(c) CE lifetimes, typically seconds to tens of seconds, extend from less than a second to greater than $`10^3`$ seconds. (There is also some indication of an association of greater total energy release with longer CE lifetimes.)
(d) Baryon loading. The energy released from the CE of a GRB source carries with it at most only a tiny baryon load of mass $`\begin{array}{c}<\hfill \\ \hfill \end{array}10^4M_{}`$.
(e) The birth rate of GRB sources $`\begin{array}{c}>\hfill \\ \hfill \end{array}10^6`$/galaxy/yr (see, for example, Böttcher & Dermer 2000).
(f) There is a very great variability among observed GRB events: durations, time scales within a burst, and pulse shape structures, sub-burst numbers, etc., vary so much that one cannot really specify a typical GRB.
The shortest time scales of (b) together with the total energy emission (a) suggest a CE formation involving stellar collapse to a neutron star or to a black hole, or a very tight binary of such collapsed objects, or as part of some exotic supernova which would form such objects. However, the lifetimes (c), baryon loading (d), the commonly observed repeated widely separated fluctuations (b), and perhaps the birthrate (e) may raise special problems for such CE models. Particularly significant is why, if the CEs are collapsed objects whose periods of rotation and vibration are expected to be milliseconds, energy emission from them so often involves several timescales which can be up to $`10^6`$ times longer.
A promising way of constructing CE models based upon collapsed objects, which incorporates this needed family of relatively long time scales, begins by converting the most of the released collapse energy into rotational energy of the collapsed objects. The subsequent transfer of that energy to emitted power in a form useful for ultimate $`\gamma `$-ray production may then be accomplished relatively slowly. It is generally necessary to have CE magnetic fields $`B\begin{array}{c}>\hfill \\ \hfill \end{array}10^{15}`$ G to extract the rotationally stored energy fast enough. Such a CE model was long ago proposed by Usov (1992). A millisecond spin-period pulsar with a magnetic field $`B10^{15}`$ G was assumed to be formed from an accretion induced collapse of a strongly magnetized ($`B10^9`$ G) white dwarf. This simple CE model would be expected to have the needed energy (a), lifetime (c) and baryon loading (d) properties, but a sufficiently high birthrate (e) may be questionable and the required fluctuation property (b) does not seem to be realized.
It has been proposed more recently that very large differential rotation plays an essential role in CE models (KR). One such model has significant similarities to Usov’s proposed millisecond-period “magnetars”, but the initial white dwarf’s pre-collapse history and magnetic field strengths differ, and there are essential differences in what happens within the neutron star and on its surface. This strongly differentially rotating CE would form and evolve in the following, quite different, way.
1) A common “garden-variety” magnetic white dwarf ($`B10^6`$ G) in a tight binary is spun up to its equilibrium spin-period ($`P10^3`$ s) by an accretion disk fed by its companion.
2) The accreting white dwarf is either an evolved one (O-Ne-Mg), or a canonical (C-O) dwarf, with accretion rates such that the accreting white dwarf increases its mass, implodes before its growing stellar mass reaches $`1.4M_{}`$, and collapses to a neutron star.
3) A neutron star is then formed with an initial spin-period $`P10^3`$ seconds, a nearly canonical pulsar polar magnetic surface dipole component $`B_p10^{12}`$ G, and, most importantly, a spin-rate which increases very greatly with distances from the star’s spin-axis. It is this crucial last feature which is the reason for choosing here to discuss this particular CE model from among the previously suggested possibilities for CEs with large initial differential rotation (KR).
4) An interior toroidal field ($`B_\varphi `$) is then stably wound up from the poloidal field ($`B_p`$) by this differential rotation until $`B_\varphi 10^{17}`$ G. After that $`B_\varphi `$ is achieved, the wound-up (and probably slightly twisted) toroid’s magnetic buoyancy for the first time exceeds interior anti-buoyancy forces (from compositional stratification). The buoyant toroid pushes up to the surface by moving parallel to the spin-axis up to, and then partly penetrating the stellar surface, within about $`10^2`$ seconds after its initial release.
5) For as long as some of this magnetic field sticks out of the rapidly spinning neutron star’s surface, this will be an extreme realization of an Usov pulsar, a hyper-magnetar powered by the star’s spin energy. It is, however, extremely transient because of surface movements.
6) This surface dipole field (and higher multipoles) can survive for only a very brief time ($`10^2`$ seconds): it is continually smeared out around the spin-axis, and thus diminished by the strong on-going differential rotation shearing the surface below any protruding field. (There may also be considerable surface field reconnection after this.)
7) After the first break-out of wound-up toroidal field, surface penetration by some of it, and the resulting transient Usov pulsar, a similar wind-up of the $`B_p10^{12}`$ G may begin again as in Step 4, around the same cylinder or a somewhat slower wind-up may exist on some other cylinder). In tiehr case, a new toroid grows until its $`B_\varphi `$ reaches $`B_\varphi (max)10^{17}`$ G when another sub-burst occurs as in Steps 5-6. The characteristic interval between the first and second sub-burst would be
$$\tau _{sb}\frac{2\pi B_\varphi (max)}{(\mathrm{\Delta }\mathrm{\Omega })B_p}10\mathrm{seconds},$$
(1)
where $`\mathrm{\Delta }\mathrm{\Omega }`$ is the spin-frequency difference between the inner and outer parts of the differentially rotating neutron star.
8) The GRB source’s CE finally turns off completely when either of two stages is reached by the engine:
a) the differential rotation ($`\mathrm{\Delta }\mathrm{\Omega }`$) which drives the wind-up of $`B_\varphi `$ becomes so diminished by the conversion of the differential rotation energy into toroidal field energy that it can no longer cause build-up to the critical $`B_\varphi 10^{17}`$ G needed for a pulsar wind sub-burst, or
b) the stellar spin ($`\mathrm{\Omega }`$ of the outer region) becomes so reduced in the transient pulsar phases sustained by it that pulsar wind emission is almost extinguished even if a huge protruding field were to survive.
In this present note, we consider the above GRB source CE proposal in more detail and discuss why and how it should have all of the desired properties.
## 2 Accretion Induced Collapse of Magnetic White Dwarfs to Neutron Stars
Some white dwarfs (WDs) in tight binaries can accrete enough mass from their companions to initiate implosions (because of electron capture by nuclei) as they approach (but just before they reach) their Chandrasekhar limits. After such an implosion begins, there is a competition between energy release from nuclear fusion reactions which act to explode the star and a growing rate of electron capture—which removes pressure support and accelerates collapse. The winner in this competition, which depends upon these relative rates, determines whether such WDs end as Type Ia supernovae (no remnant star) or as neutron stars. Figure 1 shows how the ultimate fate of such accreting WDs is determined by the mass of the WD when accretion begins ($`M`$) and the steady accretion rate ($`\dot{M}`$) which brings it to the initial implosion instability (typically when the accreting stellar mass is about $`1.35M_{}`$). There are three possibilities. 1) The accreting WD has $`M`$ and $`\dot{M}`$ in the cross-hatched region. Then nova explosions continually eject at least as much mass as is accreted between these nova explosions and the implosion mass is not reached. 2) The WD’s $`M`$ and $`\dot{M}`$ begin in the unmarked region. It then ends its life by an accretion induced collapse (AIC) to a neutron star. 3) When $`M`$ and $`\dot{M}`$ are in the dotted region, the accreting WD ends in an explosion with no remnant—a Type Ia supernova (SN)—if its composition is initially C+O. If it approaches implosion with a more evolved O+Ne+Mg composition, sustained $`\dot{M}`$ ultimately causes it to collapse to a neutron star (Nomoto & Kondo 1991; see also Bailyn & Grindlay 1990).
(Below we shall consider WDs with magnetic fields $`B10^6`$ G and, mainly because of that magnetic field, spinning with periods $`P10^3`$ seconds. Neither are of much consequence in the early stages of collapse of these WDs when the ultimate fate of the WD is determined. The magnetic field energy density is approximately $`10^9`$ that of the WD’s rotational kinetic energy, which, in turn, is approximately $`10^3`$ that of its gravitational binding energy. Therefore Figure 1 should not be sensitive to a WD’s possible $`10^6`$ G field or $`10^3`$ second spin period.)
A WD with a $`B10^6`$ G has a relatively modest field among the “magnetic white dwarfs” in the local WD population. Based upon those, they might be expected to number several percent of the WD population. Because of this field, an accretion disk around such a WD, fed by mass pulled from its companion, should spin up the accreting WD to a steady state angular spin rate
$$\mathrm{\Omega }\frac{\dot{M}^{3/7}(GM)^{5/7}}{(BR^3)^{6/7}}\dot{M}_{18}^{}{}_{}{}^{3/7}10^2\mathrm{seconds}^1,$$
(2)
with $`R`$ the WD radius and $`\dot{M}_{18}`$ the WD accretion rate in $`10^{18}`$grams/second ($`\dot{M}_{18}=1`$ when $`\dot{M}=2\times 10^8M_{}`$ per year). For equation (2) to hold, it is assumed that $`\dot{M}`$ is small enough to keep the inner edge of the accretion disk above the stellar surface, i.e.,
$$\dot{M}<\left(\frac{R^5B^2}{GM}\right)^{1/2}10^{20}\mathrm{g}/\mathrm{s}=2\times 10^6M_{}\mathrm{year}^1$$
(3)
The total mass which would have to be accreted to reach the $`\mathrm{\Omega }`$ of equation (2) is about $`10^2M_{}`$. Thus before magnetic WDs with a dipole field $`B10^6`$ G accrete enough to collapse, it is reasonable to expect a good fraction of them, probably most, to have been spun-up to a period $`P_{WD}10^3`$ seconds. After they have collapsed to neutron stars with $`R10^6`$ cm, those neutron stars would then have
$$P_{NS}10^3\mathrm{seconds}.$$
(4)
If magnetic flux is conserved during the collapse, a very plausible approximation because of the short time for collapse (on the order of seconds), these millisecond period neutron stars are formed with (poloidal) fluxes
$$B_{NS}10^{12}\mathrm{G}$$
(5)
If such neutron stars are to be candidates for GRB source CE’s, their formation rate must be $`\begin{array}{c}>\hfill \\ \hfill \end{array}10^6`$yr<sup>-1</sup>-galaxy<sup>-1</sup>. Type Ia supernovae are observed to occur at a rate $`2\times 10^3`$yr<sup>-1</sup>-galaxy<sup>-1</sup>. A plausible guess for the fraction of $`B10^6`$ G exploding WDs among them $`2\times 10^2`$, if this fraction is about the same as that for such fields to be found in the local WD population. The fraction of such moderately magnetized WDs in cataclysmic variables (accreting WDs in tight binaries) is very much greater than this. (However, their number statistics are subject to significant but still unquantified selection effects.) A fraction of $`2\times 10^2`$ assumed above, for WDs which become Type Ia SNe, seems rather conservative based on present knowledge about these WDs. Then if more than only $`3\times 10^2`$ of the WDs which accrete enough to explode as Type Ia SNe had a composition, or a combination of initial $`M`$ and $`\dot{M}`$, to implode to neutron stars, the formation rate for neutron stars satisfying equation (4) and equation (5) would be enough for them to be a candidate population for CEs if the other required properties are met.
The simplest of these and most necessary to satisfy is the maximum energy requirement (a). The spin-energy of a neutron star with an average rotation rate ($`\overline{\mathrm{\Omega }}10^4/`$second) about $`10^{53}`$ ergs is difficult to compare precisely with CE requirements because of the still unknown GRB beaming and some uncertainties in the neutron star equation of state and moment of inertia. We turn next to other special properties of these particular AIC formed neutron stars which determine CE fluctuation timescales (b), lifetimes (c) and baryon loading (d) and variability among the family of these engines (e).
## 3 Initial Differential Rotation of the Neutron Star
The pressure support in a WD whose mass approaches $`1.4M_{}`$ is from extreme relativistic electrons; the star is a $`\gamma =4/3`$ polytrope. Such a star has a central density ($`\rho _c`$) very strongly peaked relative to its average density $`\overline{\rho }`$: $`\rho _c55\overline{\rho }`$ (Shapiro & Teukolsky 1983). The difference in $`\mathrm{\Omega }`$ between the inner and outer parts of the newly formed neutron star will depend on the initial composition of the imploding WD. A C-O WD, when collapse begins from $`{}_{}{}^{16}O+e^{}^{16}C+\nu `$, has $`\rho _c2\times 10^{10}`$g/cm<sup>3</sup>. A O-Ne-Mg WD, whose collapse is initiated by $`{}_{}{}^{24}Mg+e^{}^{24}Na+\nu `$, has $`\rho _c3\times 10^9`$g/cm<sup>3</sup>. If during collapse to a neutron star the angular momentum were to be conserved independently in each of the “rings” of matter circulating around the spin-axis, then the final spin-rates of rings which were originally in from the central region of the collapsing $`\gamma =4/3`$ WD are much less than those of rings collapsing in the outer regions. Roughly, the average neutron star spin
$$\overline{\mathrm{\Omega }}_{NS}\mathrm{\Omega }_{WD}\times \left(\frac{\overline{\rho }_{NS}}{\overline{\rho }_{WD}}\right)^{\frac{2}{3}}10^4/\mathrm{second},$$
(6)
where $`\overline{\rho }_{WD}`$ is the initial average WD density and $`\overline{\rho }_{NS}`$ that of the neutron star. However, for the central region of the WD
$$\rho _c(WD)55\overline{\rho }_{WD},$$
(7)
compared to the very much more modest peaking for the central region of a $`1.3M_{}`$ neutron star,
$$\rho _c(NS)5\overline{\rho }_{NS}.$$
(8)
Insofar as the pressure support of a somewhat cooled neutron star can be approximated as that of a non-relativistic neutron kinetic energy ($`\gamma =4/3`$), a neutron star’s $`\overline{\rho }/\rho _c6`$ (Shapiro & Teukolsky 1983). Additional contributions to stiffening the neutron star’s equation of state (which must be present to increase its maximum neutron star mass from the Oppenheimer-Volkoff $`0.7M_{}`$ to the observed range which is at least twice as large) reduce this ratio further. Therefore the central regions of this newly born neutron star should initially be spinning much less rapidly than most of the matter in that star by a factor of about
$$\left[\frac{\overline{\rho }_{NS}\rho _c(WD)}{\overline{\rho }_{WD}\rho _c(NS)}\right]^{2/3}10^{2/3}0.2,$$
(9)
but, other than the fact that this number is very considerably less than $`1`$, its precise value will not be important in the approximations considered below.
To the approximation that the pressure in the neustrop stellar matter depends only on density, a dynamically stable steady state is finally achieved after fluid flow adjustments give an $`\mathrm{\Omega }_{NS}`$ which depends only on the distance from the spin axis ($`r_{}`$) (the Taylor-Proudman theorem). In this idealization, the newly formed neutron stars rotates on cylinders whose angular speed, $`\mathrm{\Omega }(r_{})`$, increases strongly with increasing $`r_{}`$ because of the differences in density distribution between the $`\gamma =4/3`$ polytrope WD and the neutron star to which it implodes. A crucial question is whether this differential rotation might first have been dissipated during collapse and, if it has survived, what becomes of it in the next $`10^3`$ seconds or so.
During collapse, canonical viscous coupling between distant parts of the star (e.g., by Ekman pumping) is far too weak to be important. However, exchanges of angular momentum by transient energetic neutrino transfer needs special consideration. If this were important, it would be expected to be most efficient for neutrinos whose mean-free path ($`\lambda `$) is of order $`R_{NS}10^6`$ cm. These are emitted in canonical SNEs from the rapidly cooling neutron star remnants (e.g. the neutrinos detected from the SN 1987A explosion) mainly over a 10 second interval. In the implosion leading to the special model CE neutron stars of interest here, most of the released energy should go into stellar rotation rather than thermal heating. Thermal neutrino emission during and just after neutron star formation should, therefore, be much less. If we approximate angular momentum transfers from neutrino transport by assuming, say, a single absorption or scatter before escape (i.e., $`\lambda R_{NS}`$), then the ratio of angular momentum transfer by emitted neutrinos of total energy $`_\nu `$ to the total angular momentum of the neutron star would be $`_\nu /Mc^2`$. Because the difference in angular momentum between the inner and outer parts of the CE neutron star is comparable to the entire stellar angular momentum, the fraction of that difference which would be dissipated before neutrino cooling would also be of order $`_\nu /Mc^2`$. This ratio is certainly less than $`10^1`$ and may be much less.
Another concern is the possibility of convective overturn (Ardeljan et al. 1996), which would mix (Taylor-Proudman) cylinders rotating with different angular speeds. (Fluid movements perpendicular to $`r_{}`$ are not relevant.) But on the short timescales of interest here, where viscosity and thermal conduction are negligible, this should be strongly suppressed by the great increase of angular momentum per unit mass with increasing $`r_{}`$: $`\frac{}{r_{}}(\mathrm{\Omega }^2\rho r_{})`$ is much greater than any plausible convective force density when $`\mathrm{\Delta }\mathrm{\Omega }\mathrm{\Omega }10^4`$/second and $`kT\begin{array}{c}<\hfill \\ \hfill \end{array}10`$ MeV.
During WD collapse there is also a small transfer of angular momentum between different collapsing regions by magnetic fields which couple them. The WD’s (polar) magnetic field which connects differently spinning rings during the collapse would take a time
$$\tau _A\frac{(4\pi )^{1/2}R_{WD}}{B_{WD}}\left(\frac{\overline{\rho }_{WD}}{\rho }\right)^{1/2}10^6\left(\frac{\overline{\rho }_{WD}}{\rho }\right)^{1/2}\mathrm{seconds}$$
(10)
to transfer angular momentum between them, where $`R_{WD}`$, $`\overline{\rho }_{WD}`$ and $`B_{WD}`$ are the WD radius, density and magnetic field at the beginning of the collapse, and $`\rho `$ is the (transient) density at any stage of the collapse. This $`\tau _A`$ is far too long for the magnetic threading during collapse to be a concern in modifying differential rotation.
This leaves the one mechanism for short time scale dissipation of differential rotation which is fundamental to our model for the CEs of GRB sources. Because the differentially rotating cylinders of the newly formed neutron star are coupled by the polar magnetic field ($`B_p`$) in the stellar interior, that field will begin to wind up a toroidal one, $`B_\varphi `$. We turn next to the stability, magnitude, and termination of that wind-up.
## 4 Stability of Toroidal Field Wind-up
The initial $`\mathrm{\Omega }(r_{})`$ in neutron stars formed in this particular AIC genesis is one which grows strongly with increasing $`r_{}`$:
$$\frac{}{r_{}}\mathrm{\Omega }>0.$$
(11)
This rotating fluid is certainly very linearly stable to axisymmetric hydrodynamic perturbations, according to the Rayleigh criterion that the angular momentum $`r_{}^2\mathrm{\Omega }`$ increases outwards. Such a neutron star is also not unstable (or, at least, no instabilities have been discovered) to the so-called Tayler instabilities (Tayler 1973; see also Spruit 1999), which do exist in the special case when $`\mathrm{\Omega }/r_{}=0`$. i.e., when the star is rotating rigidly. However, important recent work in MHD stability theory has shown that powerful instabilities may exist in differentially rotating systems when they contain even relatively weak magnetic fields (Velikhov 1959, Chandresekhar 1961, Balbus & Hawley 1991). Could the differential rotation of a neutron star satisfying equation (11) be unstable to any of these magneto-rotational instabilities? Demonstrations of related MHD instabilities in differentially rotating objects have included non-axisymmetric perturbations, compressibility, and toroidal and poloidal fields (see, for example, Balbus & Hawley 1992 and Ogilvie & Pringle 1996). Magneto-rotational instabilities have been found when angular velocity decreases with $`r_{}`$ but no instabilities have been exhibited for flows satisfying equation (11). Indeed that inequality adds to stabilizing forces in those cases for which the opposite inequality causes instability. Therefore it seems plausible that strong differential rotation satisfying equation (11) will be stable even when the differentially spinning object is threaded by a weak magnetic field: it will wind up an initial poloidal field, which threads it, into a toroidal field until that field becomes unstable because of buoyancy effects, or if magnetic forces grow to exceed gravitational ones.
For the magnetic toroid to become buoyantly unstable, the buoyant forces must overcome whatever anti-buoyant stratification may exist in the star. (For a cold neutron star, the stratification is the compositional one from varying proton/neutron ratios, which can adjust via $`np+e`$ with neutrino emission being too slow to be efficient here.) The wound-up toroidal magnetic field $`B_\varphi `$ would be stable until the buoyancy force density
$$F_b\frac{B_\varphi ^2}{8\pi c_s^2}g,$$
(12)
where $`c_s`$ the speed of sound of the embedding medium, exceeds any anti-buoyancy force density. Wind-up would ultimately increase $`B_\varphi `$ until it reaches a critical value $`B_b`$ at which the buoyancy force is enough to balance the neutron star’s interior anti-buoyancy. This has been estimated to give $`B_b`$ of order $`10^{17}`$ G (KR). (Our estimates in this paper do not depend upon knowing this $`B_b`$ accurately. It is certainly less than the equipartition value for a neutron star whose gravitational binding energy $`10^1M_{NS}c^2`$: $`B_\varphi (equipartition)10^{18}`$ G. The KR estimate of $`B_b2\times 10^{17}`$ G, based solely upon anti-buoyancy forces from compositional stratification in a cooled neutron star, varies only as the square root of that force and is not changed qualitatively by inclusion of other, mainly thermal, contributions).
Thereafter, a (probably slightly twisted) toroid should rapidly rise towards the stellar surface, moving in the direction aligned with both the spin-axis $`\stackrel{}{\mathrm{\Omega }}`$ and $`\stackrel{}{\mathrm{\Omega }}/r_{}`$. The wind-up of the non-axisymmetric poloidal field into a strong toroidal component may introduce some twist into the overwhelming toroidal field. (For the released toroid to rise stably through the neutron star, some twist may be what keeps it from fragmentation by Kelvin-Helmholtz or Rayleigh-Taylor instabilities; see Tsinganos 1980). It is appropriate to emphasize that we do not have a detailed description of the wound-up toroidal bundle’s dynamical evolution after it is released by buoyancy. The buoyancy force can push up the fluid column above it (acting like a plug) to make a surface bulge which can spread horizontally, and/or the approximate axial symmetry of the wound-up toroid may be diminished by some magnetized fluid movement perpendicular to $`\stackrel{}{r}_{}`$.
## 5 Surface Field Penetration: Transient Pulsar (Hyper-Magnetar) Formation
As this toroid rises, the toroidal magnetic flux continues to increase because it is still being wound up from the poloidal component by the continuing differential rotation. As a function of time, $`t`$, measured from the moment when the field first reaches the critical strength $`B_b`$,
$$B_\varphi =B_b+tB_p\mathrm{\Delta }\mathrm{\Omega },$$
(13)
where $`\mathrm{\Delta }\mathrm{\Omega }`$ is again the characteristic difference in angular velocity across the wind-up region. (In our model the initial $`\mathrm{\Delta }\mathrm{\Omega }\mathrm{\Omega }`$.) The buoyance force density, which depends on the square of the toroidal component ($`B_\varphi `$) minus that of the anti-buoyancy force density which balances it when $`B_\varphi =B_b`$, grows nearly linearly with $`t`$:
$$F_b\frac{tB_bB_p\mathrm{\Delta }\mathrm{\Omega }}{4\pi c_s^2}g.$$
(14)
Then the buoyancy timescale for rising up to and penetrating through the stellar surface is
$$\tau _b\left(\frac{24\pi Rc_s^2\rho }{B_bB_pg\mathrm{\Delta }\mathrm{\Omega }}\right)^{1/3}.$$
(15)
For the special AIC formed neutron star, the radius $`R10^6`$ cm, the speed of sound $`c_s10^{10}`$cm/sec, $`B_b10^{17}`$ G, $`B_p10^{12}`$ G, $`g10^{14}`$cm/sec<sup>2</sup>, and $`\mathrm{\Delta }\mathrm{\Omega }10^4`$/second. Then the buoyancy, post break-free, time scale of equation (15) is of the order $`10^2`$ seconds. At the beginning of the wind up, the non-axisymmetric $`B_p`$ is not negligible relative to $`B_\varphi `$. Again during the interval $`0<t<\tau _b10^2`$ seconds when the (positive) difference between the buoyancy and anti-buoyancy forces is small, non-axisymmetric forces from $`B_p`$ may not be entirely negligible compared to other axisymmetric ones acting on the toroids. One effect of this would be a tilt to the rising wound-up toroids accomplished by slightly different forces and fluid movements in directions aligned with $`\stackrel{}{\mathrm{\Omega }}`$ (and $`\stackrel{}{\mathrm{\Omega }}/r_{}`$). Because the toroid will not be exactly axisymmetric, a part of it will ultimately poke through the surface of the star. The field that penetrates the stellar surface will not escape because of the huge conducting mass threaded by the rest of the toroid still below the surface to which it is still strongly attached. This strongly conducting mass remains gravitationally bound to the star.
Therefore $`\tau _b`$ of equation (15) is the turn-on time for the star becoming a Usov type (hyper-magnetar) pulsar with a magnetic dipole field less than, but probably comparable to $`B_\varphi 10^{17}`$ G. (Because $`\mathrm{\Omega }R/c1/3`$, more than just the dipole component of this pulsar may be important in analyzing its properties.)
## 6 Transient Pulsar Termination: Extinction of Surface Multipoles by the Surface Shear from Differential Rotation
Because of the continuing differential rotation of the star and stellar surface, a surface dipole (and all other multipole) cannot survive long beyond the characteristic differential rotation periods involved. If $`\tau _s`$ is the time it take for the differential rotation to shear out the surface dipole (multipole), the characteristic value of the transient surface field would be
$`B_{surf}B_b\times {\displaystyle \frac{\tau _s}{\tau _b}}`$ $`\mathrm{if}\tau _s<\tau _b,`$ (16)
$`B_{surf}B_b`$ $`\mathrm{if}\tau _s>\tau _b.`$ (17)
Figure 2 shows a simplified example of effects of surface shear motion suppression of surface field moments beginning from a north-polar cap of radius $`a`$ at a distance $`d`$ from the spin-axis $`\stackrel{}{\mathrm{\Omega }}`$ together with a similar south-polar cap displaced by the same $`d`$. Because of the different angular velocities of the different parts in both polar caps (increasing with $`r_{}`$), the caps will be deformed into tighter and tighter interwoven spirals extending from $`r_{}=da`$ to $`d+a`$. After many relative wind-ups between $`r_{}=da`$ and $`r_{}=d+a`$, the surface field will be entirely smeared out and cancelled. If the cap radii or distances differed slightly, very little of the surface field would survive as rings around the spin-axis leading to hugely reduced power in the pulsar wind. \[If $`ϵ`$ is a measure of the small difference in the two polar cap radii ($`a`$) or their distances from the spin-axis ($`d`$), then the asymptotic power output after extended shearing is reduced by a factor of order $`\left(\frac{ϵ}{a}\right)\left(\frac{ϵ\mathrm{\Omega }}{c}\right)^210^2`$ relative to that from the initial configuration.\]
The timescale for this surface dipole suppression to be essentially completed is
$$\tau _s\mathrm{several}\times \frac{2\pi }{\frac{\mathrm{\Omega }}{r_{}}a}\mathrm{several}\times \frac{2\pi }{\mathrm{\Omega }}\left(\frac{R}{a}\right)10^2\mathrm{seconds},$$
(18)
where we assume $`\mathrm{\Omega }/r_{}\mathrm{\Omega }/R2\pi /10^3R`$ seconds, and several$`\times R/a10`$. This estimated $`\tau _s`$ should be characteristic for the suppression of all important surface multipoles of typical surface field geometries.
Therefore a $`10^{17}`$ G toroid rises to the surface of the star in $`\tau _b10^2`$ seconds, partly penetrates that surface and expands outside the star. It survives for a time $`\tau _s`$ which is about $`10^2`$ seconds. So long as this field sticks out of the surface of this rapidly spinning neutron star, the star will have the canonical spin-down power and wind emission of a pulsar with dipole field $`B_d\begin{array}{c}<\hfill \\ \hfill \end{array}10^{17}`$ G and $`P10^3`$ seconds.
During this time, the field will be not only be smearing out into a ring but also the north and south poles of the field will be brought into closer and closer contact with each other. This can result in some reconnection. However, reconnection is not the dominant process of field destruction and facilitates field destruction only after the field is already smeared out.
## 7 Sub-burst structures: Energies, Intervals and CE durations
During that brief interval while the millisecond pulsar’s $`B_d10^{17}`$ G dipole field exists, its pulsar wind (consisting of electromagnetic energy, $`e^\pm `$, and some baryons) carries away a sub-burst energy
$$_{sb}\frac{B_d^2R^6\mathrm{\Omega }^4}{c^3}\times \tau _s10^{52}\mathrm{ergs}$$
(19)
(where it is assumed that the dipole $`B_dB_b`$). This is about that required for sub-bursts in observed GRB events. After the wound-up toroidal field breaks away (and penetrates the surface), wind-up by still existing $`B_p`$ could continue to give yet another such wind-up of $`B_\varphi `$ to $`B_b`$, and another transient pulsar, the rewind-up interval, i.e., the time between sub-bursts,
$$\tau _{sb}\left(\frac{2\pi }{\mathrm{\Delta }\mathrm{\Omega }}\right)\frac{B_b}{B_p}\frac{10\mathrm{seconds}}{(B_p)_{12}},$$
(20)
where $`(B_p)_{12}`$ is the poloidal field strength in units of $`10^{12}`$ Gauss.
In addition and perhaps even more important, depending upon the details of $`B_p`$, there would often be significant simultaneous winding up at other rates in several different cylindrical regions within the star. Then sub-bursts from the transient pulsar formation which is the result of the wind-ups would occur from them at different times and could vary enormously with a typical separation of about the $`\tau _{sb}`$ of equation (20) but often much longer or shorter than that.
These sequences of toroidal field wind-up, breaking free, surface penetration of magnetic field $`B_d`$, transient hyper-magnetar pulsar wind, and its suppression can continue only as long as $`\mathrm{\Delta }\mathrm{\Omega }`$ remains large enough to sustain yet another wind-up of $`B_\varphi `$ to $`B_b`$. Typically the energy in the remaining differential rotation within the neutron star would then have to be greater than about $`(1/10)B_b^2R^310^{51}`$ ergs, typically $`10^210^1`$ times the equipartition energy.
The ultimate energy source for the sub-bursts of pulsar wind emission is the spin-energy of the entire star, the differential rotation serving only as the key to open the pulsar wind emission gate for brief intervals. Because this total ($`E_{max}`$) has a maximum of about $`I\overline{\mathrm{\Omega }}^2/210^{53}`$ ergs, it could sustain repeated strong emission activity for a characteristic time (cf. equation (19))
$$\tau \frac{E_{max}}{_{sb}}\times \tau _{sb}\frac{10^2\mathrm{seconds}}{(B_p)_{12}}.$$
(21)
Some CE’s may have a $`\mathrm{\Delta }\mathrm{\Omega }`$ which could give rise only to a single sub-burst with a $`\tau _b10^2`$ second rise time, but the turn-off time ($`\tau _s`$) of such a CE should be stretched considerably for small $`\mathrm{\Delta }\mathrm{\Omega }`$. It is difficult to compare sub-burst turn-off times directly with observations of the emission of $`\gamma `$-rays powered by the pulsar wind sub-burst because these observed $`\gamma `$-rays are created in relativistic expanding regions so very from from the CE.
## 8 Baryon loading and beaming
The transient emission (sub-bursts) from a CE is not what is directly observed in GRBs. The emission radius of observed $`\gamma `$-rays must not be less than about $`10^{15}`$ cm if $`\gamma +\gamma e^++e^{}`$ is not to absorb those $`\gamma `$-rays far above the pair creation threshold. At these large radii, almost all the CE emitted burst energy has all been transferred by expansion into kinetic energy of its co-moving baryon load. To account for the short time scale of many sub-bursts ($``$ a second), the observed emitting region must be expanding relativistically with such a large Lorentz $`\gamma `$ that the rest mass of baryons $`\begin{array}{c}<\hfill \\ \hfill \end{array}10^4M_{}`$ for each $`10^{53}`$ ergs in bursts.
If, as indicated in the previous section, the sub-burst emission is powered by wind from a pulsar with $`\mathrm{\Omega }10^4`$/second and transient dipole $`B\begin{array}{c}<\hfill \\ \hfill \end{array}10^{17}`$ G, the maximum possible baryon outflow from the pulsar should be the maximum for a canonical pulsar with these parameters. The flow rate of nuclei with charge $`Ze`$ out from the stellar surface ($`\dot{N}_Z`$) should then not exceed the Goldreich-Julian limit which would quench that outflow:
$$\dot{N}_z\begin{array}{c}<\hfill \\ \hfill \end{array}\frac{\mathrm{\Omega }^2BR^3}{ec}10^{15}B_{17}M_{}/\mathrm{second}.$$
(22)
This baryon load is negligible in itself and also compared to what the pulsar wind, and especially the first sub-burst, would sweep up from matter around the WD beyond the pulsar.
As noted in Section 6, there may also be small contributions to CE sub-bursts from some reconnection of magnetic loops which extend up from the stellar surface. While these loops are essentially free of baryon loading above the surface (beyond the negligible Goldreich-Julian one), how much they might pull out and up with them during reconnection is far less clear. Of course, each emission sub-burst from the pulsar need not itself be a source of power for ultimate $`\gamma `$-ray emission. Some might become beam dumps for slightly faster, later, much higher Lorentz $`\gamma `$ pulsar-wind sub-burst emissions with much less baryon loading.
If most of the CE emission is in transient pulsar winds from a spinning ephemeral surface dipole (or higher multipoles) as described above, that dipole is mainly orthogonal to the neutron star spin $`\stackrel{}{\mathrm{\Omega }}`$. In the simplest models for pulsar wind emission, with only electromagnetic power in the wind from the star, the emission would be proportional to $`\mathrm{cos}^2\theta `$ with $`\theta `$ the emission angle with respect to the spin axis. Beaming in the spin direction would then be a modest $`3`$ times the emission’s angular average. Higher multipoles ($`\mathrm{\Omega }R/c1/3`$) could significantly increase this beaming.
## 9 Variability among GRB events
Details of emissions from these proposed GRB source CEs should be sensitive to initial properties of the imploding ancestral WDs. There is an important dependence on the ancestral WD’s $`\dot{M}`$ and, especially, details of its magnetic field.
(a) If the accreting WD’s dipole component $`B`$ is much less than $`10^6`$ G, its steady state spin period from equation (4) becomes so small that it would not quickly and simply form a nuclear density neutron star. Centrifugal forces would stop much of that collapse before it had evolved that far (forming what T. Gold called a “fizzler”). Ultimate formation of a neutron star would be achieved only after angular momentum had been removed \[perhaps mainly into a surrounding disk and/or, as $`P10^3`$ seconds (the maximum spin rate of an axisymmetric neutron star), through the transient formation of a Jacobi ellipsoid and subsequent powerful gravitational radiation\]. There is no obvious reason for the final distribution of differential rotation after such a genesis being the same as that of the proposed “canonical” CE from the AIC of a magnetic WD with $`B10^6`$ G. There are expected to be many more WDs with smaller $`B`$ than those with $`B10^6`$ G.
(b) If the accreting WD’s dipole $`B`$ greatly exceeds $`10^6`$ G, the steady state $`P`$ of equation (4) increases. If any of these more slowly spinning WDs were to collapse to a GRB source CE, that CE’s spin energy could not support burst events with total energy near the maximum $`10^{53}`$ ergs. However, if $`\dot{M}`$ could greatly exceed the Eddington limit (indicated in Fig. 1) from a sufficiently massive Roche Lobe overflow from the companion, $`P(NS)10^3`$ seconds might be achieved even from a $`B10^9`$ G WD, as assumed by Usov in his CE model.
(c) The strongest sensitivity of CE emission pulse structures would probably be to details of $`B_p`$, the initial polar field in the newly formed differentially rotating neutron star, because of possible magnification of small initial $`B_p`$ differences to large ones in the wound-up $`B_\varphi `$. For example, $`B_p`$ details determine the number density of toroids which begin simultaneous wind-up in different cylindrical regions around the spin-axis; these wound-up toroids overcome anti-buoyancy constraints and break free at different times. Locally, $`B_p`$ would be expected to change somewhat during these releases and any rewind-ups so CE emission pulse shapes would not be repetitive during a GRB event.
## 10 Discussion
The required properties of GRB source CEs (summarized in the Introduction) are total energy stored and emitted (a), peak power and fluctuations within a given burst event (b), CE lifetimes (c), maximum baryon loading in the CE emission (d), CE birthrate (e) and very large variability among different CEs (f). None of these seem an embarrassing problem for the proposed model CE genesis, structure and dynamics outlined in the Introduction and described in Sections 2-9. Indeed each seems a rather expected consequence. However, a crucial point which should be considered further is the absence (so far) of any demonstrated instability in the wind-up of the toroidal field for of order $`10^4`$ turns (in about $`10`$ seconds) by the much more energetic initial differential rotations in the neutron star.
A second related, but less crucial, question is the robustness of our presumption that during and after such toroidal wind up and release the initial much smaller poloidal field component of the differentially rotating neutron star is not hugely increased. If this does not turn out to be an adequate approximation, the often observed sub-burst multiplicity could still come from toroidal field wind-up and break-away in different cylindrical regions with different wind-up times rather than from long time delays for rewinding $`B_\varphi `$. Of course, because of the very great variability within the family of GRB events, neither mechanisms may hold in all, or perhaps not even in most cases, but at least one of them should certainly not be uncommon.
Finally there is our unproven assumption that large toroidal field bundles wound-up by differential rotation can overcome anti-buoyancy restraints and break free as a unit (or almost so). If, instead, buoyant toroidal field continually dribbled up and out to support a steady state in which increasing $`B_\varphi `$ from wind-up is balanced by a that loss, there would be no strong fluctuations in CE output. Instead a CE would be an Usov-like pulsar with emission decreasing monotonically after the first emission maximum is reached. This is a generic problem for many kinds of CE models. Why does the CE depart so far from a steady equilibrium that stored energy is released in huge sub-bursts (which are often separated by very many characteristic engine periods) rather than in smoother continuous steady way? Here too such a question needs further investigation.
## 11 Acknowledgments
We are happy to thank J. Applegate, E. Costa, C. Knigge, J. Pringle, M. Rees, E. Spiegel, H. Spruit, and J. Stone for informative conversations, and the Institute of Astronomy (Cambridge) and the Aspen Center for Physics for their hospitality while much of this work was begun.
|
warning/0003/astro-ph0003211.html
|
ar5iv
|
text
|
# THE HII REGION KR 140: SPONTANEOUS FORMATION OF A HIGH MASS STAR
## 1 INTRODUCTION
Massive OB stars are almost always found in clusters. In fact, accumulating observational evidence suggests that most stars, regardless of mass, actually form as members of some kind of group, cluster, or association. While there has been a lot of work to understand the processes involved in the formation of a single star (e.g., Shu, Adams, & Lizano 1987), theories of cluster formation are still in their infancy (see the recent reviews by Elmegreen et al. 2000, and Clarke, Bonnell & Hillenbrand 2000). Nevertheless, many important issues have already been identified. Foremost among these is the question of whether the formation of clusters, particularly ones with OB stars, are triggered by an external agent (Elmegreen 1992) such as expanding HII regions (Elmegreen & Lada 1977), or colliding molecular clouds (Loren 1976, 1977; Scoville, Sanders & Clemens 1986; Usami, Hanawa, & Fujimoto 1995). Despite triggered or sequential star formation being theoretically and intuitively appealing, a major problem is to determine unambiguously a cause and effect relationship, because of the long time scales (and time lags) of the processes involved. There is observational evidence for regions that have had triggered star formation, both on large scales (e.g., IC 1396, Patel et al. 1998) and on small scales (e.g., IC 1805, Heyer et al. 1996). However, there are other young star forming regions, like Taurus, where no evidence of a trigger can been found. These are often in modest-sized molecular clouds and contain only lower mass stars.
The Perseus Arm star forming regions W3/W4/W5 (Westerhout 1958) have been studied extensively over the last twenty years (e.g., Lada et al. 1978; Braunsfurth 1983; Digel et al. 1996; Normandeau, Taylor & Dewdney 1997; Heyer & Terebey 1998). They are often considered to be the archetypical examples of how the formation of massive star clusters can be triggered by the influence of other nearby clusters. For example, W3 is thought to be have been triggered by the expansion of W4 (Dickel 1980; Thronson, Lada, & Hewagama 1985; van der Werf & Goss 1990) and there is evidence that the expansion of W5 is also triggering star formation (Vallée, Hughes, & Viner 1979; Wilking et al. 1984).
In new high resolution multiwavelength (radio and mid-infrared) data of the W3/W4/W5 complex from the Canadian Galactic Plane Survey (CGPS; English et al. 1998) we have identified a star forming region containing a single O star. Figure 1 shows a 1420 MHz continuum image from the CGPS pilot project (Normandeau et al. 1997). The circled area is the HII region in question, KR 140 ($`l=133.425^{}`$, $`b=0.055^{}`$; Kallas & Reich 1980). This region appears to be completely separate from the vigorous star formation going on nearby in W3, although it is in the same Perseus arm molecular complex (see context in CO images in Heyer et al. 1998). What is unusual is that this massive star seems to have been formed spontaneously.
In this paper we present and analyze the multiwavelength data on KR 140 in order to quantify the properties of this region of spontaneous massive star formation. As described in § 2 these data have sufficient resolution (1) to resolve this HII region for the first time and reveal a fairly symmetrical structure. Based on first impressions, we thought that KR 140 might prove to be a “textbook” spherical HII region. Instead, we find that KR 140 is likely a bowl-shaped HII region (an example of a blister geometry: Israel 1978, Yorke et al. 1989; § 3 and § 4.2.2, § 6.4). We have found that the HII region is kept ionized by an O8.5V(e) star, VES 735, and is at a distance of 2.3$`\pm `$0.3 kpc from the Sun (Kerton et al. 1999).
We analyze the ionized component of KR 140 in § 4 and make various estimates of the age in § 5. The dust, molecular, and atomic components of KR 140 are examined in § 6, § 7, and § 8, respectively. In § 9 we discuss star formation in KR 140 in the context of the Perseus arm and the possible accompanying cluster.
## 2 THE MULTIWAVELENGTH DATA
KR 140 was initially cataloged in a 1420 MHz radio continuum survey of the northern Galactic plane with the Effelsberg 100 m telescope by Kallas & Reich (1980). At the available resolution of 9 KR 140 was barely resolved, with a reported diameter of about 11. Although KR 140 was measured in other subsequent single-dish surveys (§ 4.2) it had never been examined with a radio interferometer to provide the angular resolution necessary for the present study.
The data analyzed here include 1420 MHz ($`\lambda `$ 21 cm) and 408 MHz ($`\lambda `$ 74 cm) continuum images from the Dominion Radio Astrophysical Observatory (DRAO) Synthesis Telescope (Roger et al. 1973; Veidt et al. 1985). An HI 21-cm line data cube is also available, with spectral resolution 2.64 km s<sup>-1</sup> and channel spacing 1.65 km s<sup>-1</sup>. These DRAO data were obtained during the pilot project of the CGPS (Normandeau et al. 1997), and made into a full 8$`{}_{}{}^{}\times 5^{}`$ mosaic of the W3/W4/W5 star forming regions.
We also made use of a CO ($`J=10`$, 115 GHz, $`\lambda `$ 2.60 mm, spectral resolution 0.98 km s<sup>-1</sup>) data cube from the Five Colleges Radio Astronomy Observatory (FCRAO) obtained in the complementary survey described by Heyer et al. (1998). To examine the dust components we processed Infrared Astronomical Satellite (IRAS) data to make HIRES mosaics at 12, 25, 60, and 100 $`\mu `$m. Independent processing of the entire Galactic plane data at 60 and 100 $`\mu `$m has been released as the IRAS Galaxy Atlas (IGA, Cao et al. 1997) and we have completed a complementary project at 12 and 25 $`\mu `$m (Mid-Infrared Galaxy Atlas, MIGA, Kerton & Martin 2000).
An important feature for subsequent analysis with these data sets is the common relatively high angular resolution achieved. The DRAO synthesis telescope images have a resolution of $`1^{}\times 1.14^{}`$ at 1420 MHz (the 408 MHz DRAO data have proportionately 3.5 times lower resolution), and the CO images are at a resolution of 50<sup>′′</sup> (beam sampled). The HIRES images have non-circular beams of size about 2 at 100 $`\mu `$m, 1 at 60 $`\mu `$m, and somewhat less at 12 and 25 $`\mu `$m. For detailed intercomparison, pairs or groups of images have been convolved to the same beam shape.
## 3 A MODEL FOR THE 3-DIMENSIONAL MORPHOLOGY
Analysis of an HII region benefits from a knowledge of its three-dimensional structure, but often the observed two-dimensional morphology is too complex to interpret. Fortunately, this is not the case for KR 140. Figure 2 shows the four HIRES images of KR 140 at 12, 25, 60, and 100 $`\mathrm{\mu m}`$. As discussed in § 6, the 100 $`\mathrm{\mu m}`$ image best describes the dust distribution in KR 140. The 100 $`\mathrm{\mu m}`$ HIRES image (Figure 2d) shows that most of the dust in the nebula has been swept into a shell-like structure and that we are observing a region which is close to being circularly (axially) symmetric. In the 1420 MHz image (Figure 3) KR 140 is also seen to be fairly circular with a central depression. As highlighted by the contour map, there are three high intensity areas – two “eyes” and a “mouth”. However, where the “nose” would be is a lower intensity area. The simplest interpretation appears to be a three-dimensional structure with a central hole (§ 5.2.2). We have identified VES 735 as the exciting star of this HII region (Kerton et al. 1999); it lies projected close to this lower intensity area of HII emission, though not coincident with it.
Since we have also classified the exciting star, we can draw further conclusions about the geometrical structure of KR 140 from a global analysis (§ 5 in Kerton et al. 1999). Even after taking into account the emission from dust that was not observed by IRAS, we find that the total luminosity of warm dust in the nebula is far less then the bolometric luminosity of the ionizing star VES 735 and that the covering factor of the dust is only about 0.4 – 0.5 (see § 6.4). Furthermore, the total radio flux at 1420 MHz is lower than expected for an ionization bounded nebula surrounding VES 735, with a similar implied covering factor (§ 4.2.2). This together with geometrical information to be discussed (e.g., § 4.1) implies that KR 140 is a bowl-shaped HII region (cf. Roger & Irwin 1982), rather than a classical Strömgren sphere, a not unexpected morphology when the O star is close to the edge of its parent molecular cloud (Yorke et al. 1983). This model is perhaps the simplest geometry that is consistent with the derived covering factors, but our data cannot rule out other, more complicated geometries, such as a broken shell or shredded and dissipated Strömgren sphere. Observations of a champagne flow (Tenorio-Tagle 1979), perhaps via Fabry-Perot imaging, could help establish that KR 140 is truly a blister region.
The circular symmetry observed in the dust shell (and in the ionized gas) implies that we are observing KR 140 almost face on (i.e., the opening of the bowl is oriented almost along the line-of-sight). For example, see the simulated radio maps of the R2 model by Yorke et al. (1983). The CO data and extinction measurements to VES 735 both suggest that the molecular cloud is behind the star, with the inferred opening toward us (§ 7.1).
## 4 KR 140 FROM THE PERSPECTIVE OF THE RADIO CONTINUUM
### 4.1 Size and Structure
The basic radio morphology appears to be that of a limb-brightened hemispherical shell (most of the free-free emission would be from the dense and thin ionization-bounded zone). The last contour shown (at $`T_b`$ = 5.0 K) in Figure 3 has a diameter of 8.5 which for a distance $`d=2.3`$ kpc corresponds to a physical diameter of about 5.7 pc. It is expected that neutral material could extend well beyond the HII region, and indeed there is both dust (§ 6.2) and CO gas (§ 7.1) around KR 140. The free-free surface brightness falls off most quickly (the radio contours are more tightly spaced) in the west-northwest part of the nebula. A sharper ionization front would occur in more dense (pre-existing) material.
Other internal detail might be interpreted as variations in distance between the star and the background ionization front, with lower surface brightness features originating from sectors further from the star; such an interpretation of radio surface brightness has been used to construct a topological map of the ionized surface of the Orion molecular cloud behind the Trapezium (Wen & O’Dell 1995). If we use the 1420 MHz surface brightness at the pixels near the projected position of VES 735 and the known properties of this O8.5V exciting star, the distance that the star must be from the back of the nebula turns out to be only 2.14 pc. Together with the above estimate of the diameter across the line of sight, this is consistent with the idea that KR 140 is similar to a hollowed-out hemispherical bowl.
### 4.2 Physical Properties
In this section we use the 1420 MHz data to derive a number of physical properties of the nebula. For ease of reference, Table 1 summarizes these properties, some of the properties of the exciting star, and other quantities that we later derive from the infrared and CO data.
#### 4.2.1 Emission Measure, Radio Flux, Electron Density, and Ionized Mass
The KR 140 emission is optically thin at 1420 MHz with a peak brightness temperature about 10 K. For an optically-thin nebula, the specific intensity (surface brightness) $`I_\nu `$ and brightness temperature $`T_{b\nu }`$ for thermal radiation are (Osterbrock 1989):
$$I_\nu 2\nu ^2kT_{b\nu }/c^2=j_\nu n_in_e𝑑sj_\nu E,$$
(1)
where $`j_\nu `$ is the free-free emissivity, $`n_i`$ is the ion density, $`n_e`$ is the electron density, $`ds`$ measures distance along the line of sight, and the integral $`E`$ is called the emission measure. For $`j_\nu `$ of a pure hydrogen nebula (a good approximation; see § 4.2.2 for models which include ISM abundances), $`E=5.77\times 10^2T_{b\nu }(T_e/7500\mathrm{K})^{0.35}(\nu /1420\mathrm{M}\mathrm{H}\mathrm{z})^{2.1}`$ cm<sup>-6</sup> pc. The peak emission measure in the map is about 6000 cm<sup>-6</sup> pc and after background subtraction $`E=2050`$ cm<sup>-6</sup> pc. After subtracting the local Galactic background we find the total flux ($`F_\nu =I_\nu 𝑑\mathrm{\Omega }`$) to be $`2.35\pm 0.05`$ Jy at 1420 MHz. This is in good agreement with the value $`2.3\pm 0.2`$ Jy reported by Kallas & Reich (1980). Becker, White, & Edwards (1991) detected this HII region in a 6 cm survey with the NRAO 91 m. Their source, designated \[BWE91\]0216+6053, had $`F_\nu =2.5`$ Jy. Using the same telescope and frequency Taylor & Gregory (1983) and Gregory & Taylor (1986) record this source as GT 0216+608 in their survey; in their analysis they treated all sources as point sources and so find systematically low fluxes for extended nebulae like this (648 mJy for KR 140).
In the optically thin limit the theoretical expectation for the spectral index for thermal radiation (defined as $`I_\nu \nu ^\alpha `$) is $`\alpha 0.1`$ (Oster 1961; Gordon 1988). The spectral index was measured between 1420 and 408 MHz. First, the distribution across the nebula was determined using L. Higgs’ “specmap” routine (Zhang & Higgs 1997), which evaluates the spectral index for the spatially variable component of the emission (this approach has the advantage of circumventing uncertainties in background subtraction in each image); for this procedure, we convolved the 1420 MHz image to the lower resolution of the 408 MHz image. The typical value of $`\alpha `$ is $`0.02\pm 0.08`$. We confirmed this by measuring the total background-subtracted flux at 408 MHz as well. Considering there is a 10% uncertainty in the flux calibration in the 408 MHz pilot project data at the position of KR 140 (Taylor 1999), the derived spectral index is in satisfactory agreement with theory.
Depending on details of the geometry, the appropriate pathlength for estimating the density would be of order the observed radius; with this choice we find an rms $`n_e=27`$ cm<sup>-3</sup>. The mass of radio-observable ionized gas (allowing for He) is 163 (27 cm<sup>-3</sup>/$`n_e`$$`M_{}`$. For comparison, the exciting star of KR 140 has a mass of 25 $`M_{}`$ (Kerton et al. 1999).
#### 4.2.2 Number of Ionizing Photons
In ionization equilibrium the number of ionizations (where H is the predominant species) that occur each second in the nebula is equal to the number of (H<sup>+</sup>) recombinations each second, both locally and globally. An ionization-bounded nebula occurs when there is sufficient material to intercept all of the ionizing photons ($`\lambda 912`$ Å), $`Q(H^0)`$, emitted by the star per second (Osterbrock 1989). Formally, the spatially-integrated radio emission of an ionization-bounded nebula is an effective calorimeter for $`Q(H^0)`$:
$$fQ(H^0)=n_en_{\mathrm{H}^+}\alpha _B𝑑V=\alpha _BF_\nu d^2/j_\nu $$
(2)
where $`\alpha _B`$ is the case B hydrogen recombination coefficient of H<sup>+</sup> ($`3.3\times 10^{13}`$ cm<sup>3</sup> s<sup>-1</sup> at 7500 K; Storey & Hummer 1995), $`dV`$ is the volume element, and we have used $`dV=dsd\mathrm{\Omega }\times d^2`$ and Equation (1) to find the right hand side. Here $`f`$ is the product of a number of correction factors including $`f_{ci}`$, the all-important covering factor by an ionization front: if the star is only partially surrounded by gas then a fraction of the ionizing photons will escape and a lower $`F_\nu `$ will result.
Equation 2 is based on recombination and emissivity for a pure H nebula at a fixed temperature and so $`f`$ is made up of a number of factors which modify the expected radio continuum flux for a given $`Q(H^0)`$. We used the spectral synthesis code Cloudy (Ferland et al. 1998) to analyze the observations of KR 140 in the radio, particularly the correction factors which comprise $`f`$. We constructed a number of toy models for differential comparisons, using a 35,500 K blackbody with $`\mathrm{log}[Q(H^0)]=48.06`$ and $`n_e=50`$ cm<sup>-3</sup>; see Table 2. The differential results are appropriate for any low-density HII region heated by a late O-type star.
The temperature of the nebula enters in the equation through the $`\alpha _B`$ and $`j_\nu `$ factors, although the effect on the ratio is very small. Comparing models 1 and 5 we estimate $`f_T=0.99`$.
The effect of the addition of other elements to the nebula is to increase the total radio flux. This effect is primarily due to the presence of He<sup>+</sup> increasing the effective free-free emissivity (He has a tiny effect on the amount of ionization of H; e.g., Osterbrock 1989). Comparing models 1 and 4, or models 2 and 5 we estimate that $`f_{He}=1.07`$.
Dust competes with the gas for the absorption of ionizing photons and thus, when present in the nebula, will reduce the observed radio $`F_\nu `$ for a given $`Q(H^0)`$. Comparing models 1 and 2 (with grains) to model 3 (without grains) we derive $`f_{dust}=0.78`$ for the typical ISM grains used in model 1 and $`f_{dust}=0.93`$ for dark-cloud dust, similar to the dust found in Orion, used in model 2. The dark-cloud dust is not as efficient as typical ISM dust in absorbing short-wavelength photons and so has less of an effect on the emergent radio flux.
The combination of these factors, except for the covering factor, leads to a factor of $`f_{model}=0.91\pm 0.08`$ depending upon the grain composition (i.e., $`f=0.91f_{ci}`$).
In our study of the exciting star VES 735 (Kerton et al. 1999) we used the measured radio flux at 1420 MHz ($`2.35\pm 0.05`$ Jy) to calculate $`\mathrm{log}[fQ(H^0)]=48.05\pm 0.11`$, and then estimated $`f_{ci}0.40.5`$ based on the known spectral type of VES 735 (O8.5 V(e); $`\mathrm{log}[Q(H^0)]48.45`$, Panagia 1973). We have not attempted to put any formal error estimate on the covering factor, but given the uncertainties involved, both in the stellar properties and in estimating the various $`f`$ factors, this result is certainly consistent with the idea of KR 140 being an open bowl-shaped region. As demonstrated in § 6.4, an analysis of the IR data leads to the same conclusion independently.
## 5 THE AGE OF KR 140
The approach we have adopted here is first to date the HII region using data we have on the exciting star VES 735. With that age in mind, we then investigate the dynamics of the region, the goal being to show that certain scenarios for the evolution of KR 140, such as it being a blister region, are at least consistent with the age suggested from the exciting star. This approach is similar to that adopted by Dorland et al. (1986) in their study of the Rosette Nebula.
### 5.1 Stellar Content
The idea of using the stellar content of a HII region to measure the age of the nebulosity was first attempted by Hjellming (1968). Basically one plots the evolutionary tracks of stars with various masses in the $`\mathrm{log}(L/L_{})`$ vs. $`\mathrm{log}T_{eff}`$ plane. One obtains $`T_{eff}`$ from the spectral type of the exciting star and $`\mathrm{log}(L/L_{})`$ from the radio flux (much like $`Q(H^0)`$) and then the position in the plane determines an age for the star and thus the HII region. Clearly the effectiveness of this technique depends strongly on the quality of the calibration between theoretical and observational quantities as well as the quality of stellar models, which have vastly improved in the thirty years since this technique was introduced. In the early work the primary result was to indicate whether the HII regions were ionization or density bounded, and that many HII regions required additional, unobserved, sources of ionization.
Here we follow a slightly different technique compared to Hjellming (1968) in order to avoid any uncertainties associated with the covering factor and structure of the nebula. We instead use the absolute magnitude ($`M_V`$) as a measure of $`\mathrm{log}(L/L_{})`$, which is possible because we have a good estimate of the distance to the star and the extinction along the line of sight (Kerton et al. 1999). Figure 4 plots the stellar evolution models of Schaerer & de Koter (1997) for a 20, 25 and 40 $`M_{}`$ star along with the value determined for VES 735. The observed values for VES 735 are consistent with a 25 $`M_{}`$ star with an age of a few million years away from the ZAMS (see Figure 4), where the age would be only of order $`10^5`$ years. Of course, the spectral type alone gives us a simple upper limit to the age of KR 140: for an O8.5 V star the main sequence lifetime is $`6\times 10^6`$ (Chieffi et al. 1998).
### 5.2 Dynamical Models: Spherically-symmetric
#### 5.2.1 Strömgren Sphere
The simplest description of a HII region is that of the formation and expansion of a ionized ball of pure hydrogen at constant temperature in an uniform medium of constant density (Strömgren 1939). We summarize this only as a point of departure and contrast. The evolution of a Strömgren sphere starts with a formation phase where the O star ionizes a region of space around it to the radius ($`R_s`$) given by:
$$R_s=\left(\frac{3Q(H^0)}{4\pi n_e^2\alpha _B}\right)^{1/3}.$$
(3)
The initial rapid expansion to a radius of $`r_i`$ occurs in time $`t=(1/n_e\alpha _B)\mathrm{ln}(1(r_i/R_s)^3`$) (valid to $`r_i/R_s0.98`$, Osterbrock 1989). The time to create the initial Strömgren sphere is about $`10^5/n_e`$ (cm<sup>-3</sup>) y, instantaneous compared to other timescales. For KR 140, the observed radius ($`R_{obs}`$) is 2.85 pc. Taking $`\mathrm{log}[Q(H^0)]=48.45`$ for an O8.5 V(e) star like VES 735 (Panagia 1973) and $`n_{H_2}=100`$ cm<sup>-3</sup> for the original molecular cloud, $`R_s=1.2`$ pc. The initial HII region is overpressured compared to the surroundings and will expand into a uniform medium according to a $`t^{4/7}`$ law:
$$r_i/r_o=\left(1+\frac{7C_{II}t}{4r_o}\right)^{4/7},$$
(4)
where $`C_{II}`$ is the sound speed in the ionized medium, $`r_o`$ is the initial radius and $`r_i`$ is the radius at time $`t`$. With the densities quoted above the HII region will evolve from $`R_s`$ to $`R_{obs}`$ on a timescale of 10<sup>5</sup> years, which is improbably short. Increasing the initial density makes $`R_s`$ smaller and forces the pressure expansion stage to be longer. A value of $`n_{H_2}=500`$ cm<sup>-3</sup> will push the timescale to 10<sup>6</sup> years. However, this value is not consistent with our observations of the molecular material (§ 7), so the Strömgren sphere is not an appropriate dynamical model for KR 140.
#### 5.2.2 Stellar Winds
In the radio image a local minimum is evident near the position of the central star (see Figure 1 in Kerton et al. 1999). One interpretation is that this is a wind blown bubble around the O8.5 V(e) star. The apparent radius of the central hole is 1.4 pc. We used the model of Castor et al. (1975) for the size of a circumstellar shell:
$$R=28\left(\frac{\dot{M}_6V_{2000}^2}{n}\right)^{1/5}t_6^{3/5}\mathrm{pc}$$
(5)
where $`R`$ is the shell’s radius (pc), $`\dot{M}_6`$ is the mass loss rate ($`10^6`$ $`M_{}`$ yr<sup>-1</sup>), $`V_{2000}`$ is the wind velocity (2000 km s<sup>-1</sup>) $`n`$ is the gas density (cm<sup>-3</sup>), and $`t_6`$ is the time ($`10^6`$ years). Stellar wind properties of the B2 model of Schaerer & de Koter (1997) were used. For a wide range of densities we obtain timescales on the order of only 10<sup>4</sup> years. This is far too low a timescale and clearly this model is not an appropriate description of VES 735. Either the mass loss rate adopted is much too high, or the assumption of spherical symmetry is not valid.
### 5.3 Dynamical Models: Open Geometry
Evidence summarized thus far (§ 4.2.2 and § 6.4) indicates an open geometry with a covering factor of about 0.5 for both dust and ionized material. In this case the high pressure in the interior of the initial Strömgren sphere or stellar wind bubble is relieved, slowing down the expansion. Models have been developed to investigate the evolution of an HII region at the edge of a molecular cloud (see review of early models by Yorke (1986), and Comerón (1997) for a recent example). While complex numerical modeling is possible, Franco et al. (1994) have shown that the expansion of the ionization bounded side of such a blister HII region is well described by
$$r_i/r_o=\left(1+\frac{5C_{II}t}{2r_o}\right)^{2/5}.$$
(6)
This is based upon mass conservation between material being ejected in the flow and molecular material being eroded off the cloud. One can envisage the evolution of a blister as consisting of three stages: the initial rapid formation stage, a pressure driven expansion stage, and finally a blow-out stage. The relative length of time of the latter two stages depends upon the distance of the star from the edge of the cloud and the density structure of the cloud. One very important point is that an O star very close ($`1R_s`$) to the edge of a cloud will very quickly develop a covering factor of $`0.5`$ in order 10<sup>5</sup> years and will maintain this covering factor over the lifetime of the O star (Yorke et al. 1983, 1989). We assume that the star formed very close to the edge of the cloud, thus ignoring the pressure driven expansion stage. Using Equation 6, the age of the region is of order a million years when $`R_s=0.64`$ pc. For the stellar properties of VES 735 this requires $`n_{H_2}280`$ cm<sup>-3</sup>. This is somewhat encouraging as it does not require as vast a difference between the properties of the observed molecular cloud and the putative initial conditions.
## 6 KR 140 IN THE INFRARED
### 6.1 Morphology and Emission Mechanisms
Stars that are forming and evolving in the ISM interact with the interstellar dust component by heating, redistributing, and possibly destroying it. The dust can be heated by at least three distinct mechanisms: direct radiation from the central star, reprocessed radiation from the ionized gas, and diffuse radiation from the interstellar radiation field. Radiation pressure from the star acts on the dust in the ionized zone (Spitzer 1978) which causes the dust (and gas) to be pushed away from the star. Some forms of dust like polycyclic aromatic hydrocarbons (PAHs) are destroyed in intense ultraviolet fields. A morphological study of the infrared emission from KR 140 is therefore important to fully understand the energetics and effects on the environment.
As mentioned in § 2, IRAS scans of KR 140 at 12, 25, 60, and 100 $`\mu `$m were processed by the HIRES software to generate maps of about 1 resolution. The new beam shapes are somewhat elliptical and so point sources will be visibly stretched; however, the larger scale morphology of the observed dust emission from the nebula will not be greatly affected by the asymmetric beam.
The intensity of the dust emission, $`I_{dust}`$, at a particular frequency, $`\nu `$, from a distance increment $`ds`$ along the line of sight has the form
$$I_{dust}=N_{dust}\pi a^2B_\nu (T_{dust})Q_\nu (a,T_{dust})ds,$$
(7)
where $`N_{dust}`$ is the number density of grains, $`B_\nu (T_{dust})`$ is the Planck function at a temperature $`T_{dust}`$, $`a`$ is the radius of the particles, and $`Q_\nu (a,T_{dust})`$ is the absorption efficiency factor.
Figure 2 shows the four HIRES images of KR 140, with the 12 and 25 $`\mu `$m images convolved to the 1420 MHz resolution, and overlaid by 1420 MHz contours (Figure 2a & b). The 12 and 25 $`\mu `$m emission is spread well outside the radio contours of KR 140 (this is taken up in § 6.3). Features that are common to all four images in Figure 2 are the bright arcs on either side of the nebula (easiest to see in Figure 2c) that extend outside the radio contours. We interpret this as the limb-brightened warm dust shell around KR 140. Note that there could be cooler dust further out around the KR 140 complex that will not have been detected in the IRAS bands. We account for the energetics of this cooler dust in our models of KR 140 (§ 6.4).
### 6.2 Temperature and Column Density
Equation 7 can be used to calculate a mean dust temperature for each pixel of IR emission. Taking the intensity ratio for any two frequencies, $`\nu _1`$ and $`\nu _2`$, yields:
$$\frac{I_{\nu _2}}{I_{\nu _1}}=\left(\frac{\nu _2}{\nu _1}\right)^{3+\beta }\left(\frac{e^{h\nu _1/kT}1}{e^{h\nu _2/kT}1}\right),$$
(8)
where $`T`$ is the $`T_{dust}`$ of Eq. 7, and the $`\left(\nu _2/\nu _1\right)^{3+\beta }`$ factor is from the frequency-dependent part of $`Q_\nu (a,T)`$, where $`\beta `$ depends on the type of dust. The most common components proposed for interstellar dust are silicate and graphite (or some related carbonaceous material), which, at long wavelengths, have $`\beta 2`$. This value of $`\beta `$ is close to what is generally observed in the ISM (Lagache et al. 1998).
The 60 $`\mu `$m and 100 $`\mu `$m data were used to calculate the dust temperature map, as they are the bands where classical grain emission dominates. It is likely that non-equilibrium heating of very small dust grains (VSGs) contributes some of the observed 60 $`\mu `$m flux and so the derived temperatures are probably slight overestimates of the true grain temperatures (Boulanger et al. 1988). The images were brought to the same resolution and background subtracted and the IPAC analysis program ‘cttm’ was used to compute a dust temperature map. An azmuthially-averaged radial cut of this map is shown in Figure 5. The temperature distribution was sampled every $`1^{}`$ in radius and every $`10^{}`$ azimuthally. The temperatures in the central region are around 31 K. Nearer the edge of the HII region, the dust temperatures drop to around 28 K. However, a line of sight passing near the star also has cooler dust in the background and the foreground, which would lower the apparent temperature observed.
A calculation of dust temperature from first principles was done as a check on the empirical values output from ‘cttm’. In this calculation, a single silicate or graphite dust grain with a radius of 0.1 $`\mu `$m (a typical interstellar size, see Kim, Martin & Hendry (1994)) was placed at a distance of 3.0 pc from the center of the nebula, corresponding to a dust grain within the KR 140 dust shell. The luminosity of the O8.5V(e) exciting star is about 10<sup>5</sup> $`L_{}`$ (Panagia 1973). Making use of Planck-averaged absorption factors from Laor and Draine (1993), we find that if $`\tau _{UV}1`$, the calculated dust temperature is about 28 K, in agreement with the temperature deduced empirically. Note, however, that this calculation did not take into account the recombination-line photons emitted by the ionized gas. Furthermore, some of the free-free photons and collisionally-excited cooling lines are emitted not in the ultraviolet, but in the optical (Osterbrock 1989), where the dust absorption efficiency is somewhat lower.
The optical depth $`\tau _\nu `$ from dust is defined to be
$$\tau _\nu =N_{dust}\pi a^2Q_\nu (a,T_{dust})𝑑s,$$
(9)
and can be calculated by dividing $`I_\nu `$ (Equation 7) by the Planck function $`B_\nu (T)`$, assuming a constant $`T`$ along the line of sight, adopted from the temperature map. Figure 6 displays a dust optical depth map at 100 $`\mu `$m. The values range from about 0.0004 in the middle of the nebula to about 0.002 in the bright northwest rim. These values show that the dust is transparent to its own 100 $`\mu `$m emission. The minimum in the center of the nebula and the ring-shaped appearance is most simply interpreted as limb brightening in a thick shell of dust with the highest column density along the northwest rim. The latter is the same region of the nebula where the 1420 MHz radio contours fall off most steeply (§ 4.1), and where there is no CO emission (§ 7.1).
The optical depth in the ultraviolet $`\tau _{UV}`$ can be gauged using extinction curves from the literature. We are interested in the radial as opposed to line of sight optical depth. Judging the thickness of the shell from one of the arcs in Figure 6 gives a path length of about 7 $`\times 10^{18}`$ cm. If we assume the dust is associated with gas at a molecular density of 100 cm<sup>-3</sup>, the molecular column density of this gas is about 7 $`\times 10^{20}`$ cm<sup>-2</sup>. Using the extinction curves of Kim, Martin & Hendry (1994), we find that $`\tau _\nu 2.8`$ at 1100 Å (far in the ultraviolet), whereas at 5500 Å (in the optical), $`\tau _\nu 0.6`$. The estimated optical depths show that the dust in the arcs has a high enough radial optical depth to absorb most of the incident ultraviolet photons. However, since KR 140 has a covering factor of 0.4 – 0.5 (§ 4.2.2 and § 6.4), the total infrared luminosity of the nebula will be correspondingly less than the bolometric luminosity of VES 735 (§ 6.4).
### 6.3 The 12 and 25 $`\mu `$m Emission
It has been known for over a decade that there is excess emission within the 12 $`\mu `$m IRAS passband (e.g., Boulanger, Baud & van Albada 1985). Onaka et al. (1996) found that more than 70% of the 12 $`\mu `$m diffuse interstellar emission detected by IRAS is emitted in spectral features attributed to polycyclic aromatic hydrocarbons (PAHs; Léger & Puget 1984; Allamamdola, Tielens & Barker 1985). The 12 $`\mu `$m image of KR 140 is very instructive. There is very little 12 $`\mu `$m flux within the radio contours of KR 140, implying that PAHs are destroyed there in the intense ultraviolet radiation field. However, there is a large amount of diffuse flux outside the radio contours, especially to the west (Figure 2*a*). The PAH emission is a tracer of the photodissociation region around a nebula (Giard et al. 1994; Bregman et al. 1995; Fig. 10).
Unique among the HIRES images of KR 140, the 25 $`\mu `$m image (Fig. 2b; Fig. 7) shows a bright spot near the center of the nebula, on a transition between a radio peak (the “left eye”) and the deepest depression. The IRAS Point Source Catalog (Joint IRAS Science Working Group 1988) lists this feature as IRAS 02165+6053, and it also has been identified with VES 735 (Bidelman 1988). Figure 7 shows that IRAS 02165+6053 and VES 735 are practically coincident. Dust closer to the star would tend to be warmer, contributing to the spot if not pushed away by radiation pressure.
The entire nebula looks hotter and more extended at 25 $`\mu `$m than it would be for equilibrium emission from normal-sized grains. Most of the 25 $`\mu `$m emission is probably contaminated by non-equilibrium emission from very small grains (VSGs; Sellgren 1984) which have absorbed a UV photon and have had their temperatures instantaneously rise to $`10^3`$ K. Although these small grains make up a tiny fraction of the mass in the dust distribution, they make up a good fraction of the number distribution and absorb a significant fraction of the near-ultraviolet radiation. This VSG emission is therefore a major source of uncertainty when analyzing the 25 $`\mu `$m image.<sup>1</sup><sup>1</sup>1Further evidence for contamination in the 25 $`\mu `$m IRAS passband was presented by Cox (1990). He suggests that an iron oxide emission line falls within this IRAS passband which would result in even greater excess 25 $`\mu `$m emission. Note that the “temperature-spiking” phenomenon also occurs outside the ionized zone, far from the exciting star.
### 6.4 Infrared Models
HII regions are some of the most luminous objects in the Galaxy when observed in the infrared, especially at wavelengths longer than 60 $`\mu `$m. In fact, if the dust shell covers 4$`\pi `$ steradians around the exciting star, the infrared luminosity should be a good measure of the star’s bolometric luminosity. To account for the true extent of the dust, we can define a covering factor of dust ($`f_{cd}`$). Often it is simply assumed that $`L_{ir}=L_{bol}`$; however, this is not correct even for $`f_{cd}=1`$. Some of the radiation from the star and nebula is at long enough wavelengths to avoid being absorbed by the dust. Using Cloudy we found that for models with $`f_{cd}=1`$, $`\mathrm{log}\left(L_{ir}/L_{bol}\right)0.1`$.
To calculate the integrated infrared flux from KR 140, the background in each of the four HIRES images was fitted and subtracted, and the total flux from the nebula was measured in each band. In order to estimate $`L_{bol}`$ using these data, we used Cloudy to simulate the emission from large classical grains, roughly matching the fluxes at 60 and 100 $`\mu `$m. This approach allows us to account for emission from grains with a range of temperatures and thus unobserved emission at long wavelengths. First, we constructed models with $`f_{cd}=1`$. We find $`\mathrm{log}(L_{FIR})=37.90`$. Since the resulting spectrum misses most of the observed 12 and 25 $`\mu `$m flux (which is caused by temperature spiking of VSG’s and the excitation of PAH molecules) we converted the observed 12 and 25 $`\mu `$m fluxes to luminosities using tophat approximations to the IRAS passbands (Emerson 1988) and added the results to $`L_{FIR}`$. With this addition we find $`\mathrm{log}(L_{IR})=38.00`$. Correcting from $`L_{ir}`$ to $`L_{bol}`$ we find $`\mathrm{log}\left(L_{bol}/L_{}\right)=4.52`$. This is significantly below what would be expected for any late O main sequence star (e.g., Panagia 1973). We interpret this low apparent $`L_{bol}`$ as being due to a covering factor of 0.4 – 0.5.
We know that VES 735 is an O8.5 V(e) star. Figure 8 demonstrates that one can reproduce the observed $`L_{ir}`$ using the appropriate stellar parameters and a covering factor of about 0.5. A single temperature (T = 28.25 K) $`\nu ^2B_\nu `$ spectrum is shown for comparison; note the wider model curve caused by the range of dust temperatures contributing to the emission.
Using an alternative calibration of stellar parameters which includes wind-blanketed models (Kerton 1999), we find $`f_{cd}0.5`$ and $`f_{ci}0.7`$. These results are still consistent with a blister model for KR 140.
## 7 KR 140 IN A MOLECULAR CLOUD
### 7.1 CO Signature
A slit spectrum of VES 735 allowed us to measure the radial velocity of the nebular H$`\alpha `$ line within 30<sup>′′</sup> of VES 735 to be $`46\pm 2.1`$ km s<sup>-1</sup>with respect to the Local Standard of Rest \[LSR\] (Kerton et al. 1999). The differential radial velocity between VES 735 and the nebular line was measured to be $`+2.0\pm 2.2`$ km s<sup>-1</sup>.
An expected effect of the evolution of the HII region is photodissociation of molecular gas both inside the HII region and in the immediate surrounding interstellar medium (ISM). The CO data traces the molecular gas content of the ISM, so an examination of the CO data cube ought to reveal a lack of emission near the ionized gas velocity of KR 140. Indeed, a distinct CO hole was found within the radio contours of KR 140 over the velocity range $``$45.53 km s<sup>-1</sup>to $``$47.16 km s<sup>-1</sup>(LSR), which corresponds to three channels of width 0.813 km s<sup>-1</sup> in the CO cube. At more negative velocities, CO emission from the parent molecular cloud fills in the 1420 MHz contours.
Figure 9 shows the sum of these three channels overlayed with the 1420 MHz continuum contours. A well defined ring structure is clearly seen. Interestingly, the ring does not extend all the way around the nebula: there is no CO emission in the north-west. This is the same area where the 1420 MHz contours are falling off more sharply (§ 4.1), there is a bright infrared arc (§ 6.1), and a bright HI feature (§ 8). The relationship between 12 $`\mu `$m emission, which is a good tracer of the PDR, and the CO emission is shown in Figure 10. The north-west peak in the 12 $`\mu `$m emission corresponds to the region where there is no CO emission. To the east, as one moves away from the HII region the 12 $`\mu `$m peak occurs first followed by the peak in the CO emission.
### 7.2 Density and Mass of the Molecular Cloud
To estimate the mass of the molecular cloud we integrated the CO cube over the whole velocity range of the cloud ($``$45.5 km s<sup>-1</sup>to $``$52.8 km s<sup>-1</sup>), and measured the total surface brightness of the molecular cloud. We can use the empirical $`X`$ factor to convert from CO surface brightness to molecular hydrogen column density. We have used $`X=(1.9\pm 0.5)\times 10^{20}`$ cm<sup>-2</sup> (K km s<sup>-1</sup>)<sup>-1</sup> (Strong & Mattox 1996), where the quoted uncertainty is to take into account the range of $`X`$ values measured for a variety of different clouds, to obtain a mean column density of N(H<sub>2</sub>)$`=2.2\times 10^{21}`$ cm<sup>-2</sup>. Assuming a path length of 7 pc (the north-south radial extent) the density of the molecular cloud is roughly $`n_{H_2}=100`$ cm<sup>-3</sup>, typical of a giant molecular cloud (Blitz 1993).
By integrating the column density over the face of the cloud, we estimate the mass of the cloud to be $`(4.4\pm 1.6)\times 10^3M_{}`$. This mass includes a He correction factor of 1.36, and the error bar includes the distance uncertainty. The CO images show that the parent molecular cloud of KR 140 has been greatly disrupted by the nebula and the exciting star VES 735 and so this mass will be an underestimate of the cloud’s initial mass.
#### 7.2.1 Estimating the Mass of the Original Molecular Cloud
One way to estimate the mass of the original molecular cloud would be to assume that KR 140 is indeed a blister HII region as our data suggests. Yorke et al. (1989) have modeled blister HII regions with one exciting O star and found that mass loss rate of material through the blister is $`35\times 10^3M_{}`$ yr<sup>-1</sup>. However, these authors modeled an O6 star in a cloud with $`n_{H_2}=500`$ cm<sup>-3</sup>, which is not an accurate description of either VES 735 or its parental cloud. Unfortunately, Yorke et al. (1989) give no indication of how to scale their mass loss rate for different values of $`Q(H^0)`$ or $`n_{H_2}`$. A analytic estimate of the amount of mass lost from a blister HII region is given by Whitworth (1979). His Equation (41) has the desirable property that it agrees with the results of Yorke et al. (1989) for the values used in their models; therefore, it might have the correct scaling for $`Q(H^0)`$ and $`n_{H_2}`$. However, his calculations were based on a cylindrical geometry which is not a very realistic model for a more bowl-shaped spherically symmetric region such as KR 140.
In § 5.3, we made use of the blister evolution formula given by Franco et al. (1994). These authors considered a simple symmetric blister region and were also able to estimate the cloud evaporation rate:
$$\dot{M}\pi R_S^2m_p2n_{H_2}C_{II}\left(1+\frac{5C_{II}t}{2R_S}\right)^{1/5},$$
(10)
where $`R_S`$ is the Strömgren radius and $`m_p`$ is the proton mass. In § 5.3, we found that an initial cloud density of 280 cm<sup>-3</sup> and $`R_S=0.64`$ pc would match the observed radius of KR 140 in a time of 1 Myr (consistent with the age of VES 735). Using these values in the above equation gives a mass loss rate of $`2.9\times 10^4M_{}`$ yr<sup>-1</sup>. Therefore, about 290 $`M_{}`$ of molecular material has been eroded from the cloud after 1 Myr. Including the effects of dust lowers $`R_S`$, increases the age, and lowers $`\dot{M}`$, resulting in a very similar amount of erosion (280 $`M_{}`$).
Taking into account the mass lost through erosion and the present ionized mass, we estimate the mass of the original molecular cloud in which VES 735 formed to be about $`4.9\times 10^3M_{}`$, so this cloud would be classified as a dwarf molecular cloud (Elmegreen 1985). This seems to be a remarkably small cloud mass from which an O star has formed. Williams & McKee (1997) performed a statistical analysis of many nearby OB associations and molecular clouds, and found that it was more likely that clouds with masses greater than 10$`{}_{}{}^{5}M_{}^{}`$ should form O stars than clouds with lower masses. Other HII regions have been found with small molecular cloud masses (Hunter et al. 1990), but O star formation with a cloud mass below 10<sup>4</sup> is considered rare (Elmegreen 1985). Of course we have been attracted to this cloud because of the effect of VES 735 and not to other parts of the W3 cloud complex (Heyer et al. 1998) which have not formed O stars.
### 7.3 Interpretation of the Velocity
Stepping through the velocity channels in the CO data cube near the velocity of KR 140 shows that at higher (least negative) velocities ($`40`$ km s<sup>-1</sup>) there is no emission, then the CO emission comes in at the southern end of the HII region and spreads northward, before filling in the 1420 MHz contours at a velocity of $``$48.78 km s<sup>-1</sup>. Perhaps the best way to illustrate these data is by examining the cube in velocity-latitude space. Figure 11 shows such a CO image of KR 140, averaged over the longitude range 133.366 to 133.477. Note the “hole” at the velocity of the ionized gas and VES 735.
The question arises as to the relative radial position of VES 735 and KR 140 with respect to the molecular gas. In general, it is difficult to use only the CO radial velocity information to determine absolute distances to clouds or GMCs. However, in this direction, Galactic rotation causes radial velocity to become more negative with increasing distance. On the face of it, that would place the molecular cloud on the far side of the star and nebula. However, this is possibly too naive an interpretation because if the cloud is only 10s of pc in radial extent (like its dimension in the plane of the sky), then the average shear would be too small to explain the large range in velocity (unless there were a large enhancement from a density wave).
To address the question of relative position, we made use of both the HI and CO data cubes to make an estimate of $`A_V`$ produced by gas with velocities out to $`45`$ km s<sup>-1</sup>. As in § 7.2 we used the conversion factor $`X=(1.9\pm 0.5)\times 10^{20}`$ cm<sup>-2</sup> (K km s<sup>-1</sup>)<sup>-1</sup> to convert I(CO) to N(H<sub>2</sub>). We integrated both data cubes to obtain the atomic column densities. At the position of VES 735 we found N(HI) = $`4.5\times 10^{21}`$ cm<sup>-2</sup> assuming the emission is optically thin, and N(H<sub>2</sub>) $`=3.1\times 10^{21}`$ cm<sup>-2</sup>, the latter mostly from local rather than Perseus arm gas. Summing the contributions of atomic and molecular hydrogen we obtain the total hydrogen column density, N<sub>H</sub> = 2N(H<sub>2</sub>) + N(HI)$`=1.1\times 10^{22}`$ cm<sup>-2</sup>. The total visual extinction is computed using the standard conversion factor, A$`{}_{V}{}^{}=5.3\times 10^{22}`$ mag cm<sup>-2</sup> (Bohlin, Savage & Drake 1978). At the position of VES 735 we obtain A$`{}_{V}{}^{}=5.7\pm 0.9`$, where most of the uncertainty comes from the uncertainty in the $`X`$ factor (using the $`X`$ factors of Digel et al. 1996 gives a lower $`A_V`$). This value of $`A_V`$ compares favorably with the values around $`5.7`$ derived by Kerton et al. (1999) using a number of methods (e.g., $`BV`$ colour, H$`\alpha `$ emission measure, DIBs). However, if the cloud’s molecular column density (summed over $`45`$ to $`53`$ km s<sup>-1</sup>) of N(H<sub>2</sub>) $`=2.2\times 10^{21}`$ cm<sup>-2</sup> (§ 7.2) is included in the extinction calculation, then the $`A_V`$ rises to 8.0, which is quite inconsistent with the observed $`A_V`$ to VES 735. We also made maps of the predicted extinction over the surface of the nebula for comparison with the extinction map derived by comparing H$`\alpha `$ surface brightness with radio emission (Kerton et al. 1999). Again, inclusion of the extinction from gas in the molecular cloud produces more than two magnitudes too much extinction. This is strong evidence that both VES 735 and the ionized gas of KR 140 lie on the near side of the molecular cloud gas which has velocities $`45`$ to $`53`$ km s<sup>-1</sup>.
The extinction measurements combined with Figure 11 suggest that KR 140 could be a blister HII region on the near side of the molecular cloud. This simple geometrical interpretation of KR 140 runs into difficulty if a systematic champagne flow has developed with gas flowing away from the parent cloud at up to the sound speed of $`10`$ km s<sup>-1</sup>. At face value, our H$`\alpha `$ velocity implies that the ionized gas is redshifted with respect to the molecular gas, which means the HII region should be on the far side of the molecular cloud. However, it can be noted that H$`\alpha `$ is considered to be a poor line for estimating the velocity field of champagne flows (Israel 1978; Yorke et al. 1984). Modeling of line profiles in champagne flows also shows that the velocity field can be quite random and of low amplitude, depending on geometry (Yorke et al. 1984). Recall that the geometry here is certainly not plane-parallel and probably more like a broken shell. Furthermore, our measurement sampled only gas within a projected distance of 30<sup>′′</sup> of VES 735. Fabry-Perot data of the whole nebula in a line other than a hydrogen line would be a good way to obtain a better picture of the velocity structure of the ionized gas in KR 140.
The question then arises as to the origin of the radial velocity spread within the molecular cloud. If this is a single, gravitationally bound cloud, then this is just the virial velocity, and its mass can be estimated from
$$M_{vir}=\frac{5R\sigma ^2}{\alpha _{vir}G},$$
(11)
where $`R`$ is the length scale for the cloud, $`\sigma `$ is the linewidth, and $`\alpha _{vir}`$ is the virial parameter and includes the effects due to surface pressure, magnetic fields, and nonuniform densities (Bertoldi & McKee 1992). As in § 7.2, we take $`R=7`$ pc, and measure a typical linewidth to be 5 km s<sup>-1</sup>(see Fig. 11). We take $`\alpha _{vir}`$ to be 1.1 (Williams & McKee 1997). With these values $`M_{vir}=1.9\times 10^5M_{}`$ which is well over an order of magnitude larger than our measured CO mass of $`4.4\times 10^3M_{}`$.
This result implies that the velocity field has another origin and can significantly alter the cloud structure over time. In fact, recalling that a velocity of 1 km s<sup>-1</sup>corresponds to 1 pc in a million years, the “crossing time”, $`2R/\sigma `$, is only 3 Myr, comparable to our above estimates of the age of VES 735 and the HII region. From the radial arrangement derived from extinction, it seems that the background molecular material is moving toward KR 140. Perhaps the parent cloud giving rise to VES 735 was the result of a converging flow or “cloud collision” that is still ongoing (relative motion in the plane of the sky would be expected too).
## 8 HI SIGNATURE
In theory the presence of HII regions and their surrounding molecular material should be easily “visible” in HI data sets: in both the ionized and molecular regions there should be a deficit of HI emission due to a lack of neutral atomic hydrogen. In addition, if the geometry is favorable, neutral atomic hydrogen associated with the photodissociation region (PDR) at the ionized-molecular interface should be visible. In practice it is actually very difficult to observe HI features that are unambiguously associated with an HII region in raw HI data sets. This is primarily due to velocity confusion along the line of sight caused by the “turbulent” motion of the atomic gas. Since the velocity of the neutral atomic gas can be many times the typical channel width in the HI data sets velocity becomes a much poorer proxy for physical distance than in a CO data set, which is probing a species with lower velocity dispersion (heavier, cooler, and less turbulent). The situation is especially problematic for HII regions in the Galactic plane where one has to look through a large column of atomic gas towards the HII region. For KR 140 we are looking through the local arm and part way into the Perseus arm of the Galaxy. A preliminary reconnaissance of the HI cube confirmed that the complex HI emission structure along the line of sight makes seeing any HI signal associated directly with KR 140 extremely difficult. Nevertheless, some simple processing of the HI data cube does bring out some features associated with the region.
In order to exclude local HI emission (which is assumed to have a relatively smooth spatial structure) and to enhance the dynamic range of the resulting channel maps we constructed a median-subtracted data cube (Joncas et al. 1992; Joncas et al. 1985). In this technique a median spectrum is calculated for the data cube and then subtracted from each spectrum making up the cube. The resulting channel maps thus can contain both negative and positive values indicating deviations relative to the median base level.
Figure 12 shows channel maps of the median-subtracted cube over the velocity range $`43.40`$ to $`54.95`$ km<sup>-1</sup> bracketing the ionized gas velocity of $`46\pm 2.0`$ km s<sup>-1</sup> (Kerton et al. 1999). We could not detect any features definitely associated with KR 140 in the channel maps outside of this velocity range.
Examining these maps we note the following three features. First, there is a noticeable deficit of HI seen in velocity channels $`46.70`$ to $`51.65`$ km s<sup>-1</sup> outside of the HII region. There is excellent positional agreement between this deficit in HI and the observed position of the CO emission. The deficit can be simply interpreted as being caused by a lack of HI emission in the molecular material surrounding KR 140. Second, we also see a drop in the HI emission occurring within the HII region in the $`43.40`$ and $`45.05`$ km s<sup>-1</sup> channels. This deficit is most likely associated with the ionized gas in KR 140 as suggested by their spatial correspondence. Figure 13 presents spatially averaged HI and CO spectra for an area just outside of the HII region to the north-east and the area inside the HII region. The anticorrelation between HI emission and the presence of CO and HII is evident. Third, there is an enhancement of HI emission at (133.36,0.1583), seen best in the velocity channel $`46.70`$ km s<sup>-1</sup>. This could be low velocity dispersion material associated with the PDR. Atomic material in the PDR is expected to have a low velocity and thus should be seen in channels corresponding to CO emission.
## 9 THE ENVIRONMENT OF KR 140
### 9.1 Spontaneous Massive Star Formation
In the context of the Perseus Arm star formation activity, the KR 140 complex seems to be unique. Figure 1 shows that KR 140 is isolated from the massive and violent star formation that is ongoing around it. This isolation is evidence to us that KR 140 is an example of massive star formation in our Galaxy that is *untriggered*, at least in the sense used in the context of sequential star formation. In none of our data sets does there appear any evidence for a trigger of the star formation in KR 140. The exciting cluster of W4, OCl 352, is about 60 pc away from KR 140 for a cluster distance of 2.35 kpc (Massey, Johnson & DeGioia-Eastwood 1995). For a sound speed of 0.6 km s<sup>-1</sup>(isothermal speed in H<sub>2</sub> at 100 K), the time for a signal to reach KR 140 would be about 90 Myr, much greater than our estimated age of KR 140 of a few million years or the age of OCl 352. From W3, a signal would take about 60 Myr, but W3 is itself much less than 10<sup>5</sup> years old (Kawamura & Masson 1998). The supernova remnant HB3 is also in this complex (Normandeau et al. 1997), but from Figure 1, the edge of the remnant is nowhere near to KR 140. Thus, KR 140 does not seem to be triggered by a neighboring HII region or a nearby supernova remnant, unless the impulse came along the line of sight. This conclusion is complemented by the overall smoothness of the KR 140 nebula; while the nebula does show density inhomogeneities, it is not far removed from a circular shape. Thus, any perturbation that might have triggered the star formation within KR 140 must have been a large scale phenomenon with a characteristic scale of $`10`$ pc. This kind of triggering might be more consistent with triggering via a spiral density wave (e.g., Elmegreen 1994, 1995) or by colliding molecular clouds. Of course, we cannot rule out those kinds of triggers, but the observational evidence would be difficult to find.
It is interesting to contrast KR 140 to another star forming region that has been studied with multiwavelength data, the Gemini OB1 molecular complex (Carpenter, Snell & Schloerb 1995a, 1995b). Within the molecular complex, these authors find young star clusters (from near infrared data) and a number of dense cores (as identified by CS observations) associated with IRAS point sources. Carpenter et al. suggest that the arc-shaped morphologies of these cores have been formed by swept-up gas from expanding HII regions, and that they would form the massive star clusters in the region. The other lower infrared luminosity sources (most likely belonging to lower mass cores) in the Gem OB1 complex are not found to be correlated with any arc-shaped structures or filaments in the molecular gas, and they are not adjacent to any HII regions; in fact they are spread almost randomly around the complex. Carpenter et al. therefore conclude that induced star formation is the prominent mode of formation for massive stars in the Gem OB1 molecular complex.
If KR 140 is indeed untriggered, then it seems to be an unusual form of spontaneous star formation since current observations suggest that the isolated mode of star formation is generally associated with low mass star forming regions as seen in Gem OB1 or even in Taurus.
### 9.2 IRAS Point Sources and Protostars
As has been seen in other sites in the Galaxy, an HII region can initiate star formation via the expansion of its ionization front (e.g., Elmegreen & Lada 1977). Other stars might also have formed spontaneously. Evidence for other star forming regions within KR 140 may be sought by examining the IRAS point source catalog. In addition to the IRAS point source 02165+6053 that is cross-referenced to VES 735, there are six other IRAS point sources in the area in and around KR 140 (see Table 3). The crosses in Figure 14 show the positions of these six IRAS point sources. The circled crosses are the sources discussed here that seem most likely to be associated with the KR 140 complex.
The source IRAS 02160+6057 identified with the north-west dust arc has been the subject of two molecular line investigations. Wouterloot & Brand (1989) identified it as a potential star-forming area via its IRAS colors (see their paper for the exact selection criteria), and examined it (along with about 1300 other IRAS point sources) for CO emission. They found a CO feature in that direction at a velocity of $``$49.7 km s<sup>-1</sup>(LSR), which corresponds to CO in the associated background molecular cloud, and assigned it the catalog number \[WB89\]417. This point source was then observed in an H<sub>2</sub>O line by Wouterloot, Brand & Fiegle (1993), but they were unable to detect any emission. This line of sight has one one of the highest column densities in the KR 140 HII region, and so it is possible that a protostar could be forming there as a result of the expansion of the HII region. However, examination of the HIRES images shows no point-like features in the dust arc, and none are found in follow-up submillimetre observations with SCUBA (Kerton et al. 2000). Therefore, IRAS 02160+6057 is most likely simply part of the KR 140 dust shell. Likewise, IRAS 02168+6052 appears to just be part of the eastern side of the dust shell; it has not been the subject of any molecular line observations.
The source IRAS 02157+6053 also seems to have been an identification of some part of the dust shell. In the submillimetre there is more structure in this region; this object appears to be a molecular core rather than a protostar (Kerton et al. 2000).
The point source IRAS 02171+6058 is identified with the dust feature to the north of the KR 140 complex. It was included in a CS(2–1) survey by Bronfman, Nyman & May (1996) of IRAS point sources that have colors characteristic of ultracompact HII regions; however, they were unable to make a detection. Lyder & Galt (1997) observed this source along with other ultracompact HII region candidates in a search for methanol maser emission. Again, they were unable to detect any maser emission from IRAS 02171+6058. These non-detections do not rule out the possibility that the source is a protostar. In fact there is a submillimetre detection and all evidence seems consistent with a B5V star (Kerton et al. 2000). It is outside the radio contours of HII region and is therefore difficult to interpret as being triggered.
A photographic survey of the W3 and W4 region turned up Bright InfraRed Stars (BIRS; Elmegreen 1980) which are brighter in I band than in R band. There are five of these stars in the KR 140 region of the sky. Table 4 gives their coordinates along with their R and I magnitudes and the positions are indicated by triangles on Figure 14. Elmegreen (1980) estimated that if these stars were at the distance of the Persues arm then they might be deeply embedded massive early or pre-main sequence stars, or even giants or supergiants. From examining the overall distribution of the BIRS stars (Figure 1 in Elmegreen 1980), we find there is an overabundance of BIRS around KR 140 as compared to the average. It is difficult to say whether or not these BIRS stars are physically associated with the KR 140 complex. However, BIRS 128 on the eastern side of the nebula seems to lie at a well defined edge of the dust shell, and within the CO shell.
### 9.3 The Gas Cloud
We have been able to estimate some properties of the original molecular cloud which spawned the KR 140 HII region. Originally having $`M4.9\times 10^3M_{}`$ this cloud would then be classified as a dwarf molecular cloud. These are quite numerous in the ISM. The column density of $`2.2\times 10^{21}`$ cm<sup>-2</sup> is about the same as the median column density of $`2\times 10^{21}`$ cm<sup>-2</sup> in the molecular gas in the outer Galaxy (Heyer et al. 1998). Once VES 735 formed in the molecular cloud, it had a tremendous impact on the environment around it. It has incorporated 25 $`M_{}`$, ionized about 160 $`M_{}`$, and a further $`290M_{}`$ of cloud material has been removed since the ionization front broke out of the front part of the molecular cloud. The remnant cloud is of fairly low density, $`n_{H_2}100`$ cm<sup>-3</sup>.
The velocity width of the molecular cloud is too large for the cloud to be bound, but it is unlikely to arise from velocity shear due to differential Galactic rotation. It is possible to interpret the velocity structure as a large scale flow of material towards KR 140. If true, this flow could have had an impact on the star formation history of this region.
We can use our estimate of the original molecular cloud mass to estimate the mass of the stellar cluster that likely formed along with VES 735. Hunter et al. (1990) observed clouds in our mass range and estimated cluster masses using a Miller-Scalo initial mass function (IMF; Miller & Scalo 1979). Using their data we find the correlation shown in Figure 15: $`\mathrm{log}M_{cluster}=(0.254\pm 0.111)\mathrm{log}M_{cld}+(2.306\pm 0.449)`$. With the above cloud mass we estimate a cluster mass of about $`1.7\times 10^3M_{}`$. Of course, since we know the mass of the most massive star we can use a Miller-Scalo IMF to estimate the mass of the cluster directly. Using Equation 12 in Elmegreen (1983) we estimate a cluster mass of about $`1.6\times 10^3M_{}`$, consistent with the above “empirical” estimate. Using these cluster and cloud masses we estimate a star formation efficiency (SFE = $`M_{cluster}/(M_{cluster}+M_{cld})`$) of about 25% which is typical for clouds of this size (Hunter et al. 1990), although values can range over two orders of magnitude from cloud to cloud (Williams & McKee 1997). This low value of the SFE implies that the cluster within KR 140 will not be gravitationally bound.
Near infrared observations (such as ones in the ongoing 2MASS survey) should be able to detect a number of pre-main sequence stars from this cluster around VES 735. It would be interesting to compare the IMF in this region to the IMF in regions where the star formation was triggered. Which, if either, is a good description of the IMF in the Galaxy?
## 10 CONCLUSIONS
We have utilized our multiwavelength data set (all at a resolution of about 1) to study not only the physics of the HII region itself, but also, since the data are from the larger CGPS survey, to study KR 140 in the context of the overall picture of star formation in the Perseus spiral arm. We find no evidence for a mechanism that triggered the formation of the O8.5V(e) star, VES 735, and its (largely unseen) cluster. We therefore conclude that this region formed spontaneously out of its parent molecular cloud, independent of the more vigorous star formation in W3 and W4 nearby.
Our data of KR 140 are consistent with the model of a bowl-shaped region viewed close to face on. Extinction measurements to the exciting star, VES 735, and nebula show that the HII region is quite likely on the near side of its molecular cloud. We have not observed any champagne flow, and cannot rule out other geometries. KR 140 has an age less than a few million years. We have estimated that the original molecular cloud had a mass of $`4.9\times 10^3M_{}`$ and an average density about 100 cm<sup>-3</sup>, which classifies it as a dwarf molecular cloud. This makes KR 140 even more unusual as it is a rare example of an O star that has formed in a cloud with a mass less than 10$`{}_{}{}^{4}M_{}^{}`$. There is tentative evidence that the molecular material is undergoing a large scale flow towards KR 140. Follow-up observations are needed to pursue this idea. There are four IRAS point sources associated with the KR 140 complex, one of which is a possible protostar candidate and another a molecular core. Near infrared observations of KR 140 are needed to find and study the young cluster that likely formed along with VES 735.
We would like to thank Gary Ferland for his assistance with the use of Cloudy and Doug Johnstone for helpful discussions. We also acknowledge useful comments by the referee, Peter Barnes. This research made use of the SIMBAD data base, operated at CDS, Strasbourg, France. This research was supported by the Natural Sciences and Engineering Research Council of Canada. D.R.B. participated originally via the Physics Co-op program of the University of Victoria.
|
warning/0003/cond-mat0003167.html
|
ar5iv
|
text
|
# Multiparticle Biased DLA with surface diffusion: a comprehensive model of electrodeposition
## I Introduction
Quasi-two-dimensional (quasi-2D) electrochemical deposition (ECD) has become one of the most widely studied pattern forming processes since its recognition as a paradigm of non-local, non-equilibrium growth processes . Within this general context, a great deal of work has been devoted in the past fifteen years to experimental and theoretical studies of quasi-2D ECD. A first group of works deals mainly with pattern formation, its main results concerning “phase diagrams” of morphologies , ECD as Laplacian growth process , dynamic morphological transitions , etc. All these studies aim to understanding the principles underlying the rich variety of morphologies observed, ranging from dendritic to fractal. In addition to this line of research, there is a second one whose main interest is the existence of universality and scale invariance in the roughness of the deposits produced . From all this and related research, it is now believed that complex structures with different morphologies arise from quasi-2D ECD due to the interplay of different transport mechanisms, such as cation diffusion, electromigration, fluid convection, and surface diffusion (SD) . However, the combined effect of all these factors leads to a very complex process, and it is becoming increasingly apparent that ECD is not well understood yet. In particular, the detailed role of surface diffusion (SD) is still an open question that hinders our understanding of both the morphologies and the scaling of ECD aggregates.
Much of the work mentioned in the above paragraph has been motivated by the quest to find a universal model to help understand ECD phenomena. The first model formulated with that purpose was the famous computer algorithm known as Diffusion Limited Aggregation (DLA) , in which a particle diffuses on a lattice and attaches to the growing aggregate at the place where it first hits. It is not difficult to observe (see and references therein) that this simple model represents the zero concentration, quasi-static limit of ECD. Therefore, its validity as a general description of ECD is rather restricted because it does not include most of the effects involved in the process. However, DLA has played a seminal rôle as a source of inspiration both for continuum approaches —which predict some high-current properties but take into account neither the influence of the applied voltage nor the electrolyte concentration— and for more sophisticated computer models, basically modifications of DLA (see, e.g., and also the paragraph below), which are more or less phenomenological and concentrate on changes in morphology, thus being unable to explain the underlying mechanisms yielding those patterns.
In this paper, we report on the results of detailed numerical studies of Multiparticle Biased Diffusion Limited Aggregation (MBDLA) supplemented with SD. MBDLA is a model in the family of multiparticle DLA models , in which a finite number of random walkers, possibly with constant concentration, is introduced instead of the single walker of DLA. Thus, the excluded volume interaction among the walkers leads to several of the effects neglected in DLA. As its main ingredient, MBDLA includes, in addition, a preferential bias (which had been first studied in the context of single-particle, DLA models by Meakin ) of the walkers towards the cathode to mimic the electric field: In this form, the model was successfully introduced in to study the influence of the applied electric field on the composition of magnetic, amorphous CoP alloys grown by ECD at constant current. The main virtue of MBDLA is that it is a mesoscopic model embedded on a two-dimensional square lattice, but it reproduces the mean fraction of Co and P in two-cation species ECD, as well as the qualitative morphology of the product electrodeposits. In fact, the agreement between MBDLA and ECD experiments is quantitative, as the electrical current intensity and the experiment time can be directly related to the simulation parameters . Therefore, we are confident that MBDLA is a good starting point to study the relevance of SD in ECD and, specifically, its influence on the shape of the aggregates and their dynamic scaling. Scaling properties of MBDLA without SD were briefly reported in .
The report of our results is organized as follows: We describe our model in Sec. II, where a brief introduction to the physics and chemistry of ECD is followed by a detailed account of the rules governing MBDLA. Section III reports on our numerical results, such as morphological patterns and roughness scaling. After physically showing that SD has to be included, we introduce three different rules for SD are carefully considered and compared to experiments, allowing to identify the proper way to introduce SD in the model. Finally, we conclude in Sec. IV with a discussion of our results which will allow us to suggest a reasonably approximate picture of ECD phenomena. A few technical details about one the rules for SD are given in an Appendix.
## II The Model
### A Basic facts about ECD
Prior to describing in detail what MBDLA is, and in order to motivate and to better understand the model rules, we will briefly summarize the basic physics and chemistry of ECD, by collecting the equations commonly accepted to govern its main features (see, e.g., for further details). Generally speaking, ECD experiments involve two species, named cations and anions, moving in an incompressible viscous fluid. In very many cases, ECD takes place in quasi-2D cells with parallel electrodes. The cations move towards the cathode and the anions towards the anode. The basic equations for the concentrations of both species are as follows:
$`{\displaystyle \frac{C}{t}}`$ $`=`$ $`J_c,`$ (2)
$`{\displaystyle \frac{A}{t}}`$ $`=`$ $`J_a,`$ (3)
$`𝐉_c`$ $`=`$ $`D_cC+\mu _c𝐄C+𝐯C,`$ (4)
$`𝐉_a`$ $`=`$ $`D_aA\mu _a𝐄A+𝐯A,`$ (5)
where $`C`$ and $`A`$ are the cation and anion concentration respectively, $`D_{c,a}`$ the cationic and anionic diffusion coefficients, $`\mu _{c,a}`$ their mobilities, $`𝐯`$ the fluid velocity field, and $`𝐄`$ the electric field along the cell. The latter is related to cation and anion concentration via the Poisson equation
$$E=^2\varphi =e(z_cCz_aA)/\epsilon ,$$
(6)
where $`\varphi `$ is the applied potencial, $`ez_c`$ and $`ez_a`$ are the cation and anion electric charges, respectively, and $`\epsilon `$ the dielectric permittivity of the fluid. Generally speaking, matter balance across the interface leads to an interface velocity proportional to the flux of cations and therefore, in the absence of any other limiting process, proportional to the current density as well. In addition, except for the region close to the cathode, we may assume electro-neutrality , which in turn implies that the cation mean velocity is constant. The incompressible Navier-Stokes equations determine the velocity $`𝐯`$ of the solvent. Fluid convection is always present in ECD experiments, but in many instances it can be small enough to be safely neglected, as has been shown in actual experiments .
When the cations arrive at the cathode, they reduce irreversibly and an aggregate of neutral particles begins to grow. The particles on the surface aggregate are transported all along it due to local chemical potential gradients. The resulting particle current conserves the number of particles on the surface and is given approximately by by (see for a detailed discussion):
$$𝐉_s_s\kappa ,$$
(7)
where $`𝐉_s`$ is the particle current along the surface, $`\kappa `$ the interface local curvature at each site, and $`_s`$ the gradient taken along the surface. Roughly speaking, SD tends to reduce the interface local curvature. Finally, we note that the mean concentration of charge carriers in the bath is constant as new cations are formed at the anode upon arrival of the anions .
### B Definition and rules of MBDLA
In this section we will define MBDLA through its evolution rules, for which we take into account the physical equations presented in the previous section. At this point, we do not consider SD, whose need will be justified in the next Section, and consequently we postpone the discussion of the rules to implement SD as well. Thus, MBDLA is a cellular automaton defined on a two-dimensional square lattice of horizontal dimension $`L_x`$ and vertical dimension $`L_y`$ (with lateral periodic boundary conditions and reflective boundary condition at the top; for the conditions at the bottom, see below), in which a number of random walkers (cations) are randomly distributed with concentration $`c`$. The bottom of the lattice is chosen to be the cathode. We do not consider the anion dynamics, but we implicitly introduce it by the creation of particles and by charge electro-neutrality .
The initial condition evolves in time as follows: Every time step a walker is chosen and moved to one of its four neighboring sites with probabilities taken from a finite differences scheme of Eqs. (2) and (4) : probability $`1/(4+p)`$ to move either left, right or upwards, and probability $`(1+p)/(4+p)`$ to move down, i.e., towards the cathode. The parameter $`p`$ is referred to as the bias; in galvanostatic conditions it can be quantitatively related to the electric current density in the physical system as shown in . Let us stress here that our present choice for the probabilities is different from that reported in and , but we have checked that the results hardly differ with those presented in this paper. The main reason for this new election is that, with the new rules, the bias $`p`$ ranges from $`0`$ to $`\mathrm{}`$, that is, from pure multiparticle DLA to ballistic deposition, whereas the rules in the mentioned references allow for a range in $`p`$ from $`0`$ to $`0.25`$, and the ballistic deposition limit cannot be reached (although $`p=0.25`$ is rather close already, see ). After a destination site has been chosen, the particle moves if that node of the lattice is empty; if not, we select another particle and repeat the destination selection procedure. Once the particle has been moved, if the new position has any nearest neighbor site belonging to the aggregate, the present position of the walker is added to the aggregate (the cathode or bottom boundary at the initial stage) with probability $`s`$ (and is able to diffuse over the aggregate surface if that aggregate position has just one nearest neighbor belonging to it, see the following section); otherwise it stays there (and is able to move again) with probability $`1s`$. We term $`s`$ the sticking probability; it is related to the chemical activation energy the cation needs to stick to the aggregate. As particles are added to the aggregate, other ones are created at the top of the lattice keeping the mean cation concentration $`c`$ constant, which in fact simulates an infinitely high system (experimentally this means that the distance between electrodes is much larger than their lateral dimension); consequently, the flux of particles is constant at every stage of the simulation.
As we have already pointed out, the model parameters are related to the physical factors influencing the problem. Indeed, the choice of jump probabilities for the random walkers in the bath provides a recurrence relation which is a discretized version of the continuous equations (2) and (4). Therefore, the drift velocity $`\mu _𝐜𝐄`$ is proportional to the bias $`p`$. When a finite number of walkers is considered with concentration $`c`$, we must take into account the excluded volume, so the effective diffusion coefficient and the effective drift velocity in the simulations are proportional to $`1c`$ (in a mean field approach) . It is important to note that when $`c0`$, i.e., the bath is formed by one particle alone (as in DLA), the aggregate develops tall branches which grow at the expense of short ones due to screening effects. Therefore, in the low current limit a morphological instability appears that is not always present in ECD experiments. The finite concentration and the hard core interaction among random walkers simulates the cation pressure on the aggregate, so $`c`$ is an essential ingredient in the understanding of the formation of electrodeposits and to prevent these instabilities (of Laplacian character) from dominating the whole growth process.
One important task is the definition of simulation time step. In , comparison with the experiments in allowed to show that the physical time and the simulation time measured in number of Monte Carlos trials were simply proportional to each other. For this reason, we have stuck to the definition of the time step in as a Monte Carlo trial, i.e., the time needed for a particle to jump, both if the particle does jump or if it does not. Notwithstanding, we have tried other time steps definitions, such as the Monte Carlo step being defined as the mean time for every random walker to jump at least once, but the results are basically the same. Some authors define the time step for solid-on-solid growth models as the mean time needed to complete an aggregate layer, but as we will show below, ECD electrodeposits do not grow with constant velocity, and therefore the mean interface height does not grow linearly with time. We thus believe that, in the ECD context, this time unit would be rather artificial and hence we have not used it. In fact, as we will show below, the work reported in this paper provides further evidence in favor of our choice (see the discussion of the experiments in in Sec. III B below).
## III Numerical Results
### A Morphologies
We begin the summary of our results by discussing the morphologies generated by MBDLA with and without SD and the influence of the different rules for SD on them. In addition, we want to compare our computer generated morphologies to the available experimental data. We take as a reference the comprehensive experimental work of Trigueros et al. , who reported a systematic experimental study of different growth regimes at constant applied voltage conditions. Their work gave rise to a diagram of morphologies divided into different regions in which similar morphologies were obtained as a function of the applied voltage and the electrolyte concentration. It is important to realize that, in galvanostatic conditions, there is no linear correspondence between voltage and electric current of ions, and therefore, comparison between our morphologies and those reported by these authors can only be qualitative. No similar taxonomy work has been performed for constant current conditions. Although the diagram in is quite complex, it encloses a full variety of morphologies under the label compact. Some authors have studied electrodeposits systems within this regime, and hereafter we will also refer to them. Finally, a recent work by Schilardi et al. provides exhaustive information on the asymptotic ECD regimes, which have not been considered anywhere else; hence, their research will also be compared to ours along the paper.
As we have already mentioned, from the model perspective we can compare the bias $`p`$ with the electric current density , and $`c`$ with the electrolyte concentration, even though the two latter magnitudes are not exactly coincident, i.e., an electrolyte concentration equal to $`0.1`$ M does not mean $`c=0.1`$. We will see below that the results are not very sensitive to the specific value of $`c`$ insofar it is not very small, and thus the difference between actual and model concentrations is not very relevant. The sticking probability, $`s`$, and the diffusion parameters, namely $`l`$, $`\lambda `$ and $`r`$ (or equivalently $`\tau _d`$), cannot be directly tuned in an experiment, although it is reasonable to expect that changes in the experimental conditions will correspondingly modify these parameters. How much are they modified is something we will learn through our computer simulations.
#### 1 Bias vs. Sticking Probability without SD
Figure 1 shows a diagram of morphologies obtained with $`0p5`$ and $`0.01s0.5`$ without SD, with a particle concentration $`c=0.05`$. We have included these results for two reasons: First, there has been no previous report on MBDLA morphologies, except for a brief discussion in ; and, second, we need to discuss them in order to understand later what is the effect of SD on MBDLA morphologies. It is clear from Fig. 1 that increasing the bias or decreasing the sticking probability yields denser aggregates, the ones obtained for $`p=0`$ and $`s=1`$ (bottom right) being multiparticle DLA-like as expected (compare to ). This phenomenon is related to the stabilizing effect of the parameters $`p`$ and $`s`$, which can be theoretically demonstrated . Indeed, the higher the values of $`p`$, the larger the flux of particles reaching the interface in the direction perpendicular to the cathode. This reduces the probability for a cation to stick laterally to a branch and the screening effects due to the Laplacian field. On the other hand, the electric field combined with the reduction of the sticking probability tends to fill the interface valleys. This first result, namely the fact that increasing the electric current leads to denser aggregates, is similar to the results reported by Trigueros et al. , who observed densification of the aggregates with increasing applied voltage. In particular, we can qualitatively compare the morphological changes obtained by varying the bias $`p`$ for a fixed $`s=0.5`$ in Fig. 1, with those provided by experimental voltage variations (see Fig. 2 in ). We conclude that high voltages (or in general, high density currents) yield denser aggregates. So the bias $`p`$ is an essential ingredient in any realistic ECD model.
As a second step in our study, we have monitored other relevant quantities which in turn can be experimentally measured, in order to obtain additional information aside from qualitative morphological comparisons. Figure 2 shows the local concentration of particles in the bath, still without SD, at equal time intervals. We have plotted the concentration profiles in the stationary regime, i.e., after the instability occurs (see below). Thus, the mean number of attached particles per unit time (or equivalently, the mean interface velocity) is constant. Léger et al. have reported experimental evidence consistent with this stationary behavior (see e.g. Fig. 5 in ). We thus see that MBDLA agrees well with their findings, i.e., the stationary concentration of particles in the bulk obeys approximately the equation :
$$C(z)=c_a+(c_0c_a)e^{(zz_0)u/D_c},$$
(8)
where $`z`$ is the vertical coordinate, $`z_0`$ is the interface mean position, $`c_a`$ is the concentration at the anode, $`c_0`$ the concentration at the interface, $`u=\mu _c𝐄`$, and $`D_c`$ is the bulk diffusion coefficient. As shown in Fig. 2, this function provides a good fit of our data. In Fig. 3, we plot a fit of Eq. (8) (dashed line) to the simulation results, showing a good collapse of the bulk particle density outlines for different times. The small deviations close to the aggregate are due to the interface roughness. The ratio $`D/u`$ is called diffusion length. In our fits, this length turns out to be about $`15`$ lattice spacings, that is, about $`2`$ or $`3`$ times the lateral width of the branches for the chosen parameters. This result provides another of check the physical validity of our model, as we can compare the length obtained from the fit with that taken from Ref. . In this paper, the diffussion length is of order $`0.5`$ $`mm`$, about two times the typical branch lateral width (of order $`1`$ $`mm`$), so we may conclude that the diffusion length obtained from our model is physically consistent.
The inset in Fig. 3 shows the mean concentration front position $`z_0`$ vs. time, demonstrating that, in the stationary state, MBDLA leads to a constant velocity of the advancing front as in the experiments.
#### 2 Physical Relevance of SD
The previous subsection shows that MBDLA without SD successfully reproduces some ECD experiments, in particular, under galvanostatic conditions with not very small electric current density. However, within the MBDLA model it is impossible to understand the unexpected compactification of aggregates in low voltage experiments or the columnar-like growth found in other situations. Unfortunately, MBDLA aggregates are always ramified at low bias. In , a phenomenological explanation of compactification was proposed by noticing that the reduction of $`s`$ leads to more compact aggregates. Therefore, it was proposed there that $`p`$ and $`s`$ should be related by a monotonous function, the simplest case being that of a linear relationship. With this procedure, reducing the bias leads to a corresponding decrease in the sticking probability, and hence to compact aggregates at low bias. However, this is an ad hoc assumption that cannot be experimentally tested, whereas its theoretical justification is not very clear. Besides that, this approximation does not reproduce other morphologies, as that reported by López-Salvans et al. or Kahanda et al. . In view of this, it became increasingly clear that there was some crucial ingredient missing in MBDLA, and the most obvious candidate was of course SD.
At this point, it is instructive to consider carefully the work by Kahanda et al. . According to their results, as the absolute value of the overpotential decreases the aggregate becomes denser, and it is formed by several columns which are thicker at the top than at the bottom. We interpret this as a hint on the relevance of SD: If, when a particle arrives at the top of a column, it diffuses along the aggregate interface, and if the diffusion length is shorter than the column perimeter, the particle will not reach the base of the pillar or another column, with the result of a characteristic inverted triangle structure. The onset of similar triangle structures has also been reported by Pastor and Rubio . We thus came to the conclusion that it was necessary to include SD in MBDLA in order to shed further light on the nontrivial coupling of the different transport mechanisms.
#### 3 Implementation of SD in MBDLA
We have implemented SD in MBDLA in three different ways, all of them starting when a particle in the bulk (the electrolytic solution) sticks to the aggregate but has just one neighbor. We have first tried two simple irreversible rules (other similar rules yield equivalent results, so we do not include them here for brevity), named rules A and B, and an irreversible one, named rule C:
Rule A: The newly incorporated particle jumps always in the same direction, either left or right parallel to the cathode, until it reaches a site with at least $`2`$ neighbors or completes $`l`$ jumps. This rule is similar, but not identical, to the one studied in for ballistic deposition with surface diffusion.
Rule B: In this second rule, we allow the particle to perform a random walk over the aggregate surface until it increases its coordination number, with a constant probability $`\lambda `$ to be permanently stuck to its current position (this is the so called mortal random walker ).
The last rule is characterized by Arrhenius-like jump probabilities and, what is more important, by simultaneous bulk diffusion and SD:
Rule C: This rule allows several particles to diffuse simultaneously. When a particle arrives to a coordination $`1`$ site, it sticks and jumps to one of its two nearest neighbors on the aggregate with probability $`p_n=\mathrm{exp}[E_0+(n1)E_a]`$, where $`E_{0,a}`$ are adimensional activation energies, and $`n`$ the coordination number of the target position. If the particle new position has $`2`$ or $`3`$ neighbors, it attaches to the aggregate irreversibly. Otherwise, we label the particle as a SD particle, and we allow it to take further steps. Thus, we have two kinds of diffusing particles: Particles in the bulk, distributed homogeneously with concentration $`c`$; and particles that diffuse over the aggregate surface. With probability $`r`$ we choose a bulk particle which evolves with its characteristic rules, and with probability $`1r`$ a particle on the surface which jumps to one of its nearest neighbors as we have just described for the first jump. This rule is close in spirit to the collective diffusion rules employed in studies of kinetic roughening in molecular beam epitaxy (MBE) , and in particular to MBE models beyond the solid-on-solid approximation .
The main difference between rules A and B with respect to rule C is that the latter introduces a characteristic time scale $`\tau _d=r^1`$, while in the other cases diffusion is instantaneous; then, the diffusing particle is not affected by the overall particle dynamics. As we will show below, Rule C is the only one which actually reproduces the influence of SD on the aggregates scaling and morphology. In this respect, it is important to advance that we have found that Arrhenius-like probabilities by itself are not enough to model SD: Variants of rule C with those probabilities and without the characteristic time, i.e., SD kept instantaneous, lead once again to results similar to those of rule B. All the results presented were obtained with $`E_0=3`$ and $`E_a=1`$. We have chosen these values to have jump probabilities smaller than $`1`$, but other sets of parameters yield similar results which we omit for brevity. Finally, another interesting point is that the probabilities in Rule C allow to trivially introduce temperature in the model by simply identifying $`E_{0,a}E_{0,a}^{}/k_BT`$.
Rule A, by definition, introduces a diffusion length $`l`$, but if $`l1`$ the particle jumps practically always lead to an increment of its coordination, as may be seen in Fig. 4, where some morphologies are shown for different values of $`l`$. The inverted triangle structure typical of the experiments by Kahanda et al. is reproduced with this simple rule. Nevertheless, the tops of the pillars are unrealistically flat; another problem is that decreasing $`p`$ does not lead yet to a compact aggregate regime. Rule A is therefore not appropriate. In the case of Rule B, the diffusion length is introduced indirectly by means of the attachment probability $`\lambda `$ (see Appendix for details). The mean diffusion length can be shown to be given by $`l_D=1/(2\lambda ^{1/2})`$. Morphologies obtained with this rule are plotted in Fig. 5. Once again, and in spite of the fact that rule B allows the particles to diffuse randomly over the aggregate, the columns developed during the growth turn out completely flat at the top, and the option of rule B was excluded as well.
These pitfalls (and similar ones found by using Arrhenius-like jump probabilities, which we skip for brevity) led us to the conclusion that instantaneous SD (or limited mobility rules, in the terminology of ) is a too drastic one approximation for ECD. Taking into account that the main distinctive feature of MBDLA is its non-local character, interactions of particles diffusing along the surface with newly deposited particles are expected to be relevant. Guided by these ideas, we propose rule C, which incorporates this coupling by introducing time scales for both bulk diffusion and SD. A sample of the aggregates generated by MBDLA with rule C is shown in Fig. 6. The difference with the other two rules is immediately apparent from the plot: This more realistic rule does induce the creation of pillars as we pointed out above, this time similar to those reported by Kahanda et al. and Pastor and Rubio which are rough at the top. Moreover, the compactification of the aggregates at low currents appears naturally, as can be noticed by following the sequence of aggregates appearing on the same row (same value of $`r`$): Decreasing the current leads initially to less dense aggregates, until further reduction of the current gives rise to more compact aggregates. Remarkably, there is no need to change the sticking probability by hand as in MBDLA without SD or with Rules A and B. This allows us to eliminate one model parameter, the sticking probability, which we take to be $`s=1`$ from now on.
So far, we have seen that, while simple SD rules provide good results in some solid-on-solid simulation models, the complex dynamics of Laplacian systems does not allow the particles to instantaneously diffuse; rather, we must allow several particles to interact before they become permanently stuck to the aggregate. Roughly speaking, the flux of particles arriving at the aggregate defines a characteristic time $`\tau _p`$ (typically inversely proportional to the flux, i.e., to $`p`$). Once the particles have arrived to the aggregate, they diffuse until they reach a site with coordination larger than one, or equivalently, until the particle meets another diffusing particle, thus forming a dimer on the interface which cannot move anymore. A large flux of particles arriving at the interface (large $`p`$) will increase the probability of formation of those dimers, and the particles can hardly diffuse. The situation is not so simple when $`p`$ is small. On one hand, the deposition mean time $`\tau _p`$ is large, but on the other hand, the particles hardly experience the applied electric current, so the probability of attachment to a column wall before getting to the bottom of the aggregate increases. Thus the Laplacian instability is amplified leading to a compact structure formed by columns and grooves. This kind of instability has been observed in low current galvanostatic experiments . The aggregate is therefore denser but if the diffusion time is not long enough the interface is unstable. It is remarkable that this simple picture in terms of time scales allows to understand the relevance of SD in ECD experiments.
A final important remark we would like to mention is that, when the diffusion probability $`r`$ is about $`0.99`$, we have observed some evidences of what could be a morphological transition (and the subsequent change in the branches) similar to those reported by López-Salvans et al. . However, as we want to concentrate in this paper on MBDLA with SD as a generic model for all regimes of ECD experiments, we postpone a more careful study of this possibility to future work, where we will pursue the appearance of this phenomenon for different model parameters (such as $`p`$ or $`r`$).
#### 4 Electrolyte concentration
To conclude the analysis of MBDLA parameters, we show the effect of the electrolyte concentration, $`c`$. Figure 7 exhibits the morphological changes in patterns with different $`c`$ values ranging from $`0.01`$ to $`0.1`$ for different bias $`p`$ without SD. Note that when $`c0`$ the low current limit is exactly the DLA growth model . Therefore, we should keep a finite value of $`c`$ in order to diminish the unavoidable DLA characteristic instability. The results contained in the figure allow us to conclude that, insofar $`c`$ is not very small, the morphologies obtained with MBDLA do not depend strongly on the concentration, and therefore the fact that there is no direct correspondence between physical and simulated concentrations is not a drawback of the model.
### B Dynamic Scaling
The previous subsection shows that the inspection of the morphologies is a valuable method to check the validity and relevance of the model rules. Indeed, the unrealistically flat aggregates obtained with diffusion rules A and B unable them and motivate the investigation of the more realistic, non-instantaneous rule C for SD. However, in order to exploit the main virtues of MBDLA with SD and to compare with other relevant models and experiments, we must take some quantitative criteria, for example, the analysis of the interface surface roughening. To this end, let us define some functions related to the height of the aggregate at spatial position $`x`$ at time $`t`$, given by the scalar field $`h(x,t)`$. We will also review their basic features before discussing MBDLA properties.
The global width (or roughness) $`W(L,t)`$ is nothing but the rms fluctuations of the height variable $`h(x,t)`$ around its mean value $`\overline{h}_L(t)=(1/L)_xh(x,t)`$:
$$W^2(L,t)=\frac{1}{L}\underset{x}{}\left[h(x,t)\overline{h}_L(t)\right]^2,$$
(9)
where angular brackets stand for ensemble average.
Generally speaking, in many growth models, starting from $`h(x,0)=0`$ the width satisfies the dynamic scaling hypothesis of Family-Vicsek :
$$W(l,t)\{\begin{array}{ccc}t^\beta \hfill & \mathrm{if}\hfill & tL^z,\hfill \\ L^\alpha \hfill & \mathrm{if}\hfill & tL^z.\hfill \end{array}$$
(10)
The roughness exponent $`\alpha `$, the dynamic exponent $`z`$, and their ratio (growth exponent) $`\beta =\alpha /z`$, identify the universality class the model belongs to.
In the study of kinetic roughening the height-height correlation function is frequently used :
$$C^2(l,t)=\frac{1}{L}\underset{x}{}\left(h(x+l,t)h(x,t)\right)^2,$$
(11)
where ,
$$C(l,t)\{\begin{array}{ccc}t^\beta \hfill & \mathrm{if}\hfill & tl^z,\hfill \\ t^{(\alpha \alpha _{loc})/z}l^{\alpha _{loc}}\hfill & \mathrm{if}\hfill & tl^z,\hfill \end{array}$$
(12)
where $`\alpha _{loc}`$ is the so called local roughness exponent. Another important function related to the height variable $`h`$ is the power spectrum:
$$S(k,t)=\widehat{h}(k,t)\widehat{h}(k,t),$$
(13)
where $`\widehat{h}(k,t)=L^{1/2}_x[h(x,t)\overline{h}_L(t)]\mathrm{exp}(\mathrm{i}kx)`$. $`S(k,t)`$ displays a behavior consistent with the scaling form
$$S(k,t)=k^{(2\alpha +1)}s(kt^{1/z}),$$
(15)
where
$$s(u)=\{\begin{array}{ccc}u^{2\theta }\hfill & \mathrm{if}\hfill & u1,\hfill \\ u^{2\alpha +1}\hfill & \mathrm{if}\hfill & u1.\hfill \end{array}$$
(16)
The exponent $`\theta `$ takes different values depending on the type of scaling exhibited by the model. For instance, for the so called intrinsic anomalous scaling we have $`\theta =\alpha \alpha _{loc}`$, whereas $`\theta 0`$ for Family-Vicsek scaling (including super-roughening, i.e., $`\alpha 1`$). Note that this implies $`\alpha =\alpha _{loc}`$.
To apply these ideas to MBDLA characterization, a few remarks are in order. Although, in some cases, MBDLA develops ramified aggregates leading to multivalued interfaces, i.e., interfaces with overhangs, it has been demonstrated that the interface of the active zone in DLA simulations (the aggregate sites with larger probability of arrival) corresponds to that constructed by taking the topmost site $`h(x,t)`$ at every horizontal position $`x`$. This construction does not ensure that the measured exponents are free of interpretations , but the exponents are consistent with theoretical and experimental data . The reduction of the sticking probability $`s`$ yields denser aggregates, and overhangs do not appear at any stage of the simulation for low $`s`$ values. Besides that, if SD is present the aggregates are also more compact. In all these cases the function $`h(x,t)`$ is identical to the aggregate outline and consequently the results do not have any interpretation problem.
The main scaling features of MBDLA without SD were already reported in . Therefore, here we will briefly summarize them to facilitate comparison with results including SD, and refer the reader to for the details. Without SD, MBDLA displays three temporal regimes: At early times the global width $`W(L,t)`$ features $`\beta =0.5`$, this value being simply due to shot noise. This stage corresponds to times at which the lateral correlation length is of the order of the lattice spacing. After this noisy transient, short and large length scales are governed by different dynamics because the bulk Laplacian field produces nonlocal effects (screening or shadowing among branches). Consequently, the local and the global roughness exponents, $`\alpha _{loc}`$ and $`\alpha `$, are different and the interface is not self-affine. The growth exponent, $`\beta `$, is larger than that of noise ($`\beta >1/2`$) because some isolated branches begin to grow independently from each other, which can be understood as a signature of the Laplacian instability. As a consequence, the interface width grows rapidly as compared with the noise fluctuations. At later times, branches spread by lateral growth and impinge upon each other. Eventually, the system reaches an asymptotic regime characterized by the Kardar-Parisi-Zhang (KPZ) universality class exponents ($`\alpha =1/2`$, $`\beta =1/3`$, $`z=3/2`$). The KPZ equation is the paradigmatic growth model without SD, and it is given by the stochastic partial differential equation
$$\frac{h}{t}=\nu ^2h+\frac{\lambda _0}{2}\left(h\right)^2+\eta (x,t),$$
(17)
where $`\nu `$ and $`\lambda _0`$ are constants and $`\eta (x,t)`$ is a Gaussian white noise with:
$$\eta (x,t)=0,$$
(19)
$$\eta (x,t)\eta (x^{},t^{})=2D\delta (xx^{})\delta (tt^{}).$$
(20)
As mentioned above, the definition of the interface function $`h(x,t)`$ neglecting overhangs might cast some doubts on the validity of the exponents reported in . To confirm our results, we have measured the excess velocity produced by tilting the initial substrate and imposing helicoidal boundary conditions . The inset in Fig. 8 shows that this mean velocity is well fitted by a parabola, as expected for KPZ behavior. It is important to stress that identical results are obtained using the jump rules in .
Interestingly, Schilardi et al. report experiments with large currents (equivalent to the large values of the bias $`p`$) in excellent agreement with our model. They observe the same three time regimes: An initial transient with a behavior which could not be measured due to the resolution of the experimental device; a second transient with $`\beta >1`$ characterized by the growth of isolated branches, and a third asymptotic regime at which the interface is characterized by KPZ exponents. A plot of the mean interface velocity vs. time is also given, showing a crossover from the unstable regime to the stable one in accordance with MBDLA predictions as can be seen in Fig. 8. The global width crosses over from the instability ($`\beta >0.5`$) to $`W(L,t)t^{1/3}`$ at the time pointed out by the arrow. Note that MBDLA cannot yield $`\beta `$ larger than one, because of its discrete growth rules. This would mean that the interface width grows faster than the interface mean height. Finally, the evolution of the morphology during the experiment is also the same in MBDLA and in the experiment, as seen by comparing Fig. 9, taken from , and Fig. 10, obtained in our simulations.
We now consider the scaling behavior of MBDLA+SD for the different diffusion rules. As we have pointed out in the preceding section, large values of the diffusion length $`l`$ (rule A) generate flat aggregates. This means that $`\beta 0`$ as $`l\mathrm{}`$ at early times. Fig. 11 shows the lack of universality in the growth exponent $`\beta `$: It can be seen in this plot that $`\beta `$ decreases with $`l`$ as we expected. The same happens with rule B: As in the case of rule A, the growth exponent $`\beta `$ depends strongly on the attachment probability, $`\lambda `$ (rule B). As depicted in Fig. 11, the dependence is similar to that of model A since the diffusion length is proportional to $`\lambda ^1`$.
The scaling behavior in MBDLA with SD given by rule C is more complicated. We can recognize three different kinds of behavior, which we summarize as follows:
$`\mathbf{0.05}𝐫\mathbf{0.25}`$. The characteristic diffusion time is long, and particles diffuse rather fast along the surface (let us recall that they are picked with probability $`1r`$ at every Monte Carlo trial) without much interaction with particles arriving from the bulk, thus yielding compact aggregates, except if $`p0.05`$, because then the Laplacian field creates pillars and grooves. After a short transient the global width grows slowly and, independently of the applied current, the roughness exponents are compatible with those of the Edwards-Wilkinson universality class, whose defining equation is :
$$\frac{h}{t}=\nu ^2h+\eta (x,t).$$
(21)
Figure 12 shows the global width collapse obtained by rescaling the simulation time. The plot not only shows the Edwards-Wilkinson growth exponent, but also the $`r`$ independence of the results on a wide range of simulation parameters. Note that the collapsing time step is the one defined for MBDLA without SD divided by the characteristic diffusion time $`\tau _d=r^1`$. Figure 13 shows the collapsed power spectrum using $`\alpha =1/2`$ and $`z=2`$ (and consequently $`\beta =1/4`$) consistent with (13) with $`\theta =0`$ for Edwards-Wilkinson exponents. It is important to note that this kind of dynamic scaling has been observed in two-dimensional ECD experiments . Finally, we have to mention that the restriction $`r>0.05`$ is only due to the extremely long computational times needed to study the model for such small values of $`r`$.
$`\mathbf{0.3}𝐫\mathbf{0.7}`$. For large $`p`$, the interface is compact and grows with constant velocity. The scaling is similar to that of the preceding case. When $`p0`$, initially the interface is rough and the growth exponent $`\beta `$ is in the range $`0.350.40`$ (see Fig. 14). Some experiments have reported similar interfaces at early stages of growth : Specifically, they obtained exponents consistent with the linear MBE growth model universality class ($`\alpha =3/2`$, $`\beta =3/8=0.375`$ and $`z=4`$) that is, their interfaces could be described by the equation :
$$\frac{h}{t}=K^4h+\eta (x,t).$$
(22)
Note that for this model $`\alpha >1`$, so the interfaces generated with Eq. (22) are super-rough. In our case this short regime ceases when the mean interface height, $`\overline{h}(t)`$, is about $`8`$ to $`10`$ monolayers and the global width, $`W(L,t)`$, is about $`1`$. This is compatible with the referred experiments except that we do not observe the super-rough power spectrum. Actually, in our case the tail of $`S(k)`$ presents a time shift at large wave vectors (Fig. 15) which is incompatible with the behavior obtained for Eq. (22). However, the basic phenomena, such as the value of the effective $`\beta `$ and the onset of the instability, are in good agreement with the experiments. After this transient, the aggregates are still compact and develop some grooves (see Fig. 6). When these grooves appear, the growth exponent $`\beta `$ rises dramatically due to the large slopes produced between grooves. Figure 16 summarizes all this by showing the variation of $`\beta `$ with time.
$`\mathbf{0.85}𝐫`$. Finally, when the diffusion time is short, three completely different situations are found as a function of the current $`p`$. For very large $`p`$, cations become ballistically driven to the aggregate and the unstable transient tends to dissappear (in fact, the $`p\mathrm{}`$ limit is the Ballistic Depostion discrete model, which is well known to belong to the KPZ universality class ). When $`p1`$ the aggregate grows as MBDLA without SD with similar parameters, except that in this case the aggregate mean density rises. That is, we succesively detect a noisy initial transient, the instability associated with the growing branches, and the KPZ asymptotic limit due to the lateral growth of the branches. The interfaces within the unstable regime (an example of which is shown in Fig. 17) are not self-affine but present intrinsic anomalous scaling : Figure 18 shows the power spectrum for $`r=0.85`$, $`p=4`$ and $`s=1`$. Figure 19 shows the collapse of the power spectrum and Fig. 20 the collapse of the height-height correlation function, $`C(l,t)`$, achieved in both cases for $`\alpha =1.78`$, $`\alpha _{loc}=0.49`$, $`z=2.51`$ and $`\beta =0.71`$.
For intermediate $`p`$ values (between $`0.25`$ and $`1`$, for almost every $`r`$), the aggregate is formed by several compact thin branches which grow vertically and parallel to one another. In this case, the notions of a rough surface or dynamic scaling are meaningless. Finally, for small $`p`$ some compact branches grow at the expenses of the others, so typically, one or two branches grow more than the others. As in the case of the parallel branches, it is meaningless to talk about interface roughening.
## IV Discussion and Conclusions
Our first conclusion is that MBDLA is a simple computational model which incorporates in a natural way some of the basic mechanisms involved in ECD experiments. The original model was already known to be in good agreement with some experiments . In this paper, we have provided much more evidence showing that MBDLA explains some of the morphological changes due mainly to the applied electric current and, what is more important, it predicts the recently observed KPZ scaling behavior in the high current limit (for which SD is not too relevant) and observed also at low currents . Before this regime is reached, there is an unstable transient within which MBDLA interfaces present intrinsic anomalous scaling. We believe this type of scaling is due to SD not being able to communicate different portions of the interface fast enough, so that they grow independently from one another. This is analogous of the anomalous scaling occurring in the non-linear surface diffusion equation studied in . In our case, the different portions feature a value of the roughness exponent $`\alpha _{loc}0.5`$, similarly to the interface subject to columnar disorder studied in .
Secondly, the main point of our paper is that, as we have seen, MBDLA without SD cannot explain low current experiments in which the characteristic dense branching aggregates of high current experiments are replaced by compact and column-like aggregates. Our working hypothesis was that the latter kind of patterns are due to the competition between the Laplacian field of the cations in the disolution and the SD current on the interface. Thus, our ECD model, which we wanted to improve as to explain, at least qualitatively, the complete ECD phenomenology, should incorporate a new rule for the diffusion of the adatoms attached to the aggregate. Hence, we have tried out some SD rules similar to those often used in growth models for molecular beam epitaxy . We have verified that instantaneous diffusion rules, namely rules that “freeze” the bulk particles while the most recently attached particle finds its way through the surface, do not lead to correct results in the low current limit, and produce very unrealistic, flat-topped morphologies. We have then been forced to conclude that the nonlocal character of MBDLA demands a diffusion rule which couples the overall cation dynamics: This is the rule we have named C. It introduces a characteristic diffusion time $`\tau _d=r^1`$ which competes with the time scale related to the net flux of particles arriving to the interface (which, in fact, is proportional to the applied electric current density). With this SD rule, the morphologies at low, medium and high currents are compatible with those observed by Trigueros et al. for low, medium and high applied voltages respectively. This diffusion time $`\tau _d`$ cannot be controlled from the experimental point of view, but fortunately there are wide ranges of parameters over which the simulated morphologies hardly change, which means that the description of the experiments provided by MBDLA with SD is robust and does not need an incontrollable parameter to be tuned. We have also compared MBDLA with SD with the experiments reported by Pastor and Rubio , which characterize the product interfaces by the MBE exponents. MBDLA seems to reproduce the latter behavior for very short times and short length scales as can be seen in Figs. 14 and 15, but these results are not too significative, as they are not as accurate as we would need to make any strong claim, and could be due to the appearance of a characteristic short length scale. There is a difficulty in this respect, as MBDLA scaling is intrinsically anomalous whereas the results in support standard, super-rough scaling. With the data presently at hand we have to conclude that MBDLA with SD does not describe quantitatively all the aspects of the very low current regime, but the fact it does describe most of them and, above all, the compactification of the aggregates, makes us confident that MBDLA with SD is a very good general model for ECD.
To conclude, we note that the model presented has the basic ingredients of ECD phenomena: Diffusion, electromigration and surface diffusion, but for this reason, we have to pay a big price in terms of computational time. MBDLA without SD is a very time-consuming model, and the diffusion rules make the analysis and the simulations a patience exercise. It has certainly been an improvement to find that SD rule C allows us to skip the sticking probability parameter, thus resducing the parameter space, but even then if averages of relevant quantities over large ensembles are required, a great deal of computational resources would be needed. Of course, this disadvantage can be removed by a careful reprogramming of the algorithm, but that is another line of research. As our goal was to identify the most important factors involved in ECD, we do believe that despite the computing limitations of the model, MBDLA with SD is a powerful tool to repoduce some unclear features of this kind of growth experiments, and has helped understand what are the most relevant transport properties and how they couple in different parameter regions. We hope that this work suggests further experiments to find out whether MBDLA with SD is the complete, general model for ECD or else if there are still regions which need separate modelling.
###### Acknowledgements.
It is a pleasure to thank Miguel Angel Rodríguez and Francesc Sagués for their help through many discussions on the modeling of ECD and on the subtleties of dynamic scaling during the last few years, as well as for a critical reading of the manuscript. In addition, we are thankful to Javier Buceta, Juanma Pastor, Javier de la Rubia, and Miguel Angel Rubio for sharing with us their results and expertise on ECD at very low currents. We would also like to thank Roberto Salvarezza and Leticia Schilardi for their kind permission to reproduce Fig. 9 from . We are indebted to José Cuesta and Ricardo Brito for helpful suggestions about the simulation rules and to the members of GISC for their interest in this work. Work at GISC has been supported by DGES (Spain) Grant No. PB96-0119 and CAM Grant No. 07N/0034/98.
## Surface Diffusion Rule B
In Rule B, surface diffusion starts when a particle has first arrived to the aggregate and has attached to it (with probability $`s`$). The particle jumps with equal probability to one of its two nearest neighboring sites on the aggregate until it increases its coordination number. The particle has an additional probability, $`\lambda `$, of being permanently attached. This kind of particle is usually termed a mortal random walker . The random walk is performed between two absorbing boundaries, namely, a couple of sites with higher coordination (two or three). One could try to determine a priori the total number $`N`$ of jumps the particle has to perform in each realization drawing such number from the probability for the particle to take $`N`$ steps on a flat line if it avoids to stick $`N1`$ times and die at the $`N`$th jump. This probability can be easily calculated to be given by
$$P_N(\lambda )=\lambda (1\lambda )^{N1}.$$
(23)
However, the absorbing boundaries disallow this procedure. In any case, we have compared the simulation results by allowing the particle to perform an actual mortal random walk, and to perform a simple random of $`N`$ steps given by Eq. (23). Both results are hardly different. Thus, we can approximately calculate from Eq. (23) the mean and variance of the maximum number of jumps, given by:
$`\overline{N}={\displaystyle \frac{1}{\lambda }},`$ (24)
$`\sigma _N={\displaystyle \frac{\sqrt{1\lambda }}{\lambda }}.`$ (25)
For a flat interface, the particle mean position would be $`0`$ but its variance would be
$$\sigma _N=\overline{N}^{1/2}/2=1/(2\lambda ^{1/2}),$$
(26)
which provides the characteristic diffusion length $`l_D=1/(2\lambda ^{1/2})`$.
|
warning/0003/cond-mat0003403.html
|
ar5iv
|
text
|
# References
MAGNETIC BEHAVIOR OF A MIXED ISING FERRIMAGNETIC MODEL IN AN OSCILLATING MAGNETIC FIELD.
G. M. Buendía, E. Machado
Departamento de Física. Universidad Simón Bolívar,
Apartado 89000, Caracas 1080, Venezuela.
ABSTRACT
The magnetic behavior of a mixed Ising ferrimagnetic system on a square lattice, in which the two interpenetrating square sublattices have spins $`\sigma `$ ($`\pm 1/2`$) and spins $`S`$ ($`\pm 1,0`$), in the presence of an oscillating magnetic field has been studied with Monte Carlo techniques. The model includes nearest and next-nearest neighbor interactions, a crystal field and the oscillating external field. By studying the hysteretic response of this model to an oscillating field we found that it qualitatively reproduces the increasing of the coercive field at the compensation temperature observed in real ferrimagnets, a crucial feature for magneto-optical applications. This behavior is basically independent of the frequency of the field and the size of the system. The magnetic response of the system is related to a dynamical transition from a paramagnetic to a ferromagnetic phase and to the different temperature dependence of the relaxation times of both sublattices.
INTRODUCTION
The behavior of ferrimagnetic compounds in the presence of oscillatory fields have long been used for technological applications such as high-density magneto-optical recording , but little is known about the mechanisms responsible for this behavior. In a ferrimagnet the different temperature dependencies of the sublattice magnetizations raise the possibility of the appearance of compensation temperatures: temperatures below the critical point, where the total magnetization is zero . It has been shown experimentally that the coercive field is very strong at the compensation point favoring the creation of small, stable, magnetic domains . This temperature dependence of the coercivity near the compensation point can be applied to writing and erasing in high-density magneto-optical recording media, where the temperature changes are achieved by local heating the films by a focused laser beam. It has been shown that magneto-optic thin films with compensation temperatures higher than room temperatures can attain a direct overwrite capability . As far as we know there has been only very few crude attempts to reproduce theoretically the increase of the coercivity near the compensation point using mean-field approaches .Recently new classes of magnets are being synthesized with molecular organic chemistry techniques . Biocompatible, organic materials, optically transparent, with spontaneous moments at room temperature are not far from reality. Ferrimagnetic ordering seems to play a fundamental role in some of these materials. Ferrimagnetic compounds called Prussian blue analogs, with a critical temperature of 240 K have been reported . Organometallic compounds as the amorphous $`V(TCNE)_xy(solvent)`$ where TCNE is tetracyanoethylene are believed to have ferrimagnetic structure and ordering temperatures as high as 400 . Some of these compounds have compensation temperatures near 30 K . Most of these compounds have been synthesized by assembling molecular building blocks of different magnetic moments in such a way that adjacent magnetic moment are antiparallel . Since real ferrimagnets have extremelly complicated structures mixed Ising models have been introduced as simple systems that can show ferrimagnetic behavior and may show compensation points when their Hamiltonian includes second-neighbor interactions . In this article, we present a Monte Carlo study of a mixed Ising spin system, where spins that can take the values $`\pm 1/2`$, and spins that can take the values $`\pm 1,0`$, are nearest neighbors on a 2-dimensional square lattice and interact antiferromagnetically. Spins of the same type are next-nearest neighbors. We analyze the magnetic response of this system in the presence of an oscillating magnetic field. From these studies, we determine the dynamic order parameter, the coercive field and their variation with temperature, frequency and amplitude of the applied field, and size of the system. The results reproduce the rapid increase of the coercitivity at the compensation temperature. The dynamical order parameter calculations suggest that the model exhibits a phase transition between a paramagnetic and a ferromagnetic region. A similar result was observed by a mean-field study of a simpler version of this model . Mean-field approaches and Monte Carlo simulations indicates the presence of a dynamical phase transition in a kinetic Ising model . However the distinctive behavior of the coercive field at the compensation temperature in ferrimagnets seems to be related to the different relaxation times of the sublattices.
THE MIXED ISING MODEL
Our model consists of two interpenetrating square sublattices. One sublattice has spins $`\sigma `$ that can take two values $`\pm 1/2`$, the other sublattice has spins $`S`$ that can take three values, $`\pm 1,0`$. Each $`S`$ spin has only $`\sigma `$ spins as nearest neighbors and vice versa.
The Hamiltonian of the model is given by,
$$=J_1\underset{nn}{}\sigma _iS_jJ_2\underset{nnn}{}\sigma _i\sigma _k+D\underset{j}{}S_j^2H(t)\left(\underset{i}{}\sigma _i+\underset{j}{}S_j\right)$$
(1)
where the $`J`$’s are exchange interaction parameters, $`D`$ is the crystal field and $`H`$ is an oscillating magnetic field of the form,
$$H(t)=H_0cos(\omega t)$$
(2)
where $`\omega `$ is the frequency of the external field, its period is given by, $`\mathrm{\Theta }=2\pi /\omega `$. The $`J`$’s, $`D`$ and $`H_0`$ are all in energy units. We choose $`J_1`$$`=`$$`1`$ such that the coupling between nearest neighbors is antiferromagnetic.
Previous results with Monte Carlo and Transfer-Matrix techniques have shown that the $`J_1`$$``$$`D`$ model ($`J_2`$ and $`H`$ are equal to zero) does not have a compensation temperature. These studies show that a compensation temperature is induced by the presence of the next-nearest-neighbor (nnn) ferromagnetic interaction, $`J_2`$, between the $`\pm 1/2`$ spins. The minimum strength of the $`J_2>0`$ interaction for a compensation point to appear depends on the other parameters of the Hamiltonian .
MONTE CARLO CALCULATIONS
We use standard importance sampling techniques to simulate the model described by Eq. (1) on a $`L\times L`$ square lattice with periodic boundary conditions. Configurations are generated by randomly choosing spins on the lattice and flipping them one at a time according to a heat bath algorithm. In each complete sweep through the lattice $`L\times L`$ sites are visited. Each Monte Carlo step per spin is associated to a time interval, $`\tau _S`$, such that the frequency of the external field can be written as,
$$\omega =\frac{2\pi }{(NMCS)\tau _S}$$
(3)
where $`NMCS`$ is the number of Monte Carlo steps per spin necessary to cover an entire cycle of the field. To perform the simulations we arbitrarily choose $`\tau _S`$ to be one, such that $`\mathrm{\Theta }=NMCS`$. Our program calculates the sublattice magnetizations per site at the time $`t`$ defined as,
$$M_1(t)=\frac{2}{L^2}\underset{j}{}S_j(t),M_2(t)=\frac{2}{L^2}\underset{i}{}\sigma _i(t)$$
(4)
and the total magnetization per spin at the time $`t`$, $`M(t)=\frac{1}{2}\left[M_1(t)+M_2(t)\right]`$. The averages are taken over all configurations, the sums over $`j`$ are over all sites with $`S`$ spins, and the sums over $`i`$ are over all sites with $`\sigma `$ spins. Each sum has $`L^2/2`$ terms.
The compensation temperature is defined as the temperature below the critical, $`T_{\mathrm{comp}}<T_{\mathrm{crit}}`$, where the two sublattice magnetizations cancel each other such that the total magnetization is zero, i.e.,
$$|M_1(T_{\mathrm{comp}})|=|M_2(T_{\mathrm{comp}})|$$
(5)
and
$$\mathrm{sign}[M_1(T_{\mathrm{comp}})]=\mathrm{sign}[M_2(T_{\mathrm{comp}})].$$
(6)
To characterize the time behavior we calculate the dynamical order parameter $`Q`$ defined as,
$$Q=\frac{2\pi }{\omega }M(t)𝑑t.$$
(7)
The closed integral implies that the integral is performed over a cycle of the external magnetic field.
RESULTS
The value of Q is calculated by averaging its values over 100 cycles of the external field, once the system is in its stationary state. Most of the measurements were done for a $`L=40`$ lattice. Lattices of different sizes were used to study the finite-size effects. In Fig. 1 we show a hysteresis loop, $`M(t)`$ vs $`H(t)`$, for a particular combination of parameters in the Hamiltonian. The coercive field $`H_c`$ is defined as the minimum value of the external field needed for the total magnetization to go to zero, as is indicated in the figure. In Fig. 2 and Fig. 3 we show the coercive field vs the temperature for oscillating fields of several amplitudes, $`H_0`$. In the same figures we also plot the total magnetization for the equivalent system subject to a constant field of magnitude $`H_0`$. Notice that the compensation temperature, defined as the point where the total magnetization is zero, previous verification that Eq. 5 and 6 are satisfied, increases with the magnetic field, whereas the temperature at which the magnetization becomes discontinuous does the opposite. At a certain field, which amplitude depends on the parameters of the Hamiltonian, both temperatures become equal and for any field of larger amplitude there is no more compensation point, as can be seen in Fig. 3(b). From the figures it is clear that the coercive field increases in the vicinity of the compensation temperature where it reaches its maximum . These results are summarized in Fig. 4. As expected, the maximum value of the coercive field at the compensation temperature is given by $`H_0`$.
It is interesting to notice the asymmetric behavior of the coercive field around the compensation point. In the low temperature region, $`T<T_{\mathrm{comp}}`$, the coercive field decreases with increasing $`T`$ until it reaches a minimum, after which grows rapidly reaching its maximum at $`T_{\mathrm{comp}}`$, when $`T>T_{\mathrm{comp}}`$ the coercive field decreases. Notice that for small values of $`H_0`$ there is a range of temperatures for which the coercive field is not defined (see Fig. 4). This behavior of the coercive field has been observed experimentally . This result can be understood by looking at Fig. 5 were it is shown how the hysteresis loop changes with the temperature. As the temperature increases the loop moves in such a way that the coercive field increases until it reaches its maximum, after which, if the temperature keeps increasing, the loop stays below (or above) the $`M=0`$ axis without crossing it, meaning that the applied field is not strong enough to flip the spins. If we look at Fig.6 where we plot the coercive field and the dynamical order parameter vs the temperature we see that there is a dynamical phase transition between a paramagnetic region, $`Q0`$ and a ferromagnetic region $`Q0`$, the region where the coercive field is not defined is well into the ferromagnetic phase where the magnetization does not changes sign.
By changing the size of the system and studying the behavior of the coercive field, see Fig. 7, we notice that there are finite size effects, particularly evident for small systems ($`L<20`$). However, for larger systems, the location of the peak of the coercive field around the compensation temperature seems to be independent on the size of the system. For small systems ($`L<20`$) the peak of the coercive field appears before the system reaches its compensation temperature. Also the coercive field seems to decrease more rapidly for the larger systems.
In Fig. 8 we present some results that show the dependence of the coercive field with the size of the system, these results agree qualitatively with the experimental behavior of magnetic films and nanostructured Fe and Ni samples for which the coercitivity depends on the average size of the grain. The size dependence of the coercive field is very similar to the size dependence of the switching field of a kinetic Ising model (field at which magnetization reversal is thermally induced on experimental timescales for given temperatures and system sizes), which behavior has been shown to be strongly dependent on the modes by which the system decays .
Next, we explore how the results depend on the frequency. In Fig. 9 we show how the coercive field vs the temperature changes for different values of the frequency of the external field. We found a quite different response to the frequency of the magnetic field depending on the dynamical phase of the system. In the paramagnetic phase ($`Q0`$, see Fig. 6) the coercive field is larger for systems driven by fields with higher frequency, but in the ferromagnetic phase ($`Q0`$) just the opposite happens, as can been seen in Fig. 9. This behavior is related to the temperature dependence of the relaxation time in the different regions . In the ferromagnetic phase we must take into account the relaxation time of both sublattices the $`\sigma `$ and the $`S`$, whereas in the paramagnetic phase only the $`\sigma `$ one is relevant because the $`S`$ lattice follows the field with almost no delay . We also notice in Fig.9 that when the field has a high frequency the maximum value of the coercive field (that occurs at the compensation temperature) does not reach the field amplitude, i.e., the coercive field does not reach its saturation value.
In Fig. 10 we show the behavior of the coercive field vs the inverse frequency for different values of the temperature. If the field has a long period the coercive field seems to reach a value that is independent of the frequency and depends on the temperature. Again we expect that this behavior is related to the temperature dependence of the relaxation times in the different phases.
CONCLUSIONS
We have applied a Monte Carlo algorithm to the study of the magnetic response of a mixed Ising ferrimagnetic model to an oscillating magnetic field. We found that this model gives very good qualitatively agreement with the magnetic behavior of real ferrimagnets. It shows a rapid increase of the coercive field at the compensation temperature, a crucial feature that makes ferrimagnetic compounds extremely useful for thermo-optical applications. It also reproduces qualitatively the dependence of the coercitivity with the size of the sample observed experimentally. The results show the existence of a dynamical phase transition in which the mean-period averaged magnetization, $`Q`$, changes from $`Q0`$ to a $`Q0`$. Work in progress indicates that, as recent studies shows is also the case for the kinetic Ising model , some aspects of the hysteretic response as its dependence on the frequency and amplitude of the oscillating field, depends on the metastable decay mode. Explaining the different behavior in the different regimes depending on the nucleation mechanism (i.e., single-droplet or multidroplet ) by which the system decays c̃iteMachado2.
ACKNOWLEDGMENTS
We are indebted to Mark Novotny and Per An Rikvold for many insightful comments during the course of this work. G. M. B also acknowledges the kind hospitality of the Supercomputer Computations Research Institute of Florida State University at Tallahassee, Florida.
FIGURE CAPTIONS
Figure 1. Hysteresis loop ($`J_2=6`$, $`D=1.9`$, $`k_BT=0.5`$, $`\omega =\pi /30`$). The coercive field is indicated.
Figure 2. Coercive field and magnetization vs temperature. ($`J_2=6`$, $`D=1.9`$, $`\omega =\pi /100`$). a)$`H_0=0.1`$ b)$`H_0=0.5`$. Notice the discontinuity in the coercive field.
Figure 3. Coercive field and magnetization vs temperature. ($`J_2=6`$, $`D=1.9`$, $`\omega =\pi /100`$). a)$`H_0=0.8`$ b)$`H_0=2.1`$. For this choice of parameters there is not compensation point for $`H_0>1`$.
Figure 4. Coercive field vs temperature ($`J_2=6`$, $`D=1.9`$, $`\omega =\pi /100`$). The maximum value of $`H_\mathrm{c}`$ is given by $`H_0`$.
Figure 5. Hysteresis loop ($`J_2=6`$, $`D=1.9`$, $`\omega =\pi /100`$, $`H_0=0.5`$). Notice that for high temperatures there is no coercive field (see Fig. 6).
Figure 6. Coercive field and dynamical order parameter vs temperature ($`\omega =\pi /30`$, $`H_0=0.5`$, $`J_2=6`$, $`D=1.9`$).
Figure 7. Coercive field vs temperature for different lattice sizes ($`\omega =\pi /100`$, $`H_0=0.5`$, $`J_2=6`$, $`D=1.9`$).
Figure 8. Coercive field vs lattice size ($`\omega =\pi /100`$, $`H_0=0.5`$, $`J_2=6`$, $`D=1.9`$). The lines are guides for the eye.
Figure 9. Coercive field vs temperature ($`H_0=0.5`$, $`J_2=6`$, $`D=1.9`$).
Figure 10. Coercive field vs inverse frequency ($`H_0=0.5`$, $`J_2=6`$, $`D=1.9`$). The lines are guides for the eye.
|
warning/0003/physics0003040.html
|
ar5iv
|
text
|
# 1 Introduction
## 1 Introduction
The longbase laser interferometric gravitational wave detectors are under construction at present time \[1-3\]. Their sensitivity to metric perturbation will be about $`h10^{21}`$ that corresponds to the classical regime of operation. However for future installations with projected sensitivity $`10^{22}÷10^{23}`$ the quantum features of the measurement process can play a significant role. At the same time there are no limits of principle on the accuracy of measurement of external classical force. Therefore the methods and schemes which give the possibility to overcome the quantum measurement limitations (or the so called standard quantum limit, SQL) is of vital importance for future generation of gravitational wave experiments.
There are several procedures which allow to achieve the sensitivity larger than the SQL . For example in an optimal filtration procedure for the simplest variant of the optical sensor \- a mirror attached to a mechanical resonator and illuminated with a coherent pump field - was considered. An external force acting on the mechanical oscillator displaces its equilibrium position and thus changes the phase of the reflected field. The vacuum fluctuations of the input light act on the oscillator through the radiation pressure effect and constitute the back action noise of the measuring apparatus. For such system two quadratures of reflected wave are correlated. Using correlation (phase sensitive) processing of two quadratures one can increase a signal-to-noise ratio and overcome the SQL.
However the gain in sensitivity for the schemes overcoming the SQL is usually proportional to the square root of the ratio of laser power used for pumping the interferometer and an optimal power that corresponds to the point where the sensitivity of the interferometer achieve the SQL \[4-6\]. Unfortunately the optimal power is impracticably large, about several dozens of kilowatts that restrains the experimental implementation of the technique.
The pumping with the ultrashort periodic laser pulses can be technically advantageous over a continuous wave pumping for practical realization of the schemes overcoming the SQL. Actually for a large power a problem of generating a train of short high-intensity laser pulses can be technically easier than a problem of cw light generation (when the averaged powers for two cases are equal) because in the first case the energy in laser resonator is spread over the large frequency band (and different spatial longitudinal modes) and high intensities can be produced relatively easily. At the same time the amplitude and frequency stability of the pulsed pump in the case of a mode locked laser can be at the same level as for the monochromatic pump . For example in the stability of intermode beats for the mode locked laser output was estimated as $`510^{12}`$ in 10 s.
Another consideration is that the perspectives of squeezed states generation with high nonclassicality seem more realistic for the case of short high-intensity laser pulses allowing the use of squeezed pulsed pump in displacement transducers .
Finally an analog to digital conversion is usually used in modern experiment during the processing of the output. Therefore it seems natural to take the pulsed pump at once so that the output will comprise a set of the values for appropriate variable at definite times.
The goals of this article are to consider a displacement transducer consisting of a mirror attached to a mechanical oscillator and illuminated with a train of high-intensity laser pulses, to reveal the algorithm of optimal signal processing for such transducer and to estimate the sensitivity of the scheme to a measurement of classical external force.
The model of displacement transducer and basic equations of motion is considered in section 2. The sensitivities for traditional measurement scheme and for correlative processing of the output quadratures in the case of time independent pump are estimated in sections 3 and 4 correspondingly. The pulsed pump for the displacement transducer is considered in section 5. The conclusions are in section 6.
## 2 Model for displacement transducer and transformation of quadrature components
Let consider the most simple case of optical displacement transducer - a mirror attached to a mass of a mechanical oscillator and illuminated with a train of high-intensity laser pulses. An external force displaces an equilibrium position of mechanical oscillator changing the phase of reflected wave. The variation of the reflected field phase is measured by a readout system. This model is easy to calculate and it contains at the same time all features of displacement transducers with pulsed pump. For the incident $`E_\mathrm{i}`$ and reflected $`E_\mathrm{r}`$ waves one can use the quasimonochromatic approximation
$`E_\mathrm{i}`$ $`=`$ $`(A(tx/c)+a_1)\mathrm{cos}\omega _\mathrm{p}(tx/c)a_2\mathrm{sin}\omega _\mathrm{p}(tx/c)`$
$`E_\mathrm{r}`$ $`=`$ $`(B(t+x/c)+b_1)\mathrm{cos}\omega _\mathrm{p}(t+x/c)b_2\mathrm{sin}\omega _\mathrm{p}(t+x/c)`$ (1)
where $`A(tx/c)`$ and $`\omega _\mathrm{p}`$ are an amplitude (mean value) and a frequency of the pump wave, $`a_1`$ and $`a_2`$ are the operators of the quadrature components (fluctuations) of the pump wave (vacuum for coherent state$`),B(t+x/c)`$ is an amplitude (mean value) of the reflected wave, $`b_1`$ and $`b_2`$ are the operators of the quadrature components (fluctuations) of the reflected wave. The periodic envelope function $`A(tx/c)`$ consists of a train of equally spaced pulses and the duration of each pulse is considerably larger than the period of light wave but considerably smaller than the period of the mechanical oscillator.
To obtain the equation coupling the amplitudes of the incident and reflected waves for the moving mirror one can use a transformation of electromagnetic field for moving reference frame . For a constant velocity of the mirror $`V`$ one has
$$E_\mathrm{r}=[(1V/c)/(1+V/c)]E_\mathrm{i}\mathrm{exp}(2\mathrm{i}\omega _\mathrm{p}X/c)$$
(2)
where for simplicity the reflection coefficient of the mirror is taken to be $`r1`$ and $`X`$ is the position of the mirror. Let suppose that this expression is valid also for the slowly varying velocity $`V(t)`$ and position $`X(t)`$ of the mirror and $`V(t)c`$ (the validity of equation (2) has been proved for the mirror consisting of free electrons for the general case of relativistic velocity $`V(t)`$ in ). Then in linear approximation in $`V/c`$ one can obtain from equation (2) the following expression
$$E_\mathrm{r}=(12V(t)/c2\mathrm{i}\omega _\mathrm{p}X(t)/c)E_\mathrm{i}$$
(3)
The first term in (3) is an amplitude modulation of the reflected wave due to the mirror movement and the second is a phase modulation. For slow motion of the mirror $`V\omega _\mu X`$ ($`\omega _\mu `$ is a frequency of mechanical oscillator) and the second term in brackets is considerably smaller than the third term. Therefore for the transformation of the quadrature components of the field one can obtain
$`b_1(t)`$ $`=`$ $`a_1(t)`$
$`b_2(t)`$ $`=`$ $`a_2(t)+2\omega _\mathrm{p}A(t)X(t)/c`$ (4)
For the equation of mirror motion one has
$$\ddot{X}(t)+2\delta _\mu \dot{X}(t)+\omega _\mu ^2X(t)=M^1(F_\mathrm{s}(t)+F_\mathrm{p}(t)+F_{\mathrm{th}}(t))$$
(5)
where $`M`$ and $`\delta _\mu `$ are the mass and the damping coefficient of mechanical oscillator, $`F_\mathrm{s}(t)`$ is a signal force, $`F_\mathrm{p}(t)`$ is radiation pressure force and $`F_{\mathrm{th}}(t)`$ is a force associated with the damping of the oscillator. Let suppose for simplisity that $`\delta _\mu `$ tends to zero. Then the displacement $`X(t)`$ of the mirror will consist of two parts - a signal displacement $`X_\mathrm{s}(t)`$ and a radiation pressure displacement $`X_\mathrm{p}(t)`$. For $`F_\mathrm{p}(t)`$ one has
$$F_\mathrm{p}(t)=SA(t)a_1(t)/(4\pi )$$
(6)
where $`S`$ is a cross section of the laser beam. Therefore the equations of motion for the displacement transducer are
$`b_1(t)`$ $`=`$ $`a_1(t)`$
$`b_2(t)`$ $`=`$ $`a_2(t)+2\omega _\mathrm{p}A(t)X(t)/c`$ (7)
$`\ddot{X}(t)`$ $`+`$ $`2\delta _\mu \dot{X}(t)+\omega _\mu ^2X(t)=M^1(F_\mathrm{s}(t)+SA(t)a_1(t)/(4\pi ))`$
## 3 Sensitivity for a traditional measurement scheme
For traditional measurement scheme the amplitude of the pump is constant. Therefore one can easily obtain the transformation relations for the quadratures $`b_1`$ and $`b_2`$ from equations (7)
$`b_1(\omega )`$ $`=`$ $`a_1(\omega )`$
$`b_2(\omega )`$ $`=`$ $`a_2(\omega )+\lambda \xi (\omega )A^2a_1(\omega )+A\xi (\omega )F_\mathrm{s}(\omega )`$ (8)
where $`\xi (\omega )=2\omega _\mathrm{p}G(\omega )/c`$, $`G(\omega )=\left[M(\omega ^22\delta _\mu \mathrm{i}\omega +\omega _\mu ^2)\right]^1`$ is mechanical oscillator transfer function and $`\lambda =S/(4\pi )`$.
Only quadrature $`b_2`$ contains the signal and it is this quadrature that is measured in traditional measurement scheme . This corresponds to the measurement of the phase of the reflected wave. The first term in the right hand side of equation (8) for $`b_2`$ can be treated as an additive noise and the second term as a back action noise. For small pump amplitudes the sensitivity is increasing with the increase of $`A`$ because the signal is proportional to $`A`$. However for large pump amplitudes the second term in r.h.s. of (8) becomes dominant and the sensitivity is decreasing with the increase of $`A`$. Therefore there is an optimal value of the pump amplitude and the sensitivity to external force at this pump amplitude is just the SQL .
## 4 Correlative processing of quadratures for time independent amplitude of the pump
Two quadratures of the reflected field according to equation (8) have the dependence on the amplitude fluctuations of the incident field $`a_1`$. Therefore one can expect that the sensitivity can be increased for the correlative processing of the output . Actually if one combine with appropriate weight coefficients the quadratures $`b_1`$ and $`b_2`$ of the output wave then in this combination the noise term depending on $`a_1`$ can be cancelled. This weighting can be done by a homodyne detector with appropriate choise of a local oscillator phase $`\varphi `$.
Let the field of the local oscillator have the form
$$E_\mathrm{L}(t)=A_\mathrm{L}\mathrm{cos}(\omega _\mathrm{p}t+\varphi )$$
(9)
Then the photodetector output is proportional to the following expression according to (1), (9)
$$I_{\mathrm{pd}}A_\mathrm{L}(b_1\mathrm{cos}\varphi +b_2\mathrm{sin}\varphi )$$
(10)
and at certain frequency $`\omega _\mathrm{f}`$ one can obtain
$$I_{\mathrm{pd}}A_\mathrm{L}[a_1(\omega _\mathrm{f})(\mathrm{cos}\varphi +\lambda \xi (\omega _\mathrm{f})A^2\mathrm{sin}\varphi )a_2(\omega _\mathrm{f})\mathrm{sin}\varphi +A\xi (\omega _\mathrm{f})F_\mathrm{s}(\omega _\mathrm{f})\mathrm{sin}\varphi ]$$
(11)
Therefore choosing the phase $`\varphi `$ according to the equation ($`\xi (\omega _\mathrm{f})`$ is real for $`\delta _\mu =0`$)
$$\mathrm{cos}\varphi +\lambda \xi (\omega _\mathrm{f})A^2\mathrm{sin}\varphi =0$$
(12)
one can make the photocurrent insensitive to the amplitude fluctuations $`a_1`$ of the input field at certain frequency $`\omega _\mathrm{f}`$ of the signal. In this case the increase of the pump amplitude $`A`$ results in the relative increase of the output signal at frequency $`\omega _\mathrm{f}`$ according to the equation (11) with respect to the noise level defined by $`a_2`$.
For compensation of the back action noise inside definite frequency band one has to use the time dependent local oscillator phase $`\varphi (t)`$ . In this case the optimal dependence of $`\varphi `$ on $`t`$ is defined by the displacement transducer transfer function $`\xi (\omega _\mathrm{f})`$ and by the spectrum of the external force $`F_\mathrm{s}(\omega )`$ .
So a signal-to-noise ratio is proportional to $`A^2`$ (there is no optimal power) and in principle there is no sensitivity limitation by the SQL. In real experiment when the pump power gets larger the output signal and noises get smaller according to equation (11) if the condition (12) is kept valid therefore when $`A`$ becomes greater than a certain value then the noises of photodetector electronics can limit the sensitivity. However this noises have technical character and will be neglected in the following.
Another sensitivity restriction can arise due to the damping in mechanical oscillator (mirror) . This problem is general for all supersensitive measurements. At the same time an intrinsic dissipation obtained in modern experiments for mechanical oscillator is far larger (by several orders of magnitude) than the value expected from the first principles therefore it can be treated also as a technical problem now and will not be adressed below.
It is worth to mention that the increase in sensitivity over the usual measurement scheme occurs here due to utilization of the internal squeezing (self-squeezing) of the reflected beam because of the nonlinear (quadratic) interaction of the incident light and the mirror . Actually two quadratures of the reflected beam are correlated and it is this fact that allow to use the correlative processing of the output. On the other hand the correlation of the quadratures according to equations (8) means the squeezing of the beam and the larger the correlation coefficient $`\lambda \xi (\omega )A^2`$ the larger the internal squeezing .
## 5 Sensitivity for the pulsed pump
Let consider the periodic envelope $`A(t)`$ which consists of a train of equally spaced pulses with duration $`\tau `$ and period $`T`$. The spectrum of this pump has also the form of a train of pulses in frequency domain with the distance between neighbour pulses
$$\omega _\mathrm{q}=2\pi T^1$$
(13)
For the amplitude of the pump $`A(t)`$ one can use now the expansion into the Fourier series
$$A(t)=\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}g_n\mathrm{exp}(\mathrm{i}n\omega _\mathrm{q}t)$$
(14)
and the particular form of $`A(t)`$ is defined by the set of Fourier amplitudes $`g_n`$.
The response of the displacement transducer now have many frequency components at $`\omega =n\omega _\mathrm{q},n=0,1\mathrm{}`$ according to the equations (4) and each frequency component contains the signal part besides the radiation pressure force $`F_\mathrm{p}(t)`$ have also wide spectrum (cf. (6)). So there are two problems: how to collect the signal parts from the whole spectral band of the output and how to achieve the compensation of the radiation pressure noise in the output. It is clear that the monochromatic local oscillator is inappropriate for the homodyning because quadrature $`b_1(t)`$ of the output signal contains in this case the quadrature $`a_1(t)`$ of the input noises only from one frequency and the radiation pressure force $`F_\mathrm{p}(t)`$ in expression for $`b_2(t)`$ (cf. (7)) contains $`a_1(t)`$ from all frequencies $`n\omega _\mathrm{q}`$ therefore the full compensation is impossible.
Fortunately two problems can be overcome by the use of the pulsed local oscillator with the amplitude time dependence resembling that for the pump.
For the radiation pressure displacement $`X_\mathrm{p}`$ of the mechanical oscillator one has from equations (5), (6) and (14) the following expression
$$X_\mathrm{p}(\omega )=G(\omega )F_\mathrm{p}(\omega )=\lambda G(\omega )\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}g_na_1(\omega n\omega _\mathrm{q})$$
(15)
For the quadrature transformation one can obtain instead of (8) the following equations from (4) and (14)
$`b_1(\omega )`$ $`=`$ $`a_1(\omega )`$
$`b_2(\omega )`$ $`=`$ $`a_2(\omega )+2\omega _\mathrm{p}c^1{\displaystyle \underset{k=\mathrm{}}{\overset{\mathrm{}}{}}}g_k(X_\mathrm{p}(\omega k\omega _\mathrm{q})+X_\mathrm{s}(\omega k\omega _\mathrm{q}))`$ (16)
Let suppose the local oscillator field in the form of
$$E_\mathrm{L}(t)=A_\mathrm{L}(t)\mathrm{cos}(\omega _\mathrm{p}t+\varphi )$$
(17)
where the dependence of the amplitude $`A_\mathrm{L}(t)`$ on $`t`$ is much slower than $`\mathrm{cos}\omega _\mathrm{p}`$t. Then for the envelope of the local oscillator field $`A_\mathrm{L}(t)`$ the Fourier expansion similar to (14) is valid
$$A_\mathrm{L}(t)=\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}e_n\mathrm{exp}(\mathrm{i}n\omega _\mathrm{q}t)$$
(18)
The photodetector current has now the following form
$$I_{\mathrm{pd}}A_\mathrm{L}(t)(b_1(t)\mathrm{cos}\varphi +b_2(t)\mathrm{sin}\varphi )$$
(19)
and in the frequency domain one has
$$I_{\mathrm{pd}}(\omega )\mathrm{cos}\varphi \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}e_nb_1(\omega n\omega _\mathrm{q})+\mathrm{sin}\varphi \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}e_nb_2(\omega n\omega _\mathrm{q})$$
(20)
Let consider different parts in the photodetector output. The first term in equation (20) depends only on the amplitude fluctuations of the input field according to (16)
$$\mathrm{cos}\varphi \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}e_nb_1(\omega n\omega _\mathrm{q})=\mathrm{cos}\varphi \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}e_na_1(\omega n\omega _\mathrm{q})$$
(21)
The second term in equation (20) contains the signal and the noise parts. The noise part consists of the additive noise and the back action noise and has the following expression according to (15) and (16)
$`[\mathrm{sin}\varphi {\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}e_nb_2(\omega n\omega _\mathrm{q})]_{\mathrm{noise}}=\mathrm{sin}\varphi {\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}e_na_2(\omega n\omega _\mathrm{q})+2\omega _\mathrm{p}c^1\mathrm{sin}\varphi \lambda `$
$`{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \underset{k=\mathrm{}}{\overset{\mathrm{}}{}}}e_ng_kG(\omega k\omega _\mathrm{q}n\omega _\mathrm{q})\{{\displaystyle \underset{m=\mathrm{}}{\overset{\mathrm{}}{}}}g_ma_1(\omega k\omega _\mathrm{q}n\omega _\mathrm{q}m\omega _\mathrm{q})\}`$ (22)
Let consider only the photocurrent at small frequencies $`\omega \omega _\mu `$. Then the main input into the photocurrent will be given by the resonant terms for which $`k+n=0`$. With this supposition one has from equation (22)
$`[\mathrm{sin}\varphi {\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}e_nb_2(\omega n\omega _q)]_{\mathrm{noise}}=\mathrm{sin}\varphi {\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}e_na_2(\omega n\omega _\mathrm{q})+`$
$`\mathrm{sin}\varphi \lambda \xi (\omega ){\displaystyle \underset{m=\mathrm{}}{\overset{\mathrm{}}{}}}e_mg_m{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}g_na_1(\omega n\omega _\mathrm{q})\}`$ (23)
Comparing equations (21) and (23) one can conclude that full compensation of back action noise is possible only for
$$e_n=\alpha g_n$$
(24)
where $`\alpha `$ is the same for all numbers $`n`$ so the forms of pump and local oscillator fields have to be the same (apart from the scale factor $`\alpha `$).
Let now consider the signal part of the second term in the r.h.s. of equation (20). From equations (7), (16) and (20) one has
$`[\mathrm{sin}\varphi {\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}e_nb_2(\omega n\omega _\mathrm{q})]_{\mathrm{signal}}=`$
$`\mathrm{sin}\varphi {\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}e_n\{{\displaystyle \underset{k=\mathrm{}}{\overset{\mathrm{}}{}}}g_k\xi (\omega k\omega _\mathrm{q}n\omega _\mathrm{q})F_\mathrm{s}(\omega k\omega _\mathrm{q}n\omega _\mathrm{q})\}`$ (25)
Evaluation of this expression for the condition $`k+n=0`$ gives
$$[\mathrm{sin}\varphi \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}e_nb_2(\omega n\omega _\mathrm{q})]_{\mathrm{signal}}=\mathrm{sin}\varphi \xi (\omega )F_\mathrm{s}(\omega )\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}e_ng_n$$
(26)
Combining equations (20), (23), (24) and (26) and supposing that the back action noise is compensated in the output of the photodetector one can obtain for the spectral density of noises in the photocurrent the following expression
$$N(\omega )\mathrm{sin}\varphi N_0\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}g_ng_n=\mathrm{sin}\varphi N_0P$$
(27)
where it is supposed that fluctuations at frequencies $`\omega n\omega _q,n=0,1\mathrm{}`$ are uncorrelated and have the same spectral density $`N_0`$ (this assumption is valid for not very small duration of pump pulses), $`P`$ is proportional to the time averaged power of the pulsed pump. Then for the signal-to-noise ratio $`\mu `$ one has from equations (26) and (27) the following expression
$$\mu N_0^1P_{\mathrm{}}^{\mathrm{}}\xi (\omega )F_\mathrm{s}(\omega )^2d\omega $$
(28)
This value is just equal to the signal-to-noise ratio for continuous wave pump with a power $`P`$ and correlative processing of the output (cf. equation (11)). Note that the sensitivity here is not limited by the SQL like in the case of correlative processing of quadratures for the monochromatic pump and is increasing with the increase of $`P`$.
It is worth to mention that the condition for the back action noise compensation for the pulsed pump is just the same as for the monochromatic pump (cf. equation (12)) with substitution of the $`A^2`$ with the time averaged value $`P`$. Therefore the compensation of the back action noises for the finite frequency band can be possible for the time varying phase of the local oscillator .
## 6 Conclusion
The pumping of the displacement transducer with a train of the short high-intensity laser pulses is considered. The algorithm of optimal signal processing for such transducer is revealed. It consists of the correlative processing of the output using the pulsed local oscillator with the same envelope as for the pump field (apart from the scale factor). In this case the back action noise due to the radiation pressure force can be fully compensated and the sensitivity of the scheme to a detection of a classical external force is not limited by the SQL (as for the case of correlative quadrature processing and monochromatic pump field).
The pulsed pump can be advantageous over the single frequency pumping when the nonlinear optical elements are used unside the system. Thus considerable increase in sensitivity can be achieved for a gravitational interferometric Fabry-Perot type detector with a nonlinear optical element placed in a waist of the beam . The use of the phase-conjugate mirrors in a gravitational detector of the Michelson type allows to construct the system with the parallel arms . For such systems an efficiency depends on the instant power of the light beam and can be high for the short intensive pulses.
In this article only the problem of the force detection with known spectrum is considered. The reconstruction of unknown external force acting on the displacement transducer with the pulsed pump below the standard quantum limit will be considered elsewhere.
|
warning/0003/astro-ph0003002.html
|
ar5iv
|
text
|
# RBSC-NVSS SAMPLE. I. RADIO AND OPTICAL IDENTIFICATIONS OF A COMPLETE SAMPLE OF 1500 BRIGHT X-RAY SOURCES
## 1 INTRODUCTION
The ROSAT All-Sky Survey Bright Source Catalogue (Voges et al.,, 1999, RBSC revision 1RXS) contains the first large all-sky sample of the brightest X-ray sources, analogous in many respects to the optical NGC catalog. It was derived from the soft (0.1–2.4 keV) X-ray survey performed during the first half year of the ROSAT mission in 1990/91. The catalog sky coverage is 92%, and there are 8,547 sources above its 0.1 counts s<sup>-1</sup> ($`10^{12}`$ ergs s<sup>-1</sup> cm<sup>-2</sup>) completeness limit. Bade et al., (1998) found that about one third of the RBSC sources can be reliably identified with galactic stars, while most of the rest are extragalactic. The extragalactic content of the RBSC comprises a diverse blend of (1) normal spiral galaxies whose X-ray emission is powered by stars and stellar remnants, (2) elliptical galaxies with hot gaseous halos, (3) AGN in Seyfert galaxies, elliptical galaxies, quasars, and BL Lac objects, and (4) clusters of galaxies. The large number of sources in this catalog easily permits statistical analyses of each type of X-ray object. However, the essential properties of these X-ray sources cannot be determined from the X-ray data alone—we need observations in optical and other wavebands to measure their distances, identify their energy sources, etc. Such observations are possible only for those RBSC sources whose optical counterparts have been identified. In this paper we present reliable radio and optical identifications for sources in the RBSC complete sample.
Most RBSC sources have rms positional uncertainties $`10^{\prime \prime }`$ and the sky density of faint optical objects is high, so only the nearest extragalactic X-ray sources can be optically identified by position coincidence alone (cf. Bade et al.,, 1998). Fortunately, most extragalactic RBSC sources are also radio sources in the 1.4 GHz NRAO/VLA Sky Survey (Condon et al., 1998a, , NVSS), whose sky density is low enough for identification with RBSC sources. Since the radio positions are significantly more accurate, the radio sources may be optically identified, yielding optical identifications for the corresponding X-ray sources as well. The NVSS covers the 10.3 sr of sky north of $`\delta =40^{}`$ and contains over $`1.8\times 10^6`$ sources stronger than its 2.5 mJy beam<sup>-1</sup> completeness limit. Since the NVSS was made with relatively low resolution ($`45\mathrm{}`$ FWHM), it does not discriminate against moderately extended radio sources in nearby galaxies and clusters. Its rms positional uncertainties range from $`<1\mathrm{}`$ for the $`N4\times 10^5`$ sources stronger than 15 mJy to 7$`\mathrm{}`$ for the faintest (S = 2.3 mJy) detectable sources, allowing us to make optical identifications with objects as faint as R $`21`$.
We optically identified the RBSC-NVSS sources with objects in the United States Naval Observatory catalog A2.0 (Monet et al.,, 1998, USNO). The USNO catalog contains 526,280,881 objects detected by the Precision Measuring Machine on the Palomar Optical Sky Survey I (POSS-I) blue O and red E plates, the UK Science Research Council SRC-J survey plates, and the European Southern Observatory ESO-R survey plates. The catalog was compiled from the blue/red overlaps (within 2$`\mathrm{}`$) of the detection lists generated from scans of POSS-I O and E plates centered on $`\delta >18^{}`$ and SRC-J and ESO-R plates centered on $`\delta <20^{}`$. The stated astrometric and photometric errors are about $`0\stackrel{}{\mathrm{.}}25`$ and 0.5 mag rms, respectively. The USNO catalog covers the entire sky and probes as deep as B=21 (O plates), R=20 (E plates), J=22, and F=21 for objects with appropriate colors.
Section 2 explains our method for making the identifications and assessing their reliabilities. The results are presented in Section 3.
## 2 Cross-identifications
In this section we present a sequence of increasingly powerful methods for making cross-identifications and directly evaluating the their individual reliabilities: (1) The simplest case is identification by position-coincidence between two wavebands (Sec. 2.1), associating NVSS radio sources with RBSC X-ray sources, for example. For each candidate we derive the probability that it is the correct identification (the reliability of that identification). (2) If the positions are not sufficiently accurate to guarantee reliable identifications by themselves, they may be supplemented by additional data (Sec. 2.2). To continue our example, RBSC X-ray sources have a flatter radio flux-density distribution than NVSS background sources, so radio flux densities affect identification reliabilities. The identification reliabilities derived directly in this section are similar to those obtained via the “likelihood ratio” method by Sutherland & Saunders (1992). (3) Even with the aid of additional data, the RBSC position errors are too large for making reliable X-ray identifications with faint optical objects. We show how accurate positions of NVSS radio sources associated with RBSC sources can be used to select the correct optical counterparts. The reliabilities of such multiwavelength linked position-coincidence identifications are derived in (Sec. 2.3). (4) Finally, linked cross-identifications can themselves be strengthened by applying additional constraints (e.g., optical magnitudes), and their reliabilities are obtained in Section 2.4.
### 2.1 Direct Position-coincidence Identifications
All identification programs begin with a set of identification candidates in some search area surrounding each source to be identified. This area should encompass all plausible candidates, but its exact size and shape are not critical. A larger than necessary search area increases the number of candidates which must be evaluated but does not significantly degrade identification reliability. We used search circles of 180$`\mathrm{}`$ in radius (Figure RBSC-NVSS SAMPLE. I. RADIO AND OPTICAL IDENTIFICATIONS OF A COMPLETE SAMPLE OF 1500 BRIGHT X-RAY SOURCES) to identify RBSC sources having rms positional uncertainties $`\sigma _\mathrm{x}10^{\prime \prime }20^{\prime \prime }`$. The search radius is much larger than the $`3\sigma `$ X-ray error circle to permit identification with extended and asymmetric radio sources. Each search area contains some numbers $`m0`$ of NVSS radio sources and $`l0`$ of optical objects from the USNO catalog.
We consider radio identifications of X-ray sources using only the X-ray and radio positions. Let $`P(R)`$ ($`R=1,\mathrm{},m`$) represent the probability that the $`R^{th}`$ candidate is the correct radio identification of an X-ray source and $`P(0)`$ the probability that none is. Since there is only a negligible chance that the radio identification is exterior to the search area, $`P(0)`$ is just the probability that the actual identification is too faint to be recognized as a candidate. For example, we cannot identify radio counterparts fainter than the 2.5 mJy beam<sup>-1</sup> NVSS catalog limit. Such an identification is often called an “empty field.” If we make the astronomical assumption that not more than one of the $`m`$ candidates is the correct identification, then the sum of these mutually exclusive probabilities is unity:
$$\underset{R=0}{\overset{m}{}}P(R)=1.$$
(1)
The $`R^{th}`$ candidate is the correct identification if (1) there exists a detectable identification, (2) that identification lies in the infinitesimal area element $`dA`$ containing the position of the $`R^{th}`$ candidate, and (3) the ($`m1`$) remaining candidates are unrelated sources which happen to lie in the areas $`dA`$ surrounding their positions (see Figure RBSC-NVSS SAMPLE. I. RADIO AND OPTICAL IDENTIFICATIONS OF A COMPLETE SAMPLE OF 1500 BRIGHT X-RAY SOURCES). The probability that all three independent events occur is the product of their individual probabilities, which we now evaluate.
(1) The a priori probability that there is a detectable radio identification is equal to the initially unknown fraction $`f_\mathrm{r}`$ of X-ray sources in the sample that actually have detectable radio identifications. We estimated $`f_\mathrm{r}`$ by guessing an initial value, making trial identifications for the whole source sample, replacing the initial value by the observed value, and iterating. If this procedure fails to converge rapidly, the resulting identifications are probably too unreliable to be useful. To estimate the observed value of $`f_\mathrm{r}`$ we summed the computed reliabilities of all the RBSC-NVSS identifications and divided that sum by the total number of X-ray sources. Since only a small percentage of bright galactic X-ray stars are radio sources, we calculated separate $`f_\mathrm{r}`$ values for fields containing optically bright ($`m<12`$) stars and for all other fields.
(2) The probability of finding the correct identification in the infinitesimal area $`dA`$ offset by the angle $`\varphi _{\mathrm{xr}}(R)`$ between the position of the X-ray source and the $`R^{th}`$ candidate is $`p[\varphi _{\mathrm{xr}}(R)]dA`$, where $`p[\varphi _{\mathrm{xr}}(R)]`$ is just the normalized error distribution of the measured X-ray/radio offsets. These errors include both the radio and X-ray measurement errors and may be augmented by a contribution allowing for possible astronomical offsets of extended sources. For example, the centroid of a head-tail radio source will not coincide with its parent galaxy.
(3) If the $`R^{th}`$ candidate is the correct identification, then the remaining $`(m1)`$ candidates must be unrelated sources lying in areas $`dA`$ containing their positions. The probability of finding each unrelated source is $`\rho _\mathrm{r}dA`$, where $`\rho _\mathrm{r}`$ is the mean sky density of unrelated radio candidates. For NVSS candidates,
$$\rho _\mathrm{r}=_{S_{\mathrm{min}}}^{\mathrm{}}n(S)𝑑S,$$
(2)
where $`n(S)`$ is the differential source count at 1.4 GHz and $`S_{\mathrm{min}}=2.5`$ mJy is the minimum flux density of the candidates. Possible clustering of candidates around the true identification could be addressed by increasing $`\rho _\mathrm{r}`$ from its global to local value. If there is no detectable identification, all $`m`$ candidates must be unrelated sources.
The resulting set of $`(m+1)`$ proportionalities
$`P(R)f_\mathrm{r}p[\varphi _{\mathrm{xr}}(R)]dA(\rho _\mathrm{r}dA)^{m1}\mathrm{if}R>0`$ (3a)
$`P(0)(1f_\mathrm{r})(\rho _\mathrm{r}dA)^m`$ (3b)
normalized by Equation 1 specify the $`(m+1)`$ identification reliabilities:
$`P(R)=cf_\mathrm{r}{\displaystyle \frac{p[\varphi _{\mathrm{xr}}(R)]}{\rho _\mathrm{r}}}\mathrm{if}R>0`$ (4a)
$`P(0)=c(1f_\mathrm{r}),`$ (4b)
where
$`c^1={\displaystyle \frac{f_\mathrm{r}}{\rho _\mathrm{r}}}{\displaystyle \underset{j=1}{\overset{m}{}}}p[\varphi _{\mathrm{xr}}(j)]+(1f_\mathrm{r}).`$ (4c)
Equation 4 gives the probability $`P(R)`$ that the $`R^{th}`$ of $`m`$ candidates is the correct identification and the probability $`P(0)`$ that none is. These probabilities depend on the search area only through the number $`m`$ of candidates contributing to the sum in Equation 4c. Making the search area “too big” by definition only adds candidates having negligible $`p(\varphi _{\mathrm{xr}})`$ and hence little effect in Equation 4c.
In many cases, the positional error distributions of both the source and its identification candidates are nearly circular Gaussians. \[See Condon et al., (1995) for the more complicated case of elliptical Gaussian error distributions with arbitrary orientations.\] If the X-ray source and radio candidate positions have rms uncertainties $`\sigma _\mathrm{x}`$ and $`\sigma _\mathrm{r}`$ in each coordinate, then
$$p(\varphi _{\mathrm{xr}})=\frac{1}{2\pi \sigma ^2}\mathrm{exp}\left(\frac{\varphi _{\mathrm{xr}}^2}{2\sigma ^2}\right),$$
(5)
where $`\sigma ^2\sigma _\mathrm{x}^2+\sigma _\mathrm{r}^2`$. The variance of $`\varphi _{\mathrm{xr}}`$ is $`\varphi _{\mathrm{xr}}^2=2\sigma ^2`$ so a typical identification has $`p(\varphi _{\mathrm{xr}})(2\pi \sigma ^2e)^1`$ and is highly reliable if both
$`\sigma ^2{\displaystyle \frac{f_\mathrm{r}}{2\pi e\rho _\mathrm{r}(1f_\mathrm{r})}}`$ (6a)
and
$`\sigma ^2{\displaystyle \frac{1}{2\pi e\rho _\mathrm{r}}}`$ (6b)
Equation 6a ensures that true empty fields are not misidentified with background sources, a danger if the actual identification rate is very low ($`f_\mathrm{r}1`$). Equation 6b ensures that there is no difficulty choosing the correct identification from multiple candidates. For NVSS identifications of extragalactic RBSC sources, $`f_\mathrm{r}0.61`$ and the sky density of sources stronger than the NVSS catalog limit $`S_{\mathrm{min}}2.5`$ mJy is $`\rho _\mathrm{r}1.76\times 10^5`$ sr<sup>-1</sup>, leading to the fairly weak requirement $`\sigma 10^2`$ arcsec which is easily satisfied by the RBSC positional uncertainties. The sky density of optical objects in the USNO catalog is two orders of magnitude higher, and the RBSC positional uncertainties are too large to satisfy the requirement $`\sigma 10`$ arcsec for making optical identifications by position coincidence alone, even though most RBSC sources do have optical counterparts in the USNO catalog.
### 2.2 Pairwise Identifications With Additional Constraints
Uncertain identifications based on positional coincidence alone may be strengthened or rejected by non-positional data. For example, the flux-density distribution of radio sources identified with extragalactic RBSC sources peaks well above the NVSS sensitivity limit. Figure RBSC-NVSS SAMPLE. I. RADIO AND OPTICAL IDENTIFICATIONS OF A COMPLETE SAMPLE OF 1500 BRIGHT X-RAY SOURCES shows the logarithm of the ratio of the probability $`p[S|x]`$ that a RBSC-NVSS source has flux density $`S`$ to the probability $`p[S|\overline{x}]`$ that an unrelated NVSS sources has flux density $`S`$. Most unrelated radio sources are fainter than the correct identifications (indicated in Figure RBSC-NVSS SAMPLE. I. RADIO AND OPTICAL IDENTIFICATIONS OF A COMPLETE SAMPLE OF 1500 BRIGHT X-RAY SOURCES by the positive logarithm of the ratio at all but the faintest radio flux levels), so the stronger of two radio candidates with similar $`p(\varphi _{\mathrm{xr}})`$ is the more likely identification. The reliabilities of such identifications can be calculated through the use of “likelihood ratios” (see Sutherland & Saunders,, 1992); here we obtain the identification reliabilities directly by extending the derivation of Equation 4.
Let $`p[\varphi _{\mathrm{xr}}(R),S(R)]dAdS`$ be the probability that the $`R^{th}`$ candidate is the identification lying in the area $`dA`$ surrounding its observed position and in the flux-density range $`dS`$ containing its observed flux density $`S(R)`$. The multiplicative law of probabilities states that this probability is the product of $`p[\varphi _{\mathrm{xr}}(R)]dA`$ and $`p[S(R)|x]dS`$, where $`p(S|x)`$ is the normalized ($`_{S_{\mathrm{min}}}^{\mathrm{}}p(S|x)𝑑S=1`$) flux-density distribution of radio sources identified with X-ray sources. Like $`f_\mathrm{r}`$, $`p(S|x)`$ is best estimated from the actual identification data by iteration. The probability that the $`R^{th}`$ candidate is an unrelated (X-ray quiet) radio source with flux density $`S`$ lying in area $`dA`$ is $`p[S(R)|\overline{x}]\rho _\mathrm{r}dAdS`$, where $`p(S|\overline{x})`$ is the flux-density distribution of all (background) NVSS sources. Thus Equations 3 and 4 can be replaced by:
$`P(R)f_\mathrm{r}p[\varphi _{\mathrm{xr}}(R)]{\displaystyle \frac{p[S(R)|x]}{p[S(R)|\overline{x}]}}dAdS\left\{{\displaystyle \underset{j=1}{\overset{m}{}}}p[S(j)|\overline{x}]\right\}(\rho _\mathrm{r}dAdS)^{m1}\mathrm{if}R>0`$ (7a)
$`P(0)(1f_\mathrm{r})\{{\displaystyle \underset{j=1}{\overset{m}{}}}p[S(j)]|\overline{x}]\}(\rho _\mathrm{r}dAdS)^m`$ (7b)
and
$`P(R)=cf_\mathrm{r}{\displaystyle \frac{p[\varphi _{\mathrm{xr}}(R)]}{\rho _\mathrm{r}}}{\displaystyle \frac{p[S(R)|x]}{p[S(R)|\overline{x}]}}\mathrm{if}R>0`$ (8a)
$`P(0)=c(1f_\mathrm{r})`$ (8b)
where
$`c^1={\displaystyle \frac{f_\mathrm{r}}{\rho _\mathrm{r}}}{\displaystyle \underset{j=1}{\overset{m}{}}}p[\varphi _{\mathrm{xr}}(j)]{\displaystyle \frac{p[S(j)|x]}{p[S(j)|\overline{x}]}}+(1f_\mathrm{r})`$ (8c)
Equation 8 gives the identification probabilities based on candidate flux densities as well as positional coincidence. It could be further extended to include additional continuous (e.g., spectral index) or discrete (e.g., morphological type) parameters which might prove useful for distinguishing between correct identifications and unrelated candidates.
### 2.3 Linked Position-coincidence Cross-identifications
The positional uncertainties of most RBSC sources are too large to yield reliable optical identifications with faint galaxies and quasars directly. However, the NVSS sources reliably identified with RBSC sources have sufficiently accurate radio positions that nearly all can be optically identified by position coincidence alone. The reliabilities of such linked X-ray/radio/optical position-coincidence identifications are derived below.
Let $`P(R,V)`$ be the probability that the $`R^{th}`$ radio source and the $`V^{th}`$ optically visible object are the correct identifications of an X-ray source, where $`R=0,1,\mathrm{},m`$ and $`V=0,1,\mathrm{},l`$. The values $`R=0`$ and $`V=0`$ correspond to radio and optical “empty fields,” respectively. The probabilities of these $`(m+1)(l+1)`$ mutually exclusive possible outcomes must add up to one:
$$\underset{R=0}{\overset{m}{}}\underset{V=0}{\overset{l}{}}P(R,V)=1.$$
(9)
Denote the fractions of X-ray sources in the sample having detectable ($`R>0`$) radio and optical ($`V>0`$) identifications by $`f_r`$ and $`f_v`$, respectively. Let $`p[\varphi _{\mathrm{xr}}(R)]`$ and $`p[(\varphi _{\mathrm{xv}}(V)]`$ be the probability densities of the X-ray/radio and X-ray/optical positional offsets $`\varphi `$ if the $`R^{th}`$ radio source and the $`V^{th}`$ optical object are the correct identifications. The sky densities of background radio and optical candidates are $`\rho _\mathrm{r}`$ and $`\rho _\mathrm{v}`$, respectively. Using the multiplication law $`P(R,V)=P(R)P(V|R)`$ yields:
$`P(R,V)f_\mathrm{r}p[\varphi _{\mathrm{xr}}(R)]dA(\rho _\mathrm{r}dA)^{m1}f_\mathrm{v}p[\varphi _{\mathrm{xv}}(V)]dA(\rho _\mathrm{v}dA)^{l1}\mathrm{if}R>0,V>0`$ (10a)
$`P(R,0)f_\mathrm{r}p[\varphi _{\mathrm{xr}}(R)]dA(\rho _\mathrm{r}dA)^{m1}(1f_\mathrm{v})(\rho _\mathrm{v}dA)^l\mathrm{if}R>0`$ (10b)
$`P(0,V)(1f_\mathrm{r})(\rho _\mathrm{r}dA)^mf_\mathrm{v}p[\varphi _{\mathrm{xv}}(V)]dA(\rho _\mathrm{v}dA)^{l1}\mathrm{if}V>0`$ (10c)
$`P(0,0)(1f_\mathrm{r})(\rho _\mathrm{r}dA)^m(1f_\mathrm{v})(\rho _\mathrm{v}dA)^l`$ (10d)
The normalization Equation 9 then implies:
$`P(R,V)=cf_\mathrm{r}f_\mathrm{v}{\displaystyle \frac{p[\varphi _{\mathrm{xr}}(R)]}{\rho _\mathrm{r}}}{\displaystyle \frac{p[\varphi _{\mathrm{xv}}(V)]}{\rho _\mathrm{v}}}\mathrm{if}R>0,V>0`$ (11a)
$`P(R,0)=cf_\mathrm{r}(1f_\mathrm{v}){\displaystyle \frac{p[\varphi _{\mathrm{xr}}(R)]}{\rho _\mathrm{r}}}\mathrm{if}R>0`$ (11b)
$`P(0,V)=c(1f_\mathrm{r})f_\mathrm{v}{\displaystyle \frac{p[\varphi _{\mathrm{xv}}(V)]}{\rho _\mathrm{v}}}\mathrm{if}V>0`$ (11c)
$`P(0,0)=c(1f_\mathrm{r})(1f_\mathrm{v}),`$ (11d)
where
$`c^1={\displaystyle \frac{f_\mathrm{r}f_\mathrm{v}}{\rho _\mathrm{r}\rho _\mathrm{v}}}{\displaystyle \underset{j=1}{\overset{m}{}}}p[\varphi _{\mathrm{xr}}(j)]{\displaystyle \underset{i=1}{\overset{l}{}}}p[\varphi _{\mathrm{xv}}(i)]`$
$`+{\displaystyle \frac{f_\mathrm{r}(1f_\mathrm{v})}{\rho _\mathrm{r}}}{\displaystyle \underset{j=1}{\overset{m}{}}}p[\varphi _{\mathrm{xr}}(j)]+{\displaystyle \frac{(1f_\mathrm{r})f_\mathrm{v}}{\rho _\mathrm{v}}}{\displaystyle \underset{i=1}{\overset{l}{}}}p[\varphi _{\mathrm{xv}}(i)]+(1f_\mathrm{r})(1f_\mathrm{v}).`$ (11e)
Unfortunately, applying Equation 11 does not yield good optical identifications of RBSC X-ray sources because both the X-ray positional uncertainties and the mean density $`\rho _\mathrm{v}`$ of optical candidates are large. Reliable radio identifications of X-ray sources can be made because the density $`\rho _\mathrm{r}`$ of radio candidates is much smaller, and reliable optical identifications of NVSS radio sources are possible because the NVSS positions are more accurate. If the astronomical assumption is made that the optical identifications of these radio identifications are also the optical identifications of the corresponding X-ray sources, then
$`P(R,V)f_\mathrm{r}p[\varphi _{\mathrm{xr}}(R)]dA(\rho _\mathrm{r}dA)^{m1}f_\mathrm{v}p[\varphi _{\mathrm{xv}}(V)]dA(\rho _\mathrm{v}dA)^{l1}p[\varphi _{\mathrm{rv}}(R,V)]dA(\rho _\mathrm{v}dA)^{l1}`$
$`\mathrm{if}R>0,V>0`$ (12a)
$`P(R,0)f_\mathrm{r}p[\varphi _{\mathrm{xr}}(R)]dA(\rho _\mathrm{r}dA)^{m1}(1f_\mathrm{v})(\rho _\mathrm{v}dA)^l(\rho _\mathrm{v}dA)^l\mathrm{if}R>0`$ (12b)
$`P(0,V)(1f_\mathrm{r})(\rho _\mathrm{r}dA)^mf_\mathrm{v}p[\varphi _{\mathrm{xv}}(V)]dA(\rho _\mathrm{v}dA)^{l1}(\rho _\mathrm{v}dA)^l\mathrm{if}V>0`$ (12c)
$`P(0,0)(1f_\mathrm{r})(\rho _\mathrm{r}dA)^m(1f_\mathrm{v})(\rho _\mathrm{v}dA)^l(\rho _\mathrm{v}dA)^l,`$ (12d)
where $`p[\varphi _{\mathrm{rv}}(R,V)]`$ is the probability distribution of offsets between the $`R^{th}`$ radio source and the $`V^{th}`$ optical object. The normalization Equation 9 implies
$`P(R,V)=cf_\mathrm{r}f_\mathrm{v}{\displaystyle \frac{p[\varphi _{\mathrm{xr}}(R)]}{\rho _\mathrm{r}}}{\displaystyle \frac{p[\varphi _{\mathrm{xv}}(V)]}{\rho _\mathrm{v}}}{\displaystyle \frac{p[\varphi _{\mathrm{rv}}(R,V)]}{\rho _\mathrm{v}}}\mathrm{if}R>0,V>0`$ (13a)
$`P(R,0)=cf_\mathrm{r}(1f_\mathrm{v}){\displaystyle \frac{p[\varphi _{\mathrm{xr}}(R)]}{\rho _\mathrm{r}}}\mathrm{if}R>0`$ (13b)
$`P(0,V)=c(1f_\mathrm{r})f_\mathrm{v}{\displaystyle \frac{p[\varphi _{\mathrm{xv}}(V)]}{\rho _\mathrm{v}}}\mathrm{if}V>0`$ (13c)
$`P(0,0)=c(1f_\mathrm{r})(1f_\mathrm{v}),`$ (13d)
where
$`c^1={\displaystyle \frac{f_\mathrm{r}f_\mathrm{v}}{\rho _\mathrm{r}\rho _\mathrm{v}^2}}{\displaystyle \underset{j=1}{\overset{m}{}}}p[\varphi _{\mathrm{xr}}(j)]{\displaystyle \underset{i=1}{\overset{l}{}}}p[\varphi _{\mathrm{xv}}(i)]p[\varphi _{\mathrm{rv}}(j,i)]`$
$`+{\displaystyle \frac{f_\mathrm{r}(1f_\mathrm{v})}{\rho _\mathrm{r}}}{\displaystyle \underset{j=1}{\overset{m}{}}}p[\varphi _{\mathrm{xr}}(j)]+{\displaystyle \frac{(1f_\mathrm{r})f_\mathrm{v}}{\rho _\mathrm{v}}}{\displaystyle \underset{i=1}{\overset{l}{}}}p[\varphi _{\mathrm{xv}}(i)]+(1f_\mathrm{r})(1f_\mathrm{v}).`$ (13e)
Equation 13 specifies the reliabilities of linked X-ray/radio/optical identifications made on the basis of position-coincidence alone.
### 2.4 Linked Cross-identifications With Additional Constraints
Finally, differences between the magnitude distributions of the optical counterparts to X-ray sources and background optical objects can be used to improve the reliability of the optical identifications. Figure RBSC-NVSS SAMPLE. I. RADIO AND OPTICAL IDENTIFICATIONS OF A COMPLETE SAMPLE OF 1500 BRIGHT X-RAY SOURCES gives the logarithm of the ratio of the probability $`p[\mu |x]`$ that the USNO identification of a RBSC-NVSS source has magnitude $`\mu `$ to the probability $`p[\mu |\overline{x}]`$ that an unrelated USNO object has magnitude $`\mu `$ as a function of B magnitude, and shows, in contrast to the radio, that there is very little difference between the two populations for B $`>13`$, but the ratio does become large for objects brighter than this limit. Let $`\mu (V)`$ denote the magnitude of the $`V^{th}`$ optical candidate. Let $`p(\mu |x)`$ and $`p(\mu |\overline{x})`$ be the normalized magnitude distributions of optical objects in the USNO catalog which are X-ray sources and unrelated optical objects, respectively. Let $`p(S|x)`$ be the normalized flux-density distribution of radio-detected X-ray sources and $`p(S|\overline{x})`$ be the the normalized flux-density distribution of unrelated radio sources stronger than the NVSS limit, 2.5 mJy. Then
$`P(R,V)f_\mathrm{r}p[\varphi _{\mathrm{xr}}(R)]{\displaystyle \frac{p[S(R)|x]}{p[S(R)|\overline{x}]}}dAdS\left\{{\displaystyle \underset{j=1}{\overset{m}{}}}p[S(j)|\overline{x}]\right\}(\rho _\mathrm{r}dAdS)^{m1}`$
$`\times f_\mathrm{v}p[\varphi _{\mathrm{xv}}(V)]{\displaystyle \frac{p[\mu (V)|x]}{p[\mu (V)|\overline{x}]}}dAd\mu \left\{{\displaystyle \underset{i=1}{\overset{l}{}}}p[\mu (i)|\overline{x}]\right\}(\rho _\mathrm{v}dAd\mu )^{l1}`$
$`\times p[\varphi _{\mathrm{rv}}(R,V)]dA(\rho _\mathrm{v}dA)^{l1}\mathrm{if}R>0,V>0`$ (14a)
$`P(R,0)f_\mathrm{r}p[\varphi _{\mathrm{xr}}(R)]{\displaystyle \frac{p[S(R)|x]}{p[S(R)|\overline{x}]}}dAdS\left\{{\displaystyle \underset{j=1}{\overset{m}{}}}p[S(j)|\overline{x}]\right\}(\rho _\mathrm{r}dAdS)^{m1}`$
$`\times (1f_\mathrm{v})\left\{{\displaystyle \underset{i=1}{\overset{l}{}}}p[\mu (i)|\overline{x}]\right\}(\rho _\mathrm{v}dAd\mu )^l(\rho __\mathrm{v}dA)^l\mathrm{if}R>0`$ (14b)
$`P(0,V)(1f_\mathrm{r})\left\{{\displaystyle \underset{j=1}{\overset{m}{}}}p[S(j)|\overline{x}]\right\}(\rho _\mathrm{r}dAdS)^m`$
$`\times f_\mathrm{v}p[\varphi _{\mathrm{xv}}(V)]{\displaystyle \frac{p[\mu (V)|x]}{p[\mu (V)|\overline{x}]}}dAd\mu \left\{{\displaystyle \underset{i=1}{\overset{l}{}}}p[\mu (i)|\overline{x}]\right\}(\rho _\mathrm{v}dAd\mu )^{l1}(\rho __\mathrm{v}dA)^l\mathrm{if}V>0`$ (14c)
$`P(0,0)(1f_\mathrm{r})\left\{{\displaystyle \underset{j=1}{\overset{m}{}}}p[S(j)|\overline{x}]\right\}(\rho _\mathrm{r}dAdS)^m`$
$`\times (1f_\mathrm{v})\left\{{\displaystyle \underset{i=1}{\overset{l}{}}}p[\mu (i)|\overline{x}]\right\}(\rho _\mathrm{v}dAd\mu )^l(\rho __\mathrm{v}dA)^l.`$ (14d)
The normalized probabilities are:
$`P(R,V)=cf_\mathrm{r}f_\mathrm{v}{\displaystyle \frac{p[\varphi _{\mathrm{xr}}(R)]}{\rho _\mathrm{r}}}{\displaystyle \frac{p[S(R)|x]}{p[S(R)|\overline{x}]}}{\displaystyle \frac{p[\varphi _{\mathrm{xv}}(V)]}{\rho _\mathrm{v}}}{\displaystyle \frac{p[\mu (V)|x]}{p[\mu (V)|\overline{x}]}}{\displaystyle \frac{p[\varphi _{\mathrm{rv}}(R,V)]}{\rho _\mathrm{v}}}`$
$`\mathrm{if}R>0,V>0`$ (15a)
$`P(R,0)=cf_\mathrm{r}(1f_\mathrm{v}){\displaystyle \frac{p[\varphi _{\mathrm{xr}}(R)]}{\rho _\mathrm{r}}}{\displaystyle \frac{p[S(R)|x]}{p[S(R)|\overline{x}]}}\mathrm{if}R>0`$ (15b)
$`P(0,V)=c(1f_\mathrm{r})f_\mathrm{v}{\displaystyle \frac{p[\varphi _{\mathrm{xv}}(V)]}{\rho _\mathrm{v}}}{\displaystyle \frac{p[\mu (V)|x]}{p[\mu (V)|\overline{x}]}}\mathrm{if}V>0`$ (15c)
$`P(0,0)=c(1f_\mathrm{r})(1f_\mathrm{v}),`$ (15d)
where
$`c^1=f_\mathrm{r}f_\mathrm{v}{\displaystyle \underset{j=1}{\overset{m}{}}}{\displaystyle \frac{p[\varphi _{\mathrm{xr}}(j)]}{\rho _\mathrm{r}}}{\displaystyle \frac{p[S(j)|x]}{p[S(j)|\overline{x}]}}{\displaystyle \underset{i=1}{\overset{l}{}}}{\displaystyle \frac{p[\varphi _{\mathrm{xv}}(i)]}{\rho _\mathrm{v}}}{\displaystyle \frac{p[\mu (i)|x]}{p[\mu (i)|\overline{x}]}}{\displaystyle \frac{p[\varphi _{\mathrm{rv}}(j,i)]}{\rho _\mathrm{v}}}`$
$`+f_\mathrm{r}(1f_\mathrm{v}){\displaystyle \underset{j=1}{\overset{m}{}}}{\displaystyle \frac{p[\varphi _{\mathrm{xr}}(j)]}{\rho _\mathrm{r}}}{\displaystyle \frac{p[S(j)|x]}{p[S(j)|\overline{x}]}}+(1f_\mathrm{r})f_\mathrm{v}{\displaystyle \underset{i=1}{\overset{l}{}}}{\displaystyle \frac{p[\varphi _{\mathrm{xv}}(i)]}{\rho _\mathrm{v}}}{\displaystyle \frac{p[\mu (i)|x]}{p[\mu (i)|\overline{x}]}}`$
$`+(1f_\mathrm{r})(1f_\mathrm{v}).`$ (15e)
It is with Equation 15 that we evaluated the reliabilities of the RBSC identification candidates found in Section 3.
The final identification process proceeded as follows. We started with a list of X-ray positions and errors from the RBSC. We then searched for corresponding NVSS radio and USNO optical sources within a $`3\mathrm{}`$ radius ($`R=0,\mathrm{},m`$ and $`V=0,\mathrm{},l`$ respectively). This particular search size was used because it was much larger than the 3$`\sigma `$ X-ray error circle and large enough to contain all but the most extended radio and optical sources. The positions of these sources (X-ray, radio, and optical), their positional errors, and their fluxes and magnitudes were combined with the a priori probabilities of detecting radio and optical counterparts, $`f_\mathrm{r}`$ and $`f_\mathrm{v}`$, the RBSC radio flux and optical magnitude distributions using Equation 15 to estimate the probabilities, $`P(R,V)`$. We estimated $`f_\mathrm{r}`$ and $`f_\mathrm{v}`$ initially , and then iteratively replaced these values with the actual identification rates until they converged to the values of 0.61 and 0.99. The RBSC flux and magnitude distributions were again found iteratively. The radio flux distribution of RBSC sources was compared to the distribution of background NVSS objects and the RBSC optical magnitude distribution was compared to the position dependent optical background distribution given by Bahcall & Soneira, (1980).
### 2.5 Additional Parameters and Possible Concerns
Many astronomical objects have extended X-ray, radio, or optical structures which lead to offsets between wavebands beyond those accounted for by measurement errors.
To allow for an offset between the centroid and central component of a radio source smaller than the NVSS beam, we added one tenth of the deconvolved NVSS size or upper limit in quadrature to the rms uncertainty of the NVSS centroid position. Further, if this RBSC-NVSS matched field happened to lie within the area covered by the VLA FIRST Survey (White et al.,, 1997, $`5\mathrm{}`$ resolution, 1 mJy flux limit), we then evaluated Equation 15, and replaced the NVSS sources with the FIRST sources if the reliability was better.
An attempt was also made to identify sources resolved into two or more NVSS components, such as double-lobed or head-tail radio galaxies. If two radio components in the search area were within $`3\mathrm{}`$ of one another and had a flux ratio $`<3`$, we initiated a test for a double source. This consisted of adding an artificial “source” with the combined flux of the two components and located at their radio centroid (this gave much better results than the optimally weighted centroid used in Windhorst et al. 1984) to the candidate list. We then re-evaluated Equation 15, keeping this new “source” if the reliability was greater than 0.8 or reverting back to the original reliability estimation otherwise. The rms major- and minor-axis position errors for these artificial sources were taken to be $`\frac{1}{10}`$ and $`\frac{1}{30}`$ the size of the component separation. We also found a few true double sources with flux ratios $`>3`$ which were evaluated in the same manner.
We have not corrected for extended RBSC sources. This has the effect of selecting against large scale X-ray emission not coincident with optical and radio sources, such as clusters of galaxies.
Since the USNO catalog does not list optical angular sizes, we used the AIPS Gaussian fitting task IMFIT to measure the optical diameters of of objects larger than $`0\stackrel{}{\mathrm{.}}5`$. An additional error of $`\frac{1}{15}`$ of the major axis was added in quadrature with the rms optical positional uncertainty (which we took to be $`0\stackrel{}{\mathrm{.}}5`$).
The USNO $`O`$ plate magnitudes are converted into Johnson B magnitudes using $`B=O0.119(OE)`$ (Evans,, 1988). We found that these USNO magnitudes are systematically too bright for galaxies with B $`16`$ by as much as three magnitudes in both R and B as shown by a comparison of USNO magnitudes with B magnitudes from the NASA/IPAC Extragalactic Database (NED,, 1999) (Figure RBSC-NVSS SAMPLE. I. RADIO AND OPTICAL IDENTIFICATIONS OF A COMPLETE SAMPLE OF 1500 BRIGHT X-RAY SOURCES). There is better agreement (the dispersion is $`2`$ mag) between USNO and NED magnitudes for galaxies with B $`16`$. To avoid overestimating the identification reliabilities of objects brighter than B $`16`$, we corrected the USNO blue magnitudes using the polynominal fit
$`\mathrm{B}_{\mathrm{NED}}=0.066+1.172\mathrm{B}_{\mathrm{USNO}}+0.120(\mathrm{B}_{\mathrm{USNO}})^20.014(\mathrm{B}_{\mathrm{USNO}})^3+0.000376(\mathrm{B}_{\mathrm{USNO}})^4`$
indicated by the solid curve in Figure RBSC-NVSS SAMPLE. I. RADIO AND OPTICAL IDENTIFICATIONS OF A COMPLETE SAMPLE OF 1500 BRIGHT X-RAY SOURCES. The remaining scatter ($`\pm `$ 2 mag) is not large enough to impact the reliability of our identifications since the ratio $`p(\mu |x)/p(\mu |\overline{x})`$ does not change by more than a factor of a few over the affected magnitude range (Fig. RBSC-NVSS SAMPLE. I. RADIO AND OPTICAL IDENTIFICATIONS OF A COMPLETE SAMPLE OF 1500 BRIGHT X-RAY SOURCES).
One problem associated with catalogued flux densities and magnitudes is Malmquist bias. Given some error in the flux and magnitude determination and the steep slope of the counts, the true brightness of an object will be overestimated. While this is a potential pitfall for faint flux and magnitude counts, it should not significantly affect the reliabilities calculated here for several reasons: 1) the flux and magnitude estimates are normalized by the same Malmquist biased background, 2) for the optical sources, the flat magnitude distribution (Fig. RBSC-NVSS SAMPLE. I. RADIO AND OPTICAL IDENTIFICATIONS OF A COMPLETE SAMPLE OF 1500 BRIGHT X-RAY SOURCES) carries little weight in the reliability determination anyway, and 3) the slope of the optical (and to a lesser extent the radio) RBSC source counts turns over well in advance of the survey limit, such that there are very few faint sources where the bias would be strongest.
A further complication is the high sky density of bright foreground stars at low Galactic latitudes. The likelihood that an optical star is an X-ray emitter is quite large (one third of RBSC sources are stellar in origin; (Bade et al.,, 1998)), while the likelihood that the same optical star is also a radio emitter is quite small ($`0.12\%`$, see Condon et al., (1997)). Thus a bright X-ray star near a background radio source may result in the identification of the X-ray star with that radio source. Because the density of background optical sources increases more rapidly toward fainter magnitudes than the density of RBSC-NVSS optical identifications does, a bright optical star will be assigned a greater identification reliability than a faint galaxy, even if the star is not quite as close to the radio position as the galaxy is. Weeding out X-ray objects which are likely to be stars will then decrease the chances of calling an interloper the correct identification. Thus we classified all objects brighter than B $`<13`$ mag on the basis of appearance on the Digitized Sky Survey (Lasker et al.,, 1990, DSS) as stars (with diffraction spikes, saturated point source, no diffuse halo) or galaxies. X-ray fields containing one or more bright stars but no bright galaxies within $`3\sigma `$ of the X-ray position were considered stellar. In these cases, the value of $`f_r`$ was forced to be $`0.02`$ for stellar fields rather than $`0.61`$, the extragalactic identification fraction. We may have missed a few radio stars with high proper motions because the USNO2 and radio positions were measured at different epochs.
Finally, we used the locally measured, rather than global, surface density $`\rho _\mathrm{v}`$ of background objects because many RBSC sources reside in optically overdense regions (clusters of galaxies).
## 3 Identification Results
For the parent sample, we selected the 5441 RBSC sources above the 0.1 counts s<sup>-1</sup> completeness limit with Galactic latitudes $`|b|>15^{}`$ (to reduce the number of objects which might confuse optical identification and minimize extinction at both optical and X-ray wavelengths) and J2000 $`\delta >40^{}`$, the NVSS declination limit. Of these, 1773 sources are readily identified as stars, in agreement with Bade et al., (1998). We attempted to make linked X-ray/radio/optical cross-identifications using X-ray/radio and radio/optical positional coincidence supplemented by radio flux densities and optical magnitudes, as described in Section 2.4.
We visually confirmed each match by overlaying the contours of the NVSS image and the RBSC error box on the DSS (see Figure RBSC-NVSS SAMPLE. I. RADIO AND OPTICAL IDENTIFICATIONS OF A COMPLETE SAMPLE OF 1500 BRIGHT X-RAY SOURCES). In addition, the IRAS Faint Source Catalogue (Moshir et al.,, 1992, FSC) has been matched against the RBSC (see Condon et al., 1998b, ), so we have plotted the FSC 3$`\sigma `$ error ellipses to aid identification and differentiate between physical emission processes. The nearest candidate usually has the highest calculated reliability (Equation 15), although the radio emission for some sources is extended and more care had to be taken in locating an optical counterpart. We use the term association to describe extended X-ray/radio sources which appear related but may not be spatially coincident. For instance, a cluster of galaxies may contain one or more radio galaxies and hot intercluster X-ray gas. We retained 19 identifications with low calculated reliabilities on the basis of other information, such as detection in a ROSAT pointed observation or no obvious optical counterpart within the error circle. All deviations from the standard reliability estimate are noted in Section 3.1. The reliabilities of these sources are given as ’$``$’.
Some RBSC-NVSS sources had no USNO sources listed near their NVSS positions, but objects were clearly visible on the DSS (i.e the red plate). Most were either very bright, extended galaxies or very faint objects. USNO sources were identified based on a positional coincidence of within $`2\stackrel{}{\mathrm{.}}.0`$ on the red and blue plates, so we might expect that some bright galaxies will have red and blue emission peaks separated by more than $`2\stackrel{}{\mathrm{.}}.0`$ and thus be rejected. For these we used positions and magnitudes from NED. Furthermore, the faint objects we find must only appear on the red plate (the lack of a blue plate counterpart presumably excluded them from the USNO catalog). We used the AIPS task IMFIT to measure their positions. Faint object magnitudes were estimated by comparison with nearby USNO sources.
These additional sources were added to the USNO list of optical candidates located within the search area and evaluated as part of the standard reliability calculation (see Equation 15). This additional step introduces a slight bias in the reliabilities of these new faint sources, because we have only added the faint sources in close proximity to the radio/X-ray positions, and not those from the rest of the $`3\mathrm{}`$ radius field. We can test the statistical error of this effect by simply doubling the number of local background objects; this is akin to probing roughly 2 magnitudes fainter than the USNO. While the change will depend on the exact field, we found the reliability with this background enhancement typically differed by only 1% (and at most 5%).
The Galactic absorption corrected X-ray fluxes come from the RBSC correlation catalogue (Voges et al.,, 1999), in which the RBSC count rate was converted to flux assuming a power law photon energy distribution of the form E$`^{\mathrm{\Gamma }_x+1}`$, common to AGN and clusters of galaxies. An average photon index $`\mathrm{\Gamma }_\mathrm{x}=2.3`$ was used, typical of extragalactic objects (Hasinger, Trümper & Schmidt 1991) over the ROSAT energy range. Corrections for Galactic absorption were based on hydrogen column densities N<sub>H</sub> obtained from Dickey & Lockman, (1990).
The error of the unabsorbed X-ray flux is about 30% taking into account errors in photon statistics and variable Galactic absorption. Flux estimates break down for cases when the actual $`\mathrm{\Gamma }_\mathrm{x}`$ deviates significantly from the mean value, there is additional intrinsic absorption, or when a single power law model inadequately describes the spectral shape (see Brinkmann et al., (1995) for a more detailed discussion of this problem).
A total of 1556 RBSC-NVSS sources were identified, on the criterion that the sum total of all radio identification reliabilities be greater than 0.50. The vast majority exceeded this cut-off limit a large margin. The distribution of RBSC-NVSS sources on the sky is fairly uniform (Figure RBSC-NVSS SAMPLE. I. RADIO AND OPTICAL IDENTIFICATIONS OF A COMPLETE SAMPLE OF 1500 BRIGHT X-RAY SOURCES). Roughly one half of the extragalactic RBSC sources did not have radio counterparts. Only 1 RBSC-NVSS source has no detectable optical counterpart (not surprising given the large RBSC error circles and the density of optical sources at the limiting magnitude of the POSS). Table 1 lists the 44 possible galactic stars so identified because they have $`B<`$ 13 mag, diffraction spikes, faint irregular radio contours or a spectral type from SIMBAD. Many have been independently verified with high resolution VLA images to be radio stars by Condon et al., (1997). Table 1 is arranged in the following manner:
Column 1. ROSAT BSC name, format 1RXS JHHMMSS.S+DDMMSS.
Column 2. SIMBAD Common Name, if available.
Column 3. ROSAT RBSC X-ray Flux.
Column 4. NVSS position, unless otherwise noted.
Column 5. NVSS 1.4 GHz Flux density.
Column 6. Identification Probability P(R,V) (Equation 15). Identifications with a ’$``$’ have been accepted despite having a low probability.
Column 7. SIMBAD object type, if available. Abbreviations from SIMBAD as follows:
| * | = | Star |
| --- | --- | --- |
| Al | = | Eclipsing binary of Algol type |
| bL | = | Eclipsing binary of Beta Lyr type |
| Cl | = | Star or globular cluster |
| DQ | = | Cataclysmic variable of DQ Her type |
| Fl | = | Flare star |
| Mi | = | Variable star of Mira Cet type |
| Pu | = | Pulsating variable star |
| Ro | = | Rotationally variable star |
| RS | = | Star of RS CVn type |
| sr | = | Semi-regular pulsating star |
| TT | = | T Tau-type star |
| V | = | Variable star |
| PN | = | Planetary nebula |
Column 8. SIMBAD B magnitude, if available; otherwise USNO magnitude enclosed in parentheses.
Column 9. SIMBAD spectral type, if available.
The remaining 1512 objects are extragalactic and listed in Table 3.3. Figure RBSC-NVSS SAMPLE. I. RADIO AND OPTICAL IDENTIFICATIONS OF A COMPLETE SAMPLE OF 1500 BRIGHT X-RAY SOURCES shows the cumulative fraction of identifications versus reliability. The vast majority ($`70\%`$) of identifications have reliabilities greater than 99%. Table 3.3 is arranged in the following manner:
Column 1. ROSAT RBSC name, format 1RXS JHHMMSS.S+DDMMSS.
Column 2. NED Common Name, if available.
Column 3. ROSAT BSC X-ray Flux.
Column 4. NVSS position, unless noted otherwise.
Column 5. NVSS 1.4 GHz Flux density.
Column 6. Identification Probability P(R,V) (Equation 15). Identifications with a ’$``$’ have been accepted despite having a low probability.
Column 7. B magnitude from NED, if available; otherwise from USNO as indicated by ’()’. Identifications with a ’$``$’ denote empty fields.
Column 8. Redshift, if known. From NED unless noted otherwise. Identifications with a ’$``$’ represent objects for which spectra was taken but a redshift could not be determined. We have obtained optical spectra with the KPNO 2.1 m telescope for many of the bright identifications which did not possess published redshifts. The detailed description of these spectra will be given in a subsequent paper (Bauer, Thuan & Condon 2000).
Column 9. NED optical morphology, if available.
Column 10. Spectral classification, if known. From NED unless noted otherwise in redshift column.
| AGN | = | Active galactic nuclei, type unknown |
| --- | --- | --- |
| Blazar | = | BL Lacertae type object with variable emission line spectra |
| BLLAC | = | BL Lacertae type object |
| BLRG | = | Broad-line radio galaxy |
| cD | = | Central dominant galaxy, Early-type stellar absorption continuum |
| Early | = | Early-type stellar absorption continuum |
| HII | = | Starburst galaxy with narrow emission-line spectra similar to HII regions |
| HPQ | = | High polarization QSO |
| LINER | = | Low ionization narrow emission-line region |
| LPQ | = | Low polarization QSO |
| QSO | = | Quasi-stellar object |
| Sy1-2 | = | Seyfert galaxy, emission-line spectra classified from type 1 to 2 |
There are a small number of cases having low reliabilities, but which may in fact be true identifications. They generally fall into two categories; objects needing higher radio resolution and candidates that have optical/radio matches farther than 3$`\sigma `$ from the X$``$ray source which did not have a convincing optical match. Many of the latter are likely to be associated with nearby clusters.
### 3.1 Notes on Individual Identifications
1RXS J000312.3$``$355541: Member of cluster Abell 2717.
1RXS J001031.3$`+`$105832: Member of triple galaxy III Zw 002.
1RXS J001118.5$``$285126: Member of cluster Abell 2734. Abell 2744. Diffuse radio emission; possible radio halo?
1RXS J001823.8$`+`$300357: Member of group WP 01.
1RXS J001906.0$``$202638: Cluster Abell S0026 is $`50\mathrm{}`$ from ID and overlaps within positional errors.
1RXS J002041.8$``$254307: Member of cluster Abell 0022.
1RXS J002136.8$`+`$280305: Group IV Zw 015 is $`77\mathrm{}`$ from ID and overlaps within positional errors.
1RXS J002430.5$``$292856: Extended radio emission.
1RXS J002534.9$``$330255: Greater than $`3\sigma `$ from X$``$ray position. Low X$``$ray flux and large error region imply unreliable X$``$ray position. Optical and radio information point towards a correlation. Member of cluster ABELL S0041.
1RXS J002811.6$`+`$310342: Possible ID is not at 002812.61$`+`$310315.3 (as identified in RGB). Better radio resolution would strengthen ID.
1RXS J003037.6$``$241053: Member of cluster Abell 0047.
1RXS J003413.7$``$212619: Member of group SCG 05.
1RXS J004324.9$``$203728: Member of cluster Abell 2813.
1RXS J004924.0$``$293128: Member of cluster Abell S0084.
1RXS J010721.2$`+`$322254: Greater than $`3\sigma `$ from X$``$ray position. Low X$``$ray flux and large error region imply unreliable X$``$ray position. Optical and radio information point towards a correlation. Member of group IV Zw 038.
1RXS J010750.4$``$364342: Member of cluster Abell 2871.
1RXS J010850.5$``$152438: Member of cluster Abell 0151.
1RXS J011301.1$`+`$153041: Member of cluster Abell 0160. Given the low X$``$ray flux and large error region imply unreliable X$``$ray position. Optical and radio information point towards a correlation.
1RXS J011354.9$``$314538: Greater than $`3\sigma `$ from X$``$ray position. Given the low X$``$ray flux and large error region imply unreliable X$``$ray position. Optical and radio information point towards a correlation. Pointed ASCA observation by Arimoto et al. 1997 confirms ID. Member of cluster Abell S0141.
1RXS J011513.2$`+`$405743: More reliable ID is a blend of star and galaxy on DSS.
1RXS J012058.3$``$135100: Member of cluster CID 10.
1RXS J012337.2$`+`$331506: Member of pair ARP 229.
1RXS J012507.3$`+`$084124: Member of cluster Abell 0193.
1RXS J012910.8$``$214158: Galaxy composed of two compact optical cores.
1RXS J013152.8$``$133651: Member of cluster Abell 0209.
1RXS J013632.9$`+`$390556: KPNO spectrum obtained. No features revealed amid strong blue continuum.
1RXS J013715.1$``$091145: Member of cluster CID 12.
1RXS J015244.7$`+`$360855: Member of cluster Abell 0262.
1RXS J015705.3$`+`$412029: Member of cluster Abell 0276.
1RXS J020144.2$``$021158: Member of cluster Abell 0291.
1RXS J023727.6$``$263027: Member of cluster Abell 0368.
1RXS J024122.1$``$283919: Cluster Abell 3041 is $`53\mathrm{}`$ from ID and overlaps within positional errors.
1RXS J024620.0$``$301639: Supernova SN 1992bd is $`1\mathrm{}`$ from ID and is also a valid identification.
1RXS J024937.1$``$311114: Member of cluster Abell S0301.
1RXS J030150.6$`+`$355003: Member of group UGC 02489 and cluster Abell 0407. V Zw 311 NOTES02 appears to be a cD from DSS and is likely to be true ID. Better radio resolution needed to clarify.
1RXS J031611.4$`+`$090445: KPNO spectrum obtained. No features revealed amid the flat continuum.
1RXS J033828.8$``$352701: Member of Fornax cluster.
1RXS J041325.4$`+`$102756: Member of cluster Abell 0478.
1RXS J042551.3$``$083329: Member of cluster EXO 0422$``$086.
1RXS J042700.9$``$204321: Member of cluster Abell 0490.
1RXS J043337.7$``$131528: Member of cluster Abell 0496.
1RXS J043643.8$`+`$100323: Member of cluster MS 0433.9$`+`$0957.
1RXS J043902.0$`+`$052050: Member of cluster RX J04390$`+`$0520.
1RXS J044511.8$``$155118: Member of cluster CID 25.
1RXS J044803.3$``$202741: High resolution radio position from Owen, White & Ge (1993) confirms ID. Member of cluster Abell 0514.
1RXS J044814.8$`+`$095259: Member of cluster RX J0448.2+0952.
1RXS J045859.6$``$002904: Member of cluster CID 27.
1RXS J050049.4$``$384038: Member of cluster Abell 3301.
1RXS J050120.1$``$033223: High resolution radio position from Owen, White & Ge (1993) confirms ID. Member of cluster Abell 0531.
1RXS J050507.9$``$322007: Member of cluster Abell 3313.
1RXS J051611.5$``$000807: RBSC position incorrect. X-ray observations from Ceballos & Barcons, (1996) confirm ID.
1RXS J054738.7$``$315228: Member of cluster Abell 3364.
1RXS J055040.7$``$321619: Cluster Abell S0549 ($`z=0.04030`$) is $`13\mathrm{}`$ from ID and overlaps well within positional errors.
1RXS J055711.1$``$372813: Member of cluster Abell S0555.
1RXS J060552.7$``$351759: Member of cluster Abell 3378.
1RXS J062707.5$``$352917: Member of cluster Abell 3392.
1RXS J064326.5$`+`$421420: High resolution radio position from Laurent-Muehleisen et al., (1997) confirms ID.
1RXS J070427.0$`+`$631856: Member of cluster Abell 0566.
1RXS J070907.6$`+`$483657: Member of cluster Abell 0569.
1RXS J071006.0$`+`$500243: High resolution radio position from Laurent-Muehleisen et al., (1997).
1RXS J073221.5$`+`$313750: Member of cluster Abell 0586.
1RXS J074144.8$`+`$741444: Member of cluster ZwCl 0735.7$`+`$7421.
1RXS J074238.0$`+`$092213: Member of cluster Abell 0592. Better radio resolution needed.
1RXS J081021.3$`+`$421657: Member of cluster \[VMF98\] 047.
1RXS J081112.3$`+`$700300: Member of cluster Abell 0621
1RXS J081929.5$`+`$704221: Very extended radio and optical emission associated with the dwarf galaxy Holmberg II.
1RXS J082209.5$`+`$470601: Member of cluster Abell 0646.
1RXS J083811.0$`+`$245336: Member of pair CGCG 120$``$011.
1RXS J083950.7$``$121424: High resolution radio position from Li & Jin, (1996) confirm ID as \[HB89\] 0837-120 and not \[EYG89\] 087.
1RXS J084255.9$`+`$292752: Member of cluster ZwCl 0839.9$`+`$2937.
1RXS J091037.2$`+`$332920: KPNO spectrum obtained. No features revealed among strong blue continuum.
1RXS J091949.1$`+`$334532: X-ray position offset, but Einstein observations from Fabbiano, Kim & Trinchieri (1992) confirm ID. Member of cluster Abell 0779.
1RXS J092406.1$`+`$141002: Member of cluster Abell 0795.
1RXS J094740.4$``$305655: Member of triple galaxy AM 0945$``$304.
1RXS J095821.8$``$110344: Member of cluster Abell 0907.
1RXS J100121.5$`+`$555351: Member of cluster \[YGK80\] 0957$`+`$561CLUSTER. QSO \[HB89\] 0957$`+`$561 and galaxy \[YGK81\] 097 are $`4\mathrm{}`$ and $`9\mathrm{}`$ from ID and overlap within positional errors.
1RXS J100639.1$`+`$255450: Member of cluster Abell 0923.
1RXS J101335.9$``$135108: Member of cluster RX J10136$``$1350.
1RXS J101912.1$`+`$635802: Member of pair CGCG 313$``$011.
1RXS J102228.9$`+`$500630: Member of cluster Abell 0980.
1RXS J102339.3$`+`$041117: Member of cluster ZwCl 1021.0$`+`$0426.
1RXS J102350.7$``$271522: Member of cluster Abell 3444.
1RXS J102758.9$``$064804: Member of cluster RX J1027.9-0648.
1RXS J103155.8$``$141659: Multiple redshift systems overlap within positional errors.
1RXS J103459.5$`+`$304138: Member of cluster Abell 1045.
1RXS J104043.7$`+`$395706: Member of cluster Abell 1068.
1RXS J104431.7$``$070404: Member of cluster Abell 1084.
1RXS J104651.9$``$253546: Member of cluster RX J10468$``$2535
1RXS J105344.2$`+`$492956: Member of cluster MS 1050.7$`+`$4946.
1RXS J105825.9$`+`$564716: Member of cluster Abell 1132.
1RXS J111137.2$`+`$405031: Member of cluster Abell 1190.
1RXS J111422.6$`+`$582318: High resolution radio position from Laurent-Muehleisen et al., (1997) confirms ID.
1RXS J111450.1$``$121351: Cluster Abell 1211 is $`89\mathrm{}`$ from ID and overlaps within positional errors.
1RXS J113121.4$`+`$333447: Member of cluster Abell 1290.
1RXS J113153.7$``$195543: Member of cluster Abell 1300.
1RXS J113222.4$`+`$555828: Galaxy MCG $`+`$09$``$19$``$110 ($`z=0.05130`$) is $`17\mathrm{}`$ from ID and overlaps within positional errors. Member of cluster RX J1132.3$`+`$5558.
1RXS J113448.4$`+`$490438: Member of cluster Abell 1314.
1RXS J114124.2$``$121632: Member of cluster RX J1141.4$``$1216.
1RXS J114442.2$`+`$672435: Member of cluster Abell 1366. Optical and radio information point towards a correlation. Possible X-ray cluster emission?
1RXS J114452.7$`+`$194706: Greater than $`3\sigma `$ from RBSC position. Public ROSAT HRI observations place X$``$ray source much closer to reliable ID. Association with cluster Abell 1367.
1RXS J114947.0$``$121845: Member of Abell 1391. Probable association.
1RXS J115518.9$`+`$232431: Member of cluster Abell 1413. taken from the FIRST 1.4 GHz survey (White et al.,, 1997). NVSS flux density is 5.4 mJy. Member of cluster RX J11570$`+`$2415.
1RXS J115719.0$`+`$333645: High resolution radio position from Owen, White & Ge (1993) confirms more reliable ID. Member of cluster Abell 1423.
1RXS J120511.7$`+`$392043: Member of cluster RX J1205.1$`+`$3920
1RXS J121105.4$`+`$352005: Member of cluster RX J1211.0$`+`$3520
1RXS J122752.7$`+`$632317: Member of cluster RX J1227.8$`+`$6323.
1RXS J122945.9$`+`$075927: Member of cluster Virgo S.
1RXS J123658.8$`+`$631111: Member of cluster Abell 1576.
1RXS J124349.3$``$153320: Member of cluster Abell 1603.
1RXS J125233.7$``$311605: Member of cluster RX J12525$``$3116.
1RXS J125422.2$``$290034: Member of cluster Abell 3528.
1RXS J125710.0$``$172345: Pointed ROSAT observation from Peres et al., (1998) confirms ID. Member of cluster Abell 1644.
1RXS J125921.5$``$041131: Member of cluster Abell 1651.
1RXS J130250.3$``$023041: Member of cluster Abell 1663.
1RXS J130343.6$``$241506: Member of cluster Abell 1664.
1RXS J130346.7$`+`$191635: Member of cluster Abell 1668.
1RXS J130552.6$`+`$305405: Member of cluster Abell 1677.
1RXS J130916.1$``$013658: Member of cluster MS 1306.7$``$0121.
1RXS J131129.5$``$012017: Member of cluster Abell 1689.
1RXS J131506.8$`+`$514931: Member of cluster Abell 1703.
1RXS J132016.3$`+`$330828: Member of cluster RX J1320.1$`+`$3308.
1RXS J132542.1$``$022800: No optical counterpart on DSS.
1RXS J132549.2$`+`$591937: Member of cluster Abell 1744.
1RXS J132617.4$`+`$001329: Member of cluster RX J1326.3$`+`$0013.
1RXS J132758.9$``$313027: Greater than $`3\sigma `$ from RBSC position. Optical and radio information point towards a correlation. Member of cluster Abell 3558. Possible X-ray blend of star and galaxy?
1RXS J133226.0$``$330812: Better radio resolution confirms ID.
1RXS J134104.8$`+`$395942: Member of cluster Abell 1774.
1RXS J134152.6$`+`$262230: Cluster Abell 1775 ($`z=0.06960`$) is $`72\mathrm{}`$ from ID and overlaps within positional errors.
1RXS J134730.5$``$114455: Member of cluster RX J13475$``$1145.
1RXS J134852.6$`+`$263541: Member of cluster Abell 1795.
1RXS J140102.1$`+`$025249: Member of cluster Abell 1835.
1RXS J140337.0$``$335840: Member of cluster Abell S0753.
1RXS J140728.4$``$270055: Member of cluster Abell 3581.
1RXS J141342.6$`+`$433938: Member of cluster Abell 1885.
1RXS J141357.3$`+`$711751: Association with cluster Abell 1895 ($`z=0.22492`$).
1RXS J142139.7$`+`$371743: Member of cluster Abell 1902.
1RXS J143236.0$`+`$313855: Cluster Abell 1930 ($`z=0.13130`$) is $`119\mathrm{}`$ from ID and overlaps within positional errors.
1RXS J143527.9$`+`$550747: Member of cluster Abell 1940.
1RXS J144428.4$`+`$311304: Cluster Abell 1961 ($`z=0.23200`$) is $`117\mathrm{}`$ from ID and overlaps within positional errors.
1RXS J145307.8$`+`$215333: Member of cluster Abell 1986.
1RXS J145431.4$`+`$183834: Member of cluster Abell 1991.
1RXS J145434.1$`+`$080250: Greater than $`3\sigma `$ from RBSC position. Bade et al., (1998) claim optical$``$$``$radio ID is the correct X$``$ray source. Optical and radio information point towards a correlation.
1RXS J145507.9$`+`$192025: High resolution radio position from Laurent-Muehleisen et al., (1997) confirms ID.
1RXS J145715.4$`+`$222026: Member of cluster MS 1455.0$`+`$2232.
1RXS J145904.1$``$084254: Better radio resolution needed to confirm ID.
1RXS J150020.7$`+`$212213: Member of cluster Abell 2009.
1RXS J151056.3$`+`$054431: Member of cluster Abell 2029.
1RXS J151127.2$`+`$062153: Greater than $`3\sigma `$ from RBSC position. Optical and radio information point towards a correlation. Member of cluster Abell 2033. Possible X-ray cluster emission?
1RXS J151642.4$`+`$070058: Member of cluster Abell 2052.
1RXS J151845.3$`+`$061340: High resolution radio position from Laurent-Muehleisen et al., (1997) confirms ID. Member of cluster Abell 2055.
1RXS J152151.0$`+`$074221: Member of cluster MKW 03s.
1RXS J152305.7$`+`$083550: ID greater than $`3\sigma `$ of RBSC position. There is a strong optical/radio candidate (CGCG 077-097) located within the cluster Abell 2063.
1RXS J153950.3$`+`$304305: Member of cluster Abell 2110.
1RXS J154009.4$`+`$141116: High resolution radio position from Laurent-Muehleisen et al., (1997) confirms ID (a star?).
1RXS J155611.0$`+`$662123: Association with cluster Abell 2146.
1RXS J160435.1$`+`$174333: High resolution X-ray position from Burstein et al., (1997) confirms ID.
1RXS J160456.8$`+`$235604: Member of cluster AWM 4.
1RXS J160740.7$`+`$254106: Optical ID is offset from strong radio source (FIRST ID), possibly indicative of a misidentification with a foreground star.
1RXS J161546.9$``$060841: Radio emission is very diffuse, from multiple components. Better radio resolution needed. Association.
1RXS J161711.4$`+`$063816: No object found within $`3\sigma `$ of X$``$ray source $``$ X$``$ray errors too conservative? Identified as RGB J1617$`+`$066. Optical and radio information point towards a correlation. High resolution radio position from Laurent-Muehleisen et al., (1997).
1RXS J162032.0$`+`$295321: Member of cluster Abell 2175.
1RXS J162100.4$`+`$254547: Member of Abell 2177.
1RXS J162837.7$`+`$393249: Member of cluster Abell 2199.
1RXS J163246.8$`+`$053423: Member of cluster Abell 2204
1RXS J164022.1$`+`$464215: Member of cluster Abell 2219
1RXS J164238.9$`+`$272621: Member of cluster Abell 2223.
1RXS J170242.5$`+`$340336: Member of cluster Abell 2244
1RXS J170941.2$`+`$342529: Member of cluster Abell 2249.
1RXS J171519.5$`+`$572430: Member of cluster CID 71.
1RXS J171706.8$`+`$293117: Member of cluster RBS 1634.
1RXS J171718.9$`+`$422652: Member of cluster RX J17173$`+`$4227.
1RXS J171746.9$`+`$194057: Association with cluster Abell 2254.
1RXS J171810.9$`+`$563955: Member of cluster ZwCl 1717.9$`+`$5636.
1RXS J172009.3$`+`$263727: Member of cluster RX J17201$`+`$2637.
1RXS J172226.7$`+`$320752: Member of cluster Abell 2261.
1RXS J173257.0$`+`$403635: Association with cluster Abell 2272.
1RXS J173301.7$`+`$434533: Member of cluster CID 72.
1RXS J174414.3$`+`$325925: Member of cluster ZwCl 1742.1$`+`$3306.
1RXS J175004.4$`+`$470037: High resolution radio position from Laurent-Muehleisen et al., (1997) confirms ID.
1RXS J175441.9$`+`$680334: Cluster ZwCl 1754.5$`+`$6807 ($`z=0.07700`$) is $`53\mathrm{}`$ from ID and overlaps within positional errors.
1RXS J175706.9$`+`$535130: Member of cluster Abell 2292.
1RXS J182157.4$`+`$642051: High resolution radio position from Laurent-Muehleisen et al., (1997) confirms ID.
1RXS J203445.3$``$354921: Member of cluster Abell 3695.
1RXS J210027.7$``$170913: Better radio resolution needed to confirm ID.
1RXS J210707.7$``$252643: Cluster Abell 3744 is nearby and provides a considerable amount of diffuse radio emission as well. Better radio resolution needed.
1RXS J212706.7$``$120927: Greater than $`3\sigma `$ of RBSC position. Optical and radio information point towards a correlation.
1RXS J213414.4$``$311737: Better radio resolution needed to confirm ID.
1RXS J213504.7$`+`$085807: Optical ID is actually a composite of a bright star and another fainter source.
1RXS J214015.3$``$233946: Member of cluster MS 2137.3$``$2353.
1RXS J215336.1$`+`$174111: Member of cluster Abell 2390.
1RXS J215541.5$`+`$123144: Member of cluster Abell 2396
1RXS J221020.8$``$121040: Member of cluster Abell 2420.
1RXS J221614.9$``$092033: Greater than $`3\sigma `$ of RBSC position. Optical and radio information point towards a correlation. Possible X-ray cluster emission?
1RXS J224919.5$``$372724: Member of cluster Abell S1065.
1RXS J225019.2$`+`$105406: Member of cluster Abell 2495.
1RXS J225334.1$``$334309: Member of cluster Abell 3934.
1RXS J230714.9$``$151322: Member of cluster Abell 2533.
1RXS J231607.5$``$202739: Member of cluster Abell 2566.
1RXS J231712.8$`+`$184237: Member of group HCG 094.
1RXS J231829.5$`+`$184246: Greater than $`3\sigma `$ of RBSC position. Optical and radio information point towards a correlation. Member of cluster Abell 2572. Possible X-ray cluster emission?
1RXS J232125.9$``$231230: Member of cluster Abell 2580.
1RXS J232519.4$``$120741: Member of cluster Abell 2597.
1RXS J233631.0$`+`$210848: Member of cluster Abell 2626.
1RXS J233641.8$`+`$235526: Member of cluster Abell 2627.
1RXS J233827.6$`+`$270045: Greater than $`3\sigma `$ of RBSC position. Optical and radio information point towards a correlation. Member of cluster Abell 2634. Pointed ROSAT HRI observations by Sakelliou & Merrifield 1998 show X-ray emission is a blend of cluster and individual galaxy emission.
J234741.2$``$280829: X-ray emission appears to be from bright knot in radio jet. Optical and radio information point towards a correlation.
1RXS J235034.7$`+`$292924: Cluster MS 2348.0$`+`$2913 ($`z=0.09500`$) is $`22\mathrm{}`$ from ID and overlaps within positional errors. KPNO spectrum obtained. Strong absorption features at $`z=`$-$`0.00020`$, possibly a foreground star amid the cluster. Better optical and radio resolution needed to clarify.
1RXS J234741.2$``$280829: ID greater than $`3\sigma `$ of RBSC position. There is a strong optical/radio candidate nearby. Member of cluster Abell 4038.
### 3.2 Previous RBSC Misidentifications?
1RXS J000856.1$`+`$411034: Better radio resolution needed. Possible X-ray ID does not seem to be B3 0006$`+`$408 (as identified in RGB), but rather a bright, diffuse object at 00h08m56.0s $`+`$41d10m09.1s.
1RXS J001144.5$`+`$322441: Radio ID is 87GB 000913.4$`+`$320913 (as identified in RGB), and X$``$ray ID is cluster Abell 0007.
1RXS J220220.8$`+`$035306: Optical/X$``$ray ID greater than $`3\sigma `$ from radio position (NVSS and RGB). Background radio object misidentified with B=13.2 X-ray star?
### 3.3 A simple, well-defined sample of AGN
The RBSC-NVSS sample is defined by the following criteria:
* RBSC count rate $``$ 0.1 counts s<sup>-1</sup>
* $`f_r2.5mJy`$
* $`\delta 40^{}`$
* $`|b|15^{}`$
Columns 1 and 2 of Table 3 summarize the results of the selection process. We should note here, that while the method we employ here is quite robust, the large number of RBSC sources and the loose positional errors dictate that there will be some spurious spatial coincidences. To estimate the number of false matches we might have, we shifted the X-ray positions by $`8\mathrm{}`$ and performed the identification procedure again. The results are listed in column 3 of Table 3, where we find that the expected percentage of spurious matches is around 3$`\%`$ for reliabilities above 0.50. This fraction could be further reduced by comparing the flux density and magnitude distributions of the true and spurious objects; most spurious sources lie near the NVSS and USNO survey limits, while RBSC-NVSS sources do not.
Forthcoming papers in this series will present the radio, optical, and X-ray properties of the sample and optical identifications and classifications of a subset of objects. Using a number of criteria, we find that this sample of 1512 extragalactic objects is comprised almost entirely of AGN, making this the largest, complete sample of its kind. It represents a major step forward in the identification of RBSC objects and contains a large sample of both radio-loud and radio-quiet X-ray objects (previous surveys of this type have typically sampled only the radio-loud population).
FEB thanks W. Brinkmann and J. Siebert for their helpful comments and assistance with the RBSC. He also acknowledges support from a National Radio Astronomy Observatory Jansky Pre-doctoral Fellowship. JJB acknowledges support from the NSF through grant AST 9320547. This research has made use of data obtained through the High Energy Astrophysics Science Archive Research Center Online Service, provided by the NASA/Goddard Space Flight Center, the NASA/IPAC Extragalactic Database (NED) which is operated by the Jet Propulsion Laboratory, Caltech, under contract with the National Aeronautics and Space Administration and the SIMBAD database, operated at CDS, Strasbourg, France. The Digitized Sky Surveys were produced at the Space Telescope Science Institute under US government grant NAG W-2166.
|
warning/0003/nucl-th0003040.html
|
ar5iv
|
text
|
# Projectile structure effects in the Coulomb breakup of one-neutron halo nuclei
## 1 Introduction
It has now been well established that close to the neutron drip line, there exist nuclei having one, or some times two, very loosely bound valence neutrons extending too far out in space with respect to a dense charged core . The properties of these neutron halo nuclei have been reviewed by several authors (see e.g. ). The halo systems, which involve a new form of nuclear matter, are characterized by large reaction and Coulomb dissociation cross sections . Moreover, in the breakup reactions induced by these nuclei, the angular distributions of neutrons measured in coincidence with the core nuclei are strongly forward peaked and the parallel momentum distributions of the core fragment have very narrow widths . Due to their strikingly different properties as compared to those of the stable nuclei, such systems provide a stringent test of the nuclear structure models developed for the latter.
The Coulomb dissociation is a significant reaction channel in the scattering of halo nuclei from a stable heavy target nucleus. It provides a convenient tool to investigate their structure. For instance, it would place constraints on their electric dipole response . Of course, in the Serber type of models , the breakup cross sections are directly related to the momentum space wave function of the projectile ground state. The studies of the Coulomb dissociation of weakly bound nuclei are also of interest due to their application in determining the cross sections of the astrophysically interesting radiative capture reactions at solar temperatures .
The Coulomb dissociation of halo nuclei has been investigated by several authors recently, using a number of different theoretical approaches. A semiclassical coupled channel formalism has been used by authors of Ref. , while in Refs. the relative motion of the core and the valence particle is described by a time dependent Schrödinger equation. The results within these approaches depend on the range of the impact parameter associated with the straight line trajectories used to describe the motion of the projectile in the field of target nuclei. However, in these studies the emphasis was on investigating the dynamics of the Coulomb dissociation, and not the structure of the projectile ground state which was assumed to have some very simple zero range (ZR) form. Similar assumption for the projectile structure was also made in other semiclassical and prior form distorted wave Born approximation (DWBA) calculations .
On the other hand, the post form DWBA theory of the breakup reactions incorporates the details of the ground state structure of the projectile in the breakup amplitude . However, in an earlier application of this theory to calculate the breakup of the halo nuclei, the simplifying approximation of the zero range interaction was used. This approximation necessarily excludes the use of this theory to describe the breakup of non-$`s`$ – wave projectiles. Moreover its applicability is questionable even for $`s`$ – wave projectiles at higher beam energies . Therefore, to investigate the details of the projectile structure through breakup reactions within this theory, the inclusion of the finite range effects is necessary. It may be noted that in a recently formulated adiabatic model of the Coulomb breakup reactions, where it is assumed that the projectile excitation is predominantly to the states of the low internal energy, the details of the ground state wave function also enters in the transition amplitude.
In this paper, we present a theoretical model to describe the Coulomb breakup of a projectile within the framework of the post form DWBA where finite range effects are included, approximately, via a local momentum approximation (LMA) . The exact treatment of the finite range effects within this theory, although desirable, is much too complicated as it would lead to six dimensional integrals involving functions which are oscillatory asymptotically. The LMA leads to two simplifying features. First, it factorises the dynamics of the reaction from the structure effects of the projectile and second, it results in an amplitude where the term describing the dynamics of the process is the same as that evaluated in the ZR approximation. We present the application of this theory to the Coulomb breakup of neutron rich nuclei <sup>11</sup>Be and <sup>15,17,19</sup>C on a number of heavy targets. We attempt to put constraints on the ground state structure of these nuclei by analyzing almost all the available data on the energy, angular and longitudinal momentum distributions of fragments using various configurations of the projectile ground state.
Our formalism is presented in section 2. In section 3 we present and discuss the results of our calculations for various observable for the reactions mentioned above. The summary and the conclusions of our work are described in section 4. The validity of the approximate method used by us to include the finite range effects is presented in appendix A.
## 2 Formalism
We consider the reaction $`a+tb+c+t`$, where the projectile $`a`$ breaks up into fragments $`b`$ (charged) and $`c`$ (uncharged) in the Coulomb field of a target $`t`$. The coordinate system chosen is shown in Fig. 1.
The position vectors satisfy the following relations
$`𝐫`$ $`=`$ $`𝐫_i\alpha 𝐫_1,\alpha ={\displaystyle \frac{m_c}{m_c+m_b}}`$ (1)
$`𝐫_c`$ $`=`$ $`\gamma 𝐫_1+\delta 𝐫_i,\delta ={\displaystyle \frac{m_t}{m_b+m_t}},\gamma =(1\alpha \delta )`$ (2)
The exact post form $`T`$ \- matrix for this case is
$`T`$ $`=`$ $`\chi _b^{()}(𝐤_b,𝐫)\mathrm{\Phi }_b(\xi _b)\chi _c^{()}(𝐤_c,𝐫_c)\mathrm{\Phi }_c(\xi _c)|V_{bc}(𝐫_1)|\mathrm{\Psi }_a^{(+)}(\xi _a,𝐫_1,𝐫_i),`$ (3)
where $`\chi ^{}s`$ are the distorted waves for relative motions of $`b`$ and $`c`$ with respect to $`t`$ and the center of mass (c.m.) of $`b+t`$ system respectively, and $`\mathrm{\Phi }^{}s`$ are the internal state wave functions of the concerned particles with internal coordinates $`\xi `$. $`\mathrm{\Psi }_a^{(+)}(\xi _a,𝐫_1,𝐫_i)`$ is the exact three-body scattering wave function of the projectile with a wave vector $`𝐤_a`$, with outgoing wave boundary condition. $`𝐤_b`$, $`𝐤_c`$ are Jacobi wave vectors of $`b`$ and $`c`$ respectively in the final channel of the reaction. $`V_{bc}(𝐫_1)`$ is the interaction between $`b`$ and $`c`$. The charged fragment $`b`$ interacts with the target by a point Coulomb interaction and hence $`\chi _b^{()}(𝐤_b,𝐫)`$ is a Coulomb distorted wave with incoming wave boundary condition. For pure Coulomb breakup, the interaction between the target and uncharged fragment $`c`$ is zero and hence $`\chi _c^{()}(𝐤_c,𝐫_c)`$ is replaced by a plane wave.
In the distorted wave Born approximation (DWBA), we write
$`\mathrm{\Psi }_a^{(+)}(\xi _a,𝐫_1,𝐫_i)=\mathrm{\Phi }_a(\xi _a,𝐫_1)\chi _a^{(+)}(𝐤_a,𝐫_i),`$ (4)
where $`\mathrm{\Phi }_a(\xi _a,𝐫_1)`$ represents the bound state wave function of the projectile having its radial and angular parts as $`u_{\mathrm{}}(r_1)`$ and $`Y_\mathrm{}m(\widehat{𝐫}_1)`$ respectively, which are associated with the relative motion of $`b`$ and $`c`$. $`\chi _a^{(+)}(𝐤_a,𝐫_i)`$ is the Coulomb distorted scattering wave describing the relative motion of the c.m. of the projectile with respect to the target with outgoing wave boundary condition. The assumption inherent in Eq. (4) is that the breakup channels are very weakly coupled and hence this coupling needs to be treated only in the first order. The transition amplitude (written in the integral form) is now given by
$`T={\displaystyle }`$ $`d\xi d𝐫_1d𝐫_i\chi _b^{()}(𝐤_b,𝐫)\mathrm{\Phi }_b^{}(\xi _b)e^{i𝐤_c.𝐫_c}\mathrm{\Phi }_c^{}(\xi _c)V_{bc}(𝐫_1)`$ (5)
$`\times `$ $`\mathrm{\Phi }_a(\xi _a,𝐫_1)\chi _a^{(+)}(𝐤_a,𝐫_i).`$
The integrals over the internal coordinates $`\xi `$ give
$`{\displaystyle 𝑑\xi \mathrm{\Phi }_b^{}(\xi _b)\mathrm{\Phi }_c^{}(\xi _c)\mathrm{\Phi }_a(\xi _a,𝐫_1)}`$
$`={\displaystyle \underset{\mathrm{}mj\mu }{}}\mathrm{}mj_c\mu _c|j\mu j_b\mu _bj\mu |j_a\mu _ai^{\mathrm{}}u_{\mathrm{}}(r_1)Y_\mathrm{}m(\widehat{𝐫}_1).`$ (6)
In Eq. (6) $`\mathrm{}`$ is the relative motion angular momentum between $`b`$ and $`c`$ . This is coupled to spin of $`c`$ and the resultant $`j`$ is coupled to spin of (the inert core) $`b`$ to get the spin of $`a`$ ($`j_a`$). Using Eq. (6), the $`T`$-matrix can be written as
$`T`$ $`=`$ $`{\displaystyle \underset{\mathrm{}mj\mu }{}}\mathrm{}mj_c\mu _c|j\mu j_b\mu _bj\mu |j_a\mu _ai^{\mathrm{}}\widehat{\mathrm{}}\beta _\mathrm{}m(𝐤_b,𝐤_c;𝐤_a),`$ (7)
where
$`\widehat{\mathrm{}}\beta _\mathrm{}m(𝐤_b,𝐤_c;𝐤_a)={\displaystyle }`$ $`d𝐫_1d𝐫_i\chi _b^{()}(𝐤_b,𝐫)e^{i𝐤_c.𝐫_c}V_{bc}(𝐫_1)`$ (8)
$`\times \varphi _a^\mathrm{}m(𝐫_1)\chi _a^{(+)}(𝐤_a,𝐫_i),`$
with $`\beta _\mathrm{}m`$ being the reduced $`T`$ – matrix and $`\widehat{\mathrm{}}=\sqrt{2\mathrm{}+1}`$. We have written $`\varphi _a^\mathrm{}m(𝐫_1)=u_{\mathrm{}}(r_1)Y_\mathrm{}m(\widehat{𝐫}_1)`$.
It may be noted that the reduced amplitude $`\beta _\mathrm{}m`$ involves a six dimensional integral which makes its evaluation quite complicated. The problem gets further acute due to the fact that the integrand involves three scattering waves which have oscillatory behavior asymptotically. Therefore, several approximate methods have been used in the literature to avoid the evaluation of six dimensional integrals. In the zero range approximation (ZRA) one assumes
$`V_{bc}(𝐫_1)\varphi _a^\mathrm{}m(𝐫_1)`$ $`=`$ $`D_0\delta (𝐫_1),`$ (9)
where $`D_0`$ is the usual zero range constant. This approximation reduces the six dimensional integral in Eq.(8) to a three-dimensional one. The corresponding amplitude is written as
$`\beta _{00}^{ZR}`$ $`=`$ $`D_0\chi _b^{()}(𝐤_b,𝐫_i)e^{i\delta 𝐤_c.𝐫_i}|\chi _a^{(+)}(𝐤_a,𝐫_i).`$ (10)
In Eq. (10), the details of the projectile structure enter in the reaction amplitude only as a multiplicative constant $`D_0`$. However, ZRA necessarily restricts the relative motion between $`b`$ and $`c`$ in the projectile $`a`$ to $`s`$ – state only. Even for such cases, this approximation may not be satisfied for heavier projectiles and at higher beam energies .
Baur and Trautmann (BT) have proposed an alternative approximate scheme in which the projectile c.m. coordinate in the corresponding distorted wave in Eq. (8) is replaced by that of the core-target system, i.e. $`𝐫_i𝐫`$. With this approximation the amplitude $`\beta _\mathrm{}m`$ splits into two factors each involving a three dimensional integral
$`\widehat{\mathrm{}}\beta _\mathrm{}m^{BT}`$ $`=`$ $`e^{i𝐤_c.𝐫_1}|V_{bc}|\varphi _a^\mathrm{}m(𝐫_1)\chi _b^{()}(𝐤_b,𝐫)e^{i\delta 𝐤_c.𝐫}|\chi _a^{(+)}(𝐤_a,𝐫).`$ (11)
The first term, known as the form factor, depends on the structure of the projectile through its ground state wave function $`\varphi _a^\mathrm{}m(𝐫_\mathrm{𝟏})`$. The second term involves the dynamics of the reaction. This amplitude (which will be referred to as the BT amplitude), used originally to study the deuteron breakup at sub-Coulomb energies , was applied to the calculations of the Coulomb breakup of halo nuclei in . This approximation, which allows the application of the theory to non-$`s`$ – wave projectiles, may seem to be justified if the c.m of the $`b+c`$ system is shifted towards $`b`$ (which is indeed the case if $`m_bm_c`$). However, the neglected piece of $`𝐫_i`$ ($`\alpha 𝐫_1`$) occurs in association with the wave vector $`𝐤_a`$, whose magnitude may not be negligible for the reactions at the higher beam energies. Therefore, contribution coming to the amplitude from the neglected part of $`𝐫_i`$ may still be substantial.
An approximate way of taking into account the finite range effects in the post form DWBA theory is provided by the local momentum approximation . The attractive feature of this approximation is that it leads to the factorization of the amplitude $`\beta _\mathrm{}m`$ similar to that obtained in the BT case. We use this approximation to write the Coulomb distorted wave of particle $`b`$ in the final channel as
$`\chi _b^{()}(𝐤_b,𝐫)`$ $`=`$ $`e^{i\alpha 𝐊.𝐫_1}\chi _b^{()}(𝐤_b,𝐫_i).`$ (12)
Eq. (12) represents an exact Taylor series expansion about $`𝐫_i`$ if $`𝐊(=i_{𝐫_i})`$ is treated exactly. However, this is not done in the LMA scheme. Instead, the magnitude of the local momentum here is taken to be
$`K(R)=\sqrt{{\displaystyle \frac{2m}{\mathrm{}^2}}(EV(R))},`$ (13)
where $`m`$ is the reduced mass of the $`bt`$ system, $`E`$ is the energy of particle $`b`$ relative to the target in the c.m. system and $`V(R)`$ is the Coulomb potential between $`b`$ and the target at a distance $`R`$. Therefore, the local momentum $`𝐊`$ is evaluated at some distance $`R`$, and its magnitude is held fixed for all the values of $`𝐫`$. As is discussed in appendix A, the results of our calculations are almost independent of the choice of the direction the local momentum. Therefore, we have taken the directions of $`𝐊`$ and $`𝐤_𝐛`$ to be the same in all the calculations presented in this paper. It may be noted that in the calculations presented in Ref. , the LMA was applied to the Coulomb distorted wave of the projectile channel, where making a choice of the direction of the local momentum is some what complicated due to the fact that the directions of both the fragments in the final channel has to be brought into the consideration. Detailed discussion on the validity of the local momentum approximation is presented in appendix A.
Substituting Eq. (12) into Eq. (8) we get the following factorized form of the amplitude
$`\widehat{\mathrm{}}\beta _\mathrm{}m^{FRDWBA}`$ $`=`$ $`e^{i(\gamma 𝐤_c\alpha 𝐊).𝐫_1}|V_{bc}|\varphi _a^\mathrm{}m(𝐫_1)`$ (14)
$`\times `$ $`\chi _b^{()}(𝐤_b,𝐫_i)e^{i\delta 𝐤_c.𝐫_i}|\chi _a^{(+)}(𝐤_a,𝐫_i).`$
Eq. (14) (which will be referred to as the FRDWBA amplitude in the following) looks like Eq (11) of the BT theory but with the very important difference that the form factor is now evaluated at the momentum transfer ($`\gamma 𝐤_c\alpha 𝐊`$), and not at the valence particle momentum $`𝐤_𝐜`$. The two momenta could be quite different for cases of interest in this paper. The second term, involving the dynamics part of the reaction, is the same in both the cases. Therefore, the breakup amplitude obtained in BT approximation differs from that of FRDWBA by a factor
$`F_r={\displaystyle \frac{\beta _\mathrm{}m^{BT}}{\beta _\mathrm{}m^{FRDWBA}}}={\displaystyle \frac{e^{i𝐤_c.𝐫_1}|V_{bc}|\varphi _a^\mathrm{}m(𝐫_1)}{e^{i(\beta 𝐤_c\alpha 𝐊).𝐫_1}|V_{bc}|\varphi _a^\mathrm{}m(𝐫_1)}}`$ (15)
Recently, a theory of the Coulomb breakup has been developed within an adiabatic (AD) model , where one assumes that the excitation of the projectile is such that the relative energy ($`E_{bc}`$) of the $`bc`$ system is much smaller than the total incident energy, which allows $`E_{bc}`$ to be replaced by the constant separation energy of the fragments in the projectile ground state. It was shown in that under these conditions the wave function $`\mathrm{\Psi }_a^{(+)}(\xi _a,𝐫_1,𝐫_i)`$ has an exact solution as given below
$`\mathrm{\Psi }_a^{(+),AD}(\xi _a,𝐫_1,𝐫_i)`$ $`=`$ $`\mathrm{\Phi }_a(\xi _a,𝐫_1)e^{i\alpha 𝐤_a.𝐫_1}\chi _a^{(+)}(𝐤_a,𝐫)`$ (16)
Substituting $`\mathrm{\Psi }_a^{(+),AD}`$ for $`\mathrm{\Psi }_a^{(+)}`$ in Eq. (3) leads to the reduced amplitude:
$`\widehat{\mathrm{}}\beta _\mathrm{}m^{AD}`$ $`=`$ $`e^{i(𝐤_c\alpha 𝐤_a).𝐫_1}|V_{bc}|\varphi _a^\mathrm{}m(𝐫_1)\chi _b^{()}(𝐤_b,𝐫)e^{i\delta 𝐤_c.𝐫}|\chi _a^{(+)}(𝐤_a,𝐫).`$ (17)
This amplitude differs from those of BT as well as FRDWBA only in the form factor part (first term), which is evaluated here at the momentum transfer $`(𝐤_c\alpha 𝐤_a)`$. Eq. (17) can also be obtained in the DWBA model by making a local momentum approximation to the Coulomb distorted wave in the initial channel of the reaction, and by evaluating the local momentum at $`R=\mathrm{}`$ with its direction being the same as that of the projectile . The adiabatic model does not make the weak coupling approximation of the DWBA. However, it necessarily requires one of the fragments (in this case $`c`$) to be chargeless. In contrast, the DWBA with the LMA could, in principle, be applied to cases where both the fragments $`b`$ and $`c`$ are charged (of course the dynamical part of the amplitude can not be expressed in terms of the simple bremsstrahlung integral in this case as discussed below). While the effect of nuclear breakup in the adiabatic model description of elastic scattering of the loosely bound projectile has been calculated in Ref. , the nuclear part of the amplitude for breakup reactions is yet to be calculated within this model. At the same time, the nuclear breakup cross section has not been calculated within the FRDWBA theory also (although it is straightforward to do so). These calculations will be presented in a future publication.
The triple differential cross section of the reaction is given by
$`{\displaystyle \frac{d^3\sigma }{dE_bd\mathrm{\Omega }_bd\mathrm{\Omega }_c}}`$ $`=`$ $`{\displaystyle \frac{2\pi }{\mathrm{}v_a}}\rho (E_b,\mathrm{\Omega }_b,\mathrm{\Omega }_c){\displaystyle \underset{\mathrm{}m}{}}|\beta _\mathrm{}m|^2,`$ (18)
where $`\rho (E_b,\mathrm{\Omega }_b,\mathrm{\Omega }_c)`$ is the appropriate three-body phase space factor, given by
$`\rho (E_b,\mathrm{\Omega }_b,\mathrm{\Omega }_c)`$ $`=`$ $`{\displaystyle \frac{h^6m_bm_cm_tp_bp_c}{m_t+m_cm_c\frac{𝐤_𝐜.(𝐤_𝐚𝐤_𝐛)}{k_c^2}}},`$ (19)
with $`𝐤_a,𝐤_b`$ and $`𝐤_c`$ being evaluated in the appropriate frame of reference. $`v_a`$ is the relative velocity of the projectile in the initial channel. $`p`$ in Eq. (19) are the linear momenta which are related to wave numbers $`k`$ by $`p=\mathrm{}k`$.
Substituting the following expressions for the Coulomb distorted waves
$`\chi _b^{()}(𝐤_b,𝐫_i)`$ $`=`$ $`e^{\pi \eta _b/2}\mathrm{\Gamma }(1+i\eta _b)e^{i𝐤_b.𝐫_i}{}_{1}{}^{}F_{1}^{}(i\eta _b,1,i(k_br_i+𝐤_b.𝐫_i)),`$ (20)
$`\chi _a^{(+)}(𝐤_a,𝐫_i)`$ $`=`$ $`e^{\pi \eta _a/2}\mathrm{\Gamma }(1+i\eta _a)e^{i𝐤_a.𝐫_i}{}_{1}{}^{}F_{1}^{}(i\eta _a,1,i(k_ar_i𝐤_a.𝐫_i))`$ (21)
in Eq. (14), one gets for the triple differential cross section
$`{\displaystyle \frac{d^3\sigma }{dE_bd\mathrm{\Omega }_bd\mathrm{\Omega }_c}}={\displaystyle \frac{2\pi }{\mathrm{}v_a}}\rho (E_b,\mathrm{\Omega }_b,\mathrm{\Omega }_c){\displaystyle \frac{4\pi ^2\eta _a\eta _b}{(e^{2\pi \eta _b}1)(e^{2\pi \eta _a}1)}}|I|^24\pi {\displaystyle \underset{\mathrm{}}{}}|Z_{\mathrm{}}|^2.`$ (22)
In Eqs. (20 – 22) $`\eta `$’s are the Coulomb parameters for the concerned particles. In Eq. (22) $`I`$ is the Bremsstrahlung integral which can be evaluated in a closed form:
$`I`$ $`=`$ $`i[B(0)\left({\displaystyle \frac{dD}{dx}}\right)_{x=0}(\eta _a\eta _b)_2F_1(1i\eta _a,1i\eta _b;2;D(0))`$ (23)
$`+`$ $`\left({\displaystyle \frac{dB}{dx}}\right)_{x=0}{}_{2}{}^{}F_{1}^{}(i\eta _a,i\eta _b;1;D(0))]`$
where
$`B(x)={\displaystyle \frac{4\pi }{k^{2(i\eta _a+i\eta _b+1)}}}\left[(k^22𝐤.𝐤_a2xk_a)^{i\eta _a}(k^22𝐤.𝐤_b2xk_b)^{i\eta _b}\right],`$ (24)
$`D(x)={\displaystyle \frac{2k^2(k_ak_b+𝐤_a.𝐤_b)4(𝐤.𝐤_a+xk_a)(𝐤.𝐤_b+xk_b)}{(k^22𝐤.𝐤_a2xk_a)(k^22𝐤.𝐤_b2xk_b)}}`$ (25)
with $`𝐤=𝐤_a𝐤_b\delta 𝐤_c`$. $`Z_{\mathrm{}}`$ contains the projectile structure information and is given by
$`Z_{\mathrm{}}={\displaystyle 𝑑r_1r_1^2j_{\mathrm{}}(k_1r_1)V_{bc}(𝐫_1)u_{\mathrm{}}(r_1)},`$ (26)
with $`k_1=|\gamma 𝐤_c\alpha 𝐊|`$.
## 3 Calculations and discussions
Apart from the distance at which local momentum is calculated (which is taken to be 10 $`fm`$ as discussed in appendix A) and its direction (described already in the previous section), the only other input to our calculations is the radial part of the projectile ground state wave function. The forms for this as chosen by us for the projectiles considered in this paper are discussed in the following.
### 3.1 Calculations for reactions induced by <sup>11</sup>Be
For the ground state of <sup>11</sup>Be, we have considered the following configurations : (a) a $`s`$ – wave valence neutron coupled to $`0^+`$ <sup>10</sup>Be core (<sup>10</sup>Be$`(0^+)1s_{1/2}\nu `$) with a one-neutron separation energy ($`S_{n^{10}Be}`$) of 504 keV and a spectroscopic factor (SF) of 0.74 , (b) a $`d`$ – wave valence neutron coupled to $`2^+`$ <sup>10</sup>Be core (<sup>10</sup>Be$`(2^+)0d_{5/2}\nu `$) with a binding energy of 3.872 MeV (which is the sum of energy of the excited 2<sup>+</sup> core (3.368 MeV) and $`S_{n^{10}Be}`$) and (c) an admixture of these two configurations with spectroscopic factors of 0.74 and 0.17 respectively . In each case, the single particle wave function, is constructed by assuming the valence neutron-<sup>10</sup>Be interaction to be of Woods-Saxon type whose depth is adjusted to reproduce the corresponding value of the binding energy with fixed values of the radius and diffuseness parameters (taken to be 1.15 $`fm`$ and 0.5 $`fm`$ respectively). For configuration (a), this gives a potential depth of 71.03 MeV, a root mean square (rms) radius for the valence neutron of 6.7 $`fm`$ and a rms radius for <sup>11</sup>Be of 2.91 $`fm`$ when the size of the <sup>10</sup>Be core is taken to be 2.28 $`fm`$ . In some cases we have also used a wave function for the $`s`$ – state of <sup>11</sup>Be calculated within a dynamical core polarization (DCP) model .
#### 3.1.1 Energy distributions of the fragments
In Fig. 2, we present the results of our FRDWBA calculations for the double differential cross section ($`d^2\sigma /dE_nd\mathrm{\Omega }_n`$) as a function of the neutron energy at the neutron angle of $`1^{}`$, in the breakup of <sup>11</sup>Be on a gold target at the beam energy of 44 MeV/nucleon. The experimental data are taken from . The results obtained with configurations (a), (b) and (c) for the ground state of <sup>11</sup>Be are shown by dotted, dashed and full lines respectively. The results of configuration (b) are plotted after multiplying the actual numbers by a factor of 1000. The cross sections obtained with configurations (a) and (c) are indistinguishable from each other. Thus these two configurations produce almost identical results for the Coulomb dissociation of <sup>11</sup>Be. In the following, we have used configuration (a) for the ground state state of <sup>11</sup>Be in all the calculations.
In Fig. 3, we present a comparison of our calculation with the data (taken from ) for the energy distribution of the neutron emitted in the same reaction as in Fig. 2 at two forward angles. Calculations performed within both FRDWBA and AD model are shown in this figure. The same configuration for the <sup>11</sup>Be ground state has been used in both the cases. The beam energy in this experiment varies between 36.9 – 44.1 MeV/nucleon. To take into account this spread, we have performed calculations at both its upper (44 MeV/nucleon) (solid line) and lower ends (37 MeV/nucleon) (dotted line). Although these data have large statistical errors, the calculations performed at 44 MeV/nucleon are in better agreement with the experimental values. Thus our calculations may serve to remove the uncertainty in the data in this regard. It should also be noted that the AD model calculations over-predict the experimental cross sections in the peak region.
As discussed earlier, the cross sections obtained in the BT theory differs from that of the FRDWBA by the modulus square of the factor $`F_r`$ as defined by Eq. (15). In fig. 4, we have shown the beam energy dependence of $`|F_r|^2`$ for the same reaction as in Fig. 2, for a set of forward angles of the outgoing fragments ($`\theta _b=1^{}`$, $`\theta _c=1^{}`$ and $`\varphi _c=1^{}`$). We can see that this quantity is close to unity only at the sub-Coulomb beam energies (of course at higher beam energies it crosses twice the line representing the value 1). Therefore, the BT and FRDWBA calculations are expected to produce similar results at very low incident energies. Depending upon the beam energy, the BT results can be larger or smaller than those of the FRDWBA.
A comparison of the results of the FRDWBA (solid line) and the BT (dotted line) calculations for the same reaction as in Fig. 2 is presented in Fig. 5. The bombarding energy in this case is 44 MeV/nucleon. We can see that the BT results are larger than those of the FRDWBA in almost entire region of neutron energies. This is to be expected from the results shown in Fig. 4, where the quantity $`|F_r|^2`$ is larger than unity at this beam energy. It may be noted that the FRDWBA results provide a reasonable good description of the data, particularly in the peak region while BT calculations overestimate them.
In Fig. 6, we show the energy distributions of the charged fragment <sup>10</sup>Be, calculated within the FRDWBA (upper part) and the AD model (lower part) for the breakup of <sup>11</sup>Be on the Au target at the beam energy of 42 MeV/nucleon at three angles of 1, 2 and 3. A noteworthy feature of these results is that the peak position in both the calculations is very close to the energy corresponding to the beam velocity. This suggests that the charged fragment gets almost no post acceleration in the final channel. This feature of the breakup of halo nuclei (which is in contrast to the case of the breakup of stable isotopes ), was first noted in Ref. , and was corroborated later on by the authors of Refs. . The reason for not observing the post acceleration effects, as put forward by authors of Ref. , is that due to their very small binding energies the halo nuclei break up far away from the distance of closest approach, which reduces greatly the effect of Coulomb repulsion on the charged outgoing fragment. This argument has received support from a recent calculation of the Coulomb dissociation potential of the deuteron which is shown to have a considerably large value even outside the charge density of the target nucleus. A separate reason for not observing this effect has been put forward by the authors of Ref. , according to which these effects are small as the collision time is much less than the characteristic time for the disintegration of the halo.
However, these quasiclassical arguments have been questioned by the authors of Ref. . Esbensen, Bertsch and Bertulani , who include the higher order processes in the Coulomb dissociation of halo nuclei <sup>11</sup>Li and <sup>11</sup>Be by solving the three-dimensional time-dependent Schrödinger equation, and find a magnitude for the post acceleration effect which is quite appreciable for the <sup>9</sup>Li fragment in the breakup of <sup>11</sup>Li and relatively somewhat smaller for the <sup>10</sup>Be fragment in the breakup of <sup>11</sup>Be. In the semiclassical calculations of Baur, Bertulani and Kalassa , where breakup is assumed to take place at some classical distance, it is predicted that post acceleration effects should manifest itself in the increase of the average momentum of the charged fragment with the scattering angle. Indeed, the earlier measurements of Nakamura et al. is consistent with this observation.
In a recent experiment , the momentum of the <sup>10</sup>Be core from the breakup of <sup>11</sup>Be has been measured with sufficient precision to verify the previously reported post acceleration effect. We have made a comparison of our FRDWBA calculations with this data in Fig. 7, where we show the calculated and experimental <sup>10</sup>Be average momenta (defined as $`p_b\frac{d^2\sigma }{dp_bd\mathrm{\Omega }_b}/\frac{d^2\sigma }{dp_bd\mathrm{\Omega }_b}`$) as a function of laboratory angle. In the semiclassical picture of Ref. , the impact parameter decreases with the increase of the scattering angle, thereby making the Coulomb repulsion effects on the charged fragment stronger. Therefore, the post acceleration should show up in the increase of this average momentum with the increase of the scattering angle. However, in both the experimental data of as well as our calculations no such increase is observed. Thus, neither the data of Ref. nor our calculations support the results found in Ref. . Therefore, the semiclassical picture presented in should be viewed with caution.
#### 3.1.2 Neutron Angular Distribution
The measured neutron angular distribution in the exclusive <sup>11</sup>Be + $`A`$ $``$ <sup>10</sup>Be + n +$`A`$ reaction below the grazing angle is very narrow and is shown to be dominated by the Coulomb breakup process. This reflects the narrow width of the transverse momentum distribution of the valence neutron in the ground state of <sup>11</sup>Be, which is consistent with the presence of a neutron halo structure in <sup>11</sup>Be. In the top half of Fig. 8, we compare the calculated and measured exclusive neutron angular distribution $`d\sigma /d\mathrm{\Omega }_n`$ as a function of the neutron angle $`\theta _n`$ for the above reaction on a Au target at the beam energy of 41 MeV/nucleon. Calculations (where integrations over the core fragment energy is done in the range of 390 to 430 MeV, which contributes most to the cross section) performed within the FRDWBA (solid line) and the AD model (dotted line) using the same configuration for the <sup>11</sup>Be ground state are shown in this figure. Also shown (dashed line) here is the FRDWBA calculation performed with the <sup>11</sup>Be ground state wave function obtained in the DCP model. We note that while the FRDWBA and the AD model results agree with each other well below 12, the difference between the two models starts becoming prominent as the angle increases beyond this value, with the latter developing a dip around 25. At small neutron angles, the FRDWBA calculation done with the DCP wave function overestimates somewhat the measured neutron angular distribution.
To understand the origin of the dip in the AD model calculations of the neutron angular distribution, we have plotted the quantity $`G(=d\theta _b\mathrm{sin}\theta _b|Z_{\mathrm{}}|^2`$) as a function of $`\theta _n`$ in the lower half of Fig. 8 for both the FRDWBA (solid line) and the AD (dashed line) cases. The energy of fragment $`b`$ corresponds to the beam velocity. We can see that in the AD case $`G`$ becomes very small around 25, which corresponds to a node in its form factor at the momentum transfer related to this angle.
#### 3.1.3 Relative energy spectrum of fragments
The relative energy spectrum for the breakup of <sup>11</sup>Be on a Pb target at 72 MeV/nucleon is shown in Fig. 9. The top half shows the results obtained with the FRDWBA (solid line) and the AD model (dotted line) using configuration (a) with the single particle wave function for the <sup>11</sup>Be ground state, while the bottom half depicts the FRDWBA results obtained with the single particle and DCP wave functions. We see that, while both the FRDWBA and the AD model calculations reproduce the peak value of the data well, the FRDWBA calculations done with the DCP wave function overestimate it. On the other hand, none of the calculations is able to explain the data at higher relative energies. This can be attributed to the fact that nuclear breakup effects, which can contribute substantially at higher relative energies (for $`E_{rel}`$ $`>`$ 0.6 MeV), are not included in these calculations. Of course, the authors of Ref. claim that their data have been corrected for these contributions. However, the procedure adopted by them for this purpose is inadequate. They obtained the nuclear breakup contribution on the Pb target, by scaling the cross sections measured on a carbon target. This scaling procedure is unlikely to be accurate for reactions induced by halo nuclei due to the presence of a long tail in their ground state. In a full quantum mechanical theory, both Coulomb and nuclear breakup contributions should be calculated on the same footing and corresponding amplitudes should be added coherently to get the cross sections.
Calculations done using the BT theory (dashed line in the upper part of Fig. 9) underestimates the data considerably. This difference between the FRDWBA and the BT results can again be traced to the behavior of $`|F_r|^2`$ in Fig. 4, which is smaller than unity at the beam energy of 72 MeV/nucleon of this reaction.
We would like to remark that our post form FRDWBA results for the pure Coulomb breakup contribution to the relative energy spectra for the <sup>11</sup>Be agrees quite well with a recent calculation of the breakup reaction in a non-perturbative approach where the time evolution of the projectile is calculated by solving a time-dependent Schrödinger equation . On the other hand, similar perturbative calculations performed previously overestimate the data in the peak region.
#### 3.1.4 Momentum distribution of the core
The neutron halo structure is reflected in the narrow width of the parallel momentum distribution (PMD) of the charged breakup fragments emitted in breakup reactions induced by the halo nuclei. This is because the PMD has been found to be least affected by the reaction mechanism and therefore, a narrow PMD can be related to a long tail in the matter distribution in the coordinate space via Heisenberg’s uncertainty principle.
In Fig. 10 we present the PMD of the <sup>10</sup>Be fragment emitted in the breakup of <sup>11</sup>Be on U and Ta targets at 63 MeV/nucleon beam energy. Calculations performed within both the FRDWBA and AD model formalisms using the configuration (a) are presented in this figure. The calculated cross sections are normalized to match the peak of the data points (which are given in arbitrary units) , the normalization constant being the same for both cases. The full width at half maximum ( FWHM) for the U and Ta targets are 44 MeV/c and 43 MeV/c respectively in both the FRDWBA and the AD cases. These agree well with the averaged experimental value of 43.6$`\pm `$1.1 MeV/c and also with those calculated in . The very narrow widths of the parallel momentum distributions signal the presence of a neutron halo structure in <sup>11</sup>Be.
It may be noted that the calculations performed with configuration (c) gives results identical to that obtained with configuration (a). The PMD calculated with a pure $`d`$ – wave configuration is too small in magnitude and too wide in width.
#### 3.1.5 Coulomb part of total one-neutron removal cross section
For the breakup of <sup>11</sup>Be on a Au target at the beam energy of 41 MeV/nucleon, the values of the pure Coulomb total one-neutron removal cross section ($`\sigma _n^C`$) are found to be 1.91 b and 1.94 b for the FRDWBA and the AD model respectively, using configuration (a) and the single particle wave function for the ground state of <sup>11</sup>Be. The corresponding values of $`\sigma _n^C`$ for the breakup of this projectile on the Pb target at the beam energy of 72 MeV/nucleon are 1.25 b and 1.29 barn respectively. The experimental values for the total one-neutron removal cross section ($`\sigma _n`$) for these two reactions are reported to be 2.5$`\pm `$0.5 b and 1.8 $`\pm `$ 0.4 b respectively. The difference between $`\sigma _n^C`$ and the experimental value of $`\sigma _n`$ can be attributed to the nuclear breakup effects. Incidentally, FRDWBA calculations performed with the DCP wave function for the <sup>11</sup>Be ground state leads to much larger values of 2.82 b and 1.76 b for $`\sigma _n^C`$ for the two cases respectively.
### 3.2 Results for <sup>19,15,17</sup>C
In this section,we shall compare results of our calculations with the available data on the breakup of the neutron rich carbon nuclei <sup>19,15,17</sup>C on heavier targets.
There is a large uncertainty in the value of the last neutron separation energy in <sup>19</sup>C ($`S_{n^{18}C}`$) with quoted values varying between 160 – 530 keV . It has recently been shown that most of the available data on the Coulomb dissociation of <sup>19</sup>C can be satisfactorily explained within the adiabatic model of Coulomb breakup with the one-neutron separation energy of 530 keV. The ground state spin-parity of <sup>19</sup>C has been quoted as $`1/2^+,3/2^+`$ and $`5/2^+`$ . The relativistic mean field (RMF) as well as shell model calculations using Warburton-Brown effective interaction predict the spin-parity of the ground state of this nucleus to be $`1/2^+`$.
We use single particle wave functions which are constructed by assuming a Woods-Saxon interaction between the valence neutron and the charged core whose depth (for fixed values of the radius and diffuseness parameters) is adjusted to reproduce the binding energies of the nuclei under investigation. The valence neutron binding energies, searched potential depths ($`V_{depth}`$) and calculated rms radii of the projectile with different configurations for the ground state for each isotope are summarized in Table 1.
In this work, we consider two situations: (i) different binding energies (530 keV and 160 keV) with the same configuration for <sup>19</sup>C ground state (<sup>18</sup>C $`(0^+)1s_{1/2}\nu `$), and (ii) different configurations (<sup>18</sup>C $`(0^+)1s_{1/2}\nu `$ and <sup>18</sup>C $`(0^+)0d_{5/2}\nu `$) for <sup>19</sup>C ground state with the same binding energy (530 keV). We have considered single particle wave functions in all the cases, except for the 160 keV case where we have additionally considered a DCP wave function for the $`s`$ – state . It may be noted that for all the single particle wave functions, we have used a spectroscopic factor of 1.0.
In Fig. 11, we present the PMD (calculated within the FRDWBA formalism) of the <sup>18</sup>C fragment in the breakup of <sup>19</sup>C on a Ta target at the beam energy of 88 MeV/nucleon. We have normalized the peaks of the calculated PMDs to that of the data (given in arbitrary units) (this also involves coinciding the position of maxima of the calculated and experimental PMDs). As can be seen from the upper part of this figure, the experimental data clearly favor $`S_{n^{18}C}`$ = 0.53 MeV with the $`s`$ – wave n-<sup>18</sup>C relative motion in the ground state of <sup>19</sup>C. The results obtained with the $`s`$ – wave configuration within the simple potential and DCP models (with the same value of $`S_{n^{18}C}`$) are similar to each other.
In the lower part of Fig. 11, we have shown the results obtained with the $`d`$ – wave relative motion for this system (with $`S_{n^{18}C}`$ = 0.53 MeV) and have compared it with that obtained with a $`s`$ – wave relative motion with the same value of the binding energy. As can be seen, the FWHM of the experimental PMD is grossly over-estimated by the $`d`$ – wave configuration. The calculated FWHM with the $`s`$ – state configuration (with $`S_{n^{18}C}`$ = 530 keV) is 40 MeV/c, which is in excellent agreement with the experimental value of 41$`\pm `$3 MeV/c . Thus these data favor a configuration <sup>18</sup>C($`0^+)1s_{1/2}`$, with a one-neutron separation energy of 0.530 MeV for the ground state of <sup>19</sup>C. These results are in agreement with those of Ref. . The narrow width of the PMD provides support to the presence of a one-neutron halo structure in <sup>19</sup>C.
The width of the PMD is not expected to change with beam energy . To see this, we have calculated the PMD of <sup>18</sup>C in the Coulomb breakup of <sup>19</sup>C on a heavy target (Ta) at the high beam energy of 914 MeV/nucleon. The corresponding results are shown in Fig. 12. The FWHM of the distribution in this case is 42 MeV/c which is similar to that obtained above at a lower beam energy. However, FWHM of the PMD of the <sup>18</sup>C fragment measured in the breakup of <sup>19</sup>C on a carbon target at the same beam energy, has been found to be of 69$`\pm `$4 MeV/c. It would be interesting, therefore, to repeat measurement at this energy with a heavier target to check our observation.
In Fig. 13, we have shown the results of our FRDWBA calculations for the the relative energy spectrum for the breakup of <sup>19</sup>C on Pb at 67 MeV/nucleon beam energy. The experimental data is taken from . The angular integration for the <sup>19</sup>C center of mass is done up to the grazing angle of 2.5. It can be seen that in this case also the best agreement with the data (near the peak position) is obtained with the $`s`$ – wave configuration with $`S_{n^{18}C}`$ = 0.53 MeV. Calculations done with the $`d`$ – state configuration for both 530 keV and 160 keV one-neutron separation energy fails to reproduce the data. At the same time, those performed with the $`s`$ – state configuration but with $`S_{n^{18}C}`$ = 0.16 MeV overestimates the data by at least an order of magnitude and also fail to reproduce the its peak position. However, as in the case of <sup>11</sup>Be, our calculations underestimate the relative energy spectrum for larger values of relative energies. In this case also, the proper consideration of the nuclear breakup effects is necessary to explain the data in this region. We would also like to remark that the correction for these effects made in the data by scaling the cross sections for the breakup of <sup>19</sup>C on a carbon target is unlikely to be accurate due to the same reasons as stated in case of <sup>11</sup>Be.
Table 2 summarizes the FRDWBA results of $`\sigma _n^C`$ for the breakup of <sup>19</sup>C on different heavy targets at several beam energies. The experimental value of $`\sigma _n`$ for the breakup of <sup>19</sup>C on Ta at the beam energy of 88 MeV/nucleon is 1.1$`\pm `$0.4 b . It is seen from this table that only with the $`s`$ – wave configuration of the <sup>19</sup>C ground state with $`S_{n^{18}C}`$ = 0.53 MeV, the calculated cross sections come closer to the experimental data. In this context it must be kept in mind that $`\sigma _n`$ also include contributions from the nuclear breakup effects, which is not included in our present calculations.
We next consider the breakup of <sup>15</sup>C which has a relatively larger value for the one-neutron separation energy (1.2181 MeV) and a ground state spin-parity of $`1/2^+`$ . This can be obtained from two configurations: a $`1s_{1/2}`$ neutron coupled to a <sup>14</sup>C $`(0^+)`$ core and a $`0d_{5/2}`$ neutron coupled to a <sup>14</sup>C $`(0^+)`$ core. One could also considered a <sup>14</sup>C $`(2^+)`$ core and $`0d_{5/2}`$ neutron coupling to get a $`1/2^+`$ ground state for <sup>15</sup>C, but it would raise the one-neutron separation energy to about 7.01 MeV, which is highly unfavorable for the formation of a halo. We, therefore, do not consider this configuration in our calculations.
In Fig. 14, we present the results of our FRDWBA calculations for the PMD of the <sup>14</sup>C fragment in the breakup of <sup>15</sup>C on the Ta target at the beam energy of 84 MeV/nucleon. The experimental data is taken from . As before, our calculations are normalized to the peak of the data. The $`s`$ – state configuration for the ground state of <sup>15</sup>C gives a FWHM of 62 MeV/c, while with the $`d`$ – state configuration it comes out to be 140 MeV/c. Therefore, the experimental value for the FWHM (67$`\pm `$1 MeV/c) favors the former configuration. Hence our results provide support to the existence of a halo structure in <sup>15</sup>C. This nucleus provides an example of the one-halo system with the largest one-neutron separation energy, known so far.
<sup>17</sup>C has a lower one-neutron separation energy (729 keV) as compared to that of <sup>15</sup>C. It would be interesting to see if it also has a halo structure, which seems probable if one considers only the binding energies. The quoted ground state spin-parities for this nucleus are $`1/2^+,3/2^+`$ and $`5/2^+`$ . RMF calculations predict it to have a value of $`3/2^+`$. We consider four possible ground state configurations for this nucleus and calculate the parallel momentum distributions of the <sup>16</sup>C fragment in the breakup of <sup>17</sup>C on a Ta target at 84 MeV/nucleon beam energy within our FRDWBA formalism. The FWHMs of the PMD obtained with different configurations are listed in Table 3.
It is evident from this table that the $`s`$ – state configurations predict a narrow width for the PMD, providing support to the existence of a halo structure in this nucleus. The experimental data , however, is available only for the breakup of <sup>17</sup>C on a light target (Be) at 84 MeV/nucleon, which gives a FWHM of 145$`\pm `$5 MeV/c. Since the PMD is mostly unaffected by the reaction mechanism , it is quite likely that the experimental FWHM will be the same also for the breakup of this nucleus on a heavier target. Therefore, the results shown in Table 3 seem to provide support to a $`d`$ – wave configuration for the ground state of <sup>17</sup>C . Hence, the existence of a one-neutron halo structure is quite improbable in <sup>17</sup>C. However, to arrive at a more definite conclusion in this regard the data on the breakup of <sup>17</sup>C on a heavy target is quite desirable.
## 4 Summary and conclusions
In this paper, we have performed calculations for the Coulomb breakup of the neutron rich nuclei <sup>11</sup>Be and <sup>15,17,19</sup>C, which have a single valence neutron loosely bound with a stable core. We used a theory developed within the framework of the post form distorted wave Born approximation where finite range effects have been included approximately by using a local momentum approximation on the Coulomb distorted wave of the outgoing charged fragment. Within this method, the breakup amplitude is expressed as a product of factors describing separately the projectile structure and the dynamics of the reaction. This factored form of the breakup amplitude can also be obtained within an adiabatic model which makes the approximation that the strongly excited core-valence particle relative energies in the Coulomb breakup are small. However, unlike the post form DWBA, the adiabatic model does not use the weak coupling approximation to describe the center of mass motion of the fragments with respect to the target.
Both these theories allow the use of realistic wave functions for the ground state of the projectile. Furthermore, unlike the semiclassical and quantum mechanical theories using the zero range approximation which can be applied only to $`s`$ – wave projectiles, these methods are applicable to projectiles with any relative orbital angular momentum structure between their fragments. This provides an opportunity to probe the structure of the ground state of the projectile, by comparing the predictions of these theories with the data for the breakup observable. We have calculated, the energy, angular and parallel momentum distributions of the fragments emitted in the breakup reaction of these nuclei using different configurations for their ground state, By making comparisons of the calculated cross sections with the available experimental data an effort has been made to put constraints on their ground state structure.
All the observable calculated by us are sensitive to the ground state configuration of the projectile. We find that for <sup>11</sup>Be, a $`s`$ – wave configuration (<sup>10</sup>Be$`(0^+)1s_{1/2}\nu `$), with a spectroscopic factor of 0.74 for its ground state provides best agreement with the experimental data in all the cases. In our study, it is not possible to distinguish between this configuration and the one proposed recently where there is an admixture of the $`s`$ – wave and a $`d`$ – wave configuration, (<sup>10</sup>Be$`(2^+)0d_{5/2}\nu `$), with spectroscopic factors of 0.74 and 0.17 respectively for the ground state of this nucleus. For almost all the observables, there is a general agreement between the FRDWBA and adiabatic model results even in the absolute magnitude in the region where Coulomb breakup is expected to be dominant mode (ie. below the grazing angle). This provides additional support to our choice of the parameters associated with the local momentum in our FRDWBA calculations. It may be noted that the approximation of Baur and Trautmann, which also leads to the factored form for the breakup amplitude, gives results which are very different from those obtained with the FRDWBA and AD model formalisms. The BT approximation fails to explain the data in most of the cases studied here.
In the case of neutron angular distributions for the breakup of <sup>11</sup>Be on the gold target at the beam energy of 41 MeV, there is a dip around 25 in the adiabatic model calculations, which is not seen in the corresponding FRDWBA results. It may be noted that this region was excluded in the results shown in Ref. . This dip can be traced back to the fact that the form factor in the adiabatic model has a node at the momentum transfer corresponding to this angle. In any case, this region of the angular distribution is expected to get substantial contribution from the nuclear breakup effects (and also from the Coulomb-nuclear interference terms). Therefore, full implication of this dip can become clear only after nuclear breakup effects are included in these models.
For the <sup>19</sup>C case, the results for the PMD of <sup>18</sup>C and the relative energy spectrum of the $`n`$ \+ <sup>18</sup>C system show that the most probable ground state configuration of <sup>19</sup>C is (<sup>18</sup>C $`(0^+)1s_{1/2}\nu `$) with a one-neutron separation energy of 530 keV and a spectroscopic factor of 1. Our FRDWBA calculations agree well with the experimental data of the MSU group . By performing the calculations at GSI energies we note that the width of the PMD is independent of the beam energy. It would be interesting to perform the GSI experiment with a heavier target to check this prediction. Our FRDWBA results on <sup>19</sup>C are in excellent agreement with those of the adiabatic model .
We find that the most probable configuration for <sup>15</sup>C is a $`s`$ – wave valence neutron coupled to the <sup>14</sup>C core and that for <sup>17</sup>C is a $`d`$ – wave valence neutron coupled to a <sup>16</sup>C core. Both the experimental and the calculated FWHM of the PMD for the <sup>14</sup>C core in the breakup of <sup>15</sup>C are small and they agree well with each other. This provides support to the existence of a one-neutron halo structure in <sup>15</sup>C. On the other hand, in the case of <sup>17</sup>C the value of this quantity for the <sup>16</sup>C core is probably closer to that of a stable isotope. Therefore the existence of a halo structure in <sup>17</sup>C appears to be unlikely. Interestingly the one-neutron separation energies of <sup>15</sup>C and <sup>17</sup>C are 1.2181 and 0.729 MeV respectively. So the binding energy of the valence neutron as well as its configuration with respect to the core together decide whether a nucleus has halo properties or not.
## 5 ACKNOWLEDGMENTS
The authors are thankful to Horst Lenske for providing them the wave functions of the dynamical polarization model for various projectiles. One of the authors (RS) wishes to acknowledge several fruitful discussions with Dr. M. A. Nagarajan on the local momentum approximation method.
## Appendix A Validity of the local momentum approximation
As discussed in , a condition for the validity of the local momentum approximation is that the quantity
$`\eta (r)={\displaystyle \frac{1}{2}}K(r)|dK(r)/dr|^1`$ (27)
evaluated at a representative distance $`R`$ should be much larger than $`R_{bc}`$, which is of the order of the range of the interaction $`V_{bc}`$. To check this, we show in Fig. A.1, $`\eta (r)`$ (upper half) and $`K(r)`$ (the magnitude of the local momentum) (lower half) as a function of $`r`$, for the breakup reactions <sup>19</sup>C + Ta $``$ <sup>18</sup>C + $`n`$ \+ Ta at the beam energy of 88 MeV/nucleon (left side) and <sup>11</sup>Be + Au $``$ <sup>10</sup>Be + $`n`$ \+ Au at the beam energy of 41 MeV/nucleon (right side). We see that for $`r`$ $`>`$ 8 $`fm`$, $`\eta (r)`$ is several orders of magnitude larger than rms radii of the halo in both the cases. Therefore, the above condition is well satisfied.
From the lower half of Fig. A.1, we note that the value of $`K(r)`$ remains constant for distances larger than 10 $`fm`$. Due to the peripheral nature of the breakup reaction, this region contributes maximum to the cross section. Therefore, our choice of a constant magnitude for the local momentum evaluated at 10 $`fm`$ is well justified. In fact, we noted that as $`R`$ is increased from 5 to 10 $`fm`$ the calculated cross sections vary by at the most 10$`\%`$, and with a further increase the variation is less than 1$`\%`$, in all the cases considered in this paper.
In the application of the LMA in the description of the heavy ion induced transfer reactions, it was noted that the calculated cross sections were more or less unaffected by the choice of the direction of the local momentum. However, in Ref. , some dependence of the breakup cross sections on this direction has been reported. To study the sensitivity of our FRDWBA results on the direction of $`𝐊`$, we have performed calculations for three cases where we take the angles of the local momentum ($`d_1`$) parallel to those of $`𝐤_b`$, ($`d_2`$) parallel to the direction corresponding to the half of the angles of $`𝐤_b`$ and ($`d_3`$) parallel to the beam direction (zero angles).
In Table A.1, we show the results for $`\sigma _n^C`$ for the breakup of <sup>11</sup>Be and <sup>19</sup>C (on targets and at beam energies as indicated therein), for these three choices of the direction of $`𝐊`$. We see that between the cases $`d_1`$ to $`d_3`$, the variation in the values of $`\sigma _n^C`$ is less than 5$`\%`$ for <sup>11</sup>Be, and less than 2$`\%`$ for <sup>19</sup>C. In part A of Fig. A.2, we show the energy distribution of the neutron for the same reaction as described in Fig. 2 The results obtained with cases $`d_1`$, $`d_2`$ and $`d_3`$ are shown by solid, dashed and dotted lines respectively. We note that energy distributions calculated with these choices differ from each other only in the peak region; the variation between them is of the order of only 5$`\%`$. The results for the neutron angular distribution (for the reaction reported in Fig. 8) is of the similar nature.
In part B of Fig. A.2, we show the relative energy spectrum for the same reaction as shown in Fig. (13), for the choices ($`d_1`$), ($`d_2`$) and ($`d_3`$). In this case we observe almost no variation in the calculated cross sections. Similarly, We have noted no dependence of the calculated widths of the parallel momentum distributions of heavy fragments on the direction of $`𝐊`$, in all the reactions investigated in this paper. Therefore, the dependence of various cross sections for the reaction studied in this paper, on the direction of the local momentum is either very minor or almost negligible. The measurements done so far are not able to distinguish the small differences that we see here in some cases. Therefore, we have performed all our calculations in this paper by using $`\widehat{𝐊}`$ = $`\widehat{𝐤}_b`$.
|
warning/0003/math0003217.html
|
ar5iv
|
text
|
# Upper bound for the Weil-Petersson volumes
## 0. Introduction
After Wolpert in computed the cohomology of the moduli space of Riemann surfaces as a graded vector space, the question of computing the cohomology ring structure (aka the intersection theory) on the moduli arose. The problem has been intensively studied since then. Witten’s paper is a good source for available techniques and ideas. Witten’s conjecture, which later became Kontsevich’s theorem , shows that the intersection numbers satisfy a certain KdV equation.
However, the problem of getting explicit numerical results still remains, as the recursive computations become exceedingly complicated as the genus and number of punctures grow. Carel Faber computed some low-genus intersection numbers in and has obtained numerous results in other papers.
Recently Zograf and Manin and Zograf have obtained quite explicit generating functions for the Weil-Petersson volumes, and computed the asymptotics of the volume growth for genus being fixed and the number of punctures growing to infinity.
In this paper we use a completely different set of tools, namely the decorated Teichmüller theory, to obtain an explicit asymptotic upper bound for the Weil-Petersson volumes for a fixed number of punctures and the genus growing to infinity.
## 1. Decorated Teichmüller Theory
Let us recall the notations and relevant constructions. All of these come from Penner’s work : this is just a brief summary.
Let $`_{g,n}`$ denote the moduli space of Riemann surfaces of genus $`g`$ with $`n`$ punctures — it has complex dimension $`3g3+n`$. Let $`\omega _{WP}`$ denote the two-form of the Weil-Petersson scalar product on $`_{g,n}`$. It can be extended to a closed current on the Deligne-Mumford compactification $`\overline{_{g,n}}`$ of the moduli space. Taking its highest power produces a volume form on $`_{g,n}`$, integrating which over $`\overline{_{g,n}}`$ gives the Weil-Petersson volume $`\mathrm{vol}_{WP}(_{g,n})`$, which is the principal object of our study. We will only be concerned with hyperbolic punctured surfaces, i.e. the case when $`2g+n3`$ and $`n>0`$.
An ideal triangulation of a punctured surface is a triangulation of the surface with vertices only at punctures. We can straighten an ideal triangulation so that every edge of it is a geodesic arc for the hyperbolic metric on the surface. From Euler characteristics considerations it follows that such a triangulation consists of $`V:=4g4+2n`$ triangles, and has $`N:=6g6+3n`$ edges.
The decoration of a punctured Riemann surface is an addition of a horocycle around each puncture. More rigorously, on the uniformizing hyperbolic plane we take a horocycle around a preimage of a puncture, and consider its projection to the Riemann surface. For a decorated Riemann surface with an ideal geodesic triangulation define the $`\lambda `$-length of arc $`e`$ of the triangulation to be $`\lambda (e)=\sqrt{2e^\delta }`$, where $`\delta `$ is the (signed) hyperbolic distance from the point where $`e`$ intersects the horocycle around one its end to the point of intersection with the horocycle at the other end. It turns out (, theorem 3.1) that for a fixed triangulation the $`\lambda `$-lengths establish a homeomorphism of the decorated Teichmüller space $`\stackrel{~}{𝒯_{g,n}}`$ and $`_+^{6g6+3n}`$.
An embedded graph is a graph embedded in a Riemann surface. Combinatorially it can be represented as a usual graph endowed with a cyclic order of the edges around each vertex. Given an ideal triangulation of a surface, taking its Poincaré dual produces a trivalent embedded graph on the surface, which we denote $`\mathrm{\Gamma }`$. The graph $`\mathrm{\Gamma }`$ has $`N`$ edges and $`V`$ vertices. We define $`\lambda (e)`$, the $`\lambda `$-length of an edge $`e`$ of the graph, to be the $`\lambda `$-length of its Poincaré dual geodesic arc of the ideal triangulation.
## 2. Moduli space description
Our goal is to use Penner’s description of the moduli space in the $`\lambda `$-length coordinates to obtain an explicit upper bound of the Weil-Petersson volumes.
For any edge $`e`$ of the trivalent embedded graph $`\mathrm{\Gamma }`$ let $`f_i`$ and $`g_i`$ be the adjacent edges of the graph at $`i`$’s end of $`e`$. Then define the associated simplicial coordinate to be
$$X_e:=\underset{i=1}{\overset{2}{}}\frac{\lambda (f_i)}{\lambda (e)\lambda (g_i)}+\frac{\lambda (g_i)}{\lambda (e)\lambda (f_i)}\frac{\lambda (e)}{\lambda (f_i)\lambda (g_i)}.$$
Further, let $`\rho _i`$ be the sum of simplicial coordinates of the edges in a path around puncture number $`i`$, where the path has to go to “the next edge to the left” at each vertex in the cyclic order corresponding to the embedding of the graph. Notice that if we take any edge and start going to “the next edge to the left” from it, we will end up with a loop around some puncture. Since you can “go to the left” in two directions, i.e. since it matters in which way you start going, we have
$$\rho :=\underset{i=1}{\overset{n}{}}\rho _i=2\underset{i=1}{\overset{N}{}}X_{e_i}=2\frac{\lambda (e)}{\lambda (f)\lambda (g)},$$
where the sum is over all triples of edges having a common vertex, including the possible renamings of $`e`$, $`f`$ and $`g`$.
In these notations Penner proves the following result:
###### Fact 2.1 (, 3.2.1 and 3.4.3).
Let $`\omega `$ be a top-dimension differential form on $`_{g,n}`$. Then
(2.1)
$$_{\overline{}_{g,n}}\omega =\underset{[\mathrm{\Gamma }]}{}\frac{1}{\mathrm{Aut}\mathrm{\Gamma }}_{D(\mathrm{\Gamma })}\pi ^{}(\omega ),$$
where $`\pi :\stackrel{~}{𝒯}_{g,n}_{g,n}`$ is the forgetful projection, $`[\mathrm{\Gamma }]`$ denotes the isomorphism class of an embedded trivalent graph $`\mathrm{\Gamma }`$, and $`D(\mathrm{\Gamma })`$ is the domain in the decorated Teichmüller space given by
(2.2)
$$D(\mathrm{\Gamma })=\left\{\rho _i=1;X_e>0\right|e\mathrm{\Gamma },i=1\mathrm{}n\}$$
in the $`\lambda `$-coordinates corresponding to the embedded graph $`\mathrm{\Gamma }`$.
## 3. Weil-Petersson Volume form
In further computations for simplicity we drop the $`\lambda `$’s and simply write $`e`$ for $`\lambda (e)`$, if no confusion is possible.
In $`\lambda `$-length coordinates on the moduli the Weil-Petersson two-form is given by (, theorem A.2)
(3.1)
$$\omega _{WP}=2\underset{v\mathrm{\Gamma };e,f,gv}{}\frac{\mathrm{d}e}{e}\frac{\mathrm{d}f}{f}+\frac{\mathrm{d}f}{f}\frac{\mathrm{d}g}{g}+\frac{\mathrm{d}g}{g}\frac{\mathrm{d}e}{e}.$$
The Weil-Petersson volume form is the $`3g3+n`$’s external power of $`w_{WP}`$. In general, letting $`I`$ be a multi-index, and denoting the exclusion of factors by a hat, it would be the sum
(3.2)
$$\omega _{WP}^{(3g3+n)}=\underset{|I|=n}{}a_Id\mathrm{ln}\lambda _1\mathrm{}\widehat{d\mathrm{ln}\lambda _I}\mathrm{}d\mathrm{ln}\lambda _N.$$
###### Proposition 3.1.
In the above notations, $`|a_I|2^N`$.
###### Proof.
Use the expression (3.1) for the two-form to straightforwardly take an exterior power and compute the total number of summands of one kind (with fixed $`I`$). Suppose the product contains some $`d\mathrm{ln}e`$; then it must come in pair with one of the adjoining edges — let it be $`f_i`$ in the above notations, so there are four choices. Then $`g_i`$ must come (if it is in the sum) with one of the two edges at its other end, and then the other edge at its other end must come with still another, and so on. Thus we have one factor of four, and many factors of two. Each new factor of four appears if we encounter an edge in $`I`$, so the total number of summands of one kind is at most $`2^|I|2^{3g3+n}=2^{3g3+2n}`$. Recalling the $`2`$ in the two-form, $`|a_I|2^{3g3+2n}2^{3g3+n}=2^N`$. ∎
For the case of one puncture Penner explicitly computes the Weil-Petersson volume form to be (, theorem 6.1.2)
(3.3)
$$\omega _{WP}^{(3g2)}=\pm 2^{4g2}\underset{i=1}{\overset{N}{}}(1)^id\mathrm{ln}\lambda _1\mathrm{}\widehat{d\mathrm{ln}\lambda _i}\mathrm{}d\mathrm{ln}\lambda _N.$$
Thus we are dealing with a form which is singular when the $`\lambda `$ coordinates approach zero. However, this is not a problem:
###### Proposition 3.2.
In the domain of integration $`D(\mathrm{\Gamma })`$ (formula 2.1) for any edge $`e\mathrm{\Gamma }`$ we have $`\lambda (e)>4`$ in $`D`$.
###### Proof.
We can construct at most two paths around punctures “going to the left” including the edge $`e`$ — a path can be constructed by deciding at which end of $`e`$ we start building it. Let one of these paths go through edges $`f_1,e,f_2`$, and the other — through $`g_1,e,g_2`$; let $`\rho _i`$ and $`\rho _j`$ be the sums of the simplicial coordinates of the edges in these paths. Then using the formula for $`\rho `$ in terms of $`\alpha `$-lengths of sectors (, lemma 3.4.2) we see that
$$\rho _i>\frac{2}{e}\left(\frac{g_1}{f_1}+\frac{g_2}{f_2}\right)\mathrm{and}\rho _j>\frac{2}{e}\left(\frac{f_1}{g_1}+\frac{f_2}{g_2}\right).$$
Since $`\rho _i=\rho _j=1`$ in the domain $`D`$, we get
$$2>\frac{2}{e}\left(\frac{f_1}{g_1}+\frac{g_1}{f_1}+\frac{f_2}{g_2}+\frac{g_2}{f_2}\right)>\frac{8}{e},$$
so that $`e>4`$. If there were only one path going through $`e`$, i.e. if $`i=j`$, we would have $`1=\rho _i>\frac{8}{e}`$, and thus $`e>8`$, which is even better. ∎
Thus in the domain of integration $`D(\mathrm{\Gamma })`$ the $`\lambda `$-lengths are bounded below. However, $`D`$ can have limit points at infinities of $`\lambda `$\- lengths, and since the integral of $`d\mathrm{ln}x`$ does not converge at infinity, we have to deal with this problem in great detail, using the triangle inequality.
## 4. Triangle Inequality
The problem with the decorated Teichmüller theory is that the domain of intgeration $`D(\mathrm{\Gamma })`$ cannot be simply described in $`\lambda `$-lengths. Our success will come from the following observation:
###### Theorem 4.1 (Triangle inequality: , lemma 5.2).
Let $`e`$, $`f`$, and $`g`$ be three edges of the graph $`\mathrm{\Gamma }`$ having a common vertex. Then in the domain of integration $`D(\mathrm{\Gamma })`$ the triangle inequality between them holds:
(4.1)
$$ef+g.$$
###### Proof.
Assume for contradiction that $`e>f+g`$. Note that
(4.2)
$$\frac{g}{f}+\frac{f}{g}\frac{e^2}{fg}<2e>f+g$$
by clearing the denominators and extracting full squares. Similarly
(4.3)
$$\frac{e}{f}+\frac{f}{e}\frac{g^2}{ef}>2g<|ef|,$$
and both inequalities hold if $`e>f+g`$.
Denote by $`f_1`$ and $`g_1`$ the edges at the other end of $`e`$. From above it then follows that
$$0<eX_e=\frac{g_1}{f_1}+\frac{f_1}{g_1}\frac{e^2}{f_1g_1}+\frac{g}{f}+\frac{f}{g}\frac{e^2}{fg}<\frac{g_1}{f_1}+\frac{f_1}{g_1}\frac{e^2}{f_1g_1}2$$
using inequality (4.2). But this is just the inequality (4.3) for edges $`e`$, $`f_1`$ and $`g_1`$, and thus $`e<|f_1g_1|`$. Suppose $`f_1>g_1`$; then it follows that $`f_1>g_1+e`$, and we can apply a similar argument to the edges $`f_2`$ and $`g_2`$ at the other end of $`f_1`$ to obtain $`f_1<|f_2g_2|`$. Continuing this process inductively, and assuming at each step that $`f_i>g_i`$, we end up with an infinite strictly increasing sequence of edges $`e<f_1<\mathrm{}<f_n<\mathrm{}`$, which is rather hard to achieve on a finite graph. Thus assuming that one triangle inequality among $`\lambda `$-lengths fails, we have arrived at a contradiction. ∎
To demonstrate the power of the triangle inequality we prove a simple corollary:
###### Proposition 4.2.
Let $`\rho `$ be twice the sum of all simplicial coordinates of all edges in the graph $`\mathrm{\Gamma }`$, as before. Denote the minimal edge in the graph by $`\mu `$. If the triangle inequalities are satisfied, then $`\rho <\frac{8V}{\mu }`$.
###### Proof.
For any vertex $`v`$ denote by $`e_v`$, $`f_v`$ and $`g_v`$ the edges containing it, with $`e_v`$ being the maximal among the three. Then using the triangle inequalities we have
$$\frac{\rho }{2}=\underset{v}{}\frac{e_v}{f_vg_v}+\frac{f_v}{e_vg_v}+\frac{g_v}{e_vf_v}\underset{v}{}\frac{f_v+g_v}{f_vg_v}+\frac{1}{f_v}+\frac{1}{g_v}\underset{v}{}\frac{4}{\mu }$$
## 5. Stoke’s theorem
In section 3 we obtained an expression for the Weil-Petersson form in terms of $`\lambda `$-lengths. However, this expression has multiple summands, each omitting $`n`$ variables. In this section we use the Stoke’s theorem to combine the integrals of the summands into one integral of the highest degree form over a domain in $`^N`$.
###### Proposition 5.1.
Let $`\omega `$ be an $`(Nn)`$-form in $`^N`$. Then
$$\underset{D(\mathrm{\Gamma })}{}\omega =\pm \underset{0\rho _i1,X_e>0}{}d\rho _1\mathrm{}\mathrm{d}\rho _n\omega $$
###### Proof.
We apply the Stoke’s theorem multiply. Indeed,
$$\underset{D}{}\omega =\underset{D}{}\rho _1\omega =\pm \underset{X_e>0,\rho _2=\mathrm{}=\rho _n=1,0\rho _11}{}d\rho _1\omega =$$
$$=\underset{X_e>0,\rho _2=\mathrm{}=\rho _n=1,0\rho _11}{}\rho _2d\rho _1\omega =\mathrm{}=\pm \underset{0\rho _i1,X_e>0}{}d\rho _1\mathrm{}\mathrm{d}\rho _n\omega .$$
Now we need to deal with the $`\mathrm{d}\rho _i`$ factors. We prove the following
###### Theorem 5.2.
For $`I`$ being some set of indices with $`|I|=Nn`$,
$$\left|\underset{0\rho _i1,X_e>0}{}d\rho _1\mathrm{}\mathrm{d}\rho _n\underset{iI}{}\frac{\mathrm{d}e_i}{e_i}\right|\left|\underset{0\rho _i1,X_e>0}{}n!\rho ^n\underset{i=1}{\overset{N}{}}\frac{\mathrm{d}e_i}{e_i}\right|$$
###### Proof.
Indeed, recall (, section 3.3.4) the definition of the $`\alpha `$-lengths $`\alpha (e,v):=\frac{e}{fg}`$ in the usual notations. Lemma 3.4.2 in states that $`\rho _i`$ is twice the sum of $`\alpha `$-lengths of the sectors it traverses. What matters for us is that it is a sum of some $`\alpha `$-lengths. For any $`\alpha `$-length we have
$$\frac{\alpha (e,v)}{e}=\frac{1}{e}\alpha (e,v)\mathrm{and}\frac{\alpha (e,v)}{f}=\frac{1}{f}\alpha (e,v).$$
Thus for $`\rho _i=2\underset{j=1}{\overset{n}{}}\alpha (f_i,v_i)`$ we have
$$\left|\frac{\rho _i}{e}\right|==\frac{2}{e}|\pm \alpha (f_j,v_j)|\frac{2}{e}\alpha (f_j,v_j)=\frac{\rho _i}{e}<\frac{\rho }{e}.$$
Applying this trick to each of $`\mathrm{d}\rho _i`$ yields the theorem. ∎
Combining this theorem with the bound on the coefficients of the Weil-Petersson from theorem 3.1, noting that there are $`\left(\genfrac{}{}{0pt}{}{N}{n}\right)<N^n/n!`$ summands in the Weil-Petersson form, and enlarging $`D`$ to the domain $`0\rho n,X_e>0e`$, we get the following
###### Corollary 5.3.
The integral of the Weil-Petersson volume over the domain of integration is bounded by
$$\underset{D}{}\omega _{WP}^{(3g3+n)}<2^NN^n\left|\underset{0\rho <n,X_e>0}{}\rho ^n\underset{i=1}{\overset{N}{}}\frac{\mathrm{d}e_i}{e_i}\right|$$
## 6. Triangle inequality combinatorics
Now we proceed to show how the triangle inequalities lead to a converging integral as an upper bound for the Weil-Petersson volume. We develop an algorithm for inductive estimation of the integral of the Weil-Petersson volume form over the domain where the triangle inequalities hold.
###### Definition 6.1.
Two edges $`e`$ and $`f`$ of the graph are called linked if they are adjacent at a vertex $`v`$, and the third edge $`g`$ at $`v`$ is the minimal of the three (not necessarily strictly). Notice that from the triangle inequalities $`e<f+g`$ and $`f<e+g`$ it then follows that $`\frac{1}{2}\frac{e}{f}2`$.
###### Definition 6.2.
A chain is a sequence of edges in which every two consecutive ones are linked. We will only be interested in maximal chains — the ones which are not a part of any longer chain. Such a chain must either form a loop in the graph, or end by two edges which are minimal at their outer-end vertices.
###### Definition 6.3.
Define a wheel to be an ordered sequence of (maximal) chains $`c_1,\mathrm{},c_m`$ such that for any $`i`$ there is at least one chain among $`c_1,\mathrm{},c_{i1}`$ ending at a vertex inside the chain $`c_i`$. In further considerations, we will only be interested in maximal wheels, the ones which cannot be enlarged any further. Basically this means that the ends of all chains in the wheel belong to other chains already included in the wheel, so that there are no edges “sticking out” of the wheel.
For further computations we split the domain of integration into at most $`3^V`$ parts by deciding which two edges at each vertex are linked. We then use chains and wheels to introduce some order on the set of edges, and to integrate inductively. In this we will be aided by the following technical lemmas. The notations are as in the definition of linking: $`e`$, $`f`$ and $`g`$ are the three edges at some vertex, among which $`g`$ is minimal. We are working in the domain $`\mathrm{\Delta }`$ where the triangle inequalities hold, noting that $`\mathrm{\Delta }D(\mathrm{\Gamma })`$ by theorem 4.1.
###### Lemma 6.4.
For fixed $`e`$ we have $`_\mathrm{\Delta }\frac{\mathrm{d}f}{f}\mathrm{ln}4`$, where $`_\mathrm{\Delta }`$ denotes integration over the possible values of $`f`$ within domain $`\mathrm{\Delta }`$ for fixed $`e`$.
###### Proof.
Since $`e`$ and $`f`$ link, we have $`e/2<f<2e`$. Thus
$$_\mathrm{\Delta }\frac{\mathrm{d}f}{f}_{e/2}^{2e}\frac{\mathrm{d}f}{f}=\mathrm{ln}4.$$
###### Corollary 6.5.
Let $`e_1\mathrm{}e_m`$ be a chain. If we fix (the $`\lambda `$-length of) an edge $`e_i`$, then
$$_\mathrm{\Delta }\underset{ji}{}\frac{de_j}{e_j}(\mathrm{ln}4)^{m1}$$
###### Proof.
Start from the ends of the chain, and apply the lemma to eliminate edges one by one, coming from the ends towards $`e_i`$. ∎
###### Lemma 6.6.
If we fix the edge $`g`$, then the integral over the linked edges $`e`$ and $`f`$ can be estimated as $`_\mathrm{\Delta }\frac{\mathrm{d}e\mathrm{d}f}{ef}<2`$
###### Proof.
Split the integral into two parts depending on whether $`e>f`$ or $`f>e`$. The computation for them is identical; if $`f<e`$, we have
$$\underset{\mathrm{\Delta }\{f<e\}}{}\frac{\mathrm{d}e\mathrm{d}f}{ef}\underset{g}{\overset{\mathrm{}}{}}\frac{\mathrm{d}f}{f}\underset{f}{\overset{f+g}{}}\frac{\mathrm{d}e}{e}=\underset{g}{\overset{\mathrm{}}{}}\frac{\mathrm{d}f}{f}\mathrm{ln}\left(1+\frac{g}{f}\right)<\underset{g}{\overset{\mathrm{}}{}}\frac{\mathrm{d}f}{f}\frac{g}{f}=1.$$
###### Proposition 6.7.
Starting from two linked edges $`e`$ and $`f`$, construct a wheel consisting of edges $`e_1\mathrm{}e_m`$, where the list includes $`e`$ and $`f`$ themselves. Then for fixed $`g`$
$$_\mathrm{\Delta }\underset{i=1}{\overset{m}{}}\frac{\mathrm{d}e_i}{e_i}(\mathrm{max}(\mathrm{ln}4,\sqrt{2}))^m$$
###### Proof.
The wheel is a collection of chains $`c_1\mathrm{}c_k`$. Keep the two edges of $`c_k`$ in between which some chain $`c_i`$ ends (existent by definition of a wheel), and integrate the other ones out using corollary 6.5 — by this we pick up a factor of $`\mathrm{ln}4`$ for each edge. Then use lemma 6.6 above to integrate out the last two remaining edges of the chain $`c_k`$ — here we pick up a factor of $`2`$ for two edges, i.e. $`\sqrt{2}`$ per edge. Performing induction in $`k`$ finishes the proof. ∎
If a wheel were the whole graph, we would be able to estimate the integral using the above proposition. However, if the wheel is not the whole graph, we need to be able to link it to the rest of the graph. Thus we will need the following
###### Lemma 6.8.
In the usual notations for fixed $`e`$ we have
$$_\mathrm{\Delta }\frac{\mathrm{d}f\mathrm{d}g}{fg}<\frac{8}{3}$$
###### Proof.
Indeed, recall that $`g>4`$ by proposition 3.2. Thus
$$\underset{4}{\overset{e}{}}\frac{\mathrm{d}g}{g}\underset{\mathrm{max}(g,eg)}{\overset{g}{}}\frac{\mathrm{d}f}{f}+\underset{4}{\overset{g}{}}\frac{\mathrm{d}g}{g}\underset{e}{\overset{e+g}{}}\frac{\mathrm{d}f}{f}\underset{4}{\overset{e/2}{}}\frac{\mathrm{d}g}{g}\mathrm{ln}\frac{e}{eg}+\underset{4}{\overset{e}{}}\frac{\mathrm{d}g}{g}\mathrm{ln}\left(1+\frac{g}{e}\right)$$
Since $`\mathrm{ln}(1+x)<x`$ for $`x>0`$, for the second summand we have
$$\underset{4}{\overset{e}{}}\frac{\mathrm{d}g}{g}\mathrm{ln}(1+\frac{g}{e})<\underset{4}{\overset{e}{}}\frac{\mathrm{d}g}{g}\frac{g}{e}=\frac{e4}{e}<1.$$
For the first summand we compute
$$\underset{4}{\overset{e}{}}\frac{\mathrm{d}g}{g}\mathrm{ln}\frac{e}{eg}=\underset{4}{\overset{e}{}}\frac{\mathrm{d}g}{g}\mathrm{ln}\left(1\frac{g}{e}\right)=\underset{4}{\overset{e}{}}\frac{\mathrm{d}g}{g}\underset{n=1}{\overset{\mathrm{}}{}}\frac{g^n}{ne^n}=$$
$$=\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{ne^n}\underset{4}{\overset{e}{}}g^{n1}dg=\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{ne^n}\frac{e^n4^n}{n}<\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{n^2}<\frac{5}{3}.$$
Combining the above estimates, we get the lemma. ∎
###### Theorem 6.9.
Let $`\mu `$ be the minimal edge of $`\mathrm{\Gamma }`$, and let $`e_1\mathrm{}e_{N1}`$ be all the other edges of the graph. Then for a fixed value of $`\mu `$ and a fixed choice of the two linking edges at each vertex we have
$$_\mathrm{\Delta }\underset{i=1}{\overset{N1}{}}\frac{\mathrm{d}e_i}{e_i}<\left(\frac{8}{3}\right)^{(N1)/2}$$
###### Proof.
Construct a wheel $`w_1`$ starting from edge $`\mu `$. If this wheel is not the whole graph, consider a chain $`c_{1,1}`$ ending on $`w_1`$ by at least one end. If it ends on $`w_1`$ by the other end also, consider another such chain $`c_{1,2}`$ and so on, until we either exhaust the graph, or get a chain $`c_{1,m_1}`$ which has an end not on $`w_1`$. Then construct a wheel $`w_2`$ at the other end of the chain $`c_{1,m_1}`$. If the union of these two wheels and the chains constructed is still not the whole graph, we repeat the process.
As a result, we decompose the graph into a disjoint union of wheels $`w_1\mathrm{}w_k`$ and chains $`c_{i,j}`$ for $`ik`$ and $`jm_i`$ such that the chains $`c_{i,j}`$ have both ends on $`w_i`$ for $`i=k`$ or $`j<m_i`$, and that $`c_{i,m_i}`$ connects $`w_i`$ and $`w_{i+1}`$ for $`i<k`$. Then use corollary 6.5 to eliminate all edges of chains $`c_{k,i}`$ except the terminal ones, which end at $`w_k`$. Using lemma 6.8, we can then include these terminal edges of $`c_{k,i}`$’s while implementing the proof of proposition 6.7 — integrating out the edges of $`w_k`$ one by one. Doing this, we get a factor of $`8/3`$ for eliminating two edges, instead of the smaller factor of $`(\mathrm{ln}4)^2`$, which we were getting originally.
At the last step of integrating over the edges of $`w_k`$ we use the edge at the end of $`c_{k1,m_{k1}}`$ for lemma 6.6, and thus reduce the problem to $`k1`$ wheels. Induction in $`k`$ then yields the desired result. ∎
Now we combine all the above estimates to finally obtain
###### Theorem 6.10.
In the above notations,
$$_{D(\mathrm{\Gamma })}\omega _{WP}^{(3g3+n)}<2^N3^VN^n\left(\frac{8}{3}\right)^{(N1)/2}(2V)^n$$
###### Proof.
Combining the results of corollary 5.3 and proposition 4.2, we see that the integral in question is bounded above by
$$2^NN^n\left|\rho ^n\frac{\mathrm{d}\mu }{\mu }\underset{i=1}{\overset{N1}{}}\frac{\mathrm{d}e_i}{e_i}\right|<2^NN^n\left|\left(\frac{8V}{\mu }\right)^n\frac{\mathrm{d}\mu }{\mu }\underset{i=1}{\overset{N1}{}}\frac{\mathrm{d}e_i}{e_i}\right|.$$
Using theorem 6.9, we can integrate out all variables except $`\mu `$, by acquiring an extra factor of $`(8/3)^{(N1)/2}`$. Remembering the factor of $`3^V`$ for choosing the minimal edge at each vertex, our final upper bound becomes
$$2^N3^VN^n\left(\frac{8}{3}\right)^{(N1)/2}(8V)^n\underset{4}{\overset{\mathrm{}}{}}\frac{\mathrm{d}\mu }{\mu ^{n+1}}=2^N3^VN^n\left(\frac{8}{3}\right)^{(N1)/2}(2V)^n$$
Using our explicit knowledge of the Weil-Petersson volume form, in the case of one puncture we get
$$\underset{D}{}\omega _{WP}^{(3g3+n)}<2^{4g2}3^V(8/3)^{(N1)/2}N.$$
## 7. Conclusion
For the case of one puncture, combining our estimates with Penner’s asymptotic computation of the number of cells in formula 2.1 (being $`\frac{(2g)!}{N}\left(e/6\right)^{2g}`$), and estimating $`1/\mathrm{Aut}\mathrm{\Gamma }`$ from above by one, we finally get
$$\mathrm{vol}_{WP}(_{g,1})<(2g)!\mathrm{\hspace{0.17em}2}^{4g2}3^{4g2}(\mathrm{ln}4)^{6g3}\left(\frac{6}{e}\right)^{2g}=:c_1^g(2g)!,$$
which has the same leading order infinity as Penner’s lower bound
$$\mathrm{vol}_{WP}(_{g,1})>\left(\frac{8e^2}{9}\right)^{2g}\frac{(2g)!}{2(6g3)^2}=:c_2^g(2g)!,$$
where $`c_i`$ are some explicit constants. Thus we have proven that
$$c_2^g(2g)!<\mathrm{vol}_{WP}(_{g,1})<c_1^g(2g)!\mathrm{and}\underset{g\mathrm{}}{lim}\frac{\mathrm{ln}\mathrm{vol}_{WP}(_{g,1})}{g\mathrm{ln}g}=2.$$
Intuitively, for more than one puncture the number of graphs should grow with genus in the same way as for one puncture, since all the punctures are far away from the additional handles being added, and their number should not matter.
Rigorously, let $`T(g,n)`$ denote the set of isomorphism classes of ideal triangulations of a surface of genus $`g`$ with $`n`$ punctures. Then we prove
###### Proposition 7.1.
There is a following upper bound on the number of triangulations:
$$|T(g,n)|<\frac{N^2}{2}|T(g,n1)<\mathrm{}<\frac{N^{2n2}}{2^{n1}}|T(g,1)|<\frac{(2g)!N^{2n3}}{2^{n1}}\left(\frac{6}{e}\right)^{2g}$$
###### Proof.
For $`n>1`$ we construct a relation $`\varphi T(g,n)\times T(g,n1)`$ in the following way. Consider two distinct punctures $`p_1`$ and $`p_2`$ connected by an edge $`e`$ of a triangulation $`xT(g,n)`$ — if such did not exist, i.e. if all edges emanating from a puncture went back to the puncture itself, it would not be a triangulation of the surface. Shrinking $`e`$ to a point, and collapsing triangles on both sides of $`e`$ into arcs of a triangulations, thus identifying $`p_1`$ and $`p_2`$, produces a new triangulation $`yT(g,n1)`$. We define $`\varphi `$ to be the set of all pairs $`(x,y)`$ obtained in such a way.
Now consider some $`yT(g,n1)`$. Any graph in $`xT(g,n)`$ such that $`(x,y)\varphi `$ can be reconstructed from $`y`$ by “blowing up” a pair of edges emanating from a vertex to triangles. Since there are $`(6g6+3(n1))(6g6+3(n1)1)/2<N^2/2`$ ways to choose a pair of edges of $`y`$, there are at most $`N^2/2`$ triangulations $`xT(g,n)`$ such that $`(x,y)\varphi `$. The argument works for all $`y`$, and thus $`|T(g,n)|<N^2|T(g,n1)|/2`$. Applying this argument until we decrease the number of punctures to one, and then utilizing Penner’s asymptotic computation for that case finishes the proof. ∎
Combining all our estimates, for $`gn`$ we get
$$\mathrm{vol}_{WP}(_{g,n})<\frac{(2g)!N^{2n3}}{2^{n1}}\left(\frac{6}{e}\right)^{2g}2^N3^VN^n\left(\frac{8}{3}\right)^{N/2}(2V)^n<C^g(2g)!,$$
where $`c`$ is any constant greater than $`2^{17}3^3/e^2`$ (notice that is is independent of $`n`$, since the $`g^n`$ has a lower growth order), and thus
$$\underset{g\mathrm{},n\mathrm{fixed}}{lim}\frac{\mathrm{ln}\mathrm{vol}_{WP}(_{g,n})}{g\mathrm{ln}g}2.$$
## Acknowledgements
The author would like to thank Professor Yum-Tong Siu for suggesting the problem and for many invaluable discussions. Without Professor Siu’s continuous support this work would have been impossible. We would also like to thank the University of Hong Kong for hospitality during November 1999, when this work was finalized.
|
warning/0003/gr-qc0003048.html
|
ar5iv
|
text
|
# Double Scalar–Tensor Gravity Cosmologies
## 1 Introduction
Extended gravity theories have recently assumed a prominent role in theoretical physics investigations since any unification scheme, as supergravity or superstrings, in the weak energy limit, or viable early universe cosmological models, as extended inflation, seem to base on them.
Besides, any effective theory, where quantum fields are taken into account in a curved space–time, results in nonminimal couplings between geometry and matter (scalar) fields . Furthermore, notwithstanding the fact that Einstein’s general relativity is experimentally tested with high degree of accuracy, from solar system tests to binary pulsar observational data, it has become a peremptory necessity to consider alternative theories of gravity. The issues is now: What kind of theory? The plethora of them is overwhelming from higher–order gravity theories, to Kaluza–Klein multidimensional theories, to induced–gravity theories, to gauge theories with torsion. Several of them seem to be consistent with some quantum gravity effect in the weak energy limit but, till now, no one can be universally considered the ”full” quantum gravity theory.
Despite of this shortcoming, most of them have remarkable physical implications from cosmology to particle physics. A particularly relevant role is played in inflationary cosmology where, from the early Starobinsky model to the more recent hyperextended models, non–standard theories have been widely used. In fact, the goal of every inflationary model is to generate a brief period in which the scalar factor of the universe, $`a(t)`$, increases superluminally, i.e. $`a(t)>t`$. If $`a(t)`$ grows by $`e^{60}`$ or more during this period, the horizon, flatness and monopole problems can be resolved. In addition, inflation generates energy density fluctuations which may be seeds for large scale structure formation. This features can be obtained in several alternative theories of gravity.
However, designing a detailed microphysical model that accomplishes all of these goals has proven to be extremely difficult and several times the extended theories have to be adjusted to reach some partial goal.
As a general scheme, one needs a grand entrance into inflation which is a mechanism which drives the universe into a false vacuum phase.The large, positive vacuum density acts as an effective cosmological constant which triggers a period of a (quasi) de Sitter expansion. Then one needs a graceful exit: a mechanism capable of terminating the inflationary expansion, reheat the universe to a high temperature, and restore a Friedman expansion. This second issue can result highly problematic due to the fine–tuning of parameters which one needs to connect in the isotropy of the cosmic microwave background (e.g. the so called ”big bubble problem” ).
Cosmological models deduced from nonminimally coupled theories of gravity (e.g. Brans–Dicke or induced gravity) have provided schemes capable of escaping such difficulties. Extended inflationary models could accomplish several goals which earlier models failed and, in particular, they cure some shortcomings of ”old” inflation. Extended inflationary models also start when the universe is trapped in a false vacuum state by a large energy barrier during a first–order phase transition . The failure of old inflation was that the universe could never escape the false vacuum state since the rate of tunneling through the barrier remains small compared to the inflationary expansion rate . Extended models, in the various versions, avoid the same failure by introducing mechanisms so that the tunneling rate exceeds the expansion rate and, hence, the transition to the true vacuum can be completed. For a comprehensive review, see .
The key ingredient on which extended inflation lies is the relation between the tunneling or bubble nucleation rate, $`\lambda `$, and the expansion rate (the Hubble parameter) $`H`$. The ratio is the dimensionless ”bubble nucleation rate” $`ϵ=\lambda /H^4`$. The false vacuum can be percolated by true vacuum bubbles only if $`ϵ`$ exceeds a critical value, $`ϵ_{crit}0.2`$. In old inflation, $`ϵ`$ is time–independent since $`\lambda `$ and $`H`$ (which depends on the false vacuum energy) do not vary during inflation .
Two situations are possible: i) $`ϵ<ϵ_{crit}`$, in which case the true vacuum bubbles never percolate and the universe inflates forever, or ii) $`ϵ>ϵ_{crit}`$, in which case the true vacuum percolates, but so quickly that there is insufficient inflation to solve cosmological problems. A way to bypass this shortcoming is to avoid bubble nucleation altogether. New inflation and chaotic inflation utilize this approach. For example, in new inflation, the energy barrier disappears altogether as the universe supercools and the universe evolves slowly but continuously from the false to true vacuum phase . However, the model has to be fine–tuned.
Extended inflation models employ an alternative approach asking for the time variation of $`ϵ`$. Initially, $`ϵ`$ is much less than $`ϵ_{crit}`$ to achieve sufficient inflation, but then it grows during inflation to a value $`ϵ>ϵ_{crit}`$ so that the phase transition can be completed. Being $`\lambda `$ fixed by the form of the self interacting potential, the only quantity on which one can act is $`H`$. Cosmological solution with a time–varying $`H`$ can be easily obtained by using modified theories of gravity like scalar–tensor theories. A simple extended inflationary model can be contructed by using the Brans–Dicke theory . The gravitational action is
$$𝒜=d^4x\sqrt{g}\left[\varphi R\omega \left(\frac{\varphi ^\mu \varphi _\mu }{\varphi }\right)+_m\right],$$
(1.1)
where $`\varphi `$ is the Brans–Dicke scalar field, $`\omega `$ is a dimensionless parameter and $`_m`$ is the matter Lagrangian density including all the non–gravitational fields (from now on, we shall use natural units $`\mathrm{}=c=k_B=8\pi G=1`$). The cosmological Friedman–Robertson–Walker (FRW) equations are
$$H^2=\frac{\rho _m}{3\varphi }\frac{k}{a^2}+\frac{\omega }{6}\left(\frac{\dot{\varphi }}{\varphi }\right)^2H\left(\frac{\dot{\varphi }}{\varphi }\right),$$
(1.2)
$$\ddot{\varphi }+3H\dot{\varphi }=\frac{\rho _m3p_m}{3+2\omega },$$
(1.3)
where $`H=\dot{a}/a`$, $`\rho _m`$ and $`p_m`$ are, respectively, the energy and pressure densities of matter, and finally, $`k=0,\pm 1`$. For $`p_m=\rho _m`$ and $`k=0`$ we have
$$\varphi (t)=\left(1+\frac{\chi t}{\alpha }\right)^2,$$
(1.4)
$$a(t)=\left(1+\frac{\chi t}{\alpha }\right)^{\omega +1/2},$$
(1.5)
where $`\chi `$ is an integration constant connected to the energy of false vacuum, $`\alpha =(3+2\omega )(5+6\omega )/12`$.
Immediately we see that for $`\chi t<\alpha `$, $`\dot{\varphi }0`$, we have a de Sitter solution. If, for example, $`\omega >90`$, one obtains the 60 $`e`$–foldings necessary to solve cosmological problems of standard model. When $`\chi t>\alpha `$, $`a(t)`$ evolves as a power–law expansion and $`Ht^1`$. This feature allows the successful graceful exit since $`ϵ>ϵ_{crit}`$. Of course, the main ingredients are the variation of Newton constant and the coupling of the geometry to the scalar field. In other words, the model succeeds because the effective Newton constant $`G_{eff}=\varphi ^1`$ is decreasing, then $`H`$ is decreasing and $`ϵ`$ is increasing.
The main flaw of this model is related to the expected value of the parameter $`\omega `$. In order to restore Einstein’s general relativity, we should have $`\omega \mathrm{}`$ , then the value of $`\omega `$ in constrained by the classical tests of general relativity: light deflection and time–delay experiments require $`\omega >500`$ while the bounds on the anisotropy of the microwave background radiation give $`\omega 30`$ . In conclusion, $`\omega `$ must be a function of time in order to obtain viable models. A pure Brans–Dicke theory is not able to yield realistic models and we have to introduce, at least, a function $`\omega =\omega (\varphi )`$ in order to overcome the above difficulties. Several proposal have been done to improve the early extended inflationary model and, in some of them also $`\lambda `$ is assumed to vary . The strict condition which implies $`\lambda =`$constant is a feature connected to Brans–Dicke models as it is shown in .
In the so called hyperextended inflation , one can assume
$$\omega (\varphi )=\omega _0+\omega _m\varphi ^m,$$
(1.6)
to improve the Brans–Dicke model . If $`m=5`$, the microwave background bounds are satisfied . Alternatively, assuming the coupling
$$\varphi =F(\phi ),\omega (\varphi )=\frac{F(\phi )}{2(dF/d\phi )^2},$$
(1.7)
one can get successful implementations without fine–tuning the initial conditions. In particular, if $`F(\phi )`$ is a sixth order polynomial, the big bubble problem is avoided since the model is independent of the bubble–size distribution .
Other approaches give interesting results. For example, it is possible to include a first- or second-order potential $`V(\varphi )`$ for the Brans–Dicke field $`\varphi `$ as in induced gravity theories . In this case, the potential places constraints on the percolation time–scale in the graceful exit and, furthermore, it can give rise to multiple episodes of inflation which may reveal extremely useful for large scale structure formation (e.g. super-cluster, cluster and galaxies).
Another way to escape the $`\omega `$–parameter constraints is to consider a curvature–coupled inflation . Also in this case, extended inflation results enhanced since, in some sense, the roles of Brans–Dicke field and inflaton are mixed. This feature allows to satisfy the solar system constraints on $`\omega `$, to avoid the big–bubble problem, to construct models with double inflationary episodes.
A more sophisticated way to bypass the graceful exit problem can be obtained by coupling first–order phase transitions to curvature–squared inflation . The mechanism (getaway inflation) is based on a nonminimally coupled higher–order gravity theory where terms like $`\phi ^2R^2`$ appear in the usual gravity–inflaton action. Their role is to produce an inflationary phase of the background which has a classical end. At the same time, a stage of bubble production via semiclassical tunneling occurs allowing useful spectra for large scale structure formation.
A final remark concerns the role which the Brans–Dicke scalar could have for dark matter in extended inflation. Its oscillations in the various models could account for the discrepancy between the dynamical estimate of the density of matter in the universe, $`\mathrm{\Omega }0.2`$, and the prediction of inflation, $`\mathrm{\Omega }=1`$ which seems to be confirmed by the BOOMERANG experiment .
All these arguments and several more make extremely interesting to search for cosmological solutions useful for extended inflation. A first investigation in this sense is in where general scalar–tensor theories of gravitation were studied in order to ”model” useful extended inflationary behaviours. Intermediate inflationary universes with expansion scale factor of the form
$$a(t)=a_0\mathrm{exp}t^p,0<p<1,$$
(1.8)
were found. These models allow to succeed in realizing phase transition and graceful exit.
More recently, Modak and Kamilya derived exact cosmological solutions by the so called Noether Symmetry Approach in scalar–tensor gravity theories discussing the role of the coupling function $`\omega (\phi )`$ connected to the Noether symmetry. They improved the approach in , where symmetries and solutions were found for theories with a scalar field nonminimally coupled to gravity, by introducing a second scalar field (the inflaton) as in extended inflationary models. Exponentially expanding solutions, in asymptotic region, were found and this feature does not allow to solve the graceful exit problem also if general relativity was asymptotically recovered.
In this paper we want to discuss, by Noether Symmetry Approach, a further generalization taking into account two nonminimally coupled scalar fields and their self interaction potentials. In this way, the roles of the Brans–Dicke field and the inflaton are mixed and both fields are taken on the same ground. This fact could be coherent with the stochastic approach for the fundamental laws of nature since the role of the fields is not attributed a priori .
The paper is organized as follows. In Sect. 2, we discuss the double scalar–tensor action, derive the equations of motion, the point-like FRW Lagrangian and the cosmological equations. Sect. 3 is devoted to the Noether Symmetry Approach which has to be improved for the double field case since the configuration space results enlarged. The summary of found symmetries is given considering also the subcases where one and not two nonminimal couplings are present. The cosmological solutions are given in Sect. 4 while the graceful exit problem is discussed in Sect. 5. Conclusions are drawn in Sect.6.
## 2 Double Scalar-Tensor Action and Equations of Motion
The most general action in four dimensions, where gravity is nonminimally coupled to two scalar fields noninteracting between them, is
$$𝒜=d^4x\sqrt{g}\left[F(\phi )R+G(\psi )R+\frac{1}{2}\phi _\mu \phi ^\mu V(\phi )+\frac{1}{2}\psi _\mu \psi ^\mu W(\psi )\right],$$
(2.1)
where we have not specified the four functions $`F(\phi )`$, $`V(\phi )`$, $`G(\psi )`$, and $`W(\psi )`$. This action generalizes those used till now to construct extended inflationary models<sup>1</sup><sup>1</sup>1We point out that the more general action is
$$𝒜=d^4x\sqrt{g}\left[F(\phi ,\psi )RV(\phi ,\psi )+A(\phi ,\psi )(\phi )^2+B(\phi ,\psi )(\psi )^2\right],$$
but for the purpose of the paper, we will confine ourselves to the action (2.1). The Brans–Dicke action (1.1) can be immediately recovered by using the transformations (1.7). In our units, the standard Newton coupling is recovered in the limit $`F(\phi )+G(\psi )1/2`$. The field equations can be derived by varying with respect to $`g_{\mu \nu }`$
$$[F(\phi )+G(\psi )]\left(R_{\mu \nu }\frac{1}{2}g_{\mu \nu }R\right)=T_{\mu \nu }^{(\phi )}+T_{\mu \nu }^{(\psi )}.$$
(2.2)
In the right hand side of (2.2) there is the effective stress–energy tensor containing the nonminimal coupling contributions, the kinetic terms and the potentials of the scalar fields $`\phi `$ and $`\psi `$, that is
$$T_{\mu \nu }^{(\phi )}=\frac{1}{2}\phi _{;\mu }\phi _{;\nu }+\frac{1}{4}g_{\mu \nu }\phi _{;\alpha }\phi ^{;\alpha }\frac{1}{2}g_{\mu \nu }V(\phi )g_{\mu \nu }\mathrm{}F(\phi )+F(\phi )_{;\mu \nu }$$
(2.3)
and analogously,
$$T_{\mu \nu }^{(\psi )}=\frac{1}{2}\psi _{;\mu }\psi _{;\nu }+\frac{1}{4}g_{\mu \nu }\psi _{;\alpha }\psi ^{;\alpha }\frac{1}{2}g_{\mu \nu }W(\psi )g_{\mu \nu }\mathrm{}G(\psi )+G(\psi )_{;\mu \nu }.$$
(2.4)
$`\mathrm{}`$ is the d’Alambertian operator. The variation with respect to $`\phi `$ and $`\psi `$ gives the Klein–Gordon equations
$$\mathrm{}\phi R\left(\frac{dF}{d\phi }\right)+\frac{dV}{d\phi }=0,$$
(2.5)
and
$$\mathrm{}\psi R\left(\frac{dG}{d\psi }\right)+\frac{dW}{d\psi }=0.$$
(2.6)
Their sum is equivalent to the contracted Bianchi identities . Let us now take into account a FRW metric of the form
$$ds^2=dt^2a^2\left[\frac{dr^2}{1kr^2}+r^2d\mathrm{\Omega }^2\right]$$
(2.7)
and substitute it into the action (2.1). Integrating by parts and eliminating the boundary terms, we get the point-like Lagrangian
$$=6\frac{dF}{d\phi }a^2\dot{a}\dot{\phi }+6Fa\dot{a}^26kFa+\frac{a^3\dot{\phi }^2}{2}a^3V(\phi )+6\frac{dG}{d\psi }a^2\dot{a}\dot{\psi }+6Ga\dot{a}^26kGa+\frac{a^3\dot{\psi }^2}{2}a^3W(\psi ).$$
(2.8)
The Euler–Lagrange equations, corresponding to the cosmological Einstein equations are
$$[F+G]\left[2\frac{\ddot{a}}{a}+\left(\frac{\dot{a}}{a}\right)^2+\frac{k}{a^2}\right]+2\left[\dot{\phi }\frac{dF}{d\phi }+\dot{\psi }\frac{dG}{d\psi }\right]\left(\frac{\dot{a}}{a}\right)+$$
(2.9)
$$+\left[\dot{\phi }^2\frac{d^2F}{d\phi ^2}+\ddot{\phi }\frac{dF}{d\phi }+\dot{\psi }^2\frac{d^2F}{d\psi ^2}+\ddot{\psi }\frac{dF}{d\psi }\right]\frac{1}{2}\left[\frac{1}{2}(\dot{\phi }^2+\dot{\psi }^2)(V+W)\right]=0,$$
$$6[F+G]\left(\frac{\dot{a}}{a}\right)^2+6\left[\dot{\phi }\frac{dF}{d\phi }+\dot{\psi }\frac{dG}{d\psi }\right]\left(\frac{\dot{a}}{a}\right)+\frac{6k}{a^2}[F+G]+\frac{1}{2}(\dot{\phi }^2+\dot{\psi }^2)+V+W=0,$$
(2.10)
$$\ddot{\phi }+3\left(\frac{\dot{a}}{a}\right)\dot{\phi }+6\left[\frac{\ddot{a}}{a}+\left(\frac{\dot{a}}{a}\right)^2+\frac{k}{a^2}\right]\left(\frac{dF}{d\phi }\right)+\frac{dV}{d\phi }=0,$$
(2.11)
$$\ddot{\psi }+3\left(\frac{\dot{a}}{a}\right)\dot{\psi }+6\left[\frac{\ddot{a}}{a}+\left(\frac{\dot{a}}{a}\right)^2+\frac{k}{a^2}\right]\left(\frac{dG}{d\psi }\right)+\frac{dW}{d\psi }=0,$$
(2.12)
where Eq. (2.10) is the energy constraint corresponding to the $`(0,0)`$–Einstein equation. Let us now go to solve the system (2.9)–(2.12) by using the Noether Symmetry Approach. The solutions strictly depend on the form of the functions $`F,G,V`$ and $`W`$. By the Noether symmetries it is possible to select these functions so that the system (2.10)–(2.12) can be reduced and then integrated.
## 3 Selecting Couplings and Potentials by the Noether Symmetries
Given an undefined extended gravity theory, the existence of a Noether symmetry can select the form of the coupling and the scalar field potential or the form of the higher–order Lagrangian density, e.g. $`f(R,\mathrm{}R)`$ .
At the same time, as we will show in the next section, the symmetry allows to reduce the dynamical system by a cyclic variable making it easier to solve. Taking into account the Lagrangian (2.8), its configuration space is three–dimensional, $`Q=\{a,\phi ,\psi \}`$. In the language of quantum cosmology, it can be identified with a minisuperspace . The tangent space on which the Lagrangian (2.8) is defined is $`TQ=\{a,\dot{a},\phi ,\dot{\phi },\psi ,\dot{\psi }\}`$ so that the lift vector $`X`$, the infinitesimal generator of symmetry, is
$$X=\alpha \frac{}{a}+\beta \frac{}{\phi }+\gamma \frac{}{\psi }+\frac{d\alpha }{dt}\frac{}{\dot{a}}+\frac{d\beta }{dt}\frac{}{\dot{\phi }}+\frac{d\gamma }{dt}\frac{}{\dot{\psi }},$$
(3.1)
where $`\alpha ,\beta ,\gamma `$ are functions of $`a,\phi ,\psi `$. A Noether symmetry exists if the condition
$$L_X=0,$$
(3.2)
is realized. $`L_X`$ is the Lie derivative wit respect to $`X`$. Properly speaking, Eq. (3.2) corresponds to the contraction of the vector $`X`$ with the Lagrangian (2.8). The constant of motion connected to the Noether symmetry is nothing else but
$$\mathrm{\Sigma }_0=i_X\vartheta _{}$$
(3.3)
where
$$\vartheta _{}=\frac{}{\dot{a}}da+\frac{}{\dot{\phi }}d\phi +\frac{}{\dot{\psi }}d\psi $$
(3.4)
is the Cartan one–form given by a Lagrangian $``$ and $`i_X`$ is the contraction with respect to $`X`$. The relation between Eqs. (3.2) and (3.3) can be easily seen if the vector is generally expressed as
$$X=\alpha ^i\frac{}{q^i}+\frac{d\alpha ^i}{dt}\frac{}{\dot{q}^i}.$$
(3.5)
Using the Euler–Lagrange equations, it can be shown that
$$\frac{d}{dt}\left(\alpha ^i\frac{}{\dot{q}^i}\right)=L_X.$$
(3.6)
If the Noether symmetry exists, Eq. (3.6) gives (3.2). In the Hamiltonian formalism, by a Legendre transformation, we get
$$=\dot{a}\pi _a+\dot{\phi }\pi _\phi +\dot{\psi }\pi _\psi ,$$
(3.7)
where $`\pi _q/\dot{q}`$, $`q=\{a,\phi ,\psi \}`$ are the conjugate momenta. The phase–space vector for the symmetry is now
$$\mathrm{\Gamma }=\dot{a}\frac{}{a}+\dot{\phi }\frac{}{\phi }+\dot{a\psi }\frac{}{\psi }+\ddot{a}\frac{}{\dot{a}}+\ddot{\phi }\frac{}{\dot{\phi }}+\ddot{\psi }\frac{}{\dot{\psi }},$$
(3.8)
and a Noether symmetry exists if
$$L_\mathrm{\Gamma }=0.$$
(3.9)
The conserved quantity (3.3) and the Hamiltonian (3.7) gives the Poisson brackets
$$\{\mathrm{\Sigma }_0,\}=0.$$
(3.10)
Our issue is now to determine the functions $`F,G,V,W`$ by this Noether symmetry technique. We shall adopt the Lagrangian formalism. The condition (3.2) gives the system of partial differential equations
$$F\left(\alpha +2a\frac{\alpha }{a}\right)+G\left(\alpha +2a\frac{\alpha }{a}\right)+a\left(\frac{dF}{d\phi }\right)\left(\beta +a\frac{\beta }{a}\right)+a\left(\frac{dG}{d\psi }\right)\left(\gamma +a\frac{\gamma }{a}\right)=0,$$
(3.11)
$$3\alpha +12\left(\frac{dF}{d\phi }\right)\frac{\alpha }{\phi }+2a\frac{\beta }{\phi }=0,$$
(3.12)
$$3\alpha +12\left(\frac{dG}{d\psi }\right)\frac{\alpha }{\psi }+2a\frac{\gamma }{\psi }=0,$$
(3.13)
$$a\beta \frac{d^2F}{d\phi ^2}+\left(2\alpha +a\frac{\alpha }{a}+a\frac{\beta }{\phi }\right)\left(\frac{dF}{d\phi }\right)+2\frac{\alpha }{\phi }F+\frac{a^2}{6}\frac{\beta }{a}+a\frac{\gamma }{\phi }\left(\frac{dG}{d\psi }\right)+2\frac{\alpha }{\phi }G=0,$$
(3.14)
$$a\gamma \frac{d^2G}{d\psi ^2}+\left(2\alpha +a\frac{\alpha }{a}+a\frac{\gamma }{\psi }\right)\left(\frac{dG}{d\psi }\right)+2\frac{\alpha }{\psi }G+\frac{a^2}{6}\frac{\gamma }{a}+a\frac{\beta }{\psi }\left(\frac{dF}{d\phi }\right)+2\frac{\alpha }{\psi }F=0,$$
(3.15)
$$6\frac{\alpha }{\phi }\left(\frac{dG}{d\psi }\right)+6\frac{\alpha }{\psi }\left(\frac{dF}{d\phi }\right)+\frac{\beta }{\psi }a+\frac{\gamma }{\phi }a=0,$$
(3.16)
$$6k\left[\alpha F+\beta \left(\frac{dF}{d\phi }\right)a+\alpha G+\gamma \left(\frac{dG}{d\psi }\right)a\right]+a^2\left[3\alpha V+\beta a\left(\frac{dV}{d\phi }\right)+3\alpha W+\gamma a\left(\frac{dW}{d\psi }\right)\right]=0,$$
(3.17)
obtained by equating to zero the second degree coefficients in $`\dot{a}`$, $`\dot{\phi }`$, $`\dot{\psi }`$. The number of these equations is $`1+n(n+1)/2`$, where $`n`$ is the dimension of the configuration space $`Q`$. The system (3.11)–(3.17) is the straightforward generalization of the system (5.30)–(5.34) in for the case of two scalar fields.
The integration of (3.11)–(3.17) gives as a result the functions $`\alpha (a,\phi ,\psi )`$, $`\beta (a,\phi ,\psi )`$, $`\gamma (a,\phi ,\psi )`$, $`F(\phi )`$, $`V(\phi )`$, $`G(\psi )`$, $`W(\psi )`$. Solutions are not unique and the various cases which we have found are summarized in Table I.
In Table II, being $`G(\psi )=0`$, the cases where only one nonminimal coupling is present are summarized. The value of the spatial curvature constant $`k`$ is also given. The quantities $`F_0`$, $`G_0`$, $`F_0^{}`$, $`\mathrm{\Lambda }_{1,2}`$ are constants. Using these results, the dynamical system (2.9)–(2.12) can be reduced since, as we shall show below, a change of variables can be found where a cyclic coordinate is present. This feature allows to integrate more simply the dynamics.
## 4 The Cosmological Solutions
The existence of a Noether symmetry gives, in any case, a cyclic variable so that the transformation
$$(a,\dot{a},\phi ,\dot{\phi },\psi ,\dot{\psi })(w,\dot{w},u,\dot{u},\dot{z})$$
(4.1)
is always possible. If more than one symmetry exists, more than one cyclic variable can be present . In a geometric language, it is always possible to choose a new set $`q_i=q_i(Q_k)`$ $`i,k=1,2,3,`$ adapted to the foliation given by $`X`$
$$i_XdQ_3=1,i_XdQ_j=0,j=1,2,$$
(4.2)
where $`i_X`$, as before, is the contraction given by $`X`$ and $`dQ_j=(Q_j/q_i)dq_i`$. Explicitly, in our case, Eqs. (4.2) become
$$\alpha \frac{w}{a}+\beta \frac{w}{\phi }+\gamma \frac{w}{\psi }=0,$$
(4.3)
$$\alpha \frac{u}{a}+\beta \frac{u}{\phi }+\gamma \frac{u}{\psi }=0,$$
(4.4)
$$\alpha \frac{z}{a}+\beta \frac{z}{\phi }+\gamma \frac{z}{\psi }=1,$$
(4.5)
where $`z`$ is the cyclic variable. However, the transformation (4.3)–(4.5) are specified as soon as the functions $`\alpha `$, $`\beta `$, $`\gamma `$ are given. As an example, let us take into account Case 1 in Table I. Cases 2 and 3 of Table I and 1,2,3 of Table II can be deduced from it. The system (4.3)–(4.5) is solved by the choice of the new variables
$$w=a^3\psi ^2,u=a^3\phi ^2,z=\mathrm{ln}a.$$
(4.6)
For the scalar–field functions, as we said, we choose Case 1 in Table I. Lagrangian (2.8) becomes
$$=6F_0\dot{z}(\dot{w}2\dot{z}w)+\frac{1}{8w}(\dot{w}3\dot{z}w)^2V_0w+6G_0\dot{z}(\dot{u}2\dot{z}u)+\frac{1}{8w}(\dot{u}3\dot{z}u)^2W_0u$$
(4.7)
while the equations of motion are
$$6(8G_01)\ddot{z}+3(32G_03)\dot{z}^2+2\left(\frac{\ddot{u}}{u}\frac{\dot{u}^2}{2u^2}\right)+8W_0=0,$$
(4.8)
$$6(8F_01)\ddot{z}+3(32F_03)\dot{z}^2+2\left(\frac{\ddot{w}}{w}\frac{\dot{w}^2}{2w^2}\right)+8V_0=0,$$
(4.9)
$$(8F_01)\dot{w}(32F_03)\dot{z}w+(8G_01)\dot{u}(32G_03)\dot{z}u=\mathrm{\Sigma }_0$$
(4.10)
$$6(8F_01)\dot{z}\dot{w}3(32F_03)\dot{z}^2w+\frac{\dot{w}^2}{w}+8V_0w+6(8G_01)\dot{z}\dot{u}3(32G_03)\dot{z}^2u+\frac{\dot{u}^2}{u}+8W_0u=0,$$
(4.11)
where, clearly, $`z`$ is the cyclic variable and $`\mathrm{\Sigma }_0`$ is the constant of motion connected to $`z`$. With respect to the system (2.9)–(2.12), system (4.8)–(4.11) is reduced and it is highly symmetric due to the functions $`F(\phi ),G(\psi ),V(\phi ),W(\psi )`$ selected by the Noether symmetry.
Since the role of the two fields is completely symmetric, we can suppose that, depending on the value of the parameters, one can select two regimes of physical interest. For example, in the first case, dynamics is $`\phi `$ ($`u`$)–dominated, in the second case, it is $`\psi `$ ($`w`$)–dominated. The relation among the initial data is given by the energy condition (4.11). We get in the $`\phi `$ ($`u`$)–dominated regime
$$u(t)=c_1\mathrm{exp}[\lambda _1t]+c_2\mathrm{exp}[\lambda _1t],$$
(4.12)
$$z(t)=z_1\mathrm{arctan}\left(\sqrt{\frac{c_1}{c_2}}\mathrm{exp}[\lambda _1t]\right)+z_2\mathrm{ln}|c_1\mathrm{exp}[\lambda _1t]+c_2\mathrm{exp}[\lambda _1t]|+z_0,$$
(4.13)
for $`c_1c_2>0`$, and
$$z(t)=z_1\text{arctanh}\left(\sqrt{|\frac{c_1}{c_2}|}\mathrm{exp}[\lambda _1t]\right)+z_2\mathrm{ln}|c_1\mathrm{exp}[\lambda _1t]+c_2\mathrm{exp}[\lambda _1t]|+z_0,$$
(4.14)
for $`c_1c_2<0`$, while in the $`\psi `$ ($`w`$)–dominated regime
$$w(t)=c_3\mathrm{exp}[\lambda _2t]+c_4\mathrm{exp}[\lambda _2t],$$
(4.15)
$$z(t)=z_3\mathrm{arctan}\left(\sqrt{\frac{c_3}{c_4}}\mathrm{exp}[\lambda _2t]\right)+z_4\mathrm{ln}|c_3\mathrm{exp}[\lambda _2t]+c_4\mathrm{exp}[\lambda _2t]|+z_0,$$
(4.16)
for $`c_3c_4>0`$, and
$$z(t)=z_3\text{arctanh}\left(\sqrt{|\frac{c_3}{c_4}|}\mathrm{exp}[\lambda _2t]\right)+z_4\mathrm{ln}|c_3\mathrm{exp}[\lambda _2t]+c_4\mathrm{exp}[\lambda _2t]|+z_0,$$
(4.17)
for $`c_3c_4<0`$. The constants $`z_0`$, $`c_i,z_j`$, $`i,j=1,\mathrm{},4`$ are integration constants while
$$\lambda _1=\frac{V_0(32F_03)}{4F_0(12F_01)},\lambda _2=\frac{W_0(32G_03)}{4G_0(12G_01)},$$
(4.18)
assume the role of cosmological constants depending on the couplings and the potentials.
Inverting the relations (4.6), we get
$$a(t)=a_0\left(c_1\mathrm{exp}[\lambda _1t]+c_2\mathrm{exp}[\lambda _1t]\right)\mathrm{exp}\left[z_1\mathrm{arctan}\left(\sqrt{\frac{c_1}{c_2}}\mathrm{exp}[\lambda _1t]\right)\right],$$
(4.19)
$$\phi (t)=\phi _0\left(c_1\mathrm{exp}[\lambda _1t]+c_2\mathrm{exp}[\lambda _1t]\right)^2\left(\mathrm{exp}\left[z_1\mathrm{arctan}\left(\sqrt{\frac{c_1}{c_2}}\mathrm{exp}[\lambda _1t]\right)\right]\right)^3,$$
(4.20)
for $`c_1c_2>0`$, and
$$a(t)=a_0\left(c_1\mathrm{exp}[\lambda _1t]+c_2\mathrm{exp}[\lambda _2t]\right)\mathrm{exp}\left[z_1\text{arctanh}\left(\sqrt{|\frac{c_1}{c_2}|}\mathrm{exp}[\lambda _1t]\right)\right],$$
(4.21)
$$\phi (t)=\phi _0\left(c_1\mathrm{exp}[\lambda _1t]+c_2\mathrm{exp}[\lambda _2t]\right)^2\left(\mathrm{exp}\left[z_1\text{arctanh}\left(\sqrt{\frac{c_1}{c_2}}\mathrm{exp}[\lambda _1t]\right)\right]\right)^3,$$
(4.22)
for $`c_1c_2<0`$ in the $`\phi `$–dominated regime.
The situation is analogous in the $`\psi `$–dominate regime, but the constants $`\lambda _1`$, $`z_1`$, $`c_{1,2}`$ and $`\phi _0`$ have to be substituted with $`\lambda _2`$, $`z_2`$, $`c_{3,4}`$ and $`\psi _0`$.
It is easy to see hat we have two inflationary eras. Their durations are ruled by the parameter $`\lambda _1,\lambda _2`$ which strictly depends on the strength of the couplings. Another interesting particular solution is
$$a(t)=a_0e^{z_0t},\phi (t)=\phi _0\mathrm{exp}\left\{\frac{4F_0z_0}{8F_01}t\right\},\psi (t)=\psi _0\mathrm{exp}\left\{\frac{4G_0z_0}{8G_01}t\right\},$$
(4.23)
where
$$V_0=\frac{z_0^2}{8(8F_01)}F_0\left(F_0\frac{1}{12}\right)\left(F_0\frac{1}{10.6}\right)$$
(4.24)
$$W_0=\frac{z_0^2}{8(8G_01)}G_0\left(G_0\frac{1}{12}\right)\left(G_0\frac{1}{10.6}\right).$$
(4.25)
Also here the inflationary behaviour is clear.
By using similar arguments, we can analyze Case 2 in Tab.I Here the potential terms cancel each other in the Lagrangian (4.7). We get the power–law solution
$$a(t)=a_0t^n,\phi (t)=\phi _0t^n,\psi (t)=\psi _0t^n,$$
(4.26)
which is particularly useful for extended inflation being $`n`$ an arbitrary constant depending on the initial conditions. A similar situation holds in Case 4, which is a minimally coupled case with two fields and two cosmological constants. For $`k=0`$ and $`\mathrm{\Lambda }_1=|\mathrm{\Lambda }|_2`$, one finds
$$a(t)=a_0t^n,\phi (t)=\frac{\mathrm{\Sigma }_0}{a_0^3(13n)}t^{13n}+\phi _0,\psi (t)=\frac{\mathrm{\Sigma }_1}{a_0^3(13n)}t^{13n}+\psi _1,$$
(4.27)
for $`n1/3`$, and
$$a(t)=a_0t^{1/3},\phi (t)=\frac{\mathrm{\Sigma }_0}{a_0^3}\mathrm{ln}t+\phi _0,\psi (t)=\frac{\mathrm{\Sigma }_1}{a_0^3}\mathrm{ln}t+\psi _0,$$
(4.28)
when $`n=1/3`$. In this case, two Noether symmetries are present and they assign the value of gravitational constant being
$$F_0+G_0=\frac{3(\mathrm{\Sigma }_0^2+\mathrm{\Sigma }_1^2)}{4a_0^6}.$$
(4.29)
Let us now analyse, in detail, Case 5. Without loosing of generality, we can assume $`F_0^{}=F_0=0`$, $`\gamma _0=1`$ and studying the couplings $`F(\phi )=\phi ^2/12`$ and $`G(\psi )=G_0`$. Lagrangian (2.8) becomes
$$=\left(\frac{\phi ^2}{2}+6G_0\right)a\dot{a}^2+a^2\dot{a}\phi \dot{\phi }+a^3\left(\frac{\dot{\phi }^2}{2}\mathrm{\Lambda }\right)+\frac{a^3\dot{\psi }}{2}$$
(4.30)
where $`\mathrm{\Lambda }=\mathrm{\Lambda }_1+\mathrm{\Lambda }_2`$. Clearly $`\psi `$ is the cyclic variable. Eqs. (4.3)–(4.5) are satisfied by
$$w=a,u=a\phi \psi ,z=a\phi .$$
(4.31)
With the further change
$$\chi =zu,$$
(4.32)
Lagrangian (4.30) reads
$$=\left(6G_0\dot{w}^2+\frac{\dot{z}^2}{2}\right)w+\left(\frac{\dot{\chi }^2}{2}\mathrm{\Lambda }\right)w^3,$$
(4.33)
where two cyclic variables appear. The dynamical system is
$$6G_0(2\ddot{w}w+\dot{w}^2)=\frac{\dot{z}^2}{2}+3w^2\left(\frac{\dot{\chi }^2}{2}\mathrm{\Lambda }\right),$$
(4.34)
$$\dot{\chi }w^3=\mathrm{\Sigma }_1,$$
(4.35)
$$\dot{z}w=\mathrm{\Sigma }_2,$$
(4.36)
$$\dot{z}^2+12G_0\dot{w}^2+\dot{\chi }^2w^2+2\mathrm{\Lambda }w^2=0,$$
(4.37)
whose general solution is given by the elliptic integral
$$\frac{w^2dw}{\sqrt{A_1w^6+A_2w^2+A_3}}=\pm t,$$
(4.38)
where
$$A_1=\frac{A_1}{6G_0},A_2=\frac{\mathrm{\Sigma }_2^2}{6G_0},A_3=\frac{\mathrm{\Sigma }_1^2}{12G_0}.$$
(4.39)
In the particular case where $`A_2=0`$, we get the explicit solution
$$a(t)=a_0\sqrt[3]{\frac{\mathrm{\Sigma }_1^2}{2\mathrm{\Lambda }}\mathrm{sinh}(\pm 3\sqrt{A_1}t)},$$
(4.40)
$$\phi (t)=\frac{\phi _0}{\sqrt[3]{\frac{\mathrm{\Sigma }_1^2}{2\mathrm{\Lambda }}\mathrm{sinh}(\pm 3\sqrt{A_1}t)}},$$
(4.41)
$$\psi (t)=\frac{\mathrm{cosh}(\pm 3\sqrt{A_1})t}{3\sqrt{A_1}\mathrm{sinh}^2(\pm \sqrt{A_1}t)}\frac{1}{6\sqrt{A_1}}\mathrm{ln}\mathrm{tanh}\left(\frac{\pm 3\sqrt{A_1}t}{2}\right)+\psi _0,$$
(4.42)
which asymptotically gives a de Sitter behaviour. A last interesting case is 8 in Tab.I, which can be assigned by the functions
$$F=F_0,V(\phi )=\mathrm{\Lambda },k=1,$$
$$G(\psi )=G_0+\frac{\psi ^2}{16}+\left(\frac{\phi _0^2}{12\psi _0^4}\right)\psi ^4,W(\psi )=\frac{\psi ^4}{4\psi _0^2}\mathrm{\Lambda },$$
(4.43)
being $`G(\psi )`$ and $`W(\psi )`$ free for the Noether symmetry. The model is relevant for hyper-extended inflation (see e.g. ). Power law solutions like in are easily found.
The cases in Tab.II are essentially subcases of those discussed above.
## 5 Inflation and graceful exit
As we discussed in Introduction, the goal to get a sufficient inflationary period and then to exit from it, without imposing any particular fine tuning, can be achieved by assuming the variation of the bubble nucleation rate $`ϵ`$. However, we are taking into account a first order phase transition after which we recover a Friedman stage . In principle, being $`ϵ=\lambda /H^4`$, we can expect the variation of both $`\lambda `$ and $`H`$. The form of $`\lambda `$ strictly depends on the form of the theory and, as it is discussed in , it is time independent toward the late times, if we are dealing with a Brans–Dicke theory. In our cases, by using (1.7), we get that most of the couplings selected by the existence of Noether symmetry can be recast in a Brans–Dicke–like form.
In spite of the variation of the effective gravitational coupling, it is reasonable to assume $`\lambda `$ to be approximatively constant , so that the mechanism of the graceful exit can be essentially connected to the variation of $`H`$.
However, this argument does not work for more general classes of theories, as hyperextended inflation , where $`\omega (\varphi )`$ is not a constant. Among the cases in Tabs. I and II, we have also couplings of the form
$$\omega (\varphi )=\frac{F(\phi )}{2(dF/d\phi )^2}=\frac{\frac{1}{12}\phi ^2+F_0^{}\phi +F_0}{2\left[\frac{1}{6}\phi +F_0^{}\right]^2},$$
(5.1)
see e.g. the cases 5 in Tab.I and 4 in Tab.II. This situation deserves more attention since we can distinguish a regime where we match a sort of hyperextended inflation $`(\phi 0)`$ and a regime where the extended inflationary scheme is recovered $`(\phi \mathrm{})`$. In any case, the microwave background bounds have to be satisfied, as discussed in . Furthermore, taking into account double–field models, the contributions to $`\lambda `$ come from $`\phi `$ and $`\psi `$. Being both fields nonminimally coupled, and from the forms of couplings selected by the Noether symmetry, we are dealing with a double Brans–Dicke–like theory where the extended inflationary mechanism is improved. Looking at the solutions of previous section, we can have double inflationary stages ruled by the parameters of couplings and self–interaction potentials (see e.g. (4.19)–(4.22) or (4.26)). This situation is extremely interesting since ”very” large scale structure and large scale structure can be selected by these inflationary phases. In fact, we can have ”two” first–order phase transitions and then ”two” bubble nucleations where the size of bubbles is given by the coupling parameters. In other words, we can expect two graceful exits given by the superposition of two extended inflationary phases. To be more precise, at a given time $`t>t_0`$ after nucleation, the ”comoving” bubble radius is
$$r(t,t_0)=_{t_0}^t𝑑ta(t)^1,$$
(5.2)
while the ”physical” size of the bubble is
$$(t,t_0)=a(t)r(t,t_0).$$
(5.3)
When $`t\mathrm{}`$, the form of $`a(t)`$ selects the size of the bubble. In the cases (4.19),(4.22) and (4.23), this size is finite since the asymptotic behaviour is $`a(t)\mathrm{exp}H_0t`$. For power law behaviours, the growth of the bubble size is linear.
Besides, we have a variation of the Hubble parameter in most of the cases we dealt with: in (4.19), (4.21), and (4.40), it converges to a constant for $`t\mathrm{}`$, in (4.23), it is exactly a constant, in the other cases, it is $`Ht^1`$.
Graceful exit is achieved if, being $`\lambda `$ a constant, $`ϵ`$ is less than $`ϵ_{crit}`$ during inflation and, after bubble nucleation, $`ϵ>ϵ_{crit}`$. In our cases, $`H`$ is the key parameter which governs the behaviour of $`ϵ`$. For power–law solutions, as for standard extended inflationary models, the graceful exit is easily recovered (see Eqs(4.26) and (4.27)). In the asymptotically exponential cases, the parameter $`ϵ`$ goes to a constant for $`t\mathrm{}`$ and the graceful exit problem has no solution. In fact, the function $`H`$, calculated from (4.19) and (4.21), is a sort of step function with two different constant values at $`t\pm \mathrm{}`$. The related $`H^4`$ has a singularity in the origin which does not allow a graceful exit from inflation. The situation for the solution (4.40) is similar.
## 6 Conclusions
In this paper, we derived exact cosmological solutions in double scalar–tensor gravity theories by the general approach of searching for Noether symmetries. This work generalizes those in . The couplings and the potentials of both scalar fields are connected with the existence of the symmetries, and the solutions of dynamics furnish power law or de Sitter evolutions. As a consequence, in all the above cases it is easy to calculate the bubble nucleation rate $`ϵ=\lambda /H^4`$ to test if one succeeds in graceful exit. Depending on the value of intervening parameters, this can be accomplished in several cases.
Furthermore, being in principle, both fields nonminimally coupled and self–interacting with a potential, their role is mixed and it is not possible to distinguish, a priori, a Brans–Dicke field and an inflaton field as in other extended inflationary models. This distinction seems, in our opinion, rather artificial in view of a stochastic approach to the fundamental interactions where the effective role of the various fields is distinguishable only in the low energy limit (see and reference therein) and there is no reason why a field should interact with the gravitational field and the other one not.
Another remark deserves double inflation which is ruled by the parameters of the theory and, then, by the Noether symmetry. As it is well known this feature is of extreme interest in perturbation theory since it can furnish the seeds for the formation of structures at large and at very large scales. As we have seen, it is very common in our approach and it could contribute to the enhancement of a successful extended inflation.
Finally, we want to stress the fact that the standard Newtonian coupling can be recovered in several of the above models, being
$$G_{eff}=\frac{1}{2[F(\phi )+G(\psi )]},$$
(6.1)
so that as soon as $`F(\phi )F_0`$ and/or $`G(\psi )G_0`$, general relativity is restored (in our units $`F_0+G_01/2`$) and both fields can contribute to its recovering. This means that in an accurate setting of the models, one could succeed both in graceful exit and in recovering of standard gravity.
Table I – Symmetries in double nonminimally coupled models.
| N. | $`\alpha `$ | $`\beta `$ | $`\gamma `$ | $`F(\phi )`$ | $`G(\psi )`$ | $`V(\phi )`$ | $`W(\psi )`$ | $`k`$ |
| --- | --- | --- | --- | --- | --- | --- | --- | --- |
| 1 | $`a`$ | $`3\phi /2`$ | $`3\psi /2`$ | $`F_0\phi ^2`$ | $`G_0\psi ^2`$ | $`V_0\phi ^2`$ | $`W_0\psi ^2`$ | 0 |
| 2 | $`a`$ | $`3\phi /2`$ | $`3\psi /2`$ | $`F_0\phi ^2`$ | $`G_0\psi ^2`$ | $`\mathrm{\Lambda }`$ | $`\mathrm{\Lambda }`$ | 0 |
| 3 | $`a`$ | $`3\phi /2`$ | $`3\psi /2`$ | $`F_0\phi ^2`$ | $`G_0\psi ^2`$ | 0 | $`W_0\psi ^2`$ | 0 |
| 4 | 0 | 0 | $`\gamma _0`$ | $`F_0`$ | $`G_0`$ | $`\mathrm{\Lambda }_1`$ | $`\mathrm{\Lambda }_2`$ | $`k`$ |
| 5 | 0 | $`1/a`$ | $`\gamma _0`$ | $`\frac{1}{12}\phi ^2+F_0^{}\phi +F_0`$ | $`G_0`$ | $`\mathrm{\Lambda }_1`$ | $`\mathrm{\Lambda }_2`$ | 0 |
| 6 | 0 | $`1/a`$ | $`\gamma _0`$ | $`\frac{1}{12}\phi ^2+F_0^{}\phi +F_0`$ | $`G_0`$ | $`0`$ | $`\mathrm{\Lambda }`$ | 0 |
| 7 | 0 | $`1/a`$ | $`\gamma _0`$ | $`\frac{1}{12}\phi ^2`$ | $`G_0`$ | $`\mathrm{\Lambda }_1`$ | $`\mathrm{\Lambda }_2`$ | 0 |
| 8 | 0 | $`\beta _0`$ | 0 | $`F_0`$ | $`G(\psi )`$ | $`0,\mathrm{\Lambda }`$ | $`W(\psi )`$ | $`k`$ |
| 9 | 0 | $`1/a`$ | 0 | $`\frac{1}{12}\phi ^2+F_0^{}\phi +F_0`$ | $`G(\psi )`$ | $`0,\mathrm{\Lambda }`$ | $`W(\psi )`$ | 0 |
Table II – Symmetries in single nonminimally coupled models.
| N. | $`\alpha `$ | $`\beta `$ | $`\gamma `$ | $`F(\phi )`$ | $`V(\phi )`$ | $`W(\psi )`$ | $`k`$ |
| --- | --- | --- | --- | --- | --- | --- | --- |
| 1 | $`a`$ | $`3/2\phi `$ | $`3\psi /2`$ | $`F_0\phi ^2`$ | $`V_0\phi ^2`$ | $`W_0\psi ^2`$ | 0 |
| 2 | $`a`$ | $`3\phi /2`$ | $`3\psi /2`$ | $`F_0\phi ^2`$ | $`0`$ | $`W_0\psi ^2`$ | 0 |
| 3 | $`a`$ | $`3\phi /2`$ | $`3\psi /2`$ | $`F_0\phi ^2`$ | $`\mathrm{\Lambda }`$ | $`\mathrm{\Lambda }`$ | 0 |
| 4 | 0 | $`1/a`$ | $`\gamma _0`$ | $`\frac{1}{12}\phi ^2+F_0^{}\phi +F_0`$ | 0 | $`\mathrm{\Lambda }`$ | 0 |
| 5 | 0 | $`1/a`$ | $`\gamma _0`$ | $`\frac{1}{12}\phi ^2+F_0^{}\phi +F_0`$ | $`\mathrm{\Lambda }_1`$ | $`\mathrm{\Lambda }_2`$ | 0 |
| 6 | 0 | 0 | $`\gamma _0`$ | $`F_0`$ | $`\mathrm{\Lambda }_1`$ | $`\mathrm{\Lambda }_2`$ | $`k`$ |
Acknowledgments
The authors thank the referee for the useful comments. Research supported by MURST fund 40% and 60% art. 65 D.P.R. 382/80. GL acknowledges UE (P.O.M. 1994/1999) for financial support.
|
warning/0003/hep-th0003157.html
|
ar5iv
|
text
|
# Supersymmetric black hole solutions with 𝑅²-interactions SPIN-00/10, ITP-UU-00/07, hep-th/0003157
## 1 Introduction
In this note we briefly describe a class of extremal static multi-center black hole solutions arising in four-dimensional $`N=2`$ supergravity theories with terms quadratic in the Weyl tensor. These configurations preserve $`N=1`$ supersymmetry. They are determined in terms of harmonic functions associated with the electric and magnetic charges carried by the black holes. We refer to an upcoming publication for a detailed description of the construction of these solutions.
## 2 Supersymmetry transformation <br>rules
The $`N=2`$ supergravity theories that we consider are based on vector multiplets and hypermultiplets coupled to the supergravity fields and contain the standard Einstein-Hilbert action as well as terms quadratic in the Riemann tensor. To describe such theories in a transparent way we make use of the superconformal multiplet calculus , which incorporates the gauge symmetries of the $`N=2`$ superconformal algebra. The corresponding high degree of symmetry allows for the use of relatively small field representations. One is the Weyl multiplet, whose fields comprise the gauge fields corresponding to the superconformal symmetries and a few auxiliary fields. The other are abelian vector multiplets and hypermultiplets, as well as a general chiral supermultiplet. The latter will be treated as independent in initial stages of the analysis but at the end will be expressed in terms of the fields of the Weyl multiplet. Some of the additional (matter) multiplets will provide compensating fields which are necessary in order that the superconformal action becomes gauge equivalent to a Poincaré supergravity theory. The compensating fields bridge the deficit in degrees of freedom between the Weyl multiplet and the Poincaré supergravity multiplet. For instance, the graviphoton, represented by an abelian vector field in the Poincaré supergravity multiplet, is provided by an $`N=2`$ superconformal vector multiplet.
It is possible to analyze the conditions for residual $`N=1`$ supersymmetry directly in this superconformal setting, postponing a transition to Poincaré supergravity till the end. This implies in particular that our intermediate results are subject to local scale transformations. Only towards the end we will convert to expressions that are scale invariant. We will use this strategy in the following in order to construct black hole solutions with residual $`N=1`$ supersymmetry. This is exactly the same strategy we employed when we determined $`N=2`$ supersymmetric backgrounds in the presence of $`R^2`$-interactions .
The superconformal algebra contains general-coordinate, local Lorentz, dilatation, special conformal, chiral SU(2) and U(1), supersymmetry ($`Q`$) and special supersymmetry ($`S`$) transformations. The gauge fields associated with general-coordinate transformations ($`e_\mu ^a`$), dilatations ($`b_\mu `$), chiral symmetry ($`𝒱_{\mu j}^i,A_\mu `$) and $`Q`$-supersymmetry ($`\psi _\mu ^i`$), are realized by independent fields. The remaining gauge fields of Lorentz ($`\omega _\mu ^{ab}`$), special conformal ($`f_\mu ^a`$) and $`S`$-supersymmetry transformations ($`\varphi _\mu ^i`$) are dependent fields. They are composite objects, which depend in a complicated way on the independent fields . The corresponding curvatures and covariant fields are contained in a tensor chiral multiplet, which comprises $`24+24`$ off-shell degrees of freedom; in addition to the independent superconformal gauge fields it contains three auxiliary fields: a Majorana spinor doublet $`\chi ^i`$, a scalar $`D`$ and a selfdual Lorentz tensor $`T_{abij}`$ (where $`i,j,\mathrm{}`$ are chiral SU(2) spinor indices). We summarize the transformation rules for some of the independent fields of the Weyl multiplet under $`Q`$\- and $`S`$-supersymmetry and under special conformal transformations, with parameters $`ϵ^i`$, $`\eta ^i`$ and $`\mathrm{\Lambda }_\mathrm{K}^a`$, respectively,
$`\delta e_\mu ^a`$ $`=`$ $`\overline{ϵ}^i\gamma ^a\psi _{\mu i}+\mathrm{h}.\mathrm{c}.,`$
$`\delta \psi _\mu ^i`$ $`=`$ $`2𝒟_\mu ϵ^i\frac{1}{8}T^{abij}\gamma _{ab}\gamma _\mu ϵ_j\gamma _\mu \eta ^i,`$
$`\delta b_\mu `$ $`=`$ $`\frac{1}{2}\overline{ϵ}^i\varphi _{\mu i}\frac{3}{4}\overline{ϵ}^i\gamma _\mu \chi _i\frac{1}{2}\overline{\eta }^i\psi _{\mu i}+\mathrm{h}.\mathrm{c}.`$ (2.1)
$`+\mathrm{\Lambda }_\mathrm{K}^ae_\mu ^a,`$
$`\delta A_\mu `$ $`=`$ $`\frac{1}{2}i\overline{ϵ}^i\varphi _{\mu i}+\frac{3}{4}i\overline{ϵ}^i\gamma _\mu \chi _i+\frac{1}{2}i\overline{\eta }^i\psi _{\mu i}+\mathrm{h}.\mathrm{c}.,`$
$`\delta T_{ab}^{ij}`$ $`=`$ $`8\overline{ϵ}^{[i}R(Q)_{ab}^{j]},`$
$`\delta \chi ^i`$ $`=`$ $`\frac{1}{12}\gamma _{ab}D\text{/}T^{abij}ϵ_j+\frac{1}{6}R(𝒱)_{ab}{}_{j}{}^{i}\gamma _{}^{ab}ϵ^j`$
$`\frac{1}{3}iR(A)_{ab}\gamma ^{ab}ϵ^i+Dϵ^i+\frac{1}{12}T_{ab}^{ij}\gamma ^{ab}\eta _j,`$
where $`𝒟_\mu `$ are derivatives covariant with respect to Lorentz, dilatational, U(1) and SU(2) transformations, whereas $`D_\mu `$ are derivatives covariant with respect to all superconformal transformations. The quantities $`R(Q)_{\mu \nu }^i`$, $`R(A)_{\mu \nu }`$ and $`R(𝒱)_{\mu \nu }_j^i`$ are supercovariant curvatures related to $`Q`$-supersymmetry, U(1) and SU(2) transformations. We suppress terms of higher order in the fermions throughout this paper, as we will be dealing with a bosonic background.
Let us now turn to the abelian vector multiplets, labelled by an index $`I=0,1,\mathrm{},n`$. For each value of the index $`I`$, there are $`8+8`$ off-shell degrees of freedom, residing in a complex scalar $`X^I`$, a doublet of chiral fermions $`\mathrm{\Omega }_i^I`$, a vector gauge field $`W_\mu ^I`$, and a real SU(2) triplet of scalars $`Y_{ij}^I`$. Under $`Q`$\- and $`S`$-supersymmetry the fields $`X^I`$ and $`\mathrm{\Omega }_i^I`$ transform as follows:
$`\delta X^I`$ $`=`$ $`\overline{ϵ}^i\mathrm{\Omega }_i^I,`$
$`\delta \mathrm{\Omega }_i^I`$ $`=`$ $`2D\text{/}X^Iϵ_i+\frac{1}{2}\epsilon _{ij}(F_{\mu \nu }^I\frac{1}{4}\epsilon _{kl}T_{\mu \nu }^{kl}\overline{X}^I)\gamma ^{\mu \nu }ϵ^j`$ (2.2)
$`+Y_{ij}^Iϵ^j+2X^I\eta _i,`$
where the quantity $`F_{\mu \nu }^I`$ denotes the anti-selfdual part of the abelian field strength $`F_{\mu \nu }^I=2_{[\mu }W_{\nu ]}^I`$.
The covariant quantities of the vector multiplet constitute a reduced chiral multiplet. A general chiral multiplet comprises $`16+16`$ off-shell degrees of freedom and carries an arbitrary Weyl weight $`w`$ (corresponding to the Weyl weight of its lowest component). The covariant quantities of the Weyl multiplet also constitute a reduced chiral multiplet, denoted by $`W^{abij}`$, whose lowest-$`\theta `$ component is the tensor $`T^{abij}`$. From this multiplet one may form a scalar (unreduced) chiral multiplet $`W^2=[W^{abij}\epsilon _{ij}]^2`$ which has Weyl and chiral weights $`w=2`$ and $`c=2`$, respectively .
In the following, we will also allow for the presence of an arbitrary chiral background superfield , whose component fields will be indicated with a caret. We denote its bosonic component fields by $`\widehat{A}`$, $`\widehat{B}_{ij}`$, $`\widehat{F}_{ab}^{}`$ and by $`\widehat{C}`$. Here $`\widehat{A}`$ and $`\widehat{C}`$ denote complex scalar fields, appearing at the $`\theta ^0`$\- and $`\theta ^4`$-level of the chiral background superfield, respectively, while the symmetric complex SU(2) tensor $`\widehat{B}_{ij}`$ and the anti-selfdual Lorentz tensor $`\widehat{F}_{ab}^{}`$ reside at the $`\theta ^2`$-level. The fermion fields at level $`\theta `$ and $`\theta ^3`$ are denoted by $`\widehat{\mathrm{\Psi }}_i`$ and $`\widehat{\mathrm{\Lambda }}_i`$. Under $`Q`$\- and $`S`$-supersymmetry $`\widehat{A}`$ and $`\mathrm{\Psi }_i`$ transform as
$`\delta \widehat{A}`$ $`=`$ $`\overline{ϵ}^i\widehat{\mathrm{\Psi }}_i,`$
$`\delta \widehat{\mathrm{\Psi }}_i`$ $`=`$ $`2D\text{/}\widehat{A}ϵ_i+\frac{1}{2}\epsilon _{ij}\widehat{F}_{ab}\gamma ^{ab}ϵ^j+\widehat{B}_{ij}ϵ^j`$ (2.3)
$`+2w\widehat{A}\eta _i,`$
where $`w`$ denotes the Weyl weight of the background superfield. Eventually this multiplet will be identified with $`W^2`$ in order to generate the $`R^2`$-terms in the action. This identification implies the following relations , which we will need in due time,
$`\widehat{A}`$ $`=`$ $`(\epsilon _{ij}T_{ab}^{ij})^2,`$
$`\widehat{\mathrm{\Psi }}_i`$ $`=`$ $`16\epsilon _{ij}R(Q)_{ab}^jT^{klab}\epsilon _{kl},`$
$`\widehat{B}_{ij}`$ $`=`$ $`16\epsilon _{k(i}R(𝒱)^k{}_{j)ab}{}^{}T_{}^{lmab}\epsilon _{lm}`$
$`64\epsilon _{ik}\epsilon _{jl}\overline{R}(Q)_{ab}^kR(Q)^{lab},`$
$`\widehat{F}^{ab}`$ $`=`$ $`16(M)_{cd}{}_{}{}^{ab}T_{}^{klcd}\epsilon _{kl}`$
$`16\epsilon _{ij}\overline{R}(Q)_{cd}^i\gamma ^{ab}R(Q)^{jcd},`$
$`\widehat{\mathrm{\Lambda }}_i`$ $`=`$ $`32\epsilon _{ij}\gamma _{ab}R(Q)_{cd}^j(M)_{cd}^{ab}`$
$`+16((S)_{abi}+3\gamma _{[a}D_{b]}\chi _i)T^{klab}\epsilon _{kl}`$
$`64R(𝒱)_{ab}{}_{}{}^{k}{}_{i}{}^{}\epsilon _{kl}^{}R(Q)_{ab}^l,`$
$`\widehat{C}`$ $`=`$ $`64(M)_{cd}^{}{}_{}{}^{ab}(M)_{cd}^{}^{ab}`$ (2.4)
$`+32R(𝒱)_{ab}^{}{}_{}{}^{k}{}_{l}{}^{}R(𝒱)_{ab}^{}{}_{}{}^{l}_{k}^{}`$
$`32T^{ijab}D_aD^cT_{cbij}`$
$`+128\overline{}(S)_i^{ab}R(Q)_{ab}^i`$
$`+384\overline{R}(Q)^{abi}\gamma _aD_b\chi _i.`$
We refer to for a precise definition of the various curvature tensors. The derivatives $`D_a`$ are superconformally covariant.
In the presence of a chiral background superfield, the coupling of the abelian vector multiplets to the Weyl multiplet is encoded in a function $`F(X^I,\widehat{A})`$, which is holomorphic and homogenous of degree two,
$`X^IF_I+w\widehat{A}F_{\widehat{A}}=2F,`$
$`F_I=_{X^I}F,F_{\widehat{A}}=_{\widehat{A}}F.`$ (2.5)
The field equations of the vector multiplets are subject to equivalence transformations corresponding to electric-magnetic duality, which will not involve the fields of the Weyl multiplet and of the chiral background. As is well-known, two complex $`(2n+2)`$-component vectors can be defined which transform linearly under the SP$`(2n+2;𝐑)`$ duality group, namely
$`\left(\begin{array}{c}X^I\\ \\ F_I(X,\widehat{A})\end{array}\right)\text{ and}\left(\begin{array}{c}F_{\mu \nu }^{+I}\\ \\ G_{\mu \nu I}^+\end{array}\right).`$ (2.6)
The first vector has weights $`w=1`$ and $`c=1`$, whereas the second one has zero Weyl and chiral weights. The field strengths $`G_{\mu \nu I}^\pm `$ are defined as follows:
$`G_{\mu \nu I}^+=\overline{F}_{IJ}F_{\mu \nu }^{+J}+𝒪_{\mu \nu I}^+,`$
$`G_{\mu \nu I}^{}=F_{IJ}F_{\mu \nu }^J+𝒪_{\mu \nu I}^{},`$ (2.7)
where
$`𝒪_{\mu \nu I}^+=\frac{1}{4}(F_I\overline{F}_{IJ}X^J)T_{\mu \nu ij}\epsilon ^{ij}+\widehat{F}_{\mu \nu }^+\overline{F}_{I\widehat{A}}.`$
(2.8)
They appear in the field equations of the vector fields. Eventually we will solve the Bianchi identities, $`𝒟^\mu (F^{}F^+)_{\mu \nu }^I=0`$, and the field equations, $`𝒟^\mu (G^{}G^+)_{\mu \nu }^I=0`$, for a given configuration of magnetic and electric charges in a static spacetime geometry with the chiral background turned on. These charges, which will be denoted by $`(p^I,q_I)`$, comprise a symplectic vector.
Next, let us introduce a particular spinor that transforms inhomogenously under $`S`$-supersymmetry transformations. This spinor is given by
$`\zeta _i^V\left(\mathrm{\Omega }_i^I{\displaystyle \frac{}{X^I}}+\widehat{\mathrm{\Psi }}_i{\displaystyle \frac{}{\widehat{A}}}\right)𝒦`$ (2.9)
$`=i\mathrm{e}^𝒦\left[(\overline{F}_I\overline{X}^JF_{IJ})\mathrm{\Omega }_i^I\overline{X}^IF_{I\widehat{A}}\widehat{\mathrm{\Psi }}_i\right],`$
where we introduced the symplectically covariant factor (with $`w=2`$ and $`c=0`$),
$$\mathrm{e}^𝒦=i\left[\overline{X}^IF_I(X,\widehat{A})\overline{F}_I(\overline{X},\overline{\widehat{A}})X^I\right],$$
(2.10)
which resembles (but is not equal to) the Kähler potential in special geometry. It can be shown, using the results contained in , that $`\zeta _i^V`$ transforms covariantly under symplectic reparametrizations. Under $`Q`$\- and $`S`$-supersymmetry $`\zeta _i^V`$ transforms as (ignoring higher-order fermionic terms)
$`\delta \zeta _i^V`$ $`=`$ $`\mathrm{e}^𝒦𝒟\text{/}\mathrm{e}^𝒦ϵ_i+2i𝒜\text{/}ϵ_i\frac{1}{2}i\epsilon _{ij}_{\mu \nu }^{}\gamma ^{\mu \nu }ϵ^j`$ (2.11)
$`+\mathrm{e}^𝒦[(\overline{F}_I\overline{X}^J\overline{F}_{IJ})N^{IK}F_{KA}\widehat{B}_{ij}`$
$`(\overline{F}_I\overline{X}^JF_{IJ})N^{IK}\overline{F}_{KA}\epsilon _{ik}\epsilon _{jl}\widehat{B}^{kl}]ϵ^j`$
$`+2\eta _i,`$
where
$`𝒜_\mu `$ $`=`$ $`\frac{1}{2}e^𝒦\left(\overline{X}^J\underset{\mu }{\overset{}{𝒟}}F_J\overline{F}_J\underset{\mu }{\overset{}{𝒟}}X^J\right),`$
$`_{\mu \nu }^{}`$ $`=`$ $`\mathrm{e}^𝒦\left(\overline{F}_IF_{\mu \nu }^I\overline{X}^IG_{\mu \nu I}^{}\right).`$ (2.12)
In arriving at (2.11) we have used the field equations for the auxiliary fields $`Y_{ij}^I`$ , which is necessary for $`\delta \zeta _i^V`$ to take a symplectically covariant form.
We note that $`\delta \zeta _i^V`$ is not the only spinor that can be constructed which transform inhomogenously under $`S`$-supersymmetry transformations. Another such spinor, which we denote by $`\zeta _i^H`$, is constructed out of the hypermultiplet fermions . It transforms as follows under $`Q`$\- and $`S`$-supersymmetry,
$`\delta \zeta _i^\mathrm{H}=\frac{1}{2}\chi ^1𝒟\text{/}\chi \epsilon _{ij}ϵ^j+k\text{/}_{ij}ϵ^j+\epsilon _{ij}\eta ^j,`$ (2.13)
where $`\chi `$ denotes the hyper-Kähler potential and where $`k_{\mu ij}`$ denotes a quantity that is symmetric in $`i,j`$ , but whose explicit form is not important here.
Since the $`\zeta _i`$ transform inhomogenously under $`S`$-supersymmetry, they can act as compensators for this symmetry. This observation is relevant when constructing supersymmetric backgrounds, where one requires (some of) the $`Q`$-supersymmetry variations of the spinors (as well as of derivatives of the spinors) to vanish modulo a uniform $`S`$-transformation. This can conveniently be done by considering $`S`$-invariant spinors, constructed by employing $`\zeta _i`$. Relevant examples of such spinors are, for instance, $`\mathrm{\Omega }_i^IX^I\zeta _i^V`$ and $`\widehat{\mathrm{\Psi }}_iw\widehat{A}\zeta _i^V`$.
## 3 The ansatz
We seek to construct static multi-center black hole solutions with $`N=1`$ residual supersymmetry. For the line element we make the ansatz
$`ds^2=e^{2g(\stackrel{}{x})}dt^2+e^{2f(\stackrel{}{x})}d\stackrel{}{x}^2.`$ (3.1)
We impose the following restriction on the $`Q`$-supersymmetry transformation parameter $`ϵ_i`$,
$`ϵ_i=h\epsilon _{ij}\gamma _0ϵ^j,`$ (3.2)
where $`h(\stackrel{}{x})`$ denotes a phase factor of chiral weight $`c=1`$. The condition (3.2) is covariant with respect to SU(2) and spatial rotations. The multi-center solutions that we wish to construct have the feature that, when the centers are made to coincide, they lead to one-center solutions that are invariant with respect to SU(2) and spatial rotations. The latter satisfy condition (3.2).
In addition to (3.2), we impose that $`𝒜_\mu =0`$ as well as $`_{\mu \nu }^{}=0`$.
We denote the magnetic and electric charges associated to each center by $`(p_A^I,q_{AI})`$. In the geometry (3.1) the Bianchi identities and field equations for the vector fields are solved by
$`F_{ti}^IF_{ti}^{+I}`$ $`=`$ $`ie^{gf}_iH^I,`$
$`G_{tiI}^{}G_{tiI}^+`$ $`=`$ $`ie^{gf}_iH_I,`$ (3.3)
where $`H^I`$ and $`H_I`$ denote harmonic functions given by
$`H^I`$ $`=`$ $`{\displaystyle \underset{A}{}}\left(h_A^I+{\displaystyle \frac{p_A^I}{|\stackrel{}{x}\stackrel{}{x}_A|}}\right),`$
$`H_I`$ $`=`$ $`{\displaystyle \underset{A}{}}\left(h_{AI}+{\displaystyle \frac{q_{AI}}{|\stackrel{}{x}\stackrel{}{x}_A|}}\right).`$ (3.4)
Here $`h_A^I`$ and $`h_{AI}`$ denote integration constants. We now identify $`\widehat{A}`$ with $`(\epsilon _{ij}T^{abij})^2`$ so that we are dealing with black hole solutions in the presence of $`R^2`$-interactions.
## 4 Static multi-center solutions
It will be convenient to use rescaled variables $`Y^I=\mathrm{e}^{𝒦/2}\overline{\mathrm{\Sigma }}X^I`$ and $`\mathrm{{\rm Y}}=\mathrm{e}^𝒦\overline{\mathrm{\Sigma }}^2\widehat{A}`$. Here $`\overline{\mathrm{\Sigma }}`$ is taken to have weights $`c=1`$ and $`w=0`$ so that $`Y^I`$ and $`\mathrm{{\rm Y}}`$ have vanishing chiral and Weyl weights. Then, from (2.10) and from (2.12), we obtain
$`|\mathrm{\Sigma }|^2`$ $`=`$ $`i\left[\overline{Y}^IF_I(Y,\mathrm{{\rm Y}})\overline{F}_I(\overline{Y},\overline{\mathrm{{\rm Y}}})Y^I\right],`$
$`A_\mu `$ $`=`$ $`\frac{i}{2}_\mu \mathrm{log}\frac{\mathrm{\Sigma }}{\overline{\mathrm{\Sigma }}}𝒜_\mu ^Y,`$ (4.1)
$`𝒜_\mu ^Y`$ $`=`$ $`\frac{1}{2}\frac{1}{|\mathrm{\Sigma }|^2}\left(\overline{Y}^J\underset{\mu }{\overset{}{}}F_J\overline{F}_J\underset{\mu }{\overset{}{}}Y^J\right).`$
Using (3.1), we find that the vanishing of the $`Q`$-supersymmetry variation (subject to (3.2)) of the various $`S`$-invariant spinors yields a number of restrictions on the $`N=1`$ background, as follows. In a Poincaré frame (where $`b_\mu =0`$ and $`𝒦=\mathrm{const}.`$) we find that
$`e^{2g}`$ $`=`$ $`e^{2f}=\mathrm{e}^𝒦|\mathrm{\Sigma }|^2,`$
$`\mathrm{e}^{𝒦/2}\overline{\mathrm{\Sigma }}T_{ti}^{}`$ $`=`$ $`4_if,\mathrm{{\rm Y}}=64(_if)^2,`$
$`F_{ti}^I`$ $`=`$ $`_i\left[e^{2f}(Y^I+\overline{Y}^I)\right],`$
$`𝒜_\mu ^Y`$ $`=`$ $`0.`$ (4.2)
The symplectic vector $`(Y^I,F_I(Y,\mathrm{{\rm Y}}))`$ is determined in terms of the symplectic vector $`(H^I,H_I)`$ as follows:
$`_i\left(\begin{array}{c}Y^I\overline{Y}^I\\ \\ F_I(Y,\mathrm{{\rm Y}})\overline{F}_I(\overline{Y},\overline{\mathrm{{\rm Y}}})\end{array}\right)=i_i\left(\begin{array}{c}H^I\\ \\ H_I\end{array}\right).`$
Thus, we see that as one approaches the individual centers ($`|\stackrel{}{x}\stackrel{}{x}_A|0`$) the scalar fields $`Y^I`$ are entirely determined in terms of the charges associated to these individual centers. This behaviour, namely that the values of the $`Y^I`$ near the centers are independent of the constants $`h_A^I`$ and $`h_{AI}`$ which determine the values of the $`Y^I`$ far away from the centers, is the same that has been observed without $`R^2`$-interactions .
In addition, we find that $`h=\overline{\mathrm{\Sigma }}/|\mathrm{\Sigma }|`$, whereas the field $`D`$ is determined to be $`D=\frac{1}{3}R`$. The SU(2) curvature $`R(𝒱)_{\mu \nu }^i_j`$, and hence $`\widehat{B}_{ij}`$, vanishes. Since $`𝒜_\mu ^Y=0`$, also the U(1) curvature $`R(A)_{\mu \nu }`$ vanishes.
The solution given above describes a static multi-center black hole with residual $`N=1`$ supersymmetry in the presence of $`R^2`$-interactions. It approaches flat Minkowski spacetime at spatial infinity. Setting $`\mathrm{e}^𝒦|\mathrm{\Sigma }|_{\mathrm{}}^2=1`$ expresses its ADM mass as $`M_{\mathrm{ADM}}=\mathrm{e}^𝒦_A(p_A^IF_I(Y_{\mathrm{}})q_{AI}Y_{\mathrm{}}^I)`$.
In the case of one center, the solution interpolates between two $`N=2`$ supersymmetric vacua : flat spacetime at spatial infinity and Bertotti-Robinson spacetime at the horizon. When switching off $`R^2`$-interactions this solution agrees with the one constructed in . Its macroscopic entropy is given by
$`𝒮=\pi r^2\left[|\mathrm{\Sigma }|^2+4\mathrm{Im}\left(\mathrm{{\rm Y}}F_\mathrm{{\rm Y}}(Y,\mathrm{{\rm Y}})\right)\right]|_{r=0}.`$ (4.4)
## 5 Outlook
The static multi-center solution (4.2) can now be used as the starting point for computing the metric on the moduli space of four-dimensional BPS black holes in the presence of $`R^2`$-interactions. In the absence of $`R^2`$-interactions, it was found that the moduli metric of electrically charged BPS black holes is determined in terms of a moduli potential $`\mu `$ given by $`\mu =d^3xe^{4f}`$. It was furthermore established that for small black hole separations the associated one-dimensional Lagrangian describing the slow-motion of these BPS black holes exhibits an enhanced superconformal symmetry. It was suggested that it should be possible to reproduce the macroscopic entropy of BPS black holes by performing a state counting in this superconformal quantum mechanics model. This would imply that the degeneracy of states of such a model is encoded in the moduli potential $`\mu `$. In view of the formula (4.4) for the macroscopic entropy one thus expects that $`\mu `$ will receive corrections steming from $`R^2`$-interactions. This is indeed likely to be the case, since the one-dimensional Lagrangian describing the slow-motion of the black holes will now be derived from a four-dimensional action containing $`R^2`$-interactions. Schematically, since the four-dimensional Einstein-Maxwell action gives rise to a moduli potential $`\mu =d^3xe^{4f}`$, the term $`\mathrm{Im}F_{\widehat{A}}R|T|^2`$, which appears in the four-dimensional Lagrangian with $`R^2`$-interactions, suggests a correction to the moduli potential of the form $`d^3xe^{2f}\mathrm{Im}(\mathrm{{\rm Y}}F_\mathrm{{\rm Y}})`$. In view of (4.4), this suggests that in the presence of $`R^2`$-terms the moduli potential $`\mu `$ will be given by $`\mu =d^3xe^{2f}[e^{2f}+4\mathrm{Im}(\mathrm{{\rm Y}}F_\mathrm{{\rm Y}}(Y,\mathrm{{\rm Y}}))]`$. This feature is currently under investigation .
## Acknowledgements
This work was supported in part by the European Commission TMR programme ERBFMRX-CT96-0045.
|
warning/0003/cs0003018.html
|
ar5iv
|
text
|
# Description of GADEL
## General Info
The system works on SUN/Solaris or PC/Linux with Sicstus Prolog3.7.x and C. It is written in Prolog and it generates a Sicstus library in C. The current version is about 3000 lines of Prolog.
## Description of the System
### Default Logic and Genetic Algorithms
Default Logic has been introduced by Reiter (?) in order to formalize common sense reasoning from incomplete information, and is now recognized as one of the most appropriate framework for non monotonic reasoning. In this formalism, knowledge is represented by a default theory $`(W,D)`$ where $`W`$ is a set of first order formulas representing the sure knowledge, and $`D`$ a set of *default rules* (or defaults). A *default* $`\delta =\frac{\alpha :\beta _1,\mathrm{},\beta _n}{\gamma }`$ is an inference rule providing conclusions relying upon given, as well as absent information meaning “if the *prerequisite* $`\alpha `$ is proved, and if for all $`i=1,\mathrm{},n`$ each *justification* $`\beta _i`$ is individually consistent (in other words if nothing proves its negation) then one concludes the *consequent* $`\gamma `$”. From a default theory $`(W,D)`$ one tries to build some extensions, that are maximal sets of plausible conclusions. Reiter has given the following pseudo iterative characterization of an extension $`E`$: we define
* $`E_0=W`$
* and for all $`k0`$,
$`E_{k+1}`$ $`=`$ $`Th\left(E_k\right)\{\gamma {\displaystyle \frac{\alpha :\beta _1,\mathrm{},\beta _n}{\gamma }}D,`$
$`\alpha E_k,\neg \beta _iE,i=1,\mathrm{},n\}`$
then, $`E`$ is an extension of $`(W,D)`$ iff $`E=_{k=0}^{\mathrm{}}E_k`$. The computation of an extension is known to be $`\mathrm{\Sigma }_2^pcomplete`$ (?). Even if the system DeRes (?) has very good performance on certain classes of default theories, there is no efficient system for general extension calculus. The aim of the present work is to describe the first implementation of the GADEL system (*Genetic Algorithms for DEfault Logic*) which applies Genetic Algorithms principles to propositional default reasoning (?).
Based on the principle of natural selection, Genetic Algorithms have been quite successfully applied to combinatorial problems such as scheduling or transportation problems. The key principle of this approach states that, species evolve through adaptations to a changing environment and that the gained knowledge is embedded in the structure of the population and its members, encoded in their chromosomes. If individuals are considered as potential solutions to a given problem, applying a genetic algorithm consists in generating better and better individuals. A genetic algorithm consists of the following components:
* a representation of the potential solutions in a chromosome, in most cases, a string of bits representing its genes,
* an initial population,
* an evaluation function which rates each potential solution w.r.t. the given problem,
* genetic operators that define the composition of the offsprings : two different operators will be considered : crossover allows to generate two new chromosomes (the offsprings) by crossing two chromosomes of the current population (the parents), mutation arbitrarily alters one or more genes of a selected chromosome,
* parameters : population size $`p_{size}`$ and probabilities of crossover $`p_c`$ and mutation $`p_m`$.
and an iteration process:
* evaluate each chromosomes,
* order the population according to evaluation rates and select the bests chromosomes,
* perform crossover and mutation on pairs of randomly selected chromosomes,
* repeat this full process until a user-defined number of populations has been explored.
The best chromosome of each population w.r.t. the evaluation function represents the current best solution to the problem.
Before a detail description of the GADEL system, it is necessary to give some arguments about the choice of the implementation language:
* In the previous section we have presented the most common version of GA but in fact each part of the system can take various forms. To develop easily a GA system the implementation language must be very flexible and modular.
* A GA system is an iterative system. The implementation language must be efficient.
* Default Logic is based on classical logic. To develop easily a GA system about Default Logic, the implementation language must be logic and symbolic.
From above, the choice of the implementation language is obvious: the most popular and efficient of the logic programming language, Prolog. We assume in the rest of the paper a minimal knowledge of Prolog. The article is organized as follows: section 2 presents the genetic algorithms aspects of GADEL, section 3 provides the process of compilation of a default theory to a prolog program, section 4 focuses on the evaluation function of GADEL and section 5 describes our experiments w.r.t. other existing systems.
### GADEL: a GA system
#### Representation and semantics.
Our purpose is to construct an extension of a given default theory $`(W,D)`$. For each default $`\frac{\alpha :\beta _1,\mathrm{},\beta _n}{\gamma }`$ we encode in the chromosome the prerequisite $`\alpha `$ and all justifications $`\beta _1,\mathrm{},\beta _n`$ conjointly. Given a set of defaults of size $`n`$ the chromosome will be of size $`2n`$. A candidate extension $`CE(G)`$ associated to a chromosome $`G`$ is :
$$CE(G)=Th(W\left\{\begin{array}{c}\gamma _i\frac{\alpha _i:\beta _i^1,\mathrm{},\beta _i^{k_i}}{\gamma _i}D\hfill \\ andG|_{2i1}=1andG|_{2i}=0\hfill \end{array}\right\})$$
#### Population tree.
According to the principles of Genetic Algorithms, we now consider a population of individuals representing candidate extensions. Usually chromosomes are strings of bits and population are sets of chromosomes. We have chosen a new representation for the population: binary trees. A population is defined inductively on the set of constructors $`\mathrm{\Lambda }`$, $`zero`$, $`one`$ and $`zeroone`$ of arity, respectively, 0, 1, 1, and 2. The two advantages of this representation are its compactness and unicity of each chromosome.
For example, the population $`\{0101,1000,1010\}`$ is represented by the tree in Figure 1.
### Compilation of default theory
A default theory is the given of a set of (propositional) formulas $`W`$ and a set of defaults $`D`$. Prerequisite, conclusion and justifications of a default are all (propositional) formulas. So the GADEL system needs a classical theorem prover. It must be efficient because it is applied on each chromosome at each new population. The obvious choice is to compile all these sets of formulas into a set of clauses. We have chosen the compilation to a disjunctive Prolog program.
#### Small introduction to disjunctive logic programming.
The theoretical basis of Prolog is the SLD-resolution for Horn clauses<sup>1</sup><sup>1</sup>1A Horn clause is a clause with at most one positive literal. It is not possible to directly insert disjunctive clauses in a Prolog program<sup>2</sup><sup>2</sup>2A clause is disjunctive if it contains at least two positive literals. Disjunctive logic programming (resp. disjunctive Prolog) is an “extension” of Horn logic programming (resp. Prolog) which allows disjunctions in the heads of definite<sup>3</sup><sup>3</sup>3A clause is definite if it contains one and only one positive literal clauses (resp. Prolog clauses). A way to handle disjunctive clauses is the case-analysis principle: a set of clauses $`\{CC^{},C_1,\mathrm{},C_m\}`$ is unsatisfiable if and only if the two sets of clauses $`\{C,C_1,\mathrm{},C_m\}`$ and $`\{C^{},C_1,\mathrm{},C_m\}`$ are also unsatisfiable. Disjunctive clauses $`(h_1\mathrm{}h_p\neg b_1\mathrm{}\neg b_n)`$ ($`h_k`$, $`1kp`$ and $`b_k`$, $`1kn`$ atoms) are then written $`(h_1\mathrm{}h_pb_1\mathrm{}b_n)`$. We have chosen the SLOU-resolution approach (?) (case-analysis and SLD-reduction) as the theoretical basis of SLOU Prolog, our implementation of disjunctive Prolog.
The strategy of SLOU Prolog applies case-analysis by necessity: the case-analysis is used only when a head $`h_i`$ of a disjunctive clause $`(h_1\mathrm{}h_pb_1\mathrm{}b_n)`$ is useful in the proof. This strategy needs that a negative clause<sup>4</sup><sup>4</sup>4A clause is negative if it contains only negative literals $`\neg b_1\mathrm{}\neg b_n`$ is written $`falseb_1\mathrm{}b_n`$ and $`false`$ becomes the goal. This choice assumes that almost all the clauses are definite (Prolog) clauses.
#### From set of defaults to disjunctive Prolog program.
Let us give a default $`\frac{\alpha _i:\beta _i^1,\mathrm{},\beta _i^{k_i}}{\gamma _i}`$, the $`i^{th}`$ of the default theory and a chromosome $`G`$. If $`G|_{2i1}=1`$ and $`G|_{2i}=0`$ then the default is supposed to be applied and $`\gamma `$ must be added to the candidate extension. Hence $`G`$ has to be a parameter of the theorem prover. In order to calculate $`(CE(G)\alpha _i)`$ or $`(j(CE(G)\neg \beta _i^j))`$, $`i`$ and $`(i,j)`$ also have to be parameters of the theorem prover. So a propositional variable $`h`$ is compiled in a disjunctive Prolog atom $`h(I,G)`$. We can now describe the compilation of the three parts of a default rule.
* Compilation of a conclusion: In order to calculate $`CE(G)`$, we must add $`\gamma _i`$ to the candidate extension if $`G|_{2i1}=1`$ and $`G|_{2i}=0`$. We first normalize $`\gamma _i`$ in a set of disjunctive clauses $`\{h_1\mathrm{}h_pb_1\mathrm{}b_n\}`$ and compile it in a disjunctive Prolog definition:
$`\{h_1(I,G);\mathrm{};h_p(I,G):`$
$`G|_{2I1}=1,G|_{2I}=0,b_1(I,G),\mathrm{},b_n(I,G)\}`$
* Compilation of a prerequisite: The function $`f`$ of the GADEL evaluation function compares $`G|_{2i1}`$ with $`CE(G)\alpha _i`$ (see section Evaluation function of GADEL). To prove it with a disjunctive logic program we prove $`CE(G),\neg \alpha _ifalse`$. We first normalize $`\neg \alpha _i`$ in a set of clauses $`\{h_1\mathrm{}h_pb_1\mathrm{}b_n\}`$ and compile it in a disjunctive Prolog definition:
$`\{h_1(I,G);\mathrm{};h_p(I,G):`$
$`I=i,b_1(I,G),\mathrm{},b_n(I,G)\}`$
* Compilation of justifications: The function $`f`$ of the GADEL evaluation function compares $`G|_{2i}`$ with $`j(CE(G)\neg \beta _i^j)`$. To prove it with a disjunctive logic programs we prove $`(j(CE(G),\beta _i^jfalse))`$. We first normalize $`\beta _i^j`$ in a set of clauses $`\{h_1\mathrm{}h_p`$$`b_1\mathrm{}b_n\},j`$, and compile it in a disjunctive Prolog definition:
$`\{h_1(I,G);\mathrm{};h_p(I,G):`$
$`I=(i,j),b_1(I,G),\mathrm{},b_n(I,G)\}`$
* Compilation of a formula $`\omega W`$: We first normalize $`\omega `$ in a set of clauses $`\{h_1\mathrm{}h_p`$ $`b_1\mathrm{}b_n\}`$ and compile it in a disjunctive Prolog definition:
$`\{h_1(I,G);\mathrm{};h_p(I,G):`$
$`b_1(I,G),\mathrm{},b_n(I,G)\}`$
#### From disjunctive Prolog to Prolog.
During the execution of a disjunctive Prolog program, some clauses are dynamic: case-analysis creates from a disjunctive clause $`(h_1\mathrm{}h_pb_1\mathrm{}b_n)`$ $`p`$ new clauses $`(h_1b_1\mathrm{}b_n)`$ and $`\{h_k\}_{1<kp}`$. A disjunctive clause is called usable if it has not been splitted in this set of clauses otherwise it is called unusable. A new clause is usable if it is the result of a case-analysis. To realize the case-analysis principle in Prolog, one needs to extend each predicate with a program continuation $`P`$ (a difference list) that handles those dynamic clauses. A disjunctive Prolog clause $`(h_1(I,G);\mathrm{};h_p(I,G):`$ $`b_1(I,G),\mathrm{},b_n(I,G))`$ is then compiled in a set of Prolog clauses as follows:
$`h_1(I,G,P):`$
$`usable\mathrm{\_}\mathrm{?}(h_1b_1\mathrm{}b_n,P),`$
$`b_1(I,G,P),\mathrm{},b_n(I,G,P)`$
;
$`usable\mathrm{\_}\mathrm{?}(h_1\mathrm{}h_pb_1\mathrm{}b_n,P),`$
$`unusable\mathrm{\_}!(h_1\mathrm{}h_pb_1\mathrm{}b_n,P),`$
$`\{`$ $`same\mathrm{\_}assumptions(P,P_l),`$
$`usable\mathrm{\_}!(h_l,P_l),false(I,G,P_l),\}_{1<lp}`$
$`usable\mathrm{\_}!(h_1b_1\mathrm{}b_n,P),`$
$`b_1(I,G,P),\mathrm{},b_n(I,G,P).`$
$`\{h_k(I,G,P):`$
$`usable\mathrm{\_}\mathrm{?}(h_k,P)`$
;
$`usable\mathrm{\_}\mathrm{?}(h_1\mathrm{}h_pb_1\mathrm{}b_n,P),`$
$`unusable\mathrm{\_}!(h_1\mathrm{}h_pb_1\mathrm{}b_n,P),`$
$`\{`$ $`same\mathrm{\_}assumptions(P,P_l),`$
$`usable\mathrm{\_}!(h_l,P_l),false(I,G,P_l),\}_{1lp,lk}`$
$`usable\mathrm{\_}!(h_kb_1\mathrm{}b_n,P).\}_{1<kp}`$
### Evaluation function of GADEL
The evaluation function is the heart of the GADEL system. It rates each chromosome given a default theory compiled in a disjunctive prolog program.
#### Evaluation of pair of genes.
For a default $`\delta _i=\frac{\alpha _i:\beta _i^1,\mathrm{},\beta _i^{k_i}}{\gamma _i}`$, an intermediate evaluation function $`f`$ is defined in Table 1. Given the two positions $`G|_{2i1}`$ and $`G|_{2i}`$ in the chromosome associated to the default $`\delta _i`$, the first point is to determine w.r.t. these values if this default is supposed to be involved in the construction of the candidate extension and then to check if this application is relevant.
We only illustrate $`f`$ on the case $`G|_{2i1}=1`$ and $`G|_{2i}=0`$ of the Table 1 (with penalty $`p_2`$ and a default with only one justification):
$`f(I,(G,Eval\mathrm{\_}G_{>I}),(G,Eval\mathrm{\_}G_I)):`$
$`G|_{2I1}=1,G|_{2I}=0,`$
$`false(I,G,\mathrm{\_}P),\%\%(CE(G)\alpha _i)=true`$
$`false((I,1),G,\mathrm{\_}P),\%\%(CE(G)\neg \beta _i)=true`$
$`Eval\mathrm{\_}G_IisEval\mathrm{\_}G_{>I}+p_2.`$
#### Evaluation of chromosome and population.
The evaluation of a chromosome is the total sum of the evaluations for each pair of genes. Our evaluation function is calculated directly over the population tree by a depth-first traversal. The result is a set of pairs of a chromosome and its evaluation. During the traversal, the construction of the current evaluated chromosome is prefixed and the evaluation of the genes is postfixed.
$`evaluation(Current\mathrm{\_}Pop,Evaluated\mathrm{\_}Pop):`$
$`eval\mathrm{\_}alpha(0,\mathrm{},Current\mathrm{\_}Pop,Evaluated\mathrm{\_}Pop).`$
$`eval\mathrm{\_}alpha(I,G,one(Subtree),Evaluations):`$
$`G|_{2I1}=1,`$
$`eval\mathrm{\_}beta(I,G,Subtree,Evaluations).`$
$`eval\mathrm{\_}beta(I,G,zero(Subtree),Evaluations^{}):`$
$`G|_{2I}=0,`$
$`I^{}isI+1,`$
$`eval\mathrm{\_}alpha(I^{},G,Subtree,Evaluations),`$
$`map(f,I,Evaluations,Evaluations^{}).`$
## Applying the System
### Methodology
Methodology for using GADEL is the same as using default logic as a framework for knowldege representation.
### Specifics
The semantics of our system is the Reiter’s propositional default logic.
### Users and Useability
The GADEL system takes default theory as Prolog facts in an input file. Classical formulas of default theory are arbitrary formulas with conjunctions, disjunctions and negations (noted resp. $`\&\&`$, $`||`$ and $`!`$). A default is a triplet composed of a prerequisite, a list of justifications and a conclusion. GADEL is a framework for non monotonic reasoning systems. To extend GADEL to an other system, one needs to redefine the evaluation function.
## Evaluating the System
### Benchmarks
We define two kinds of benchmarks: a taxonomic default theory “*people*” described in Table 2 and the well known Hamiltonian cycle problem in Table 5 as it has been described and encoded in (?).
### Comparison
DeRes and GADEL are compared on our two kinds of benchmarks. CPU times given are in seconds on a SUN E3000 $`(2\times 250Mhz)`$. The parameters of the genetic algorithm are $`p_c=0.8`$ and $`p_m=0.1`$.
#### GADEL:1/DeRes:0.
Table 3 gives results about the people default theory. Each line corresponds to the common part of $`W_{people}`$ augmented with one of the specified formula of the first column. The second column gives $`p_{size}`$ the initial number of chromosomes in the population, the third one is the average number of populations needed to find an extension. The last two columns give CPU times. (?) describes the very good performances of DeRes on some kind of default theories: the stratified ones. But it is also noticed that for a non stratified default theory the performance of DeRes are not enough to deal with a non very few number of defaults. Results given in this table shows that DeRes has a lot of difficulties with our taxonomic people example (even if the local prover is used). Conversely the number of populations are quite small for GADEL (even if the time is not so good: all the implementation is written in Prolog).
Table 4 gives results about $`W_{people}\{man\}`$ with different sizes of populations (200 tests for each size of population). The second line gives $`p_{size}`$ the initial number of chromosomes in the population. The third one gives the time spent for one complete computation of a new population. The fourth one gives the average number of populations needed to find an extension. The last one gives the average time to find an extension. These results demonstrate that the size of the population must be balanced by the time spent for one complete computation of a new population. The increase of the population size does not necessarily increase the efficiency of the genetic algorithm. Finally, Figure 2 presents for $`W_{people}\{man\}`$ with $`N=17`$ the number of tests w.r.t. the number of populations needed to obtain an extension. This figure suggests to stop computation after 6 populations and to restart with a brand new one since $`80\%`$ of tests end after 6 populations at most.
#### GADEL:1/Deres:1.
GADEL has poor performances on Hamiltonian problems. We think that it is because we do not take into account the groundedness (?) into our evaluation function. As a matter of fact, in the Hamiltonian problem, a solution is exactly one “chain”<sup>5</sup><sup>5</sup>5We say that $`\delta `$ is chained to $`\delta ^{}`$ if the prerequisite of $`\delta ^{}`$ is deducible from $`W`$ and the consequent of $`\delta `$. of defaults, but, there is a lot of potential solutions (whose evaluation is null) based on two, or more, chains of defaults. The only criterion to discard these candidate extensions is the groundedness property that they do not satisfy. Conversely, in people example, a solution is a set of non conflicting defaults, but at most four defaults are chained together, and so the groundedness property is less important to reach a solution. We are now testing some new evaluation functions in order to take into account this criterion.
#### Other systems.
We have also in mind that in the area of logic programming and non monotonic reasoning there exist others systems (Smodels (?), DLV (?)) able to compute stable models of extended logic program. Since this task is equivalent to compute an extension of a default theory it seems interesting to compare GADEL to these systems. But, even if DLV has the advantage to accept formulas with variables which are instantiated before computation, this system does not accept theories like our people example. On its part, Smodels does not deal with this default theory because it can not be represented by a normal logic program without disjunction. Because we have the objective to deal with every kind of propositional formulas, GADEL spends a lot of time in theorem proving and it seems not realistic to compare it with those two systems. But it will be very inter-resting to work on GADEL’s architecture in order to improve its performances on particular subclasses of default theories.
### Problem Size
The system is a prototype which can handle non stratified theories with about one hundred defaults.
## Conclusion
In this paper, we have described the first implementation of our system GADEL whose goal is to compute extensions of every kind of finite propositional Reiter’s default theories. Our new approach, using principles of genetic algorithms, seems to be relevant as it is illustrated by our experimental results. But this present work is a first approach and we have in mind many improvements as : more accurate definition of the evaluation function, using reparation techniques, local search heuristics.
|
warning/0003/math0003074.html
|
ar5iv
|
text
|
# Relating the Connes-Kreimer and Grossman-Larson Hopf algebras built on rooted trees
## 1 Introduction
In , Dirk Kreimer discovered the striking fact that the process of renormalization in quantum field theory may be described, in a conceptual manner, by means of certain Hopf algebras (which depend on the chosen renormalization scheme). A toy model was studied in detail by Alain Connes and Dirk Kreimer in ; the Hopf algebra which occurs, denoted by $`H_R`$, is the polynomial algebra in an infinity of indeterminates, one for each rooted tree, but with a non-cocommutative comultiplication. Some operators, denoted by $`N`$ and $`L`$, have been defined on $`H_R`$. The first one is the natural growth operator, which acts as a derivation; it defines some elements $`\delta _k`$, for $`k1`$, which provide the link between $`H_R`$ and another Hopf algebra, introduced in by Alain Connes and Henri Moscovici in a completely different context, namely in noncommutative geometry. The operator $`L`$ is a solution of the “Hochschild equation”, and the pair $`(H_R,L)`$ is characterized as the solution of a universal problem in Hochschild cohomology. It was proved also in that $`H_R`$ is in duality with the universal enveloping algebra of a certain Lie algebra $`^\mathcal{1}`$, which has a linear basis indexed by all (non-empty) rooted trees. Let us note that this Lie algebra $`^\mathcal{1}`$ appeared also, very recently, in , in the context of pre-Lie algebras and the operad of rooted trees.
In this note we would like to draw the attention to another Hopf algebra built on rooted trees, introduced ten years ago by Robert Grossman and Richard Larson in (see also their survey ). This Hopf algebra (denoted by $`A`$ in what follows) has a $`linear`$ basis consisting of all (non-empty) rooted trees, a noncommutative product, and is a cocommutative graded connected Hopf algebra, hence, by the Milnor-Moore theorem, it is the universal enveloping algebra of the Lie algebra of its primitive elements, which may also be described explicitly: $`P(A)`$ has a linear basis consisting of all rooted trees whose root has exactly one child. Using these properties of $`A`$, Grossman and Larson gave a Hopf algebraic proof of the classical result of Cayley on the number of rooted trees. The construction of the Hopf algebra $`A`$ is motivated also by some ideas concerning differential operators and differential equations, in particular the Runge-Kutta method in numerical analysis (see ). Let us note also that actually the construction of Grossman and Larson is slightly more general: they associate a Hopf algebra to any family of trees which satisfy a certain list of axioms. Among these families are: the family of all rooted trees (this is the one who gives the Hopf algebra $`A`$) and the families of all ordered, heap-ordered and respectively labelled rooted trees. The construction of the Hopf algebra associated to such a family is similar to that of $`A`$.
As noted in , , the Hopf algebra $`H_R`$ may also be related to the Runge-Kutta method and the Butcher group, so it is very likely that there is a relation between $`H_R`$ and $`A`$. As we shall see, this relation is best expressed by the fact that the Hopf algebras $`A`$ and $`U(^\mathcal{1})`$ are isomorphic, which in turn is a consequence of a Lie algebra isomorphism between $`P(A)`$ and $`^\mathcal{1}`$. This isomorphism is given by sending a rooted tree $`tP(A)`$ to the rooted tree in $`^\mathcal{1}`$ obtained by deleting the root of $`t`$.
We believe that this relation between $`A`$ and $`H_R`$ may be useful for a better understading of both these Hopf algebras. On one hand, the advantage of the isomorphism between $`A`$ and $`U(^\mathcal{1})`$, which allows one to work with $`A`$ instead of $`U(^\mathcal{1})`$, is clear, since we know on $`A`$ a very explicit $`linear`$ basis, more manageable than the PBW basis of $`U(^\mathcal{1})`$. On the other hand, some known results for one of the Hopf algebras $`A`$ and $`H_R`$ may serve as a motivation and inspiration for obtaining similar results for the other. We shall make here a first step in this direction, by studying two natural operators on $`A`$. The first is the natural growth operator $`N`$ (defined exactly as the one introduced by Connes and Kreimer for $`H_R`$), which will turn out to be a $`coderivation`$ on $`A`$; the sequence $`\{x_k\}_{k0}`$ defined by $`N`$ will turn out to generate a commutative cocommutative Hopf subalgebra of $`A`$, isomorphic to the polynomial algebra in one indeterminate, with its usual Hopf algebra structure. A nice feature of $`N`$ (considered on $`A`$) is that, for any rooted tree $`t`$, $`N(t)`$ may be described as the product (in $`A`$) between the rooted tree with two vertices and $`t`$. The second one, denoted by $`M`$, is in some sense dual to the operator $`L`$ on $`H_R`$: we shall prove that $`M`$ is a derivation (the right $`A`$-module structure on $`A`$ being the one induced by $`\epsilon `$) and that the transpose of $`M`$ gives a solution of the Hochschild equation on the finite dual Hopf algebra $`A^0`$.
## 2 The relation between $`H_R`$ and $`A`$
Throughout, $`k`$ will be a fixed field of characteristic zero and all algebras, linear spaces etc. will be over $`k`$; unadorned $``$ means $`_k`$.
We start by recalling some facts from , , , to which we refer for the terminology and more details (the reader will find also in some nice pictures of rooted trees and the operations which may be performed with them).
A $`rooted`$ $`tree`$ $`t`$ is a connected and simply-connected set of oriented edges and vertices such that there is exactly one distinguished vertex with no incoming edges, called the $`root`$ of $`t`$. Every edge connects two vertices. The $`fertility`$ $`f(v)`$ of a vertex $`v`$ is the number of edges outgoing from $`v`$, that is the number of $`children`$ of $`v`$. A $`forest`$ is a finite set of rooted trees.
We denote by $`B_{}`$ the operator which assigns to a rooted tree $`t`$ a forest, by removing the root of $`t`$, and by $`B_+`$ the operator which maps a forest consisting of $`n`$ rooted trees $`t_1,\mathrm{},t_n`$ to a new rooted tree $`t`$ which has a root $`r`$ with fertility $`f(r)=n`$ which connects to the $`n`$ roots of $`t_1,\mathrm{},t_n`$. Obviously, $`B_+(B_{}(t))=B_{}(B_+(t))=t`$ for any rooted tree $`t`$. We also set $`B_{}(e)=\mathrm{}`$, $`B_+(\mathrm{})=e`$, where $`e`$ is the rooted tree with only one vertex and $`\mathrm{}`$ is the empty tree.
If $`t`$ is a rooted tree, an $`elementary`$ $`cut`$ is a cut of $`t`$ at a single chosen edge, and an $`admissible`$ $`cut`$ is a set of elementary cuts such that any path from any vertex of $`t`$ to the root of $`t`$ contains at most one elementary cut. If $`c`$ is an admissible cut, we denote by $`|c|`$ the number of elementary cuts of $`c`$. If we perform an admissible cut $`c`$ in a rooted tree $`t`$, we obtain a forest, denoted by $`P^c(t)`$, consisting of the $`cut`$ $`branches`$ of $`t`$, and a $`trunk`$, denoted by $`R^c(t)`$, which is the branch which remains (it is the one which contains the root of $`t`$).
We can now define the Connes-Kreimer Hopf algebra over $`k`$, denoted by $`H_R`$. As an algebra, it is the polynomial algebra in an infinity of indeterminates, one for each (non-empty) rooted tree (we denote also by $`t`$ the indeterminate corresponding to the rooted tree $`t`$). The unit is denoted by 1 (it corresponds to the empty tree). The comultiplication $`\mathrm{\Delta }`$ is defined by:
$$\mathrm{\Delta }(1)=11$$
$$\mathrm{\Delta }(t)=1t+t1+\underset{c}{}P^c(t)R^c(t)$$
for any rooted tree $`t`$, where the sum is over all admissible cuts $`c`$ of $`t`$ (with $`|c|1`$) and $`P^c(t)`$ is the monomial corresponding to the forest $`P^c(t)`$ (as a general rule, we identify any forest with its monomial). An alternative recursive description of $`\mathrm{\Delta }(t)`$ is
$$\mathrm{\Delta }(t)=t1+(idB_+)(\mathrm{\Delta }(B_{}(t))$$
The counit is given by
$$\epsilon (1)=1$$
$$\epsilon (t)=0$$
for any rooted tree $`t`$. The antipode is given iteratively by
$$S(1)=1$$
$$S(t)=t\underset{c}{}S(P^c(t))R^c(t)$$
for any rooted tree $`t`$, where the sum is over all admissible cuts $`c`$ of $`t`$ (with $`|c|1`$).
Define now the operator $`N`$ on $`H_R`$ (the natural growth operator) which maps a rooted tree $`t`$ with $`n`$ vertices to a sum $`N(t)`$ of $`n`$ rooted trees $`t_i`$, each having $`n+1`$ vertices, by attaching one more outgoing edge and vertex to each vertex of $`t`$ (the root remains the same under this operation). On products of rooted trees, $`N`$ acts, by definition, as a derivation. For any $`k1`$, define elements $`\delta _kH_R`$ by $`\delta _1=e`$, $`\delta _{k+1}=N(\delta _k)`$, that is $`\delta _{k+1}=N^k(e)`$ for all $`k1`$. We recall that we have denoted by $`e`$ the tree with one vertex (in order to be consistent with ; in by $`e`$ is denoted the unit of $`H_R`$). These elements are very important in , because they provide the link between $`H_R`$ and the Connes-Moscovici Hopf algebra introduced in . They generate a (non-cocommutative) Hopf subalgebra of $`H_R`$.
Define now the operator $`L:H_RH_R`$ as the unique linear map satisfying the condition $`L(t_1\mathrm{}t_m)=t`$ for any rooted trees $`t_1,\mathrm{},t_m`$, where $`t`$ is the rooted tree obtained by connecting a new root to the roots of $`t_1,\mathrm{},t_m`$. Obviously it agrees with the map $`B_+`$ introduced above. It was shown in that the operator $`L`$ satisfies the so-called “Hochschild equation”
$$\mathrm{\Delta }L=L1+(idL)\mathrm{\Delta }$$
(this is a 1-cocycle condition) and the pair $`(H_R,L)`$ has the following universal property: if $`(H_1,L_1)`$ is a pair with $`H_1`$ a commutative Hopf algebra and $`L_1:H_1H_1`$ a linear map satisfying the Hochschild equation on $`H_1`$, then there exists a unique Hopf algebra map $`\rho :H_RH_1`$ such that $`L_1\rho =\rho L`$.
Let $`^\mathcal{1}`$ be the linear span of the elements $`Z_t`$, indexed by all (non-empty) rooted trees. Define an operation on $`^\mathcal{1}`$ by
$$Z_{t_1}Z_{t_2}=\underset{t}{}n(t_1,t_2;t)Z_t$$
where the integer $`n(t_1,t_2;t)`$ is the number of admissible cuts $`c`$ of $`t`$ with $`|c|=1`$ such that the cut branch is $`t_1`$ and the trunk is $`t_2`$ (note that this operation is $`not`$ associative). Define then a bracket on $`^\mathcal{1}`$ by
$$[Z_{t_1},Z_{t_2}]=Z_{t_1}Z_{t_2}Z_{t_2}Z_{t_1}$$
Then it was proved in that $`(^\mathcal{1},[,])`$ is a Lie algebra and moreover there is a Hopf duality between $`H_R`$ and the universal enveloping algebra of $`^\mathcal{1}`$.
We recall now from , the structure of the Grossman-Larson Hopf algebra, which will be denoted in what follows by $`A`$. It has a linear basis consisting of all (non-empty) rooted trees. The unit is the tree $`e`$ with only one vertex. The multiplication on the basis is given as follows: let $`t_1`$ and $`t_2`$ be two rooted trees, let $`s_1,\mathrm{},s_r`$ be the children of the root of $`t_1`$, let $`n`$ be the number of vertices of $`t_2`$; then there are $`n^r`$ ways to attach the $`r`$ subtrees of $`t_1`$ which have $`s_1,\mathrm{},s_r`$ as roots to the tree $`t_2`$ by making each $`s_i`$ the child of some vertex of $`t_2`$. The product $`t_1t_2`$ is defined to be the sum of these $`n^r`$ rooted trees (note that this product is $`not`$ commutative).
The coalgebra structure of $`A`$ is given as follows. If $`t`$ is a rooted tree whose root has children $`s_1,\mathrm{},s_r`$, the coproduct $`\mathrm{\Delta }(t)`$ is the sum of the $`2^r`$ terms $`t_1t_2`$, where the children of the root of $`t_1`$ and the children of the root of $`t_2`$ range over all $`2^r`$ possible partitions of the children of the root of $`t`$ into two subsets. If $`t=e`$, then $`\mathrm{\Delta }(e)=ee`$. The counit $`\epsilon `$ is given by $`\epsilon (e)=1`$, $`\epsilon (t)=0`$ if $`te`$. Obviously $`\mathrm{\Delta }`$ is cocommutative. Moreover, $`A`$ is a graded connected bialgebra, the component $`A_n`$ of degree $`n`$ having as basis all trees with $`n+1`$ vertices. Being a graded connected cocommutative bialgebra, A is a Hopf algebra and by the Milnor-Moore theorem $`A`$ is the universal enveloping algebra of its primitives, $`AU(P(A))`$, where $`P(A)=\{aA/\mathrm{\Delta }(a)=ea+ae\}`$. There is an explicit description of $`P(A)`$: it has a basis consisting of all rooted trees whose root has exactly one child.
We can state now the result which expresses the relation between the Hopf algebras $`H_R`$ and $`A`$.
###### Proposition 2.1
The Lie algebras $`^\mathcal{1}`$ and $`P(A)`$ are isomorphic, hence $`A`$ is isomorphic to $`U(^\mathcal{1})`$ as Hopf algebras.
Proof: Define $`\phi :P(A)^\mathcal{1}`$ as the unique linear map which on the basis of $`P(A)`$ acts as follows: if $`tP(A)`$ is a rooted tree, then $`\phi (t)=Z_{B_{}(t)}`$. Recall that $`B_{}(t)`$ is the rooted tree obtained by deleting the root of $`t`$ (here it is a tree, since the root of $`t`$ has exactly one child). Obviously, $`\phi `$ is a linear isomorphism, its inverse being the map $`\psi :^\mathcal{1}𝒫(𝒜)`$, $`\psi (Z_t)=B_+(t)`$ for any rooted tree $`t`$.
It remains to prove that $`\phi `$ is a Lie algebra map. Let $`t_1,t_2P(A)`$ be two rooted trees. In $`P(A)`$, we have $`[t_1,t_2]=t_1t_2t_2t_1`$. By the definition of the multiplication of $`A`$, we obtain that
$$t_1t_2=B(t_1t_2)+\underset{i}{}t^i$$
$$t_2t_1=B(t_2t_1)+\underset{j}{}T^j$$
where $`B(t_1t_2)`$ is the rooted tree obtained by identifying the roots of $`t_1`$ and $`t_2`$, and each $`t^i`$ is obtained by identifying the root of $`t_1`$ with a vertex of $`t_2`$, except for the root of $`t_2`$ (and similarly for $`t_2t_1`$), so all the rooted trees $`t^i`$ and $`T^j`$ are in $`P(A)`$. Obviously $`B(t_1t_2)=B(t_2t_1)`$, hence
$$[t_1,t_2]=\underset{i}{}t^i\underset{j}{}T^j$$
We obtain that $`\phi ([t_1,t_2])=_i\phi (t^i)_j\phi (T^j)`$. From the definition of the operation $``$ on $`^\mathcal{1}`$, it is easy to see that
$$\underset{i}{}\phi (t^i)=\phi (t_1)\phi (t_2)$$
$$\underset{j}{}\phi (T^j)=\phi (t_2)\phi (t_1)$$
hence we obtain
$$\phi ([t_1,t_2])=[\phi (t_1),\phi (t_2)]$$
that is $`\phi `$ is a Lie algebra map.
Define now the natural growth operator $`N`$ on $`A`$, by the same formula as the one defined by Connes and Kreimer on $`H_R`$, that is $`N`$ is the linear map $`N:AA`$ such that, for any rooted tree $`t`$, $`N(t)`$ is the sum of the rooted trees obtained from $`t`$ by attaching one more outgoing edge and vertex to each vertex of $`t`$. The properties of $`N`$ are collected in the following
###### Proposition 2.2
$`(1)N(e)b=N(b)`$ for all $`bA`$.
$`(2)N(ab)=N(a)b`$ for all $`a,bA`$.
$`(3)N^k(e)b=N^k(b)`$ for all $`k1,bA`$.
$`(4)N^k(e)N(e)=N^{k+1}(e)=N(e)N^k(e)`$ for all $`k1`$.
$`(5)N`$ is a coderivation, that is, for all $`bA`$, we have
$$\mathrm{\Delta }(N(b))=(idN+Nid)(\mathrm{\Delta }(b))$$
Proof: $`(1)`$ obviously, for any rooted tree $`t`$, we have $`N(e)t=N(t)`$, from which we obtain $`N(e)b=N(b)`$ for all $`bA`$.
$`(2)`$ using $`(1)`$, we have $`N(ab)=N(e)ab=N(a)b`$.
$`(3)`$ follows easily by induction, and $`(4)`$ follows from $`(3)`$.
$`(5)`$ by $`(1)`$, we obtain that $`\mathrm{\Delta }(N(b))=\mathrm{\Delta }(N(e)b)=\mathrm{\Delta }(N(e))\mathrm{\Delta }(b)`$. We shall use the $`\mathrm{\Sigma }`$-notation, so we write $`\mathrm{\Delta }(b)=b_{(1)}b_{(2)}`$; since $`N(e)`$ is a primitive element, we have then:
$$\mathrm{\Delta }(N(b))=(eN(e)+N(e)e)(b_{(1)}b_{(2)})=$$
$$b_{(1)}N(e)b_{(2)}+N(e)b_{(1)}b_{(2)}=$$
$$b_{(1)}N(b_{(2)})+N(b_{(1)})b_{(2)}=$$
$$(idN+Nid)(\mathrm{\Delta }(b))$$
Now, for any $`k0`$, define the element $`x_k=N^k(e)A`$. These elements are analogous to the elements $`\delta _k`$ of .
###### Proposition 2.3
The elements $`x_k`$ have the following properties:
$`(1)x_mx_n=x_nx_m=x_{m+n}`$
$`(2)x_0=e`$
$`(3)\epsilon (x_m)=\delta _{0,m}`$
$`(4)\mathrm{\Delta }(x_m)=_{i=0}^m\left(\begin{array}{c}m\\ i\end{array}\right)x_ix_{mi}`$
$`(5)S(x_m)=(1)^mx_m`$
for all $`m,n0`$, where $`S`$ is the antipode of $`A`$. Moreover, the elements $`\{x_k\}_{k0}`$ are linearly independent. Hence, the subspace of $`A`$ generated by the elements $`x_k`$ with $`k0`$ is a commutative cocommutative Hopf subalgebra of $`A`$, isomorphic (as a Hopf algebra) to the polynomial algebra $`k[X]`$ with its usual Hopf algebra structure.
Proof: $`(1)`$ follows easily from the previous proposition; $`(2)`$ and $`(3)`$ are obvious; $`(4)`$ and $`(5)`$ follow by induction, using the facts that $`x_{m+1}=N(e)x_m`$ and $`N(e)`$ is primitive. Since $`x_kA_k`$ for all $`k0`$ (hence $`deg(x_m)deg(x_n)`$ if $`mn`$) it follows that the elements $`\{x_k\}_{k0}`$ are linearly independent. The isomorphism between the Hopf subalgebra given by the elements $`x_k`$ and $`k[X]`$ is determined by $`x_kX^k`$ for all $`k0`$.
Define now the $`k`$-linear map $`M:AA`$, by $`M(e)=0`$ and $`M(t)=tN(e)`$ for all rooted trees $`te`$. As we shall see, this operator $`M`$ is in some sense dual to the operator $`L`$ on $`H_R`$.
###### Proposition 2.4
$`(1)M(b)=(b\epsilon (b))N(e)`$ for all $`bA`$.
$`(2)M(ab)=aM(b)+M(a)\epsilon (b)`$ for all $`a,bA`$
(that is, $`M`$ is a derivation from $`A`$ into the bimodule $`A`$, where the left $`A`$-module structure of $`A`$ is the one given by multiplication and the right $`A`$-module structure is given via $`\epsilon `$).
$`(3)\mathrm{\Delta }(M(b))=(idM+Mid)(\mathrm{\Delta }(b))+(b\epsilon (b))N(e)+N(e)(b\epsilon (b))`$
for all $`bA`$.
Proof: $`(1)`$ Let $`bA`$, written as $`b=b_0e+_{t_ie}b_it_i`$, with $`b_0,b_ik`$. Then we have:
$$M(b)=b_0M(e)+b_iM(t_i)=b_it_iN(e)=$$
$$b_it_iN(e)+b_0N(e)b_0N(e)=bN(e)b_0N(e)=(b\epsilon (b))N(e)$$
$`(2)`$ using $`(1)`$, we have $`M(ab)=(ab\epsilon (a)\epsilon (b))N(e)`$, and
$$aM(b)+M(a)\epsilon (b)=a(b\epsilon (b))N(e)+(a\epsilon (a))\epsilon (b)N(e)=$$
$$(ab\epsilon (b)a+\epsilon (b)a\epsilon (a)\epsilon (b))N(e)=(ab\epsilon (a)\epsilon (b))N(e)$$
$`(3)`$ write $`\mathrm{\Delta }(b)=b_{(1)}b_{(2)}`$; then we have:
$$\mathrm{\Delta }(M(b))=\mathrm{\Delta }((b\epsilon (b))N(e))=\mathrm{\Delta }(b)\mathrm{\Delta }(N(e))\epsilon (b)\mathrm{\Delta }(N(e))=$$
$$b_{(1)}N(e)b_{(2)}+b_{(1)}b_{(2)}N(e)\epsilon (b)N(e)N(e)\epsilon (b)=$$
$$M(b_{(1)})b_{(2)}+\epsilon (b_{(1)})N(e)b_{(2)}+b_{(1)}M(b_{(2)})+$$
$$+b_{(1)}\epsilon (b_{(2)})N(e)\epsilon (b)N(e)N(e)\epsilon (b)=$$
$$(idM+Mid)(\mathrm{\Delta }(b))+(b\epsilon (b))N(e)+N(e)(b\epsilon (b))$$
Recall from , that we can associate to $`A`$ the so-called $`finitedual`$, which is also a Hopf algebra, and which consists of the elements $`fA^{}`$ such that $`Ker(f)`$ contains a cofinite ideal of $`A`$. If we denote by $`m`$ the multiplication of $`A`$, then an element $`fA^{}`$ belongs to $`A^0`$ if and only if $`m^{}(f)A^{}A^{}`$, which in turn is equivalent to the fact that there exist some elements $`f_i,f_i^{}A^{}`$, with $`i`$ in some finite set, such that $`f(ab)=f_i(a)f_i^{}(b)`$ for all $`a,bA`$. Moreover, we have that $`m^{}(A^0)A^0A^0`$ and the comultiplication of $`A^0`$ is $`\mathrm{\Delta }=m^{}|A^0`$.
Now, let $`fA^0`$, $`m^{}(f)=f_if_i^{}`$, with $`f_i,f_i^{}A^0`$. By using the condition $`M(ab)=aM(b)+M(a)\epsilon (b)`$ satisfied by $`M`$, we can compute:
$$f(M(ab))=f(M(a))\epsilon (b)+f(aM(b))$$
for all $`a,bA`$, which may be rewritten as
$$f(M(ab))=f(M(a))\epsilon (b)+f_i(a)f_i^{}(M(b))$$
that is
$$M^{}(f)(ab)=M^{}(f)(a)\epsilon (b)+f_i(a)M^{}(f_i^{})(b)$$
hence $`M^{}(f)A^0`$. So, we have $`M^{}(A^0)A^0`$, and we denote by $`M^0`$ the restriction of $`M^{}`$ to $`A^0`$. Also, for $`fA^0`$, since $`M^{}(f),\epsilon ,f_i,M^{}(f_i^{})A^0`$ for all $`i`$, we obtain finally that in $`A^0`$ the following relation holds:
$$\mathrm{\Delta }(M^0(f))=M^0(f)\epsilon +(idM^0)(\mathrm{\Delta }(f))$$
which expresses the fact that $`M^0`$ is a solution for the Hochschild equation on $`A^0`$. Hence, since $`A^0`$ is commutative (because $`A`$ is cocommutative), by the universal property of the pair $`(H_R,L)`$ we obtain the following
###### Proposition 2.5
There exists a unique Hopf algebra map $`\rho :H_RA^0`$ such that $`M^0\rho =\rho L`$.
|
warning/0003/nucl-th0003074.html
|
ar5iv
|
text
|
# References
Dynamical model for Pion $``$ Nucleon Bremsstrahlung
A. Mariano<sup>a,b</sup> and G. López Castro<sup>a</sup>
<sup>a</sup> Departamento de Física, Centro de Investigación y de Estudios Avanzados del IPN
A.P. 14$``$740 México 07000 D.F.
<sup>b</sup> Departamento de Física, Facultad de Ciencias Exáctas, Universidad Nacional de La Plata
cc.67, 1900 La Plata, Argentina
Abstract
A dynamical model based on effective Lagrangians is proposed to describe the bremsstrahlung reaction $`\pi N\pi N\gamma `$ at low energies. The $`\mathrm{\Delta }(1232)`$ degrees of freedom are incorporated in a way consistent with both, electromagnetic gauge invariance and invariance under contact transformations. The model also includes the initial and final state rescattering of hadrons via a T-matrix with off-shell effects. The $`\pi N\gamma `$ differential cross sections are calculated using three different T-matrix models and the results are compared with the soft photon approximation, and with experimental data. The aim of this analysis is to test the off-shell behavior of the different T-matrices under consideration.
PACS numbers: 25.80.Ek, 13.60.-n, 13.75.-r
I. INTRODUCTION
In order to extract resonant parameters of the nucleon resonances (N) from the $`\gamma N\pi N`$ reaction, it is important to evaluate the background contribution to isolate the resonant peak. An important contribution to the background of this photo-production reaction is provided by the final state rescattering (FSI) of the $`\pi N`$ system . Consequently, we require knowledge of the T-matrix in the off-momentum-shell regime(off-shell) to describe this rescattering process. That is, we need information about the amplitude $`T(\stackrel{}{q}{}_{}{}^{},\stackrel{}{q};z(\stackrel{}{q}))`$ with $`|\stackrel{}{q}||\stackrel{}{q}{}_{}{}^{}|`$, where $`z`$ which is the total energy of the $`\pi N`$ system is a function of the relative momentum $`|\stackrel{}{q}|`$ of the initial state. This particular rescattering amplitude is more properly called the half-off-shell T-matrix. As is known, the T-matrix can be generated by solving an integral equation of the Lippman-Schwinger or Bethe-Salpeter type by iteration of a pure phenomenological potential, or from an effective potential based on meson-exchange models . In all cases the so called ‘realistic’ interactions are fitted to reproduce the phase shifts in elastic $`\pi N`$ scattering, which only depends on the on-shell ($`|\stackrel{}{q}|=|\stackrel{}{q}{}_{}{}^{}|`$) values of the relative momenta. Thus elastic scattering is not useful to constrain the off-shell behavior of the T-matrix, i.e., interactions which yield similar results for elastic scattering may have different behavior in the off-shell regime.
Another reaction where the $`\pi N`$ off-shell T-matrix is required is $`\pi N\pi N\gamma `$ bremsstrahlung. This process has been studied within the Soft-Photon Approximation (SPA). The soft-photon amplitude, defined by the first two terms of the soft-photon expansion, depends only on the electromagnetic constants of $`\pi `$ and $`N`$ and on the corresponding $`\pi N`$ elastic scattering amplitude (on-shell $`\pi N`$ T-matrix) . It should be pointed out that Low’s original amplitude fails to describe the $`\pi ^\pm N\gamma `$ process near the resonance region . Nevertheless the SPA for the $`\pi N\gamma `$ reaction, as implemented by Ref., describes well the experimental data near the $`\mathrm{\Delta }`$ resonance region and can provide a determination of the $`\mathrm{\Delta }^{++}`$ magnetic moment .
Because the soft-photon amplitude depends only on the on-shell T-matrix, the information on the off-shell behavior of the T-matrix can be tested by adding the contributions to the radiative $`\pi N`$ scattering within the framework of a specific dynamical model. The purpose of the present paper is to check the off-shell behavior of three different T-matrices for $`\pi N`$ rescattering in the reaction $`\pi N\pi N\gamma `$. A similar analysis of these FSI effects using different rescattering off-shell amplitudes has been done with the aim to determine the form factors of nucleons . It was found that pion photo-production reactions are very sensitive to the off-shell behavior of the $`\pi N`$ interaction, and also that there are certain inconsistencies in fixing phenomenological form factors, to match photo-production current and FSI effects.
For this purpose, we use a dynamical model to describe the $`\pi N\pi N\gamma `$ reaction. The gauge invariant electromagnetic current is constructed explicitly, with vertexes and propagators derived from the relevant hadronic and electromagnetic Lagrangians. We also include two-body meson exchange currents, and the full energy-momentum dependence of the $`\pi N`$ T-matrix which exhibits its off-shell behavior. Finally we implement this model with different T-matrices in order to compare their different off-shell dependencies.
This paper is organized as follows. In section II we will construct the gauge-invariant amplitude for radiative $`\pi N`$ scattering. In section III we give a summary of the corresponding results obtained in the SPA approximation in order to make comparisons with our dynamical model. The Lagrangians and propagators used to construct the gauge invariant current for our process are provided in section IV. Finally the results and conclusions are given in section V.
II. GAUGE INVARIANT BREMSSTRAHLUNG AMPLITUDE
In the pion-nucleon bremsstrahlung process we deal with a problem of the scattering by two potentials : the strong pion-nucleon and the electromagnetic interactions. The cross section for $`\pi N\pi N\gamma `$ process reads
$`d\sigma `$ $`=`$ $`{\displaystyle \frac{d\stackrel{}{k}}{\omega _\gamma }\frac{d\stackrel{}{q}_f}{\omega (\stackrel{}{q}_f)}\frac{d\stackrel{}{p}_f}{E(\stackrel{}{p}_f)}(2\pi )^4\delta ^4(p_i+q_ip_fq_fk)}`$ (1)
$`\times {\displaystyle \frac{1}{2}}{\displaystyle \underset{ϵ_\lambda ,ms_f,ms_i}{}}\left|{\displaystyle \frac{m_N^2}{2\sqrt{2}}}M_{\pi N\gamma ,\pi N}(ϵ_\lambda ,k;q_f,p_f,ms_f;q_i,p_i,ms_i)\right|^2,`$
where $`q=(\omega ,\stackrel{}{q})`$, $`p=(E,\stackrel{}{p})`$ and $`k=(\omega _\gamma ,\stackrel{}{k})`$ denote pion, nucleon and photon four-momenta, respectively; $`ms`$ is the nucleon’s spin projection and $`ϵ_\lambda `$ indicates the polarization four-vector of the photon. The subindexes $`i,f`$ refer to initial and final state quantities.
The Lorentz invariant amplitude<sup>1</sup><sup>1</sup>1Throughout this paper, $`M`$ will denote the amplitude generated by the operator $`\widehat{M}`$,ie., $`M=\overline{u}|\widehat{M}|u`$. $`M_{\pi N\gamma ,\pi N}`$ explicitly reads
$`M_{\pi N\gamma ,\pi N}=\overline{u}(\stackrel{}{p}_f,ms_f)|\widehat{M}_{\pi N\gamma ,\pi N}(ϵ_\lambda ,k;q_f,p_f;q_i,p_i)|u(\stackrel{}{p}_i,ms_i),`$ (2)
where $`u(\stackrel{}{p},ms)`$ denote nucleon Dirac spinors, and the amplitude operator $`\widehat{M}_{\pi N\gamma ,\pi N}`$ is obtained from the coupled channel Bethe-Salpeter equation for the $`\pi N\gamma `$ system as follows (we consider electromagnetic interactions at the lowest order):
$`\widehat{M}_{\pi N\gamma ,\pi N}`$ $`=`$ $`\widehat{V}_{\pi N\gamma ,\pi N}`$ (3)
$`+i{\displaystyle \frac{dq^4}{(2\pi )^4}\left[\widehat{V}_{\pi N\gamma ,\pi N}(q)\widehat{G}_{\pi N}(q)\widehat{M}_{\pi N,\pi N}(q)+\widehat{M}_{\pi N,\pi N}(q)\widehat{G}_{\pi N}(q)\widehat{V}_{\pi N\gamma ,\pi N}(q)\right]}`$
$`+i^2{\displaystyle \frac{dq^4}{(2\pi )^4}\frac{dq^4}{(2\pi )^4}\left[\widehat{M}_{\pi N,\pi N}(q^{})\widehat{G}_{\pi N}(q^{})\widehat{V}_{\pi N\gamma ,\pi N}(q^{},q)\widehat{G}_{\pi N}(q)\widehat{M}_{\pi N,\pi N}(q)\right]}.`$
The symbol $`𝑑q^4`$ will indicate integration over intermediate four-momenta variables.
In terms of the above operator amplitude, the T-matrix, defined as
$$\widehat{T}(q_f,p_f;q_i,p_i)=\frac{1}{(2\pi )^3}\widehat{M}_{\pi N,\pi N}(q_f,p_f;q_i,p_i),$$
(4)
satisfies the integral equation
$`\widehat{T}`$ $`=`$ $`\widehat{U}+i{\displaystyle \frac{dq^4}{(2\pi )^4}\widehat{U}(q)\widehat{G}(q)\widehat{T}(q)},`$ (5)
$`\widehat{U}`$ $`=`$ $`{\displaystyle \frac{1}{(2\pi )^3}}\widehat{V}_{\pi N,\pi N},`$
$`\widehat{G}`$ $`=`$ $`(2\pi )^3\widehat{G}_{\pi N}.`$
In the previous equations $`\widehat{V}_{ij}`$ denote $`\widehat{M}`$-matrix elements corresponding to the irreducible Feynman diagrams for each process, while $`\widehat{G}_i`$ is the product of Feynman propagators of intermediate particles.
Following the Thompson’s prescription , we can represent the above integrals in a three-dimensional form as follows (we set in the center of mass frame of the $`\pi N`$ system)
$$\widehat{T}(\stackrel{}{q}{}_{}{}^{},\stackrel{}{q},z)=\widehat{U}(\stackrel{}{q}{}_{}{}^{},\stackrel{}{q})+d^3\stackrel{}{q}^{\prime \prime }\widehat{U}(\stackrel{}{q}{}_{}{}^{},\stackrel{}{q}^{\prime \prime })\widehat{G}_{TH}(z,\stackrel{}{q}^{\prime \prime })\widehat{T}(\stackrel{}{q}^{\prime \prime },\stackrel{}{q},z),$$
(6)
with,
$$\widehat{G}_{TH}(z,\stackrel{}{q}^{\prime \prime })=\frac{m_N}{2\omega (\stackrel{}{q}^{\prime \prime })E(\stackrel{}{q}^{\prime \prime })}\frac{\underset{ms^{\prime \prime }}{}|u(\stackrel{}{q}^{\prime \prime },ms^{\prime \prime })\overline{u}(\stackrel{}{q}^{\prime \prime },ms^{\prime \prime })|}{zz^{\prime \prime }+i\eta },$$
(7)
where $`z^{\prime \prime }=E(\stackrel{}{q}^{\prime \prime })+\omega (\stackrel{}{q}^{\prime \prime })`$. In the above expressions $`\widehat{G}_{TH}`$ denotes the Thompson propagator replacing the full $`\widehat{G}_{\pi N}`$ Feynman propagator which, as a consequence of the three-dimensional reduction, eliminates the propagation of antiparticles and puts intermediate particles on their mass-shell. The kernel function $`\widehat{U}(\stackrel{}{q}{}_{}{}^{},\stackrel{}{q})`$ contains all the $`\pi N`$-interaction irreducible diagrams to be iterated in the T-matrix calculation, but usually only second-order contributions are kept.
The electromagnetic current $`\widehat{V}_{\pi N\gamma ,\pi N}`$ can be broken into two pieces,
$$\widehat{V}_{\pi N\gamma ,\pi N}\widehat{V}_{\pi N\gamma ,\pi N}^{(1)}+\widehat{V}_{\pi N\gamma ,\pi N}^{(2)},$$
(8)
where the upper indices denote one- and two-body contributions, respectively, which are obtained by coupling the photon to all the internal lines of $`\widehat{U}`$. As is known, the operator $`\widehat{V}_{\pi N\gamma ,\pi N}^{(2)}`$ must be added to the electromagnetic current in order to satisfy the electromagnetic gauge invariance of the total amplitude , while the one-body amplitude $`V_{\pi N\gamma ,\pi N}^{(1)}`$ vanishes for free hadrons. Both contributions to the total amplitude are illustrated in Fig. 1.
Let us now discuss some problems related to gauge invariance. The bremsstrahlung amplitude can be directly computed from Eqs.(2-4), by making a reduction that replaces $`\widehat{G}_{\pi N}`$ by $`\widehat{G}_{TH}`$ in Eq.(3) . This procedure, however, introduces some inconveniences. First, $`V_{\pi N\gamma ,\pi N}`$ is not gauge invariant by itself because it involves the addition of zero and second order terms in the hadronic vertexes, and thus $`M_{\pi N\gamma ,\pi N}`$ is not manifestly gauge invariant. Second, the three-dimensional reduction destroys any possibility of obtaining gauge invariance, because we loose the propagation of antiparticles<sup>2</sup><sup>2</sup>2As an example let us consider the amplitude for photon emission off a charged pion. This contribution involves the product of the electromagnetic vertex $`\mathrm{\Gamma }_\pi `$ and the pion propagator $`\mathrm{\Delta }_\pi `$ in the following form: $`\widehat{\mathrm{\Gamma }}_\pi (q\pm k,q).ϵ_\lambda ^{}\mathrm{\Delta }_\pi (q\pm k)=\pm e(2q\pm k).ϵ_\lambda ^{}\frac{1}{(q\pm k)^2m_\pi ^2}`$. When we replace $`ϵ_\lambda k`$ to verify gauge-invariance we obtain, $`\widehat{\mathrm{\Gamma }}_\pi (q\pm k,q).ϵ_\lambda ^{}\mathrm{\Delta }_\pi (q\pm k)=\pm e`$. The three dimensional reduction replaces the full propagator $`\mathrm{\Delta }_\pi (q)`$ by $`\mathrm{\Delta }_\pi ^+(q)=\frac{1}{2\omega (\stackrel{}{q})}\frac{1}{(q^0\omega (\stackrel{}{q}))}`$ and such a relation is no longer fulfilled..
We can follow an alternative procedure that makes the bremsstrahlung amplitude manifestly gauge invariant and where the Thompson reduction does not introduce the problems mentioned above. If we substitute Eqs. (4) and (5) into Eq.(3), only in the one-body component of the amplitude $`M_{\pi N\gamma ,\pi N}^{(1)}M_{\pi N\gamma ,\pi N}(\widehat{V}_{\pi N\gamma ,\pi N}^{(1)})`$, we can isolate the lowest order nonzero contribution of the one-body current. After the three dimensional reduction, the total amplitude can be rewritten as
$$M_{\pi N\gamma ,\pi N}\left[\stackrel{~}{V}_{\pi N\gamma ,\pi N}+\stackrel{~}{M}_{\pi N\gamma ,\pi N}^{pre}+\stackrel{~}{M}_{\pi N\gamma ,\pi N}^{post}+\stackrel{~}{M}_{\pi N\gamma ,\pi N}^{double}\right],$$
(9)
with
$`\stackrel{~}{V}_{\pi N\gamma ,\pi N}`$ $`=`$ $`\overline{u}(\stackrel{}{q}{}_{}{}^{}\stackrel{}{k}/2,ms_f)|\widehat{\stackrel{~}{V}}_{\pi N\gamma ,\pi N}(ϵ_\lambda ,\stackrel{}{k},\stackrel{}{q}{}_{}{}^{},\stackrel{}{q})|u(\stackrel{}{q},ms_i),`$
$`\stackrel{~}{M}_{\pi N\gamma ,\pi N}^{pre}`$ $`=`$ $`{\displaystyle }dq^{\prime \prime 3}\overline{u}(\stackrel{}{q}{}_{}{}^{},ms_f)|\widehat{T}^{()}(\stackrel{}{q}{}_{}{}^{},\stackrel{}{q}^{\prime \prime },z^{})\widehat{G}_{TH}(z^{},\stackrel{}{q}^{\prime \prime })\widehat{\stackrel{~}{V}}_{\pi N\gamma ,\pi N}(ϵ_\lambda ,\stackrel{}{k},\stackrel{}{q}^{\prime \prime },\stackrel{}{q})|u(\stackrel{}{q}+\stackrel{}{k}/2),ms_i,`$
$`\stackrel{~}{M}_{\pi N\gamma ,\pi N}^{post}`$ $`=`$ $`{\displaystyle }dq^{\prime \prime 3}\overline{u}(\stackrel{}{q}{}_{}{}^{}\stackrel{}{k}/2,ms_f)|\widehat{\stackrel{~}{V}}_{\pi N\gamma ,\pi N}(ϵ_\lambda ,\stackrel{}{k},\stackrel{}{q}{}_{}{}^{},\stackrel{}{q}^{\prime \prime })\widehat{G}_{TH}(z,\stackrel{}{q}^{\prime \prime })\widehat{T}(\stackrel{}{q}^{\prime \prime },\stackrel{}{q},z)|u(\stackrel{}{q},ms_i),`$
$`\stackrel{~}{M}_{\pi N\gamma ,\pi N}^{double}`$ $`=`$ $`{\displaystyle }dq^{\prime \prime 3}{\displaystyle }dq^{\prime \prime \prime 3}\overline{u}(\stackrel{}{q}{}_{}{}^{}\stackrel{}{k}/2,ms_f)|\widehat{T}^{()}(\stackrel{}{q}{}_{}{}^{}\stackrel{}{k}/2,\stackrel{}{q}^{\prime \prime },z^{})`$ (10)
$`\widehat{G}_{TH}(z^{},\stackrel{}{q}^{\prime \prime })\widehat{\stackrel{~}{V}}_{\pi N\gamma ,\pi N}(ϵ_\lambda ,\stackrel{}{k},\stackrel{}{q}^{\prime \prime },\stackrel{}{q}^{\prime \prime \prime })\widehat{G}_{TH}(z,\stackrel{}{q}^{\prime \prime \prime })\widehat{T}(\stackrel{}{q}^{\prime \prime \prime },\stackrel{}{q},z)|u(\stackrel{}{q},ms_i)`$
where the current
$`\widehat{\stackrel{~}{V}}_{\pi N\gamma ,\pi N}=i\widehat{V}_{\pi N\gamma ,\pi N}^{(1)}\widehat{G}\widehat{U}+i\widehat{U}^{}\widehat{G}\widehat{V}_{\pi N\gamma ,\pi N}^{(1)}+\widehat{V}_{\pi N\gamma ,\pi N}^{(2)},`$ (11)
generates a gauge invariant Born amplitude $`\stackrel{~}{V}_{\pi N\gamma ,\pi N}`$, that involves all possible ways of attaching a photon to the $`\pi N`$ scattering amplitude $`U`$, and contains the full propagator $`\widehat{G}\widehat{G}_{\pi N}`$. The operator $`\widehat{T}^{()}(z^{})`$, where $`z^{}=z+\omega _\gamma `$, obeys Eq.(6) if we change $`\eta \eta `$ in Eq.(7).
The superscript pre (post) in Eq. (9) indicates that the photon is emitted before (after) the action of the T-matrix, while the superscript double refers to a double-scattering term where the photon is emitted from internal lines between two T-matrices. In the above equations, the Born, pre and double amplitudes were evaluated in the initial center of mass frame ($`\stackrel{}{q}_f=\stackrel{}{q}{}_{}{}^{}\stackrel{}{k}/2,\stackrel{}{p}_f=\stackrel{}{q}{}_{}{}^{}\stackrel{}{k}/2`$, $`\stackrel{}{q}_i=\stackrel{}{p}_i=\stackrel{}{q}`$), while the other (post) amplitude was evaluated in the corresponding final frame ($`\stackrel{}{q}_f=\stackrel{}{p}_f=\stackrel{}{q}^{}`$, $`\stackrel{}{q}_i=\stackrel{}{q}+\stackrel{}{k}/2,\stackrel{}{p}_i=\stackrel{}{q}+\stackrel{}{k}/2`$). The different terms in Eq. (9) are illustrated in Fig. 2a ($`\stackrel{~}{V}_{\pi N\gamma ,\pi N}`$) and in Fig. 2b (remaining terms).
III.SOFT PHOTON APPROXIMATION
In this section we present a brief review of the soft-photon approximation (SPA) to the radiative $`\pi N`$ scattering process. This will be helpful for the introduction of the notation and for later comparison with our dynamical model. Within the SPA the T-matrices can be represented as an expansion in powers of the photon energy $`\omega _\gamma `$. Using the TETA (two energy - two angle) kinematics we can write this expansion as follows:
$`\widehat{T}(s,t,\mathrm{\Delta })`$ $`=`$ $`\widehat{T}(s,t,m^2)+{\displaystyle \frac{\widehat{T}}{\mathrm{\Delta }}}\times {\displaystyle \frac{\mathrm{\Delta }}{\omega _\gamma }}\omega _\gamma +O(\omega _\gamma ^2),`$ (12)
where $`(s,t,\mathrm{\Delta })`$ indicates one of the most convenient set of variables among $`(s_i,t_p,\mathrm{\Delta }_{q_i})`$, $`(s_i,t_q,\mathrm{\Delta }_{p_i})`$, $`(s_f,t_p,\mathrm{\Delta }_{q_f})`$, and $`(s_f,t_q,\mathrm{\Delta }_{p_f})`$ to describe the process. These Lorentz-invariant variables are defined as
$`s_i=(q_i+p_i)^2,s_f=(q_f+p_f)^2,t_p=(p_ip_f)^2,t_q=(q_iq_f)^2,`$
$`\mathrm{\Delta }_{q_i}=(q_ik)^2,\mathrm{\Delta }_{q_f}=(q_f+k)^2,\mathrm{\Delta }_{p_i}=(p_ik)^2,\mathrm{\Delta }_{p_f}=(p_f+k)^2.`$ (13)
If we set $`k=0`$ in the previous equations, we get $`\mathrm{\Delta }_{q_i}=\mathrm{\Delta }_{q_f}=m_\pi ^2`$ and $`\mathrm{\Delta }_{p_i}=\mathrm{\Delta }_{p_f}=m_N^2`$, i.e. they reduce to the particle’s masses in every case. Therefore, when $`k0`$, the $`\mathrm{\Delta }`$ variables provide a convenient set to measure the off-shell character of the intermediate particles, and the derivatives in Eq.(12) account for off-shell effects in the T-matrix.
Within the SPA, the total bremsstrahlung amplitude can be split into external ($`E`$) and internal ($`I`$) contributions:
$`M_{\pi N\gamma ,\pi N}M_{\pi N\gamma ,\pi N}^E+M_{\pi N\gamma ,\pi N}^I,`$ (14)
where we can identify
$`M_{\pi N\gamma ,\pi N}^E`$ $``$ $`\stackrel{~}{M}_{\pi N\gamma ,\pi N}^{pre}(\widehat{\stackrel{~}{V}}_{\pi N\gamma ,\pi N}\widehat{V}_{\pi N\gamma ,\pi N}^{(1)})+\stackrel{~}{M}_{\pi N\gamma ,\pi N}^{post}(\widehat{\stackrel{~}{V}}_{\pi N\gamma ,\pi N}\widehat{V}_{\pi N\gamma ,\pi N}^{(1)}),`$
and the internal contribution $`M_{\pi N\gamma ,\pi N}^I`$ can be obtained “by imposing” the gauge invariance condition
$$M_{\pi N\gamma ,\pi N}(ϵ_\lambda ^\mu =k^\mu )=0.$$
(15)
In this approach, the total amplitude depends only on the static electromagnetic properties of the external particles and on the elastic $`\pi N`$ scattering amplitude (see Eq. (16) below) . Therefore, the total amplitude at this order eliminates any dependence on off-shell effects and model-dependent contributions.
Up to terms of $`O(\omega _\gamma ^0)`$, the total amplitude is given by:
$`M_{\pi N\gamma ,\pi N}`$
$`\overline{u}(\stackrel{}{p}_+^{},m_{s_f})|\{\widehat{e}_\pi [{\displaystyle \frac{q_{}^{}ϵ_\lambda }{q_{}^{}k}}{\displaystyle \frac{(q_{}^{}+p_+^{})ϵ_\lambda }{k(q_{}^{}+p_+^{})}}]\widehat{T}(s_i,t_p,m_\pi ^2)`$
$`+[{\displaystyle \frac{(\widehat{e}_Np_+^{}\widehat{R}(p_+^{}))ϵ_\lambda }{p_+^{}k}}{\displaystyle \frac{(\widehat{e}_N(q_{}^{}+p_+^{})\widehat{R}(p_+^{}))ϵ_\lambda }{k(q_{}^{}+p_+^{})}}]\widehat{T}(s_i,t_q,m_N^2)\}|u(\stackrel{}{p},m_{s_i})`$
$`\overline{u}(\stackrel{}{p}^{},m_{s_f})|\{\widehat{e}_\pi \widehat{T}(s_f,t_p^{},m_\pi ^2)[{\displaystyle \frac{q_+ϵ_\lambda }{q_+k}}{\displaystyle \frac{(q_++p_{})ϵ_\lambda }{k(q_++p_{})}}]`$
$`+\widehat{T}(s_f,t_q^{},m_N^2)[{\displaystyle \frac{(\widehat{e}_Np_{}\widehat{R}(p_{}))ϵ_\lambda }{p_{}k}}{\displaystyle \frac{(\widehat{e}_N(q_++p_{})\widehat{R}(p_{}))ϵ_\lambda }{k(q_++p_{})}}]\}|u(\stackrel{}{p}_{},m_{s_i}),`$
(16)
where
$`q(\omega (\stackrel{}{q}),\stackrel{}{q}),p(E(\stackrel{}{q}),\stackrel{}{q}),q^{}(\omega (\stackrel{}{q}{}_{}{}^{}),\stackrel{}{q}{}_{}{}^{}),p^{}(E(\stackrel{}{q}{}_{}{}^{}),\stackrel{}{q}{}_{}{}^{}),`$
$`q_{}^{}(\omega (\stackrel{}{q}{}_{}{}^{}),\stackrel{}{q}{}_{}{}^{}),p_+^{}(E(\stackrel{}{q}{}_{+}{}^{}),\stackrel{}{q}{}_{+}{}^{}),q_+(\omega (\stackrel{}{q}_+),\stackrel{}{q}_+),p_{}(E(\stackrel{}{q}_{}),\stackrel{}{q}_{}),`$
$`t_p=(p_+^{}p)^2,t_q=(q_+^{}q)^2,t_p^{}=(p_{}p^{})^2,t_q^{}=(q_+q^{})^2,`$
$`\widehat{R}_\mu (x){\displaystyle \frac{1}{4}}\widehat{e}_N[\overline{)}k,\gamma _\mu ]+{\displaystyle \frac{\widehat{\kappa }_N}{8m_N}}\{[\overline{)}k,\gamma _\mu ],\overline{)}x\},\stackrel{}{q}_\pm =\stackrel{}{q}\pm \stackrel{}{k}/2,\stackrel{}{q}{}_{\pm }{}^{}=\stackrel{}{q}{}_{}{}^{}\pm \stackrel{}{k}/2,`$ (17)
and $`\widehat{e}_\pi =e𝒯_z`$, $`\widehat{e}_N=e(1+\tau _z)/2`$, and $`\widehat{\kappa }_N=\kappa _p(1+\tau _z)/2+\kappa _n(1\tau _z)/2`$ denote the charge and anomalous magnetic moment operators of pions and nucleons. As anticipated, Eq.(16) shows that within the SPA the $`M_{\pi N\gamma ,\pi N}`$ amplitude depends only on the elastic T-matrix, because derivative terms of $`\widehat{T}`$ cancels in the addition of internal and external contributions. Let us emphasize that any dependence on the structure of internal contributions (in particular, the dependence of off-shell effects) are of higher order in $`\omega _\gamma `$ and must be included explicitly in the amplitude in a gauge invariant way. This is the purpose of the forthcoming section.
IV. DYNAMICAL MODEL
In this section we compute the bremsstrahlung amplitude along the lines developed in Eqs. (9-11), using a potential $`\widehat{U}`$ obtained from effective Lagrangians, and three specific models for the T-matrix to describe the $`\pi N`$ rescattering.
The three models for the T-matrix will be called OBQA, SEP and NEW, respectively. The OBQA version for the T-matrix interaction, is based on a model that includes $`\pi `$ and $`\rho `$ mesons exchange through a correlated $`2\pi `$ exchange potential. The SEP model for $`\pi N`$ rescattering is generated by a phenomenological separable potential. Finally, the NEW model is obtained from the exchange of $`\pi `$ and (sharp) $`\rho `$ mesons.
The operator $`\widehat{U}`$ is constructed from a Lagrangian density that includes the nucleon ($`N`$), the $`\mathrm{\Delta }`$-isobar, and the $`\pi `$, $`\rho `$ and $`\sigma `$ mesons in the following form
$$\widehat{}_{hadr}=\widehat{}_{\pi NN}+\widehat{}_{N\mathrm{\Delta }\pi }+\widehat{}_{\rho NN}+\widehat{}_{\rho \pi \pi }+\widehat{}_{\sigma NN}+\widehat{}_{\sigma \pi \pi }$$
(18)
where the individual terms are given by
$`\widehat{}_{\pi NN}(x)`$ $`=`$ $`\left({\displaystyle \frac{f_{\pi NN}}{m_\pi }}\right)\overline{\psi }(x)\gamma _5\stackrel{}{\tau }.(\overline{)}\stackrel{}{\pi }(x))\psi (x),`$
$`\widehat{}_{N\mathrm{\Delta }\pi }(x)`$ $`=`$ $`\left({\displaystyle \frac{f_{N\mathrm{\Delta }\pi }}{m_\pi }}\right)\overline{\psi }_\mathrm{\Delta }^\mu (x)\stackrel{}{T}^{}.(_\mu \stackrel{}{\pi }(x))\psi (x)+h.c.,`$
$`\widehat{}_{\rho NN}(x)`$ $`=`$ $`{\displaystyle \frac{1}{2}}g_\rho \overline{\psi }(x)\left[\gamma _\mu {\displaystyle \frac{\kappa _\rho }{2m_N}}\sigma _{\mu \nu }^\nu \right]\stackrel{}{\tau }.\stackrel{}{\rho }^\mu (x)\psi (x),`$
$`\widehat{}_{\rho \pi \pi }(x)`$ $`=`$ $`g_{\rho \pi \pi }\stackrel{}{\rho }_\mu (x).(\stackrel{}{\pi }(x)^\mu \stackrel{}{\pi }(x)),`$
$`\widehat{}_{\sigma NN}(x)`$ $`=`$ $`g_\sigma \overline{\psi }(x)\psi (x)\sigma (x),`$
$`\widehat{}_{\sigma \pi \pi }(x)`$ $`=`$ $`\left({\displaystyle \frac{g_{\sigma \pi \pi }}{2m_\pi }}\right)\sigma (x)(_\mu \stackrel{}{\pi }(x)).(^\mu \stackrel{}{\pi }(x)).`$ (19)
The isotopic fields $`\psi (x)`$ and $`\psi _\mathrm{\Delta }^\mu (x)`$ denote the $`N`$ and $`\mathrm{\Delta }`$ baryons, respectively, while $`\stackrel{}{\pi }(x)`$, $`\stackrel{}{\rho }^\mu (x)`$ and $`\sigma (x)`$ denote the pion, $`\rho `$-meson and $`\sigma `$-meson fields. The arrow over the meson fields refers to the isospin space. $`\stackrel{}{T}^{}`$ stands for the isospin $`1/2`$ to $`3/2`$ transition operator, and $`f_{\pi NN}`$, $`g_\rho `$( $`\kappa _\rho `$), $`g_\sigma `$, $`f_{N\mathrm{\Delta }\pi }`$, $`g_{\rho \pi \pi }`$, and $`g_{\sigma \pi \pi }`$ are the corresponding coupling constants.
The electromagnetic currents can be obtained from the Lagrangian density
$$\widehat{}_{elec}=\widehat{}_{\gamma NN}+\widehat{}_{\gamma \pi NN}+\widehat{}_{\gamma \pi \pi }+\widehat{}_{\gamma \mathrm{\Delta }\mathrm{\Delta }}+\widehat{}_{\gamma \pi N\mathrm{\Delta }}+\widehat{}_{\gamma \rho \pi \pi }+\widehat{}_{\gamma \sigma \pi \pi }$$
(20)
with
$`\widehat{}_{\gamma NN}(x)`$ $`=`$ $`e\overline{\psi }(x)\left[\widehat{e}_N\gamma _\mu {\displaystyle \frac{\widehat{\kappa }}{2m_N}}\sigma _{\mu \nu }^\nu \right]A^\mu (x)\psi (x),`$
$`\widehat{}_{\gamma \pi NN}(x)`$ $`=`$ $`e\left({\displaystyle \frac{f_{\pi NN}}{m_\pi }}\right)\overline{\psi }(x)\gamma _5\gamma _\mu \left[\stackrel{}{\tau }\times \stackrel{}{\pi }(x)\right]_3\psi (x)A^\mu (x),`$
$`\widehat{}_{\gamma \pi \pi }(x)`$ $`=`$ $`e\left[\stackrel{}{\pi }(x)\times _\mu \stackrel{}{\pi }(x)\right]_3A^\mu (x),`$
$`\widehat{}_{\gamma \mathrm{\Delta }\mathrm{\Delta }}(x)`$ $`=`$ $`e\overline{\psi }_\mathrm{\Delta }^\nu (x)\widehat{\mathrm{\Gamma }}_{\nu \mu \alpha }\psi _\mathrm{\Delta }^\mu (x)A^\alpha (x),`$
$`\widehat{}_{\gamma \pi N\mathrm{\Delta }}(x)`$ $`=`$ $`e\left({\displaystyle \frac{f_{N\mathrm{\Delta }\pi }}{m_\pi }}\right)\overline{\psi }_\mathrm{\Delta }^\mu (x)\left[\stackrel{}{\pi }(x)\times \stackrel{}{T}^{}\right]_3\psi (x)A_\mu (x)+h.c.,`$
$`\widehat{}_{\gamma \rho \pi \pi }(x)`$ $`=`$ $`eg_\rho \{[\stackrel{}{\pi }(x).\stackrel{}{\pi }(x)\pi _3(x)\pi _3(x)]\rho _3^\nu (x)`$
$`\pi _3(x)[\stackrel{}{\pi }(x).\stackrel{}{\rho }^\nu (x)\pi _3(x)\rho _3^\nu (x)]\}A_\nu (x),`$
$`\widehat{}_{\gamma \sigma \pi \pi }(x)`$ $`=`$ $`2e\left({\displaystyle \frac{g_{\sigma \pi \pi }}{2m_\pi }}\right)\sigma (x)\left[\stackrel{}{\pi }(x)\times _\mu \stackrel{}{\pi }(x)\right]_3A^\mu (x).`$ (21)
The electromagnetic vertex operator of the $`\mathrm{\Delta }`$-isobar is given by
$`\widehat{\mathrm{\Gamma }}_{\nu \mu \alpha }`$ $`=`$ $`\widehat{e}_\mathrm{\Delta }\left[(\gamma _\alpha g_{\nu \mu }{\displaystyle \frac{1}{3}}\gamma _\alpha \gamma _\nu \gamma _\mu {\displaystyle \frac{1}{3}}\gamma _\nu g_{\mu \alpha }+{\displaystyle \frac{1}{3}}\gamma _\mu g_{\nu \alpha }){\displaystyle \frac{\widehat{\kappa }_\mathrm{\Delta }}{2m_\mathrm{\Delta }}}\sigma _{\alpha \beta }k_\beta g_{\nu \mu }\right],`$
where $`F_{\mu \nu }=_\mu A_\nu (x)_\nu A_\mu (x)`$, $`A_\nu (x)`$ being the electromagnetic four-potential. $`\widehat{\kappa }_\mathrm{\Delta }`$ and $`\widehat{e}_\mathrm{\Delta }`$ are the anomalous magnetic moment<sup>3</sup><sup>3</sup>3We restrict ourselves to the $`\mathrm{\Delta }^{++}`$ contribution, the only one for which we have experimental information on $`\kappa _\mathrm{\Delta }`$ . and charge operators whose action upon the Rarita-Schwinger field give as eigenvalues the corresponding values of these properties of the $`\mathrm{\Delta }`$-isobar.
The propagators for the $`\pi (\sigma ),\rho ,N`$ and $`\mathrm{\Delta }`$ hadrons obtained from the above Lagrangians are given, respectively, by
$`\mathrm{\Delta }_{\pi ,\sigma }(q)`$ $`=`$ $`{\displaystyle \frac{1}{q^2m_{\pi ,\sigma }^2+i\eta }},`$
$`D_\rho ^{\mu \nu }(q)`$ $`=`$ $`{\displaystyle \frac{g^{\mu \nu }q^\mu q^\nu /m_\rho ^2}{q^2m_\rho ^2+i\eta }},`$
$`S(q)`$ $`=`$ $`{\displaystyle \frac{\overline{)}q+m_N}{q^2m_N^2+i\eta }}`$
$`G^{\mu \nu }(q)`$ $`=`$ $`{\displaystyle \frac{\overline{)}q+m_\mathrm{\Delta }}{q^2m_\mathrm{\Delta }^2+i\eta }}\left[g^{\mu \nu }+{\displaystyle \frac{1}{3}}\gamma ^\mu \gamma ^\nu +{\displaystyle \frac{2}{3}}{\displaystyle \frac{q^\mu q^\nu }{m_\mathrm{\Delta }^2}}{\displaystyle \frac{1}{3}}{\displaystyle \frac{q^\mu \gamma ^\nu q^\nu \gamma ^\mu }{m_\mathrm{\Delta }}}\right]`$ (23)
$`{\displaystyle \frac{2}{3m_\mathrm{\Delta }^2}}(q^2m_\mathrm{\Delta }^2)\left[{\displaystyle \frac{}{}}(\gamma ^\mu q^\nu \gamma ^\nu q^\mu )+(\overline{)}q+m_\mathrm{\Delta })\gamma ^\mu \gamma ^\nu \right]`$
Observe that we have kept the off-shell part of $`G^{\mu \nu }(q)`$ in Eqs. (23); without this term it is impossible to get gauge invariance consistently using simultaneously the vertex given in Eq.(LABEL:35) . As was mentioned previously, the potential operator $`\widehat{U}`$ can be computed by using these Feynman rules. The scattering amplitude $`U`$ is depicted in Fig.3, while the amplitude $`\stackrel{~}{V}_{\pi N\gamma ,\pi N}`$ can be obtained by coupling the photon to all diagrams in $`U`$ as shown in fig.4.
Good convergence properties of the scattering equations given in Eqs. (10) can be obtained by introducing hadronic form factors, which are supposed to describe the composite nature of hadrons. It is a common practice to use different parametrizations of the form factors for different T-matrices. For example, the OBQA model uses monopole and dipole forms with cutoff parameters ranging from $`\mathrm{\Lambda }=12001600MeV`$ . In the case of the SEP interaction different form factors are introduced for each partial wave component , while in the NEW model form factors usually advocated are of monopolar form with $`\mathrm{\Lambda }=13002300MeV`$ . However, the introduction of form factors replacing point vertexes in $`\stackrel{~}{V}_{\pi N\gamma ,\pi N}`$ spoils the gauge invariance of the total amplitude. Fortunately, the gauge invariance of the amplitude can be recovered by using the method of Gross and Riska which, however, does not yield unique electromagnetic couplings to hadrons. Therefore, we follow the more simple prescription of using a common form factor of monopole type
$$f(\stackrel{}{q}^{\prime \prime })=\frac{\mathrm{\Lambda }^2}{\mathrm{\Lambda }^2+\stackrel{}{q}^{\prime \prime 2}},$$
(24)
where the scale $`\mathrm{\Lambda }`$ can be adjusted at a given incident energy for each T-matrix model.
V. NUMERICAL RESULTS AND CONCLUSIONS
The differential cross section $`d\sigma /d\mathrm{\Omega }_\pi d\mathrm{\Omega }_\gamma d\omega _\gamma `$ to be compared with experimental data can be obtained from Eq.(1), where the amplitude $`M_{\pi N\gamma ,\pi N}`$ is calculated from Eqs.(9-11). The dynamical model approximation (DMA) advocated in the present paper contains the following steps. The current operator $`\widehat{\stackrel{~}{V}}_{\pi N\gamma ,\pi N}`$ is computed from the effective Lagrangians given in Eqs.(18) and (20), with the propagators obtained in Eqs.(23), and the monopole form factor given in Eq.(24) to have good convergence of the intermediate momentum integrals.
As is known, double scattering terms have a significant contribution in the case of $`protonproton`$ bremsstrahlung, mainly in the end point region of $`\omega _\gamma `$ . In the present work we will neglect $`\stackrel{~}{M}_{\pi N\gamma ,\pi N}^{double}`$ in Eq.(10), because the numerical calculation of the three-dimensional integrals requires an enormous computational effort. Nevertheless, we keep double scattering-like contributions in post and pre amplitudes coming from current components $`i\widehat{U}^{}\widehat{G}\widehat{V}_{\pi N\gamma ,\pi N}^{(1)}`$ and $`i\widehat{V}_{\pi N\gamma ,\pi N}^{(1)}\widehat{G}\widehat{U}`$, respectively. In order to compare the approaches provided by the DMA and SPA, we fix $`M_{\pi N\gamma ,\pi N}`$ to coincide quantitatively at low photon energies.
For illustration purposes, we will implement the DMA approach with the OBQA, SEP and NEW T-matrices for the specific example of $`\pi ^+p\pi ^+p\gamma `$. In this case, we have contributions coming from the diagrams depicted in Figs. 4a (with intermediate $`\mathrm{\Delta }`$), 4b (with intermediate $`N`$), and 4c. The different coupling constants and masses needed to evaluate $`\widehat{V}_{\pi N\gamma ,\pi N}`$ were taken from the model II of ref., from , and from , and are displayed in table 1. For direct pole diagrams we use bare masses and coupling constants since they get dynamically dressed in the T-matrix scattering Eq.(5). In the OBQA case, we replace the $`2\pi `$ correlated exchange potential by $`\pi `$ and $`\rho `$ sharp mass exchange terms, since they do not lead to sizable differences as shown in ref.. In the SEP case we use the same coupling constants and masses, because the scattering potential is not generated from a dynamical model.
We will compare the theoretical predictions to the experimental cross sections measured by Nefkens (EXP I) and Meyer (EXP II), which have been reported for different kinematical configurations. In EXP I the pions were detected at fixed angles $`\theta _\pi =50.5^0,\varphi _\pi =180^0`$, for three different energies of incident pions (269, 298 and 324 MeV), and the photons were detected at various $`\theta _\gamma ,\varphi _\gamma `$ angles in the range of energies $`\omega _\gamma =0150`$ MeV. In EXP II $`\theta _\pi `$ ranges from 55<sup>0</sup> to 95<sup>0</sup>, the incident energy of the pion is fixed at 299 MeV and the photons were detected at angles $`\theta _\gamma =120^0,\varphi _\gamma =0^0`$ with energies in the range $`\omega _\gamma =0140`$ MeV. In EXP II the resulting cross section measurements were averaged over $`\theta _\pi `$.
Our results for the cross sections in the DMA approach, for the respective ansatz of the T-matrix, are shown in Figs. 5 and 6. The SPA and experimental values are also plotted for comparison. The predictions of the DMA approach shown in fig. 5 using the three models of the T-matrix, are compared to the results of EXP I for photon angles given by $`G_{14}(\theta _\gamma =103^0,\varphi _\gamma =180^0)`$. Since the parameters entering T-matrices are usually quoted to reproduce the elastic phase shifts, the cutoff parameters $`\mathrm{\Lambda }`$ were fixed in order to have good coincidence between the predictions of the different interactions and the SPA at low $`\omega _\gamma `$ values. We get (in MeV units) $`\mathrm{\Lambda }_{OBQA}=750,700,600`$, $`\mathrm{\Lambda }_{SEP}=700,600,500`$, and $`\mathrm{\Lambda }_{NEW}=550,500,450`$, for incident pion energies $`T_{lab}=269,298,324`$ respectively. Observe that the value of $`\mathrm{\Lambda }`$ found in the case of the SEP interaction for $`T_{lab}=298MeV`$ is consistent with the one previously found for pion photo-production at $`T_{lab}=300MeV`$, while the OBQA values are roughly consistent with the form factors used in ref., which correspond to a monopole form-factor with $`\mathrm{\Lambda }800MeV`$. On the other hand, the resulting cross sections for the conditions of EXP II are shown in Fig. 6. In this case, the differential cross section is averaged over $`\theta _\pi =55^075^0`$, and $`\theta _\pi =75^095^0`$, for $`T_{lab}=299MeV`$. However, we keep for consistency the values of $`\mathrm{\Lambda }`$ obtained at $`T_{lab}=298MeV`$ in EXP I.
As we can observe, the SPA reproduces very well the experimental cross section for EXP I in the whole range of measured photon energies, while it gives values somewhat below those obtained in EXP II. According to ref. , the agreement with EXP II might be improved if the emission of the photon from the $`\mathrm{\Delta }^{++}`$( seventh graph in fig.3a) is included explicitly as a piece of the internal amplitude.
In almost all the cases, the predictions of the DMA lies above the experimental cross section and the SPA for energies $`\omega _\gamma >20MeV`$. One of the reasons for this may be the use of an overall form factor to cure the gauge invariance problems. The total bremsstrahlung amplitude is built up, as can be seen from Eqs. (9) and (10), by adding different components. It is not expected that the common form factor works as satisfactorily upon adding up these components as it does the one used to generate the individual T-matrices, which change their values from vertex to vertex. The comparison between the results of the SPA and DMA schemes shows that the additional off-shell effects, added coherently to the lowest order contributions, may have important contributions since they do not cancel exactly as the derivative terms of the T-matrix appearing from the soft-photon expansion.
From fig.5 we can check that the SEP interaction provides the closest results to the experimental cross section with deviations starting for $`\omega _\gamma >40MeV`$. This indicates that the dynamical model involved in the SEP interaction gives the smallest off-shell effects. On the other hand, the strongest off-shell effects appear in the OBQA model. This conclusion agrees with a previous study on observables in pion photo-production experiments. Similar conclusions can be drawn from Fig. 6. Note, however, that the SPA lies somewhat below the experimental results in this case, which seems to indicate the necessity of additional dynamical degrees of freedom.
As was discussed in section III, the off-shell contributions to the external and internal amplitudes within the SPA cancel each other, thus we cannot study these effects within this approximation. In addition since we get gauge invariance in the SPA by adjusting the internal amplitude, the gauge-invariant electromagnetic currents remain hidden. In the DMA approach, these cancellations must occur explicitly between the different components of the amplitude (the so-called Born,pre,post contributions). The departure of the different T-matrices from the SPA can be used to estimate the size of unbalanced off-shell terms and provide a test of their off-shell behaviors. Also, since the gauge-invariant electromagnetic current is constructed explicitly from effective Lagrangians, we can use the radiative $`\pi N`$ reaction to study the relevance of the degrees of freedom and the parameters involved in this dynamical model.
Finally one may wonder about the relevance of additional contributions not included in the present formulation of effective Lagrangians, for example the $`a_1`$, $`\omega `$ mesons. A preliminary estimation of these effects shows that, in the region of photon energies considered in this work $`\omega _\gamma =0150MeV`$, their contribution to the amplitude are suppressed by an order of magnitude with respect to the degrees of freedom considered here. In addition it may be required to include the double-scattering terms, in order to provide a better fit of experimental data. These considerations, together with a partial wave analysis of the studied T-matrices and a more detailed analysis of each particular contribution to the electromagnetic current are beyond the scope of the present paper.
ACKNOWLEDGMENTS
The work of A. Mariano was supported in part by Conacyt (México) through the Fondo de Cátedras Patrimoniales de Excelencia Nivel II, and Conicet (Argentina). He is also grateful to K. Nakayama for the valuable support in making this work. G. López Castro was partially supported by Conacyt under contract 32429-E.
FIGURE CAPTIONS
Fig.1 One- and two-body contributions to the bremsstrahlung current amplitude.
Fig.2 (a) Gauge-invariant bremsstrahlung current amplitude. (b) Post-, pre- and double-scattering amplitude contributions (se Eq. (9)).
Fig.3 Born amplitude corresponding to the $`\pi N`$ potential operator $`\widehat{U}`$. The first diagram denote the nucleon-pole, and the $`\mathrm{\Delta }`$-pole corresponds to the third diagram. The fifth and sixth diagrams correspond to $`\rho `$ and $`\sigma `$ mesons exchange.
Fig.4 The gauge-invariant amplitude obtained by coupling a photon to the Born terms in fig.3, together with the two-meson exchange currents (fifth, sixth and seventh graphs in each line): (a) Contributions obtained from the N and $`\mathrm{\Delta }`$ direct-pole diagrams in fig.3; (b) Terms generated by the cross-pole diagrams in fig.3, and (c) diagrams obtained from $`\rho `$ and $`\sigma `$ exchange contributions in fig.3.
Fig.5 $`\pi N\gamma `$ cross section for $`T_{lab}=269`$, $`298`$ and $`324MeV`$ and $`G_{14}(\theta _\gamma =103^0,\varphi _\gamma =180^0)`$ in EXPI, calculated in the DMA for the different T-matrices. We also include the SPA cross section and the measured values.
Fig.6 Same as fig.5 for $`T_{lab}=299MeV`$. The cross sections are averaged over the angles $`\theta _\pi =5575^0`$ (upper plot) and $`\theta _\pi =7595^0`$(lower plot) in EXPII.
|
warning/0003/cond-mat0003065.html
|
ar5iv
|
text
|
# Modeling of the Magnetic Susceptibilities of the Ambient- and High-Pressure Phases of (VO)_𝟐P_𝟐O_𝟕
## I Introduction
A resurgence of research on the magnetic properties of low-dimensional quantum spin systems has occurred over the last decade. This work was mainly initially undertaken to understand the relationships between the magnetic properties of layered cuprates containing Cu<sup>+2</sup> spin-1/2 square lattice layers and the high superconducting transition temperatures of the doped materials. This goal has also spawned much research on related one- and two-dimensional (1D and 2D) spin systems. Indeed, the new subfield of spin ladder physics was created as a result of these efforts. The basic $`n`$-leg spin ladder consists of a planar arrangement of $`n`$ parallel vertical spin chains (the ladder “legs”) with (nonfrustrating) horizontal nearest neighbor couplings between adjacent chains, i.e., across the ladder “rungs”.
Self-doped two-leg spin ladders are realized in the compound $`\mathrm{NaV}_2\mathrm{O}_5`$, in which the V atoms are crystallographically equivalent and the oxidation state of the V atoms is +4.5, resulting formally in a mixed-valent $`d^{0.5}`$ system. However, the material is a semiconductor rather than a metal. Theoretical studies have indicated that the reason for this is that the one $`d`$ electron per two V atoms is localized on the respective V-O-V rung of the ladder due to the on-site Coulomb repulsion on each rung. This in turn led to the hypothesis that each rung acts as a spin-1/2 site, in which case the compound is expected to
behave magnetically like a spin $`S=1/2`$ antiferromagnetic (AF) Heisenberg chain. Additionally, the observation of a spin dimerization transition below 34 K, accompanied by a lattice distortion, is consistent with this scenario. These and other aspects of the thermal and magnetic behaviors of $`\mathrm{NaV}_2\mathrm{O}_5`$ were recently examined in detail in a combined theoretical and experimental study. As part of this study, an accurate function was generated for the theoretical magnetic susceptibility $`\chi `$ versus temperature $`T`$ of the spin $`S=1/2`$ AF alternating-exchange ($`J_1`$ and $`J_2`$ with $`J_2J_1`$) Heisenberg chain over the entire range $`0\alpha 1`$ of the alternation parameter $`\alpha J_2/J_1`$.
The availability of this high-accuracy theoretical $`\chi (k_\mathrm{B}T/J_1,\alpha )`$ function now allows one to accurately and precisely test the consistency of proposals for the occurrence of alternating-exchange chains in specific compounds by comparing the observed $`\chi (T)`$ with that expected theoretically. In particular, in the present work we use this method to test the consistency of recently
proposed alternating-exhange chain models for the ambient- and high-pressure phases of vanadyl pyrophosphate, $`(\mathrm{VO})_2\mathrm{P}_2\mathrm{O}_7`$, and we obtain the exchange constants and spin gaps in the respective chains. The influence of interchain coupling on the derived intrachain exchange constants and spin gaps is investigated using a mean-field approximation for the interchain coupling. The accuracies of the spin gaps obtained using fits to the low-$`T`$ $`\chi (T)`$ data by theoretical low-$`T`$ approximations are determined. The results of these studies are compared with previously reported estimates of the exchange couplings and spin gaps in the two phases and with the dispersion relations measured for the ambient-pressure
phase by inelastic neutron scattering.
The history of the study of the magnetic properties of $`(\mathrm{VO})_2\mathrm{P}_2\mathrm{O}_7`$ is interesting and extensive. In the following two sections we give brief overviews of the previous work on the ambient-pressure (AP) and high-pressure (HP) phases of this compound, respectively, and then present the plan for the remainder of the paper.
### A AP-$`\mathbf{(}\mathrm{𝐕𝐎}\mathbf{)}_\mathrm{𝟐}𝐏_\mathrm{𝟐}𝐎_\mathrm{𝟕}`$
The V$`{}_{}{}^{+4}d_{}^{1}`$ ambient-pressure phase AP-$`(\mathrm{VO})_2\mathrm{P}_2\mathrm{O}_7`$, sometimes abbreviated in the recent literature as “VOPO”, is an industrial catalyst for the selective oxidation of $`n`$-butane to maleic anhydride. This compound was found to have an orthorhombic crystal structure (containing four inequivalent types of V atoms) which can be viewed crystallographically as containing $`S=1/2`$ two-leg ladders, where the rungs of the ladder lie along the $`c`$-axis and the legs are oriented along the $`a`$-axis, as sketched in Fig. 1. A single crystal x-ray diffraction structural study claimed that the previous structural studies were incorrect and that the structure is monoclinic with eight inequivalent V atoms in the unit cell, although the overall structural features and the unit cell dimensions were found to be very similar to those of the previously proposed orthorhombic structure. However, a recent study of a polycrystalline sample using both x-ray and neutron diffraction Rietveld refinements and transmission electron microscopy confirmed the orthorhombic structure and ruled out the monoclinic structure; it was suggested that whether the orthorhombic or monoclinic structure occurs in a particular sample may depend on the exact composition and the details of sample synthesis. There have been two conventions used in the literature for designating the $`b`$ and $`c`$ axes in which these two axes are mutually interchanged; we will adhere to the convention in Fig. 1, as in Ref. , for which the approximate lattice parameters are listed in the figure caption.
The magnetic susceptibility versus temperature $`\chi (T)`$ of AP-$`(\mathrm{VO})_2\mathrm{P}_2\mathrm{O}_7`$ was found to exhibit an energy gap (“spin gap”) for magnetic excitations. The $`\chi (T)`$ was initially fitted to high precision by the prediction for the $`S=1/2`$ AF alternating-exchange Heisenberg chain, with the exchange constants $`J_1`$ and $`J_2`$ and alternation parameter $`\alpha J_2/J_1`$ given in Table I. A fit by the spin ladder model as suggested from the crystal structure (see Figs. 1 and 2) was not possible at that time (1987) due to lack of theoretical predictions for $`\chi (T)`$ of this model. It was also speculated that even though the crystallographic features suggest a spin ladder model, the actual magnetic interactions might turn out to correspond to those of alternating-exchange spin chains. When $`\chi (T)`$ calculations for the spin ladder model were eventually done, it was found that the same experimental $`\chi (T)`$ data set could be fitted by the spin ladder model to the same high precision as for the very different alternating-exchange chain model. The existence of a spin gap was subsequently confirmed and its value estimated from NMR measurements of the <sup>31</sup>P Knight shift $`{}_{}{}^{31}K(T)`$ and nuclear spin-lattice relaxation rate $`{}_{}{}^{31}(1/T_1)(T)`$, from $`\chi (T)`$ (Ref. ) and inelastic neutron scattering measurements on polycrystalline samples at low temperatures, and most recently from elastic constant measurements on a single crystal in pulsed magnetic fields up to 50 T at $`T=1.6`$ K and from Raman scattering intensity measurements on single crystals at low temperatures, as listed in Table I.
The above inelastic neutron scattering measurements on polycrystalline AP-$`(\mathrm{VO})_2\mathrm{P}_2\mathrm{O}_7`$ reportedly confirmed the spin ladder model and rejected the alternating-exchange chain model by a comparison of the observed spin gap \[43(2) K\] with the values 45.7 K and 57 K predicted for the two respective models from the sets of exchange constants determined from respective fits to the $`\chi (T)`$ data. However, subsequent inelastic neutron scattering results on polycrystalline samples and especially on a collection of about 200 oriented small single crystals proved that AP-$`(\mathrm{VO})_2\mathrm{P}_2\mathrm{O}_7`$ is not a spin-ladder compound. The strongest dispersion of the one-magnon spectra of the single crystals was found to be along the $`c`$ axis, i.e., in the direction of the structural ladder rungs, and the coupling in the direction of the ladder legs was found to be weakly ferromagnetic. Thus, perhaps surprisingly, the superexchange coupling path V-O-P-O-V along the $`c`$ axis, coupling the structural two-leg ladders as shown in Fig. 2, is much stronger than the shorter V-O-V coupling along the ladder legs parallel to the $`a`$-axis. These results were interpreted in terms of an alternating-exchange chain model with the chains running along the $`c`$ axis, with weak coupling between the chains, thus confirming the above speculation in Ref. . Subsequent $`\chi (T)`$ data for powder and single crystal samples have been interpreted in terms of the alternating-exchange chain model, with exchange parameters listed in Table I.
The spin gaps of 43(2) K (Ref. ) and 40.4(4) K (Ref. ) found from inelastic neutron scattering measurements on powder samples are both significantly larger than the one-magnon spin gap of 36.2(3) K determined from the neutron scattering measurements on single crystals. Since powder samples have usually been found to show relatively high levels of paramagnetic impurities and/or defects, this comparison suggests that the larger spin gaps in the powder samples may be a real effect arising from termination of the spin chains by defects. One would indeed expect finite segments of alternating-exchange chains to exhibit larger spin gaps than for the infinite chain.
We have also listed in Table I the intradimer exchange constant determined from $`\chi (T)`$ (Ref. ), <sup>31</sup>P NMR Knight shift $`{}_{}{}^{31}K(T)`$ (Ref. ), and inelastic neutron scattering measurements of polycrystalline samples of the $`d^1`$ $`S=1/2`$ vanadium dimer compound
$`(\mathrm{VO})\mathrm{HPO}_4\frac{1}{2}\mathrm{H}_2\mathrm{O}`$ (vanadyl hydrogen phosphate hemihydrate), which is a precursor for the synthesis of, and has structural similarities to, AP-$`(\mathrm{VO})_2\mathrm{P}_2\mathrm{O}_7`$. In particular, the neutron scattering study of this compound confirmed the importance and strength of the V-O-P-O-V superexchange pathway.
From the above neutron scattering measurements on single crystals of AP-$`(\mathrm{VO})_2\mathrm{P}_2\mathrm{O}_7`$, a second spin gap at a larger energy of 67 K was found in addition to the gap of 36 K for coherent one-magnon propagation along the $`c`$ axis. The two spin gaps cannot both arise from one-magnon excitations in an isolated alternating-exchange chain and and the larger one was suggested to arise from neutron scattering from two-magnon triplet bound states of such chains, although the scattered neutron intensity was larger than expected from scattering from such states. A 2D model incorporating both nonfrustrating ($`J_a`$, see Fig. 1) and diagonal frustrating ($`J_\times `$) AF interactions between alternating-exchange chains was subsequently proposed. The intrachain and interchain ex-
change constants were determined by a (very good) fit to the one-magnon dispersion relation (including the one-magnon spin gap of 36 K) measured in the inelastic neutron scattering measurements and by using the results of high-$`T`$ $`\chi (T)`$ measurements. The two alternating AF intrachain exchange constants so determined are listed in Table I. The two AF interchain interactions were predicted to be rather large: $`J_a/J_1=0.203`$ and $`J_\times /J_1=0.255`$. However, the $`\chi (T)`$ calculated using this model in Ref. was found to be in poor agreement with experimental single-crystal $`\chi (T)`$ data even at the highest measured temperature $`T200`$ K, which is significantly above the temperature ($`70`$ K) at which $`\chi (T)`$ shows a broad maximum and where the prediction would be expected to be quite accurate. On the other hand, we find in Sec. III B 2 below that their $`\chi (T)`$ prediction is in quite good agreement with experimental data if the comparison is done in a somewhat different way. Additional calculations indicated that the frustrating AF interchain interaction stabilizes two-magnon bound states and thus supported the conjecture that the higher energy mode at 67 K is a triplet two-magnon bound state. The same 2D model containing alternating-exchange chains and frustrating AF interchain couplings and also a model containing ferromagnetic interchain couplings were studied in Ref. where the former model was found to agree better with the experimental neutron scattering data. From a Raman scattering study of single crystals, the spin-phonon interaction was suggested to be responsible for the formation of the two-magnon triplet bound states identified in the neutron scattering experiments, rather than arising from frustrating interchain interactions.
Electron spin resonance (ESR) intensity versus temperature $`I(T)`$ data obtained for a single crystal were interpreted in terms of an isolated AF spin dimer model, yielding a spin gap of 67 K, corresponding to the higher-energy gap seen in the inelastic neutron scattering data. In this ESR study, transitions between the Zeeman levels within the one-magnon triplet band were reportedly not observed. These $`I(T)`$ data were subsequently re-interpreted within the above 2D coupled alternating-exchange chain model as arising from transitions between the Zeeman levels within the one-magnon band and good agreement between the theoretical prediction for $`I(T)`$ and the experimental $`I(T)`$ data was found.
On the other hand, recent <sup>31</sup>P and <sup>51</sup>V NMR and magnetization versus applied magnetic field $`M(H)`$ measurements at high fields and low temperatures have indicated that there are two magnetically distinct types of alternating-exchange V chains in AP-$`(\mathrm{VO})_2\mathrm{P}_2\mathrm{O}_7`$, interpenetrating with each other, each with its own spin gap. The two spin gaps inferred for the two types of chains, 35(2) K from $`{}_{}{}^{31}K(T)`$ measurements and 68(2) K from <sup>51</sup>V $`K(T)`$ measurements, and 33(1) K and 62(3) K from $`M(H)`$ measurements at 1.3 K, agree well with the above two spin gaps found from the neutron scattering measurements, respectively, thus providing an alternate explanation for the larger of the two spin gaps. Using additional information from the neutron scattering one-magnon dispersion relation measurements, the two alternating exchange constants in each chain have been estimated and are listed in Table I. This model is not supported by the Raman scattering results. However, recent unpublished inelastic neutron scattering measurements on a large single crystal at $`T=2`$ K in zero and high magnetic fields show that the two types of magnetic excitations found earlier split in a magnetic field according to expectation for two triplet bands, results which constitute independent evidence for the validity of the two-chain model for AP-$`(\mathrm{VO})_2\mathrm{P}_2\mathrm{O}_7`$. In addition, two previously undetected magnetic excitations termed “shadow bands” were found which are thought to arise from the staggered alignment of successive V dimers along the V alternating-exchange chains within the structure (see the top panel of Fig. 2).
At present there is thus no universal agreement about a Hamiltonian which can self-consistently explain the various experimental measurements probing the magnetism in AP-$`(\mathrm{VO})_2\mathrm{P}_2\mathrm{O}_7`$, although all of the models considered recently contain $`S=1/2`$ AF alternating-exchange Heisenberg chains as an essential element and a consen-
sus is emerging that the two-chain model can explain many of the observed properties.
### B HP-$`\mathbf{(}\mathrm{𝐕𝐎}\mathbf{)}_\mathrm{𝟐}𝐏_\mathrm{𝟐}𝐎_\mathrm{𝟕}`$
The high-pressure phase HP-$`(\mathrm{VO})_2\mathrm{P}_2\mathrm{O}_7`$ was recently synthesized by heating polycrystalline AP-$`(\mathrm{VO})_2\mathrm{P}_2\mathrm{O}_7`$ for 1/2 h at 700 C under a pressure of 2 GPa. As shown in Fig. 3 and by comparison with Fig. 2, HP-$`(\mathrm{VO})_2\mathrm{P}_2\mathrm{O}_7`$ has a simpler structure than AP-$`(\mathrm{VO})_2\mathrm{P}_2\mathrm{O}_7`$. In HP-$`(\mathrm{VO})_2\mathrm{P}_2\mathrm{O}_7`$, all V atoms are crystallographically equivalent but the same basic structure as in the ambient-pressure phase was found. The similarities between the two structures suggest that HP-$`(\mathrm{VO})_2\mathrm{P}_2\mathrm{O}_7`$ also contains $`S=1/2`$ AF alternating-exchange chains, but of a single type. Indeed, high-field $`M(H)`$ measurements at 1.3 K revealed a single spin gap of $`23`$ K, consistent with this hypothesis. Modeling of $`\chi (T)`$ data below 30 K was carried out using the low-$`T`$ approximation in Eq. (24) below for the spin susceptibility of a 1D spin system with a spin gap, yielding a similar spin gap of 27 K. The $`\chi (T)`$ data above 30 K were analyzed using the spin susceptibility of the $`S=1/2`$ alternating-exchange Heisenberg chain model, yielding the exchange constants listed in Table I and the same spin gap $`\mathrm{\Delta }/k_\mathrm{B}=27`$ K. These estimates of the gap value are similar to the one-magnon gap of $`36`$ K in AP-$`(\mathrm{VO})_2\mathrm{P}_2\mathrm{O}_7`$ found from the inelastic neutron scattering and other measurements discussed above.
Equation (24) has also been previously used to fit $`\chi (T)`$ data for other $`S=1/2`$ 1D compounds, but to our knowledge all such studies, with the exception of a study of the $`\mathrm{Cu}^{+2}`$ $`S=1/2`$ two-leg Heisenberg spin ladders in $`\mathrm{SrCu}_2\mathrm{O}_3`$, have assumed that $`A`$ and the spin gap $`\mathrm{\Delta }`$ are independently adjustable parameters when fitting experimental $`\chi (T)`$ data. We discuss in Sec. II below that $`A`$ is uniquely related to $`\mathrm{\Delta }`$ and is not an independently adjustable parameter for a given type of 1D spin lattice. On the other hand, if one assumes that the spin lattice in a material has a spin gap but the type of 1D spin lattice is unknown, or if the fit is not carried out in the low-temperature limit, then $`A`$ and $`\mathrm{\Delta }`$ would have to be treated as independently adjustable parameters. In the present work we evaluate the accuracy of using Eq. (24) to determine the spin gap of the 1D $`S=1/2`$ alternating-exchange Heisenberg chain from $`\chi (T)`$ data for HP-$`(\mathrm{VO})_2\mathrm{P}_2\mathrm{O}_7`$ by comparing the spin gap obtained using this approximation with the spin gap obtained by modeling the same data set using the accurate theoretical prediction for the spin susceptibility of the $`S=1/2`$ AF alternating-exchange Heisenberg chain model.
Most recently, single crystals of HP-$`(\mathrm{VO})_2\mathrm{P}_2\mathrm{O}_7`$ have been grown. The anisotropic spectroscopic splitting factors ($`g`$ factors) of the $`S=1/2`$ V<sup>+4</sup> ions in a single crystal were determined from ESR measurements and the anisotropic magnetic susceptibilites of a single crystal were measured. Here we perform detailed modeling of these $`\chi (T)`$ data using the $`g`$ factors determined from ESR and using the $`S=1/2`$ AF alternating-exchange Heisenberg chain model for the spin susceptibility. We determine the exchange constants within the chains and from these we obtain an estimate of the spin gap. An independent estimate of the spin gap is obtained by modeling only the low-temperature data.
### C Plan of the Paper
The plan for the remainder of this paper is as follows. In Sec. II a summary is given of the theory that we will need to carry out and discuss the modeling described above. General considerations for fitting experimental data by theoretical predictions for the spin susceptibility are discussed in Sec. III A. New and literature $`\chi (T)`$ data for AP-$`(\mathrm{VO})_2\mathrm{P}_2\mathrm{O}_7`$ are presented and fitted in Sec. III B. High-precision fits to the $`\chi (T)`$ of a high-purity powder sample are presented in Sec. III B 1. In this section, we show how the fitted exchange constants and spin gap(s) of the model of alternating-exchange chain(s) vary depending on whether a single-chain or two-chain model is used to fit the data, and on whether the $`g`$ value is fixed or allowed to vary during the fits. The influences of interchain couplings on the exchange constants and spin gaps inferred from modeling $`\chi (T)`$ for the powder sample are quantitatively determined in Sec. III B 2 using a molecular field theory for the interchain couplings, where we also compare our derived interchain couplings with the corresponding theoretical predictions of Uhrig and Normand which were obtained using a one-chain model. The $`\chi (T)`$ data for two single crystals of AP-$`(\mathrm{VO})_2\mathrm{P}_2\mathrm{O}_7`$ are analyzed using the two-chain model in Sec. III B 3. In Sec. III C we test our predicted dispersion relations for the two proposed chains by comparison with the results of inelastic neutron scattering measurements. The $`\chi (T)`$ data for HP-$`(\mathrm{VO})_2\mathrm{P}_2\mathrm{O}_7`$ are presented and modeled in Sec. III D. The most accurate and precise $`\chi (T)`$ data for this phase were obtained for a sample of crushed crystals. These data are analyzed using the one-chain model in Sec. III D 1. The influences of interchain couplings on the derived exchange constants and spin gap of this sample are considered in Sec. III D 2. Our evaluation of the accuracy of the spin gap obtained using a theoretical low-$`T`$ approximation to the spin susceptibility of a 1D $`S=1/2`$ spin system with a spin gap, as previously used to analyze powder $`\chi (T)`$ data for HP-$`(\mathrm{VO})_2\mathrm{P}_2\mathrm{O}_7`$, is given in Sec. III D 3. The anisotropic $`\chi (T)`$ data for a single crystal of HP-$`(\mathrm{VO})_2\mathrm{P}_2\mathrm{O}_7`$ are modeled in Sec. III D 4. The powder average of these single crystal $`\chi (T)`$ data are modeled in Sec. III D 5 by two low-$`T`$ approximations for the spin susceptibility to obtain an independent estimate of the spin gap. A summary of our modeling results and our conclusions are given in Sec. IV.
## II Theory
The Hamiltonian for the $`S=1/2`$ alternating-exchange Heisenberg chain is written in three equivalent ways as
$``$ $`=`$ $`{\displaystyle \underset{i}{}}J_1𝑺_{2i1}𝑺_{2i}+J_2𝑺_{2i}𝑺_{2i+1}`$ (2)
$`=`$ $`{\displaystyle \underset{i}{}}J_1𝑺_{2i1}𝑺_{2i}+\alpha J_1𝑺_{2i}𝑺_{2i+1}`$ (3)
$`=`$ $`{\displaystyle \underset{i}{}}J(1+\delta )𝑺_{2i1}𝑺_{2i}+J(1\delta )𝑺_{2i}𝑺_{2i+1},`$ (4)
where
$`J_1`$ $`=`$ $`J(1+\delta )={\displaystyle \frac{2J}{1+\alpha }},`$ (6)
$`\alpha `$ $`=`$ $`{\displaystyle \frac{J_2}{J_1}}={\displaystyle \frac{1\delta }{1+\delta }},`$ (8)
$`\delta `$ $`=`$ $`{\displaystyle \frac{J_1}{J}}1={\displaystyle \frac{J_1J_2}{2J}}={\displaystyle \frac{1\alpha }{1+\alpha }},`$ (10)
$`J`$ $`=`$ $`{\displaystyle \frac{J_1+J_2}{2}}=J_1{\displaystyle \frac{1+\alpha }{2}},`$ (12)
with AF couplings $`J_1J_20`$, $`0(\alpha ,\delta )1`$. The uniform undimerized chain corresponds to $`\alpha =1,\delta =0`$ and $`J_1=J_2=J`$, whereas the isolated dimer has $`\alpha =0,\delta =1,J_2=0`$ and intradimer exchange interaction $`J_1`$. The form of the Hamiltonian in Eq. (4), in which the appropriate variables are $`\delta `$ and the average exchange constant along the chain $`J`$ instead of $`\alpha `$ and the maximum exchange constant $`J_1`$ as in Eq. (3), is often used for compounds in which the spin dimerization is weak and/or for systems showing a second-order spin dimerization transition with decreasing temperature such as occurs in spin-Peierls systems.
The spin gap $`\mathrm{\Delta }`$ for magnetic spin excitations from the $`S=0`$ ground state to the lowest-lying $`S=1`$ triplet excited states for the alternating-exchange chain is uniquely related to the alternation parameter $`\alpha `$ and the larger exchange constant $`J_1`$ (or equivalently to $`\delta `$ and $`J`$). The ratio $`\mathrm{\Delta }(\alpha )/J_1`$ for the $`S=1/2`$ AF alternating-exchange Heisenberg chain was computed to high ($`1\%`$) accuracy for $`0\alpha 0.9`$, in $`\alpha `$ increments of 0.1, using multiprecision methods by Barnes, Riera and Tennant. They found that their calculations could be parametrized very well by the simple expression
$$\frac{\mathrm{\Delta }(\alpha )}{J_1}(1\alpha )^{3/4}(1+\alpha )^{1/4},$$
(14)
or equivalently
$$\frac{\mathrm{\Delta }(\delta )}{J}2\delta ^{3/4}.$$
(15)
An expression for $`\mathrm{\Delta }(\delta )/J`$ which is thought to be more accurate ($`\pm 0.0002`$) over the entire range $`0\delta 1`$, obtained by fitting numerical $`\mathrm{\Delta }(\delta )/J`$ data by a generalized form of Eq. (15), is
$$\frac{\mathrm{\Delta }(\delta )}{J}=2\delta ^{y(\delta )},$$
(17)
where the exponent $`y(\delta )`$ is given by
$`y(\delta )=y(1)`$ $`+`$ $`n_1\mathrm{tanh}\left[{\displaystyle \frac{\mathrm{ln}\delta }{m_1}}\mathrm{ln}\left({\displaystyle \frac{\mathrm{ln}\delta }{m_2}}\right)\right]`$ (18)
$`+`$ $`n_2\mathrm{tanh}^2\left[{\displaystyle \frac{\mathrm{ln}\delta }{m_1}}\mathrm{ln}\left({\displaystyle \frac{\mathrm{ln}\delta }{m_2}}\right)\right]`$ (19)
with parameters
$`y(1)=0.74922,n_1=0.00776,n_2=0.00685,`$
$$m_1=3.3297,m_2=2.2114.$$
(20)
The expression for $`\mathrm{\Delta }(\delta )/J`$ in Eqs. (II) can be transformed into an expression for $`\mathrm{\Delta }(\alpha )/J_1`$ using the conversion expressions (II).
For notational convenience, we define the reduced spin susceptibility $`\chi ^{}`$, reduced temperature $`t`$ and reduced spin gap $`\mathrm{\Delta }^{}`$ as
$$\chi ^{}\frac{\chi ^{\mathrm{spin}}J_1}{Ng^2\mu _\mathrm{B}^2},t\frac{k_\mathrm{B}T}{J_1},\mathrm{\Delta }^{}=\frac{\mathrm{\Delta }}{J_1},$$
(21)
where $`\chi ^{\mathrm{spin}}`$ is the spin susceptibility, $`N`$ is the number of spins, $`k_\mathrm{B}`$ is Boltzmann’s constant and $`\chi ^{}`$ depends on both $`t`$ and the alternating exchange parameter $`\alpha J_2/J_1`$.
An accurate but unwieldy two-dimensional function $`\chi ^{}(t,\alpha )`$ for the $`S=1/2`$ AF alternating-exchange Heisenberg chain has been derived for the entire range $`0\alpha 1`$ of the alternation parameter by a global fit to numerical quantum Monte Carlo (QMC) simulations and transfer-matrix density-matrix renormalization group (TMRG) and Bethe ansatz $`\chi ^{}(t,\alpha )`$ calculations, which we will not reproduce here but will explicitly use to model $`\chi (T)`$ data for both AP- and HP-$`(\mathrm{VO})_2\mathrm{P}_2\mathrm{O}_7`$. The absolute accuracy of this function for $`0\alpha 1`$ and $`0.01t`$ is estimated to be $`2\times 10^4`$, which corresponds to $`0.1`$ % of the susceptibility at the broad maximum. For practical purposes of fitting experimental $`\chi (T)`$ data, this function can be considered to be exact for $`t0.01`$.
Troyer, Tsunetsugu, and Würtz have derived a general expression for the low-$`T`$ limit of $`\chi ^{}(t)`$ for a one-dimensional spin system with a spin gap, assuming that (i) the one-magnon dispersion relation is nondegenerate (apart from the Zeeman degeneracy), (ii) the lowest magnetic excited states are one-magnon $`S=1`$ triplet excitations, (iii) $`k_\mathrm{B}T\mathrm{\Delta }`$ and $`k_\mathrm{B}T`$ one-magnon bandwidth, and (iv) the one-magnon dispersion relation $`E(k)`$ is parabolic near the minimum according to (in the present notation)
$$\epsilon _k\frac{E(k)}{J_1}=\mathrm{\Delta }^{}+c^{}(ka)^2,$$
(22)
where $`k`$ is the wavevector in the direction of the 1D system and $`a`$ is the (average) nearest-neighbor spin-spin distance. These assumptions hold for the present case of the $`S=1/2`$ AF alternating-exchange Heisenberg chain except for the limit $`\alpha =0`$ as discussed below and for $`\alpha =1`$ for which $`\mathrm{\Delta }=0`$. With these four assumptions, $`\chi ^{}(t)`$ is given by
$$\chi ^{}(t)=\frac{A}{\sqrt{t}}\mathrm{e}^{\mathrm{\Delta }^{}/t},(t\mathrm{\Delta }^{},\mathrm{bandwidth}/J_1)$$
(24)
with
$$A=\frac{1}{2\sqrt{\pi c^{}}}.$$
(25)
The dimensionless dispersion parameter $`c^{}`$ in Eq. (22) has a unique relationship to the reduced spin gap $`\mathrm{\Delta }^{}`$ for any given 1D spin system. For example, according to the model of Ref. , this relationship for the $`S=1/2`$ AF alternating-exchange Heisenberg chain gives the value of the parameter $`A`$ in Eq. (24) as
$$A=\frac{\sqrt{\mathrm{\Delta }^{}}}{\sqrt{2\pi }f(\mathrm{\Delta }^{})},$$
(26)
which in turn yields the low-$`T`$ limit of $`\chi ^{}(t)`$ in Eq. (24) as
$$\chi ^{}(t)=\frac{1}{\sqrt{2\pi }f(\mathrm{\Delta }^{})}\left(\frac{\mathrm{\Delta }^{}}{t}\right)^{1/2}\mathrm{e}^{\mathrm{\Delta }^{}/t},(t\mathrm{\Delta }^{})$$
(28)
where the dimensionless function $`f(\mathrm{\Delta }^{})`$ is the solution of
$$\mathrm{E}\left[\frac{f^2(\mathrm{\Delta }^{})}{\mathrm{\Delta }_{}^{}{}_{}{}^{2}}\right]=\frac{\pi }{2\mathrm{\Delta }^{}}$$
(29)
and E($`x`$) is the complete elliptic integral of the second kind.
From Eq. (26), the parameter $`A`$ in Eq. (24) is not an independently adjustable parameter but instead is a unique function of the reduced spin gap $`\mathrm{\Delta }^{}`$, as was also previously inferred for two-leg spin ladders. In addition, we see from Eq. (28) that the two independent parameters of $`\chi ^{}(t,\alpha )`$ can be written for low temperatures as $`\chi ^{}(t/\mathrm{\Delta }^{},\mathrm{\Delta }^{})\chi ^{}(k_\mathrm{B}T/\mathrm{\Delta },\mathrm{\Delta }/J_1)`$. Finally, and perhaps most importantly, the high-temperature limit of the low-temperature regime in which $`\chi ^{}(t)`$ is closely approximated by Eqs. (26) is of order $`\mathrm{\Delta }^{}/10`$. At such low temperatures $`\chi (T)`$ is immeasurably small, and hence the spin gap obtained by analyzing experimental $`\chi (T)`$ data for various compounds up to temperatures corresponding to a sizable fraction of $`\mathrm{\Delta }^{}`$ using Eq. (24) leads to fitted spin gap values which may be significantly different from the actual spin gaps.
The result for the low-$`t`$ limit of $`\chi ^{}(t)`$ in Eq. (24) is not valid for the isolated dimer, which is one limit of the alternating-exchange chain with $`\alpha =J_2=0`$, because the assumption (iii) that $`k_\mathrm{B}T`$ one-magnon bandwidth required for that equation to hold is violated at all finite temperatures. For the isolated dimer, the one-magnon bandwidth is identically zero and the reduced spin gap is $`\mathrm{\Delta }^{}=1`$. The reduced spin susceptibility is given exactly by
$$\chi ^{,\mathrm{dimer}}(t)=\frac{1}{t(3+\mathrm{e}^{1/t})}$$
(31)
with the low-temperature limit
$$\chi ^{,\mathrm{dimer}}(t)=\frac{1}{t}\mathrm{e}^{1/t}(t1).$$
(32)
The temperature dependence of the prefactor to the exponential term in Eq. (32) is different from that in Eq. (24). As discussed in Ref. , a crossover occurs with decreasing temperature at low temperatures in the effective prefactor from a $`1/t`$ dependence to a $`1/\sqrt{t}`$ dependence if $`0<\alpha 1`$.
As noted above, the form for the low-$`T`$ behavior of the spin susceptibility in Eq. (24) is valid only at very low temperatures. Many years ago, Bulaevskii found that his numerical values of $`\chi ^{}(t,\alpha )`$ for the $`S=1/2`$ AF alternating-exchange Heisenberg chain, computed from an analytic theory based on the Hartree-Fock approximation, could be fitted over a relatively large temperature interval $`0.033t1/4`$ by
$$\chi ^{}(t)=\frac{A}{t}\mathrm{e}^{\mathrm{\Delta }^{}/t},$$
(33)
and he tablulated $`A`$ and $`\mathrm{\Delta }^{}`$ versus the alternation parameter $`\alpha `$. A recent extensive numerical study of his theory confirmed that the numerical predictions of his theory in the above-cited low-temperature range are fitted better by the form (33) than by (24). In addition, this study showed that although the fitting parameter $`\mathrm{\Delta }^{}(\alpha )`$ approximately follows the actual spin gap of Bulaevskii’s theory, significant discrepancies occur. Finally, a detailed numerical comparison of the prediction of Bulaevskii’s theory for $`\chi ^{}(t,\alpha )`$ with QMC simulations and TMRG calculations of this quantity showed that Bulaevskii’s theory is unsuitable for accurately extracting $`\alpha `$ values from experimental $`\chi (T)`$ data when $`\alpha 1`$.
For Heisenberg spin lattices consisting of identical spin subsystems with susceptibility $`\chi _0^{}(t)`$ which are weakly coupled to each other, the molecular field theory (MFT) prediction for the reduced spin susceptibility $`\chi ^{}(t)`$ in the paramagnetic state of the system is
$$\chi ^{}(t)=\frac{\chi _0^{}(t)}{1+\lambda \chi _0^{}(t)},$$
(35)
or equivalently
$$\frac{1}{\chi ^{}(t)}=\frac{1}{\chi _0^{}(t)}+\lambda ,$$
(36)
where the MFT coupling constant $`\lambda `$ is given by
$$\lambda =\underset{j}{\overset{}{}}\frac{J_{ij}}{J^{\mathrm{max}}},$$
(37)
the prime on the sum over $`j`$ signifies that the sum is only taken over exchange bonds $`J_{ij}`$ from a given spin $`𝑺_i`$ to spins $`𝑺_j`$ not in the same spin subsystem, and $`J^{\mathrm{max}}`$ is the exchange constant in the system with the largest magnitude. By definition, the expression for $`\chi _0^{}(t)`$ does not contain any of these $`J_{ij}`$ interactions which are external to a subsystem. Within MFT, Eqs. (II) are correct at each temperature in the paramagnetic state not only for bipartite AF spin systems, but also for any system containing subsystems coupled together by any set of FM and/or AF Heisenberg exchange interactions.
## III Modeling of Experimental $`𝝌\mathbf{(}𝑻\mathbf{)}`$ Data
### A Introduction
We fitted the $`\chi (T)`$ data per mole of V spins-1/2 in (VO)<sub>2</sub>P<sub>2</sub>O<sub>7</sub> by the general expression
$$\chi (T)=\chi _0+\frac{C_{\mathrm{imp}}}{T\theta _{\mathrm{imp}}}+\chi ^{\mathrm{spin}}(T),$$
(39)
with
$$\chi _0=\chi ^{\mathrm{core}}+\chi ^{\mathrm{VV}}$$
(40)
and
$`\chi ^{\mathrm{spin}}(T)`$ $`=`$ $`{\displaystyle \frac{N_\mathrm{A}g^2\mu _\mathrm{B}^2}{J_1}}\chi ^{}(t)`$ (41)
$`=`$ $`\left(0.3751{\displaystyle \frac{\mathrm{cm}^3\mathrm{K}}{\mathrm{mol}}}\right){\displaystyle \frac{g^2}{J_1/k_\mathrm{B}}}\chi ^{}\left({\displaystyle \frac{k_\mathrm{B}T}{J_1}}\right),`$ (42)
where $`N_\mathrm{A}`$ is Avogadro’s number, $`\mu _\mathrm{B}`$ is the Bohr magneton, $`k_\mathrm{B}`$ is Boltzmann’s constant and $`g`$ is the spectroscopic splitting factor ($`g`$ factor) appropriate to a particular direction of the applied magnetic field with respect to the crystal axes.
The first term $`\chi _0`$ in Eq. (39), according to Eq. (40), is the sum of the nearly isotropic orbital diamagnetic atomic core contribution $`\chi ^{\mathrm{core}}`$ and the anisotropic orbital paramagnetic Van Vleck contribution $`\chi ^{\mathrm{VV}}`$ which are normally nearly independent of $`T`$. Using the values $`\chi ^{\mathrm{core}}=8,12`$ and $`47\times 10^6`$ cm<sup>3</sup>/mol for V<sup>+4</sup>, O<sup>-2</sup> and $`(\mathrm{PO}_4)^3`$, respectively, we obtain
$$\chi ^{\mathrm{core}}=6.1\times 10^5\frac{\mathrm{cm}^3}{\mathrm{mol}\mathrm{V}}$$
(43)
for (VO)<sub>2</sub>P<sub>2</sub>O<sub>7</sub>. The second term in Eq. (39) is an extrinsic impurity Curie-Weiss term with impurity Curie constant $`C_{\mathrm{imp}}`$ and Weiss temperature $`\theta _{\mathrm{imp}}`$ which gives rise to a low-temperature upturn in $`\chi (T)`$ which is not predicted by theory for the third term, the intrinsic spin susceptibility $`\chi ^{\mathrm{spin}}(T)`$, and is assumed to arise from finite chain segments containing an odd number of spins, impurity phase intergrowths in the crystals, paramagnetic impurity phases and/or defects. The $`C_{\mathrm{imp}}`$ and $`\theta _{\mathrm{imp}}`$ parameters can be anisotropic if the paramagnetic impurity principal directions are fixed with respect to the crystal axes, as can occur in a single crystal of a material such as studied here in Sec. III D 4, rather than being randomly distributed.
Unless otherwise stated, we assume that the spin susceptibility $`\chi ^{\mathrm{spin}}(T)`$ in Eq. (39), written in terms of $`\chi ^{}(t,\alpha )`$ in Eq. (42), is the intrinsic spin susceptibility per mole of spins-1/2 in an AF alternating-exchange Heisenberg chain. The explicit expression for $`\chi ^{}(t,\alpha )`$ of this chain is given in Ref. . In Eq. (42), $`J_1`$ is the larger of the two ($`J_1`$ and $`J_2`$ with $`J_1>J_2>0`$) AF alternating exchange constants along the alternating-exchange chain, as denoted in Sec. II above.
One of the parameters entering the calculated spin susceptibility $`\chi ^{\mathrm{spin}}(T)`$ in Eq. (42) is the $`g`$ value of the V magnetic moments. Measurements of the anisotropic $`g`$ values of the spins in both AP-(VO)<sub>2</sub>P<sub>2</sub>O<sub>7</sub> and HP-(VO)<sub>2</sub>P<sub>2</sub>O<sub>7</sub> have been carried out using ESR measurements and the results are listed in Table II. The significant differences between the $`g`$ values of the two phases of $`(\mathrm{VO})_2\mathrm{P}_2\mathrm{O}_7`$ reflect the differences in the local bonding of the V atoms with the coordinating O atoms in the two structures. Comparison of the average $`g`$ value for HP-$`(\mathrm{VO})_2\mathrm{P}_2\mathrm{O}_7`$ with those for V in the “trellis layer” compounds $`R\mathrm{V}_2\mathrm{O}_5`$ with $`R`$ = Ca, Mg and Na, also shown in Table II, suggests that the local crystalline electric field (CEF) at the V sites in HP-$`(\mathrm{VO})_2\mathrm{P}_2\mathrm{O}_7`$ is closer to that in these compounds than to the CEF in AP-$`(\mathrm{VO})_2\mathrm{P}_2\mathrm{O}_7`$.
A measure of the goodness of a fit to experimental $`\chi (T)`$ data is the statistical $`\chi ^2`$ per degree of freedom
$$\frac{\chi ^2}{\mathrm{DOF}}\frac{1}{N_\mathrm{p}P}\underset{i=1}{\overset{N_\mathrm{p}}{}}[\chi (T_i)\chi _{\mathrm{fit}}(T_i)]^2,$$
(44)
where $`N_\mathrm{p}`$ is the number of data points in the data set and $`P`$ is the number of independent fitting parameters. This is the quantity that is minimized during our nonlinear least-squares fits to experimental $`\chi (T)`$ data. An additional measure of the quality of a fit is the relative rms deviation $`\sigma _{\mathrm{rms}}`$ of the fit from the data, given by
$$\sigma _{\mathrm{rms}}^2\frac{1}{N_\mathrm{p}}\underset{i=1}{\overset{N_\mathrm{p}}{}}\left[\frac{\chi (T_i)\chi _{\mathrm{fit}}(T_i)}{\chi (T_i)}\right]^2.$$
(45)
The fits were all carried out on a 400 MHz Macintosh G3 computer using the software Mathematica 3.0.
### B Magnetic Susceptibility of AP-(VO)$`_\mathrm{𝟐}`$P$`_\mathrm{𝟐}`$O$`_\mathrm{𝟕}`$
#### 1 Powder sample
The purpose of the present section is to test consistency with experimental $`\chi (T)`$ data of the model of Yamauchi and co-workers for AP-(VO)<sub>2</sub>P<sub>2</sub>O<sub>7</sub>, discussed in the Introduction, in which this compound is proposed to consist magnetically of equal numbers of two independent types of isolated $`S=1/2`$ AF alternating-exchange Heisenberg chains with spin gaps of about 35 and 68 K, respectively.
The $`\chi (T)`$ of a polycrystalline (“powder”) sample of AP-(VO)<sub>2</sub>P<sub>2</sub>O<sub>7</sub> of mass 172.2 mg and with a moss-green color was measured from 2 to 350 K in a magnetic field of 1 kG and the results are shown as the open circles in Fig. 4. The details of the sample preparation will be presented elsewhere. The color of the sample indicates that it is stoichiometric with a vanadium oxidation state very close to +4. The sample was a cylinder of 4 mm diameter and 7 mm length. There was no difference between field-cooled and zero-field-cooled $`\chi (T)`$ measurements. The quality of the sample, judging from the very small Curie-Weiss upturn at low temperatures due to magnetic impurities and/or defects, is better than previously reported for any powder sample and is about the same as recently reported for a high-quality single crystal as shown in Figs. 7 and 8(b) below.
We will describe in detail our modeling results for this sample to indicate how the parameters and the quality of fit change for various types of fits. Similar variations were found from modeling $`\chi (T)`$ data for two AP-(VO)<sub>2</sub>P<sub>2</sub>O<sub>7</sub> single crystals below. To make contact with previous modeling of $`\chi (T)`$ for this material, we first fitted the data by Eqs. (III A) assuming that $`\chi ^{\mathrm{spin}}(T)`$ is due to a single type of $`S=1/2`$ AF alternating-exchange Heisenberg chain, where the $`g`$ value is either fixed at the powder-averaged value $`g=1.969`$ determined from ESR measurements for AP-(VO)<sub>2</sub>P<sub>2</sub>O<sub>7</sub> as shown in Table II, yielding “Fit 1”, or allowed to vary during the fit (“Fit 2”). Throughout
the modeling in this section, we use the expression (II) to determine the spin gap from the fitted exchange constants for an alternating-exchange chain. The parameters obtained from each fit are shown in Table III, together with the statistical $`\chi ^2`$/DOF and $`\sigma _{\mathrm{rms}}`$ for each fit. The defect and/or impurity Curie constant is equivalent to the contribution of about 1.6 mol % with respect to V of spins 1/2 with $`g=2`$. The values of the alternating exchange constants $`J_1`$ and $`J_2`$ and the spin gap $`\mathrm{\Delta }`$ are respectively similar in each fit and are about the same as previously estimated from similar fits to $`\chi (T)`$ data for this compound, respectively (see Table I). Next, we fitted the $`\chi (T)`$ data using the above model of Yamauchi and co-workers for AP-(VO)<sub>2</sub>P<sub>2</sub>O<sub>7</sub>, but where we constrained the spin gaps to be 36.2 and 66.7 K as found from the inelastic neutron scattering measurements, and again either fixed $`g=1.969`$ (“Fit 3”) or allowed $`g`$ to vary during the fit (“Fit 4”). In order to enforce the constraint on the two spin gaps using the expression for $`\mathrm{\Delta }(\delta )/J`$ in Eqs. (II), it was more convenient to use $`J`$ and $`\delta `$ as the independent parameters in $`\chi ^{\mathrm{spin}}(T)`$ during our least-squares fits instead of $`J_1`$ and $`\alpha `$. The parameters obtained from the two fits are shown in Table III, together with other parameters derived from the fitted ones. As can be seen from the values of the $`\chi ^2/`$DOF and $`\sigma _{\mathrm{rms}}`$, the qualities of the two fits show dramatic improvements over those of Fits 1 and 2 where only a single type of alternating-exchange chain was assumed in the model-
ing. However, the $`g`$ value obtained from Fit 4 is somewhat larger than expected.
Finally, we fitted the same data set in Fig. 4 using the above model of Yamauchi and co-workers for AP-(VO)<sub>2</sub>P<sub>2</sub>O<sub>7</sub>, where we again either fixed $`g=1.969`$ (“Fit 5”) or allowed $`g`$ to vary during the fit (“Fit 6”), but where we did not constrain the fitting parameters $`J_1`$ and $`\alpha `$ of the two independent chains to yield the respective spin gaps found from the inelastic neutron scattering measurements. The fitted parameters are listed for each fit in Table III, together with the statistical $`\chi ^2`$/DOF and $`\sigma _{\mathrm{rms}}`$ and the derived $`\mathrm{\Delta }`$ for each fit. We checked that the identical fitted parameters are obtained for Fit 5 independent of whether the starting parameters are the fitted parameters of Fit 1, for which the exchange constants in the two chains are identical, or of Fit 3 for which they are different. Fit 6 is shown as the solid curve in Fig. 4 and the fitted $`\chi ^{\mathrm{spin}}(T)`$ is shown as the dashed curve.
The fitted and derived parameters for Fits 5 and 6 in Table III exhibit a number of important features. First, the qualities of Fits 5 and 6 to the data are far superior to those of Fits 1 and 2. Second, the values of the alternation parameters and spin gaps for the two independent isolated chains of the model did not converge to the same respective values for the two chains, but rather are clearly differentiated. Third, the fitted $`g`$ value from Fit 6 is identical within the respective errors with the powder averaged $`g`$ value in Table II determined from ESR measurements. Fourth, and perhaps most importantly, the two spin gaps derived from the respective exchange constants for the two chains are respectively nearly identical to the two spin gaps found from the single-crystal inelastic neutron scattering measurements by Garrett et al. and with the values inferred previously by Yamauchi and co-workers from high-field magnetization measurements and a subset of the NMR measurements. It seems very unlikely that the two spin gaps we deduce from this model could be so close to those determined from other independent measurements without the model being essentially correct. The exchange constants and spin gaps we derived from $`\chi (T)`$ data for the same powder sample in $`H=50`$ kG, data which are not otherwise discussed here, are identical within the respective errors to those we obtained above for $`H=1`$ kG. Finally, we will see in Sec III D 1 below that when the two-chain model is used to extract the exchange constants within the proposed single-chain high-pressure phase HP-(VO)<sub>2</sub>P<sub>2</sub>O<sub>7</sub>, essentially the same exchange constants and spin gaps are obtained for both chains of the model. This result indicates that our fitting procedure can clearly differentiate between pairs of chains that have the same or different spin gaps, respectively.
We conclude that our analysis of $`\chi (T)`$ is precisely consistent with the model of Yamauchi and co-workers for the nature of the important spin interactions in AP-(VO)<sub>2</sub>P<sub>2</sub>O<sub>7</sub>. Our values of the spin gaps of the two independent isolated alternating-exchange chains of the model are in good agreement with those determined from their high-field magnetization and NMR measurements and with the two values determined from the inelastic neutron scattering measurements, respectively.
#### 2 MFT analysis of interchain coupling
The neutron scattering measurements on single crystals showed unambiguously that interchain coupling $`J_a`$ along the $`a`$ axis of the structure, perpendicular to the alternating-exchange chains and parallel to the legs of the structural two-leg ladders (see Fig. 1), is not negligible. However the ratio $`|J_a/J_1|`$ was estimated from fits to the data to be only 2–3 %, where $`J_1`$ is the larger of the two exchange couplings along the alternating-exchange chains running along the $`c`$ axis. Another estimate can be obtained as one-half the ratio of the average total dispersion of the two presumed one-magnon bands in the $`a`$ axis direction \[16(4) K\] to the one-magnon excitation energy along the direction of the alternating-exchange chains at the zone boundary (180 K), yielding a slightly larger $`|J_a/J_1|4.5(10)`$ %. This interchain coupling along the $`a`$ axis was of course ignored in the fits to the experimental $`\chi (T)`$ in the previous section. Here we obtain an estimate of the strength of this interchain coupling $`J_a`$ from analysis of the powder $`\chi (T)`$ data presented in the previous section. In the absence of accurate calculations of $`\chi ^{}(t)`$ for this case, we will utilize the prediction of MFT given in Eqs. (II) for the influence of the interchain coupling on $`\chi (T)`$.
In order to apply Eqs. (II) to the present modeling framework in which two distinct types of alternating-exchange chains (1) and (2) are assumed to be present, one must appropriately define the “isolated subsystem” discussed in Sec. II. Here, an isolated subsystem consists of one set of the two chains (1) and (2). Thus, the reduced susceptibility of our isolated subsystem is
$`\chi _0^{}`$ $``$ $`{\displaystyle \frac{\chi J_1^{\mathrm{max}}}{Ng^2\mu _\mathrm{B}^2}}`$ (46)
$`=`$ $`{\displaystyle \frac{1}{2}}\left[{\displaystyle \frac{J_1^{\mathrm{max}}}{J_1^{(1)}}}\chi _{\mathrm{chain}}^{}(t^{(1)},\alpha ^{(1)})+{\displaystyle \frac{J_1^{\mathrm{max}}}{J_1^{(2)}}}\chi _{\mathrm{chain}}^{}(t^{(2)},\alpha ^{(2)})\right]`$ (47)
where
$$t^{(1)}\frac{k_\mathrm{B}T}{J_1^{(1)}},t^{(2)}\frac{k_\mathrm{B}T}{J_1^{(2)}},J_1^{\mathrm{max}}\mathrm{max}[J_1^{(1)},J_1^{(2)}].$$
(48)
Then the MFT coupling constant $`\lambda `$ in Eqs. (II) is, according to Eq. (37), the average interchain coupling in the $`a`$ axis direction of a spin in one of the two distinct alternating-exchange chains with all spins in the respective adjacent alternating-exchange chains. Assuming that a spin in each chain is coupled to two nearest
neighbor spins in the $`a`$ direction by exchange constant $`J_a`$, one obtains
$$\lambda =2\frac{J_a}{J_1^{\mathrm{max}}}.$$
(49)
We fitted the data in Fig. 4 by Eqs. (III A), assuming $`g=1.969`$ for all of the fits and using Eqs. (35), (42), and (48), where $`t\mathrm{min}(t^{(1)},t^{(2)})`$, to determine the spin susceptibility $`\chi ^{\mathrm{spin}}(T)`$. The resulting fitted parameters $`J_1^{(1,2)}`$ and $`\alpha ^{(1,2)}`$ for the chains (1) and (2), the rms deviation $`\sigma _{\mathrm{rms}}`$ of the fit from the data, and the spin gaps $`\mathrm{\Delta }^{(1,2)}`$ for the chains (1) and (2) derived using Eqs. (II), are plotted versus $`\lambda `$ in Fig. 5 for $`0.4\lambda 0.4`$ in $`\lambda `$ increments of 0.05, where positive (negative) values of $`\lambda `$ correspond to AF (ferromagnetic FM) coupling $`J_a`$. Not plotted in Fig. 5 are the fitted values of $`\chi _0`$, $`C_{\mathrm{imp}}`$, and $`\theta _{\mathrm{imp}}`$, which for $`\lambda =0.4`$, 0 and 0.4 are $``$7.0(1), $``$1.99(2), 3.8(2)$`\times 10^5`$ cm<sup>3</sup>/mol V, 4.65(7), 4.26(1), 3.4(1)$`\times 10^3`$ cm<sup>3</sup> K/mol V, and $``$4.0(1), $``$3.89(2), $``$3.0(2) K, respectively. From Fig. 5, a pronounced minimum (0.095 %) occurs in $`\sigma _{\mathrm{rms}}`$ at $`\lambda 0.05`$, which corresponds according to Eq. (49) to a FM $`J_a`$ with $`|J_a/J_1^{(1)}|0.025`$. This is quantitatively consistent with the above-cited estimates of this ratio based on a one-chain model with anisotropic (in spin space) spin interactions for the one-magnon dispersion relations observed by inelastic neutron scattering in Ref. . The fit to the data for $`\lambda =0.05`$ is shown as the solid curve in Fig. 6, where $`\chi _0=2.50(2)\times 10^5`$ cm<sup>3</sup>/mol V, $`C_{\mathrm{imp}}=0.00428(1)`$ cm<sup>3</sup> K/mol V, $`\theta _{\mathrm{imp}}=3.87(2)`$ K, $`J_1^{(1)}/k_\mathrm{B}=139.5(2)`$ K, $`\alpha ^{(1)}=0.8597(4)`$, $`J_1^{(2)}/k_\mathrm{B}=123.7(2)`$ K, $`\alpha ^{(2)}=0.6319(6)`$, $`\mathrm{\Delta }^{(1)}/k_\mathrm{B}=37.5(2)`$ K, and $`\mathrm{\Delta }^{(2)}/k_\mathrm{B}=66.1(2)`$ K. In deriving the spin gap for each chain using Eqs. (II) we have implicitly assumed that the spin gap is unaffected by the interchain couplings.
A somewhat more precise estimate of $`\lambda `$ is obtained by allowing this parameter to vary during the fit. The fit parameters and derived spin gaps of the two chains (1) and (2) for the best fit are
$`\chi _0=2.37(4)\times 10^5{\displaystyle \frac{\mathrm{cm}^3}{\mathrm{mol}\mathrm{V}}},C_{\mathrm{imp}}=0.00427(1){\displaystyle \frac{\mathrm{cm}^3\mathrm{K}}{\mathrm{mol}\mathrm{V}}},`$
$`\theta _{\mathrm{imp}}=3.87(2)\mathrm{K},\lambda =0.037(4),`$
$$\frac{J_1^{(1)}}{k_\mathrm{B}}=137.1(7)\mathrm{K},\alpha ^{(1)}=0.855(2),$$
(51)
$`{\displaystyle \frac{J_1^{(2)}}{k_\mathrm{B}}}=124.8(4)\mathrm{K},\alpha ^{(2)}=0.634(1),`$
$`{\displaystyle \frac{\chi ^2}{\mathrm{DOF}}}=0.89\left(10^6{\displaystyle \frac{\mathrm{cm}^3}{\mathrm{mol}\mathrm{V}}}\right)^2,\sigma _{\mathrm{rms}}=0.099\%,`$
$`{\displaystyle \frac{\mathrm{\Delta }^{(1)}}{k_\mathrm{B}}}=37.8(5)\mathrm{K},{\displaystyle \frac{\mathrm{\Delta }^{(2)}}{k_\mathrm{B}}}=66.4(3)\mathrm{K}.`$
The spin gaps are similar to and may be compared with those for Fit 5 in Table III for the same data, in which $`g=1.969`$ was also assumed but where $`\lambda =0`$. From Eq. (49) which assumes a nearest neighbor interchain coordination number of 2, we obtain the average interchain coupling strength
$$\frac{J_a}{k_\mathrm{B}}=\frac{\lambda J_1^{(1)}}{2k_\mathrm{B}}2.5\mathrm{K}.$$
(52)
As noted in the Introduction, Uhrig and Normand proposed a model for AP-(VO)<sub>2</sub>P<sub>2</sub>O<sub>7</sub> in which only one type of alternating-exchange chain occurs and in which the (AF) interchain couplings are given by
$$\frac{J_1}{k_\mathrm{B}}=124\mathrm{K},\alpha =0.793,\frac{J_a}{J_1}=0.203,\frac{J_\times }{J_1}=0.255,$$
(54)
where these parameters were obtained by fitting the one-magnon inelastic neutron scattering dispersion relation data for the lower band and using the observed Weiss temperature obtained by fitting experimental $`\chi (T)`$ data at high temperatures by a Curie-Weiss law (as predicted by MFT). Assuming that the interchain couplings do not affect the spin gap, Eqs. (II) yield
$$\frac{\mathrm{\Delta }}{k_\mathrm{B}}=44.0\mathrm{K}.$$
(55)
This is about 16 % larger than our estimate for $`\mathrm{\Delta }^{(1)}`$ in Eqs. (51). Of course, since their exchange constants were determined by fitting their theory to the experimental neutron scattering data, their spin gap is the observed value ($`36`$ K) and not that in Eq. (55). The discrepancy arises because they find that the spin gap does depend on the interchain couplings, as further discussed in Sec. III C below.
Since the interchain spin coordination number for each of the interchain couplings is 2, the value of the MFT interchain coupling constant predicted by Eq. (37) is
$$\lambda =2\left(\frac{J_a}{J_1}+\frac{J_\times }{J_1}\right).$$
(56)
Inserting the parameters of Uhrig and Normand in Eq. (54) into this equation yields $`\lambda =0.916`$. Our fit parameters and their variations with $`\lambda `$ in Fig. 5 argue against this very large AF value of $`\lambda `$. To further illustrate the discrepancy within MFT between this one-chain theory and the experimental $`\chi (T)`$ data, shown as the dotted curve in Fig. 6 is the predicted $`\chi (T)`$ using $`g=1.969`$, the $`\chi _0`$, $`C_{\mathrm{imp}}`$, and $`\theta _{\mathrm{imp}}`$ values obtained for $`\lambda =0.05`$ in Sec. III B 2, and the MFT prediction for the spin susceptibility in Eqs. (II), where $`\chi _0^{}(t)`$ is that of the isolated alternating-exchange chain for which the exchange constants estimated by Uhrig and Normand in Eq. (54) were used. The relative deviation of the prediction from the data is $`\sigma _{\mathrm{rms}}=5.6`$ %, which is about 60 times larger than the $`\sigma _{\mathrm{rms}}`$ obtained using the two-chain model for $`\lambda =0.05`$. The agreement of both theoretical predictions with the data at high temperatures is expected and in fact is required for either model, since the MFT is most accurate at high temperatures where it yields the Curie-Weiss law. The significant differences between the predictions of the two models only become apparent at the lower temperatures.
In summary, our high-precision fits to the $`\chi (T)`$ data using the model of two independent chains, in which nearest neighbor chains along the $`a`$ axis are coupled using MFT, indicate that the average interchain coupling is weakly ferromagnetic, in agreement with the analysis of neutron scattering data by Garrett et al. using a one-chain model and in disagreement with the one-chain model of Uhrig and Normand with strong AF interchain couplings.
#### 3 Single crystals
In this section we analyze the anisotropic $`\chi (T)`$ data for two single crystals of AP-(VO)<sub>2</sub>P<sub>2</sub>O<sub>7</sub>. The data for a small dark green crystal of mass 2.0(1) mg (“crystal 1”), measured in a magnetic field of 5 T using a Quantum Design SQUID magnetometer, have not been reported previously. There was no discernable difference between field-cooled and zero-field-cooled $`\chi (T)`$ data for this crystal. The data for crystal 2 were reported by Prokofiev et al. in Fig. 4(a) of Ref. and were measured in a magnetic field of 2 T. An overview of the anisotropic $`\chi (T)`$ data for the two crystals is shown in Fig. 7, where the $`\chi (T)`$ of the powder sample in Fig. 4 is shown for comparison as the dashed curve. The data for the two crystals are in agreement on a coarse scale. The powder averages of the data for both crystals lie above the data for the powder sample for $`T25`$ K, although in the case
of crystal 1 this difference is not significant since the uncertainty in the crystal mass is about 5 %. Crystal 1 shows a larger Curie-Weiss-type paramagnetic defect and/or impurity contribution at low temperatures than crystal 2. Quantitative differences are seen between the anisotropic $`\chi (T)`$ for the two crystals. In particular, above about 30 K the $`\chi _a(T)`$ of crystal 1 agrees with $`\chi _{b,c}(T)`$ of crystal 2 and $`\chi _{b,c}(T)`$ of crystal 1 agrees with $`\chi _a(T)`$ of crystal 2. These qualitative anisotropy differences cannot arise from inaccuracy in, e.g., the crystal masses, which would only affect the respective ordinate scale.
The $`a`$-, $`b`$\- and $`c`$-axis $`\chi (T)`$ data for each crystal in Fig. 7 were fitted simultaneously using Eqs. (III A) by writing $`\chi `$ as a diagonal tensor. We assumed the two-chain model for the spin susceptibility of each crystal where the values of $`J_1`$ and $`\alpha `$ for each chain are the same for all three crystal directions, the anisotropic $`g`$ values are the same for both chains, and allowed $`\chi _0`$, $`C_{\mathrm{imp}}`$ and $`\theta _{\mathrm{imp}}`$ to be different for each chain and for the three field directions, for a total of 16 fitting parameters. The 4D fits obtained for crystals 1 and 2 are shown as the sets of three solid curves in Figs. 8(a) and 8(b), respectively, and the fitted
parameters for both crystals are given in Table IV where the goodnesses of fit for the two crystals are also listed. The anisotropic spin susceptibilities $`\chi _\alpha ^{\mathrm{spin}}(T)`$ were derived using Eq. (39), i.e., by subtracting the respective $`\chi _0`$ and defect and/or impurity Curie-Weiss terms from the $`\chi (T)`$ fit function, and are plotted versus temperature for crystals 1 and 2 as the two sets of three dashed curves in Figs. 8(a) and 8(b), respectively. From Eqs. (III A), the only source of anisotropy in $`\chi ^{\mathrm{spin}}(T)`$ is the anisotropy in the $`g`$ factor. Also listed in Table IV are the spin gaps computed using Eqs. (II) for the two distinct alternating-exchange spin chains in each crystal.
Several features of the data in Table IV are of note. First, as qualitatively expected from Fig. 7, the concentration of paramagnetic defects and/or impurities in crystal 1 is about a factor of two larger than in crystal 2. Second, the spin gaps of the two chains in each of the crystals 1 and 2 are consistent within the error bars with each other and with those found in the above section for the high-quality powder sample of AP-(VO)<sub>2</sub>P<sub>2</sub>O<sub>7</sub> as listed in Table III. The large error bars on the fitted parameters for crystal 2 arise in large part because the resolution in $`\chi `$ for the data above 20 K is only $`1\times 10^5`$ cm<sup>3</sup>/mol V, which corresponds to a relative resolution of, e.g., 1 % at 20 K and 0.5 % at 70 K. Third, the fitted $`g_\alpha `$ values are similar to, but differ in detail from, the corresponding ESR values for AP-(VO)<sub>2</sub>P<sub>2</sub>O<sub>7</sub> in Table II. These discrepancies between the respective $`g_\alpha `$ values may originate at least in part from the large Curie-Weiss contribution in crystal 1 and the low resolution of the $`\chi `$ data for crystal 2.
### C Dispersion Relations of AP-(VO)$`_\mathrm{𝟐}`$P$`_\mathrm{𝟐}`$O$`_\mathrm{𝟕}`$
The one-magnon dispersion relation $`E(k_c)/J_1`$ in the chain direction for the isolated $`S=1/2`$ AF alternating-exchange Heisenberg chain, with and without a frustrating second-neighbor coupling, was recently calculated to tenth order in $`\alpha `$ by Knetter and Uhrig. Thus it is possible to make a direct and accurate comparison of the dispersion relations predicted for the two proposed alternating-exchange chains on the basis of our exchange constants in AP-(VO)<sub>2</sub>P<sub>2</sub>O<sub>7</sub> with those determined by Garrett et al. using inelastic neutron scattering (INS) measurements. For this comparison, we first use the exchange constants for the two chains from Fit 6 in Table III determined from our fit to the $`\chi (T)`$ data for the high-purity powder sample using the isolated chain model. The predicted dispersion relations for the two chains are shown as the two solid curves in Fig. 9 where the experimental INS data $`()`$ are also plotted. Also shown as dashed curves are the dispersion relations predicted for the two chains using the intrachain exchange constants in Eq. (51) found from our fit to the same $`\chi (T)`$ data by the same two-chain model but where the chains are coupled along the $`a`$ axis using MFT. The range of our
prediction for the dispersion relation of each chain is thus approximately given by the region between the respective pair of solid and dashed curves.
Our predicted dispersion relations for the two chains agree well with the experimental inelastic neutron scattering data in Fig. 9 except near the zone boundary at $`k_c=\pi /c`$. Our prediction is that two peaks in the scattered neutron intensity versus energy at this wavevector should be seen with an energy separation of about 20–40 K ($``$2–3 meV), contrary to the single data point at this wavevector in Fig. 9. However, the error bar on the data point in Fig. 9 at $`k_c=\pi /c`$ is only the accuracy of determining the position of the centroid of the neutron scattering peak and is not a direct measure of the width of the peak. After most of the fits to the $`\chi (T)`$ data described in this paper and the determinations of the exchange constants were completed, we learned that recent unpublished INS measurements on a large crystal of AP-(VO)<sub>2</sub>P<sub>2</sub>O<sub>7</sub> indeed show two peaks at this wavevector with an energy splitting of about 2 meV, which partially confirms our prediction. We also predict from Fig. 9 that the dispersion relations of the two chains should cross within the intermediate wavevector regime. We are not aware of experimental INS data that address this aspect of the dispersion relations.
We note that the two sets of exchange constants in Table IV for the two chains in crystal 1 also predict an energy splitting of the neutron peak at $`k_c=\pi /c`$, but in this case the predicted dispersion relations of the two chains do not cross. We consider the experimental data and modeling for the high-purity powder sample to be more reliable than for this crystal, for reasons discussed in Sec. III B 3 above, and hence expect that the dispersion relations of the two chains will ultimately be observed to cross.
A quantitative estimate of the interchain couplings $`J_a`$ and $`J_\times `$ can be obtained by fitting the observed dispersion relation for each of the two chains in the $`a`$ direction, shown in Fig. 10, by the prediction of Uhrig and Normand
$`E(k_a,k_c`$ $`=`$ $`0)=\mathrm{\Delta }(J_1,\alpha ,\mu _+,\mu _{})`$ (58)
$`+`$ $`J_1{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}a_n(\alpha ,\mu _+,\mu _{})[\mathrm{cos}(nk_aa)1]`$ (59)
where $`\mu _\pm (J_a\pm J_\times )/J_1`$. To third order in $`\alpha `$, $`\mu _{}`$, and/or $`\mu _+`$ their general dispersion relation yields
$`a_1={\displaystyle \frac{\mu _{}}{4}}\left[4+\alpha (2+\mu _+)\alpha ^2\mu _{}^2\right],`$
$$a_2=\frac{\mu _{}^2}{8}(2+3\alpha +2\mu _+),a_3=\frac{\mu _{}^3}{8}.$$
(60)
We found that the data in Fig. 10 could not be fitted by Eqs. (10) assuming $`J_\times =0`$ and $`J_a/J_10.019`$ as inferred from $`\lambda 0.037`$ in Eqs. (5). However, we can still retain this experimentally determined value of $`\lambda `$ by allowing $`J_\times `$ to be nonzero. In particular, from Eq. (56) one obtains $`\mu _+=\lambda /2`$. Therefore we set $`\mu _+=0.019`$ and used the experimentally determined values of $`J_1`$ and $`\alpha `$ for each chain in Eq. (51) in the fit to the respective dispersion relation. This left $`\mathrm{\Delta }`$ and $`\mu _{}`$ as the only adjustable parameters. We did not use the experimentally determined $`\mathrm{\Delta }`$ values in Eqs. (5) because the total width of the experimental dispersion relation for each chain is relatively small and the small difference between the experimental $`\mathrm{\Delta }`$ value and the neutron scattering result would give a significant systematic shift to the fit (see also below). The fits are shown as the solid curves in Fig. 10 for chains (1) and (2), respectively, for which the parameters are
$`{\displaystyle \frac{J_a^{(1)}}{J_1}}=0.035,{\displaystyle \frac{J_\times ^{(1)}}{J_1}}=+0.016,{\displaystyle \frac{\mathrm{\Delta }^{(1)}}{k_\mathrm{B}}}=36.2\mathrm{K},`$
$$\frac{J_a^{(2)}}{J_1}=0.032,\frac{J_\times ^{(2)}}{J_1}=+0.013,\frac{\mathrm{\Delta }^{(2)}}{k_\mathrm{B}}=65.7\mathrm{K}.$$
(61)
Thus we find that a small but finite AF value of the interchain interaction $`J_\times `$ is necessary to fit the dispersion data perpendicular to the alternating-exchange chains if we retain the fitted $`\lambda `$ value in Eqs. (5). Since the nearest-neighbor interchain interaction $`J_a`$ is ferromagnetic, the next-nearest-neighbor AF interchain interaction $`J_\times `$ is not a geometrically frustrating interaction.
On the other hand, an equally good and nearly identical fit for each chain as shown in Fig. 10 can be obtained assuming that $`J_\times =0`$ if we relax the above condition on $`\lambda `$. In this case we still use the exchange constants in Eqs. (5) but we set $`\mu _{}=\mu _+\mu `$ in the fit to the transverse dispersion data for each chain, yielding the same respective gap values as in Eqs. (61), but where $`\mu ^{(1)}=0.050(2)`$ and $`\mu ^{(2)}=0.044(2)`$, so that the $`J_a/J_1=\mu `$ and $`\lambda =2\mu 0.10`$ values are larger in magnitude than in Eqs. (61).
The dispersion relation parallel to a chain ($`||c`$) calculated by Uhrig and Normand, in which the $`J_a`$ and $`J_\times `$ interchain couplings are included to third order in $`\alpha `$, $`\mu _{}`$, and/or $`\mu _+`$, yields
$`E(k_a`$ $`=`$ $`0,k_c)=\mathrm{\Delta }+[E_0(k_a=0,k_c)E_0(0,0)]`$ (63)
$`+`$ $`J_1{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}b_n(\alpha ,\mu _+,\mu _{})[\mathrm{cos}(nk_cc)1]`$ (64)
with
$$b_1=\frac{\mu _{}\alpha }{2}\left[1\frac{\alpha }{4}+\frac{\mu _+\mu _{}}{2}\right],b_2=\frac{3\mu _{}\alpha ^2}{16},$$
(65)
where $`E_0(0,k_c)`$ is the dispersion relation for the isolated alternating-exchange chain to tenth order in $`\alpha `$. The resulting dispersion relations calculated for the two chains using Eqs. (10), with the the intrachain exchange constants and spin gaps in Eqs. (5) and the interchain couplings given in Eqs. (61), are shown as the dotted curves in Fig. 9. These are respectively very similar to those for the MFT-coupled chain parameters already plotted as the dashed curves in Fig. 9.
An inconsistency in our fit to the transverse dispersion relation for each chain is that the calculated spin gap is smaller than the observed and fitted one in Fig. 10. The spin gap in Eq. (59) is given by the third order expansion of Uhrig and Normand as
$`\mathrm{\Delta }=J_1\{{\displaystyle \frac{\mathrm{\Delta }_0(\alpha )}{J_1}}`$ $`+`$ $`{\displaystyle \frac{\mu _{}}{4}}\left[4+\alpha (2+\mu _+)\alpha ^2\right]`$ (66)
$`+`$ $`{\displaystyle \frac{\mu _{}^2}{8}}(42\alpha +\mu _+){\displaystyle \frac{\mu _{}^3}{8}}\},`$ (67)
where $`\mathrm{\Delta }_0(\alpha )/J_1`$ is the spin gap in Eqs. (II) or (II) in the absence of interchain coupling. For the exchange constants in Eqs. (51) and (61) used to fit the transverse dispersion relations in Fig. 10 for chains (1) and (2), from Eq. (67) we obtain $`\mathrm{\Delta }^{(1)}/k_\mathrm{B}=29.2`$ K and $`\mathrm{\Delta }^{(2)}/k_\mathrm{B}=59.8`$ K, respectively. These spin gaps are each significantly smaller than the fitted ones in Eqs. (61), respectively. These discrepancies arise because the interchain couplings change the spin gap, contrary to our implicit assumption when we fitted the experimental $`\chi (T)`$ data using the MFT-coupled chain model, so the $`J_1`$ and $`\alpha `$ intrachain exchange parameters for each chain derived from the MFT fit to these data must be considered in the present model to be effective values. The degree to which the effective exchange constants differ from the actual ones is difficult to evaluate. The combined analysis we have done of the susceptibility and dispersion relations is as rigorous as can be done without having in hand an accurate theoretical expression for the spin susceptibility of the system which includes the influence of interchain couplings and concommitant changes in the two spin gaps.
### D Magnetic Susceptibility of HP-(VO)$`_\mathrm{𝟐}`$P$`_\mathrm{𝟐}`$O$`_\mathrm{𝟕}`$
#### 1 Crushed crystals
We begin our modeling of $`\chi (T)`$ for HP-(VO)<sub>2</sub>P<sub>2</sub>O<sub>7</sub> using data shown in Fig. 11 which we obtained for a 72.2 mg sample of crushed green transparent single crystals. These data are expected to be more accurate and yield more reliable values of the exchange constants and spin gap than the fits to the data for a powder and for a very small single crystal discussed in the following two sections, respectively. The data were modeled using Eqs. (III A) in which $`\chi ^{}(t)`$ is the theoretical reduced susceptibility for the $`S=1/2`$ AF alternating-exchange Heisenberg chain as proposed by Azuma et al. The fit
is shown as the solid curve in Fig. 11 and the fitted $`\chi ^{\mathrm{spin}}(T)`$ is shown as the dotted curve. The fitted $`g`$, $`J_1`$, and $`\alpha `$ values are listed in the figure, along with the spin gap computed using Eqs. (II). The other parameters of the fit are shown in Table V. The fitted $`g`$ value is very close to the powder average value 1.958 in Table II obtained from ESR measurements. The impurity Curie constant is equivalent to the contribution of 1.3 mol % with respect to V of spins 1/2 with $`g=2`$.
We used the data set in Fig. 11 to estimate typical nonstatistical errors that may arise when using the two-chain model to fit the $`\chi (T)`$ data for AP-(VO)<sub>2</sub>P<sub>2</sub>O<sub>7</sub> in the above two sections. In addition, if the present one-chain model for HP-(VO)<sub>2</sub>P<sub>2</sub>O<sub>7</sub> is appropriate, then a fit by the two-chain model should yield very similar exchange constants and spin gaps for the two chains of the model, which ideally would be respectively identical for the two chains. The parameters of the two-chain fit are compared with those of the above single-chain fit in Table V. We see that the fitted parameters of the two chains using the two-chain model are the same within the limits of error with each other and with the parameters of the single-chain model, respectively. This result indirectly confirms that the large differences between the exchange constants and spin gaps found above for the two chains in AP-(VO)<sub>2</sub>P<sub>2</sub>O<sub>7</sub> are reliable.
#### 2 MFT analysis of interchain coupling
There are no data available for the strength of the interchain coupling in HP-(VO)<sub>2</sub>P<sub>2</sub>O<sub>7</sub>. We estimate this coupling in the same way as in Sec. III B 2 for AP-(VO)<sub>2</sub>P<sub>2</sub>O<sub>7</sub>. The most precise and accurate $`\chi (T)`$ data available for HP-(VO)<sub>2</sub>P<sub>2</sub>O<sub>7</sub> are those for the crushed crystal sample presented and discussed in the previous section. We determined the fitting parameters $`\chi _0`$, $`C_{\mathrm{imp}}`$, $`\theta _{\mathrm{imp}}`$, $`J_1`$, and $`\alpha `$ as a function of the MFT interchain coupling constant $`\lambda `$ over the range $`0.5\lambda 0.5`$, where the fixed $`g=1.958`$ was assumed in all these one-chain model fits. The $`\sigma _{\mathrm{rms}}`$ and the $`J_1`$ and $`\alpha `$ parameters are plotted versus $`\lambda `$ in Fig. 12, along with the spin gap $`\mathrm{\Delta }`$ determined from $`J_1`$ and $`\alpha `$ using Eqs. (II). At $`\lambda =`$0.5, 0, and 0.5, the fitted parameters $`\chi _0`$, $`C_{\mathrm{imp}}`$, and $`\theta _{\mathrm{imp}}`$ were, respectively, $``$6.1(2), $``$0.86(9), 3.7(2)$`\times 10^5`$ cm<sup>3</sup>/mol V, 5.52(9), 4.96(4), 4.53(8)$`\times 10^3`$ cm<sup>3</sup> K/mol V, and $``$2.33(9), $``$2.03(4), $``$1.81(8) K. From the top panel of Fig. 12, the $`\sigma _{\mathrm{rms}}`$ shows an approximately parabolic variation with $`\lambda `$, with a minimum at $`\lambda 0.05`$, indicating a weak ferromagnetic interchain coupling as also deduced above for AP-(VO)<sub>2</sub>P<sub>2</sub>O<sub>7</sub>.
We next allowed $`\lambda `$ to vary during the fit to determine a more precise value. The fit parameters and derived spin gap of the alternating-exchange chain for the best fit are
$`\chi _0=1.4(2)\times 10^5{\displaystyle \frac{\mathrm{cm}^3}{\mathrm{mol}\mathrm{V}}},C_{\mathrm{imp}}=0.00502(4){\displaystyle \frac{\mathrm{cm}^3\mathrm{K}}{\mathrm{mol}\mathrm{V}}},`$
$`\theta _{\mathrm{imp}}=2.06(4)\mathrm{K},\lambda =0.055(14),`$
$$\frac{J_1}{k_\mathrm{B}}=136.2(3)\mathrm{K},\alpha =0.8748(5),$$
(69)
$`{\displaystyle \frac{\chi ^2}{\mathrm{DOF}}}=0.34\left(10^5{\displaystyle \frac{\mathrm{cm}^3}{\mathrm{mol}\mathrm{V}}}\right)^2,\sigma _{\mathrm{rms}}=0.59\%,`$
$`{\displaystyle \frac{\mathrm{\Delta }}{k_\mathrm{B}}}=33.8(2)\mathrm{K}.`$
The spin gap is identical with that obtained for the one-chain fit in Table V for the same data, in which $`g=1.958`$ was also assumed but where $`\lambda =0`$, and the other parameters are also very similar, respectively. From Eq. (49), we obtain the average interchain coupling strength
$$J_a=\frac{\lambda J_1}{2k_\mathrm{B}}3.7\mathrm{K}.$$
(70)
#### 3 Powder sample: low-$`T`$ fits
The $`\chi (T)`$ of a powder sample of HP-(VO)<sub>2</sub>P<sub>2</sub>O<sub>7</sub> was previously reported by Azuma et al., shown as the open circles in Fig. 13. A fit of the data up to 30 K by Eqs. (III A), where $`\chi ^{}(t)`$ is the low-$`T`$ approximation for the spin susceptibility of a gapped 1D $`S=1/2`$ spin system in Eq. (24), yielded the spin gap $`\mathrm{\Delta }/k_\mathrm{B}=27`$ K. In this fit, the parameter $`A`$ in Eq. (24) was treated as an independently adjustable parameter. In this section we carry out a precise fit of the same data set by Eqs. (III A) using the accurately known $`\chi ^{}(t,\alpha )`$ spin susceptibility prediction for the $`S=1/2`$ AF alternating-exchange Heisenberg chain. The present fit was carried out in order to compare the fitted parameters respectively obtained from the two types of fits to the same $`\chi (T)`$ data set for the same sample.
We fitted the $`\chi (T)`$ data in Fig. 13 using Eqs. (III A), where $`\chi ^{}(t,\alpha )`$ is that for the $`S=1/2`$ AF alternating-exchange Heisenberg chain given in Ref. , and with the $`g`$ value fixed at the spherically-averaged value 1.958 determined from single-crystal anisotropic ESR measurements (see Table II). The resulting fit is shown as the solid curve in Fig. 13, and the fitted spin susceptibility is shown as the dashed curve. The parameters of the fit are
$`\chi _0=3.8(3)\times 10^5{\displaystyle \frac{\mathrm{cm}^3}{\mathrm{mol}\mathrm{V}}},`$
$`C_{\mathrm{imp}}=0.0121(2){\displaystyle \frac{\mathrm{cm}^3\mathrm{K}}{\mathrm{mol}\mathrm{V}}},\theta =2.6(1)\mathrm{K},`$
$$\frac{J_1}{k_\mathrm{B}}=141.6(4)\mathrm{K},\alpha =0.849(2).$$
(71)
The $`C_{\mathrm{imp}}`$ value in Eq. (71) is rather large, equivalent to the contribution of 3.2 mol% of spins-1/2 with respect to V and with $`g=2`$. Therefore, the fitted prefactor $`1/J_1`$ to $`\chi ^{}(t)`$ in Eq. (42) could be too small by about 3% if the magnetic species in the impurity phase is $`S=1/2`$ V<sup>+4</sup>. On the other hand, if the impurities/defects have a spin larger than 1/2, the influence on the fitted parameters could be much smaller. In order to test this possible influence, we next allowed the $`g`$ value to be an adjustable parameter in the fit, which has the effect of allowing the amount of V in the HP-(VO)<sub>2</sub>P<sub>2</sub>O<sub>7</sub> phase relative to that in the impurity phase to be variable. The parameters obtained,
$`\chi _0=0.0(12)\times 10^5{\displaystyle \frac{\mathrm{cm}^3}{\mathrm{mol}\mathrm{V}}},`$
$`C_{\mathrm{imp}}=0.0125(3){\displaystyle \frac{\mathrm{cm}^3\mathrm{K}}{\mathrm{mol}\mathrm{V}}},\theta =2.7(1)\mathrm{K},`$
$$g=2.00(1),\frac{J_1}{k_\mathrm{B}}=143.9(8)\mathrm{K},\alpha =0.854(2),$$
(72)
are very similar to those obtained in the above fit with fixed $`g`$. The fitted value of the prefactor $`g^2/(J_1/k_\mathrm{B})`$ to $`\chi ^{}(t)`$ in Eq. (42) increased, as anticipated, from 2.71 K<sup>-1</sup> to 2.77 K<sup>-1</sup>. Although the fitted $`g`$ value increased slightly from the value used in the first fit, it is
still close to 2. We conclude that the magnetic impurities and/or defects giving rise to the Curie-Weiss term have little influence on the fitted exchange constants in the HP-(VO)<sub>2</sub>P<sub>2</sub>O<sub>7</sub> phase in the sample.
Taking $`J_1/k_\mathrm{B}=142`$ K and $`\alpha =0.85`$ from Eq. (71), Eq. (14) yields the spin gap $`\mathrm{\Delta }/k_\mathrm{B}=40`$ K, about 50% larger than the above value of 27 K estimated by Azuma et al. by fitting the same data up 30 K (i.e., up to a $`T\mathrm{\Delta }/k_\mathrm{B}`$) using the low-$`T`$ approximation to the spin susceptibility in Eq. (24). Thus from the respective values of $`\mathrm{\Delta }`$, we find that the error arising from estimating the gap value by fitting low-$`T`$ $`\chi (T)`$ data using the low-$`T`$ approximation for the spin susceptibility is about 50% in this case. The temperature range over which the low-$`T`$ approximation is fitted to the experimental data is expected to influence this error (see Sec. III D 5).
#### 4 Single crystal
The $`\chi (T)`$ data reported by Saito et al. for a 0.26 mg single crystal of HP-(VO)<sub>2</sub>P<sub>2</sub>O<sub>7</sub> are plotted in Fig. 14. By comparison of these data with those for the powder sample in Fig. 13, the Curie-Weiss impurity and/or defect contribution to $`\chi (T)`$ for the crystal at low temperatures is seen to be significantly smaller than for the powder sample. From our fit below, we find that it is in fact about a factor of three smaller. This much smaller impurity contribution enhances the reliability of the fitted exchange constants and the derived spin gap obtained from modeling the data for the single crystal. At higher temperatures, the magnitude of the powder-averaged $`\chi (T)`$ for the single crystal is very similar to the $`\chi (T)`$ for the powder, as would have been expected.
The modeling of the single-crystal $`\chi (T)`$ data was carried out in a similar way as for the two crystals of AP-(VO)<sub>2</sub>P<sub>2</sub>O<sub>7</sub> in Sec. III B 3, except that here we use a one-chain model instead of a two-chain model for the spin susceptibility. All 885 data points of the $`a`$-, $`b`$-, and $`c`$-axis $`\chi (T)`$ data sets in Fig. 14 were fitted simultaneously using Eqs. (III A) by writing $`\chi `$ as a diagonal tensor and using the three fixed $`g`$ values determined for fields along the three principal axis directions from ESR measurements as given in Table II above, respectively. With eleven fitting parameters, the data to parameter ratio is 80. The four-dimensional fit obtained is shown as the set of three solid curves in Fig. 14, and the fitted parameters are given in
Table VI along with the goodnesses of fit. The spin susceptibilities for the three crystal directions are plotted versus temperature in Fig. 14 as dashed curves.
The average of the three fitted $`C_{\mathrm{imp}}`$ values is $`3.9\times 10^3\mathrm{cm}^3\mathrm{K}/\mathrm{mol}\mathrm{V}`$, which is equivalent to the contribution of 1.0 mol % with respect to V of paramagnetic species with $`S=1/2`$ and $`g=2`$. This contribution is about a factor of three smaller than that found for the powder sample in the previous section. The average of the three $`\theta _{\mathrm{imp}}`$ values is $`2.0`$ K, about the same as for the powder sample. The negative sign of $`\theta _{\mathrm{imp}}`$ may indicate AF interactions between the defect and/or impurity magnetic moments. The $`\theta _{\mathrm{imp}}`$ can also arise from single-impurity-ion CEF effects, and/or as a reflection of partial $`T`$-dependent paramagnetic saturation of the paramag-
netic impurities at low temperatures in the fixed field of the measurements. The anisotropic $`\chi ^{\mathrm{VV}}`$ values are derived from Eqs. (40) and (43) using the fitted anisotropic $`\chi _0`$ values in Table VI. The resulting three $`\chi ^{\mathrm{VV}}`$ values are listed in Table VI.
From the above global fit to the anisotropic $`\chi (T)`$ data for all three field directions, we find $`J_1/k_\mathrm{B}=131.6(1)`$ K and $`\alpha =0.8709(5)`$. The average of the two exchange constants $`J_1`$ and $`J_2`$ according to Eq. (12) is $`J/k_\mathrm{B}=123.1(1)`$ K, and the alternation parameter expressed in the form of $`\delta `$ is obtained using our $`J_1`$ and $`\alpha `$ values and Eq. (10) as $`\delta =0.0690(3)`$. Using our fitted $`\alpha `$ and $`J_1`$ parameters, Eq. (14) yields the spin gap $`\mathrm{\Delta }/k_\mathrm{B}=33.2(2)`$ K. Using the derived $`J`$ and $`\delta `$ parameters, the more accurate Eqs. (II) predict the spin gap to be $`\mathrm{\Delta }/k_\mathrm{B}=33.4(2)`$ K,
which is the same to within the error bars as the first estimate.
A summary of all the fitted and derived quantities obtained in this section from our modeling of the $`\chi (T)`$ measurements for single crystal HP-$`(\mathrm{VO})_2\mathrm{P}_2\mathrm{O}_7`$ is given in Table VI.
#### 5 Powder average of single crystal $`\chi (T)`$: low-$`T`$ fits
Additional fits to the powder-averaged single crystal $`\chi (T)`$ data were carried out at low temperatures to obtain an estimate of the spin gap which is independent of the model for the gapped spin susceptibility. The powder average was used in order to reduce the number of parameters needed to fit the data. We used the general fit expressions (III A) in which the spin susceptibility $`\chi ^{}(t)`$ is given by the low-$`T`$ approximations in Eq. (24) (“Fit 1”) or Eq. (33) (“Fit 2”), and where the prefactor parameter $`A`$ was fitted independently of the spin gap. The only difference between Fits 1 and 2 is the exponent of $`1/t`$ in the prefactor to the exponential in the expression for the low-$`T`$ $`\chi ^{}(t)`$, which is 1/2 for Fit 1 and 1 for Fit 2.
The powder-averaged single-crystal $`\chi (T)`$ data from Fig. 14 were fitted from 2 K up to a maximum temperature $`T^{\mathrm{max}}`$, and the fitted parameters obtained from Fits 1 and 2 are plotted versus $`T^{\mathrm{max}}`$ in Fig. 15 as open and filled circles, respectively. Also shown in Fig. 15 are the statistical variances ($`\chi ^2`$ per degree of freedom) obtained from both types of fits. The variances for both fits are quite similar and both have a minimum for $`T^{\mathrm{max}}25`$ K. However, the impurity Curie constant and Weiss temperature for Fit 1 in Fig. 15 are strongly and nonmonotonically dependent on $`T^{\mathrm{max}}`$ in contrast to the corresponding dependences for Fit 2. In addition, the values of the fitted $`\chi _0`$ values for Fit 1 in Fig. 15 are all strongly negative. Since $`\chi _0`$ cannot be more negative than $`\chi ^{\mathrm{core}}`$ as estimated above in Eq. (43), because the $`\chi ^{\mathrm{VV}}`$ in Eq. (40) is necessarily positive, the Fit 1 fits for all the fitted $`T^{\mathrm{max}}`$ values are unphysical and hence the other parameters obtained using Fit 1 are most likely also highly inaccurate.
Shown in Fig. 16 are the two optimum fits obtained for $`T^{\mathrm{max}}=25`$ K for Fits 1 and 2, respectively, along with the respective fitted spin susceptibilities $`\chi ^{\mathrm{spin}}(T)`$. The $`\chi ^{\mathrm{spin}}(T)`$ for the optimum Fit 1 is highly unlikely, as are the fit parameters as just noted. On the other hand, $`\chi ^{\mathrm{spin}}(T)`$ and the fit parameters for the optimum Fit 2 are reasonable. The values of the fitted parameters for the optimum Fit 2 with $`T^{\mathrm{max}}=25`$ K in Fig. 15 are
$`\chi _0=0(4)\times 10^5{\displaystyle \frac{\mathrm{cm}^3}{\mathrm{mol}\mathrm{V}}},`$
$`C_{\mathrm{imp}}=0.0038(4){\displaystyle \frac{\mathrm{cm}^3\mathrm{K}}{\mathrm{mol}\mathrm{V}}},\theta =1.8(3)\mathrm{K},`$
$$A=0.118(2)\frac{\mathrm{cm}^3\mathrm{K}}{\mathrm{mol}\mathrm{V}},\frac{\mathrm{\Delta }}{k_\mathrm{B}}=33.8(2)\mathrm{K}.$$
(73)
The first three parameters are very close to the corresponding powder-averaged anisotropic single crystal values in Table VI which were obtained in the above section assuming the $`S=1/2`$ AF alternating-chain Heisenberg model for the spin susceptibility, and the two spin gaps are nearly identical. The agreement between the spin gaps found from the two independent fits supports the applicability of this spin Hamiltonian to the spin system in HP-$`(\mathrm{VO})_2\mathrm{P}_2\mathrm{O}_7`$.
## IV Summary and Conclusions
We have carried out detailed modeling studies of the magnetic susceptibilities $`\chi (T)`$ of powder and single crystal samples of the ambient- and high-pressure phases of (VO)<sub>2</sub>P<sub>2</sub>O<sub>7</sub>. The major goal of the modeling was to determine whether the recent proposals of a two-chain model for AP-(VO)<sub>2</sub>P<sub>2</sub>O<sub>7</sub> (Refs. and ) and a single-chain model for HP-(VO)<sub>2</sub>P<sub>2</sub>O<sub>7</sub> (Ref. ) are consistent with the respective experimental $`\chi (T)`$ data. Using the high-accuracy theoretical $`\chi ^{}(t,\alpha )`$ function for isolated $`S=1/2`$ AF alternating-exchange Heisenberg chains in Ref. , high precision tests of these models were possible.
The $`\chi (T)`$ data for each phase were first analyzed using the AF alternating-exchange chain model for isolated chains. We found that the proposed models are strongly supported by our high-precision fits to the $`\chi (T)`$ data for each phase, from which the exchange constants and spin gap of each type of chain in each phase were determined. We then considered the case of coupled chains. The influences of interchain couplings on the values of the intrachain exchange constants and the spin gap of each type of chain in the two phases were evaluated from additional fits to the $`\chi (T)`$ data where the interchain coupling was treated in the molecular field approximation. For both phases, we find that the interchain molecular field coupling constant is weakly ferromagnetic with a value $`\lambda 0.05`$, in agreement with Ref. and in disagreement with Ref. . Assuming that the interchain coordination number is two and using the values of $`\lambda `$ and the intrachain exchange constants, the interchain exchange coupling constant along the $`a`$ axis direction is computed to be $`J_a/k_\mathrm{B}3.0(5)`$ K in both phases of (VO)<sub>2</sub>P<sub>2</sub>O<sub>7</sub>. Thus although our modeling of $`\chi (T)`$, from which the intrachain exchange constants and spin gaps were derived, did not explictly incorporate the influence of the magnon dispersion perpendicular to the chain direction, we believe that our mean-field treatment effectively captures most of these effects on $`\chi (T)`$ since the interchain coupling is found to be very weak compared to the intrachain couplings. This supposition is confirmed by our fit to the low-$`T`$ powder-averaged data for a single crystal of HP-(VO)<sub>2</sub>P<sub>2</sub>O<sub>7</sub> by a model-independent low-$`T`$ approximation for the spin susceptibility of a 1D spin system, which yielded a spin gap of 33.8(2) K that is identical within the error bars to the spin gap of 33.4(2) K obtained from a fit to the three complete anisotropic $`\chi (T)`$ data sets for the crystal using the $`\chi ^{}(t)`$ spin susceptibility function for the alternating-exchange chain. The good agreement of the respective spin gaps with those in Table I obtained from inelastic neutron scattering and NMR measurements also supports the magnetic models that we have used for the two phases.
According to the usual simple model for $`d`$ orbitals of transition metal atoms in a distorted octahedral crystalline electric field, the value of the Van Vleck susceptibility $`\chi ^{\mathrm{VV}}`$ increases as the deviation of the $`g`$ value from the free-electron value of 2 increases. Thus from the $`g`$ values determined by ESR for HP-(VO)<sub>2</sub>P<sub>2</sub>O<sub>7</sub> in Table II, one would predict that $`\chi _a^{\mathrm{VV}}>\chi _c^{\mathrm{VV}}\chi _b^{\mathrm{VV}}`$ for this phase. This expectation is borne out by the values of $`\chi ^{\mathrm{VV}}`$ in Table VI for a high-quality single crystal of this phase. This agreement further supports our conclusion that the $`\chi (T)`$ data are consistent with the presence of a single type of $`S=1/2`$ AF alternating-exchange Heisenberg chain in this phase. The powder average of our $`\chi ^{\mathrm{VV}}`$ values in Table II is close to the value $`6\times 10^5\mathrm{cm}^3/\mathrm{mol}\mathrm{V}`$ estimated from $`\chi (T)`$ and $`{}_{}{}^{31}K(T)`$ NMR measurements by Kikuchi et al. for AP-$`(\mathrm{VO})_2\mathrm{P}_2\mathrm{O}_7`$.
Additional confirmation of the two-chain model for AP-(VO)<sub>2</sub>P<sub>2</sub>O<sub>7</sub> is the agreement we find between our predicted one-magnon dispersion relations in the chain direction for the two chains with the results of inelastic neutron scattering measurements at small and large wavevectors. In the intermediate wavevector regime, our calculated dispersion relations of the respective chains predict that they should cross. To our knowledge, there are no relevant inelastic neutron scattering data yet with which to test this prediction. With the caveat given in the next paragraph, our final estimates of the intrachain exchange constants and of the spin gaps of the respective alternating-exchange chains in the two phases are given in Table I, where the error bars on each quantity take our mean-field modeling of interchain interactions into account.
By fitting the experimental dispersion relations perpendicular to the two chains in AP-(VO)<sub>2</sub>P<sub>2</sub>O<sub>7</sub> of Garrett et al. by the theoretical predictions of Uhrig and Normand which incorporate the influence of interchain couplings, both of the couplings $`J_a`$ and $`J_\times `$ were found to be small but nonnegligible. In addition, the theoretical dispersion relations show that these couplings change the spin gap from that of an isolated chain with the same intrachain exchange constants, whereas our modeling of the experimental $`\chi (T)`$ data including the influence of the interchain coupling in a mean-field approximation implicitly assumed that the interchain couplings do not change the spin gap. Using our interchain exchange constants obtained from fitting the experimental dispersion relations perpendicular to the chains using their theory and using our intrachain exchange constants obtained from modeling the experimental $`\chi (T)`$ data, the spin gap of each chain was calculated using their theory to be significantly smaller than the actual spin gap for each chain. Thus the intrachain exchange constants we obtain from the mean-field treatment of the interchain coupling should be considered to be effective values within this model. An improved evaluation of the exchange constants from $`\chi (T)`$ data will only be possible using a theoretical expression for $`\chi (T)`$ which incorporates the effects of the interchain couplings on the two spin gaps.
The spin gap of HP-(VO)<sub>2</sub>P<sub>2</sub>O<sub>7</sub> obtained from analyzing $`\chi (T)`$ data using the low-$`T`$ approximation $`\chi ^{}(t)=(A/\sqrt{t})\mathrm{exp}(\mathrm{\Delta }^{}/t)`$ \[Eq. (24)\] is found to be different than obtained using the above high-accuracy theoretical $`\chi ^{}(t,\alpha )`$ function for alternating-exchange chains to analyze the same data. For example, from a comparison of the spin gap obtained previously for a powder sample of HP-$`(\mathrm{VO})_2\mathrm{P}_2\mathrm{O}_7`$ using this approximation with the spin gap we obtained from a fit to the same data set using the accurate $`\chi ^{}(t,\alpha )`$ function, we infer that the error involved in determining the spin gap using this low-$`T`$ approximation is about 50 % in this case. Similar discrepancies have been found previously when analyzing $`\chi (T)`$ data for 1D spin systems in a similar way. In the compound $`\mathrm{SrCu}_2\mathrm{O}_3`$, for example, the spin gap of the $`S=1/2`$ Cu<sup>+2</sup> two-leg ladders within the $`\mathrm{Cu}_2\mathrm{O}_3`$ trellis layers obtained by fitting $`\chi (T)`$ data up to temperatures $`T\mathrm{\Delta }/k_\mathrm{B}`$ using Eq. (24) (assuming that $`A`$ is an independently adjustable parameter) yielded $`\mathrm{\Delta }/k_\mathrm{B}=420`$ K, whereas inelastic neutron scattering measurements on this compound yielded $`\mathrm{\Delta }/k_\mathrm{B}380`$ K. On the other hand, we found that the low-$`T`$ approximation $`\chi ^{}(t)=(A/t)\mathrm{exp}(\mathrm{\Delta }^{}/t)`$, in which the power of $`t`$ in the prefactor to the exponential is modified, can yield much more accurate values of the spin gap.
Our AF exchange constants in Table I along the alternating-exchange V chains in the two phases of $`(\mathrm{VO})_2\mathrm{P}_2\mathrm{O}_7`$ are of the same order as the nearest-neighbor exchange interactions estimated experimentally and theoretically between the V ions in the two-leg ladder compound $`\mathrm{MgV}_2\mathrm{O}_5`$, but are much smaller than the value of 660–670 K found for the V-V coupling across the ladder rungs in isostructural $`\mathrm{CaV}_2\mathrm{O}_5`$. Korotin et al. have inferred theoretically that the large differences between the exchange constants in the latter two compounds arise from the stronger tilting of the $`\mathrm{VO}_5`$ square pyramids in $`\mathrm{MgV}_2\mathrm{O}_5`$ as compared to $`\mathrm{CaV}_2\mathrm{O}_5`$. The conventional empirical rules for estimating the strengths of nearest-neighbor superexchange interactions in oxides are strongly violated in $`\mathrm{CaV}_2\mathrm{O}_5`$ and also in cuprate spin ladder compounds, as extensively discussed in Ref. . A similar analysis of the exchange coupling strengths in the two phases of $`(\mathrm{VO})_2\mathrm{P}_2\mathrm{O}_7`$ would be informative and perhaps quite relevant to a more general evaluation of this issue.
###### Acknowledgements.
We thank S. E. Nagler and G. S. Uhrig for helpful discussions and correspondence, and our recent collaborators in Ref. on work which made the present study possible. We are grateful to H. Schwenk for the anisotropic $`\chi (T)`$ data for a single crystal of AP-$`(\mathrm{VO})_2\mathrm{P}_2\mathrm{O}_7`$ in Fig. 4(a) of Ref. , designated as “crystal 2” in the present paper, to S. E. Nagler for sending the dispersion relation data for AP-$`(\mathrm{VO})_2\mathrm{P}_2\mathrm{O}_7`$ in Fig. 3 of Ref. , and to C. Knetter for sending the expansion coefficients for the dispersion relation of the frustrated alternating-exchange chain in Ref. prior to publication. Ames Laboratory is operated for the U.S. Department of Energy by Iowa State University under Contract No. W-7405-Eng-82. The work at Ames Laboratory was supported by the Director for Energy Research, Office of Basic Energy Sciences. This work was partly supported by CREST (Core Research for Evolutional Science and Technology) of Japan Science and Technology Corporation (JST) and Grant-in-Aid for Scientific Research of the Ministry of Education, Science, Sports and Culture, Japan.
|
warning/0003/math0003043.html
|
ar5iv
|
text
|
# Between Sobolev and Poincaré Research partially supported by KBN Grant 2 P03A 043 15
## Abstract
Let $`a[0,1]`$ and $`r[1,2]`$ satisfy relation $`r=2/(2a).`$ Let $`\mu (dx)=c_r^n\mathrm{exp}((|x_1|^r+|x_2|^r+\mathrm{}+|x_n|^r))dx_1dx_2\mathrm{}dx_n`$ be a probability measure on the Euclidean space $`(R^n,).`$ We prove that there exists a universal constant $`C`$ such that for any smooth real function $`f`$ on $`R^n`$ and any $`p[1,2)`$
$$E_\mu f^2(E_\mu |f|^p)^{2/p}C(2p)^aE_\mu f^2.$$
We prove also that if for some probabilistic measure $`\mu `$ on $`R^n`$ the above inequality is satisfied for any $`p[1,2)`$ and any smooth $`f`$ then for any $`h:R^nR`$ such that $`|h(x)h(y)|xy`$ there is $`E_\mu |h|<\mathrm{}`$ and
$$\mu (hE_\mu h>\sqrt{C}t)e^{Kt^r}$$
for $`t>1,`$ where $`K>0`$ is some universal constant.
Let us begin with few definitions.
###### Definition 1
Let $`(\mathrm{\Omega },\mu )`$ be a probability space and let $`f`$ be a measurable, square integrable non-negative function on $`\mathrm{\Omega }.`$ For $`p[1,2)`$ we define the $`p`$-variance of $`f`$ by
$$Var(p)_\mu (f)=_\mathrm{\Omega }f(x)^2\mu (dx)(_\mathrm{\Omega }f(x)^p\mu (dx))^{2/p}=E_\mu f^2(E_\mu f^p)^{2/p}.$$
Note that $`Var(1)_\mu (f)=D_\mu ^2(f)=Var_\mu (f)`$ coincides with classical notion of variance, while
$$\underset{p2^{}}{lim}\frac{Var(p)_\mu (f)}{2p}=\frac{1}{2}(E_\mu f^2\mathrm{ln}(f^2)E_\mu f^2\mathrm{ln}(E_\mu f^2))=\frac{1}{2}Ent_\mu (f^2),$$
where $`Ent_\mu `$ denotes a classical entropy functional (see \[L\] for a nice introduction to the subject).
###### Definition 2
Let $``$ be a non-negative functional on some class $`𝒞`$ of non-negative functions from $`L^2(\mathrm{\Omega },\mu ).`$ We will say that $`f𝒞`$ satisfies
* the Poincaré inequality with constant $`C`$
if $`Var_\mu (f)C(f),`$
* the logarithmic Sobolev inequality with constant $`C`$
if $`Ent_\mu (f^2)C(f),`$
* the inequality $`I_\mu (a)`$ (for $`0a1`$) with constant $`C`$
if $`Var(p)_\mu (f)C(2p)^a(f)`$ for all $`p[1,2).`$
###### Lemma 1
For a fixed $`f𝒞`$ and $`p[1,2)`$ let
$$\phi (p)=\frac{Var(p)_\mu (f)}{1/p1/2}.$$
Then $`\phi `$ is a non-decreasing function.
Proof. Hölder’s inequality yields that $`\alpha (t)=t\mathrm{ln}(E_\mu f^{1/t})`$ is a convex function for $`t(1/2,1].`$ Hence also $`\beta (t)=e^{2\alpha (t)}=(E_\mu f^{1/t})^{2t}`$ is convex and therefore $`\frac{\beta (t)\beta (1/2)}{t1/2}`$ is non-decreasing on $`(1/2,1].`$ Observation that
$$\phi (p)=\frac{\beta (1/2)\beta (1/p)}{1/p1/2}$$
completes the proof. $`\mathrm{}`$
###### Corollary 1
For $`f𝒞`$ the following implications hold true:
* $`f`$ satisfies the Poincaré inequality with constant $`C`$
if and only if $`f`$ satisfies $`I_\mu (0)`$ with constant $`C,`$
* if $`f`$ satisfies the logarithmic Sobolev inequality with constant $`C`$
then $`f`$ satisfies $`I_\mu (1)`$ with constant $`C,`$
* if $`f`$ satisfies $`I_\mu (1)`$ with constant $`C`$
then $`f`$ satisfies the logarithmic Sobolev inequality with constant $`2C,`$
* if $`f`$ satisfies $`I_\mu (a)`$ with constant $`C`$ and $`0\alpha a1`$
then $`f`$ satisfies $`I_\mu (\alpha )`$ with constant $`C.`$
Proof.
* To prove the first part of Corollary 1 it suffices to note that $`pVar(p)_\mu (f)`$ is a non-increasing function.
* The second part of Corollary 1 follows easily from the fact that
$$\underset{p2^{}}{lim}\frac{Var(p)_\mu (f)}{2p}=\frac{1}{2}Ent_\mu (f^2).$$
* To prove the third part of Corollary 1 use Lemma 1 and note that for $`p[1,2)`$ we have
$$\frac{Var(p)_\mu (f)}{2p}=\frac{\phi (p)}{2p}\frac{lim_{p2^{}}\phi (p)}{2}=Ent_\mu (f^2).$$
* The last part of statement is trivial. $`\mathrm{}`$
Corollary 1 shows that inequalities $`I_\mu (a)`$ interpolate between Poincaré and logarithmic Sobolev inequalities. Note that $`I_\mu (a)`$ for $`a<0`$ would be equivalent to the Poincaré inequality and the only functions satisfying $`I_\mu (a)`$ for $`a>1`$ would be the constant functions (because in this case $`I_\mu (a)`$ would imply the logarithmic Sobolev inequality with constant $`0`$). Therefore restriction to $`a[0,1]`$ is natural.
###### Definition 3
Given probability space $`(\mathrm{\Omega },\mu ),`$ a class $`𝒞L_+^2(\mathrm{\Omega },\mu )`$ and non-negative functional $``$ on $`𝒞`$ we will say that a pair $`(\mu ,)`$ satisfies $`I(a)`$ (respectively the Poincaré or the logarithmic Sobolev) inequality if every $`f𝒞`$ satisfies $`I_\mu (a)`$ (resp. the Poincaré or the logarithmic Sobolev) inequality with constant C (for these particular $`\mu `$ and $``$). For the sake of brevity we will assume that $`\mu `$ identifies probability space and $``$ carries information about $`𝒞.`$
An obvious modification of Corollary 1 for pairs $`(\mu ,)`$ follows. In some cases we can establish the precise relation between best possible constants in $`I(1)`$ and logarithmic Sobolev inequalities.
Let $`m:(a,a)R`$ be an even, strictly postive continuous density of some probability measure $`\mu `$ on $`(a,a),`$ where $`0<a\mathrm{}`$ and assume that $`_a^ax^2m(x)𝑑x<\mathrm{}.`$ For $`fC_0^{\mathrm{}}(a,a)`$ put
$$(Lf)(x)=xf^{}(x)u(x)f^{\prime \prime }(x),$$
where $`u(x)=\frac{_x^atm(t)𝑑t}{m(x)}0.`$ General theory (see \[KLO\] for detailed references and some related results) yields that $`L`$ can be extended to a positive definite self-adjoint operator (denoted by the same symbol), defined on a dense subspace $`Dom(L)`$ of $`L^2((a,a),\mu ),`$ whose spectrum $`\sigma (L)`$ is contained in $`\{0\}[1,\mathrm{}).`$ Moreover $`P_t=e^{tL}`$ ($`t0`$) is a Markov semigroup with invariant measure $`\mu .`$ Put $`(f)=L^{1/2}f_2^2`$ (we accept $`(f)=+\mathrm{}`$ for $`f`$ which do not belong to $`Dom(L^{1/2})`$) and take $`𝒞=L_+^2((a,a),\mu ).`$
###### Lemma 2
Under the above assumptions the following equivalence holds true:
$`(\mu ,)`$ satisifes the inequality $`I(1)`$ with constant $`C`$
if and only if
$`(\mu ,)`$ satisfies the logarithmic Sobolev inequality with constant $`2C.`$
Proof. If $`(\mu ,)`$ satisifes the inequality $`I(1)`$ with constant $`C`$ then by Corollary 1 it satisfies the logarithmic Sobolev inequality with constant $`2C.`$ Now let us assume that $`(\mu ,)`$ satisfies the logarithmic Sobolev inequality with constant $`2C.`$ Then for any $`fL^2((a,a),\mu )`$ we have
$$Ent_\mu (f^2)=Ent_\mu (|f|^2)2C(|f|)2C(f)$$
(the last inequality is a well known property of Dirichlet forms of Markov semigroups - see for example Theorem 1. 3. 2 of \[D\]). Therefore classical hypercontractivity result \[G\] yields
$$P_{t(p)}f_2f_p,$$
where $`t(p)=\frac{C}{2}\mathrm{ln}(\frac{1}{p1})`$ for $`p[1,2);`$ if $`p=1`$ then we put $`t(p)=\mathrm{}`$ and $`P_{\mathrm{}}(f)=E_\mu f.`$ Hence
$$Efe^{2t(p)L}f(Ef^p)^{2/p}$$
or equivalently
$$Ef^2(Ef^p)^{2/p}Ef(Ide^{2t(p)L})f$$
for any $`f𝒞.`$ Now it suffices to prove that for any $`\lambda \sigma (L)`$ we have
$$1e^{2t(p)\lambda }(2p)C\lambda ,$$
i.e.
$$1(2p)C\lambda (p1)^{C\lambda }.$$
For $`\lambda =0`$ and $`p(1,2)`$ the inequality is trivial. It is known that if $`(\mu ,)`$ satisfies the logarithmic Sobolev inequality with constant $`2C`$ then (under the assumptions of Lemma 2) $`C1`$ \- to see this consider the logarithmic Sobolev inequality for functions of the form $`f(x)=|1+\epsilon x|`$ with $`\epsilon `$ tending to zero (this is a special case of more general observation which says that, for functionals $``$ satisfying certain natural conditions, if $`(\mu ,)`$ satisfies the logarithmic Sobolev inequality with constant $`2C`$ then it also satisfies the Poincaré inequality with constant $`C`$). We can restrict our considerations to the case $`\lambda 1`$ since $`\sigma (L)\{0\}[1,\mathrm{}).`$ Therefore $`(p1)^{C\lambda }`$ is a convex function of $`p`$ and to prove that
$$h(p)=(p1)^{C\lambda }+(2p)C\lambda 10$$
for $`p[1,2)`$ it suffices to check that $`h(2)=h^{}(2)=0`$ which is obvious. The case $`p=1`$ (omitted when $`\lambda =0`$ because $`(p1)^{C\lambda }`$ was not well defined) follows easily since the function $`p(Ef^p)^{2/p}`$ is continuous for $`p[1,2].`$ $`\mathrm{}`$
###### Corollary 2
If $`\mu `$ is a $`𝒩(0,1)`$ Gaussian measure on real line, $`(f)=E_\mu (f^{})^2`$ and $`𝒞`$ is a class of non-negative smooth functions then $`(\mu ,)`$ satisfies $`I(1)`$ with constant 1.
Proof. If $`\mu `$ is a $`𝒩(0,1)`$ Gaussian measure and operator $`L`$ is defined as before then
$$E_\mu fLf=E_\mu (f^{})^2.$$
The assertion follows from Lemma 2 and well known fact (\[G\]) that Gaussian measures satisfy the logarithmic Sobolev inequality with constant 2. $`\mathrm{}`$
###### Remark 1
Method used in Lemma 2 seems applicable also in more general situation (see \[O\] for possible directions of generalization). Let us mention just one interesting application. If $`\mathrm{\Omega }=\{1,1\},`$ $`\mu (\{1\})=\mu (\{1\})=1/2`$ and $`(f)=(\frac{f(1)f(1)}{2})^2`$ then $`(\mu ,)`$ satisfies $`I(1)`$ with constant 1.
###### Remark 2
Let $`\mu `$ be a non-symmetric two-point distribution on $`\{1,1\}`$, $`\mu (\{1\})=1\mu (\{1\})=\alpha `$ with $`\alpha (0,1/2)(1/2,1)`$. Then for any $`p[1,2)`$ and any $`f:\{1,1\}R_+`$ the inequality
$$E_\mu f^2(E_\mu f^p)^{2/p}C_\alpha (p)(f(1)f(1))^2$$
holds with
$$C_\alpha (p)=\frac{\alpha ^{12/p}(1\alpha )^{12/p}}{\alpha ^{2/p}(1\alpha )^{2/p}}$$
and the constant cannot be improved.
Proof (sketch). To check the optimality of $`C_\alpha (p)`$ put $`f(1)=\alpha ^{2/p}`$ and $`f(1)=(1\alpha )^{2/p}`$. To prove the inequality observe that for $`p(1,2)`$, $`\phi (y)=((1+\sqrt{y})^p+(1\sqrt{y})^p)^{2/p}`$ is a strictly convex function of $`y(0,1)`$, since
$$\phi ^{^{}}(y)=[(1+\sqrt{y})^p+(1\sqrt{y})^p]^{\frac{2}{p}1}\frac{(1+\sqrt{y})^{p1}(1\sqrt{y})^{p1}}{\sqrt{y}}$$
$$=\left(2\underset{k=0}{\overset{\mathrm{}}{}}\left(\genfrac{}{}{0pt}{}{p}{2k}\right)y^k\right)^{\frac{2}{p}1}2\underset{k=0}{\overset{\mathrm{}}{}}\left(\genfrac{}{}{0pt}{}{p1}{2k+1}\right)y^k$$
is clearly increasing (note that $`\left(\genfrac{}{}{0pt}{}{p}{2k}\right)`$ and $`\left(\genfrac{}{}{0pt}{}{p1}{2k+1}\right)`$ are positive for $`k=0,1,\mathrm{}`$). Hence for each $`y_0(0,1)`$ and $`p(1,2)`$ there exist unique real numbers $`A`$ and $`B`$ such that
$$\phi (y^2)=((1+y)^p+(1y)^p)^{2/p}A+By^2\text{ for all }y(1,1)$$
with equality holding for $`|y|=y_0`$ only. By the homogenity we may assume that $`f(1)=(1\alpha )^{1/p}(1+y)`$ and $`f(1)=\alpha ^{1/p}(1y)`$. Putting $`y_0=\frac{(1\alpha )^{1/p}\alpha ^{1/p}}{(1\alpha )^{1/p}+\alpha ^{1/p}}`$, using the above inequality after some elementary, but a little involved computations one proves the assertion. $`\mathrm{}`$
###### Definition 4
Let us denote by $`\mathrm{\Phi }`$ the class of all continuous functions $`\phi :[0,\mathrm{})R`$ having strictly positive second derivatve and such that $`1/\phi ^{\prime \prime }`$ is a concave function. Let us additionally include in $`\mathrm{\Phi }`$ all functions $`\phi `$ of the form $`\phi (x)=ax+b,`$ where $`a`$ and $`b`$ are some real constants.
Although it is not obvious, functions belonging to $`\mathrm{\Phi }`$ form a convex cone. There are some interesting questions connected with the class $`\mathrm{\Phi }`$ and its generalizations but we postpone them till the end of the note.
###### Lemma 3
For any $`\phi \mathrm{\Phi }`$ and $`t[0,1]`$ the function $`F_t:[0,\mathrm{})\times [0,\mathrm{})R`$ defined by
$$F_t(x,y)=t\phi (x)+(1t)\phi (y)\phi (tx+(1t)y)$$
is non-negative and convex.
Proof. Non-negativity of $`F_t`$ is an easy consequence of convexity of $`\phi .`$ Obviously $`F_t`$ is continuous on $`[0,\mathrm{})\times [0,\mathrm{})`$ and twice differentiable on $`(0,\mathrm{})\times (0,\mathrm{}).`$ Therefore it suffices to prove that $`HessF_t`$ (second derivative matrix) is positive definite on $`(0,\mathrm{})\times (0,\mathrm{}).`$ We skip the trivial case of $`\phi `$ being an affine function. Note that from the positivity of $`\phi ^{\prime \prime }`$ and the concavity of $`1/\phi ^{\prime \prime }`$ it follows that
$$\frac{1}{\phi ^{\prime \prime }(tx+(1t)y)}\frac{t}{\phi ^{\prime \prime }(x)}+\frac{1t}{\phi ^{\prime \prime }(y)}\frac{t}{\phi ^{\prime \prime }(x)}.$$
Therefore
$$\frac{^2F_t}{x^2}(x,y)=t\phi ^{\prime \prime }(x)t^2\phi ^{\prime \prime }(tx+(1t)y)0.$$
In a similar way we prove that $`\frac{^2F_t}{y^2}(x,y)0.`$ Now it is enough to prove that $`det(HessF_t)0`$ i.e. that
$$\frac{^2F_t}{x^2}(x,y)\frac{^2F_t}{y^2}(x,y)(\frac{^2F_t}{xy}(x,y))^2$$
which is equivalent to
$$(t\phi ^{\prime \prime }(x)t^2\phi ^{\prime \prime }(tx+(1t)y))((1t)\phi ^{\prime \prime }(y)(1t)^2\phi ^{\prime \prime }(tx+(1t)y))$$
$$(t(1t)\phi ^{\prime \prime }(tx+(1t)y))^2$$
or
$$\phi ^{\prime \prime }(x)\phi ^{\prime \prime }(y)t\phi ^{\prime \prime }(y)\phi ^{\prime \prime }(tx+(1t)y)+(1t)\phi ^{\prime \prime }(x)\phi ^{\prime \prime }(tx+(1t)y).$$
After dividing by $`\phi ^{\prime \prime }(x)\phi ^{\prime \prime }(y)\phi ^{\prime \prime }(tx+(1t)y)`$ the last inequality follows from concavity of $`1/\phi ^{\prime \prime }`$ and the proof is complete. $`\mathrm{}`$
###### Lemma 4
For a non-negative real random variable $`Z`$ defined on probability space $`(\mathrm{\Omega },\mu )`$ and having finite first moment, and for $`\phi \mathrm{\Phi }`$ let
$$\mathrm{\Psi }_\phi (Z)=E_\mu \phi (Z)\phi (E_\mu Z).$$
Then for any non-negative real random variables $`X`$ and $`Y`$ defined on $`(\mathrm{\Omega },\mu )`$ and having finite first moment, and for any $`t[0,1]`$ the following inequality holds:
$$\mathrm{\Psi }_\phi (tX+(1t)Y)t\mathrm{\Psi }_\phi (X)+(1t)\mathrm{\Psi }_\phi (Y);$$
in other words $`\mathrm{\Psi }_\phi `$ is a convex functional on the convex cone of integrable non-negative real random variables defined on $`(\mathrm{\Omega },\mu ).`$
Proof. Let us note that (under notation of Lemma 3)
$$\mathrm{\Psi }_\phi (tX+(1t)Y)t\mathrm{\Psi }_\phi (X)(1t)\mathrm{\Psi }_\phi (Y)=$$
$$(E_\mu \phi (tX+(1t)Y)tE_\mu \phi (X)(1t)E_\mu \phi (Y))$$
$$(\phi (tE_\mu X+(1t)E_\mu Y)t\phi (E_\mu X)(1t)\phi (E_\mu Y))$$
$$=E_\mu F_t(X,Y)F_t(E_\mu X,E_\mu Y)=E_\mu F_t(X,Y)F_t(E_\mu (X,Y)).$$
We are to prove that it is a non-negative expression and this follows easily from Jensen inequality. For the sake of clarity we present a detailed argument.
Let $`x_0=E_\mu X`$ and $`y_0=E_\mu Y.`$ Lemma 3 yields that $`F_t`$ is convex, so that there exist constants $`a,b,cR`$ such that
$$F_t(x,y)ax+by+c$$
for any $`x,y[0,\mathrm{})`$ and
$$F_t(x_0,y_0)=ax_0+by_0+c.$$
Therefore
$$E_\mu F_t(X,Y)E_\mu (aX+bY+c)=ax_0+by_0+c=F_t(x_0,y_0)=F_t(E_\mu X,E_\mu Y)$$
and the proof is finished. $`\mathrm{}`$
###### Lemma 5
Let $`(\mathrm{\Omega }_1,\mu _1)`$ and $`(\mathrm{\Omega }_2,\mu _2)`$ be probability spaces and let $`(\mathrm{\Omega },\mu )=(\mathrm{\Omega }_1\times \mathrm{\Omega }_2,\mu _1\mu _2)`$ be their product probability space. For any non-negative random variable $`Z`$ defined on $`(\mathrm{\Omega },\mu )`$ and having finite first moment and for any $`\phi \mathrm{\Phi }`$ the following inequality holds true:
$$E_\mu \phi (Z)\phi (E_\mu Z)E_\mu ([E_{\mu _1}\phi (Z)\phi (E_{\mu _1}Z)]+[E_{\mu _2}\phi (Z)\phi (E_{\mu _2}Z)]).$$
Proof. For $`\omega _2\mathrm{\Omega }_2`$ let $`Z_{(\omega _2)}`$ be a non-negative random variable defined on $`(\mathrm{\Omega }_1,\mu _1)`$ by the formula
$$Z_{[\omega _2]}(\omega _1)=Z(\omega _1,\omega _2).$$
By Lemma 4 used for the probability space $`(\mathrm{\Omega }_1,\mu _1)`$ and Jensen inequality used for the family of random variables $`(Z_{[\omega _2]})_{\omega _2\mathrm{\Omega }_2}`$ (this time we skip the detailed argument which the reader can easily repeat after the proof of Lemma 4) we get
$$E_{\mu _2}(E_{\mu _1}\phi (Z)\phi (E_{\mu _1}Z))E_{\mu _1}\phi (E_{\mu _2}Z)\phi (E_{\mu _1}(E_{\mu _2}Z))$$
which is equivalent to the assertion of Lemma 5. $`\mathrm{}`$
By an easy induction argument we obtain
###### Corollary 3
Let $`(\mathrm{\Omega }_1,\mu _1),(\mathrm{\Omega }_2,\mu _2),\mathrm{},(\mathrm{\Omega }_n,\mu _n)`$ be probability spaces and let $`(\mathrm{\Omega },\mu )=(\mathrm{\Omega }_1\times \mathrm{\Omega }_2\times \mathrm{}\times \mathrm{\Omega }_n,\mu _1\mu _2\mathrm{}\mu _n)`$ be their product probability space. Let $`Z`$ be any integrable non-negative real random variable defined on $`(\mathrm{\Omega },\mu ).`$ Then for any $`\phi \mathrm{\Phi }`$ the following inequality holds:
$$E_\mu \phi (Z)\phi (E_\mu Z)\underset{k=1}{\overset{n}{}}E_\mu (E_{\mu _k}\phi (Z)\phi (E_{\mu _k}Z)).$$
Let us observe that the function $`\phi `$ defined by $`\phi (x)=x^{2/p}`$ belongs to the class $`\mathrm{\Phi }`$ if $`p[1,2].`$ Therefore by applying Corollary 3 to the random variable $`Z=f^p,`$ where $`fL_+^2(\mathrm{\Omega },\mu ),`$ we obtain
###### Corollary 4
Under the notation of Corollary 3 for any $`fL_+^2(\mathrm{\Omega },\mu )`$ we have
$$E_\mu f^2(Ef^p)^{2/p}\underset{k=1}{\overset{n}{}}E_\mu (E_{\mu _k}f^2(E_{\mu _k}f^p)^{2/p}).$$
This sub-additivity property of functional $`Var(p)_\mu `$ immediately yields the following
###### Corollary 5
Assume that pairs $`(\mu _1,_1),`$$`(\mu _2,_2),\mathrm{}`$ $`(\mu _n,_n)`$ satisfy the inequality $`I(a)`$ with some constant $`C.`$ Let $`\mu =\mu _1\mu _2\mathrm{}\mu _n`$ and $`(f)=E_\mu (_1(f_1)+_2(f_2)+\mathrm{}+_n(f_n)),`$ where
$$f_i(x)=f(x_1,\mathrm{},x_{i1},x,x_{i+1},\mathrm{},x_n)$$
for given $`x_1,\mathrm{},x_{i1},x_{i+1},\mathrm{},x_n.`$ Class $`𝒞`$ can be chosen in any way which assures that $`f𝒞`$ implies $`f_i𝒞_i,`$ for example $`𝒞=𝒞_1𝒞_2\mathrm{}𝒞_n.`$ Then the pair $`(\mu ,)`$ also satisfies the inequality $`I(a)`$ with constant $`C.`$
The case we will concentrate on is $`(f)=E_\mu f^2.`$
###### Proposition 1
Let $`\mu _1,`$$`\mu _2,\mathrm{}`$$`\mu _n`$ be probability measures on $`R.`$ Let $`C>0`$ and $`a[0,1].`$ Assume that for any smooth function $`f:R[0,\mathrm{})`$ the inequality
$$E_{\mu _i}f^2(E_{\mu _i}f^p)^{2/p}C(2p)^aE_{\mu _i}(f^{})^2$$
holds true for $`p[1,2)`$ and $`i=1,2,\mathrm{}n.`$ Then for $`\mu =\mu _1\mu _2\mathrm{}\mu _n`$ the inequality
$$E_\mu f^2(E_\mu f^p)^{2/p}C(2p)^aE_\mu f^2,$$
where $``$ denotes standard Euclidean norm, is satisfied for $`p[1,2)`$ and any smooth function $`f:R^n[0,\mathrm{}).`$
Proof. Use Corollary 5 and note that
$$E_\mu f^2=E_\mu [(\frac{f}{x_1})^2+\mathrm{}+(\frac{f}{x_1})^2]=E_\mu [(f_1^{})^2+\mathrm{}+(f_n^{})^2]$$
$$=E_\mu [E_{\mu _1}(f_1^{})^2+\mathrm{}+E_{\mu _n}(f_n^{})^2].\mathrm{}$$
Now let us demonstrate that the inequality $`I(a)`$ for the $`(f)=E_\mu f^2`$ functional implies concentration of Lipschitz functions.
###### Theorem 1
Let $`\mu `$ be a probability measure on $`R^n.`$ Assume that there exist constants $`C>0`$ and $`a[0,1]`$ such that the inequality
$$E_\mu f^2(E_\mu f^p)^{2/p}C(2p)^aE_\mu f^2$$
is satisfied for any smooth function $`f:R^n[0,\mathrm{})`$ and $`p[1,2).`$ Let $`h:R^nR`$ be a Lipschitz function with Lipschitz constant 1, i.e. $`|h(x)h(y)|xy`$ for any $`x,yR^n,`$ where $``$ denotes a standard Euclidean norm. Then $`E_\mu |h|<\mathrm{}`$ and
* for any $`t[0,1]`$
$$\mu (hE_\mu ht\sqrt{C})e^{Kt^2}$$
* for any $`t1`$
$$\mu (hE_\mu ht\sqrt{C})e^{Kt^{\frac{2}{2a}}}$$
where $`K`$ is some universal constant.
Proof. Our proof will work for $`K=1/3`$ but we do not know optimal constants (it is also interesting what the optimal $`K`$ is for given value of parameter $`a`$). Note that it is essential part of the assumptions that we study the limit behaviour when $`p2`$. For any fixed $`p(1,2)`$ the inequality
$$E_\mu f^2(E_\mu f^p)^{2/p}C(2p)^aE_\mu f^2$$
is weaker than the Poincaré inequality with constant $`C(2p)^a`$ and therefore it cannot imply anything stronger than the exponential concentration.
We will follow the aproach of \[AS\]. Assume first that $`h`$ is bounded and smooth. Then $`h1.`$ Define $`H(\lambda )=E_\mu e^{\lambda h}`$ for $`\lambda 0.`$ Assumptions of Theorem 1 for $`f=e^{\lambda h/2}`$ give
$$H(\lambda )H(\frac{p}{2}\lambda )^{2/p}\frac{C\lambda ^2}{4}(2p)^aE_\mu h^2e^{\lambda h}\frac{C\lambda ^2}{4}(2p)^aH(\lambda ).$$
Hence
$$H(\lambda )\frac{H(\frac{p}{2}\lambda )^{2/p}}{1\frac{C}{4}(2p)^a\lambda ^2}$$
for any $`p[1,2)`$ and $`0\lambda \frac{2}{\sqrt{C}}(2p)^{a/2}.`$ Applying the same inequality for $`\frac{p}{2}\lambda `$ and iterating, after $`m`$ steps we get
$$H(\lambda )\frac{H((\frac{p}{2})^m\lambda )^{(2/p)^m}}{_{k=0}^{m1}(1\frac{C\lambda ^2}{4}(2p)^a(\frac{p}{2})^{2k})^{(2/p)^k}}.$$
Note that
$$1\frac{C\lambda ^2}{4}(2p)^a(\frac{p}{2})^{2k}(1\frac{C\lambda ^2}{4}(2p)^a)^{(p/2)^{2k}}$$
since $`(\frac{p}{2})^{2k}<1.`$ Hence
$$H(\lambda )H((\frac{p}{2})^m\lambda )^{(2/p)^m}(1\frac{C\lambda ^2}{4}(2p)^a)^{_{k=0}^{m1}(p/2)^k}.$$
As $`lim_m\mathrm{}(\frac{p}{2})^m=0`$ we get
$$\underset{m\mathrm{}}{lim}H((\frac{p}{2})^m\lambda )^{(2/p)^m}=e^{\lambda E_\mu h}.$$
Therefore
$$E_\mu e^{\lambda (hE_\mu h)}(1\frac{C\lambda ^2}{4}(2p)^a)^{\frac{2}{2p}}$$
and
$$\mu (hE_\mu ht\sqrt{C})e^{\lambda t\sqrt{C}}(1\frac{C\lambda ^2}{4}(2p)^a)^{\frac{2}{2p}}.$$
* Putting $`p=1`$ and $`\lambda =\frac{t}{\sqrt{C}}`$ we get for any $`t[0,2)`$
$$\mu (hE_\mu ht\sqrt{C})e^{t^2}(1\frac{t^2}{4})^2.$$
In particular for $`t[0,1]`$ we have $`1\frac{t^2}{4}>e^{t^2/3}`$ and
$$\mu (hE_\mu ht\sqrt{C})e^{t^2/3}.$$
* If $`t1,`$ let us put $`p=2t^{\frac{2}{2a}}`$ and $`\lambda =t^{\frac{a}{2a}}/\sqrt{C}.`$ Then we arrive at
$$\mu (hE_\mu ht\sqrt{C})e^{t^{\frac{2}{2a}}}(1\frac{1}{4})^{2t^{\frac{2}{2a}}}=(\frac{16}{9e})^{t^{\frac{2}{2a}}}$$
which completes the proof (if $`h`$ is bounded and smooth) since $`\frac{16}{9e}e^{1/3}.`$
Therefore by a standard approximation argument we prove the assertion for any bounded $`h`$ which satisfies assumptions of Theorem 1. Finally for general $`h`$ define its bounded truncations $`(h_N)_{N=1}^{\mathrm{}}`$ putting $`h_N(x)=h(x)`$ if $`|x|N`$ and $`h_N(x)=Nsgn(x)`$ if $`|x|N.`$ One can easily check that if $`h`$ satisfies the assumptions of Theorem 1 then $`|h_N|`$ is also a Lipschitz function with a Lipschitz constant $`1`$ and therefore using Theorem 1 for a bounded function $`|h_N|`$ we arrive at
$$\mu (|h_N|E_\mu |h_N|4\sqrt{C})(\frac{16}{9e})^{4^{\frac{2}{2a}}}(\frac{16}{9e})^4\frac{1}{5}.$$
Similarly
$$\mu (|h_N|E_\mu |h_N|4\sqrt{C})=\mu (|h_N|E_\mu (|h_N|)4\sqrt{C})\frac{1}{5}.$$
Hence
$$\mu (h_N|E_\mu |h_N4\sqrt{C})\frac{2}{5}$$
and
$$\mu (h|E_\mu |h_N4\sqrt{C})\frac{2}{5}+\mu (|h|>N).$$
Therefore
$$\mu (|E_\mu |h_N|E_\mu |h_M||8\sqrt{C})$$
$$\mu (h|E_\mu |h_N4\sqrt{C})+\mu (h|E_\mu |h_M4\sqrt{C})$$
$$\frac{4}{5}+\mu (|h|>N)+\mu (|h|>M)\frac{4}{5}<1$$
as $`\mathrm{min}(N,M)\mathrm{},`$ which means that the sequence $`(E_\mu |h_N|)_{N=1}^{\mathrm{}}`$ is bounded. As $`|h_N|`$ grows monotonically to $`|h|,`$ by Lebesgue Lemma we get $`E_\mu |h|<\mathrm{}`$ and $`E_\mu h_NE_\mu h`$ as $`N\mathrm{}.`$ Now an easy approximation argument completes the proof. $`\mathrm{}`$
In order to prove that the order of concentration implied by Theorem 1 cannot be improved we will need the following
###### Theorem 2
Let $`a[0,1]`$ and $`r[1,2]`$ satisfy $`r=2/(2a).`$ Put $`c_r=\frac{1}{2\mathrm{\Gamma }(1+1/r)}=\frac{r}{2\mathrm{\Gamma }(1/r)}.`$ Then $`\mu _r(dx)=c_r^n\mathrm{exp}((|x_1|^r+|x_2|^r+\mathrm{}+|x_n|^r))dx_1dx_2\mathrm{}dx_n`$ is a probability measure on $`R^n`$ and there exists a universal constant $`C>0`$ (not depending on $`a`$ or $`n`$) such that
$$E_{\mu _r}f^2(E_{\mu _r}f^p)^{2/p}C(2p)^aE_{\mu _r}f^2$$
for any smooth non-negative function $`f`$ on $`R^n`$ and any $`p[1,2).`$
Proof. Proposition 1 shows that it is enough to prove Theorem 2 in the case $`n=1.`$ Therefore the assertion easily follows from the two following propositions. $`\mathrm{}`$
###### Proposition 2
Let $`a[0,1]`$ and $`r[1,2]`$ satisfy $`r=2/(2a).`$ Put $`c_r=\frac{1}{2\mathrm{\Gamma }(1+1/r)},`$ so that $`\mu _r(dx)=c_r\mathrm{exp}(|x_1|^r)dx`$ is a probability measure on $`R.`$ Let $`\lambda (dx)=\frac{1}{2}e^{|x|}`$ be a symmetric exponential probability measure on $`R.`$ Under these assumptions the following implications hold true:
* If $`C>0`$ is a constant such that for any smooth function $`f:R[0,\mathrm{})`$ and any $`p[1,2)`$ there is
$$E_{\mu _r}f^2(E_{\mu _r}f^p)^{2/p}C(2p)^aE_{\mu _r}(f^{})^2$$
then for any smooth function $`g:R[0,\mathrm{})`$ and any $`p[1,2)`$ there is
$$_Rg(x)^2\lambda (dx)(_Rg(x)^p\lambda (dx))^{2/p}600C(2p)^a_R\mathrm{max}(1,|x|^a)g^{}(x)^2\lambda (dx).$$
* Conversely, if $`C>0`$ is such a constant that for any smooth function $`g:R[0,\mathrm{})`$ and any $`p[1,2)`$ there is
$$_Rg(x)^2\lambda (dx)(_Rg(x)^p\lambda (dx))^{2/p}C(2p)^a_R\mathrm{max}(1,|x|^a)g^{}(x)^2\lambda (dx)$$
then for any smooth function $`f:R[0,\mathrm{})`$ and any $`p[1,2)`$ there is
$$E_{\mu _r}f^2(E_{\mu _r}f^p)^{2/p}50C(2p)^aE_{\mu _r}(f^{})^2.$$
###### Proposition 3
There exists a universal constant $`C`$ such that for any $`a[0,1],`$ any $`p[1,2)`$ and any smooth function $`g:R[0,\mathrm{})`$ the following inequality holds
$$_Rg(x)^2\lambda (dx)(_Rg(x)^p\lambda (dx))^{2/p}C(2p)^a_R\mathrm{max}(1,|x|^a)g^{}(x)^2\lambda (dx).$$
We will start with proof of Proposition 2. The proof of Proposition 3 will be postponed to the end of the paper.
Proof of Proposition 2. Let us define the function $`z_r:RR`$ by
$$\frac{1}{2}_{z_r(x)}^{\mathrm{}}e^{|t|}𝑑t=c_r_x^{\mathrm{}}e^{|t|^r}𝑑t,$$
where $`c_r=\frac{r}{2\mathrm{\Gamma }(1/r)}=\frac{1}{2\mathrm{\Gamma }(1+1/r)}.`$ It is easy to see that $`z_r`$ is a homeomorphism of $`R`$ onto itself and
$$z_r^{}(x)=2c_re^{|z_r(x)||x|^r}.$$
Therefore $`z_r`$ is a $`C^1`$diffeomorphism of $`R`$ onto itself. Binding $`f`$ and $`g`$ by relation $`f(x)=g(z_r(x))`$ and using standard change of variables formula we reduce the proof of Proposition 2 to the following lemma. $`\mathrm{}`$
###### Lemma 6
Under notation introduced above
$$\frac{1}{50}\mathrm{max}(1,|x|^a)(z_r^{}(z_r^1(x)))^2600\mathrm{max}(1,|x|^a)$$
for any $`xR.`$
Proof. First let us note that $`1/3c_re/2.`$ Indeed,
$$\mathrm{\Gamma }(1/r)=_0^{\mathrm{}}x^{\frac{1}{r}1}e^x𝑑x_0^1x^{\frac{1}{r}1}𝑑x+_1^{\mathrm{}}e^x𝑑x=r+1/e.$$
Hence $`c_r\frac{r}{2r+2/e}1/3.`$ On the other hand
$$\mathrm{\Gamma }(1/r)=_0^{\mathrm{}}x^{\frac{1}{r}1}e^x𝑑x\frac{1}{e}_0^1x^{\frac{1}{r}1}𝑑x=r/e.$$
Therefore $`c_re/2.`$ Let us also notice that by obvious symmetry we can consider only the case $`x>0.`$ Now let us estimate from below $`z_r^1(1).`$ We have
$$\frac{e}{2}z_r^1(1)c_rz_r^1(1)c_r_0^{z_r^1(1)}e^{t^r}𝑑t=\frac{1}{2}_0^1e^t𝑑t=\frac{1}{2}(11/e)$$
and therefore $`z_r^1(1)\frac{e1}{e^2}1/5.`$ Note that by definition of $`z_r(x)`$ for $`x>0`$ we have
$$\frac{1}{2}e^{z_r(x)}=c_r_x^{\mathrm{}}e^{t^r}𝑑tc_r_x^{\mathrm{}}\frac{rt^{r1}}{rx^{r1}}e^{t^r}𝑑t=\frac{c_re^{x^r}}{rx^{r1}}$$
and therefore
$$z_r^{}(x)=2c_re^{z_r(x)x^r}rx^{r1}.$$
Hence also $`z_r(x)x^r`$ and $`z_r^1(x)x^{1/r}`$ for all positive $`x.`$ If $`x1/5`$ then
$$_x^{\mathrm{}}e^{t^r}𝑑t_x^{6x}e^{t^r}𝑑t\frac{1}{r(6x)^{r1}}_x^{6x}rt^{r1}e^{t^r}𝑑t=$$
$$6^{1r}\frac{e^{x^r}e^{6^rx^r}}{rx^{r1}}\frac{1}{12}\frac{e^{x^r}}{rx^{r1}},$$
since $`6^rx^rx^r+1`$ for $`x1/5`$ and $`r[1,2].`$ Therefore for $`xz_r^1(1)1/5`$ we have
$$z_r^{}(x)12rx^{r1}24x^{r1}$$
and
$$z_r(x)z_r(z_r^1(1))+12_{z_r^1(1)}^xrt^{r1}𝑑t=1+12(x^r[z_r^1(1)]^r)1+12x^r37x^r.$$
Hence $`z_r^1(x)(x/37)^{1/r}`$ for $`xz_r^1(1).`$ If $`x1`$ then $`z_r^1(x)1/5`$ and therefore
$$z_r^{}(z_r^1(x))24[z_r^1(x)]^{r1}24x^{\frac{r1}{r}}=24x^{a/2}.$$
Also if $`x1`$ then $`z_r^1(x)z_r^1(1)`$ and
$$z_r^{}(z_r^1(x))r[z_r^1(x)]^{r1}(x/37)^{\frac{r1}{r}}37^{\frac{1}{r}1}x^{a/2}\frac{1}{7}x^{a/2}.$$
This proves Lemma 6 for $`|x|1.`$ For any $`x0`$ we have
$$z_r^{}(z_r^1(x))=2c_re^{xz_r^1(x)^r}2c_r2/3.$$
We used the previously proved fact that $`z_r^1(x)x^{1/r}.`$ Now it remains only to establish upper estimate on $`z_r^{}(z_r^1(x))`$ for $`x[0,1].`$ Note that if $`xz_r^1(1)`$ then
$$c_r_x^{\mathrm{}}e^{t^r}𝑑t=\frac{1}{2}_{z_r(x)}^{\mathrm{}}e^t𝑑t\frac{1}{2}_1^{\mathrm{}}e^t𝑑t=\frac{1}{2e}$$
and therefore
$$z_r^{}(x)=\frac{2c_re^{x^r}}{2c_r_x^{\mathrm{}}e^{t^r}𝑑t}\frac{c_r}{c_r_x^{\mathrm{}}e^{t^r}𝑑t}2ec_re^28.$$
Hence $`z_r^{}(z_r^1(x))8`$ for any $`|x|1`$ and the proof is finished. $`\mathrm{}`$
###### Lemma 7
For $`s(1,2]`$ and $`x,y0`$ put
$$\rho _s(x,y)=(\frac{x^s+y^s}{2}(\frac{x+y}{2})^s)^{1/2}.$$
Then $`\rho _s`$ is a metric on $`[0,\mathrm{}).`$
Proof. Since $`k_t(a,b)=e^{(a+b)t}`$ is obviously positive definite integral kernel and $`K(a,b)=s(s1)(a+b)^{s2}=\frac{s(s1)}{\mathrm{\Gamma }(2s)}_0^{\mathrm{}}t^{1s}k_t(a,b)𝑑t`$ we get, by Schwartz inequality (applied to a scalar product defined by the kernel $`K(a,b)`$), that for any $`yx0`$ and $`zt0`$ the following inequality is true:
$`{\displaystyle _{x/2}^{y/2}}{\displaystyle _{t/2}^{z/2}}`$ $`K(a,b)dadb`$
$``$ $`({\displaystyle _{x/2}^{y/2}}{\displaystyle _{x/2}^{y/2}}K(a,b)𝑑a𝑑b)^{1/2}({\displaystyle _{t/2}^{z/2}}{\displaystyle _{t/2}^{z/2}}K(a,b)𝑑a𝑑b)^{1/2}.`$
Now, as
$$K(a,b)=\frac{^2}{ab}(a+b)^s,$$
we get by integration by parts
$$(\frac{y+z}{2})^s+(\frac{x+t}{2})^s(\frac{x+z}{2})^s(\frac{y+t}{2})^s$$
$$(x^s+y^s2(\frac{x+y}{2})^s)^{1/2}(z^s+t^s2(\frac{z+t}{2})^s)^{1/2}$$
Putting $`t=y`$ we arrive at
$$(\frac{x+y}{2})^s+(\frac{y+z}{2})^s(\frac{x+z}{2})^sy^s2\rho _s(x,y)\rho _s(y,z)$$
which is equivalent to
$$\rho _s(x,z)^2\rho _s(x,y)^2\rho _s(y,z)^22\rho _s(x,y)\rho _s(y,z).$$
Hence $`\rho _s(x,z)\rho _s(x,y)+\rho _s(y,z).`$ For $`xyz`$ we have also easily $`\rho _s(x,z)\rho _s(x,y)`$ and $`\rho _s(x,z)\rho _s(y,z),`$ so that $`\rho _s(x,y)\rho _s(x,z)+\rho _s(z,y)`$ and $`\rho _s(y,z)\rho _s(y,x)+\rho _s(x,z).`$ This completes the proof of triangle inequality for $`s<2.`$ Other metric properties of $`\rho _s`$ as well as the case $`s=2`$ are trivial. $`\mathrm{}`$
###### Remark 3
In a similar way one can prove that $`\rho _s(x,y)=|\frac{x^s+y^s}{2}(\frac{x+y}{2})^s|^{1/2}`$ is a metric on $`(0,\mathrm{})`$ for $`s(\mathrm{},0)(0,1).`$ It was pointed out to the authors by B. Maurey that Lemma 7 follows also from isometrical immersion of $`([0,\mathrm{}),\rho _s)`$ into $`L^2([0,\mathrm{}),\kappa _s^1t^{s1}dt),`$ where $`x[0,\mathrm{})`$ is sent to the function $`e^{xt}1`$ and $`\kappa _s=2^{s+1}_0^{\mathrm{}}(e^u1+u)u^{s1}𝑑u.`$
###### Lemma 8
Let $`s[1,2]`$, $`t[0,1]`$ and $`c,d,x`$ be nonnegative numbers. The following inequality holds
$$(1t)c^s+td^s((1t)c+td)^s$$
$$K[(1t)c^s+td^s+x^s((1t)c+tx)^s(td+(1t)x)^s].$$
(1)
under anyone of the following additional assumptions
* $`x`$ lies outside the open interval $`(c,d)`$ and $`K=1`$
* $`t=\frac{1}{2}`$ and $`K=2`$
* $`t\frac{1}{2}`$, $`cd`$ and $`K=12`$
Proof. Let us remind that
$$F_t(x,y)=tx^s+(1t)y^s(tx+(1t)y)^s$$
is a convex function on $`[0,\mathrm{})\times [0,\mathrm{}).`$ Note that the inequality of Lemma 8 is equivalent to
$$F_t(d,c)K[F_t(d,x)+F_t(x,c)].$$
* As
$$\frac{}{x}[F_t(d,x)+F_t(x,c)]$$
$$=s[(1t)(x^{s1}(td+(1t)x)^{s1})+t(x^{s1}(tx+(1t)c)^{s1})],$$
we see that the right-hand side of the inequality as a function of $`x`$ is increasing on $`(\mathrm{max}(c,d),\mathrm{})`$ and decreasing on $`[0,\mathrm{min}(c,d)).`$ For $`x=\mathrm{max}(c,d)`$ and $`x=\mathrm{min}(c,d)`$ the inequality is trivially satisfied with $`K=1.`$ This completes the case of $`x`$ which does not lie between $`c`$ and $`d.`$
* The second part of Lemma 8 follows easily by Lemma 7, as
$$F_{1/2}(d,c)=\rho _s(d,c)^2(\rho _s(d,x)+\rho _s(x,c))^2$$
$$2[\rho _s(d,x)^2+\rho _s(x,c)^2]=2[F_{1/2}(d,x)+F_{1/2}(x,c)].$$
* To prove the last part of the statement we will use convexity of $`F_t.`$ Since $`F_t(d,x)+F_t(x,c)F_t(\frac{d+x}{2},\frac{x+c}{2}),`$ it suffices to prove that $`F_t(d,c)12F_t(\frac{d+x}{2},\frac{x+c}{2}).`$ Thanks to the first part of Lemma 8 we can restrict our considerations to the case $`x[d,c].`$ Note that
$$\frac{}{x}F_t(\frac{d+x}{2},\frac{x+c}{2})$$
$$=\frac{s}{2}[t(\frac{d+x}{2})^{s1}+(1t)(\frac{x+c}{2})^{s1}(t(\frac{d+x}{2})+(1t)(\frac{x+c}{2}))^{s1}]0,$$
since the function $`\phi (u)=u^{s1}`$ is concave. Therefore it is enough to prove that
$$F_t(d,c)12F_t(\frac{d+c}{2},c).$$
Using the homogenity of the above formula we can reduce our task to proving that
$$F_t(1u,1)12F_t(1u/2,1)$$
for any $`u[0,1]`$ and $`t[0,1/2].`$
Using the Taylor expansion we get
$$F_t(1u,1)=t(1u)^s+1t(1tu)^s=$$
$$s(s1)u^2t(1t)\left[\frac{1}{2}+\underset{k=1}{\overset{\mathrm{}}{}}\frac{u^k}{(k+1)(k+2)}\underset{m=0}{\overset{k}{}}t^m\underset{l=1}{\overset{k}{}}(1\frac{s1}{l})\right].$$
Therefore
$$F_t(1u/2,1)\frac{1}{2}s(s1)(u/2)^2t(1t)$$
and
$$F_t(1u,1)s(s1)u^2t(1t)\left[\frac{1}{2}+2\underset{k=1}{\overset{\mathrm{}}{}}\frac{1}{(k+1)(k+2)}\right]$$
$$=\frac{3}{2}s(s1)u^2t(1t)$$
because $`_{m=0}^{\mathrm{}}t^m2`$. Hence
$$F_t(1u,1)12F_t(1u/2,1)$$
which completes the proof. $`\mathrm{}`$
###### Lemma 9
Let $`a[0,1]`$, $`0x_1<x_2`$ and $`g`$ be a smooth function on $`[x_1,x_2]`$ such that $`g(x_1)=y_1,g(x_2)=y_2`$. Then
$$_{x_1}^{x_2}\mathrm{max}(1,x^a)g^{}(x)^2𝑑\lambda (x)\frac{(y_2y_1)^2}{4(e^{x_2}e^{x_1})}\mathrm{max}(1,x_2^a).$$
(2)
Proof. By the Schwartz inequality
$$|y_2y_1|_{x_1}^{x_2}|g^{}(x)|𝑑x$$
$$(_{x_1}^{x_2}\mathrm{max}(1,x^a)g^{}(x)^2𝑑\lambda (x))^{1/2}(2_{x_1}^{x_2}\mathrm{min}(1,x^a)e^x𝑑x)^{1/2}.$$
Therefore to show (2) it is enough to prove that
$$f_1(x_2)=_{x_1}^{x_2}\mathrm{min}(1,x^a)e^x𝑑x2\mathrm{min}(1,x_2^a)(e^{x_2}e^{x_1})=f_2(x_2).$$
For $`x_22`$ this is obvious because for $`0<x<x_22`$ we have $`\mathrm{min}(1,x^a)12\mathrm{min}(1,x_2^a),`$ and for $`x2`$ we have
$$f_2^{}(x)=2x^a(e^xax^1(e^xe^{x_1}))x^ae^x=f_1^{}(x).\mathrm{}$$
###### Lemma 10
Let $`0y_1<y_2`$, $`0x_1<x_2`$ and $`g`$ is defined on $`(\mathrm{},x_2)`$ by the formula
$$g(x)=\{\begin{array}{cc}y_1\hfill & \text{for }xx_1\\ y_1+(e^xe^{x_1})\frac{y_2y_1}{e^{x_2}e^{x_1}}\hfill & \text{for }x(x_1,x_2]\end{array}.$$
Then
$$_{\mathrm{}}^{x_2}g^{}(x)^2𝑑\lambda (x)=\frac{(y_2y_1)^2}{2(e^{x_2}e^{x_1})}.$$
(3)
and for all $`p1`$
$$_{\mathrm{}}^{x_2}g(x)^p𝑑\lambda (x)\lambda (\mathrm{},x_2)[(1\frac{x_2}{2}e^{x_2})y_1^p+\frac{x_2}{2}e^{x_2}y_2^p].$$
(4)
Proof. Equation (3) follows by direct calculations. It is easy to see that $`g(x)`$ is maximal (for fixed values of $`x_2,y_1`$ and $`y_2`$) when $`x_1=0`$, so to prove (4) we may and will assume that this is the case. To easy the notation we will denote $`x_2`$ by $`x`$. First we will consider $`p=1`$. After some standard calculations (4) is equivalent in this case to
$$\frac{e^x(x1+e^x)}{(2e^x1)(e^x1)}\frac{1}{2}xe^x\text{ for all }x>0,$$
that is
$$2+3xxe^x+2e^x\text{ for all }x>0,$$
which immeditely follows from well known estimates $`e^x1x`$ and $`e^x1+x+x^2/2`$.
Now, for arbitrary $`p1`$ notice that $`g(x)=(1\theta (x))y_1+\theta (x)y_2`$ with $`0\theta (x)1.`$ Therefore we have by the convexity of $`x^p`$
$$_{\mathrm{}}^{x_2}g(x)^p𝑑\lambda (x)_{\mathrm{}}^{x_2}((1\theta (x))y_1^p+\theta (x)y_2^p)𝑑\lambda (x)$$
$$\lambda (\mathrm{},x_2)[(1\frac{x_2}{2}e^{x_2})y_1^p+\frac{x_2}{2}e^{x_2}y_2^p],$$
where the last inequality follows by the previously established case $`p=1`$. $`\mathrm{}`$
###### Lemma 11
Suppose that $`s(1,2]`$, $`t(0,1)`$, $`u=\frac{s}{4(s1)}e^{s/2(s1)}`$ and positive numbers $`a,b,c,d,\stackrel{~}{a},\stackrel{~}{c},x`$ satisfy the following conditions
$$c<x<d,c^sa,d^sb,\stackrel{~}{c}^s\stackrel{~}{a},\stackrel{~}{c}(1u)c+ux.$$
Then
$$(1t)a+tb((1t)c+td)^s$$
$$8[(1t)\stackrel{~}{a}+tb((1t)\stackrel{~}{c}+td)^s+(1t)a+tx^s((1t)c+tx)^s].$$
(5)
Proof. Without loss of generality we may assume that $`a=c^s,b=d^s,\stackrel{~}{a}=\stackrel{~}{c}^s`$. Since the function $`y(1t)y^s((1t)y+td)^s`$ is nonincreasing on $`[0,d]`$, it is enough to show that
$$(1t)c^s+td^s((1t)c+td)^s$$
$$3[(1t)((1u)c+ud)^s+td^s((1t)(1u)c+(t+(1t)u)d)^s].$$
By the homogenity we may and will assume that $`d=1`$. We are then to show that
$$f((1c))8f((1u)(1c)),$$
(6)
where
$$f(x)=(1t)(1x)^s+t(1(1t)x)^s=\underset{i=2}{\overset{\mathrm{}}{}}(1)^i\left(\genfrac{}{}{0pt}{}{s}{i}\right)(1t)(1(1t)^{i1})x^k.$$
We use the following simple observation: if $`a_i,b_i`$ are two summable sequences of positive numbers such that for any $`i>j`$, $`a_i/a_jb_i/b_j`$ then for any nondecreasing nonnegative sequence $`c_i`$
$$\frac{a_ic_i}{a_i}\frac{b_ic_i}{b_i}.$$
We apply the above to the sequences $`a_i=(1)^i\left(\genfrac{}{}{0pt}{}{s}{i}\right)(1t)(1(1t)^{i1})x^i`$, $`b_i=(i1)(1)^i\left(\genfrac{}{}{0pt}{}{s}{i}\right)`$ and $`c_i=(1u)^i`$, $`i=2,3,\mathrm{}`$ and notice that
$$h(y):=\underset{i=2}{\overset{\mathrm{}}{}}b_iy^i=1(1y)^{s1}(1+(s1)y)\text{ for }y[0,1]$$
Therefore we get
$$f((1u)x)\frac{h(1u)}{h(1)}=\left(1u^{s1}(1+(s1)(1u))\right)f(x)$$
Inequality (6) follows if we notice that
$$u^{s1}(1+(s1)(1u))su^{s1}=\frac{s^s}{4^{s1}}e^{s/2}(\frac{1}{s1})^{s1}1e^{1/2}e^{1/e}\frac{7}{8}\mathrm{}$$
###### Proposition 4
Suppose that for all $`p[1,2)`$ and all nonnegative smooth functions $`g`$ we have
$$_Rg^2𝑑\lambda (_Rg^p𝑑\lambda )^{2/p}K_1(2p)^i_R(g^{}(x))^2\mathrm{max}(1,|x|^i)𝑑\lambda (x)\text{ for }i=0,1,$$
(7)
where $`K_1`$ is a universal constant. Then for all $`p`$ and $`g`$ as above we have
$$_Rg^2𝑑\lambda (_Rg^p𝑑\lambda )^{2/p}$$
$$K_2(2p)^a_R(g^{}(x))^2\mathrm{max}(1,|x|^a)𝑑\lambda (x)\text{ for }a(0,1),$$
(8)
where $`K_232K_1`$ is some universal constant.
Proof. An easy approximation argument shows that (7) holds for any continuous function $`g`$, continuously differentiable everywhere except possibly finitely many points.
First we assume that $`g`$ is constant on $`R^{}`$ or $`R^+`$, without loss of generality say it is $`R^{},`$ and we show that (8) holds with $`K_2=16K_1.`$ Let us fix $`p[1,2)`$ and define
$$x_p=(2p)^1,y=g(x_p),t=\lambda (x_p,\mathrm{}),s=\frac{2}{p},$$
$$a=\frac{1}{1t}_{\mathrm{}}^{x_p}g^2𝑑\lambda ,b=\frac{1}{t}_{x_p}^{\mathrm{}}g^2𝑑\lambda $$
$$c=\frac{1}{1t}_{\mathrm{}}^{x_p}g^p𝑑\lambda \text{ and }d=\frac{1}{t}_{x_p}^{\mathrm{}}g^p𝑑\lambda .$$
Notice that by Hölder’s inequality we have
$$ac^s\text{ and }bd^s.$$
(9)
We will consider two cases
Case 1. $`y^p`$ lies outside $`(c,d)`$ or $`c>d`$.
We first apply inequality (7) for $`i=1`$ and a function $`gI_{(\mathrm{},x_p)}+yI_{[x_p,\mathrm{})}`$ to get
$$(1t)a+ty^2((1t)c+ty^p)^sK_1(2p)_0^{x_p}(g^{}(x))^2\mathrm{max}(1,|x|)𝑑\lambda (x)$$
$$K_1(2p)^a_0^{x_p}(g^{}(x))^2\mathrm{max}(1,|x|^a)𝑑\lambda (x).$$
In a similar way using the case of $`i=0`$ for the function $`yI_{(\mathrm{},x_p)}+gI_{[x_p,\mathrm{})}`$ we get
$$tb+(1t)y^2(td+(1t)y^p)^sK_1_{x_p}^{\mathrm{}}(g^{}(x))^2𝑑\lambda (x)$$
$$K_1(2p)^a_{x_p}^{\mathrm{}}(g^{}(x))^2\mathrm{max}(1,|x|^a)𝑑\lambda (x).$$
Notice also that
$$_Rg^2𝑑\lambda (_Rg^p𝑑\lambda )^{2/p}=(1t)a+tb((1t)c+td)^s$$
$$12\left[(1t)a+ty^2((1t)c+ty^p)^s+tb+(1t)y^2(td+(1t)y^p)^s\right]$$
$$12K_1(2p)^a_R(g^{}(x))^2\mathrm{max}(1,|x|^a)𝑑\lambda (x).$$
The middle inequality follows by Lemma 8 with $`x=y^p`$ together with estimates (9).
Case 2. $`c<y^p<d`$, we can then find $`0<x_0<x_p`$ such that $`g(x_0)=c^{1/p}`$. Define new function $`f`$ by the formula
$$f(x)=\{\begin{array}{cc}g(x)\hfill & \text{for }x>x_p\\ c^{1/p}+\frac{yc^{1/p}}{e^{x_p}e^{x_0}}(e^xe^{x_0})\hfill & \text{for }x[x_0,x_p]\\ c^{1/p}\hfill & \text{for }x<x_0.\end{array}$$
Let
$$\stackrel{~}{a}=\frac{1}{1t}_{\mathrm{}}^{x_p}f^2𝑑\lambda \text{ and }\stackrel{~}{c}=\frac{1}{1t}_{\mathrm{}}^{x_p}f^p𝑑\lambda .$$
By Lemma 9 and 10 we have
$$_Rf^{}(x)^2𝑑\lambda (x)2(2p)^a_R\mathrm{max}(1,|x|^a)g^{}(x)^2𝑑\lambda (x).$$
Therefore by (7) with $`i=0`$, used for the function $`f`$, we have
$$(1t)\stackrel{~}{a}+tb((1t)\stackrel{~}{c})+td)^s2K_1(2p)^a\mathrm{max}(1,|x|^a)g^{}(x)^2d\lambda (x).$$
We conclude as in the previous case using Lemmas 10 and 11 instead of Lemma 8.
Finally suppose that $`g`$ is arbitrary. A similar argument as in case 1 (but now with $`x_p=0`$ and $`t=1/2`$) together with the already proved case of $`g`$ constant on $`R_{}`$ or $`R_+`$ proves the assertion in this case. $`\mathrm{}`$
Proof of Proposition 3. We need only to prove that assumptions of Proposition 4 are satisfied. But in view of Proposition 2 they are equivalent to the Poincaré inequality for symmetric exponential probability measure ($`i=0`$) and the logarithmic Sobolev inequality for the centered $`𝒩(0,\sqrt{2}/2)`$ Gaussian measure ($`i=1`$) which are well known to hold with some universal constants. This completes the proof. $`\mathrm{}`$
In the end of the paper we would like to come back to the class $`\mathrm{\Phi }`$ introduced in Definition 4. It is easy to check that if Lemma 5 holds for some function $`\phi C^2((0,\mathrm{}))C([0,\mathrm{}))`$ for any $`(\mathrm{\Omega }_1,\mu _1),(\mathrm{\Omega }_2,\mu _2)`$ and any $`Z`$ then $`\phi \mathrm{\Phi }.`$ Indeed, it is even true if we restrict our consideration to $`(\mathrm{\Omega }_1,\mu _1)`$ and $`(\mathrm{\Omega }_2,\mu _2)`$ being two-point probability spaces whose atoms have $`1/2`$ measures. This gives a natural characterization of the class $`\mathrm{\Phi }.`$
One can try to generalize the definition of $`\mathrm{\Phi }.`$ Let $`U`$ be an open, convex subset of $`R^d.`$ We will say that a continuous function $`f:UR`$ belongs to the class $`C_n(U)`$ if for any probability spaces $`(\mathrm{\Omega }_1,\mu _1),\mathrm{},(\mathrm{\Omega }_n,\mu _n)`$ and any integrable random variable $`Z`$ with values in $`U,`$ defined on $`(\mathrm{\Omega },\mu )=(\mathrm{\Omega }_1\times \mathrm{}\times \mathrm{\Omega }_n,\mu _1\mathrm{}\mu _n)`$ the following inequality is satisfied:
$$\underset{K\{1,2,\mathrm{},n\}}{}(1)^{|K|}E_{K^c}f(E_KZ)0,$$
where $`E_K`$ denotes expectation with respect to $`\mu _k`$ for all $`kK.`$ One can easily see that $`C_1(U)`$ is just a set of all convex functions on $`U,`$ while $`C_2((0,\mathrm{}))`$ is closely related to the class $`\mathrm{\Phi }.`$ In fact $`fC_2((0,\mathrm{}))`$ if and only if it is an affine function or it has a continuous strictly positive second derivative such that $`1/f^{\prime \prime }`$ is a concave function. One can prove that always $`C_{n+1}(U)C_n(U)`$ and therefore it is natural to define $`C_{\mathrm{}}(U)`$ as an intersection of all $`C_n(U).`$ Then it appears that $`fC_{\mathrm{}}(U)`$ if and only if $`f`$ is given by the formula $`f(x)=Q(x,x)+x^{}(x)+y,`$ where $`Q`$ is a non-negative definite symmetric quadratic form, $`x^{}`$ is a linear functional on $`R^d`$ and $`y`$ is a constant. The above inclusions do not need to be strict. For example it is easy to see that $`C_2(R)`$=$`C_{\mathrm{}}(R).`$ It would be interesting to know some nice characterization of $`C_2(U)`$ for general $`U`$ and $`C_n((0,\mathrm{}))`$ for $`n>2.`$ It is not clear what applications of $`C_n`$ for $`n>2`$ could be found but it is easy to see that this class has some tensorization property. By now, we do not know even the answer to the following question: For which $`p[1,2]`$ does $`f(x)=x^p`$ belong to $`C_n((0,\mathrm{}))\mathrm{?}`$ We can only give some estimates.
These problems will be discussed in a separate paper.
###### Remark 4
Recently some new results were announced to the authors by F. Barthe (private communication) - he proved (using Theorem 2 above) that if a log-concave probability measure $`\mu `$ on the Euclidean space $`(R^n,)`$ satisfies inequality $`\mu (\{xR^n;x>t\})ce^{(t/c)^r}`$ for some constants $`c>0,r[1,2]`$ and any $`t>0`$ then it satisfies also inequality
$$E_\mu f^2(E_\mu f^p)^{2/p}C(c,n,r)(2p)^aE_\mu f^2$$
for any non-negative smooth function $`f`$ on $`R^n`$ and $`p[1,2),`$ where $`C(c,n,r)`$ is some positive constant depending on $`c,n`$ and $`r`$ only and $`a=22/r`$.
Acknowledgements. The article was inspired by the questions of Prof. Stanisław Kwapień and Prof. Gideon Schechtman. This work was done while the first named author was visiting Southeast Applied Analysis Center at School of Mathematics, Georgia Institute of Technology and was partially supported by NSF Grant DMS 96-32032. The research of the second named author was performed at the Weizmann Institute of Science in Rehovot, Israel and Equipe d’Analyse, Université Paris VI.
Institute of Mathematics
Warsaw University
Banacha 2
02-097 Warszawa
Poland
E-mail: rlatala@mimuw.edu.pl, koles@mimuw.edu.pl
|
warning/0003/hep-ph0003002.html
|
ar5iv
|
text
|
# Space-time structure of a bound nucleon
## 1 INTRODUCTION
Decades of matter structure studying has shown that complete information about nucleon structure cannot be obtained from free nucleon data only. First of all, it is connected with absence of a stable enough free neutron target. The attempts to solve this problem were based on utilization of nuclear targets. However, small nuclear effects, which were supposed to be negligible, led to the qualitatively different results for the structure function of a nucleon bound in a deuteron and in heavy nuclei . This phenomenon reflects untrivial difference between nuclear and nucleon structure at internucleon distances and its evolution with atomic number $`A`$. Nature of the effect was analyzed in numerous models which were produced since it was discovered (comprehensive reviews of the models can be found in ). Summing up the basic qualitative pictures, proposed in the models, one can conclude that the difference between nucleon and nuclear structure functions should originate from properties of nucleon structure and structure of nucleon-nucleon interaction. Thus, from the one hand the nuclear effects should be defined by properties the $`n`$-nucleon Green functions ($`n=1,2\mathrm{}A`$), from the other hand detailed information about nuclear effects can provide new information about nucleon structure. In this paper the general properties of relativistic bound states, which presumably lead to the EMC-effect, are discussed. It is shown that the relativistic consideration allows one to establish new regularities in nuclear structure which are important in the construction of the nucleon parton distributions.
## 2 SPACE-TIME STRUCTURE OF RELATIVISTIC BOUND STATES
The amplitude of deep inelastic scattering is defined by the imaginary part of the forward Compton amplitude:
$`W_{\mu \nu }^A(P,q)=\mathrm{Im}_{q_0}i{\displaystyle d^4xe^{iqx}A,P|\mathrm{T}\left(J_\mu \left(x\right)J_\nu \left(0\right)\right)|A,P}`$
In the framework of the Bethe-Salpeter formalism the amplitude is expressed in terms of solutions of the homogeneous Bethe-Salpeter equation
$$\chi _{\alpha ,P}^A(𝒳)=𝑑𝒵𝑑𝒵^{}S_{(n)}(𝒳,𝒵)\overline{G}_{2n}(𝒵,𝒵^{})\chi _{\alpha ,P}^A(𝒵^{})$$
(1)
and the Mandelstam vertex $`\overline{G}_{2(n+1)}^{}{}_{\mu \nu }{}^{}(𝒵,x,𝒵^{})`$ :
$$A,P|\mathrm{T}(J_\mu (x)J_\nu (0))|A,P=𝑑𝒵𝑑𝒵^{}\overline{\chi }_{\alpha ,P}^A(𝒵)\overline{G}_{2(n+1)}^{}{}_{\mu \nu }{}^{}(𝒵,x,𝒵^{})\chi _{\alpha ,P}^A(𝒵^{}).$$
(2)
The calligraphic letters denote a set of nucleons positions in the four-dimensional space – $`𝒵=z_1,\mathrm{}z_n`$. Relative positions of the nucleons are defined by four-dimensional intervals – $`r_i=_j^nz_j/nz_i.`$ This is the most unusual feature of relativistic bound states that nucleons inside it are not only divided by three-dimensional intervals but also shifted in time. The shift is defined by the zero component of the interval – relative time $`\tau _i=r_{i}^{}{}_{0}{}^{}`$.
If one considers nucleons as three-dimensional objects then the shift in time looks like unphysical feature . One should assume in this case that the shift does not affect on observables, and different values of the variables $`\tau _i`$ in Eq. (1) should lead to equivalent quasi-potential approaches . Recently it was shown, however, that the approaches are not equivalent from the point of view of relativistic covariance . In general, the covariance can be kept only in equal-time approaches where $`\tau _i=0`$. This contradiction points to existence of observable effects which could result from the shift in time of bound nucleons. Existence of such effects can be natural within the hypothesis that bound nucleons are shifted in time four-dimensional objects.
The space-time diagram for the bound state of two nucleons, which are considered to be four-dimensional objects, can be represented as it is shown in Figure 1. The Fermi motion is disregarded for simplicity. The static structure shown in diagram (a) is homogeneous in time, while diagram (b) represents an inhomogeneous in time bound nucleon structure. If one makes an instant flash of the bound state at the moment $`t_0`$ one finds that the structure is defined by one nucleon at the moment $`t_1`$ in the future and by another one at the moment $`t_2`$ in the past. If the nucleon had the static structure then the three points $`t_0,t_1,t_2`$ in Figure 1(a) would be equivalent to each other and no effect would be observed from the shift in time. In another case, the points in Figure 1(b) are not equivalent to each other due to the evolution of the nucleon structure in time. The shift in time would then show up in modifications of the observables. From the stability of bound nucleons one can infer that the evolution of the nucleon structure in time has a periodic character. The period $`\theta `$ should be comparable to the mean value of relative time absolute $`\theta |\tau |`$. If $`|\tau |\theta `$ or $`|\tau |\theta `$ then no effect from the shift in time of bound nucleons can be measured.
The established relation between the structure of the nucleon in the future and the measurements at present looks as if the causality principle is violated. However, the quantum nature of the considered objects does not allow one to perform an instant measurement. Instead one have to consider the contribution of nucleon states averaged between the time moments $`t_1`$ and $`t_2`$. This suggests a qualitative interpretation of the modification of the bound nucleon structure. One can estimate from the uncertainty relation the size of the area in which the nucleon can be localized in 4-space. The radius $`\rho `$ of the area, “nucleon localization radius”, can be related to the nucleon mass as follows: $`\rho ^2=\mathrm{\Delta }t^2\mathrm{\Delta }x^21/m_\mathrm{N}^2.`$ The shift in time results in additional uncertainty $`|\tau |=|t_2t_1|`$ in definition of the time moment of the measurement: $`\stackrel{~}{\rho }^2=(\mathrm{\Delta }t+|\tau |)^2\mathrm{\Delta }x^21/m_{\mathrm{N}}^{}{}_{}{}^{2}.`$ Since $`|\tau |>0`$ by definition, the localization radius becomes larger $`\stackrel{~}{\rho }>\rho `$ and effective mass of the nucleon becomes smaller $`m^{}<m`$. This is similar to the hypothesis of the $`x`$-rescaling model , where the Bjorken variable $`x`$ of a bound nucleon is rescaled $`\stackrel{~}{x}=x/(1ϵ/m_\mathrm{N})`$ due to the shift of nucleon mass by the rescaling parameter $`ϵ`$: $`m^{}=mϵ`$. The $`x`$-rescaling model gives qualitative description of the deviation of nuclear to nucleon structure function ratio from unity (EMC – effect). Thus we can conclude that the shift in time manifests itself in observables, the EMC effect being one of the examples.
This conclusion has been justified in general terms in the publications where the nuclear structure functions were calculated in the framework of the covariant approach based on the Bethe-Salpeter formalism. In the framework of this approach the relative time dependence has been consistently taken into account with the help of analytical properties of nucleon Green functions. The developed approach allows one to consider the theory with equal-time bound nucleons as the leading approximation in description of relativistic bound states, where the relativistic corrections, connected with the relative time dependence, are disregarded. In order to take the corrections into account, one has to keep the dependence on the relative time through the entire chain of calculation of observables. As a result, the nuclear structure function has been obtained in the form:
$$F_2^A(x_A)=\frac{d^3k}{(2\pi )^3}\underset{a,a^{}}{\overset{A1}{}}\left[\frac{E_ak_3}{E_a}F_2^a(x_a)+\frac{\mathrm{\Delta }_{a,a^{}}^A}{E_a}x_a\frac{dF_2^a(x_a)}{dx_a}\right]\mathrm{\Phi }_{a,a^{}}^A(𝐤)^2,$$
(3)
where the Bjorken variables for the nucleus $`A`$ and the nuclear fragment $`a`$ are introduced in the form $`x_\text{A}=Q^2/(2P_Aq)`$ and $`x_a=Q^2/(2p_aq)`$. The function $`\mathrm{\Phi }_{a,a^{}}^A(𝐤)^2`$ defines distribution of the nuclear fragment $`a`$ ($`a=\mathrm{N},\mathrm{D},^3\mathrm{He},\mathrm{}`$) in field of the spectator system $`a^{}`$ ($`a^{}=\mathrm{N},\mathrm{NN},\mathrm{D},\mathrm{DN},^3\mathrm{He},\mathrm{}`$). The energy of the nuclear fragment is denoted as $`E_a=\sqrt{M_a^2+𝐤^2}`$, $`M_a`$ is mass of the fragment, the coefficients $`\mathrm{\Delta }_{\mathrm{a},\mathrm{a}^{}}^A=M_A+E_\mathrm{a}+E_a^{}`$ can be interpreted as the removal energy of the corresponding nuclear fragment. The term with derivative of the nucleon structure function in Eq.(3) results from the shift in time of bound nucleons . It is clear that the result can be rewritten in the form which one obtains within the $`x`$-rescaling model:
$`F_2^A(x_A)={\displaystyle 𝑑y𝑑ϵF_2^\mathrm{N}\left(\frac{x_A}{yϵ/M_A}\right)\frac{d^3k}{(2\pi )^3}\frac{my}{E_\mathrm{N}}\delta \left(y\frac{E_\mathrm{N}k_3}{m}\right)\underset{\mathrm{a}^{}}{}\mathrm{\Phi }_{\mathrm{N},\mathrm{a}^{}}^A(𝐤)\delta \left(ϵ\mathrm{\Delta }_{\mathrm{N},a^{}}^A\right)}.`$
The term with derivative in Eq.(3) leads to the depletion of the deuteron to nucleon structure functions ratio from unity in numerical calculations. The evaluated ratio $`F_2^A(x)/F_2^\mathrm{D}(x)`$ is in good agreement with the data available for $`A=4`$. When evaluated at $`A=3`$, the ratio offers the prediction for the experiments with $`{}_{}{}^{3}\mathrm{He}`$, $`{}_{}{}^{3}\mathrm{H}`$ and D targets. On the basis on this prediction the two-stage conception of $`A`$-dependence for the nuclear structure functions was proposed. At the first stage ($`1A4`$) the pattern of the $`F_2^A/F_2^\mathrm{N}`$ changes due to the relative time effects. At the second stage the pattern does not change but the amplitude of ratio oscillation around unity grows due to nuclear density evolution . The calculations based on this conception give good description of heavy nuclear data for the ratio .
## 3 CONCLUSION
To sum up the discussion presented in this paper, the following conclusions can be made. The proposed qualitative picture shows that the shift in time of bound nucleons exposes dynamical properties of a nucleon structure. Both the clear connection between the relativistic picture and the $`x`$-rescaling model and successful description of existent EMC effect data proves that the EMC effect is manifestation of the shift in time of the bound nucleons. The relativistic calculations show that the observable effects of the shift in time are defined by the derivative with respect to $`x`$ of the nucleon structure functions, which, therefore, reflect the dynamical properties of a nucleon structure. Thus, precision data on nuclear to deuteron structure functions ratio will provide information about the derivative, what is important in construction of nucleon parton distributions.
I would like to thank A.M. Baldin, V.V. Burov, C. Ciofi degli Atti, V.A. Nikolaev and G.I. Smirnov for fruitful discussions. I am grateful to organizers of the meeting for warm hospitality and support.
|
warning/0003/astro-ph0003351.html
|
ar5iv
|
text
|
# STELLAR STRUCTURE AND EVOLUTION: DEDUCTIONS FROM HIPPARCOS
## 1 INTRODUCTION
Stars are the main constituents of the observable Universe. The temperatures and pressures deep in their interiors are out of reach for the observer, while the description of stellar plasmas requires extensive knowledge in various domains of modern physics such as nuclear and particle physics, atomic and molecular physics, thermo- and hydrodynamics, physics of the radiation and of its interaction with matter, and radiative transfer. The development of numerical codes to calculate models of stellar structure and evolution began more than forty years ago with the pioneering works of Schwarzschild (1958) and Henyey et al (1959). These programs have allowed at least the qualitative study and understanding of numerous physical processes that intervene during the various stages of stellar formation and evolution.
During the last two decades, observational data of increasingly high accuracy have been obtained as a result of 1) the coming of modern ground-based or space telescopes equipped with high-quality instrumentation and with detectors giving access to almost any possible range of wavelengths and 2) the elaboration of various sophisticated techniques of data reduction. Ground-based astrometry has progressed, while space astrometry was initiated with Hipparcos. In the meantime, CCD detectors on large telescopes opened the era of high-resolution, high signal-to-noise ratio spectroscopy while multi-color filters were designed for photometry. New fields have appeared or are under development, such as helio- and asteroseismology or interferometry. On the other hand, stellar models have been enriched by a continuously improved physical description of the stellar plasma, while the use of increasingly powerful computers has led to a gain in numerical accuracy.
The confrontation of models with observations allows testing and even validation of the input physics of the models if numerous observations of high quality are available. Fundamental returns are expected in many domains that make use of quantitative results of the stellar evolution theory such as stellar, Galactic, and extragalactic astrophysics as well as cosmology. Because of their positions, movements, or interactions with the interstellar medium stars are actors and tracers of the dynamical and chemical evolution of the Galaxy. Astrophysicists aim to determine their ages and chemical compositions precisely. For example, the firm determination of the ages of the oldest stars, halo stars or members of globular clusters, is a long-standing objective because it is one of the strongest constraints for cosmology.
Although great progress has been made, a number of observations cannot be reproduced by stellar models, which raises many questions regarding both the observations and the models. In the last few years, two scientific meetings have been explicitly devoted to unsolved problems in stellar structure and evolution (Noels et al 1995, Livio 2000). A major point of concern is that of transport processes at work in stellar interiors (transport of the chemical elements, angular momentum or magnetic fields by microscopic diffusion and/or macroscopic motions). Observations show that transport processes are indeed playing a role in stellar evolution but many aspects remain unclear (sometimes even unknown) and need to be better characterized. Another crucial point concerns the atmospheres, which link the stellar interior model to the interstellar medium and are the intermediate agent between the star and the observer. Uncertainties and inconsistencies in atmospheric descriptions generate errors in the analysis of observational data and in model predictions.
This paper is the third of the series in ARAA dedicated to the results of the Hipparcos mission; Kovalevsky (1998) presented the products of the mission and the very first astrophysical results obtained immediately after the release of the data, while Reid (1999) reviewed the implications of the Hipparcos parallaxes for the location of the main sequence (MS) in the Hertzsprung-Russell (H-R) diagram, the luminosity calibration of primary distance indicators, and the Galactic distance scale. Also, van Leeuwen (1997) presented the results of the mission, and Baglin (1999) and Lebreton (2000) discussed the impact of Hipparcos data on stellar structure and evolution.
Hipparcos has provided opportunities to study rather large and homogeneous samples of stars sharing similar properties, for instance, in terms of their space location or chemical composition. I review studies based on Hipparcos observations which (1) confirmed several elements of stellar internal structure theory, (2) revealed some problems related to the development of stellar models, and (3) yielded more precise characteristics of individual stars and clusters. In Sections 2 and 3, I discuss the recent observational (including Hipparcos) and theoretical developments from which new studies could be undertaken. In Section 4, I concentrate on the nearest stars, observed with highest precision (A-K disk and halo single or binary field stars, and members of open clusters). In Section 5, I review recent results on variable stars, globular clusters and white dwarfs based on Hipparcos data. The stars considered are mostly of low or intermediate mass, and except for white dwarfs, the evolutionary stages cover the main sequence and subgiant branch. Throughout this paper, I emphasize that the smaller error bars on distances that result from Hipparcos make the uncertainties on the other fundamental stellar parameters more evident; fluxes, effective temperatures, abundances, gravities, masses and radii have to be improved correspondingly, implying in many cases the need for progress in atmospheric description.
## 2 NEW HIGH-ACCURACY OBSERVATIONAL MATERIAL
This section presents a brief review of Hipparcos results and complementary ground-based or space observations which, if combined, provide very homogeneous and precise sets of data.
### 2.1 Space Astrometry with Hipparcos
The Hipparcos satellite designed by the European Space Agency was launched in 1989. The mission ended in 1993 and was followed by 3 years of data reduction. The contents of the Hipparcos Catalogue (Eur. Space Agency 1997) were described by Perryman et al (1997a). The data were released to the astrophysical community in June 1997. General information on the mission is given in van Leeuwen’s (1997) and Kovalevsky’s (1998) review papers.
Stars of various masses, chemical compositions and evolutionary stages located either in the Galactic disk or in the halo were observed; this was done systematically to a V-magnitude that depends on the galactic latitude and spectral type of the star, and more generally, with a limit of V$``$12.4 mag. The Hipparcos Catalogue lists positions, proper motions, and trigonometric parallaxes of 117 955 stars as well as the intermediate astrometric data, from which the astrometric solutions were derived; this allows alternative solutions for the astrometric parameters to be reconstructed according to different hypotheses (see van Leeuwen & Evans 1998).
A total of 12 195 double or multiple systems are resolved (among which about 25 % were previously classified as single stars), and 8 542 additional stars are suspected to be non-single. Detailed information on multiple systems, as described by Lindegren et al (1997), can be found in The Double and Multiple System Annex of the Catalogue.
The median accuracy on positions and parallaxes ($`\pi `$) is typically $``$1 milliarcsecond (1 mas), whereas precisions on proper motions are about 1 mas per yr. Precisions become much higher for bright stars and worsen toward the ecliptic plane and for fainter stars. The astrometric accuracy and formal precision of Hipparcos data have been investigated by Arenou et al (1995) and Lindegren (1995), and discussed by van Leeuwen (1999a): for the Catalogue as a whole, the zero-point error on parallaxes is below 0.1 mas and the formal errors are not underestimated by more than 10%. After Hipparcos about 5 200 single stars and 450 double (or multiple) stars have parallaxes known with an accuracy $`\sigma _\pi /\pi `$ better than 5%, 20 853 stars have $`\sigma _\pi /\pi `$ lower than 10% and 49 333 stars have $`\sigma _\pi /\pi `$ lower than 20% (Mignard 1997). Martin et al (1997, 1998) and Martin & Mignard (1998) determined the masses of 74 astrometric binaries with accuracies in the range 5-35%. S derhjelm (1999) obtained masses and improved orbital elements for 205 visual binaries from a combination of Hipparcos astrometry and ground-based observations; among these, 12 (20) systems have mass-errors below 5 (7.5)%.
The Hipparcos Catalogue also includes detailed and homogeneous photometric data for each star, obtained from an average number of 110 observations per star. The broad-band Hipparcos (Hp) magnitude corresponding to the specific passband of the instrument spanning the wavelength interval $``$ 350-800 nanometers (see Figure 1 in van Leeuwen et al 1997), is provided with a median precision of 0.0015 mag for Hp$`<`$9 mag. The Johnson V magnitude derived from combined satellite and ground-based observations is given with a typical accuracy of 0.01 mag. The star mapper Tycho had passbands close to the Johnson B and V bands and provided two-color $`\mathrm{B}_\mathrm{T}`$ and $`\mathrm{V}_\mathrm{T}`$ magnitudes (accuracies are typically 0.014 mag and 0.012 mag for stars with $`\mathrm{V}_\mathrm{T}<9`$).
Hipparcos provided a detailed variability classification of stars (van Leeuwen et al 1997), resulting in 11 597 variable or possibly variable stars. Among these 2 712 stars are periodic variables (970 new) including 273 Cepheids, 186 RR Lyrae, 108 $`\delta `$ Scuti or SX Phoenicis stars, and 917 eclipsing binaries.
Hipparcos was planned more than fifteen years ago, and while its development proceeded, significant progress was made in the derivation of ground-based parallaxes using CCD detectors. Parallaxes with errors less than 1.4 mas have already been obtained for a few tens of stars, and errors are expected to drop to $`\pm `$0.5 mas in the years to come (Harris et al 1997, Gatewood et al 1998). In addition, the Hubble Space Telescope (HST ) Fine Guidance Sensor observations can provide parallaxes down to V$``$15.8 with errors at the 1 mas level (Benedict et al 1994, Harrison et al 1999). However, the distances of a rather small number of stars will be measured by HST because of the limited observing time available for astrometry. The enormous advantage of Hipparcos resides in the large number of stars it dealt with, providing homogeneous trigonometric parallaxes that are essentially absolute.
### 2.2 Ground-Based Photometry and Spectroscopy
The fundamental stellar parameters (bolometric magnitude $`\mathrm{M}_{\mathrm{bol}}`$, effective temperature $`\mathrm{T}_{\mathrm{eff}}`$, surface gravity $`g`$, and chemical composition) can be determined from photometry and/or from detailed spectroscopic analysis. However, the determination largely relies on model atmospheres and sometimes uses results of interior models. Direct masses and radii can be obtained for stars belonging to binary or multiple systems. Interferometry combined with distances yields stellar diameters giving direct access to $`\mathrm{T}_{\mathrm{eff}}`$, but still for a very limited number of rather bright stars which then serve to calibrate other methods. The different methods (and related uncertainties) used to determine the fundamental stellar parameters mainly for A to K Galactic dwarfs and subgiants are briefly discussed, and improvements brought by Hipparcos are underlined.
BOLOMETRIC MAGNITUDES. Integration of UBVRIJHKL photometry gives access to the bolometric flux on Earth $`\mathrm{F}_{\mathrm{bol}}`$, at least for F-G-K stars where most energy is emitted in those bands and which are close enough not to be affected by interstellar absorption; the (small) residual flux, emitted outside the bands, is estimated from model atmospheres. Recently, Alonso et al (1995) applied the method to $``$100 F-K dwarfs and subdwarfs and obtained bolometric fluxes accurate to about 2% and, as a by-product, empirical bolometric corrections for MS stars.
If $`\mathrm{F}_{\mathrm{bol}}`$ and distances are known, $`\mathrm{M}_{\mathrm{bol}}`$ can be derived with no need for bolometric correction. The accuracy is then $`\sigma _{\mathrm{M}_{\mathrm{bol}}}=\mathrm{log}e[(2.5\frac{\sigma _{\mathrm{F}_{\mathrm{bol}}}}{\mathrm{F}_{\mathrm{bol}}})^2+(5\frac{\sigma _\pi }{\pi })^2]^{\frac{1}{2}}`$, meaning that if $`\frac{\sigma _{\mathrm{F}_{\mathrm{bol}}}}{\mathrm{F}_{\mathrm{bol}}}`$2% then $`\sigma _{\mathrm{M}_{\mathrm{bol}}}`$ is dominated by the parallax error as soon as $`\frac{\sigma _\pi }{\pi }>`$1%. In other cases, when the distance is known, $`\mathrm{M}_{\mathrm{bol}}`$ is obtained from any apparent magnitude $`m`$ and its corresponding bolometric correction BC($`m`$), derived from empirical calibrations or from model atmospheres. Up to now Hipparcos magnitudes Hp have not been used extensively, despite their excellent accuracy (0.0015 mag), because of remaining difficulties in calculating BC(Hp) (Cayrel et al 1997a).
EFFECTIVE TEMPERATURES. The InfraRed Flux Method (IRFM; Blackwell et al 1990), applicable to A-K stars proceeds in two steps. First, the stellar angular diameter $`\varphi `$ is evaluated by comparing the IR flux observed on Earth in a given band to the flux predicted by a model atmosphere calculated with the observed gravity and abundances and an approximate $`\mathrm{T}_{\mathrm{eff}}`$ (the IR flux does not depend sensitively on $`\mathrm{T}_{\mathrm{eff}}`$). Then $`\mathrm{T}_{\mathrm{eff}}`$ is obtained from the total (integrated) flux $`\mathrm{F}_{\mathrm{bol}}`$ and $`\varphi `$. Iteration of the procedure yields a “definite” value of $`\mathrm{T}_{\mathrm{eff}}`$. Using IRFM, Alonso et al (1996a) derived temperatures of 475 F0-K5 stars ($`\mathrm{T}_{\mathrm{eff}}`$ in the range 4000-8000K) with internal accuracies of $``$1.5%. The zero-point of their $`\mathrm{T}_{\mathrm{eff}}`$-scale is based on direct interferometric measures by Code et al (1976), and the resulting systematic uncertainty is $``$1%. Accuracies of $``$1% were obtained by Blackwell & Linas-Gray (1998), who applied IRFM to 420 A0-K3 stars, corrected for interstellar extinction using Hipparcos parallaxes. Both sets of results compare well, with differences below 0.12$`\pm `$1.25% for the 93 stars in common.
The surface brightness method (Barnes et al 1978) was applied by Di Benedetto (1998) to obtain a ($`\mathrm{T}_{\mathrm{eff}}`$, V-K) calibration. The calibration is based on 327 stars with high-precision K-magnitudes from the Infrared Space Observatory (ISO), Hipparcos V-magnitudes and parallaxes (the latter to correct for interstellar extinction), and bolometric fluxes from Blackwell & Linas-Gray (1998). First, the visual surface brightness $`S_V=\mathrm{V}+5\mathrm{log}\varphi `$ is calibrated as a function of (V-K) using stars with precise $`\varphi `$ from interferometry. Then for any star $`S_V`$ is obtained from (V-K), $`\varphi `$ from $`S_V`$ and V, yielding in turn $`\mathrm{T}_{\mathrm{eff}}`$ from $`\mathrm{F}_{\mathrm{bol}}`$ and $`\varphi `$. From the resulting ($`\mathrm{T}_{\mathrm{eff}}`$, V-K) calibration, Di Benedetto derived $`\mathrm{T}_{\mathrm{eff}}`$ values of 537 ISO A-K dwarfs and giants with $`\pm `$1% accuracy. The method produces results in good agreement with those of IRFM and is less dependent on atmosphere models.
Multiparametric empirical calibrations of $`\mathrm{T}_{\mathrm{eff}}`$ as a function of the color indices and eventually of metallicity $`[\mathrm{Fe}/\mathrm{H}]`$ (logarithm of the number abundances of Fe to H relative to the solar value) and gravity can be derived from the empirical determinations of the effective temperatures of the rather nearby stars. In turn, the effective temperature of any star lying in the (rather narrow) region of the H-R diagram covered by a given calibration can easily be derived (see for example Alonso et al 1996b). Empirical calibrations also serve to validate purely theoretical calibrations based on model atmospheres; these latter have the advantage of covering the entire parameter space of the H-R diagram (i.e. wide ranges of color indices, metallicities and gravities; see Section 3.2 later in this article).
Spectroscopic determination of $`\mathrm{T}_{\mathrm{eff}}`$ is based on the analysis of chosen spectral lines that are sensitive to temperature; for instance the Balmer lines for stars with $`\mathrm{T}_{\mathrm{eff}}`$ in the interval 5000-8000 K. Because of the present high quality of the stellar spectra, precisions of $`\pm `$50-80 K on $`\mathrm{T}_{\mathrm{eff}}`$, that correspond to the adjustment of the theoretical line profile to the observed one, are commonly found in the literature (Cayrel de Strobel et al 1997b, Fuhrmann 1998). This supposes that theoretical profiles are very accurate, and therefore neglects the model atmosphere uncertainties.
Popper (1998) used detached eclipsing binaries with rather good Hipparcos parallaxes, accurate radii, and measured V-flux to calibrate the radiative flux as a function of (B-V); he found good agreement with similar calibrations based on interferometric angular diameters. From the same data, Ribas et al (1998) derived effective temperatures (this required bolometric corrections) and found them to be in reasonable agreement (although systematically smaller by 2-3%) with $`\mathrm{T}_{\mathrm{eff}}`$ derived from photometric calibrations. However the stars are rather distant, which implies rather significant internal errors on $`\mathrm{M}_{\mathrm{bol}}`$ and $`\mathrm{T}_{\mathrm{eff}}`$ (a parallax error of 10% is alone responsible for a $`\mathrm{T}_{\mathrm{eff}}`$-error of 5%). In Ribas et al’s sample, only a few systems have $`\sigma _\pi /\pi `$ $`<`$10%, and because errors on radius, magnitudes, and BC also intervene, only 5 systems have $`\mathrm{T}_{\mathrm{eff}}`$ determined to better than 3%.
SURFACE GRAVITIES. If $`\mathrm{T}_{\mathrm{eff}}`$ and $`\mathrm{M}_{\mathrm{bol}}`$ are known, the radius of the star may be derived from the Stefan-Boltzmann law and the mass estimated from a grid of stellar evolutionary models, yielding in turn the value of $`g`$. This method has been applied to a hundred metal-poor subdwarfs and subgiants with accurate distances from Hipparcos (Nissen et al 1997, Fuhrmann 1998, Clementini et al 1999). Nissen et al showed that among the various sources of errors, the error on distance still dominates, but pointed out that if the distance error is lower than 20% then the error on $`\mathrm{log}g`$ may be lower than $`\pm `$ 0.20 dex.
On the other hand, $`g`$ can be determined from spectroscopy. Different gravities produce different atmospheric pressures , modifying the profiles of some spectral lines. Two methods have been widely used to estimate $`g`$. The first method is based on the analysis of the equation of ionization equilibrium of abundant species, iron, for instance. The iron abundance is determined from FeI lines that are not sensitive to gravity, and then $`g`$ is adjusted so that the analysis of FeII lines, which are sensitive to gravity, leads to the same value of the iron abundance. The accuracy in $`\mathrm{log}g`$ is in the range $`\pm `$0.1-0.2 dex (Axer et al 1994). The second method relies on the analysis of the wings of strong lines broadened by collisional damping, such as Ca I (Cayrel et al 1996) or the Mg Ib triplet (Fuhrmann et al 1997), leading to uncertainties smaller than 0.15 dex. The two methods often produce quite different results, with systematic differences of $``$0.2-0.4 dex, at least when ionization equilibria are estimated from models in local thermodynamical equilibrium (LTE).
Thévenin & Idiart (1999) have studied the effects of departures from LTE on the formation of FeI and FeII lines in stellar atmospheres, and found that modifications of the ionization equilibria resulted from the overionization of iron induced by significant UV fluxes. The nice consequence is that the gravities they inferred from iron ionization equilibrium for a sample of 136 stars spanning a large range of metallicities become very close to gravities derived either from pressure-broadened strong lines or through Hipparcos parallaxes.
ABUNDANCES OF THE CHEMICAL ELEMENTS. The spectroscopic determination of abundances of chemical elements rests on the comparison of the outputs of model atmospheres (synthetic spectra, equivalent widths) with their counterpart in the observed spectra. This requires a preliminary estimate of $`\mathrm{T}_{\mathrm{eff}}`$ and $`g`$. If high-resolution spectra are used, the line widths are very precise and the internal uncertainty in abundance determinations depends on uncertainties in $`g`$ and $`\mathrm{T}_{\mathrm{eff}}`$, on the validity of the model atmosphere, and on the oscillator strengths. Error bars in the range $`\pm `$0.05-0.15 dex are typical (Cayrel de Strobel et al 1997b, Fuhrmann 1998). Also when different sets of $`[\mathrm{Fe}/\mathrm{H}]`$ determinations are compared, the solar Fe/H ratio used as reference must be considered; values differing by $``$0.15 dex used to be found in the literature (Axer et al 1994). This has resulted in long-standing difficulties in determining the solar iron abundance from FeI or FeII lines, because of uncertain atomic data. In a recent paper, Grevesse & Sauval (1999) reviewed the problem and opted for a “low” Fe-value, $`\mathrm{A}_{\mathrm{Fe}}=7.50\pm 0.05`$ ($`\mathrm{A}_{\mathrm{Fe}}=\mathrm{log}(\mathrm{n}_{\mathrm{Fe}}/\mathrm{n}_\mathrm{H})+12`$ is the logarithm of the number density ratio of Fe to H particles), in perfect agreement with the meteoritic value.
Furthermore, if abundances are estimated from model atmospheres in LTE, perturbations of statistical equilibrium by the radiation field are neglected. Thévenin & Idiart (1999) found that in metal-deficient dwarfs and subgiants, the iron overionization resulting from reinforced UV flux modifies the line widths. They obtained differential non-LTE/LTE abundance corrections increasing from 0.0 dex at $`[\mathrm{Fe}/\mathrm{H}]`$=0.0 to +0.3 dex at $`[\mathrm{Fe}/\mathrm{H}]`$= -3.0. These corrections are indeed supported by the agreement between spectroscopic gravities and “Hipparcos” gravities discussed previously.
Helium lines do not form in the photosphere of low-temperature stars which precludes a direct determination of helium abundance. The calibration of the solar model in luminosity and radius at solar age yields the initial helium content of the Sun (Christensen-Dalsgaard 1982), while oscillation frequencies give access to the present value in the convection zone (Kosovichev et al 1992). In other stars, it is common to use the well-known scaling relation $`YY_p=Z\frac{\mathrm{\Delta }Y}{\mathrm{\Delta }Z}`$, which supposes that the helium abundance has grown with metallicity $`Z`$ from the primordial value $`Y_p`$ to its stellar birth value $`Y`$ ($`Y`$ and $`Z`$ represent abundances in mass fraction); $`\mathrm{\Delta }Y/\mathrm{\Delta }Z`$ is the enrichment factor.
$`\alpha `$-element abundances (O, Ne, Mg, Si, S, Ar, Ca, Ti) have now been widely measured in metal-deficient stars. Stars with $`[\mathrm{Fe}/\mathrm{H}]`$$``$-0.5 dex generally exhibit an $`\alpha `$-element enhancement with respect to the Sun ($`[\alpha /\mathrm{Fe}]`$) quite independent of their metallicity (Wheeler et al 1989, Mc William 1997). Recent determinations of $`[\alpha /\mathrm{Fe}]`$ in 99 dwarfs with $`[\mathrm{Fe}/\mathrm{H}]`$$`<`$-0.5 from high-resolution spectra by Clementini et al (1999) yield $`[\alpha /\mathrm{Fe}]`$=+0.26$`\pm `$0.08 dex.
## 3 RECENT THEORETICAL AND NUMERICAL PROGRESS
Recent developments in the physical description of low and intermediate mass stars are briefly presented.
### 3.1 Microscopic Physics
The understanding of stellar structure benefited substantially from the complete re-examination of stellar opacities by two groups: the Opacity Project (OP, see Seaton et al 1994) and the OPAL group at Livermore (see Rogers & Iglesias 1992). Both showed by adopting different and independent approaches, that improved atomic physics lead to opacities generally higher than the previously almost “universally” used Los Alamos opacities (Huebner et al 1977). The opacity enhancements reach factors of 2-3 in stellar envelopes with temperatures in the range $`10^510^6`$ K. With these new opacities, (1) a number of long-standing problems in stellar evolution have been solved or at least lessened and (2) finer tests of stellar structure could be undertaken. Since opacity is very sensitive to metallicity, any underlying uncertainty on metallicity may be problematic.
Great efforts have also been invested in the derivation of low-temperature opacities, including millions of molecular and atomic lines and grain absorption that are fundamental for the calculation of the envelopes and atmospheres of cool stars (Kurucz 1991, Alexander & Ferguson 1994).
OP and OPAL opacities have been shown to be in reasonable agreement (Seaton et al 1994, Iglesias & Rogers 1996); and very good agreement between OPAL and Alexander & Ferguson’s or Kurucz’s opacities is found in the domains where they overlap. Although some uncertainties remain that are difficult to quantify, the largest discrepancies between the various sets of tables do not exceed 20% and are generally well understood (Iglesias & Rogers), making opacities much more reliable today than they were ten years ago.
The re-calculation of opacities required appropriate equations of state (EOS). The MH&D EOS (see Mihalas et al 1988) is part of the OP project, while the OPAL EOS was developed at Livermore (Rogers et al 1996). In the meantime, another EOS was designed to interpret the first observations of very low-mass stars and brown dwarfs (Saumon & Chabrier 1991). OPAL and OP EOS are needed to satisfy the strong helioseismic constraints (Christensen-Dalsgaard & D ppen 1992).
### 3.2 Atmospheres
Atmospheres intervene at many levels in the analysis of observations (Section 2.2). They also provide external boundary conditions for the calculation of stellar structure and necessary relations to transform theoretical ($`\mathrm{M}_{\mathrm{bol}}`$, $`\mathrm{T}_{\mathrm{eff}}`$) H-R diagrams to color-magnitude (C-M) or color-color planes. Models have improved during the last two decades, and attention has been paid to the treatment of atomic and molecular line blanketing. The original programs MARCS of Gustafsson et al (1975) and ATLAS by Kurucz (1979) evolved toward the most recent ATLAS9 version, appropriate for O-K stars (Kurucz 1993) and NMARCS for A-M stars (see Brett 1995 and Bessell et al 1998). On the other hand; very low-mass stellar atmosphere models were developed; Carbon (1979) and Allard et al (1997) reviewed calculation details and remaining problems (such as incomplete opacity data, poor treatment of convection, neglect of non-LTE effects or assumption of plane-parallel geometry).
COLOR-MAGNITUDE TRANSFORMATIONS. Different sets of transformations (empirical or theoretical) were used to analyze the “Hipparcos” stars. Empirical transformations have been discussed in Section 2.2. The most recent theoretical transformations are compiled by Bessell et al (1998), who used synthetic spectra derived from ATLAS9 and NMARCS to produce broad-band colors and bolometric corrections for a very wide range of $`\mathrm{T}_{\mathrm{eff}}`$, $`g`$ and \[Fe/H\] values. These authors found fairly good agreement with empirical relations except for the coolest stars (M dwarfs, K-M giants).
INTERIOR/ATMOSPHERE INTERFACE. The external boundary conditions for interior models are commonly obtained from $`\mathrm{T}(\tau )`$-laws ($`\tau `$ is the optical depth) derived either from theory or full atmosphere calculation. This method is suitable for low- and intermediate-mass stars (it is not valid for masses below $``$0.6 $`\mathrm{M}_{}`$, Chabrier & Baraffe 1997). Morel et al (1994) and Bernkopf (1998) focused on the solar case where seismic constraints require a careful handling of external boundary conditions. Morel et al pointed out that homogeneous physics should be used in interior and atmosphere (opacities, EOS, treatment of convection) and showed that the boundary level must be set deep enough, in zones where the diffusion approximation is valid. Bernkopf discussed some difficulties in reproducing Balmer lines related to the convection treatment.
### 3.3 Transport Processes
CONVECTION. 3-D numerical simulations at current numerical resolution are able to reproduce most observational features of solar convection such as images, spectra, and helioseismic properties (Stein & Nordlund 1998). However, the “connection” with a stellar evolution code is not easy, and stellar models still mostly rely on 1-D phenomenological descriptions such as the mixing-length theory of convection (MLT, Böhm-Vitense 1958). The mixing-length parameter $`\alpha _{\mathrm{MLT}}`$ (ratio of the mixing-length to the pressure scale height) is calibrated so that the solar model yields the observed solar radius at the present solar age. The question of the variations of $`\alpha _{\mathrm{MLT}}`$ in stars of various masses, metallicities, and evolutionary stages remains a matter of debate (Section 4). As pointed out by Abbett et al (1997), the MLT can reproduce the correct entropy jump across the superadiabatic layer near the stellar surface, but fails to describe the detailed depth structure and dynamics of convection zones. Abbett et al found that the solar entropy jump obtained in 3-D simulations corresponds to predictions of the MLT for $`\alpha _{\mathrm{MLT}}`$ $``$1.5. Ludwig et al (1999) calibrated $`\alpha _{\mathrm{MLT}}`$ from 2-D simulations of compressible convection in solar-type stars for a broad range of $`\mathrm{T}_{\mathrm{eff}}`$ and $`g`$-values. The solar $`\alpha _{\mathrm{MLT}}`$ inferred from 3-D and 2-D simulations is close to what is obtained in solar model calibration. The $`\alpha _{\mathrm{MLT}}`$ dependence with $`\mathrm{T}_{\mathrm{eff}}`$ and $`g`$ of Ludwig et al can be used to constrain the range of acceptable variations of $`\alpha _{\mathrm{MLT}}`$ in stellar models (see Section 4).
OVERSHOOTING. Penetration of convection and mixing beyond the classical Schwarzschild convection cores (overshooting process) modifies the standard evolution model of stars of masses M$``$1.2 $`\mathrm{M}_{}`$, in particular the lifetimes (see for instance Maeder & Mermilliod 1981, Bressan et al 1981). The extent of overshooting was estimated for the first time from the comparison of observed and theoretical MS widths of open clusters (Maeder & Mermilliod 1981), which yields an overshooting parameter $`\alpha _{\mathrm{ov}}`$ $``$0.2 (ratio of overshooting distance to pressure scale height). As discussed in detail by Roxburgh (1997), $`\alpha _{\mathrm{ov}}`$ is still poorly constrained despite significant efforts made to establish the dependence of overshooting with mass, evolutionary stage, or chemical composition (see Section 4). Andersen (1991) first pointed out that the simultaneous calibration of well-known binaries (masses and radii at 1-2%) may provide improved constraints for $`\alpha _{\mathrm{ov}}`$ . A modeling of the sample of the best-known binaries indicates a trend for $`\alpha _{\mathrm{ov}}`$ to increase with mass and suggests a decrease of $`\alpha _{\mathrm{ov}}`$ with decreasing metallicity (Ribas 1999), although a larger sample would be desirable to confirm those trends. Further advances are expected from asteroseismology (Brown et al 1994, Lebreton et al 1995).
DIFFUSION OF CHEMICAL ELEMENTS. Various mixing processes may occur in stellar radiative zones (see Pinsonneault 1997). In low-mass stars, microscopic diffusion due to gravitational settling carries helium and heavy elements down to the center and modifies the evolutionary course as well as the surface abundances. It has been proved that microscopic diffusion can explain the low helium abundance of the solar convective zone derived from seismology (Christensen-Dalsgaard et al 1993). On the other hand, turbulent mixing (resulting, for instance, from hydrodynamical instabilities related to rotation, see Zahn 1992) probably inhibits microscopic diffusion. Richard et al (1996) did not find any conflict between solar models including rotation-induced mixing (to account for Li and Be depletion at the surface) and microscopic diffusion (to account for helioseismic data). More constraints are required to clearly identify (and quantify the effects of) the various candidate mixing processes; this will be illustrated in the following sections.
## 4 STUDIES OF THE BEST-KNOWN OBJECTS
Stellar model results depend on a number of free input parameters. Some are observational data (mass, chemical composition and age, the latter for the Sun only), whereas others enter phenomenogical descriptions of poorly-known physical processes (mixing-length parameter for convection, overshooting, etc). The model outputs have to be compared with the best available observational data: luminosity, $`\mathrm{T}_{\mathrm{eff}}`$ or radius, oscillation frequencies, etc. Numerous and precise observational constraints allow assessment of the input physics or give more precise values of the free parameters. They may reveal the necessity to include processes previously neglected and in the best cases to characterize them.
The model validation rests on (1) the nearest objects with the most accurate observations, (2) special objects with additional information such as stars belonging to binary systems, members of stellar clusters or stars with seismic data, and (3) large samples of objects giving access to statistical studies.
### 4.1 Stars in Binary Systems
Masses are available for a number of stars belonging to binary systems, allowing their “calibration” under the reasonable assumption that the stars have the same age and were born with the same chemical composition (Andersen 1991, Noels et al 1991). A solution is sought which reproduces the observed positions in the H-R diagram of both stars. Andersen (1991) claimed that the only systems able to really constrain the internal structure theory are those with errors lower than 2% in mass, 1% in radius, 2% in $`\mathrm{T}_{\mathrm{eff}}`$ and 25% in metallicity.
However, additional observations may sometimes cast doubts on an observed quantity previously determined with good internal accuracy. This occurred recently for the masses of stars in the nearest visual binary system $`\alpha `$ Centauri. The system has been widely modeled in the past (Noels et al 1991, Edmonds et al 1992, Lydon et al 1993, Fernandes & Neuforge 1995) with the objective of getting (among others) constraints on the mixing-length parameter. At that time, the astrometric masses were used (internal error of 1%) but the $`[\mathrm{Fe}/\mathrm{H}]`$-value was controversial, leading to various possibilities for $`\alpha _{\mathrm{MLT}}`$ -values. Today the situation is still confused: metallicity is better assessed, but new radial velocity measurements yield masses higher than those derived from astrometry (by 6-7%, Pourbaix et al 1999). The higher masses imply a reduction in age by a factor of 2 and slightly different $`\alpha _{\mathrm{MLT}}`$ -values for the two stars. However, the orbital parallax corresponding to the high-mass “option” is smaller than and outside the error bars of both ground-based and Hipparcos parallax $`\pi _{\mathrm{Hipp}}`$. Pourbaix et al noted the lack of reliability of $`\pi _{\mathrm{Hipp}}`$ given in the Hipparcos Catalogue, but since then it has been re-determined from intermediate data by S derhjelm (1999) and is now close to (and in agreement with) the ground-based parallax. More accurate radial velocity measurements are therefore needed to assess the high-mass solution.
Possible variations of $`\alpha _{\mathrm{MLT}}`$ have been investigated through the simultaneous modeling of selected nearby visual binary systems (Fernandes et al 1998, Pourbaix et al 1999, Morel et al 2000). Small variations of $`\alpha _{\mathrm{MLT}}`$ (not greater than $``$0.2) in the two components of $`\alpha `$Cen (Pourbaix et al) and $`\iota `$ Peg (Morel et al) have been suggested. Fernandes et al, who calibrated 4 systems and the Sun with the same program and input physics, found that $`\alpha _{\mathrm{MLT}}`$ is almost constant for $`[\mathrm{Fe}/\mathrm{H}]`$ in the range $`[\mathrm{Fe}/\mathrm{H}]_{}`$$`\pm `$0.3 dex and masses between 0.6 and 1.3 $`\mathrm{M}_{}`$. In this mass range the sensitivity of models to $`\alpha _{\mathrm{MLT}}`$ increases with mass (due to the increase with mass of the entropy jump across the superadiabatic layer) which makes the MS slope vary with $`\alpha _{\mathrm{MLT}}`$ . Also, I estimate from my models that a change of $`\alpha _{\mathrm{MLT}}`$ of $`\pm 0.15`$ around 1 $`\mathrm{M}_{}`$ translates into a $`\mathrm{T}_{\mathrm{eff}}`$-change of $``$40-55 K depending on the metallicity. On the other hand, with the solar-$`\alpha _{\mathrm{MLT}}`$ value the MS slope of field stars and Hyades stars is well fitted (Section 4.3). It is therefore reasonable to adopt the solar-$`\alpha _{\mathrm{MLT}}`$ value to model solar-type stars. For other stars, the situation is less clear. The calibration of $`\alpha _{\mathrm{MLT}}`$ depends on the external boundary condition applied to the model, itself sensitive to the low-temperature opacities, and on the color transformation used for the comparison with observations. Chieffi et al (1995) examined the MS and red giant branch (RGB) in metal-deficient clusters and suggested a constancy of $`\alpha _{\mathrm{MLT}}`$ from MS to RGB and a decrease with decreasing $`Z`$. They found variations of $`\alpha _{\mathrm{MLT}}`$ with $`Z`$ of $``$ 0.2-0.4, but these are difficult to assess considering uncertainties in the observed and theoretical cluster sequences. On the other hand, calibration of $`\alpha _{\mathrm{MLT}}`$ with 2-D simulations of convection gives complex results (Freytag & Salaris 1999; Freytag et al 1999). In particular, (1) for solar metallicity, $`\alpha _{\mathrm{MLT}}`$ is found to decrease when $`\mathrm{T}_{\mathrm{eff}}`$ increases above solar $`\mathrm{T}_{\mathrm{eff}}`$, and to increase slightly when stars move toward the RGB (by $``$ 0.10-0.15) and; (2) $`\alpha _{\mathrm{MLT}}`$ does not vary importantly when metallicity decreases at solar $`\mathrm{T}_{\mathrm{eff}}`$. More work is needed to go into finer details, and other calibrators of $`\alpha _{\mathrm{MLT}}`$ are required, such as binary stars in the appropriate range of mass and with various chemical compositions.
The modeling of a moderately large sample of binaries might give information on the variation of helium $`Y`$ and age with metallicity $`Z`$, of great interest for Galactic evolution studies. The combined results for six binary systems and the Sun with the same program by Fernandes et al (1998) and Morel et al (2000) show a general trend for $`Y`$ to increase with $`Z`$: $`Y`$ increases from 0.25 to 0.30 ($`\pm `$0.02) when $`Z`$ increases from 0.007 to 0.03 ($`\pm `$0.002). However, the Hyades appear to depart from this tendency (see Section 4.3).
The sample of binaries with sufficiently accurate temperatures and abundances is still too meager to allow full characterization of physical processes. Additional data are needed such as observations of binaries in clusters (see Section 4.3) or asteroseismological measurements.
### 4.2 The Nearest Disk and Halo Stars
FINE STRUCTURE OF THE H-R DIAGRAM. Highly accurate distances for a rather large number of stars in the solar neighborhood were provided by Hipparcos. This allowed the first studies of the fine structure of the H-R diagram and related metallicity effects to be undertaken.
Among an ensemble of “Hipparcos” F-G-K stars closer than 25 pc, with error on parallax lower than 5%, Lebreton et al (1997b) selected stars with $`[\mathrm{Fe}/\mathrm{H}]`$ in the range $`[1.0,+0.3]`$ from detailed spectroscopic analysis ($`\sigma _{[\mathrm{Fe}/\mathrm{H}]}`$ $``$0.10 dex, Cayrel de Strobel et al 1997b), $`\mathrm{F}_{\mathrm{bol}}`$ and $`\mathrm{T}_{\mathrm{eff}}`$ from Alonso et al (1995, 1996a) with $`\frac{\sigma _{\mathrm{F}_{\mathrm{bol}}}}{\mathrm{F}_{\mathrm{bol}}}2\%`$ and $`\frac{\sigma _{\mathrm{T}_{\mathrm{eff}}}}{\mathrm{T}_{\mathrm{eff}}}1.5\%`$ (see Section 2.2) and not suspected to be unresolved binaries. Figure 1 presents the H-R diagram of the 34 selected stars: the error bars are the smallest obtained for stars in the solar neighborhood ($`\sigma _{\mathrm{M}_{\mathrm{bol}}}`$are in the range 0.031-0.095 with an average value $``$$`\sigma _{\mathrm{M}_{\mathrm{bol}}}`$$``$0.045 mag). The sample is compared with theoretical isochrones derived from standard stellar models in Figure 2. Models cover the entire $`[\mathrm{Fe}/\mathrm{H}]`$-range. They account for an $`\alpha `$-element enhancement $`[\alpha /\mathrm{Fe}]`$=+0.4 dex for $`[\mathrm{Fe}/\mathrm{H}]`$$``$-0.5 and, for non-solar $`[\mathrm{Fe}/\mathrm{H}]`$, have a solar-scaled helium content ($`Y`$=$`Y_\mathrm{p}`$+$`Z`$$`(\mathrm{\Delta }Y/\mathrm{\Delta }Z)_{}`$). The splitting of the sample into a solar metallicity sample and a moderately metal-deficient one (Figure 2a and b) shows that:
1. The slope of the MS is well reproduced with the solar $`\alpha _{\mathrm{MLT}}`$ ,
2. Stars of solar metallicity and close to it occupy the theoretical band corresponding to their (LTE) metallicity range, while for moderately metal deficient stars there is a poor fit.
In general, stars have a tendency to lie on a theoretical isochrone corresponding to a higher metallicity than the spectroscopic (LTE) value. This trend was already noticed by Axer et al (1995) but it is now even more apparent because of the high accuracy of the data. Helium content well below the primordial helium value would be required to resolve the conflict.
This is exemplified by the star $`\mu `$ Cas A, the A-component of a well-known, moderately metal-deficient binary system that has a well-determined mass (error in mass of 8 per cent). The standard model (Figure 3) is more than 200 K hotter than the observed point and is unable to reproduce the observed $`\mathrm{T}_{\mathrm{eff}}`$ even if (reasonable) error bars are considered (Lebreton 2000). On the other hand, the mass-luminosity properties of the star are well reproduced if the helium abundance is chosen to be close to the primordial value, although the error bar in mass is somewhat too large to provide strong constraints.
Several reasons can be invoked to explain the poor fit at low metallicities:
1. Erroneous temperature-scale. 3-D model atmospheres could still change the $`\mathrm{T}_{\mathrm{eff}}`$-scale as a function of metallicity (Gustafsson 1998), but with the presently (1-D) available models it seems difficult to increase Alonso et al’s (1996a) $`\mathrm{T}_{\mathrm{eff}}`$ by as much as 200-300 K. As noted by Nissen (1998), this scale is already higher than other photometric scales, by as much as 100 K. Also, Lebreton et al (1999) verified that spectroscopic effective temperatures lead to a similar misfit.
2. Erroneous metallicities. As discussed in Section 2.2, the $`[\mathrm{Fe}/\mathrm{H}]`$-values inferred from model atmosphere analysis should be corrected for non-LTE effects. According to Thévenin & Idiart (1999) no correction is expected at solar metallicity, whereas for moderately metal-deficient stars the correction amounts to $``$0.15 dex.
3. Inappropriate interior models. In low-mass stars, microscopic diffusion by gravitational settling can make helium and heavy elements sink toward the center, changing surface abundances as well as inner abundance profiles. In metal-deficient stars this process may be very efficient for three reasons: (1) densities at the bottom of the convection zone decrease with metallicity, which favors settling; (2) the thickness of the convection zones decreases with metallicity, making the reservoir easier to empty; and (3) metal-deficient generally means older, which implies more time available for diffusion.
The two latter reasons are attractive because they qualitatively predict an increasing deviation from the standard case when metallicity decreases. As shown in Figure 3, a combination of microscopic diffusion effects with non-LTE Fe/H corrections could remove the discrepancy noted for $`\mu `$Cas A: an increase of $`[\mathrm{Fe}/\mathrm{H}]`$ by 0.15 dex produces a rightward shift of 80 K of the standard isochrone, representing about one third of the discrepancy. Additionally, adding microscopic diffusion effects, according to recent calculations by Morel & Baglin (1999), provides a match to the observed positions. Moreover, the general agreement for solar metallicity stars (Figure 2a) should remain: (1) at solar metallicities non-LTE corrections are found to be negligible, and (2) at ages of $``$5 Gyr chosen as a mean age for those (expectedly) younger stars, diffusion effects are estimated to be smaller than the error bars on $`\mathrm{T}_{\mathrm{eff}}`$ (Lebreton et al 1999).
To conclude on this point, the high-level accuracy reached for a few tens of stars in the solar neighborhood definitely reveals imperfections in interior and atmosphere models. It casts doubts on abundances derived from model atmospheres in LTE, and favors models that include microscopic diffusion of helium and heavy elements toward the interior over standard models. Also, diffusion makes the surface $`[\mathrm{Fe}/\mathrm{H}]`$-ratio decrease by $``$0.10 dex in 10 Gyr in a star like $`\mu `$ Cas (Morel & Baglin 1999), which is rather small and hidden in the observational error bars. In very old, very deficient stars, the $`[\mathrm{Fe}/\mathrm{H}]`$-decrease is expected to be larger (Salaris et al 2000), which makes the relation between observed and initial abundances difficult to establish. In the future, progress will come from the study of enlarged samples reaching the same accuracies and of the acquisition of additional parameters to constrain the models. The knowledge of masses for several binaries in a narrow mass range but large metallicity range would help to constrain the helium abundances, while access to seismological data for at least one or two stars would help to better characterize mixing processes.
STATISTICAL STUDIES. Complete H-R diagrams of stars of the solar neighborhood have been constructed by adopting different selection criteria, and have been compared to synthetic H-R diagrams based on theoretical evolutionary tracks.
* Schröder (1998) proposed diagnostics of MS overshooting based on star counts in the different regions of the “Hipparcos” H-R diagram of stars in the solar neighborhood (d$`<`$50-100 pc). In the mass range 1.2-2 $`\mathrm{M}_{}`$, convective cores are small, and it is difficult to estimate the amount of overshooting with isochrone shapes. Schröder suggested using the number of stars in the Hertzsprung gap, associated with the onset of H-shell burning, as an indicator of the extent of overshooting around 1.6 $`\mathrm{M}_{}`$; the greater the overshooting on the MS, the larger the He-burning core, and in turn the longer the passage through the Hertzsprung gap. Actual star counts favor an onset of overshooting around $``$1.7 $`\mathrm{M}_{}`$(no overshooting appears necessary below that mass), which is broadly consistent with other empirical calibrations (MS width, eclipsing binaries), but finer quantitative estimates would require more accurate observational parameters, mainly in $`\mathrm{T}_{\mathrm{eff}}`$ and $`Z`$.
* Jimenez et al (1998) compared the red envelope of “Hipparcos” subgiants ($`\sigma _\pi /\pi `$ $`<`$0.15, $`\sigma _{(\mathrm{B}\mathrm{V})}<`$0.02 mag) with isochrones to determine a minimum age of the Galactic disk of 8 Gyr, which is broadly consistent with ages obtained with other methods (white-dwarf cooling curves, radioactive dating, isochrones, or fits of various age-sensitive features in the H-R diagram). The fit is still qualitative: the metallicities of subgiants are unknown because of the inadequacy of model atmospheres in that region. For this reason, Jimenez et al investigated the isochrones in other regions, MS and clump (He core burning). They calculated the variations with mass of the clump position for a range of metallicities in the disk, and showed that stars with masses from 0.8 to 1.3 $`\mathrm{M}_{}`$(ages from 2 to 16 Gyr) all occupy a well-defined vertical branch, the red-edge of the clump. The color of this border line is sensitive to metallicity, which makes it a good metallicity indicator in old metal-rich populations.
* Ng & Bertelli (1998) revised the ages of stars of the solar neighborhood and derived corresponding age-metallicity and age-mass relations. Fuhrmann (1998) combined the \[Mg/H\]-\[Fe/H\] relation with age and kinematical information to distinguish thin and thick disk stars. Several features seem to emerge from these studies: (1) no evident age-metallicity relation exists for the youngest ($`<8`$ Gyr) thin-disk stars; some of them are rather metal-poor, and super metal-rich stars appear to have been formed early in the history of the thin disk; (2) there is an apparent lack of stars in the age-interval 10-12 Gyr which is interpreted by Fuhrmann as a signature of the thin-disk formation; and (3) beyond 12 Gyr there is a slight decrease of metallicity with increasing age for stars of the thick disk; some of them are as old as halo stars. To assess these suggestions and to assist progress in the understanding of the Galactic evolution scenario (see Fuhrmann 1998 for details), enlarged stellar samples and further improvements on age determinations are of course required.
THE SUBDWARF/SUBGIANT SEQUENCE. Hipparcos provided the very first high-quality parallaxes for a number of halo stars. Age determinations of the local halo could be undertaken, as well as comparisons with globular cluster sequences.
Among a large sample of Population II Hipparcos halo subdwarfs, Cayrel et al (1997b) extracted the best-known stars with criteria similar to those adopted by Lebreton et al (1999) for disk stars. Stars were corrected for reddening, excluding stars with E(B-V)$`>`$0.05. Prior to Hipparcos, only 5 halo stars had parallax errors smaller than 10%; now there are 17, which represents sizeable progress. The halo stars are plotted in Figure 4; subdwarfs but also subgiants are present, delineating an isochrone-like shape with a turn-off region.
To make a first estimate of the age of the local halo, Cayrel et al kept 13 stars with the lowest error bars and spanning a narrow metallicity range ($`[\mathrm{Fe}/\mathrm{H}]`$=-1.5$`\pm `$0.3), the most commonly found in the halo (Figure 5). They found that halo stars, like disk stars, are colder than the theoretical isochrone corresponding to their metallicity. The misfit was also noted by Nissen et al (1997) and Pont et al (1997) in larger samples of halo stars. The discrepancy amounts to 130 to 250 K depending on the metallicity, and comparisons indicate that it is independent of the particular set of isochrones used. Again, non-LTE corrections leading to increased $`[\mathrm{Fe}/\mathrm{H}]`$-values ($`\mathrm{\Delta }`$$`[\mathrm{Fe}/\mathrm{H}]`$=+0.2 for $`[\mathrm{Fe}/\mathrm{H}]`$$``$-1.5 according to Thévenin & Idiart 1999), added to the effects of microscopic diffusion, can be invoked to reduce the misfit. Figure 5a compares Cayrel et al’s sample with standard isochrones by PA Bergbusch & DA VandenBerg (2000, in preparation), showing that the subdwarf main sequence cannot be reproduced by isochrones computed with the LTE $`[\mathrm{Fe}/\mathrm{H}]`$-value, but increasing the metallicity (to mimic non-LTE corrections) improves the fit. Figure 5b compares the halo sequence with Proffitt & VandenBerg’s (1991) isochrones that include He sedimentation. Microscopic diffusion makes the isochrones redder, modifies their shape, and predicts a lower turn-off luminosity: the best fit with the observed sequence is achieved for an age smaller by 0.5-1.5 Gyr than that obtained without diffusion. Models by Castellani et al (1997) show that, if sedimentation of metals is also taken into account, including its effects on the matter opacity, the isochrone shift is smaller than the shift obtained with He diffusion only.
Cayrel et al (1997b) and Pont et al (1997) estimated the local halo to be 12-16 Gyr old (from standard isochrones). To improve the precision more stars with accurate parallaxes are required. Subgiants are about 100 times rarer than subdwarfs, and we have only two subgiants with $`\sigma _\pi /\pi `$ $`<`$12.5% (and no subgiant with $`\sigma _\pi /\pi `$ $`<`$5%). After Hipparcos the position of the subgiant branch is still poorly determined, which limits the accuracy on the age determination of the halo stars.
THE ZAMS POSITIONS. The sample made of Hipparcos disk and halo stars spans the whole Galactic metallicity range. Figure 6 shows the non-evolved stars ($`\mathrm{M}_{\mathrm{bol}}`$$`>`$5.5) of Figure 1 and Figure 4 along with standard isochrones of various metallicities and solar-scaled helium ($`(\mathrm{\Delta }Y/\mathrm{\Delta }Z)_{}`$=2.2). It allows a discussion of the position of the zero age main sequence (ZAMS) as a function of metallicity and implications for the unknown helium abundances.
* MS width. Although stars generally do not lie where predicted, in particular at low metallicities, the observational and theoretical MS widths are in reasonable agreement for $`\mathrm{\Delta }Y/\mathrm{\Delta }Z`$ =2.2. This qualitative agreement is broadly consistent with $`\mathrm{\Delta }Y/\mathrm{\Delta }Z`$ ratio of $``$3$`\pm `$2 derived from similar measures of the lower MS width by Pagel & Portinari (1998) and the lower limit $`\mathrm{\Delta }Y/\mathrm{\Delta }Z`$ $``$2 obtained by Fernandes et al (1996) from pre-Hipparcos MS. It also agrees with extragalactic determinations (see Izotov et al 1997) or nucleosynthetic predictions.
* Helium abundance at solar metallicities. It can be noted from Figure 6 that there are 4 stars with Fe/H close to solar on the $`[\mathrm{Fe}/\mathrm{H}]`$=0.3 isochrone. Non-LTE $`[\mathrm{Fe}/\mathrm{H}]`$-corrections are negligible at solar metallicity. Microscopic diffusion may produce a shift in the H-R diagram: for a 0.8 $`\mathrm{M}_{}`$ star of solar Fe/H at 5 Gyr the shift is small and comparable to the observational error bars (but it increases with age). These disk stars are not expected to be very old and the shift could instead indicate that their He-content is lower than the solar-scaled value. Calibration of individual objects and groups with metallicities close to solar indicate an increase of helium with metallicity corresponding to $`\mathrm{\Delta }Y/\mathrm{\Delta }Z`$ $``$2.2 from the Sun (Lebreton et al 1999) and $`\mathrm{\Delta }Y/\mathrm{\Delta }Z`$ $``$2.3$`\pm `$1.5 from visual binaries (Fernandes et al 1998) but exceptions are found, such as in the (rather young) Hyades which, although metal-rich ($`[\mathrm{Fe}/\mathrm{H}]`$=0.14), appear to have a solar or even slightly sub-solar helium content with $`\mathrm{\Delta }Y/\mathrm{\Delta }Z`$ $``$1.4 (Perryman et al 1998). Going into finer resolution would clearly require more complete data including masses for enlarged samples of non-evolved stars.
* Position of metal-deficient stars. Very few metal-deficient stars have accurate positions in the non-evolved part of the H-R diagram: a gap appears for $`[\mathrm{Fe}/\mathrm{H}]`$$``$\[-1.4,-0.3\] and only 4 subdwarfs are found below $`[\mathrm{Fe}/\mathrm{H}]`$$``$-1.4. The empirical dependence of the ZAMS location with metallicity is impossible to establish for these stars, which are expected to have practically primordial helium contents. This adds to difficulties in estimating the distances of globular clusters (Eggen & Sandage 1962; Sandage 1970, 1983; Chaboyer et al 1998).
### 4.3 Stars in Open Clusters
Hipparcos observed stars in all open clusters closer than 300 pc and in the richest clusters located between 300 and 500 pc providing valuable material for distance scaling of the Universe and for studies of kinematical and chemical evolution of the Galaxy. The absolute cluster sequences in the H-R diagram may be constructed directly from Hipparcos distances independently of any chemical composition consideration. Each sequence covers a large range of stellar masses and contains stars which can reasonably be considered to be born at the same time with similar chemical composition. Several clusters provide tests of the internal structure models for a wide range of initial parameters, in particular for different metallicities.
THE HYADES. Obtaining high-quality astrometric data for the Hyades has been crucial, for it is the nearest rich star cluster, used to define absolute magnitude calibrations and the zero-point of the Galactic and extragalactic distance scales. Individual distances (mean accuracy of 5%) and proper motions were given by Hipparcos, providing a consistent picture of the Hyades distance, structure and dynamics (Perryman et al 1998). The recent determinations of the Hyades distance modulus $`(mM)`$ are all in very good agreement while the internal accuracy was largely improved with Hipparcos:
* ground-based: $`mM=3.32\pm 0.06`$ mag (104 stars, van Altena et al 1997b)
* HST : $`mM=3.42\pm 0.09`$ mag (7 stars, van Altena et al 1997a)
* Hipparcos: $`mM=3.33\pm 0.01`$ mag (134 stars within 10 pc of the cluster center, Perryman et al 1998)
* statistical parallaxes based on Hipparcos proper motions: $`mM=3.34\pm 0.02`$ mag (43 stars, Narayanan & Gould 1999a who also showed that the systematic error on the parallaxes toward the Hyades is lower than 0.47 mas).
Greatly improved precision is seen in the H-R diagrams built with Hipparcos data combined with the best ground-based observations (Perryman et al 1998):
* Figure 7 shows 40 stars with $`\mathrm{T}_{\mathrm{eff}}`$ and $`[\mathrm{Fe}/\mathrm{H}]`$=0.14$`\pm `$0.05 from detailed spectroscopic analysis delineating the lower part of the observational MS of the cluster (Cayrel de Strobel et al 1997a).
* Figure 8 is the whole H-R diagram in the ($`\mathrm{M}_\mathrm{V}`$, B-V) plane for 69 cluster members. Known or suspected binaries, variable stars, and rapid rotators have been excluded (Perryman et al 1998). Also, Dravins et al (1997) derived dynamical parallaxes for the Hyades members from the relation between the cluster space motion, the positions and the projected proper motions; these parallaxes are more precise (by a factor of about 2) than those directly measured by Hipparcos, yielding in turn a remarkably well-defined MS sequence in the H-R diagram, narrower than that given in Figure 8 (see Figure 2 in Dravins et al 1997).
* Figure 9 shows the empirical mass-luminosity (M-L) relation drawn from the (very accurate) masses of 5 binary systems (see caption). Nine of the stars are MS stars.
Comparisons with theoretical models yield some of the cluster characteristics (Lebreton et al 1997a, Perryman et al 1998, Lebreton 2000):
1. The comparison of the lowest part of the MS (Figure 7), representing the non-evolved stars, with theoretical ZAMS corresponding to the mean observed $`[\mathrm{Fe}/\mathrm{H}]`$ yields the initial cluster helium content $`Y_H`$=0.26$`\pm `$0.02 and metallicity $`Z_H`$=0.024$`\pm `$0.04. Metallicity is the dominant source of the uncertainty on $`Y`$.
2. The comparison of the whole observed sequence with model isochrones yields the cluster age. Figure 8 shows that the optimum fit is achieved with an isochrone of $`625\pm 50`$ Myr, $`Y_H`$=0.26, $`Z_H`$=0.024 and including overshooting. The turn-off region (which in the Hyades corresponds to the instability strip of $`\delta `$ Scuti stars) is rather well represented by the 625 Myr isochrone (see also Antonello & Pasinetti Fracassini 1998). The quoted uncertainty on age only includes the contribution from visual fitting of the isochrones. Additional errors on age result from unrecognized binaries, rotating stars, color calibrations and bolometric corrections, and from theoretical models in particular through the parameterization of overshooting (Lebreton et al 1995). It is therefore reasonable to give an overall age uncertainty of at least 100 Myr.
3. In Figure 9a the observed M-L relation is compared with the theoretical isochrone of 625 Myr, $`Y_H=0.26`$, $`Z_H=0.024`$, showing an excellent agreement. The lower part of the relation is defined by the very accurate masses of the two components of vB22. This system gives additional constraints on the $`Y_H`$-value derived from ZAMS calibration. Figure 9b illustrates how the positions of the two vB22 components may be used to constrain $`Y_H`$ in the whole metallicity range allowed by observations, $`[\mathrm{Fe}/\mathrm{H}]`$=0.14$`\pm `$0.05 (see also Lebreton 2000).
Furthermore, in the turn-off region of the Hyades, 5 $`\delta `$ Scuti stars are found that are quite rapid rotators ($`v_e\mathrm{sin}i`$ in the range 80-200 $`\mathrm{km}.\mathrm{s}^1`$, see Antonello & Pasinetti Fracassini 1998). From the measurement and analysis of their oscillation frequencies and the identification of the corresponding modes by means of models (of same age and chemical composition), we should be able to derive the inner rotation profile and learn about the size of convective cores and transport processes at work in the interiors (Goupil et al 1996, Michel et al 1999). For instance, the rotation profile is related to the redistribution of angular momentum by internal motions which could be generated by meridional circulation and shear turbulence in a rotating medium (see Zahn 1992). On the other hand, such motions might induce internal mixing, and as shown by Talon et al (1997), in the H-R diagram rotational effects “mimic” overshooting (for instance, in a star of 9 $`\mathrm{M}_{}`$, a rotational velocity of $`100\mathrm{km}.\mathrm{s}^1`$ is equivalent to an overshooting of $`\alpha _{\mathrm{ov}}`$ $``$0.2).
The study and intercomparison of accurate observations of the non-pulsating and pulsating stars located in the instability strip should clearly provide deeper insight into the internal structure and properties of stars of the Hyades cluster. However, such analysis has to integrate the various complications related to rotation, such as the displacements in any photometric H-R diagram by amounts depending on the equatorial velocity and inclination (Maeder & Peytremann 1972, Pérez-Hernández et al 1999) or the splitting of oscillation frequencies, which has to be considered in the mode identification.
THE PLEIADES AND OTHER OPEN CLUSTERS. The membership of stars in nine clusters closer than 300 pc was carefully assessed by van Leeuwen (1999a) and Robichon et al (1999a). Robichon et al also studied nine rich clusters within 500 pc with more than 8 members and 32 more distant clusters. For clusters closer than 500 pc, the accuracy on the mean parallax is in the range 0.2-0.5 mas and the accuracy on the mean proper motions is of the order of 0.1 to 0.5 mas per year. Results from the two groups are in very good agreement. Platais et al (1998) looked for (new) star clusters in Hipparcos data and found one new, a nearby cluster in Carina (d=132 pc) with 7 identified members.
In order to obtain an optimal mean parallax with correct error estimates, van Leeuwen (1999a) and Robichon et al (1999a) worked with the Hipparcos intermediate data corresponding to each cluster, parallax and proper motion of the cluster center, and position of each cluster member, instead of making a straight average of the parallaxes of the cluster members. Stars in open clusters are located within a few degrees on the sky and hence were often observed in the same field of view of the satellite. A combined solution can be obtained from intermediate data, which allows angular correlations to be taken into account and the resulting parallax errors to be minimized (van Leeuwen & Evans 1998; van Leeuwen 1999a, b; Robichon et al 1999a).
Mermilliod et al (1997), Robichon et al (1997), and Mermilliod (2000) compared the sequences of the various clusters in C-M diagrams derived from different photometric systems, and found puzzling results that are at odds with the common idea that differences in metallicity fully explain the relative positions of the non-evolved parts of the MS of different clusters:
* Some clusters have different metallicities but define the same main sequence in the ($`\mathrm{M}_\mathrm{V}`$, B-V) plane (Praesepe, Coma Ber, $`\alpha `$Per, Blanco 1). For instance, Coma Ber has a quasi-solar metallicity while its sequence is similar to that of the Hyades, or of the metal-rich Praesepe.
* Some clusters sequences (Pleiades, IC 2391 and 2602) are abnormally faint with respect to others, for instance Coma Ber. The metallicity of the Pleiades as determined from spectroscopy is almost solar, and similar to that of Coma Ber, but the Pleiades sequence lies (unexpectedly) $``$ 0.3-0.4 mag below the Praesepe, Coma Ber, or Hyades sequence.
* Van Leeuwen (1999a, b) even suggested a possible (although unexpected) correlation between the age of a cluster and its position in the H-R diagram.
Prior to Hipparcos, precise trigonometric parallaxes had not been obtained for clusters except the Hyades. Distances to open clusters were evaluated through the main sequence fitting technique: the non-evolved part of the (observed) cluster sequence was compared to the non-evolved part of the (absolute) lower MS (ZAMS) of either (1) theoretical isochrones, (2) field stars or (3) Hyades after a possible correction of chemical composition differences. The magnitude differences between absolute and apparent ZAMS directly yielded the distance modulus of the cluster.
The Hipparcos distances to the 5 closest open clusters (Hyades, Pleiades, $`\alpha `$Per, Praesepe, and Coma Ber) can be compared to those recently derived from MS fitting by Pinsonneault et al (1998); they compared theoretical isochrones, translated into the C-M plane by means of Yale color calibrations, to observational data both in the ($`\mathrm{M}_\mathrm{V}`$, B-V) and ($`\mathrm{M}_\mathrm{V}`$, V-I) planes. The B-V color indice is more sensitive to metallicity than V-I (Alonso et al 1996b), so Pinsonneault et al derived as a by-product the value of the metallicity that gives the same distance modulus in the two planes and compared it to spectroscopic determinations. They judged their distance modulii to be in good agreement with Hipparcos results except for the Pleiades and Coma Ber. For Coma Ber, the problem could result from the old VRI colors used. For the Pleiades the discrepancy with Hipparcos amounts to 0.24 mag, and the $`[\mathrm{Fe}/\mathrm{H}]`$-value derived from MS fitting in the two color planes agrees with the spectroscopic determination of Boesgaard & Friel (1990), $`[\mathrm{Fe}/\mathrm{H}]`$= -0.034$`\pm `$0.024, although values in the range -0.03 to +0.13 can be found in the literature. In fact, with that metallicity the Hipparcos sequence of the Pleiades could be reproduced by classical theoretical models, provided they have a high helium content. The exact value depends on the model set and its input physics: Pinsonneault et al found $`Y`$=0.37, Belikov et al (1998) found $`Y`$=0.34 but for $`[\mathrm{Fe}/\mathrm{H}]`$=0.10 and I find Y $``$ 0.31. In any case, high helium content is only marginally supported by observations (Nissen 1976). Pinsonneault et al examined other possible origins of the discrepancy (erroneous metallicity, age-related effects, reddening) and concluded that none of them is likely to be responsible for the Pleiades discrepancy.
In parallel, Soderblom et al (1998) looked for young solar-type stars appearing as (anomalously) faint as the Pleiades. They found 50 field stars expected to be young (i.e. showing activity from Ca II H and K lines), but none of them lies significantly below the ZAMS. They also examined the subluminous stars observed by Hipparcos: they chose six stars among those lying well below the ZAMS, measured their spectroscopic metallicities, and found them to be metal-deficient with respect to the Sun with, in addition, kinematics typical of stars of a thick disk or halo population.
Soderblom et al and Pinsonneault et al concluded that, taking the Hipparcos results for the Pleiades at face value, it would be abnormal not to find stars similar to the Pleiades in the field. They inferred that the distance obtained from multi-color MS fitting is correct and accurate to about 0.05 mag, and concluded that the distance to Pleiades obtained from the analysis of Hipparcos data is possibly wrong at the 1 mas level, which is greater than the mean random error. They invoked statistical correlations between right ascension and parallax ($`\rho _{\alpha \mathrm{cos}\delta }^\pi `$) arising from the non-uniform distribution of Hipparcos observations over time (and in turn along the parallactic ellipse) which affects all stars, including clusters. Pinsonneault et al noted that in the Pleiades the brightest stars (1) are highly concentrated near the cluster center and are therefore subject to spatial correlations which gives them nearly the same parallax, (2) have smaller $`\sigma _\pi `$ than fainter stars which gives them more weight in the mean parallax, and (3) are those which have the highest values of $`\rho _{\alpha \mathrm{cos}\delta }^\pi `$ and also the highest parallaxes in the Hipparcos Catalogue. They suggest that the “true” parallax (close to that obtained through MS fitting) is obtained if the brightest stars with high $`\rho _{\alpha \mathrm{cos}\delta }^\pi `$ are excluded from the calculation.
Narayanan & Gould (1999b) determined the parallaxes of the Pleiades stars by means of Hipparcos proper motions. The resulting distance modulus has a rather large error bar ($`mM=5.58\pm 0.18`$ mag), but it is in disagreement with that derived directly from Hipparcos parallaxes ($`mM=5.36\pm 0.07`$ mag), and in agreement with that obtained through MS fitting ($`mM=5.60\pm 0.05`$ mag). Narayanan & Gould also argue that the differences between the Hipparcos trigonometric parallaxes and the parallaxes derived from Hipparcos proper motions reflect spatial correlations over small angular scales with an amplitude of up to 2 mas.
Robichon et al (1999a, b) and van Leeuwen (1999b) have subsequently derived more reliable distance estimates to these clusters and performed tests that do not support Pinsonneault et al’s conclusion. The difference between Hipparcos and MS fitting distance moduli is small for the Hyades (0.01 mag), whereas for other clusters it ranges from -0.17 mag ($`\alpha `$Per) to +0.24 mag (Pleiades). In fact, except for the Hyades, the difference is always larger than the error on MS fitting distance modulus (0.05 mag). Robichon et al showed that while the solution proposed by Pinsonneault et al improves the situation for the Pleiades, it would introduce new difficulties for Praesepe. By means of Monte-Carlo simulations of the Pleiades stars, they showed that the mean value of the Pleiades parallax does not depend on the correlations $`\rho _{\alpha \mathrm{cos}\delta }^\pi `$. They also carefully examined distant stars and clusters with high $`\rho _{\alpha \mathrm{cos}\delta }^\pi `$. Through these tests, Robichon et al made the Hipparcos distance to Coma Ber or Pleiades more secure, and did not find any obvious bias on the parallax resulting from a correlation between right ascension and parallax, either for stars within a small angular region or for the whole sky.
On the other hand, distances from MS fitting could be subject to higher error bars than quoted by Pinsonneault et al. They depend on reddening and on transformations from the ($`\mathrm{M}_{\mathrm{bol}}`$, $`\mathrm{T}_{\mathrm{eff}}`$) to the C-M plane if theoretical ZAMS are used as reference (or on metallicity corrections if empirical ZAMS are compared). Robichon et al (1999b) compared solar ZAMS from Pinsonneault et al to those I calculated both in the theoretical and in the ($`\mathrm{M}_\mathrm{V}`$, B-V) planes. They showed that while the two ZAMS are within 0.05 mag in the theoretical H-R diagram, they differ by 0.15-0.20 mag in the range B-V=0.7-0.8 in the ($`\mathrm{M}_\mathrm{V}`$, B-V) plane, simply because different C-M transformations have been applied. Also, MS fitting often relies on rather old and inhomogeneous color sources (in the separate Johnson and Kron-Eggen RI systems) requiring transformations to put all data on the same (Cape-Cousins system) scale. It would therefore be worthwhile to verify the quality and precision of these data by making new photometric measurements of cluster stars.
Let us come back to the difficult question of metallicity. As pointed out by Mermilliod et al (1997), photometric and spectroscopic approaches may produce quite different results. Metallicities have been derived recently by M Grenon (1998, private communication) from large sets of homogeneous observations in the Geneva photometric system. He obtained $`[\mathrm{Fe}/\mathrm{H}]`$=-0.112$`\pm `$0.025 for the Pleiades (quite different from published spectroscopic values) and $`[\mathrm{Fe}/\mathrm{H}]`$=0.170$`\pm `$0.010 and 0.143$`\pm `$0.008 for Praesepe and the Hyades respectively (both in agreement with spectroscopy). The observed cluster sequences obtained with Hipparcos distances for the three clusters can be roughly reproduced by theoretical models computed with the photometric metallicities (and allowing for small variations of the helium content around the solar-scaled value) and transformed to the C-M plane according to the Alonso et al (1996b) and Bessell et al (1998) calibrations (Robichon et al 1999b).
In conclusion, we point out that a detailed study of the fine structure of the H-R diagram of the Pleiades (and other clusters) requires supplementary observations (colors and abundances) and further progress in model atmospheres. Today there is no obvious solid argument against the published Hipparcos distances. In order to identify and understand the remaining discrepancies with stellar models, the entire set of observed clusters has to be considered (van Leeuwen 1999a). Furthermore, not only the positions of the sequences in the H-R diagram but also the density of stars along them have to be intercompared. For instance, the luminosity function of young clusters exhibits a particular feature (local peak followed by dip) that is interpreted as a signature of pre-MS stars and might provide information on the initial mass function and stellar formation history (see the study of the Pleiades by Belikov et al 1998). On the other hand, since the error bars on luminosity are now small with respect to errors on color indices, stronger constraints are expected from the mass-luminosity relation, as in the Hyades. Observations of binaries in clusters are urgently needed and there is hope to detect them in the future, for, as pointed out by Soderblom et al, the difficult detection and measurement of visual binary orbits in the Pleiades is within the capabilities of experiments on board HST .
## 5 RARE, FAINT, SPECIAL, OR INACCESSIBLE OBJECTS
### 5.1 Globular Clusters Through Halo Stars
Globular clusters were beyond the possibilities of Hipparcos, but the knowledge of distances to nearby subdwarfs gave distance estimates to a few of them through the MS fitting technique (Sandage 1970; Reid 1997, 1998; Gratton et al 1997; Pont et al 1997, 1998; Chaboyer et al 1998), comparing the non-evolved part of the (absolute) subdwarf main-sequence to the non-evolved part of the (observed) globular cluster sequence. Although simple, the technique has to be applied with caution:
* Only halo subdwarfs and globular clusters with the most precise data should be retained. Abundances should be accurate and on a consistent scale. Globular cluster abundances are usually determined only for giants, while recent preliminary values have been obtained for subgiants in M92 (King et al 1998). Abundance comparisons between (1) field and cluster stars and (2) dwarfs and giants have shown sometimes puzzling differences (King et al 1998, Reid 1999). Questions have been raised as to whether they are primordial or appear during evolution, but definite answers clearly require better spectra for all types of stars as well as spherical model atmospheres with better treatment of convection. The globular cluster sequence should be determined from good photometry well below the MS turn-off, and the correction for interstellar reddening has to be well estimated. Very few halo stars have parallaxes accurate enough to fix precisely the position of the ZAMS (Section 4.2).
* Biases (see e.g. Lutz & Kelker 1973; Hanson 1979; Smith 1987) resulting from the selection of the sample in apparent magnitude, parallax, and metallicity, have to be corrected for (Pont et al 1998, Gratton et al 1997); alternatively, samples free of biases must be selected, which implies retaining the very nearby stars with highly accurate parallaxes (Chaboyer et al 1998, Brown et al 1997).
* Globular cluster sequence and halo sequence should (ideally) have similar initial chemical compositions. Because of the small number of subdwarfs in each interval of metallicity, it is not possible to properly establish the variation of the observed ZAMS position with metallicity, and to correct for chemical composition differences between globular clusters and subdwarfs empirically. Chaboyer et al (1998) found it safer to limit the method to globular clusters that have their equivalent in the field with the same $`[\mathrm{Fe}/\mathrm{H}]`$ and $`[\alpha /\mathrm{Fe}]`$ content. Gratton et al (1997) and Pont et al (1998) applied theoretical color corrections to the subdwarf data to account for metallicity differences with globulars. In addition, element sedimentation might introduce further difficulties, as already mentioned in Section 4.2. As pointed out by Salaris & Weiss (2000), the present surface chemical composition of field subdwarfs no more reflects the initial one if microscopic diffusion has been efficient during evolution, while, in globular cluster giants which have undergone the first dredge-up, the chemical abundances have been almost restored to the initial ones.
* Unresolved known or suspected binaries can introduce errors in the definition of the ZAMS position. Chaboyer et al and Gratton et al excluded them whereas Pont et al applied an average correction of 0.375 mag on their position, a procedure that has been criticized by Chaboyer et al and Reid (1998).
* Evolved stars have to be excluded (there is no certainty that globular clusters and halo dwarfs have exactly the same age). From theoretical models it is estimated that stars fainter than $`\mathrm{M}_\mathrm{V}`$$``$ 5.5 are essentially unevolved.
The number of globular clusters studied by the different authors varies because of the different criteria and techniques chosen to select the subdwarfs samples. Nevertheless they all agree on the general conclusion that globular cluster distances derived from MS fitting are larger by $``$5-7% than was previously found. Chaboyer (1998) calculated an average of distances to globular clusters obtained with different methods (MS fitting, astrometry, white dwarf sequence fitting, calibration of the mean magnitudes of RR Lyrae stars in the Large Magellanic Cloud, comparison with theoretical models of horizontal branch stars, statistical parallax absolute magnitude determinations of field RR Lyrae from the Hipparcos proper motions, etc) and noted that the distance scale is larger (by 0.1 mag) than his pre-Hipparcos reference.
It is worth pointing out that the statistical parallax method alone favors a shorter distance scale (by $``$0.3 mag with respect to MS fitting result). As reviewed by Layden (1998) and by Reid (1999), the statistical parallax method was applied by independent groups who found concordant results. However, the absolute magnitudes $`\mathrm{M}_\mathrm{V}(\mathrm{RR})`$ of the halo RR Lyrae derived from the statistical parallax method (on the basis of Hipparcos proper motions and of radial velocities) are also $``$ 0.3 mag fainter than the magnitudes obtained through a method directly based on Hipparcos parallaxes (Groenewegen & Salaris 1999); these latter being in turn in good agreement with the MS fitting result. As discussed thoroughly by Reid (1999), there are several difficulties related to the $`\mathrm{M}_\mathrm{V}(\mathrm{RR})`$-calibration and to its comparison with other distance calibrations that still hinder the coherent and homogeneous understanding of the local distance scale.
In Caputo’s (1998) and Reid’s (1999) reviews the ages of globulars are discussed. After Hipparcos, ages of globular clusters are reduced by typically 2-3 Gyr, because of both larger distances and improvements in the physics of the models, mainly in the equation of state and in the consideration of microscopic diffusion. Present ages are now in the range 10-13 Gyr, which can be compared to the previous interval of 13-18 Gyr (see VandenBerg et al’s review, 1996).
Chaboyer et al (1998) claimed that the (theoretical) absolute magnitude $`\mathrm{M}_\mathrm{V}`$ and lifetimes of stars at the MS turn-off (T-O) in globular clusters are now well understood, since the physics involved is very similar to that of the Sun, which is in turn well constrained by seismology. In particular, $`\mathrm{M}_\mathrm{V}`$(T-O) is quite insensitive to uncertainties related to model atmospheres or convection modeling (see also Freytag & Salaris 1999). Chaboyer et al suggested that no significant changes (more than $``$5%) in the derived ages of globular clusters are expected from future improvements in stellar models. Conversely, distance and abundance determinations are far from definite, and the quasi-verticalness of isochrones in the T-O region makes the determination of $`\mathrm{M}_\mathrm{V}`$(T-O) difficult (see Vandenberg et al 1996). Further revision of the ages is therefore not excluded. Also a global agreement between an entire globular cluster sequence and the corresponding model isochrone is far from being reached.
Age of the Universe. The ages of the oldest objects in the Galaxy, the most metal-poor halo or globular cluster stars, provide a minimum value for the age of the Universe $`\mathrm{T}_\mathrm{U}`$. Globular cluster ages (10-13 Gyr) from comparisons of isochrones with observed $`\mathrm{M}_\mathrm{V}`$(T-O) are presently the most reliable, but two independent methods look promising:
* Thorium (half-life 14.05 Gyr) has been detected and measured by Sneden et al (1996) in an ultra-metal-poor giant, too faint for observation by Hipparcos, and the star radioactive decay age is estimated to be 15$`\pm `$4 Gyr. In the future, such observations of more stars and the possible detection of Rhenium and Uranium could provide strong constraints for $`\mathrm{T}_\mathrm{U}`$.
* Observations of (faint) white dwarfs (WD) in globular clusters are now within reach of experiments on board HST , and a lower limit to the age of WD in M4 of $``$9 Gyr has been derived from a comparison with theoretical WD cooling curves (Richer et al 1997). Future access to cooler and fainter objects will better constrain $`\mathrm{T}_\mathrm{U}`$.
According to Sandage & Tammann (1997) and Saha et al (1999) the Hubble constant $`H_0`$ should be in the range $`55\pm 5\mathrm{k}\mathrm{m}.\mathrm{s}^1.\mathrm{Mpc}^1`$, which implies $`T_U=\frac{2}{3}H_0^11113.5`$ Gyr, indicating that no strong discrepancy with the age of the oldest known stars remains.
### 5.2 Variable Stars
I shall not discuss the revisions of the distance scale based on pulsating stars (RR Lyrae, Cepheids, Miras, high-amplitude $`\delta `$ Scuti stars) because this topic has been extensively reviewed by Caputo (1998) and Reid (1999).
Both new insight as well as new questions about the physics governing pulsating stars have been generated from the combination of Hipparcos distances with asteroseismic data. When the magnitude of a star is modified and its error box reduced, the mass and evolutionary stage attributed to the star may be modified. For variable stars, a different evolutionary stage may give a drastically different eigenmode spectrum, and in turn may change the mode identification and asteroseismic analysis (see Liu et al 1997). Høg & Petersen (1997) showed that for the two double-mode, high-amplitude $`\delta `$ Scuti variables SX Phe and AI Vel, the masses derived on one hand from stellar envelope models and pulsation theory, and on the other hand from the position in the H-R diagram through stellar evolution models are in nice agreement if the Hipparcos parallaxes (accurate to 5-6%) are used. Further implications of Hipparcos distances on the understanding of $`\delta `$ Scuti stars, $`\lambda `$ Bootis, and rapidly oscillating Ap stars have been discussed in several papers (see for instance Audard et al 1998, Viskum et al 1998, Paunzen et al 1998, Matthews et al 1999), whereas the physical processes relevant to the Asymptotic Giant Branch and pulsation modeling of Miras and Long Period Variables were examined by Barthès (1998, see also references therein).
### 5.3 White Dwarfs
The white dwarf (WD) mass-radius (M-R) relation was first derived by Chandrasekhar (1931) from the theory of stars supported by the fully degenerate electron gas pressure. It has been refined by Hamada & Salpeter (1961), who calculated zero-temperature (fully degenerate) WD models of different chemical composition (He, C, Mg, Si, S, and Fe) and by Wood (1995), who calculated WD models with carbon cores and different configurations of hydrogen and/or helium layers and followed the thermal evolution of WD as they cool. Although theoretical support is strong, it has long been difficult to confirm the relation empirically because of (1) the very few available WD with measures of masses and radii, (2) the size of the error bars and (3) the intrinsic mass distribution of the WD, which concentrates them in only a small interval around 0.6 $`\mathrm{M}_{}`$ (Schmidt 1996).
The M-R relation, assuming that WD have a carbon core, is a basic underlying assumption in most studies of WD properties. It serves to determine the mass of WD, and in turn their mass distribution and luminosity function. It is important because WD feature in many astrophysical applications such as the calibration of distances to globular clusters (Renzini et al 1996) or the estimate of the age of Galactic disk and halo by means of WD cooling sequences (see Winget et al 1987, d’Antona & Mazzitelli 1990). The more precisely the M-R relation is defined by observations, the better the tests of theoretical models of WD interiors that can be undertaken. These include tests of the inner chemical composition of WD, thickness of the hydrogen envelope of DA WD, or the characterization of their strong inner magnetic fields.
Depending on the white dwarf considered, the empirical M-R relation can be obtained by different means:
1. Surface brightness method. If $`\mathrm{T}_{\mathrm{eff}}`$ and $`\mathrm{log}g`$ are determined (generally from spectroscopy), then model atmospheres allow calculation of the energy flux at the surface of the star, which when compared to the flux on Earth yields the angular diameter $`\varphi `$ (see Schmidt 1996). The radius R is obtained from the parallax and $`\varphi `$, M is deduced from R and $`\mathrm{log}g`$. This method requires high-resolution spectra and largely depends on model atmospheres.
2. Gravitational redshift. The strong gravitational field at the surface of a WD causes a redshift of the spectral lines, the size of which depends on the gravitational velocity $`v_{\mathrm{grs}}=\frac{GM}{Rc}`$ ($`c`$ is the speed of light). If $`v_{\mathrm{grs}}`$ can be measured and the gravity is known, then M and R can easily be obtained independently of the parallax. $`v_{\mathrm{grs}}`$ can only be measured in WD members of binary systems, common proper motion pairs (CPM), or clusters because the radial velocity is required to distinguish the gravitational redshift from the line shift due to Doppler effect. Also, very high-resolution spectra are needed.
3. WD in visual binary systems. Masses may be derived directly from the orbital parameters through the Kepler’s third law, provided the parallax is known. Radii are derived from the knowledge of $`\mathrm{T}_{\mathrm{eff}}`$ and distances.
More than 15 years ago, when the Hipparcos project began, uncertainties on WD ground-based parallaxes were at least 10 mas. During the last 10 years, due to great instrumental progress, parallax determinations were improved by a factor of about 2, and more accurate atmospheric parameters $`\mathrm{T}_{\mathrm{eff}}`$, $`g`$ and $`v_{\mathrm{grs}}`$ were obtained. In the meantime, Hipparcos observed 22 white dwarfs (11 field WD, 4 WD in visual binaries, and 7 in CPM systems) among which 17 are of spectral type DA. Although they are close to the faint magnitude limit of Hipparcos, the mean accuracy on their parallaxes is $`\sigma _\pi `$3.6 mas (Vauclair et al 1997).
Vauclair et al (1997) and Provencal et al (1998) studied the whole sample of WD observed by Hipparcos. The M-R relation is narrower and most points are within $`1\sigma `$ of Wood’s (1995) evolutionary models of WD with carbon cores and hydrogen surface layers. The theoretical shape is still difficult to confirm because of the lack of objects in the regions of either high or low mass. Furthermore, the error bars are still too large to distinguish fine features of the theoretical models, such as between evolutionary and zero-temperature sequences or thick and thin hydrogen envelopes (Vauclair et al 1997), except for some particular stars (Shipman et al 1997, Provencal et al 1998). Other effects such as alterations due to strong internal magnetic fields are not yet testable (Suh & Mathews 2000).
WD in binary systems. Prior to Hipparcos, Sirius B was the only star roughly located on the expected theoretical M-R relations; the others (Procyon B, 40 Eridani B and Stein 2051) were at least $`1.5\sigma `$ below the theoretical position (see Figure 1 in Provencal et al 1998). After Hipparcos, as shown by Provencal et al, the error on the radius is dominated by errors on flux and $`\mathrm{T}_{\mathrm{eff}}`$. On the other hand, the parallax error still dominates the error on mass, except for Procyon where the error on the component separation plays a major role.
* Sirius B is more precisely located on a Wood’s (1995) M-R relation for a DA white dwarf of the observed $`\mathrm{T}_{\mathrm{eff}}`$ with a thick H layer and carbon core (Holberg et al 1998). Also compatible with Wood’s thick H layer models is the position of V471 Tau, a member of an eclipsing binary system for which the Hipparcos parallax supports the view that it is a member of the Hyades (Werner & Rauch 1997, Barstow et al 1997). The mass of 40 Eri B increased by 14%. The star is now back on Hamada & Salpeter’s (1961) M-R relation for carbon cores, making it compatible with single star evolution (Figure 1 by Shipman et al 1997) and it does not appear to have a thick H layer. Sirius B (the most massive known WD) and 40 Eri B (one of the less massive) nicely anchor the high-mass and low-mass limits of the M-R relation.
* The case of Procyon B remains puzzling. The position is not compatible with models with carbon cores, and would be better accounted for with iron or iron-rich core models (Provencal et al 1997). Provencal et al (1998) also examined seven white dwarf members of CPM pairs (more distant and fainter than WD in visual binaries) with Hipparcos distances and gravitational redshift measurements. They showed that two of them also lie on theoretical M-R relations corresponding to iron cores. This is not predicted by current stellar evolution theory, and further work is required to clarify this problem.
In conclusion, better distances from Hipparcos and high-resolution spectroscopy has allowed better assessment of the theoretical M-R relation for white dwarfs, and has shown evidence for difficulties for a few objects that do not appear to have carbon cores. Future progress will come from further parallax improvements and from better $`\mathrm{T}_{\mathrm{eff}}`$, $`v_{\mathrm{grs}}`$, magnitudes, and orbital parameters for visual binaries. A better understanding of the atmospheres is required: convection plays a role in the cooler WD, additional pressure effects due to undetectable helium affect the gravity determination, and incorrect H layer thickness estimates change the mass attributed to WD. Further information coming from asteroseismology or spectroscopy would help.
## 6 FUTURE INVESTIGATIONS
Hipparcos has greatly enlarged the available stellar samples with accurate and homogeneous astrometric and photometric data. To fully exploit this new information, many studies have been undertaken (several hundred papers devoted to stellar studies based on or mentioning Hipparcos data can be found in the literature), therefore the present review could not be fully exhaustive.
The Hipparcos mission succeeded in clarifying our knowledge of nearby objects, and allowed first promising studies of rarer or farther objects. After Hipparcos, the theory of stellar structure and evolution is further anchored, and some of its physical aspects have been better characterized. For instance, new indications that the evolution of low-mass stars is significantly modified by microscopic diffusion have been provided by fine studies of the H-R diagram, and consequences for age estimates or surface abundance alterations have been further investigated. On the other hand, Hipparcos left us with intriguing results that raised new questions. For example, the unexpected position of the white dwarf Procyon B on a theoretical mass-radius relation corresponding to iron cores is still not understood.
Today, uncertainties on distances of nearby stars have been reduced significantly such that other error sources emerge to dominate, hindering further progress in the fine characterization of stellar structure. Progress on atmosphere modeling is worth being pursued, for it has implications for observational parameters (effective temperatures, gravities, abundances, bolometric corrections), theoretical models (outer boundary conditions), and color calibrations. A thorough theoretical description of transport processes (convection, diffusion) and related effects (rotation, magnetic fields) is needed to improve stellar models, as well as further improvements or refinements in microscopic physics (low-temperature opacities, nuclear reaction rates in advanced evolutionary stages).
What is now needed from the observational side is (1) enlarged samples of rare objects (distant objects, faint objects or objects undergoing rapid evolutionary phases), (2) an increased number of more “common” objects with extremely accurate data (including masses), and (3) a census over all stellar populations.
These goals should be (at least partially) achieved by future astrometric missions. The NASA Space Interferometry Mission (SIM), scheduled for launch in 2006, will have the capability to measure parallaxes to 4 $`\mu `$arcsecond and proper motions to 1-2 $`\mu `$arcsecond per year down to the $`20^{\mathrm{th}}`$ magnitude which represents a gain of three orders of magnitude with respect to Hipparcos (Peterson & Shao 1997). The ESA candidate mission GAIA is dedicated to the observation of about one billion objects down to V$``$20 mag (and typical $`\sigma _\pi `$10 $`\mu `$arcsecond at V=15 mag). GAIA will also provide multi-color, multi-epoch photometry for each object, and will give access to stars of various distant regions of the Galaxy (halo, bulge, thin and thick disk, spiral arms). It is aimed to be launched in 2009 (Perryman et al 1997b).
Asteroseismology has already proved to be a unique tool to probe stellar interiors. Space experiments are under study. The first step will be the COROT mission (aimed to be launched in 2003), designed to detect and characterize oscillation modes in a few hundred stars, including solar-type stars and $`\delta `$ Scuti stars (see Baglin 1998).
## 7 Acknowledgments
I enjoyed very much working with Ana Gómez, Roger and Giusa Cayrel, João Fernandes, Michael Perryman, Noël Robichon, Annie Baglin, Jean-Claude Mermilliod, Marie-Noël Perrin, Frédéric Arenou, Catherine Turon, and Corinne Charbonnel on the various subjects presented here. I particularly thank Roger Cayrel, Annie Baglin, Michael Perryman, and Catherine Turon for a careful reading of the original manuscript. I am also very grateful to Allan Sandage for his remarks and suggestions which helped improving the paper. I am especially grateful to Frédéric Thévenin, Hans-Günter Ludwig, Misha Haywood, Don VandenBerg, Marie-Jo Goupil, Pierre Morel, Jordi Carme and Franck Thibault for fruitful discussions and advice. The University of Rennes 1 is acknowledged for working facilities. This review has made use of NASA’s Astrophysics Data System Abstract Service mirrored at CDS (Centre des Données Stellaires, Strasbourg, France).
|
warning/0003/hep-th0003135.html
|
ar5iv
|
text
|
# Refining the anomaly consistency condition
## I Introduction
Upon quantization, symmetries of the classical action can be affected by anomalies. These anomalies have been shown to satisfy important consistency conditions following from the algebra of the symmetries . The consistency conditions have been elegantly reformulated as a cohomological problem, first in the case of the gauge anomaly in Yang-Mills theories and then also in the case of global symmetries with a closed algebra .
In the case of gauge theories, the formalism has been generalized some time ago to cover generic gauge theories, and the expression of the anomaly consistency in that context has been discussed . Only recently, the corresponding generalization in the case of global symmetries with a generic algebra has been achieved .
In the same way, one can study the constraints on the anomalies to an invariant renormalization of symmetric integrated and non integrated operators . In the cohomological reformulation, symmetric integrated or non integrated operators correspond to BRST cohomological classes in ghost number $`0`$ in the space of local functionals, respectively in the space of local functions. The non trivial anomalies that can appear can be shown to belong to the corresponding BRST cohomological classes in ghost number $`1`$.
The computation in the space of local functionals of the cohomology of the standard BRST differential $`s`$ with antifields can be reduced to the computation of a relative cohomological group in the space of local $`n`$-forms by introducing the space-time exterior derivative $`d`$. It is then related to the cohomology of $`s`$ in the space of form valued local functions through descent equations . The same holds for the BRST differential $`\overline{s}`$ associated to the extended antifield formalism.
More generally, we will call the relative cohomological groups $`H^{g,p}(s|d)`$ and $`H^{g,p}(\overline{s}|d)`$ in ghost number $`g`$ and form degree $`p`$ local BRST cohomological groups. A more detailed analysis of the descent equations shows that these groups are characterized by two integers, the length $`d`$ of their descents and the length $`l`$ of their lifts.
The purpose of this article is twofold. First, it is shown that in the extended antifield formalism, a quantum BRST differential can be constructed even for anomalous theories, because the quantum version in terms of the effective action of the extended master equation can be written as a functional differential equation. Then it is shown that both in the extended antifield formalism and in the standard antifield formalism, the anomaly appearing in the renormalization of a local BRST cohomological class with a descent of length $`d`$ and a lift of length $`l`$ are characterized by a descent which is shorter or equal to $`d`$ and a lift which is longer or equal to $`l`$.
As an application, it is shown that the anomalous master equation for Yang-Mills theories discussed in can be viewed as a particular case of the anomalous master equation of the extended antifield formalism. Then, we discuss how the Adler-Bardeen theorem for the non abelian gauge anomaly can be understood as a direct consequence of the fact that the length of the descent of the gauge anomaly is $`4`$, while the length of the descent of all the other cohomological classes coupled to the action is $`0`$. These considerations are purely cohomological, so that they do not depend on the way the gauge is fixed or on power counting restrictions. Furthermore, this aproach to the Adler-Bardeen theorem does not require the use of the Callan-Symanzik equation or assumptions on the beta functions of the theory.
The article is organized as follows. In the next section, we review the relevant features of the extended antifield formalism, with special care devoted to the derivation of the quantum BRST differential in the anomalous case. The characterization of local BRST cohomological groups according to the lengths of their descents and their lifts is explained in section III. The fourth section contains the main result on the lengths of the descents and the lifts of the anomalies. An alternative derivation of the main results clarifying the underlying mechanism is presented in section V, where differentials controling the one loop anomalies arising in the renormalization of BRST cohomological groups are introduced. Anomalous Yang-Mills theory and the Adler-Bardeen theorem are discussed in section VI.
## II The extended antifield formalism
### A Classical theory
The Batalin-Vilkovisky formalism allows to formulate the BRST differential controling the gauge symmetries under renormalization for generic gauge theories. The formalism can be extended so as to include (non linear) global symmetries with a generic algebra . This is achieved by coupling the BRST cohomological classes in negative ghost numbers with constant ghosts. There is a further extension to include the BRST cohomological classes in all the ghost numbers , which allows to take into account in a systematic way all higher order cohomological constraints due to the antibracket maps .
The completely extended formalism is obtained by first computing a basis for the local BRST cohomological classes $`H^{g,n}(s|d)`$<sup>*</sup><sup>*</sup>*We consider the local BRST cohomology in the space of Lorentz-invariant polynomials or formal power seris in the $`dx^\mu `$, the coupling constants, the fields, antifields and their derivatives. associated to the standard differential $`s`$ and coupling those classes that are not already contained in the solution of the standard master equation with the help of new independent coupling constants. This action can then be extended by terms of higher orders in the new couplings in such a way that, if we denote by $`\xi ^A`$ all the couplings corresponding to the independent BRST cohomological classes, the resulting action $`S`$ satisfies the extended master equation
$`{\displaystyle \frac{1}{2}}(S,S)+\mathrm{\Delta }_cS=0.`$ (1)
The BRST differential associated to the solution of the extended master equation is
$`\overline{s}=(S,)+\mathrm{\Delta }_c^L,`$ (2)
where $`\mathrm{\Delta }_c=\frac{^R}{\xi ^A}f^A`$, while $`\mathrm{\Delta }_c^L=()^Af^A\frac{^L}{\xi ^A}`$, with $`f^A`$ depending at least quadratically on the couplings $`\xi `$ alone. Both derivations satisfy $`(\mathrm{\Delta }_c)^2=0=(\mathrm{\Delta }_c^L)^2`$. Since there is no dependence on the fields and the antifields, $`\mathrm{\Delta }_c^L(A,B)=(\mathrm{\Delta }_c^LA,B)+()^{A+1}(A,\mathrm{\Delta }_c^LB)`$, with the appropriate version holding for the right differential $`\mathrm{\Delta }_c`$.
Besides the information on the invariance of the original action under the non trivial gauge symmetries and their commutator algebra contained in the standard solution of the master equation, the extended master equation contains the information on the antibracket algebra of all local BRST cohomological classes associated to the standard BRST differential $`s`$. These classes contain the generators of all the generalized non trivial symmetries Higher order symmetries have to be taken into account in the generic case . of the theory in negative ghost number. This is the reason why the extended master equation encodes in particular the invariance of the original action under all the non trivial global symmetries as well as their commutator algebra. In ghost number zero, the local BRST cohomological classes $`H^{g,n}(s|d)`$ contain all the non trivial generalized observables of the theoryThey correspond to local functionals built out of the original fields and their derivatives that are gauge invariant when the gauge covariant equations of motion hold, and that are linearily independent even when the gauge covariant equations of motion hold (see e.g. ).. In ghost number $`1`$, these classes contain all the anomaly candidates that could affect the standard Batalin-Vilkovisky master equation and in ghost number strictly greater than $`1`$, they contain the anomaly candidates in ghost number $`g+1`$ that could occur because the anomaly candidates in ghost number $`g`$ have been coupled to the action.
The cohomology of BRST differential $`\overline{s}`$ of the extended formalism in the space $`F`$ of $`\xi `$ dependent local functionals in the fields, the antifields and their derivatives, is isomorphic to the cohomology of
$`s_{\mathrm{\Delta }_c}=[\mathrm{\Delta }_c,]`$ (3)
in the space of graded right derivations $`\lambda =\frac{^R}{\xi ^A}\lambda ^A`$, with $`\lambda ^A`$ a function of $`\xi `$ alone, $`[,]`$ being the graded commutator for graded right derivations. If $`\mu `$ is a $`s_{\mathrm{\Delta }_c}`$ cocycle, the corresponding $`\overline{s}`$ cocycle is given by $`\mu S=\frac{^RS}{\xi ^B}\mu ^B`$. In particular, the general solution to the system of equations
$`\{\begin{array}{c}\mu S+\overline{s}C=0,\\ s_{\mathrm{\Delta }_c}\mu =0,\end{array}`$ (6)
is given by
$`\{\begin{array}{c}C=\overline{s}D+\nu S,\\ \mu =s_{\mathrm{\Delta }_c}\nu .\end{array}`$ (9)
### B Quantum theory
In the usual version of the BRST-Zinn-Justin-Batalin-Vilkovisky set-up, there are two main issues to be considered (see e.g. ): stability and anomalies.
#### 1 Stability
The problem of stability is the question if to every local BRST cohomological class $`H^{0,n}(s|d)`$ in ghost number $`0`$, there corresponds an independent coupling of the standard master equation. The extended formalism solves this problem by construction, because all standard cohomological classes have been coupled with independent couplings. Indeed, in the extended formalism, the differential controling the “instabilities”, i.e., the divergencies and/or counterterms, is the differential $`\overline{s}`$, and according to the previous section, $`\overline{s}A=0`$ implies $`A=\mu S+\overline{s}B`$, where $`\mu `$ belongs to $`H(s_{\mathrm{\Delta }_c})`$, so that the non trivial part of $`A`$ can indeed be absorbed by a modification of the couplings of the extended master equation (see for more details).
In different words, the extended antifield formalism guarantees “renormalizability in the modern sense” for all gauge theories. Of course, it will be often convenient not to couple all the local BRST cohomological classes but only a subset needed to guarantee that the theory is stable.
#### 2 Anomalous Zinn-Justin equation
In the standard set-up, the question of anomalies is mostly reduced to the computation of the local BRST cohomological group $`H^{1,n}(s|d)`$ in ghost number $`1`$ and to a discussion of the coefficients of the corresponding classes. In the presence of anomalies, there is no differential on the quantum level associated to the anomalously broken Zinn-Justin equation for the effective action. In the extended antifield formalism however, because all the local BRST cohomological classes in positive ghost numbers have been coupled to the solution of the master equation, the broken Zinn-Justin equation can be written as a functional differential equation and an associated differential exists, even in the presence of anomalies. To show this, is the object of the remainder of this subsection <sup>§</sup><sup>§</sup>§We rederive section 4 of in a more appropriate notation, insisting on the existence of the quantum BRST differential in the anomalous case and its relation to its classical counterpart $`\overline{s}`$. Note that the relation after (4.9) of that paper should read $`s_Q\frac{^R}{\xi ^A}\sigma ^A=\frac{1}{2}\frac{^R}{\xi ^A}[\sigma ,\sigma ]^A`$ instead of $`s_Q\frac{^R}{\xi ^A}\sigma ^A=0`$.
The quantum action principle applied to (1) gives
$`{\displaystyle \frac{1}{2}}(\mathrm{\Gamma },\mathrm{\Gamma })+\mathrm{\Delta }_c\mathrm{\Gamma }=\mathrm{}𝒜\mathrm{\Gamma },`$ (10)
where $`\mathrm{\Gamma }`$ is the renormalized generating functional for 1PI vertices associated to the solution $`S`$ of the extended master equation and the local functional $`𝒜`$ is an element of $`F`$ in ghost number $`1`$. Applying $`(\mathrm{\Gamma },)+\mathrm{\Delta }_c^L`$ to (10), the l.h.s vanishes identically because of the graded Jacobi identity for the antibracket and the properties of $`\mathrm{\Delta }_c`$, so that one gets the consistency condition $`(\mathrm{\Gamma },𝒜\mathrm{\Gamma })+\mathrm{\Delta }_c^L𝒜\mathrm{\Gamma }=0`$. To lowest order in $`\mathrm{}`$, this gives $`\overline{s}𝒜=0`$, the general solution of which can be writen as
$`𝒜={\displaystyle \frac{^RS}{\xi ^A}}\mathrm{\Delta }_1^A+\overline{s}\mathrm{\Sigma }_1,`$ (11)
with $`[\mathrm{\Delta }_c,\mathrm{\Delta }_1]=0`$, because of the relation between the $`\overline{s}`$ and the $`s_{\mathrm{\Delta }_c}`$ cohomologies discussed in the previous subsection.
If one now defines $`S^1=S\mathrm{}\mathrm{\Sigma }_1`$, the corresponding generating functional admits the expansion $`\mathrm{\Gamma }^1=\mathrm{\Gamma }\mathrm{}\mathrm{\Sigma }_1+O(\mathrm{}^2)`$ and satisfies $`\frac{1}{2}(\mathrm{\Gamma }^1,\mathrm{\Gamma }^1)+\mathrm{\Delta }^1\mathrm{\Gamma }^1=O(\mathrm{}^2)`$, where $`\mathrm{\Delta }^1=\mathrm{\Delta }_c+\mathrm{}\mathrm{\Delta }_1`$. On the other hand, the quantum action principle applied to $`\frac{1}{2}(S^1,S^1)+\mathrm{\Delta }^1S^1=O(\mathrm{})`$ implies $`\frac{1}{2}(\mathrm{\Gamma }^1,\mathrm{\Gamma }^1)+\mathrm{\Delta }^1\mathrm{\Gamma }^1=\mathrm{}\overline{𝒜}\mathrm{\Gamma }^1`$, for a local functional $`\overline{𝒜}`$. Comparing the two expressions, we deduce that
$`{\displaystyle \frac{1}{2}}(\mathrm{\Gamma }^1,\mathrm{\Gamma }^1)+\mathrm{\Delta }^1\mathrm{\Gamma }^1=\mathrm{}^2𝒜^{}\mathrm{\Gamma }^1,`$ (12)
for a local functional $`𝒜^{}`$. Applying now $`(\mathrm{\Gamma }^1,)+(\mathrm{\Delta }^1)^L`$, one gets as consistency condition
$`{\displaystyle \frac{1}{2}}[\mathrm{\Delta }^1,\mathrm{\Delta }^1]\mathrm{\Gamma }^1+\mathrm{}^2((\mathrm{\Gamma }^1,𝒜^{}\mathrm{\Gamma }^1)+\mathrm{\Delta }^1𝒜^{}\mathrm{\Gamma }^1)=0,`$ (13)
giving to lowest order
$`1/2[\mathrm{\Delta }_1,\mathrm{\Delta }_1]S+\overline{s}𝒜^{}=0.`$ (14)
Since $`1/2[\mathrm{\Delta }_1,\mathrm{\Delta }_1]`$ is a $`s_{\mathrm{\Delta }_c}`$ cocycle because of the graded Jacobi identity for the graded commutator, equations (6) and (9) of the previous section imply that the general solution to this equation is
$`{\displaystyle \frac{1}{2}}[\mathrm{\Delta }_1,\mathrm{\Delta }_1]+[\mathrm{\Delta }_c,\mathrm{\Delta }_2]=0,`$ (15)
$`𝒜^{}=\overline{s}\mathrm{\Sigma }_2{\displaystyle \frac{^RS}{\xi ^B}}\mathrm{\Delta }_2^B.`$ (16)
The redefinition $`S^2=S^1\mathrm{}^2\mathrm{\Sigma }_2`$ then allows to achieve
$`{\displaystyle \frac{1}{2}}(\mathrm{\Gamma }^2,\mathrm{\Gamma }^2)+\mathrm{\Delta }^2\mathrm{\Gamma }^2=\mathrm{}^3𝒜^{\prime \prime }\mathrm{\Gamma }^2,`$ (17)
for a local functional $`𝒜^{\prime \prime }`$, with $`\mathrm{\Delta }^2=\mathrm{\Delta }^1+\mathrm{}^2\mathrm{\Delta }_2`$. The reasoning can be pushed recursively to all orders with the result
$`{\displaystyle \frac{1}{2}}(\mathrm{\Gamma }^{\mathrm{}},\mathrm{\Gamma }^{\mathrm{}})+\mathrm{\Delta }^{\mathrm{}}\mathrm{\Gamma }^{\mathrm{}}=0,`$ (18)
where $`\mathrm{\Gamma }^{\mathrm{}}`$ is associated to the action $`S^{\mathrm{}}=S\mathrm{\Sigma }_{k=1}\mathrm{}^k\mathrm{\Sigma }_k`$ and $`\mathrm{\Delta }^{\mathrm{}}=\mathrm{\Delta }_c+\mathrm{}\mathrm{\Delta }_1+\mathrm{}^2\mathrm{\Delta }_2+\mathrm{}`$ satisfies $`(\mathrm{\Delta }^{\mathrm{}})^2=0`$. The associated quantum BRST differential is
$`s^q=(\mathrm{\Gamma }^{\mathrm{}},)+(\mathrm{\Delta }^{\mathrm{}})^L.`$ (19)
In the limit $`\mathrm{}`$ going to zero, it coincides with the classical differential $`\overline{s}`$.
In the extended antifield formalism, the anomalous Zinn-Justin equation can thus be written as a functional differential equation for the renormalized effective action. The derivations $`\mathrm{\Delta }_1,\mathrm{\Delta }_2,\mathrm{}`$ are guaranteed to exist due to the quantum action principles. They satisfy a priori cohomological restrictions due to the fact that the differential $`\mathrm{\Delta }^{\mathrm{}}`$ is a formal deformation with deformation parameter $`\mathrm{}`$ of the differential $`\mathrm{\Delta }_c`$. In the context of chiral Yang-Mills theories, where $`\mathrm{\Delta }_c=0`$, an anomalous master equation of the form (18) for the renormalized effective action has appeared for the first time in .
#### 3 Renormalization of local BRST cohomological classes of the extended formalism
Associated quantum extension of the classical BRST cohomological classes $`\lambda _0S`$ of the extended formalism are given by $`\lambda _0\mathrm{\Gamma }^{\mathrm{}}`$. By acting with $`\lambda _0`$ on (18), one gets
$`s^q\lambda _0\mathrm{\Gamma }^{\mathrm{}}=[\mathrm{\Delta }^{\mathrm{}},\lambda _0]\mathrm{\Gamma }^{\mathrm{}}.`$ (20)
The derivation $`[\mathrm{\Delta }^{\mathrm{}},\lambda _0]=_{nL}\mathrm{}^n\mu _n`$ on the right hand side is of order at least $`\mathrm{}`$, $`L1`$, because $`[\mathrm{\Delta }_c,\lambda _0]=0`$ and corresponds to the anomaly in the invariant renormalization of $`\lambda _0S`$. The question then is whether there exists modified quantum extension $`\lambda ^{\mathrm{}}\mathrm{\Gamma }^{\mathrm{}}`$, with $`\lambda ^{\mathrm{}}=\lambda _0+\mathrm{}\lambda _1+\mathrm{}`$, such that
$`s^q\lambda ^{\mathrm{}}\mathrm{\Gamma }^{\mathrm{}}=0.`$ (21)
This is for instance the case if $`\lambda _0`$ corresponds to the trivial cohomological class, $`\lambda _0=[\mathrm{\Delta }_c,\nu _0]`$. The searched for extension can then simply be taken to be $`\lambda ^{\mathrm{}}=[\mathrm{\Delta }^{\mathrm{}},\nu _0]`$.
Because $`[\mathrm{\Delta }^{\mathrm{}},[\mathrm{\Delta }^{\mathrm{}},\lambda _0]]=0`$, the lowest order part of the anomaly, $`\mu _L`$ is an $`s_{\mathrm{\Delta }_c}`$ cocyle, $`[\mathrm{\Delta }_c,\mu _L]=0`$. Suppose that $`\mu _L`$ is a trivial solution to this equation, $`\mu _L=[\mathrm{\Delta }_c,\lambda _L]`$. The modified quantum extension $`(\lambda _0+\mathrm{}^L\lambda _L)\mathrm{\Gamma }^{\mathrm{}}`$ allows to push the anomaly to order $`L+1`$.
Hence, to lowest order in $`\mathrm{}`$, the non trivial part of the anomaly in the renormalization of a classical cohomological class $`H^g(s_{\mathrm{\Delta }_c})`$ is constrained to belong to $`H^{g+1}(s_{\mathrm{\Delta }_c})`$. All the quantum information on this anomaly is encoded in the derivations $`\mathrm{\Delta }_1,\mathrm{\Delta }_2,\mathrm{}`$ and the whole discussion has been shifted from local functionals to derivations of functions of the coupling constants.
## III Characterization of local BRST cohomological classes
In this and the following section, we decompose the space $`\mathrm{\Omega }`$ of Lorentz-invariant polynomials or formal powers series in the $`dx^\mu `$, the couplings $`\xi ^A`$, the fields, antifields and their derivatives into the direct sum of the constants and the remaining part, $`\mathrm{\Omega }=𝐑\mathrm{\Omega }_+`$. We have $`\overline{s}\alpha =0=d\alpha `$, for a constant $`\alpha `$ and $`d\mathrm{\Omega }_+\mathrm{\Omega }_+`$. We furthermore assume that if $`\overline{s}\omega =\alpha `$ for a constant $`\alpha `$, then $`\alpha =0`$, which amounts to assuming that the equations of motions are consistent (see the discussion in chapter 9 of ). This means that $`\overline{s}\stackrel{~}{\mathrm{\Omega }}\stackrel{~}{\mathrm{\Omega }}`$. Hence, we can consider the cohomological groups $`H(\overline{s},\mathrm{\Omega }_+)`$, $`H(\overline{s}|d,\mathrm{\Omega }_+)`$ and $`H(d,\mathrm{\Omega }_+)`$.
By analyzing the cohomological groups $`H(\overline{s}|d)`$ (the space $`\mathrm{\Omega }_+`$ being always understood in the following) using descent equations, one can prove that the elements of these groups can be classified into chains of length $`r`$ with an obstruction to further lifts and chains of length $`s`$ whose lifts are unobstructed, i.e., chains with a non trivial element in degree $`n`$ (see also ; we follow here the notations of the review , where explicit proofs of the statements below can be found). More precisely, we have $`H^p(d)=0`$, $`pn1`$, and there exists a basis
$`\{[h_{i_r}^0],[\widehat{h}_{i_r}],[e_{\alpha _s}^0]\}`$ (22)
of $`H(\overline{s})`$, for $`r=0,\mathrm{},n1`$, $`s=0,\mathrm{},n`$, such that a corresponding basis of $`H(\overline{s}|d)`$ is given by
$`\{[h_{i_r}^q],[e_{\alpha _s}^p]\}`$ (23)
for $`q=0,\mathrm{}r`$ and $`p=0,\mathrm{},s`$, with
$`\begin{array}{c}\overline{s}h_{i_r}^{r+1}+dh_{i_r}^r=\widehat{h}_{i_r},\\ \overline{s}h_{i_r}^r+dh_{i_r}^{r1}=0,\\ \mathrm{}\\ \overline{s}h_{i_r}^1+dh_{i_r}^0=0,\\ \overline{s}h_{i_r}^0=0,\end{array}`$ (29)
and
$`\begin{array}{c}\mathrm{form}\mathrm{degree}e_{\alpha _s}^s=n,de_{\alpha _s}^s=0,\\ \overline{s}e_{\alpha _s}^s+de_{\alpha _s}^{s1}=0,\\ \mathrm{}\\ \overline{s}e_{\alpha _s}^1+de_{\alpha _s}^0=0,\\ \overline{s}e_{\alpha _s}^0=0.\end{array}`$ (35)
The cohomological group $`H(\overline{s})`$ can thus be decomposed into elements $`[e_{\alpha _s}^0]`$ that are bottoms of unobstructed chains of length $`s`$, elements $`[h_{i_r}^0]`$ that are bottoms of obstructed chains of length $`r`$ and obstructions $`[\widehat{h}_{i_r}]`$ to chains of length $`r`$.
For the cohomological group $`H(\overline{s}|d)`$, the element $`[h_{i_r}^l]`$ is said to be the element of level $`l`$ of a chain of length $`r`$ with obstruction; it has $`l`$ non trivial descents and $`rl`$ non trivial lifts ; while the element $`[e_{\alpha _s}^l]`$ is said to be the element of level $`l`$ of a chain of length $`s`$ without obstructions, it has $`l`$ non trivial descents and $`sl`$ non trivial lifts.
One can furthermore show that the general solution to a set of descent equations involving at most $`l`$ steps, $`\overline{s}\omega ^l+d\omega ^{l1}=0,`$ $`\overline{s}\omega ^{l1}+d\omega ^{l2}=0,\mathrm{},`$ $`\overline{s}\omega ^0=0`$, can be written in terms of the above elements as
$`\omega ^l={\displaystyle \underset{q=0}{\overset{l}{}}}{\displaystyle \underset{r=lq}{\overset{n1}{}}}\lambda _q^{i_r}h_{i_r}^{lq}+{\displaystyle \underset{p=0}{\overset{l}{}}}{\displaystyle \underset{s=lp}{\overset{n}{}}}\mu _p^{\alpha _s}e_{\alpha _s}^{lp}+\overline{s}\eta ^l+d[\eta ^{l1}+{\displaystyle \underset{r=0}{\overset{n1}{}}}\nu ^{(l)i_r}h_{i_r}^r],`$ (36)
with $`\eta ^1=0`$. This means that a $`\overline{s}`$ modulo $`d`$ cocycle which has $`l`$ non trivial descents, is a linear combination of all elements of the chains (29) and (35) which have $`l`$ or less non trivial descents. If such a linear combination is $`\overline{s}`$ modulo $`d`$ trivial, the coefficients of the linear combination must vanish, i.e.,
$`{\displaystyle \underset{q=0}{\overset{l}{}}}{\displaystyle \underset{r=lq}{\overset{n1}{}}}\lambda _q^{i_r}h_{i_r}^{lq}+{\displaystyle \underset{p=0}{\overset{l}{}}}{\displaystyle \underset{s=lp}{\overset{n}{}}}\mu _p^{\alpha _s}e_{\alpha _s}^{lp}=\overline{s}()+d().`$ (37)
implies that $`\lambda _q^{i_r}=0=\mu _p^{\alpha _s}`$.
## IV Lengths of descent and lifts of lowest order anomalies
Let us now investigate the anomalies in the BRST invariant renormalization of a chain (35) of length $`s`$ without obstructions. We follow the approach of , which consists in considering simultaneously the anomalies for a whole chain of descent equations (for a review, see ). Using the quantum action principles, one can prove (as in lemma 1 in the appendix of ) that if $`a`$ is a local form, then
$`\overline{s}a\mathrm{\Gamma }^{\mathrm{}}=s^qa\mathrm{\Gamma }^{\mathrm{}}+\mathrm{}b\mathrm{\Gamma }^{\mathrm{}},`$ (38)
for some local form $`b`$. When applied to the chain (35), we find that the quantum version of this chain is
$`\begin{array}{c}\mathrm{form}\mathrm{degree}e_{\alpha _s}^s\mathrm{\Gamma }^{\mathrm{}}=n,de_{\alpha _s}^s\mathrm{\Gamma }^{\mathrm{}}=0,\\ s^qe_{\alpha _s}^s\mathrm{\Gamma }^{\mathrm{}}+de_{\alpha _s}^{s1}\mathrm{\Gamma }^{\mathrm{}}=\mathrm{}a_{\alpha _s}^s\mathrm{\Gamma }^{\mathrm{}},\\ \mathrm{}\\ s^qe_{\alpha _s}^1\mathrm{\Gamma }^{\mathrm{}}+de_{\alpha _s}^0\mathrm{\Gamma }^{\mathrm{}}=\mathrm{}a_{\alpha _s}^1\mathrm{\Gamma }^{\mathrm{}},\\ s^qe_{\alpha _s}^0\mathrm{\Gamma }^{\mathrm{}}=\mathrm{}a_{\alpha _s}^0\mathrm{\Gamma }^{\mathrm{}},\end{array}`$ (44)
where $`a_{\alpha _s}^l\mathrm{\Gamma }^{\mathrm{}}=a_{\alpha _s}^s+O(\mathrm{})`$ for a local function $`a_{\alpha _s}^l`$. Applying $`s^q`$, we get the consistency condition,
$`\begin{array}{c}\mathrm{form}\mathrm{degree}a_{\alpha _s}^s\mathrm{\Gamma }^{\mathrm{}}=n,da_{\alpha _s}^s\mathrm{\Gamma }^{\mathrm{}}=0\\ s^qa_{\alpha _s}^s\mathrm{\Gamma }^{\mathrm{}}+da_{\alpha _s}^{s1}\mathrm{\Gamma }^{\mathrm{}}=0,\\ \mathrm{}\\ s^qa_{\alpha _s}^1\mathrm{\Gamma }^{\mathrm{}}+da_{\alpha _s}^0\mathrm{\Gamma }^{\mathrm{}}=0,\\ s^qa_{\alpha _s}^0\mathrm{\Gamma }^{\mathrm{}}=0.\end{array}`$ (50)
At lowest order in $`\mathrm{}`$, we get
$`\begin{array}{c}\mathrm{form}\mathrm{degree}a_{\alpha _s}^s=n,da_{\alpha _s}^s=0\\ \overline{s}a_{\alpha _s}^s+da_{\alpha _s}^{s1}=0,\\ \mathrm{}\\ \overline{s}a_{\alpha _s}^1+da_{\alpha _s}^0=0,\\ \overline{s}a_{\alpha _s}^0=0.\end{array}`$ (56)
Using equation (36) and the fact that the form degree of $`a_{\alpha _s}^s`$ is $`n`$, it follows that
$`a_{\alpha _s}^l={\displaystyle \underset{p=0}{\overset{l}{}}}\mu _{p}^{\beta _{sp}}{}_{\alpha _s}{}^{}e_{\beta _{sp}}^{lp}+\overline{s}\eta _{\alpha _s}^l+d[\eta _{\alpha _s}^{l1}+{\displaystyle \underset{r0}{}}\nu _{\alpha _s}^{(l)i_r}h_{i_r}^r],`$ (57)
for $`l=0,\mathrm{},s`$. This gives our first result:
The anomaly in the renormalization of an element of level $`l`$ of a chain of length $`s`$ without obstructions involves at most elements of chains of the same type with less non trivial descents and the same number of non trivial lifts.
For the anomalies for a chain with obstruction (29), we get,
$`\begin{array}{c}s^qh_{i_r}^{r+1}\mathrm{\Gamma }^{\mathrm{}}+dh_{i_r}^r\mathrm{\Gamma }^{\mathrm{}}=\widehat{h}_{i_r}\mathrm{\Gamma }^{\mathrm{}}+\mathrm{}a_{i_r}^{r+1}\mathrm{\Gamma }^{\mathrm{}},\\ s^qh_{i_r}^r\mathrm{\Gamma }^{\mathrm{}}+dh_{i_r}^{r1}\mathrm{\Gamma }^{\mathrm{}}=\mathrm{}a_{i_r}^r\mathrm{\Gamma }^{\mathrm{}},\\ \mathrm{}\\ s^qh_{i_r}^1\mathrm{\Gamma }^{\mathrm{}}+dh_{i_r}^0\mathrm{\Gamma }^{\mathrm{}}=\mathrm{}a_{i_r}^1\mathrm{\Gamma }^{\mathrm{}},\\ s^qh_{i_r}^0\mathrm{\Gamma }^{\mathrm{}}=\mathrm{}a_{i_r}^0\mathrm{\Gamma }^{\mathrm{}},\end{array}`$ (63)
We also have
$`s^q\widehat{h}_{i_r}\mathrm{\Gamma }^{\mathrm{}}=\mathrm{}\widehat{a}_{i_r}\mathrm{\Gamma }^{\mathrm{}}.`$ (64)
Applying $`s^q`$ gives $`s^q\widehat{a}_{i_r}\mathrm{\Gamma }^{\mathrm{}}=0`$ and then to lowest order, $`\overline{s}\widehat{a}_{i_r}=0`$. Applying now $`s^q`$ to the chain (63) gives
$`\begin{array}{c}s^qa_{i_r}^{r+1}\mathrm{\Gamma }^{\mathrm{}}+da_{i_r}^r\mathrm{\Gamma }^{\mathrm{}}=\widehat{a}_{i_r}\mathrm{\Gamma }^{\mathrm{}},\\ s^qa_{i_r}^r\mathrm{\Gamma }^{\mathrm{}}+da_{i_r}^{r1}\mathrm{\Gamma }^{\mathrm{}}=0,\\ \mathrm{}\\ s^qa_{i_r}^1\mathrm{\Gamma }^{\mathrm{}}+da_{i_r}^1\mathrm{\Gamma }^{\mathrm{}}=0,\\ s^qa_{i_r}^0\mathrm{\Gamma }^{\mathrm{}}=0,\end{array}`$ (70)
and to lowest order,
$`\begin{array}{c}\overline{s}a_{i_r}^{r+1}+da_{i_r}^r=\widehat{a}_{i_r},\\ \overline{s}a_{i_r}^r+da_{i_r}^{r1}=0,\\ \mathrm{}\\ \overline{s}a_{i_r}^1+da_{i_r}^1=0,\\ \overline{s}a_{i_r}^0=0,\end{array}`$ (76)
On the one hand, it follows from (36) that
$`a_{i_r}^l={\displaystyle \underset{q=0}{\overset{l}{}}}{\displaystyle \underset{r^{}=rq}{\overset{n1}{}}}\lambda _{q}^{j_r^{}}{}_{i_r}{}^{}h_{j_r^{}}^{lq}+{\displaystyle \underset{p=0}{\overset{l}{}}}{\displaystyle \underset{s=rp+1}{\overset{n}{}}}\mu _{p}^{\beta _s}{}_{i_r}{}^{}e_{\beta _s}^{lp}+\overline{s}\eta _{i_r}^l+d[\eta _{i_r}^{l1}+{\displaystyle \underset{r^{}=0}{\overset{n1}{}}}\nu _{i_r}^{(l)j_r^{}}h_{j_r^{}}^r^{}],`$ (77)
for $`l=0,\mathrm{},r`$, while on the other hand, the cohomology of $`\overline{s}`$ implies
$`\widehat{a}_{i_r}={\displaystyle \underset{r^{}=0}{\overset{n1}{}}}\alpha _{i_r}^{j_r^{}}\widehat{h}_{j_r^{}}+{\displaystyle \underset{r^{}=0}{\overset{n1}{}}}\beta _{i_r}^{j_r^{}}h_{i_r}^0+{\displaystyle \underset{s=0}{\overset{n}{}}}\gamma _{i_r}^{\alpha _s}e_{\alpha _s}^0+\overline{s}\widehat{o}_{i_r}.`$ (78)
Applying $`d`$ to (77) at $`l=r`$ gives
$`da_{i_r}^r={\displaystyle \underset{q=0}{\overset{r}{}}}{\displaystyle \underset{r^{}=rq+1}{\overset{n1}{}}}\lambda _{q}^{i_r^{}}{}_{i_r}{}^{}\overline{s}h_{i_r^{}}^{rq+1}{\displaystyle \underset{p=0}{\overset{r}{}}}{\displaystyle \underset{s=rp+1}{\overset{n}{}}}\mu _{p}^{\beta _s}{}_{i_r}{}^{}\overline{s}e_{\beta _s}^{rp+1}\overline{s}d\eta _{i_r}^r`$ (79)
$`+{\displaystyle \underset{q=0}{\overset{r}{}}}\lambda _{q}^{i_{rq}}{}_{i_r}{}^{}(\overline{s}h_{i_{rq}}^{rq+1}+\widehat{h}_{i_{rq}}).`$ (80)
Injecting now (78) and (80) into the first equation of (76) gives first of all $`\beta _{i_r}^{j_r^{}}=0=\gamma _{i_r}^{\alpha _s}`$ and also
$`\widehat{a}_{i_r}={\displaystyle \underset{q=0}{\overset{r}{}}}\lambda _{q}^{j_{rq}}{}_{i_r}{}^{}\widehat{h}_{j_{rq}}+\overline{s}\widehat{o}_{i_r},`$ (81)
and then,
$`\overline{s}(a_{i_r}^{r+1}{\displaystyle \underset{q=0}{\overset{r}{}}}{\displaystyle \underset{r^{}=rq+1}{\overset{n1}{}}}\lambda _{q}^{i_r^{}}{}_{i_r}{}^{}h_{i_r^{}}^{rq+1}{\displaystyle \underset{p=0}{\overset{r}{}}}{\displaystyle \underset{s=rp+1}{\overset{n}{}}}\mu _{p}^{\beta _s}{}_{i_r}{}^{}e_{\beta _s}^{rp+1}`$ (82)
$`{\displaystyle \underset{q=0}{\overset{r}{}}}\lambda _{q}^{i_{rq}}{}_{i_r}{}^{}h_{i_{rq}}^{rq+1}d\eta _{i_r}^r\widehat{o}_{i_r})=0,`$ (83)
so that, using the cohomology of $`\overline{s}`$,
$`a_{i_r}^{r+1}={\displaystyle \underset{q=0}{\overset{r+1}{}}}{\displaystyle \underset{r^{}=rq}{\overset{n1}{}}}\lambda _{q}^{i_r^{}}{}_{i_r}{}^{}h_{i_r^{}}^{rq+1}+{\displaystyle \underset{p=0}{\overset{r+1}{}}}{\displaystyle \underset{s=rp+1}{\overset{n}{}}}\mu _{p}^{\beta _s}{}_{i_r}{}^{}e_{\beta _s}^{rp+1}`$ (84)
$`+d[\eta _{i_r}^r+{\displaystyle \underset{r^{}=0}{\overset{n1}{}}}\nu _{i_r}^{(r+1)j_r^{}}h_{j_r^{}}^r^{}]+\widehat{o}_{i_r}+\overline{s}\eta _{i_r}^{r+1}.`$ (85)
Our second result is then:
The anomaly in the renormalization of an element of level $`l`$ of a chain of length $`r`$ with obstructions involves at most elements of chains with obstructions with less non trivial descents and more non trivial lifts and elements of chains without obstructions with less non trivial descents and strictly more non trivial lifts.
Let us now rewrite (57) and (77) at $`l`$=0 as
$`a_{\alpha _s}^0=\mu _{0}^{\beta _s}{}_{\alpha _s}{}^{}e_{\beta _s}^0+\overline{s}[\eta _{\alpha _s}^0{\displaystyle \underset{r=0}{\overset{n1}{}}}\nu _{\alpha _s}^{(0)i_r}h_{i_r}^{r+1}]+{\displaystyle \underset{r=0}{\overset{n1}{}}}\nu _{\alpha _s}^{(0)i_r}\widehat{h}_{i_r},`$ (86)
$`a_{i_r}^0={\displaystyle \underset{r^{}=r}{\overset{n1}{}}}\lambda _{0}^{j_r^{}}{}_{i_r}{}^{}h_{j_r^{}}^0+{\displaystyle \underset{s=r+1}{\overset{n}{}}}\mu _{0}^{\beta _s}{}_{i_r}{}^{}e_{\beta _s}^0+\overline{s}[\eta _{i_r}^0{\displaystyle \underset{r^{}=0}{\overset{n1}{}}}\nu _{i_r}^{(0)j_r^{}}h_{j_r^{}}^{r^{}+1}]+{\displaystyle \underset{r^{}=0}{\overset{n1}{}}}\nu _{i_r}^{(0)j_r^{}}\widehat{h}_{j_r^{}}.`$ (87)
Combined with (81), our third result on anomalies in the renormalization of elements of $`H(\overline{s})`$ is accordingly:
The anomaly in the renormalization of obstructions to chains of length $`r`$ involves at most obstructions to shorther chains; the anomaly in the renormalization of bottoms of unobstructed chains of length $`s`$ involves at most bottoms of unobstructed chains of the same length and obstructions to chains of all possible lengths; the anomaly in the renormalization of bottoms of obstructed chains of length $`r`$ involves at most bottoms of obstructed chains of greater length, bottoms of unobstructed chains of strictly greater length and obstructions to chains of all possible lengths.
## V Differentials associated to one loop anomalies
A different, more compact, way to formulate and prove the results of section III and IV is to use the exact couple describing the descent and the associated spectral sequence .
Indeed, the diagram
$`\begin{array}{ccc}H(\overline{s}|d)\stackrel{𝒟}{}H(\overline{s}|d)& & \\ i_0l_0& & \\ H(\overline{s})& & \end{array}`$ (91)
can be shown to be exact at all corners. The various maps are defined as follows: $`i_0`$ is the map which consists in regarding an element of $`H(\overline{s})`$ as an element of $`H(\overline{s}|d)`$, $`i_0:H(\overline{s})H(\overline{s}|d)`$, with $`i_0[a]=[a]`$. It is well defined because every $`\overline{s}`$ cocycle is a $`\overline{s}`$ cocycle modulo $`d`$ and every $`\overline{s}`$ coboundary is a $`\overline{s}`$ coboundary modulo $`d`$. The descent homomorphism $`𝒟:H^{k,l}(\overline{s}|d)H^{k+1,l1}(\overline{s}|d)`$ with $`𝒟[a]=[b]`$, if $`\overline{s}a+db=0`$ is well defined because of the triviality of the cohomology of $`d`$ in form degree $`pn1`$. Finally, the map $`l_0:H^{k+1,l1}(\overline{s}|d)H^{k+1,l}(\overline{s})`$ is defined by $`l_0[a]=[da]`$. It is well defined because the relation $`\{\overline{s},d\}=0`$ implies that it maps cocycles to cocycles and coboundaries to coboundaries.
Associated to such an exact couple $`(H(\overline{s}|d),K_0=H(\overline{s}))`$, one can associate in a standard way derived exact couples $`(𝒟^rH(\overline{s}|d),K_r)`$,
$`\begin{array}{ccc}𝒟^rH(\overline{s}|d)\stackrel{𝒟}{}𝒟^rH(\overline{s}|d)& & \\ i_rl_r& & \\ K_r,& & \end{array}`$ (95)
and a spectral sequence $`K_{r+1}=H(d_r,K_r)`$, with $`K_0H(\overline{s})`$. The maps of these exact couples are defined recursively as follows: the map $`d_{r1}=l_{r1}i_{r1}`$ can be shown to be a differential, the map $`i_r`$ is the map induced by $`i_{r1}`$ in $`K_r`$, while $`l_r𝒟[a]=l_{r1}[a]`$.
Explicitly, the differential $`d_0:K_0K_0`$ is defined by $`d_0[a]=[da]`$, where $`\overline{s}a=0`$. An element $`k_rK_r`$ is identified with the equivalence class $`[a]_r`$ of an element $`[a]H(\overline{s})`$, where $`[a]_r[a^{}]`$ if $`[a][a^{}]_{q=0}^{r1}\mathrm{im}d_q`$. The relations $`d_qk_r=0`$, $`q=0,\mathrm{},r1`$ mean that $`k_r`$ is a bottom that can be lifted at least $`r`$ times, i.e., there exist $`c_1,\mathrm{},c_{r+1}`$ such that $`\overline{s}a=0,da+\overline{s}c_1=0,\mathrm{},dc_{r1}+\overline{s}c_r=0`$. Then, the differential $`d_r`$ is defined by $`d_rk_r=[dc_r]_r`$.
Because there are no forms of form degree higher than $`n`$, $`𝒟^{n+1}H(\overline{s}|d)=0`$ and $`d_n0`$ so that the construction stops at $`r=n`$.
The space of local forms $`\mathrm{\Omega }`$ is decomposed as $`\mathrm{\Omega }=E_0G\overline{s}G𝐑`$, with $`E_0K_0=H(\overline{s})`$. If we define $`E_r,F_{r1}E_{r1}`$ through $`E_{r1}=E_{r1}E_rd_{r1}F_{r1}`$ with $`E_rK_r`$, we get the decomposition
$`E_0=F_0\mathrm{}F_{n1}E_nd_{r1}F_{r1}\mathrm{}d_0F_0.`$ (96)
The $`e_{\alpha _s}^0`$ are elements of a basis of $`E_n`$ that can be lifted $`s`$ times before hitting form degree $`n`$, i.e., that are of form degree $`ns`$, while $`\widehat{h}_{i_r}`$ and $`h_{i_r}^0`$ are elements of a basis of $`d_rF_r`$ and $`F_r`$ respectively. This sums up the results of section III.
Let us now define the linear map
$`\delta _0:H^g(\overline{s})H^{g+1}(\overline{s}),`$ (97)
$`\delta _0[a]=[b],\text{where}s^qa\mathrm{\Gamma }^{\mathrm{}}=\mathrm{}b\mathrm{\Gamma }^{\mathrm{}}.`$ (98)
The map associates to a given BRST cohomological class the non trivial order $`\mathrm{}`$, i.e., 1 loop contribution of its anomaly.
The map is well defined, because the consistency condition implies that $`\overline{s}b=0`$, and if $`a=\overline{s}c`$, $`a\mathrm{\Gamma }^{\mathrm{}}=s^qc\mathrm{\Gamma }^{\mathrm{}}+\mathrm{}d\mathrm{\Gamma }^{\mathrm{}}`$, so that $`s^qa\mathrm{\Gamma }^{\mathrm{}}=\mathrm{}s^qd\mathrm{\Gamma }^{\mathrm{}}`$, meaning that $`b=\overline{s}d`$, so that the map does not depend on the choice of the representative. Furthermore, this map is a differential
$`\delta _0^2=0.`$ (99)
Indeed, if $`[a]=\delta _0[c]`$, we have $`a\mathrm{\Gamma }^{\mathrm{}}=\frac{1}{\mathrm{}}s^qc\mathrm{\Gamma }^{\mathrm{}}`$. It follows that $`s^qa\mathrm{\Gamma }^{\mathrm{}}=0`$. A BRST cohomological class which is a $`\delta _0`$-cocycle has no 1-loop anomaly, while a BRST cohomological class which is a $`\delta _0`$-coboundary is the 1-loop anomaly of some other BRST cohomological class.
We thus have two differentials in $`K_0=H(\overline{s})`$, $`d_0`$ introduced above and $`\delta _0`$. These differentials anticommute,
$`\{d_0,\delta _0\}=0.`$ (100)
Indeed, if $`s^qa\mathrm{\Gamma }^{\mathrm{}}=\mathrm{}b\mathrm{\Gamma }^{\mathrm{}}`$, $`d_0\delta _0[a]=d_0[b]=[db]`$, while $`\delta _0d_0[a]=\delta _0[da]`$, and $`s^q(da\mathrm{\Gamma }^{\mathrm{}})=s^qd(a\mathrm{\Gamma }^{\mathrm{}}+\mathrm{}c\mathrm{\Gamma }^{\mathrm{}})=ds^q(a\mathrm{\Gamma }^{\mathrm{}}+\mathrm{}c\mathrm{\Gamma }^{\mathrm{}})=\mathrm{}d(b\mathrm{\Gamma }^{\mathrm{}}+s^qc\mathrm{\Gamma }^{\mathrm{}})`$, so that $`\delta _0[da]=[d(b+\overline{s}c)]=[db]`$.
The relation $`d_0\delta _0[a]=\delta _0d_0[a]`$ means:
* if $`[a]`$ belongs to $`\mathrm{im}d_0`$, $`[a]=d_0[b]`$, then $`\delta _0[a]=d_0\delta _0[b]`$, i.e., if $`[a]`$ represents an obstruction to the lift of an element $`[b]`$, its anomaly represents minus the obstruction to the lift of the anomaly of $`[b]`$,
* if $`d_0[a]=0`$, then $`d_0\delta _0[a]=0`$, i.e., if $`[a]`$ can be lifted, then so does its anomaly $`\delta _0[a]`$,
* the anomaly in a bottom $`[a]`$ of $`K_0`$ that cannot be lifted is minus the anomaly of the corresponding obstruction, up to elements that can be lifted.
If we organize the space $`E_0K_0=H(\overline{s})`$ as $`E_0=F_0E_1d_0F_0`$, with $`E_1K_1`$, we have shown that 1-loop contribution to the anomaly of an element in one of these subspaces belongs to the same subspace or to a subspace that stands to the right. Together with the last point of the previous list, this sums up the results for the elements of $`H(\overline{s})`$, i.e., for the obstructions and the bottoms contained in (81), (86) and (87) restricted to $`r=0`$.
In the same way, these results for all $`r`$ and $`s`$ follow from the fact that $`\delta _0`$ induces a well-defined differential (also called $`\delta _0`$ in the following) in the spaces $`K_r`$, anticommuting with $`d_r`$, $`\delta _0:K_rK_r`$ with $`\{\delta _0,d_r\}=0`$.
Indeed, suppose this result to be true for $`K_0,\mathrm{},K_{r1}`$, $`d_0,\mathrm{},d_{r1}`$. An element $`[a]_rK_r`$ satisfies $`\overline{s}a=0`$, $`da+\overline{s}c_1=0`$, $`\mathrm{},`$ $`dc_{r1}+\overline{s}c_r=0`$. This implies $`s^qa\mathrm{\Gamma }^{\mathrm{}}=\mathrm{}b\mathrm{\Gamma }^{\mathrm{}}`$, $`da\mathrm{\Gamma }^{\mathrm{}}+s^qc_1\mathrm{\Gamma }^{\mathrm{}}=\mathrm{}f_1\mathrm{\Gamma }^{\mathrm{}}`$, $`\mathrm{},`$ $`dc_{r1}\mathrm{\Gamma }^{\mathrm{}}+s^qc_r\mathrm{\Gamma }^{\mathrm{}}=\mathrm{}f_r\mathrm{\Gamma }^{\mathrm{}}`$. Applying $`s^q`$ gives to lowest order $`\overline{s}b=0`$, $`db+\overline{s}f_1=0`$, $`\mathrm{},`$ $`df_{r1}+\overline{s}f_r=0`$, so that $`[b]_r`$ is well defined. Suppose now that $`[a]_r=d_0[g_0]_0+\mathrm{}+d_{r1}[g_{r1}]_{r1}`$. Anticommutativity of $`\delta _0`$ with $`d_0,\mathrm{},d_{r1}`$ then implies that $`\delta _0[a]_r=0`$. Hence, $`\delta _0`$ does not depend on the representative and is well defined in $`K_r`$. Finally, $`d_r\delta _0[a]_r=d_r[b]_r=[df_r]_r`$, while $`\delta _0d_r[a]_r=\delta _0[dc_r]_r`$, and $`s^q(dc_r)\mathrm{\Gamma }^{\mathrm{}}=s^qd(c_r\mathrm{\Gamma }^{\mathrm{}}+\mathrm{}c^{}\mathrm{\Gamma }^{\mathrm{}})=d(s^qc_r\mathrm{\Gamma }^{\mathrm{}}+\mathrm{}s^qc^{}\mathrm{\Gamma }^{\mathrm{}})`$, so that $`\delta _0[dc_r]_r=[d(f_r+\overline{s}c^{})]_r=d_r\delta _0[a]_r`$.
The results (81), (86) and (87) can then be summarized by the statement that an anomaly in one of the subspaces of the decomposition (96) must belong to the same subspace or to one that stands to the right; furthermore, the part of the anomaly of an element of $`F_i`$ in $`F_i`$ is minus the part of the anomaly of the corresponding element of $`d_iF_i`$ in $`d_iF_i`$.
In order to recover the results for elements of $`H(\overline{s}|d)`$, we define $`\mathrm{\Delta }_0`$ to be the equivalent of $`\delta _0`$ for modulo $`d`$ BRST cohomological classes, $`\mathrm{\Delta }_0[a]=[b]`$, for $`[a],[b]H(\overline{s}|d)`$, where $`\overline{s}a+dm=0`$, $`\overline{s}m+du=0`$, $`s^qa\mathrm{\Gamma }^{\mathrm{}}+d(m\mathrm{\Gamma }^{\mathrm{}})=\mathrm{}b\mathrm{\Gamma }^{\mathrm{}}`$, $`s^qm\mathrm{\Gamma }^{\mathrm{}}+d(u\mathrm{\Gamma }^{\mathrm{}})=\mathrm{}n\mathrm{\Gamma }^{\mathrm{}}`$. Indeed, the map is well defined because the consistency condition implies to lowest order $`\overline{s}b+dn=0`$, while if $`a=\overline{s}c+dg`$, we have $`m=\overline{s}g+du`$, so that $`s^qa\mathrm{\Gamma }^{\mathrm{}}+dm\mathrm{\Gamma }^{\mathrm{}}=\mathrm{}b\mathrm{\Gamma }^{\mathrm{}}`$ gives $`s^q(s^qc\mathrm{\Gamma }^{\mathrm{}}+dg\mathrm{\Gamma }^{\mathrm{}}+\mathrm{}f\mathrm{\Gamma }^{\mathrm{}})+d(s^qg\mathrm{\Gamma }^{\mathrm{}}+du\mathrm{\Gamma }^{\mathrm{}}+\mathrm{}v\mathrm{\Gamma }^{\mathrm{}})=\mathrm{}b\mathrm{\Gamma }^{\mathrm{}}`$, which implies $`b=\overline{s}f+dv`$ as it should.
The following properties are straightforward to check: $`[\mathrm{\Delta }_0,𝒟]=0`$, $`l\mathrm{\Delta }_0=\delta _0l`$, $`i_0\delta _0=\mathrm{\Delta }_0i_0`$. One says (see e.g. , Chapter VIII.9) that $`(\mathrm{\Delta }_0,\delta _0)`$ is a mapping of the exact couple $`(H(\overline{s}|d),H(\overline{s}))`$.
The previous result, that $`\delta _0`$ induces well defined maps in the spaces of the spectral sequence anticommuting with the differentials $`d_r`$, follows directly from the way the spectral sequence is associated to an exact couple. The relation between (57), (77) at $`l=0`$ and (86), (87) is summarized by $`i_0\delta _0=\mathrm{\Delta }_0i_0`$ ; the relations between (57), (77) at different values of $`l`$ are summarized by $`[\mathrm{\Delta }_0,𝒟]=0`$ ; finally, the relation between (77) at $`l=r`$ and (81) is summarized by $`l\mathrm{\Delta }_0=\delta _0l`$.
Note that in this case, we have furthermore the property that $`\mathrm{\Delta }_0`$ is a differential, $`\mathrm{\Delta }_0^2=0`$.
Remark: It follows from the above analysis that the relevant property of the differentials $`d_r`$ is $`\{\delta _0,d_r\}=0`$. This means that analogous results that constrain the anomalies to belong to particular subspaces of $`H(\overline{s})`$ or $`H(\overline{s}|d)`$ can be derived if one can find maps $`\lambda _0:H(\overline{s})H(\overline{s})`$, respectively $`\mathrm{\Lambda }_0:H(\overline{s}|d)H(\overline{s}|d)`$ such that $`[\delta _0,\lambda _0]=0`$, respectively $`[\mathrm{\Delta }_0,\mathrm{\Lambda }_0]=0`$.
## VI Adler-Bardeen theorem revisited
We now apply the ideas of the extended antifield formalism to standard Yang-Mills theory. In this case, it is sufficient for our purpose to couple the local BRST cohomology classes in ghost number $`0`$ and ghost number $`1`$, because this will be enough, under some assumptions stated explicitly below, to guarantee stability and to control the anomalies. The starting point action contains from the beginning a coupling to the Adler-Bardeen anomaly as in , with additional couplings to possibly higher dimensional gauge invariant operators, if one does not want to restrict oneself to the power counting renormalizable case . More precicley, the starting point is the action
$`S_\rho ={\displaystyle d^4x[\frac{1}{4g^2}F_I^{\mu \nu }F_{\mu \nu }^I+L_{\mathrm{matter}}^{\mathrm{kin}}(\psi ^i,D_\mu \psi ^i)]}`$ (101)
$`+{\displaystyle d^4x[D_\mu C^IA_I^\mu +C^IT_{Ij}^i\psi ^j\psi _i^{}\frac{1}{2}C^IC^Jf_{JI}^{}{}_{}{}^{K}C_K^{}]}`$ (102)
$`+g^i{\displaystyle d^4x𝒪_i}+𝒜\rho ,`$ (103)
satisfying the master equation
$`{\displaystyle \frac{1}{2}}(S_\rho ,S_\rho )=0.`$ (104)
The Lagrangian $`L_{\mathrm{matter}}^{\mathrm{kin}}(\psi ^i,D_\mu \psi ^i)`$ is the gauge invariant extension of the kinetic terms for the matter fields $`\psi ^i`$. For simplicity, we assume the gauge group to be $`SU(3)`$. The $`𝒪_i`$ are gauge invariant local functions built out of the field strengths $`F_{\mu \nu }^I`$, the matter fields $`\psi ^i`$ and their covariant derivatives such that the $`d^4x𝒪_i`$ (which can, but need not, be assumed to be power counting renormalizable) and $`d^4x\frac{1}{4g^2}F_I^{\mu \nu }F_{\mu \nu }^I`$ are linearily independent even when the gauge covariant equations of motions hold. Finally, $`𝒜=\mathrm{Tr}[Cd(AdA+\frac{1}{2}A^3)]`$ is the Adler-Bardeen gauge anomaly, $`g`$ is the gauge coupling constant, $`g^i`$ are coupling constants for the other gauge invariant operators, while $`\rho `$ is a Grassmann odd coupling constant with ghost number $`1`$ for the Adler-Bardeen anomaly. In this particular case, $`\mathrm{\Delta }_c=0`$. This can be traced back to the fact that all representatives of the local BRST cohomological classes in ghost number $`0`$ and $`1`$ can be choosen to be independent of the antifields.
The gauge is fixed by introducing the cohomologically trivial non minimal sector consisting of the antighost $`\overline{C}^I`$ and the Lagrange multiplier $`B^I`$ and their antifields. One adds to the action (103) the term $`d^4x\overline{C}_I^{}B^I`$ and chooses an appropriate gauge fixing fermion $`\mathrm{\Psi }`$, which is used to generate an anticanonical transformation in the fields and antifields such that the propagators of the theory are well defined. The gauge fixing is irrelevant for the cohomological considerations below.
For the question of stability and anomalies, we have to analyze the cohomology $`H^{0,4}(s_\rho |d)`$ and $`H^{1,4}(s_\rho |d)`$ in the space of functions in the couplings $`g,g^i,\rho `$ with coefficients that are Lorentz invariant polynomials in the $`dx^\mu `$, the fields, the antifields and their derivatives.
In order to compute this cohomology, we decompose, as in , both the BRST differential $`s_\rho `$ and the local forms into parts independent of $`\rho `$ and parts linear in $`\rho `$. Explicitly, $`s_\rho =s_0+s_1`$, where $`s`$ is the standard BRST differential associated to the solution $`S_{\rho =0}`$ of the master equation, while $`s_1=(𝒜\rho ,)`$. The $`\rho `$ independent part of the cocycle condition $`s_\rho \omega (\rho )+d\eta (\rho )=0`$ in form degree $`4`$ gives (see e.g. for a review) $`\omega (0)=\alpha (g,g^i)d^4xF_I^{\mu \nu }F_{\mu \nu }^I+\alpha ^j(g,g^i)d^4x𝒪_i+s()+d()`$. This implies $`\omega (\rho )=\alpha (g,g^i)d^4xF_I^{\mu \nu }F_{\mu \nu }^I+\alpha ^j(g,g^i)d^4x𝒪_i+\omega _1^{}\rho +s_\rho ()+d().`$ Because $`d^4xF_I^{\mu \nu }F_{\mu \nu }^I`$ and $`d^4x𝒪_i`$ are also $`s_1`$ closed and $`\rho ^2=0`$, the cocycle condition reduces to $`s\omega _1^{}+d()=0`$, where the ghost number of $`\omega _1^{}`$ is $`1`$. It follows that $`\omega _1^{}=\lambda (g,g^i)\mathrm{Tr}[Cd(AdA+\frac{1}{2}A^3)]+s()+d()`$. Hence, the general solution of the consistency condition in the space of local functionals in ghost number $`0`$ is given by
$`\alpha (g,g^i){\displaystyle \frac{S_\rho }{g}}+\alpha ^j(g,g^i){\displaystyle \frac{S_\rho }{g^i}}+{\displaystyle \frac{^RS_\rho }{\rho }}\rho \lambda (g,g^i)+(S_\rho ,\mathrm{\Xi }_\rho )`$ (105)
for some local functional $`\mathrm{\Xi }_\rho `$ in ghost number $`1`$. This implies that the theory is stable.
Similarily, in ghost number $`1`$, the $`\rho `$ independent part of the cohomology gives as only anomaly candidate the Adler-Bardeen anomaly. There could however be a $`\rho `$ linear non trivial contribution to the anomaly, because the cohomology of $`s`$ in form degree $`4`$ and ghost number $`2`$ is not necessarily empty (see , section 12.4), contrary to the claim in . More precisely, to each $`x^\mu `$-independent, gauge and Lorentz invariant non trivial conserved current $`j_\mathrm{\Delta }=j_\mathrm{\Delta }^\mu ϵ_{\mu \nu _1\nu _2\nu _3}dx^{\nu _1}dx^{\nu _2}dx^{\nu _3}`$, there corresponds the cohomological class $`V_\mathrm{\Delta }^{2,4}=j_\mathrm{\Delta }[\mathrm{Tr}C^3]^1+\mathrm{𝑎𝑛𝑡𝑖𝑓𝑖𝑒𝑙𝑑}\mathrm{𝑑𝑒𝑝𝑒𝑛𝑑𝑒𝑛𝑡}\mathrm{𝑡𝑒𝑟𝑚𝑠}`$ in $`H^{2,n}(s|d)`$, with $`s[\mathrm{Tr}C^3]^1+\mathrm{Tr}C^3=0`$.(There could in principle be another type of antifield dependent cohomology classes in exceptional situations , which we exclude from the present considerations).
If there are such non trivial currents $`j_\mathrm{\Delta }`$, we have to change our starting point and also couple the “anomaly for anomaly candidates” $`V^{2,4}`$ from the beginning with couplings in ghost number $`2`$. But then the cohomology of $`s`$ in ghost number $`3`$ becomes relevant. There are plenty of such classes, for instance classes of the form $`d^4xI\mathrm{Tr}C^3`$, where $`I`$ are invariant functions built out of the field strengths, the matter fields and their covaraint derivatives. In this way, one is led to use the full extended antifield formalism as described in section II with all BRST cohomology classes in positive ghost number coupled from the beginning. In the case where the algebra of the non trivial symmetries associated to the currents $`j_\mathrm{\Delta }`$ is non abelian, the operator $`\mathrm{\Delta }_c`$ will be non vanishing at the classical level and involve in particular the structure constants of this algebra.
Another possibility is to try to show that the anomaly candidates $`V_\mathrm{\Delta }^{2,4}`$ do not effectively arise in the theory, by using higher order cohomological restrictions: as in the proof of the absence of similiar instabilities in the presence of abelian factors for standard Yang-Mills theories (see , Appendix A) one couples with external fields gauge invariant functions that break the symmetries associated to the currents $`j_\mathrm{\Delta }`$. In this way, one eliminates the currents $`j_\mathrm{\Delta }`$ and the associated anomaly for anomaly candidates $`V_\mathrm{\Delta }^{2,4}`$ from the extended theory. At the end of the computations, the external fields can be put to zero.
Because this discussion is not central to the argument below, we will simply assume here that the observables $`𝒪_i`$ are such that there are no non trivial currents $`j_\mathrm{\Delta }`$ and thus no anomaly for anomaly candidates $`V_\mathrm{\Delta }^{2,4}`$ in the theory. The general solution of the consistency condition in ghost number $`1`$ is then given by
$`{\displaystyle \frac{S_\rho }{\rho }}\sigma (g,g^i)+(S_\rho ,\mathrm{\Sigma }_\rho ),`$ (106)
for some local functional $`\mathrm{\Sigma }_\rho `$ in ghost number $`0`$.
By standard arguments, using in addition the same reasoning as in section II B 2, it follows from (105) and (106) that the model is renormalizable and that the renormalized generating functional for 1 particle irreducible vertex functions $`\mathrm{\Gamma }_\rho `$ satisfies
$`{\displaystyle \frac{1}{2}}(\mathrm{\Gamma }_\rho ,\mathrm{\Gamma }_\rho )+{\displaystyle \frac{^R\mathrm{\Gamma }_\rho }{\rho }}\mathrm{}\sigma (g,g^i)=0,`$ (107)
where $`\sigma (g,g^i)`$ is a formal power series in $`\mathrm{}`$. Hence, in this case $`\mathrm{\Delta }^{\mathrm{}}=\frac{^R}{\rho }\mathrm{}\sigma (g,g^i)`$ and the quantum BRST differential is $`s^q=(\mathrm{\Gamma }_\rho ,)\mathrm{}\sigma (g,g^i)^L\rho `$. (It can be shown along standard lines that the coefficient $`\sigma (g,g^i)`$ does not depend on the parameters of the gauge fixing.)
Let us now investigate the renormalization of the operators $`d^4xF_I^{\mu \nu }F_{\mu \nu }^I`$ and $`d^4x𝒪_i`$. According to the classification in section III, they are both of the type $`e_{\alpha _0}^0`$, because they are non trivial bottoms in maximal form degree $`4`$, while the Adler-Bardeen anomaly $`\mathrm{Tr}[Cd(AdA+\frac{1}{2}A^3)]`$ is of the type $`e_{\alpha _4}^4`$, as it descends to the non trivial bottom $`\mathrm{Tr}C^5`$. Because there is no $`e_{\alpha _0}^0`$ in ghost number $`1`$ (and form degree 4) and no $`h_{i_r}^r`$ in form degree $`3`$ and ghost number $`1`$, equation (57) implies that the lowest order contribution to the anomaly in the renormalization of $`d^4xF_I^{\mu \nu }F_{\mu \nu }^I`$ and of $`d^4x𝒪_i`$ is $`s_\rho `$ exact and can thus be absorbed through a BRST breaking counterterm added to these operators. This reasoning can be pushed to all orders, with the result that one can achieve, through the addition of suitable counterterms,
$`s^q[d^4xF_I^{\mu \nu }F_{\mu \nu }^I\mathrm{\Gamma }_\rho ]=0,`$ (108)
$`s^q[d^4x𝒪_i\mathrm{\Gamma }_\rho ]=0.`$ (109)
If we now apply $`\frac{g}{2}\frac{}{g}`$, respectively $`\frac{}{g^j}`$ to (107), we get on the one hand,
$`s^q[{\displaystyle d^4x}{\displaystyle \frac{1}{4g^2}}F_I^{\mu \nu }F_{\mu \nu }^I\mathrm{\Gamma }_\rho ]+{\displaystyle \frac{^R\mathrm{\Gamma }_\rho }{\rho }}\mathrm{}[{\displaystyle \frac{g}{2}}{\displaystyle \frac{\sigma (g,g^i)}{g}}]=0,`$ (110)
$`s^q[{\displaystyle d^4x𝒪_i\mathrm{\Gamma }_\rho }]+{\displaystyle \frac{^R\mathrm{\Gamma }_\rho }{\rho }}\mathrm{}[{\displaystyle \frac{\sigma (g,g^i)}{g^j}}]=0.`$ (111)
Comparing on the other hand with the integrated versions of (108) and (109), we deduce that
$`{\displaystyle \frac{\sigma (g,g^i)}{g}}=0,{\displaystyle \frac{\sigma (g,g^i)}{g^j}}=0,`$ (112)
which is an expression of the Adler-Bardeen theorem (see e.g. section 6.3.2 of ).
## VII Conclusion and outlook
In the completely extended antifield formalism, where in particular potential anomaly candidates are coupled to the starting point action with independent couplings, the whole program of algebraic renormalization can be extended to the case of theories with gauge anomalies. For instance, it is possible to write the Callan-Symanzik equation to all orders as a functional differential equation as in , by using $`\mathrm{\Delta }^{\mathrm{}}`$ instead of $`\mathrm{\Delta }_c`$. The $`\beta `$ functions can then be considered as anomaly coefficients in the renormalization of the cohomology class corresponding to the generator of dilatations. Vanishing results for $`\beta `$ functions associated to couplings of cohomological classes with a non trivial descent can then be obtained as a particular case of the discussion of section IV.
The discussion in section IV on the length of descents and lifts of BRST cohomological classes and their anomalies does not rely on the use of the extended antifield formalism. It can be done along the same lines in the standard set-up as long as one assumes the quantum theory to be anomaly free and stable, so that the standard Zinn-Justin equation $`\frac{1}{2}(\mathrm{\Gamma },\mathrm{\Gamma })=0`$ holds. A virtue of the completely extended formalism is that it allows in principle to treat the general case.
###### Acknowledgements.
This work has been partly supported by the “Actions de Recherche Concertées” of the “Direction de la Recherche Scientifique - Communauté Française de Belgique”, by IISN - Belgium (convention 4.4505.86). The author wants to thank F. Brandt and M. Henneaux for useful discussions.
|
warning/0003/nlin0003038.html
|
ar5iv
|
text
|
# The Burnett expansion of the periodic Lorentz gas
## 1 Introduction
The Lorentz gas is a model used in statistical mechanics, consisting of a point particle moving at constant velocity except for specular collisions with smooth (specifically $`C^3`$) convex fixed scatterers in $`d2`$ dimensions. The original model has randomly placed scatterers in infinite space, and is thought to have power law decay of correlations, so that the Burnett coefficients (defined below as sums of such correlations) are not generally expected to exist . Here we consider a periodic arrangement of scatterers which is equivalent to a dispersing billiard on a torus, for which it is known that two-time correlations of the discrete (collision) dynamics decay exponentially . This, together with the finite horizon condition, that is, that the time between collisions is bounded, implies the existence of the diffusion coefficient ($`D^{(2)}`$ below).
A recent paper gives a stretched exponential decay of multiple correlations , and uses this to show (again with finite horizon) that the fourth order Burnett coefficient ($`D^{(4)}`$ below) exists. Here we extend this result to all the Burnett coefficients. A common example for $`d=2`$ with a finite horizon is given by circular scatterers on a triangular lattice; for $`d>2`$ the finite horizon condition requires either nonspherical scatterers, or more than one scatterer per unit cell. The Lorentz gas and a number of extensions are discussed in Ref. .
The Burnett coefficients $`D^{(m)}`$ discussed in this paper are defined using series of correlation functions. Section 2 defines these series and gives three basic results about them. Section 3 gives the main result of this paper, the proof of convergence of these series. The series arise in a physical description of diffusion, however the derivation involves hydrodynamic approximations and interchange of limits which have not been justified rigorously for the Lorentz gas. The physical motivation together with the non-rigorous derivation of the series given below from previously stated formulas is given in the final section, together with a conjecture about the Burnett expansion.
The author is grateful for helpful discussion with N. I. Chernov, E. G. D. Cohen, J. R. Dorfman and P. Gaspard, and for the support of the Engineering Research Program of the Office of Basic Energy Sciences at the US Department of Energy, contract #DE-FG02-88-ER13847, and the Nuffield Foundation, grant NAL/00353/G.
## 2 Definitions
In the following, $`\mathit{\varphi }(x)`$ is the billiard map defined on the collision space $`M`$, consisting of points $`x=(𝐫,𝐯)M`$ for which the position $`𝐫\text{I}\text{R}^d`$ is on the boundary of one of the scatterers and the velocity following a collision $`𝐯\text{I}\text{R}^d`$ is of unit magnitude in an outward direction from the scatterer. Greek indices $`\alpha ,\beta ,\mathrm{}=1,\mathrm{},d`$ denote components of vectors and tensors in $`\text{I}\text{R}^d`$, and a dot $`𝐚𝐛`$ denotes the usual inner product $`_\alpha a_\alpha b_\alpha `$ corresponding to the Euclidean metric. We have two functions $`T:M\text{I}\text{R}`$ and $`𝐚:M\text{I}\text{R}^d`$ which describe the embedding of the collision dynamics into physical space and time, as follows. $`T(x)`$ is the time (also distance since the speed of the particle is one) between the collision at $`x`$ and the next; it is a piecewise Hölder continuous function . $`𝐚(x)`$ is the lattice translation vector associated with this free flight when the configuration variable $`𝐫`$ is unfolded onto a periodic tiling of $`\text{I}\text{R}^d`$; it is a linear combination of the lattice basis vectors $`𝐞^{(\alpha )}`$ with integer coefficients, and is a piecewise constant function. The finite horizon condition ensures that both $`T`$ and $`𝐚`$ are bounded. The average $``$ denotes integration over $`M`$ with respect to the invariant equilibrium measure. In terms of this average we define $`\mathrm{\Delta }T:M\text{I}\text{R}`$ by $`\mathrm{\Delta }T(x)=T(x)T`$ so that $`\mathrm{\Delta }T=0`$.
The billiard dynamics is time reversal invariant, that is, there exists an involution $`𝒯:MM`$ (given simply by the specular reflection law) with the property
$$\mathit{\varphi }𝒯\mathit{\varphi }=𝒯$$
(1)
In addition, $`𝒯`$ preserves the equilibrium measure, that is,
$$𝒯g=g$$
(2)
for arbitrary measurable function $`g:MM`$. The map $`𝒯`$ also satisfies
$`T𝒯\mathit{\varphi }`$ $`=`$ $`T`$ (3)
$`𝐚𝒯\mathit{\varphi }`$ $`=`$ $`𝐚`$ (4)
Thus $`𝐚=0`$.
The wave vector $`𝐤`$ is to be understood as a formal real expansion parameter with $`d`$ components (although physically we would like to interpret it as a vector with a value in $`\text{I}\text{R}^d`$). The dispersion relation $`s[𝐤]`$ is to be understood as a formal power series
$$s[𝐤]=\underset{m=2}{\overset{\mathrm{}}{}}\mathrm{i}^m\underset{\alpha _1\mathrm{}\alpha _m}{}D_{\alpha _1\mathrm{}\alpha _m}^{(m)}k_{\alpha _1}\mathrm{}k_{\alpha _m}$$
(5)
in terms of the Burnett coefficients $`D^{(m)}`$ which are assumed to be real, totally symmetric tensors of rank $`m`$. That is, an equation (specifically Eq. (15) below) involving $`s[𝐤]`$ is to be interpreted as a sequence of equations (specifically Eq. (16) below) obtained by equating coefficients of powers of $`𝐤`$. The symbol $`\mathrm{i}`$ denotes $`\sqrt{1}`$.
The existence of Burnett coefficients satisfying the equations (16) is not assumed a priori; we show in Lemma 1 below that equations (16) express the $`d(d+1)/2`$ independent components of $`D^{(2)}`$ as series not containing any of the $`D^{(m)}`$, then the $`d(d+1)(d+2)/6`$ independent components of $`D^{(3)}`$ as series containing only the $`D^{(2)}`$ and so on. Lemma 2 shows that they are indeed real, and Thm. 4 shows that the limit exists.
We define formal power series $`f`$ and $`F`$ by
$$f[𝐤]s[𝐤]\mathrm{\Delta }T+\mathrm{i}𝐤𝐚$$
(6)
$$F[𝐤]\underset{i=n}{\overset{n1}{}}f[𝐤]\mathit{\varphi }^i$$
(7)
where the dependence on $`x`$ and on the positive integer $`n`$ is suppressed in the notation; the limit $`n\mathrm{}`$ will be taken later. We have $`f=0`$ and $`F=0`$ at each order in $`𝐤`$ and for each $`n`$ as a consequence of $`\mathrm{\Delta }T=0`$ and $`𝐚=0`$ above.
We define cumulants $`Q_N[𝐤]`$ (also formal power series) for integers $`N2`$ as
$$Q_N[𝐤]=\underset{\{\nu _j\}:_jj\nu _j=N}{}(1)^{\nu 1}\frac{(\nu 1)!_jF[𝐤]^j^{\nu _j}}{_j(\nu _j!j!^{\nu _j})}$$
(8)
with $`j`$ and $`\nu _j`$ integers satisfying $`j2`$ and $`\nu _j0`$, and $`\nu =_j\nu _j`$ is the total number of correlations in the product. For example
$`Q_2`$ $`=`$ $`F^2/2`$ (9)
$`Q_3`$ $`=`$ $`F^3/6`$ (10)
$`Q_4`$ $`=`$ $`(F^43F^2^2)/24`$ (11)
$`Q_5`$ $`=`$ $`(F^510F^3F^2)/120`$ (12)
$`Q_6`$ $`=`$ $`(F^615F^4F^210F^3^2+30F^2^3)/720`$ (13)
Now $`Q_N`$ contains exactly $`N`$ powers of $`F`$, and so it contains terms $`𝐤^m`$ only for $`mN`$, and we can write it as
$$Q_N[𝐤]=\underset{m=N}{\overset{\mathrm{}}{}}\underset{\alpha _1\mathrm{}\alpha _m}{}q_{N,m;\alpha _1\mathrm{}\alpha _m}k_{\alpha _1}\mathrm{}k_{\alpha _m}$$
(14)
thus defining totally symmetric tensors $`q_{N,m}`$ for $`mN`$.
The Burnett coefficients are found by equating the formal power series on both sides of
$$s[𝐤]=\underset{n\mathrm{}}{lim}\frac{1}{2nT}\underset{N=2}{\overset{\mathrm{}}{}}Q_N[𝐤]$$
(15)
that is,
$$\mathrm{i}^mD_{\alpha _1\mathrm{}\alpha _m}^{(m)}=\underset{n\mathrm{}}{lim}\frac{1}{2nT}\underset{N=2}{\overset{m}{}}q_{N,m;\alpha _1\mathrm{}\alpha _m}$$
(16)
These equations determine the $`D^{(m)}`$ explicitly as real tensors, subject to convergence of the limit, as shown by the following two lemmas.
###### Lemma 1
The right hand side of Eq. (16) does not contain $`D^{(m^{})}`$ such that $`m^{}m`$.
Proof: We have $`N2`$, and each $`Q_N`$ contains $`N`$ powers of $`F`$, thus each term has at least 2 powers of $`F`$. Each $`F`$ has at least 1 power of $`𝐤`$, and there are $`m`$ powers of $`𝐤`$ in total, so each $`F`$ has at most $`m1`$ powers of $`𝐤`$. $`D^{(m^{})}`$ appear in $`F`$ associated with $`m^{}`$ powers of $`𝐤`$, so $`m^{}m1`$ for any $`D^{(m^{})}`$ appearing.
Remark: It is possible there are no factors of $`D^{(m^{})}`$ on the right hand side, in fact the lemma shows that this is true for $`m=2`$. The case $`m=2`$ can easily be written explicitly; Eq. (16) becomes
$$D_{\alpha \beta }^{(2)}=\underset{n\mathrm{}}{lim}\frac{1}{4nT}\underset{i=n}{\overset{n1}{}}\underset{j=n}{\overset{n1}{}}a_\alpha ^ia_\beta ^j$$
(17)
This is a discrete time version of the well-known Green-Kubo formula for the diffusion tensor (which reduces to a single diffusion coefficient in the isotropic case $`D_{\alpha \beta }^{(2)}=D\delta _{\alpha \beta }`$). An equivalent discrete time equation appears in , also for $`m=4`$.
###### Lemma 2
Despite the appearance of the imaginary number $`\mathrm{i}`$ in the above definitions, the Burnett coefficients are real if they exist.
Proof: We note from the definitions that $`s[𝐤]`$, $`f[𝐤]`$ and $`F[𝐤]`$ have pure imaginary coefficients for odd powers of $`𝐤`$ and real coefficients for even powers of $`𝐤`$. This property is preserved by addition and multiplication of power series, so it also holds for the $`Q_N[𝐤]`$. This implies that the $`q_{N,m}`$ are imaginary for odd $`m`$ and real for even $`m`$. The result follows from Eq. (16).
Before proceeding with the more technical convergence proof, we note another important result:
###### Lemma 3
$`D^{(m)}=0`$ for $`m`$ odd.
Proof: From the properties of the time reversal operator $`𝒯`$ given above, $`F^j`$ has zero contribution from any term with an odd number of $`𝐚`$ factors. The result follows by induction on $`m`$: assume that $`D^{(m^{})}=0`$ for all odd $`m^{}<m`$, then by Lemma 1 all terms in $`s[𝐤]`$ contributing to $`D^{(m)}`$ have even powers of $`𝐤`$, and from the oddness of $`𝐚`$ under time reversal, so also do the $`\mathrm{i}𝐤𝐚`$ terms. Thus $`D^{(m)}`$, which is constructed from terms with $`m`$ powers of $`𝐤`$, must be zero for $`m`$ odd.
## 3 Convergence of the series
The averages $`F^j`$ appearing in the cumulants contain summations over $`j`$ variables with range $`n`$ to $`n1`$, and could grow as fast as $`O(n^j)`$ in general. Thus each term, which is a product of such averages could grow as $`O(n^N)`$ in general. For the limit in Eq. (16) to exist, we require that the series grows only as $`O(n)`$. Although the growth of each product of correlations cannot be controlled this well, cancellations occur in constructing the cumulants. This is expressed in the following theorem which, together with Lemmas 1 and 2, implies the existence of the Burnett coefficients:
###### Theorem 4
$`q_{N,m}`$ is defined in Eqs. (7, 8, 14) for integers $`N`$ and $`m`$ satisfying $`2Nm`$. The limit
$$\underset{n\mathrm{}}{lim}\frac{1}{n}q_{N,m;\alpha _1\mathrm{}\alpha _m}$$
(18)
exists for all such $`N`$ and $`m`$ in the periodic Lorentz gas.
The structure of the proof of Thm. 4 is as follows. We state the theorem expressing stretched exponential decay of multiple correlation functions. Next, the terms appearing in (16) are written as a time ordered sum, so that this theorem can be applied. Then we show that all the terms connected by the application of the theorem have coefficients which sum to zero, so that only the stretched exponential corrections remain. Finally, a bound of $`n`$ multiplied by a polynomial is put on the number of terms at each order of the stretched exponential, so that the series divided by $`n`$ converges absolutely.
Thm. 4 is based on the following result:
###### Theorem 5
(Theorem 2 of Ref. ) Let $`i_1\mathrm{}i_k`$ and $`1tk1`$. Then
$$|f_1^{i_1}\mathrm{}f_k^{i_k}f_1^{i_1}\mathrm{}f_t^{i_t}f_{t+1}^{i_{t+1}}\mathrm{}f_k^{i_k}|C_k|i_ki_1|^2\lambda ^{|i_{t+1}i_t|^{1/2}}$$
(19)
where $`C_k>0`$ depends on the functions $`f_1,\mathrm{},f_k`$, and $`\lambda <1`$ is independent of $`k`$ and $`f_1,\mathrm{},f_k`$.
The theorem applies to piecewise Hölder continuous functions $`f_j`$ such that $`f_j=0`$ for all $`j`$ and uses notation $`f_j^if_j\mathit{\varphi }^i`$. As noted in Ref. , we expect based on Refs. that it should be possible to prove a stronger bound $`\lambda ^{|i_{t+1}i_t|}`$, but the above bound is sufficient for our purposes here.
The $`q_{N,m}`$ as defined in the previous section are finite sums of terms of the form (see Eqs. (7, 8, 14))
$$\underset{\{\nu _j\}:_jj\nu _j=N}{}(1)^{\nu 1}\frac{(\nu 1)!}{_j(\nu _j!j!^{\nu _j})}\underset{i_1\mathrm{}i_N=n}{\overset{n1}{}}f_1^{i_1}\mathrm{}f_j^{i_j}f_{j+1}^{i_{j+1}}\mathrm{}\mathrm{}\mathrm{}f_N^{i_N}$$
(20)
multiplied by constants such as the lower order Burnett coefficients. The $`f`$ here and for the remainder of this section are $`T`$ or $`𝐚`$, both of which satisfy the conditions of Thm. 5. The exact number of terms of this kind is not important; it depends on $`N`$ and $`m`$ but not $`n`$ and therefore does not affect convergence of the limit $`n\mathrm{}`$.
In order to use Thm. 5 we need to put the times $`i_p`$ in numerical order. The unrestricted sum over all the $`i_p`$ is replaced by an ordered sum $`i_1i_2\mathrm{}i_N`$ over all $`N!/S[i]`$ permutations of the $`i_p`$. $`S[i]`$ is a symmetry factor to account for the fact that some of the $`i_p`$ may be equal; the exact form is unimportant since it is a common prefactor, independent of the $`\nu _j`$. Not all $`N!`$ permutations of the correlations are distinct: it does not matter in which order the $`f_j`$ are multiplied within a correlation, or which order correlations of equal numbers of $`f_j`$ are multiplied; thus both factorials in the denominator disappear, leading to
$`{\displaystyle \underset{\{\nu _j\}:_jj\nu _j=N}{}}`$ $`(1)^{\nu 1}(\nu 1)![{\displaystyle \underset{i_1i_2\mathrm{}i_N}{}}{\displaystyle \frac{1}{S[i]}}\{f_1^{i_1}\mathrm{}f_j^{i_j}f_{j+1}^{i_{j+1}}\mathrm{}\mathrm{}\mathrm{}f_N^{i_N}`$ (21)
$`+\text{ permutations}\}]`$
The “permutations” remaining in (21) consist of the remaining $`N!/(_j\nu _j!j!^{\nu _j})1`$ rearrangements of the $`i_p`$ that are not equivalent by reordering the product of correlations or the product of $`f`$ within a correlation.
As an example, we give the expression for $`N=6`$:
$`{\displaystyle \underset{i_1i_2i_3i_4i_5i_6}{}}{\displaystyle \frac{1}{S[i]}}\{f_1^{i_1}f_2^{i_2}f_3^{i_3}f_4^{i_4}f_5^{i_5}f_6^{i_6}`$
$`[f_1^{i_1}f_2^{i_2}f_3^{i_3}f_4^{i_4}f_5^{i_5}f_6^{i_6}+f_1^{i_1}f_3^{i_3}f_2^{i_2}f_4^{i_4}f_5^{i_5}f_6^{i_6}+f_1^{i_1}f_4^{i_4}f_2^{i_2}f_3^{i_3}f_5^{i_5}f_6^{i_6}`$
$`+f_1^{i_1}f_5^{i_5}f_2^{i_2}f_3^{i_3}f_4^{i_4}f_6^{i_6}+f_1^{i_1}f_6^{i_6}f_2^{i_2}f_3^{i_3}f_4^{i_4}f_5^{i_5}+f_2^{i_2}f_3^{i_3}f_1^{i_1}f_4^{i_4}f_5^{i_5}f_6^{i_6}`$
$`+f_2^{i_2}f_4^{i_4}f_1^{i_1}f_3^{i_3}f_5^{i_5}f_6^{i_6}+f_2^{i_2}f_5^{i_5}f_1^{i_1}f_3^{i_3}f_4^{i_4}f_6^{i_6}+f_2^{i_2}f_6^{i_6}f_1^{i_1}f_3^{i_3}f_4^{i_4}f_5^{i_5}`$
$`+f_3^{i_3}f_4^{i_4}f_1^{i_1}f_2^{i_2}f_5^{i_5}f_6^{i_6}+f_3^{i_3}f_5^{i_5}f_1^{i_1}f_2^{i_2}f_4^{i_4}f_6^{i_6}+f_3^{i_3}f_6^{i_6}f_1^{i_1}f_2^{i_2}f_4^{i_4}f_5^{i_5}`$
$`+f_4^{i_4}f_5^{i_5}f_1^{i_1}f_2^{i_2}f_3^{i_3}f_6^{i_6}+f_4^{i_4}f_6^{i_6}f_1^{i_1}f_2^{i_2}f_3^{i_3}f_5^{i_5}+f_5^{i_5}f_6^{i_6}f_1^{i_1}f_2^{i_2}f_3^{i_3}f_4^{i_4}]`$
$`[f_1^{i_1}f_2^{i_2}f_3^{i_3}f_4^{i_4}f_5^{i_5}f_6^{i_6}+f_1^{i_1}f_2^{i_2}f_4^{i_4}f_3^{i_3}f_5^{i_5}f_6^{i_6}+f_1^{i_1}f_2^{i_2}f_5^{i_5}f_3^{i_3}f_4^{i_4}f_6^{i_6}`$
$`+f_1^{i_1}f_2^{i_2}f_6^{i_6}f_3^{i_3}f_4^{i_4}f_5^{i_5}+f_1^{i_1}f_3^{i_3}f_4^{i_4}f_2^{i_2}f_5^{i_5}f_6^{i_6}+f_1^{i_1}f_3^{i_3}f_5^{i_5}f_2^{i_2}f_4^{i_4}f_6^{i_6}`$
$`+f_1^{i_1}f_3^{i_3}f_6^{i_6}f_2^{i_2}f_4^{i_4}f_5^{i_5}+f_1^{i_1}f_4^{i_4}f_5^{i_5}f_2^{i_2}f_3^{i_3}f_6^{i_6}+f_1^{i_1}f_4^{i_4}f_6^{i_6}f_2^{i_2}f_3^{i_3}f_5^{i_5}`$
$`+f_1^{i_1}f_5^{i_5}f_6^{i_6}f_2^{i_2}f_3^{i_3}f_4^{i_4}]`$ (22)
$`+2[f_1^{i_1}f_2^{i_2}f_3^{i_3}f_4^{i_4}f_5^{i_5}f_6^{i_6}+f_1^{i_1}f_2^{i_2}f_3^{i_3}f_5^{i_5}f_4^{i_4}f_6^{i_6}+f_1^{i_1}f_2^{i_2}f_3^{i_3}f_6^{i_6}f_4^{i_4}f_5^{i_5}`$
$`+f_1^{i_1}f_3^{i_3}f_2^{i_2}f_4^{i_4}f_5^{i_5}f_6^{i_6}+f_1^{i_1}f_3^{i_3}f_2^{i_2}f_5^{i_5}f_4^{i_4}f_6^{i_6}+f_1^{i_1}f_3^{i_3}f_2^{i_2}f_6^{i_6}f_4^{i_4}f_5^{i_5}`$
$`+f_1^{i_1}f_4^{i_4}f_2^{i_2}f_3^{i_3}f_5^{i_5}f_6^{i_6}+f_1^{i_1}f_4^{i_4}f_2^{i_2}f_5^{i_5}f_3^{i_3}f_6^{i_6}+f_1^{i_1}f_4^{i_4}f_2^{i_2}f_6^{i_6}f_3^{i_3}f_5^{i_5}`$
$`+f_1^{i_1}f_5^{i_5}f_2^{i_2}f_3^{i_3}f_4^{i_4}f_6^{i_6}+f_1^{i_1}f_5^{i_5}f_2^{i_2}f_4^{i_4}f_3^{i_3}f_6^{i_6}+f_1^{i_1}f_5^{i_5}f_2^{i_2}f_6^{i_6}f_3^{i_3}f_4^{i_4}`$
$`+f_1^{i_1}f_6^{i_6}f_2^{i_2}f_3^{i_3}f_4^{i_4}f_5^{i_5}+f_1^{i_1}f_6^{i_6}f_2^{i_2}f_4^{i_4}f_3^{i_3}f_5^{i_5}+f_1^{i_1}f_6^{i_6}f_2^{i_2}f_5^{i_5}f_3^{i_3}f_4^{i_4}]\}`$
Here, the four terms correspond to the partitions of $`6`$ which do not contain $`1`$; in the above notation the nonzero $`\nu _j`$ are $`\{\nu _6=1\}`$ with $`6!/6!=1`$ term; $`\{\nu _2=1,\nu _4=1\}`$ with $`6!/2!4!=15`$ terms; $`\{\nu _3=2\}`$ with $`6!/2!3!^2=10`$ terms; and $`\{\nu _2=3\}`$ with $`6!/3!2!^3=15`$ terms; compare with Eq. (13).
Now we apply Thm. 5 to the largest gap, $`i_{t+1}i_t`$. Any of the largest gaps will suffice if more than one is largest. Before tackling the general case, we see how it works in the $`N=6`$ example. Notice that, whatever the value of $`t`$, the theorem combines all the above correlations to leave terms (the number of which is a function of $`N`$) bounded by $`\lambda ^{|i_{t+1}i_t|^{1/2}}`$ multiplied by powers of the time differences. Explicitly, for $`t=1`$, all terms cancel individually because $`f_j=0`$. For $`t=2`$ the $`f^6`$ term cancels with one of the $`f^2f^4`$ terms, six other $`f^2f^4`$ terms cancel with three of the $`f^2^3`$ terms and the remaining terms all split leaving an $`f`$ term. For $`t=3`$ the $`f^6`$ term cancels with one of the $`f^3^2`$ terms, and all of the others split leaving an $`f`$ term. $`t=4`$ is analogous to $`t=2`$ and $`t=5`$ is analogous to $`t=1`$.
In general we must show that the coefficient $`(1)^{\nu 1}(\nu 1)!`$ in Eq. (21) combined with the numbers of terms of various types leads to complete cancellation for all values of $`N`$. Consider a general term (ignoring the $`S[i]`$ which is the same for each term) which is unaffected by a split at time $`t`$. Each correlation contains times $`i_pt`$ or times $`i_p>t`$ but not both. Thus it can be written schematically as
$$\mathrm{}|\mathrm{}$$
(23)
where all times $`i_p`$ to the left of the bar “$`|`$” are less than or equal to $`t`$ and all times to the right of the bar are greater than $`t`$. Let there be $`A`$ correlations to the left and $`B`$ correlations to the right, so $`A+B=\nu `$.
This term will cancel (up to stretched exponential corrections) with any term which is split to the same form, if the sum of the coefficients (the $`(1)^{\nu 1}(\nu 1)!`$) is zero. The terms that are split to a given form consist of correlations that are either the same as the above, or are joined in a pairwise fashion with a correlation on the other side of the bar.
Again, an example is helpful: When $`N=8`$, a split at $`t=4`$ combines the following terms: $`6f_1^{i_1}f_2^{i_2}f_3^{i_3}f_4^{i_4}|f_5^{i_5}f_6^{i_6}f_7^{i_7}f_8^{i_8}`$ with $`2f_1^{i_1}f_2^{i_2}f_5^{i_5}f_6^{i_6}f_3^{i_3}f_4^{i_4}f_7^{i_7}f_8^{i_8}`$,
$`2f_1^{i_1}f_2^{i_2}f_7^{i_7}f_8^{i_8}f_3^{i_3}f_4^{i_4}f_5^{i_5}f_6^{i_6}`$, $`2f_1^{i_1}f_2^{i_2}f_3^{i_3}f_4^{i_4}f_5^{i_5}f_6^{i_6}f_7^{i_7}f_8^{i_8}`$, $`2f_1^{i_1}f_2^{i_2}f_3^{i_3}f_4^{i_4}f_7^{i_7}f_8^{i_8}f_5^{i_5}f_6^{i_6}`$, $`f_1^{i_1}f_2^{i_2}f_5^{i_5}f_6^{i_6}f_3^{i_3}f_4^{i_4}f_7^{i_7}f_8^{i_8}`$ and $`f_1^{i_1}f_2^{i_2}f_7^{i_7}f_8^{i_8}f_3^{i_3}f_4^{i_4}f_5^{i_5}f_6^{i_6}`$. These all cancel because $`6+2+2+2+211=0`$.
The term given in Eq. (23) has coefficient $`(1)^{\nu 1}(\nu 1)!`$. There are $`AB`$ terms with coefficient $`(1)^{\nu 2}(\nu 2)!`$ obtained by combining a single correlation on the left and the right. There are $`A(A1)B(B1)/2!`$ terms with coefficient $`(1)^{\nu 3}(\nu 3)!`$ obtained by combining two correlations on the left and the right, and so on until all $`\mathrm{min}(A,B)`$ correlations on the side with the fewest correlations have been combined. The total coefficient is thus given by
$$H(A,B)\underset{p=0}{\overset{\mathrm{min}(A,B)}{}}(1)^{A+Bp1}(A+Bp1)!\frac{A!B!}{(Ap)!(Bp)!p!}$$
(24)
To show that the coefficients cancel, we therefore need the following lemma:
###### Lemma 6
$`H(A,B)=0`$ for all positive integers $`A`$ and $`B`$.
Proof: The sum is symmetric in $`A`$ and $`B`$ so suppose that $`AB`$ without loss of generality. Then the summand is the product of a constant $`(1)^{A+B1}A!`$, an alternating binomial of degree $`B`$, that is, $`(1)^pB!/((Bp)!p!)`$ and a polynomial in $`p`$ of degree $`B1`$, that is, $`(A+Bp1)!/(Ap)!`$. We will use summation by parts to lower the degree of both until the result is zero.
We note the summation by parts formula
$$\underset{p=0}{\overset{B}{}}x_py_p=y_0\underset{p=0}{\overset{B}{}}x_p+\underset{q=1}{\overset{B}{}}(y_qy_{q1})\underset{p=q}{\overset{B}{}}x_p$$
(25)
which can be demonstrated by collecting terms on the right hand side. Now substituting $`x_p=(1)^pB!/((Bp)!p!`$ and $`y_p=(A+Bp1)!/(Ap)!`$ we can show by induction on $`q`$ from $`B`$ downwards that
$$\underset{p=q}{\overset{B}{}}x_p=\{\begin{array}{cc}(1)^q\frac{(B1)!}{(Bq)!(q1)!}& q>0\\ 0& q=0\end{array}$$
(26)
hence the first term on the right hand side of Eq. (25) vanishes. We can also simplify
$$y_qy_{q1}=(1B)\frac{(A+Bq1)!}{(Aq+1)!}$$
(27)
so Eq. (24) now reads
$$H(A,B)=(1)^{A+B1}A!(1B)\underset{q=1}{\overset{B}{}}\frac{(1)^q(B1)!}{(Bp)!(p+1)!}\frac{(A+Bq1)!}{(Aq+1)!}$$
(28)
Shifting the summation index by one we find
$$H(A,B)=(1B)H(A,B1)$$
(29)
The proof of Lemma 6 follows by noting that $`H(A,1)=0`$.
We now conclude the proof of Thm. 4. Recall that the series (20) have been rewritten in the form (21). Thm. 5 is applied to (one of) the largest gap(s) $`\mathrm{\Delta }i_{\mathrm{max}}i_{t+1}i_t`$, partitioning the terms into subsets which split into a particular form (23). Lemma 6 shows that the coefficients of all terms in a subset conspire to cancel, so that each subset is bounded by the error term in theorem 5, that is $`\lambda ^{|\mathrm{\Delta }i_{\mathrm{max}}|^{1/2}}`$ multiplied by a polynomial in the time differences.
Finally we estimate the number of terms with each value of $`\mathrm{\Delta }i_{\mathrm{max}}`$. The first time $`i_1`$ varies freely from $`n`$ to $`n1`$, having a total of $`2n`$ values. One of the time differences is equal to $`\mathrm{\Delta }i_{\mathrm{max}}`$, and the other $`k2`$ time differences can range from $`0`$ to $`\mathrm{\Delta }i_{\mathrm{max}}`$, so the total number of terms with a given $`\mathrm{\Delta }i_{\mathrm{max}}`$ is less than $`2n(k1)\mathrm{\Delta }i_{\mathrm{max}}^{(k2)}`$, in particular a polynomial in $`\mathrm{\Delta }i_{\mathrm{max}}`$ multiplied by $`n`$. Thus the series divided by $`n`$ appearing in Thm. 4 is bounded by a product of polynomial factors and the decaying stretched exponential, and hence converges absolutely. This concludes the proof of Thm. 4 and the proof of existence of Burnett coefficients.
## 4 Physical motivation and remarks
This section makes the connection between the Burnett coefficients defined in the previous sections and equations found in the physics literature. The latter equations are phenomenological and have not been shown rigorously in a limiting fashion from the Lorentz gas, and a few nonrigorous limit interchanges are made to connect them with the expressions defined in the previous sections. First we consider the dispersion relation, the equation for the Burnett coefficients, and finally the whether the dispersion relation can be used to define an analytic function.
The dispersion relation (5) with $`𝐤`$ interpreted as a real vector represents the solution of a generalised diffusion equation proposed by Burnett containing higher derivative terms that become important on small scales,
$$_t\rho =\underset{m=2}{\overset{\mathrm{}}{}}\underset{\alpha _1\mathrm{}\alpha _m}{}D_{\alpha _1\mathrm{}\alpha _m}^{(m)}_{\alpha _1}\mathrm{}_{\alpha _m}\rho $$
(30)
assuming a solution of the form
$$\rho (𝐫,t)\mathrm{exp}(s(𝐤)t+\mathrm{i}𝐤𝐫)$$
(31)
Here, $`_\alpha /r_\alpha `$. Nonlinear terms such as powers of $`_\alpha \rho `$ are excluded on physical grounds since $`\rho `$ is a projection onto real space ($`𝐫\text{I}\text{R}^d`$) of a phase space density satisfying a linear evolution equation. The phase space is a subset of $`\text{I}\text{R}^{2dM}`$ corresponding to the possible positions and velocities of $`M1`$ particles. The dispersion relation is a more robust formulation than the generalized diffusion equation (30) since the former may be supplemented by nonanalytic functions of $`𝐤`$ to account for situations (other than the periodic Lorentz gas) in which some of the Burnett coefficients do not exist.
Chapter 7 of Ref. obtains the dispersion relation from the microscopic dynamics using the equation (7.91 in this reference):
$$1=\underset{n\mathrm{}}{lim}\underset{i=n}{\overset{n1}{}}\mathrm{exp}[s(𝐤)T(\mathit{\varphi }^ix)\mathrm{i}𝐤𝐚(\mathit{\varphi }^ix)]$$
(32)
We write $`T=T+\mathrm{\Delta }T`$ as in previous sections, take out the constant factor of $`T`$, and take the logarithm to find
$$s(𝐤)=\underset{n\mathrm{}}{lim}\frac{1}{2nT}\mathrm{ln}\mathrm{exp}[F(𝐤)]$$
(33)
where $`F`$ is defined (as a power series) in Eq. (7). Now the exponential and the logarithm are expanded in power series and the resulting terms containing $`N`$ powers of $`F`$ are collected to become the cumulants $`Q_N`$ defined in Eq. (8). The cumulant form of the expansion is possibly more robust than the above equations due to the cancellations among the terms that combine to construct each cumulant.
Since it is desirable from a physical point of view to interpret $`𝐤`$ as a real variable, we conclude with the following conjecture:
###### Conjecture 7
The series (5) converges when $`𝐤𝒟\text{I}\text{R}^d`$ for some nontrivial domain $`𝒟`$, and so defines a function $`s(k)`$ in this domain.
Note that $`s(k)`$ (if it exists) is a real function as a consequence of Lemma 3, and physically is expected to be negative except at the origin (otherwise the density $`\rho `$ would grow exponentially with time); this puts further constraints on the Burnett coefficients.
Unfortunately the proof given in the previous section contains many undetermined functions of $`𝐤`$, and the Burnett coefficients are defined by a complicated recursive relation (16), so a proof is unlikely using the techniques of this paper.
There are two results that make such a result plausible. The first is that in the Boltzmann limit of a hard sphere gas, that is, a gas with many moving particles at low density and with recollisions ignored, the expansion in $`𝐤`$ (in this context called the linearized Chapman-Enskog expansion) converges . Of course, the hard sphere collisions are similar to that of the Lorentz gas, but recollisions cannot be ignored in general.
The second result is exact, but for a highly simplified (piecewise linear) system. We consider the map $`\mathit{\varphi }:\text{I}\text{R}\text{I}\text{R}`$ given by
$$\mathit{\varphi }(x)=\frac{3}{2}2x+3[x]$$
(34)
where $`[x]`$ is the greatest integer less than or equal to $`x`$. The dynamics defined by $`\mathit{\varphi }`$ is equivalent to a random walk where the particle moves with equal probability from one interval $`I_n(n1/2,n+1/2)`$ to the left, $`I_{n1}`$ or to the right, $`I_{n+1}`$. The dispersion relation $`s(k)`$ follows directly from the above phenomenological solution (31),
$$\rho (n,t)=\mathrm{exp}(st+ikn)$$
(35)
After one iteration,
$`\rho (n,1)`$ $`=`$ $`{\displaystyle \frac{1}{2}}[\mathrm{exp}(\mathrm{i}k(n1))+\mathrm{exp}(\mathrm{i}k(n+1))]`$ (36)
$`=`$ $`\mathrm{cos}k\mathrm{exp}(\mathrm{i}kn)`$ (37)
leading to
$$s(k)=\mathrm{ln}\mathrm{cos}k$$
(38)
which has a power series around $`k=0`$ with a radius of convergence equal to $`\pi /2`$.
|
warning/0003/math-ph0003004.html
|
ar5iv
|
text
|
# Periodic and almost periodic potentials in the inverse problems
## 1 Introduction
To begin with we are going to consider the inverse spectral problem for one-dimensional Schrödinger operator with periodic potential. In the late 60-s the famous discovery of the Inverse Scattering Transform for the KdV equation was done. A periodic analog of this transform was found in 1974. It is based on the solution of the following inverse spectral problem:
to describe effectively the ”isospectral manifold” of all the potentials with a given spectrum on the line (i.e. the spectrum of Schrödinger operator acting in the Hilbert space of square integrable complex-valued functions on the line $`R`$).
As everybody working in the quantum solid state physics knows, this spectrum generically is a union of infinite number of intervals (allowed bands) on the energy line $`ϵ`$. The complementary part on the energy line is also a union of infinite number of intervals (gaps or forbidden bands) whose lengths tend to zero for $`ϵ+\mathrm{}`$.
The periodic problem was solved in 1974-75 for the so-called finite-gap potentials that have only a finite number of gaps. Any periodic potential can be approximated by the finite-gap ones. This solution involves combination of the theory of Riemann surfaces and their $`\theta `$-functions, Hamiltonian dynamics of special completely integrable systems and the spectral theory of Schrödinger operator. The mathematical technique used was (and still is) unusual for the community of physicists. Later the necessity to use this kind of mathematics appeared also in other branches of mathematical and theoretical physics (for example, in the string theory, matrix models, and supersymmetric Yang-Mills theory ). It seems to the authors this technique in future will be needed to the broad community of theoretical physicists.
Integrability of the famous KdV equation $`u_t=6uu_x+u_{xxx}`$ was discovered in 1965-68 (see ) for the rapidly decreasing initial data on the line $`x`$. Exact solutions for the KdV equation expressing $`u(x,t)`$ through the inverse scattering data of the Schrödinger operator $`L=_x^2+u(x,0)`$ were found. This procedure has been called the Inverse Scattering Transform (IST). It was extended later for some other highly nontrivial $`(1+1)`$-systems including such famous systems as Nonlinear Schrödinger $`NS_\pm :i\psi _t=\psi _{xx}\pm |\psi |^2\psi `$ and Sine(Sinh)-Gordon equations $`SG:u_{xt}=\mathrm{sin}u`$ or $`u_{xt}=\mathrm{sinh}u`$. Note, that for the SG equation a large family of exact solutions was already constructed in XIX century by Bianchi, Lie and Bäcklund (see in ). Beginning with 1974, several $`(2+1)`$-dimensional physically interesting systems have been discovered as integrable by the IST procedure. The most famous of them is the KP system (see ).
It is necessary to emphasize that the IST procedure in its original form can not be applied to the solution of the periodic problem (i.e. $`u(x,t)`$ is periodic in the variable $`x`$). This problem was solved on the base of the new approach proposed in , in the works (see the surveys ). An extension of this method to $`(2+1)`$-systems was found in . New development of this approach associated with two-dimensional Schrödinger operator was started in 1976 (see ).
Complete detailed description of the solution of the periodic problem can be found in the surveys, Encyclopedia articles , and in the book . We are going to present here basic ideas of this theory in the simplest form possible. Let us point out that KdV as well as other nontrivial completely integrable by IST PDE systems are indeed completely integrable in any reasonable sense for rapidly decreasing or periodic (quasi-periodic) boundary conditions, only. In fact, even that is well established for few of them. For example, for the KdV any periodic solution can be approximated by the finite-gap solutions. This statement easily follows from the theory of finite-gap potentials if we do not try to preserve the period, i.e. in the class of all quasi-periodic finite-gap potentials. The approximation of any periodic potential by the finite-gap potentials with exactly the same period, was constructed on the base of other approach developed in . The extension of the theory of Riemann surfaces and $`\theta `$-functions to the specific class of surfaces of the infinite genus associated with periodic Schrödinger operator, was done in . This theory is a beautiful description of the infinite limit. However, it seems that all fundamental properties of $`\theta `$-functions associated with the complex continuation of variables, are lost in this limit. It is interesting to point out that analogous (but more complicated) theory of Riemann surfaces of the infinite genus was developed later in for the periodic 2D Schrödinger operators.
Outside these functional classes almost no effective information is known. Beautiful methods have been developed also for the studies of the special self-similar and ”string-type” solutions, but in most cases they lead to the very hard analytical problems associated with the famous Painleve equations and their generalizations ().
## 2 Rapidly decreasing potentials and GGKM Procedure. Bäcklund Transformations
Let us recall the basic information about the IST method for KdV. We start from the so-called Lax Representation for KdV (see ). The Heisenberg type equation for the Schrödinger operator $`L`$
$$L_t=[L,A]=LAAL,$$
(2.1)
$$L=_x^2+u,A=4_x^3+3(u_x+_xu),$$
(2.2)
is equivalent to the identity (KdV equation)
$$u_t=6uu_x+u_{xxx}.$$
(2.3)
For this reason, any KdV type equations admitting some analog of the Lax representation is called isospectral deformations. The existence of such deformations indicates the possibility of effective solution of the Inverse Scattering Problem for the operator $`L`$.
For the rapidly enough decreasing functions $`u(x,t)0,x\pm \mathrm{}`$ we define two bases of solutions ($`t`$ is fixed):
$`\varphi _\pm (x,t;k)\mathrm{exp}^{\pm ikx},`$ $`x\mathrm{},`$ $`L\varphi _\pm =\lambda \varphi _\pm ,`$ (2.4)
$`\psi _\pm (x,t;k)\mathrm{exp}^{\pm ikx},`$ $`x+\mathrm{},`$ $`L\psi _\pm =\lambda \psi _\pm ,k^2=\lambda .`$ (2.5)
By definition, monodromy matrix $`T`$ connects these two bases, $`T\varphi =\psi `$ for the column vectors $`\varphi =(\varphi _+,\varphi _{}),\psi =(\psi _+,\psi _{})`$:
$`T=\left(\begin{array}{cc}a& b\\ c& d\end{array}\right),\psi _+=a\varphi _++b\varphi _{},\psi _{}=c\varphi _++d\varphi _{}.`$ (2.8)
A conservation of the Wronskian implies that $`detT=adbc=1`$. For the real values of $`k`$ or $`\lambda >0`$ we have $`a=\overline{d},c=\overline{b}`$. Therefore, $`|a|^2|b|^2=1`$. The whole set of the so-called inverse scattering data can be extracted from the monodromy matrix $`T`$ if it is well-defined for all complex values of $`k`$. The so-called scattering matrix is constructed from $`T`$ for the real $`k`$. Its entries are the transmission coefficient $`1/a`$ and the reflection coefficient $`\overline{b}/a`$. The property $`b=0`$ for all real values of $`k`$ characterizes reflexionless or multisoliton potentials. For all rapidly decreasing potentials the matrix element $`a(k)`$ is well-defined for complex $`k`$ such that $`\mathrm{Im}k>0`$ and $`a1`$ for $`k\mathrm{},\mathrm{Im}k>0`$. There is only a finite number of purely imaginary zeroes $`a(k_n)=0`$ in this domain. They correspond to the discrete spectrum $`\lambda _n=k_n^2<0`$.
The famous result of (GGKM procedure) easily follows from the Lax representation, which implies the equation:
$`T_t=[T,\mathrm{\Lambda }],\mathrm{\Lambda }=\left(\begin{array}{cc}4ik& 0\\ 0& 4ik\end{array}\right)`$ (2.11)
This result was formulated as a set of the following GGKM formulas
$`a_t=0,b_t=8(ik)^3=c_t,d_t=0.`$ (2.12)
The latter equations give a full description of KdV dynamics in these variables because any rapidly decreasing potential can be reconstructed from the inverse scattering data. A special family of the reflexionless potentials where $`b=0`$ for real values of $`k`$, leads to the so-called multisoliton solutions for KdV equation (see .)
The multisoliton solutions can be also directly obtained with the help of the elementary substitutions (Bäclund Transformations) transforming any solution of the KdV into another solution: let $`u`$ be a solution of the KdV equation and $`v`$ be a solution of the Ricatti equation $`\alpha +u=v_x+v^2`$ with the initial value independent of time. The new function $`\stackrel{~}{u}=v_x+v^2`$ satisfies the KdV equation. Starting from the trivial solution $`u=u_0=0`$ we construct a sequence of potentials $`\stackrel{~}{u}_{n1}=u_n,n>0,`$ given by the Bäcklund Transformation. We choose parameters $`\alpha _n,\alpha _1>\alpha _2>\mathrm{}>\alpha _n>\mathrm{}`$, and take the real nonzero functions $`f_n\mathrm{},x\pm \mathrm{},f_{nxx}+u_{n1}f_n=\alpha _nf_n`$, which define $`v_n=(\mathrm{log}f_n)_x`$. Every such sequence leads to the multisoliton reflexionless potential:
$$u_0=0,u_1=\frac{2\alpha }{ch^2(\sqrt{\alpha }(xx_0)+\beta t)},\mathrm{}$$
(2.13)
In terms of the Schrödinger operator this transformation (invented by L.Euler in 1742) is called Darboux Transformation. The operator $`L`$ can be factorized
$$L=^2+u=(+v)(v).$$
Using the non-commutativity of these factors, we define the Darboux Transformation for the operator and its eigenfunction in the following way:
$$\stackrel{~}{L}=^2+\stackrel{~}{u}=(v)(+v),\stackrel{~}{\psi }=(+v)\psi .$$
(2.14)
These transformations can be considered as some kind of discrete spectral symmetries for Schrödinger operators. They preserve a spectrum of the operator $`L`$ (maybe except for one eigenfunction).
## 3 KdV Hierarchy. Integrals of motion. Hamiltonian formalism.
The local integrals for the KdV equation can be constructed with the help of the Schrödinger operator. Consider the associated Riccati equation $`v_x+v^2=uk^2`$ and find the solution for it as a formal series in the variable $`k`$:
$$v(x,k)=ik+\underset{n=1}{\overset{\mathrm{}}{}}v_n(x)(ik)^n,$$
(3.1)
where all $`v_n`$ are polynomials in the variables $`u,u_x,\mathrm{}`$ The integral along the line $`x`$ is a $`k`$-dependent constant of motion for the KdV equation
$$_t\left(v(x,k)𝑑x\right)=0,u_t=6uu_x+u_{xxx}.$$
(3.2)
For the real potentials $`u(x,t)`$ and real $`k`$ we can see that the imaginary part of $`v(x,k)ik`$ is a total derivative. The remaining quantities in the expansion :
$$(v(x,k)ik)𝑑x=\underset{n1}{}v_n(x)(ik)^n𝑑x.,$$
(3.3)
define local integrals of motion
$$I_n=c_nv_{2n+3}(x)𝑑x,n=1,0,1,2,\mathrm{},$$
(3.4)
where $`c_n`$ are constants. After the proper choice of the constants $`c_n`$ we have
$$I_1=u𝑑x,I_0=u^2𝑑x,I_1=(u_x^2/2+u^3)𝑑x,\mathrm{},$$
(3.5)
$$I_{nt}=0,u_t=6uu_x+u_{xxx}.$$
(3.6)
Let us introduce a GZF Poisson bracket on the space of functions
$$\{u(x),u(y)\}=\delta ^{}(xy)).$$
(3.7)
Then any functional $`H`$ (Hamiltonian) defines the corresponding Hamiltonian system
$$u_t=_x\left(\frac{\delta H}{\delta u(x)}\right).$$
(3.8)
For the case $`H=I_1`$ we get a trivial flow (i.e. this integral is a Casimir for the GZF bracket). For $`H=I_0`$ we are coming to the $`x`$-translations $`u_t=u_x`$. Let us call this equation $`KdV_0`$. For the case $`H=I_1`$ we have the ordinary $`KdV=KdV_1`$. Higher integrals give us the equations $`KdV_n`$ of the order $`2n+1`$ admitting the Lax representations with the same Schrödinger operator $`L`$ but with the differential operators $`A_n=(const)_x^{2n+1}+\mathrm{}`$ :
$$u_{t_n}=_x\left(\frac{\delta I_n}{\delta u(x)}\right)=[L,A_n].$$
(3.9)
In particular, $`A_0=_x,A_1=A.`$ Nice formula for all operators $`A_n`$ can be extracted from .
Let $`L=^2`$, where $`=_x+_{k1}a_k(u,u_x,\mathrm{})_x^k`$. Here all $`a_k`$ are polynomials in the variables $`u,u_x,\mathrm{}`$ and $`_x^1a=_{n0}(1)^na^{(n)}_x^{n1}`$ for the composition of the operator $`^1`$ and multiplication operator by $`a`$. By definition,
$$A_n=(^{2n+1})_+=(L^{2n+1/2})_+,$$
(3.10)
where the sign $`+`$ means omitting of all strictly negative powers of $`_x`$.
All higher $`KdV_n`$ systems can be integrated by the same IST procedure for the class of rapidly decreasing functions. In particular, GGKM equations for the scattering data (or monodromy matrix) have the form
$$T_{t_n}=[T,\mathrm{\Lambda }_n],\mathrm{\Lambda }_n=(const)\left(\begin{array}{cc}(ik)^{2n+1}& 0\\ 0& (ik)^{2n+1}\end{array}\right).$$
(3.11)
The latter result also implies that all these flows commute with each other. Hence, we get the following conclusion without any calculation:
integrals $`I_n`$ have zero Poisson brackets,
$$\{I_n,I_m\}=\frac{\delta I_n}{\delta u(x)}_x\frac{\delta I_m}{\delta u(x)}dx=0.$$
(3.12)
A generalization of the GZF Poisson bracket for the isospectral deformations of the higher order (scalar) Lax operators $`L`$ was found in .
It should be emphasized that there exists a family of local field-theoretical Poisson brackets (LM-brackets, ) describing the KdV theory:
$$B=B_{(\lambda ,\mu )}=\lambda _x+\mu (_x^3+4u_x+2u_x)=\lambda B_0+\mu B_1,$$
(3.13)
$$\{u(x),u(y)\}_{\lambda ,\mu }=B\delta (xy),$$
(3.14)
$$\{I_n,I_m\}_{\lambda ,\mu }=0.$$
(3.15)
These brackets were generalized for the higher order operators $`L`$ in .
The recurrence operator $`B_1B_0^1=C`$ generates all the right-hand sides of all higher KdV systems:
$$C(0)=u_x,C^2(0)=6uu_xu_{xxx},\mathrm{},C^n(0)=_x\left(\frac{\delta I_{n1}}{\delta u(x)}\right).$$
(3.16)
It gives also a simple proof of a very useful identity ():
$$\frac{\delta I_n}{\delta u(x)}𝑑x=(const)I_{n1}.$$
(3.17)
All these identities are local and can be used for the class of periodic functions as well. However, we shall see in the next section that the direct analog of GGKM procedure does not lead to the integration procedure. We are going to use a different approach.
## 4 Spectral theory of periodic Schrödinger operators. Finite-gap potentials
The spectral theory for the periodic potentials on the whole line $`x`$ is based on the monodromy matrix as in case of the scattering theory. However, in the periodic case with a period $`T<\mathrm{}`$, we have nothing like the selected point $`x=\mathrm{}`$ for the definition of the monodromy matrix (as it was for the rapidly decreasing case $`(T=\mathrm{})`$–see Section 2 above). Any point $`x_0`$ can be used. Let us fix initial point $`x_0`$ and choose a special basis of the solutions $`C(x,x_0,ϵ),S(x,x_0,ϵ)`$ for the spectral equation $`LC=ϵC,LS=ϵS`$ such that for $`x=x_0`$ we have:
$`\left(\begin{array}{cc}C& S\\ C_x& S_x\end{array}\right)=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)`$ (4.5)
The shift operator $`T:xx+T`$ in the basis $`C,S`$ defines the monodromy matrix
$$\widehat{T}(x_0,ϵ)=\left(\begin{array}{cc}a& b\\ c& d\end{array}\right),$$
(4.6)
$$C(x+T)=aC(x)+bS(x),S(x+T)=cC(x)+dS(x).$$
(4.7)
The key element of the periodic spectral theory is a notion of the, so-called, Bloch waves or Bloch-Floquet eigenfunctions. We present here some essential properties of these functions without any proofs. An exposition of this theory may be found in the Encyclopedia article , where the main ideas of the proofs are clearly presented for the difference Schrödinger operator (it is much simplier).
By definition, the Bloch-Floquet functions are solutions of the Schrödinger equation that are at the same time eigenfunctions of the shift operator, i.e.
$$L\psi =ϵ\psi ,T\psi (x)=\psi (x+T,ϵ)=\mathrm{exp}(\pm ip(ϵ)T)\psi (x).$$
(4.8)
We uniquely normalize $`\psi `$ by the condition $`\psi |_{x=x_0}=1`$.
For any complex number $`ϵ`$ the eigenvalues, $`w_\pm (ϵ)=\mathrm{exp}(\pm p(ϵ)T)`$, of the shift operator are defined by the characteristic equation for the monodromy matrix. ¿From the Wronskian property it follows that $`det\widehat{T}=1`$. Therefore, the characteristic equation has the form
$$w^2(\mathrm{tr}\widehat{T})w+1=0.$$
(4.9)
The multivalued function $`p(\epsilon )`$ is called quasi-momentum.
The spectrum of the Schrödinger operator on the whole line is a union of spectral zones which are segments of the real line of the variable $`ϵ`$, where the quasi-momentum is real. The latter condition is equivalent to the inequality:
$$|\mathrm{Tr}T|=2\mathrm{cos}(pT)2.$$
(4.10)
The typical graph of the function $`f(ϵ)=\mathrm{cos}(pT)`$ is presented on the fig 1. In particular, its extreme points $`f^{}(ϵ)=0`$ ”generically” are located inside of the gaps (i.e. for the open and everywhere dense set of periodic potentials we have $`|f|>1`$ at the extreme points, and there is only one extremal point in each gap). For some special cases we may have $`f=\pm 1`$ at the extreme point. Such point lies inside the spectral zone. However, generic perturbations create new small gap nearby this point — see fig 2. This point is a double point of periodic or anti-periodic spectrum with the boundary conditions $`\psi (x)=\pm \psi (x+T)`$.
The Riemann surface of the Bloch-Floquet functions is defined by the equation
$$z^2=\mathrm{cos}(p(ϵ)T),$$
(4.11)
but this surface is nonsingular only for the generic case (when there are no double points of the periodic or anti-periodic problem for the Schrödinger operator). For the large $`ϵ+\mathrm{}`$ the following asymptotics is valid for the gaps (see fig. 3).
1. The length of gaps tends to zero with its rate depending on the smoothness of the potential (this rate is exponential for the analytic potentials).
2. All gaps are located nearby the points $`ϵ_m=4\pi ^2m^2T^2`$, and the distance of gap from this point tends to zero.
A nonsingular Riemann surface of the Bloch-Floquet solutions is defined with the help of the equation
$$z^2=(ϵϵ_0)(ϵϵ_1)\mathrm{}(ϵϵ_n)\mathrm{},$$
(4.12)
where $`ϵ_0<ϵ_1<\mathrm{},`$ are simple eigenvalues of the periodic and anti-periodic spectral problems.
As it follows from the Lax representation for all higher $`KdV_n`$ systems (3.9), these boundary spectral points for the periodic operator $`L`$ in the Hilbert space $`L^2(R)`$ of the square integrable complex-valued functions on $`x`$–line are integrals of motion for the KdV hierarchy.
We can say that the Riemann surface of the Bloch-Floquet solutions as a whole is an integral of motion for the KdV hierarchy. How many potentials correspond to the same Riemann surface (i.e. have the same spectrum in $`L^2(R)`$)?
The original idea to introduce a special class of potentials for which this problem can be solved effectively was proposed in and is based on the KdV hierarchy. This idea naturally combines two different ways to describe the corresponding isospectral manifold of potentials. Both are in fact closely related to each other.
The first approach – use of the KdV hierarchy and Hamiltonian dynamics:
Let us consider a stationary equation for some linear combination of higher $`KdV_n`$ flows (3.9). It is an ordinary differential equation that can be written in the form:
$$\delta (I_n+c_1I_{n1}+\mathrm{}+c_nI_0+c_{n+1}I_1)=0.$$
(4.13)
This is a finite-dimensional Hamiltonian system with $`n`$ degrees of freedom depending on $`n+1`$ parameters $`c_1,\mathrm{},c_{n+1}`$. It is a completely integrable Hamiltonian system because there is a pencil of the commuting flows. This pencil coincides with the restriction of all higher $`KdV_m`$ on this stationary subset of functions given by the equation (4.13). Therefore, its generic nonsingular solution is expected to be a periodic or quasiperiodic function of $`x`$.
In the next section we shall construct for (4.13) some kind of the Lax representation
$$\mathrm{\Lambda }_x=[Q(x,\lambda ),\mathrm{\Lambda }(x,\lambda )],\mathrm{\Lambda }=\mathrm{\Lambda }_n+\underset{k=1}{\overset{n+1}{}}c_k\mathrm{\Lambda }_k,$$
(4.14)
using $`2\times 2`$ traceless matrices depending on the parameter $`\lambda `$ and $`(u,u_x,\mathrm{})`$ polynomially. In the today terminology some people call them loop groups. Note that in all soliton systems this $`\lambda `$-loops have very specific $`\lambda `$-dependence (polynomial, rational and in some exotic examples – elliptic functions).
Lax representation (4.14) implies that an algebraic Riemann surface defined by characteristic equation
$$det(\mathrm{\Lambda }(\lambda )zI)=P(\lambda ,z)=0,$$
(4.15)
does not depend on $`x`$ and therefore is an integral of (4.13). The same Riemann surface can be extracted from the commutativity equation $`[L,A]=0`$, which according to (3.9) is equivalent to (4.13).
It should be emphasized that the latter form of (4.13), i.e. the commutativity condition for two ordinary differential operators
$$[L,A_n+c_kA_{nk}]=[L,A]=0,L=_x^2+u,$$
(4.16)
was considered formally (i.e. locally in the variable $`x`$ without any periodicity assumptions) as a pure algebraic problem in 1920-s (see ). Even the formal algebraic Riemann surface (4.15) appeared as a relation $`P(L,A)=0`$.
According to our logic however, the corresponding system of equations in the variable $`x`$ is Hamiltonian and completely integrable. Therefore, its generic solution is quasiperiodic in $`x`$, containing a dense family of periodic solutions. We may ask about the spectrum of the corresponding operators $`L`$ in $`L^2(R)`$ and boundaries of gaps. Remarkably, they exactly coincide with the branching points of the Riemann surface defined by equation (4.15). Therefore, for the periodic potential which satisfy equation (4.13) the nonsingular spectral Riemann surface of the Bloch-Floquet solutions is an algebraic Riemann surface of genus $`g=n`$. The spectrum of such an operator contains only a finite number of gaps $`[ϵ_{2j1},ϵ_{2j}],j=1,\mathrm{},n.`$
This key step unifies the first approach with the second one (below). So the solution of the inverse spectral problem can be identified with the process of solution of some special families of completely integrable systems using Riemann surfaces.
The second approach – periodic spectral theory in the Hilbert space $`L^2(R)`$:
For any periodic Schrödinger operator $`L=^2+u(x=t_0,t_1,\mathrm{})`$ we have already defined the Riemann surface $`\mathrm{\Gamma }`$ of Bloch-Floquet solutions with the help of the monodromy matrix $`\widehat{T}`$. For the finite-gap potentials a graph of the function $`f=\mathrm{cos}(p(ϵ)T=1/2\mathrm{t}\mathrm{r}\widehat{T}`$ is highly degenerate (see fig.2 and compare with fig.1). For all real and large enough $`ϵ`$ we have $`|f|1`$.
Once again, we ask how to describe all the potentials with the same finite-gap spectrum and what are additional variables which uniquely define a finite-gap potential? It turns out that these additional spectral data are the poles of the Bloch-Floquet functions. It can by shown that the Bloch-Floquet solutions defined by (4.8) and normalized by the condition $`\psi |_{x=x_0}=1`$ have exactly one simple pole $`\gamma _j(x_0)`$ inside each gap or on its boundary. A point of the hyperelliptic surface $`\mathrm{\Gamma }`$ (4.12) can be represented as the complex number $`ϵ`$ and a sign (or a branch) of the radical $`z=\sqrt{(ϵϵ_0)(ϵϵ_1)\mathrm{}}`$. The branches of the radical coincide at the boundaries of gaps. Therefore, each gap should be considered as a closed cycle $`a_j,j=1,2,\mathrm{}n`$ (see fig. 4). Any pole of $`\psi `$ should be considered as a point $`\gamma _j`$ on the cycle $`a_j`$. The above-mentioned statement that there is only one pole in a gap means, in particular, that $`\psi `$ has no pole at the point $`(\gamma _j,+)`$ if it already has pole at $`(\gamma _j,)`$ and vice versa. Geometrically, the total set of poles $`(\gamma _1,\mathrm{},\gamma _n)`$ represents a point on the real torus $`T^g=a_1\times \mathrm{}\times a_n`$. As we shall see later, this set of poles completely determines the original potential. Different points of the real torus lead to the different potentials (if normalization point $`x_0`$ is fixed).
## 5 Periodic analog of GGKM. Zero Curvature representation for the KdV hierarchy and corollaries
As it was emphasized above, the definition of the monodromy matrix in the periodic case depends on a choice of initial point $`x_0`$. For different choices of the initial point the corresponding monodromy matrices are conjugated (because they represent the same linear transformation in different bases). Therefore, a dependence of the monodromy matrix with respect to the choice of the initial point $`x_0`$ can be described by the equation:
$$T_{x_0}=[Q(ϵ,x_0),T],Q=\left(\begin{array}{cc}0& 1\\ ϵu(x_0)& 0\end{array}\right).$$
(5.1)
In the same way, for the isospectral deformations corresponding to the KdV hierarchy (i.e, all higher $`KdV_n`$ systems) we can establish the following equations (the periodic analogs of GGKM):
$$T_{t_j}=[\mathrm{\Lambda }_j,T],\mathrm{\Lambda }_0=Q,t_0=x_0.$$
(5.2)
A compatibility condition of equations (5.2) for any pair of variables $`t_i,t_j`$ implies the following zero curvature representation for the KdV hierarchy (where periodic boundary conditions are already inessential)
$$[_{t_i}\mathrm{\Lambda }_i,_{t_j}\mathrm{\Lambda }_j]=0.$$
(5.3)
Using the Lax representations (3.9) for all higher $`KdV_n`$, we can express all the matrices $`\mathrm{\Lambda }_j`$ as polynomials in the variable $`ϵ`$ and variables $`u(x_0),u_{x_0}(x_0),\mathrm{}`$. For example, for the ordinary KdV we have (replacing $`x_0`$ by $`x`$):
$$\mathrm{\Lambda }_1=\left(\begin{array}{cc}u_x& 2u+4ϵ\\ 4ϵ^2+2ϵu+2u^2u_{xx}& u_x\end{array}\right)$$
(5.4)
The matrices $`\mathrm{\Lambda }_k`$ can be completely reconstructed from equations (5.3) and from the following properties of matrix elements:
$$\mathrm{\Lambda }_k=(const)\left(\begin{array}{cc}a_k& b_k\\ c_k& d_k\end{array}\right),$$
(5.5)
$$d_k+a_k=\mathrm{Tr}\mathrm{\Lambda }_k=0,b_k=ϵ^k+u/2ϵ^{k1}+\mathrm{},$$
(5.6)
$$det\mathrm{\Lambda }_k=(a^2+bc)=(const)R_{2k+1}(ϵ)=(const)(ϵ^{2k+1}+\mathrm{}),$$
(5.7)
$$2a_k=b_{kx},(R_{2k+1})_x=b_k.$$
(5.8)
For $`i=0`$, equations (5.3), where $`t_0=x`$, give the zero-curvature representation for the $`KdV_n`$ system:
$$_x\mathrm{\Lambda }_n_{t_n}Q=[Q,\mathrm{\Lambda }_n].$$
(5.9)
From this representation and from the periodic analog of GGKM we are coming to the following results for the stationary higher $`KdV_n`$ system (4.14):
1. The monodromy matrix $`T`$ commutes with $`\mathrm{\Lambda }`$
$`[T,\mathrm{\Lambda }]=0.`$ (5.10)
Therefore, they have common eigenvectors. It implies, in particular, that the Bloch-Floquet function $`\psi _\pm `$ is single-valued on the algebraic Riemann surface associated with the matrix $`\mathrm{\Lambda }`$;
2. The stationary higher KdV admits an $`ϵ`$-parametric Lax-type representation in the variable $`x`$ (4.14):
$`(\mathrm{\Lambda })_x=[Q,\mathrm{\Lambda }],\mathrm{\Lambda }=\mathrm{\Lambda }_n+{\displaystyle c_k\mathrm{\Lambda }_k}.`$ (5.11)
Therefore, we have a full set of conservation laws organized in form of the Riemann surface (i.e. all coefficients of the polynomial $`P`$ are $`x`$-independent):
$$det(\mathrm{\Lambda }(ϵ)zI)=P(ϵ,z)=0,$$
(5.12)
$$P(ϵ,z)=z^2R_{2n+1}(ϵ)=(const)(ϵϵ_0)\mathrm{}(ϵϵ_{2n}).$$
(5.13)
These points $`ϵ_j`$ are exactly the boundaries of gaps for periodic potentials because the Riemann surfaces of the matrices $`T`$ and $`\mathrm{\Lambda }`$ coincide. The matrix element $`b_{12}`$ for the matrix $`\mathrm{\Lambda }`$ determines another set of points
$`b_{12}(ϵ)=(const)(ϵ\gamma _1)\mathrm{}(ϵ\gamma _n).`$ (5.14)
These points coincide with projections of the zeroes of Bloch wave $`\psi _\pm `$ as the functions of the variable $`x`$ or poles as the functions of $`x_0`$ (see the next section).
An original approach to the solution of the inverse spectral problem was based on the use of the following trace type formula for the potential:
$$u(x)/2+const=\gamma _1(x)+\mathrm{}\gamma _n(x).$$
(5.15)
¿From (5.11) and (5.14), “Dubrovin equations” defining dynamics in $`x`$ of the points $`\gamma _j`$ can be derived. They have the form
$$\gamma _{jx}=\frac{\sqrt{R_{2n+1}}(\gamma _j)}{_{kj}(\gamma _j\gamma _k)}$$
(5.16)
These equations (see ) can be linearized by the, so-called, Abel transformation.
Consider the first kind differentials on the Riemann surface (i.e. holomorphic 1-forms without poles anywhere - even at the infinity). The basic first kind forms are
$`\omega _j={\displaystyle \frac{ϵ^{j1}dϵ}{\sqrt{R_{2n+1}(ϵ)}}},j=1,\mathrm{},n.`$ (5.17)
It it is convenient to choose a normalized basis taking linear combinations $`\mathrm{\Omega }_s=_jv_{sj}\omega _j`$ such that
$$_{a_j}\mathrm{\Omega }_s=\delta _{sj},$$
for all the gaps $`a_j`$.
Fix a set of paths $`\kappa _j`$ from the point $`P_0=\mathrm{}`$ to the points $`P_j`$ on the Riemann surface. Abel transformation is defined by the formula
$`A_p(P_1,\mathrm{},P_n)={\displaystyle \underset{j}{}}{\displaystyle _{\kappa _j}}\mathrm{\Omega }_p.`$ (5.18)
Equations (5.16) after the Abel transform become linear:
$$A_{px}=U_q=\frac{1}{2\pi }_{b_q}𝑑p.$$
(5.19)
Here the closed paths $`b_q`$ are ”canonically conjugated” to the paths $`a_j`$ (see fig. 5). It means that the intersection numbers are $`a_jb_q=\delta _{jq}`$. The differential $`dp`$ is equal to $`p_ϵdϵ`$ where $`p(ϵ)`$ is a multi-valued quasimomentum. It is a second kind differential (i.e. meromorphic 1-form on the surface $`\mathrm{\Gamma }`$ with a pole of order 2 at the point $`\mathrm{}`$ with negative part $`dk=(dw)w^2`$ in the local coordinate $`w=k^1`$ near $`\mathrm{}`$), and therefore has the form
$$dp=\frac{(ϵ^{n+1}+_{j>0}v_jϵ^{n1j})dϵ}{\sqrt{R_{2n+1}(ϵ)}}.$$
(5.20)
All the coefficients $`v_j`$ can be found from the normalization condition:
$$_{a_j}𝑑p=0,j=1,\mathrm{},n.$$
(5.21)
Following the classical XIX century theory of Riemann surfaces, formulas (5.15-5.21) lead to some expression of the potential through the $`\theta `$-functions avoiding the calculation of the eigenfunction $`\psi _\pm `$ (see ). The most beautiful $`\theta `$-functional formula for the potential was obtained in the work .
However, in the next section we shall not follow this approach for solving the inverse problem which has been outlined above, but use another approach proposed for the first time in the article (see appendix, based on the idea of A.Its): it is possible to calculate all the family of the Bloch-Floquet functions. The original approach has been used later in the problems, where the effective calculation of eigenfunction is impossible (like in the cases of higher rank commuting OD linear operators (see for example ).
## 6 Solution of the periodic and quasiperiodic inverse spectral problems. Baker-Akhiezer functions
We are going to solve the inverse spectral problem for the finite-gap potentials following the scheme proposed in and based on the concept of the Baker-Akhiezer functions. These functions are uniquely defined by their analytical properties on the spectral Riemann surface. As we shall see later this scheme is evenly applicable to the solution of inverse problems in two-dimensional case where the corresponding analytical properties naturely generalize the analytical properties of the Bloch-Floquet solutions for finite-gap Schrödinger operators.
Let us start with the following real inverse spectral data:
1. Riemann surface $`\mathrm{\Gamma }`$ given in the form
$$z^2=(ϵϵ_0)\mathrm{}(ϵϵ_{2n}),$$
where all numbers $`ϵ_j`$ are real;
2. a set of $`n`$ real points $`\gamma _j\mathrm{\Gamma }`$ such that there is exactly one point $`\gamma _j`$ on the cycle $`a_j`$.
Below we shall for brevity identify the point $`\gamma _j`$ as a point on $`\mathrm{\Gamma }`$ with its projection on complex $`ϵ`$–plane. The condition that there is only one point on each cycle means that projection of the points satisfy the restriction:
$$ϵ_{2j1}\gamma _jϵ_{2j},j=1,\mathrm{}n,$$
(6.1)
By definition, the complex inverse spectral data are the same data where the Riemann surface (i.e. its branching points $`ϵ_k`$) and the points $`\gamma _j\mathrm{\Gamma }`$ are arbitrary complex points.
As we shall see below a generic set of algebraic-geometric spectral data leads to the explicit solution of the inverse spectral problem in terms of the Riemann $`theta`$-function. The corresponding potentials are complex meromorphic quasi-periodic functions of the variable $`x`$. For the real data described above, we are coming to the smooth (even analytic) quasiperiodic potentials with their periods expressed through the hyperelliptic integrals (see below). It is necessary to mention that there is no way to find simple and effective criteria for the potential to be periodic in terms of these data. The periods depend on the Riemann surface, only. Of course, we may write the condition that all the corresponding hyperelliptic integrals are commensurable, however this condition in transcendental. Recently, based on the results of , the effective numerical approach has been developed for the solution of this problem. It is based on the discovery of some specific dynamic systems on the set of potentials which preserve all the periods but change spectrum. Following , we start from one periodic potential and create all others using these dynamic systems.
In the last section we shall present some examples of finite-gap potentials written in terms of elliptic functions. They are periodic in the variable $`x`$ (even double-periodic as functions of the complex variable $`x`$). The first nontrivial examples different from the classical Lame potentials $`u(x)=n(n+1)\mathrm{}(x)`$ were found in . This subject was developed in the works .
We define the Baker-Akhiezer function $`\psi =\psi (x,t_1,\mathrm{}t_n;P)`$ for the parameters $`x,t_j`$ and the point $`P=(ϵ,\pm )\mathrm{\Gamma },`$ by its analytical properties on $`\mathrm{\Gamma }`$ with respect to the variable $`P`$. For the case $`t_1=t_2=\mathrm{}=0`$ these analytical properties are just the same as the analytical properties of the Bloch-Floquet solutions of the periodic finite-gap operator.
¿From pure algebraic-geometric arguments it follows that there exists a unique function $`\psi `$ such that
1. it is meromorphic on $`\mathrm{\Gamma }`$ outside the infinity and has at most simple poles at the points $`\gamma _j,j=1,\mathrm{},n`$;
2. in a neighborhood of the infinity the function $`\psi `$ has the form
$$\psi =\mathrm{exp}[xk+t_1k^3+\mathrm{}+k^{2n+1}t_n](1+\xi _1(x,t)k^1+\mathrm{}),$$
where $`k^2=ϵ`$ and therefore, $`k^1`$ is a local coordinate on the Riemann surface near infinity.
A proof of this statement is identical to the proof of the existence and uniqueness of general Baker-Akhiezer functions that were introduced in for construction of exact solutions of two-dimensional KP equation and all associated Zakharov-Shabat hierarchies.
General Baker-Akhiezer function is defined with the help of an arbitrary Riemann surface of finite genus $`n`$ instead of special hyperelliptic (i.e. two-sheeted) surfaces. We fix an arbitrary ”infinity” point $`P_0`$ on it; local coordinate $`k^1=w`$ where $`w(P_0)=0`$ ; generic set of points $`(\gamma _1,\mathrm{},\gamma _n)`$, and numbers $`(x=\tau _1,\tau _2,\mathrm{},\tau _k,\mathrm{})`$. The corresponding Baker-Akhiezer function has the same analytical properties as above but, (2) is replaced by the following:
$$\psi =\mathrm{exp}[xk+\tau _1k^2+\tau _2k^3+\mathrm{}](1+\xi _1(x,\tau )k^1+\mathrm{}).$$
(6.2)
For the KdV hierarchy we have a hyperelliptic Riemann surface, $`k^2=ϵ`$, and $`\tau _{2k}=0,\tau _{2j+1}=t_j`$.
The use of such type of functions (dependent only on a single parameter $`x`$) has been proposed by Baker in his note , for the common eigenfunctions of two commuting OD linear operators. He expected that this construction will improve the results of . He made also very interesting conjecture that this approach may seriously improve the classical theory of $`\theta `$-functions. Unfortunately, this program was not realized and was forgotten. The soliton theory appeared many decades later absolutely independently. In the 70-s it started to use such kind of functions on Riemann surfaces in the process of solution of a periodic problem for the KdV type systems and inverse spectral periodic problems. In the classical spectral theory Akhiezer () was the first to use some special cases of this function for the construction of some examples of operators on the half-line $`x0`$ with interesting spectral properties. No one of the authors before the 70-s had associated anything like that with periodic problems. Indeed, this type of functional construction on the Riemann surface was extracted in 1974 from the work .
We shall prove the existence and uniqueness of the Baker-Akhiezer function and present its exact expression through the $`\theta `$-functions later. At this moment we would like to show how the uniqueness of this function leads to the proof that $`\psi `$ is an eigenfunction for the Schrödinger operator for our special case. Apply twice the operator $`=_x`$ to this function and use analytical properties. After elementary calculation, we are coming to the formulas:
$`\psi `$ $`=`$ $`k\psi +\mathrm{exp}[kx+\mathrm{}](\xi _{1x}k^1+\mathrm{}),`$
$`^2\psi `$ $`=`$ $`k^2\psi +2\xi _{1x}\psi +O(k^1)\mathrm{exp}[kx+\mathrm{}].`$ (6.3)
¿From the latter equations, we get the equality
$$(^2k^22\xi _{1x})\psi =O(k^1)\mathrm{exp}[kx+\mathrm{}].$$
The left-hand side is globally well-defined function on the same Riemann surface because $`k^2=ϵ`$. It has the same poles independent of parameters. Up to the same exponential factor, it is of the order $`O(k^1)`$ at the infinity. Therefore, it is equal to zero due to the uniqueness of the Baker-Akhiezer function.
So we are coming to the conclusion that
$`L\psi =ϵ\psi ,L=^2+u,u=2\xi _{1x}.`$ (6.4)
We can apply the operators $`_j=_{t_j}`$ to the function $`\psi `$:
$`_j\psi =k^{2j+1}\psi +\mathrm{exp}[kx+\mathrm{}](\xi _{1j}k^1+\xi _{2j}k^2).`$ (6.5)
As in the previous case, we can easily construct a linear operator $`A_j=_x^{2j+1}+\mathrm{}`$ with the coefficients independent of $`k`$ such that $`(_jA_n)\psi =O(k^1)\mathrm{exp}(kx+\mathrm{})`$. Using global analytical properties on the Riemann surface, we deduce from this that the left-hand side is equal to zero as before. The compatibility condition of these pairs of equations is exactly a $`KdV_n`$-system for the potential $`u(x,t_1,\mathrm{},t_m)`$.
Following the works , we can prove in the same way that for the general Baker-Akhiezer function $`\psi (x,\tau ,P)`$ associated with arbitrary Riemann surface with fixed local coordinate near a puncture, the following equations are valid:
$`(_{\tau _k}L_k)\psi =0,k=2,3,\mathrm{}.`$ (6.6)
Here $`L_k=_x^k+\mathrm{}`$ are linear OD operators acting on the variable $`x`$ with coefficients depending on parameters $`\tau `$. These coefficients are differential polynomials in the coefficients of the expansion of the regular factor of the Baker-Akhiezer function at the puncture (infinity). They are uniquely defined by the condition that the expansion of the left-hand side of (6.6) at the puncture has the form $`O(k^1)\psi `$.
The compatibility conditions
$`[_{\tau _j}L_j,_{\tau _k}L_k]=0.`$ (6.7)
are equivalent to non-linear partial differential equations for the coefficients of the operators $`L_k`$. For the case $`j=2,k=3,\tau _2=y,\tau _3=z`$, the operators $`L_2`$ and $`L_3`$ have the form
$$L_2=_x^2u(x,y,t),L_3=_x^3\frac{3}{2}u_x+w(x,y,t),$$
(6.8)
where we consider the dependence of the coefficients with respect to the first three variables $`\tau _1=x,\tau _2=y,\tau _3=t`$, only . The coefficient $`u(x,y,t)`$ is equal to $`u=2\xi _{1x}(x,y,t)`$, where $`\xi _1`$ is the first coefficient of the expansion (6.2).
¿From (6.7) we get a system of two equations for two coefficients $`u(x,y,t)`$ and $`w(x,y,t)`$ which can be reduced for an equation for $`u(x,y,t)`$. The reduced equation is the famous KP equation
$$3u_{yy}=(4u_t6uu_x+u_{xxx})_x.$$
(6.9)
It appeared in physics literature for the investigation of the transversal stability of the KdV solitons (see ) and is one of the most natural physical two-dimensional analogs of KdV. Lax representation for it was found in . We call the whole set of higher systems (6.7) the KP hierarchy. We shall discuss the periodic problem for this equation at greater length in the next section.
For the special choice of the hyperelliptic Riemann surface and $`k=\pm \sqrt{ϵ}`$ we have $`k^{2s}=ϵ^s`$, where the function $`ϵ`$ is well-defined globally as a meromorphic function on the Riemann surface $`\mathrm{\Gamma }`$. Therefore, we may represent globally the corresponding Baker-Akhiezer function in the form
$`\psi =\stackrel{~}{\psi }\mathrm{exp}(\tau _2k^2+\tau _4k^4+\mathrm{}),`$ (6.10)
where $`\stackrel{~}{\psi }`$ does not depend on the parameters $`\tau _{2j}`$. So in this case all the KP hierarchy reduces to the KdV hierarchy.
The relationship of this constructions with commuting OD linear operators is as follows. Let $`f(P)`$ be a meromorphic function on the Riemann surface $`\mathrm{\Gamma }`$ with one pole at the puncture $`P_0`$. Its negative part written in the parameter $`k`$ is some polynomial $`q(k)=q_1^l+q_2k^{l1}+\mathrm{}+q_lk`$. Apply the operator $`A_f=q_1L_l+q_2L_{l1}+\mathrm{}+q_1_x`$ to the function $`\psi `$. By the definition of the operators $`L_k`$, we can see that
$`(A_ff)\psi =O(k^1)\psi .`$ (6.11)
We conclude as before, that the difference $`A_f\psi f(P)\psi =0`$.
For any pair of functions $`f`$ and $`g`$ on the Riemann surface $`\mathrm{\Gamma }`$ with poles at the point $`\mathrm{}`$, we get a pair of the commuting OD scalar linear operators $`A_f,A_g`$ such that $`A_fA_g=A_gA_f`$. In a special case of the hyperelliptic Riemann surface $`z^2=(ϵϵ_0)\mathrm{}(ϵϵ_{2n})`$, we have pair of functions $`f=ϵ,g=z`$, leading to the Schrödinger operator $`L=_x^2+u`$ commuting with the second operator of order $`2n+1`$, because $`z=k^{2n+1}+\mathrm{}`$
Now we return back to the problem of existence and uniqueness of the Baker-Akhiezer function. The simplest way to prove this existence is to define this function by exact formula in terms of the $`\theta `$-function and meromorphic differentials. Let us recall first necessary information.
As in the previously discussed case of hyperelliptic curves, we introduce a basis of cycles $`a_j,b_j,j=1,\mathrm{},g,`$ on a Riemann surface $`\mathrm{\Gamma }`$ of genus $`g`$ with canonical matrix of intersections: $`a_ib_j=\delta _{ij}`$, and a basis $`\omega _i`$ of holomorphic differentials normalized by the condition
$$_{a_i}\omega _j=\delta _{ij}.$$
(6.12)
The matrix $`B=(B_{ij})`$ of $`b`$-periods of these differentials
$$B_{ij}=_{b_i}\omega _j,$$
(6.13)
is symmetric and has positively defined imaginary part. The Riemann $`\theta `$-function is a function defined with the help of this matrix by the formula:
$$\theta (z|B)=\underset{mZ^g}{}e^{2\pi i(z,m)+(Bm,m)}$$
(6.14)
where $`z=(z_1,\mathrm{},z_g)`$ is a complex $`g`$-dimensional vector, $`(m,z)`$ stands for a standard scalar product and summation is taken over all integer vectors $`m=(m_1,\mathrm{},m_g)`$. Theta-function is an entire periodic function of $`g`$ variables $`z_j`$ and has the following monodromy properties with respect to the shifts defined by vectors $`B_k`$ which are columns of the matrix of $`b`$-periods:
$$\theta (z+B_k)=e^{2\pi iz_k\pi iB_{kk}}\theta (z).$$
(6.15)
The basic vectors $`e_k`$ and the vectors $`B_k`$ define a lattice $``$ in $`C^g`$ which determines $`g`$-dimensional complex torus $`J(\mathrm{\Gamma })=C^g/`$ called Jacobian of the curve. The Abel map $`A:\mathrm{\Gamma }J(\mathrm{\Gamma })`$ is defined by the formula
$$A_k(P)=_{P_0}^P\omega _k.$$
(6.16)
Note that the vector $`A(P)`$ with coordinates $`A_k(P)`$ depends on the choice of path of integration but its ambiguity just coincides with shifts by vectors of the lattice $``$.
¿From the monodromy properties of $`\theta `$-function it follows that zeros of the multivalued function $`\theta (A(P)+Z)`$ considered as a function on $`\mathrm{\Gamma }`$ are well-defined. For a generic vector $`Z`$ this function has exactly $`g`$ zeros $`(\gamma _1,\mathrm{},\gamma _g)`$. The vector $`Z`$ can be expressed in terms of Abel transforms of these points by the formula
$$Z=\underset{j=1}{\overset{g}{}}A(\gamma _j)+K,$$
(6.17)
where $`K`$ is a vector of the Riemann constants.
Let us introduce a set of meromorphic differentials $`d\mathrm{\Omega }_i`$ that are holomorphic on $`\mathrm{\Gamma }`$ outside the puncture where they have poles of the form
$$d\mathrm{\Omega }_i=dk^i(1+O(k^{i1})),$$
(6.18)
and normalized in a usual way by the condition
$$_{a_i}𝑑\mathrm{\Omega }_j=0.$$
(6.19)
The Abelian integrals
$$\mathrm{\Omega }_i(P)=^P𝑑\mathrm{\Omega }_i$$
(6.20)
are multivalued functions on $`\mathrm{\Gamma }`$.
Let $`U_j`$ be a vector with the coordinates
$$U_{jk}=\frac{1}{2\pi i}_{b_k}𝑑\mathrm{\Omega }_j.$$
(6.21)
Then from the statements presented above, it follows directly that the formula
$$\psi (\tau ,P)=\mathrm{exp}(\underset{i}{}\tau _i\mathrm{\Omega }_i(P))\frac{\theta ((A(P)+\underset{i}{}U_i\tau _i+Z)\theta (Z)}{\theta (A(P)+Z)\theta (_iU_i\tau _i+Z)},$$
(6.22)
correctly defines a function on $`\mathrm{\Gamma }`$ which satisfies all the properties of the Baker-Akhiezer function.
Suppose now that there exists another Baker-Akhiezer function $`\psi _1`$. ¿From the definition the Baker-Akhiezer functions it follows that the ratio $`\psi _1/\psi `$ is a meromorphic function on $`\mathrm{\Gamma }`$ which is equal to 1 at the puncture and with only possible poles at the zeros of the function $`\psi `$. According to (6.22) the zeros of $`\psi `$ are zeros of the function $`\theta (A(P)+_iU_i\tau _i+Z)`$. Therefore, $`\psi `$ has $`g`$ zeros. The simplest form of the Riemann-Roch theorem (which can be considered as generalization of the Liouville theorem for Riemann surfaces) implies that a function on $`\mathrm{\Gamma }`$ with at most $`g`$ poles at a generic set of points is a constant. Therefore, $`\psi _1=\psi `$ and the existence and uniqueness of the Baker-Akhiezer function is proved.
Now according to the previously established formula $`u=2_x\xi _1`$, in order to get an exact formula for a solution of the KP hierarchy it is enough to take the first coefficient $`\xi _1(\tau )`$ of the expansion of the pre-exponential factor in (6.2) at the puncture. Finally we obtain the expression
$$u(\tau )=2_x^2\mathrm{log}\theta (\underset{i}{}U_i\tau _i+Z)+const$$
(6.23)
for the finite gap solutions of the whole KP hierarchy.
If we consider only the KP equation, we get the formula
$$u(x,y,t)=2_x^2\mathrm{log}\theta (Ux+Vy+Wt+Z)+const,$$
(6.24)
where we redenote $`x=\tau _1,y=\tau _2,t=\tau _3`$ and $`U=U_1,V=U_2,W=U_3`$.
For the case of the hyperelliptic curve the vector $`V=0`$ and we get the Its-Matveev formula for the finite-gap solutions of the KdV equation.
The formula (6.24) derived in has led to one of the most important pure mathematical applications of the theory of non-linear integrable systems. This is the solution of the famous Riemann-Shottky problem.
According to the Torrelli theorem, the matrix of $`b`$-periods of normalized holomorphic differentials uniquely defines the corresponding algebraic curve. The Riemann-Shottky problem is: to describe symmetric matrices with the positive imaginary part which are the matrices of $`b`$-periods of normalized holomorphic differentials on algebraic curves. One of the authors conjected that the function $`u(x,y,t)`$ given by (6.24) is a solution of the KP-equation iff the matrix $`B`$ that defines the theta-function is the matrix of $`b`$-periods of normalized holomorphic differentials on an algebraic curve and $`U,V,W`$ are vectors of $`b`$-periods of corresponding normalized meromorphic differentials with the only pole at a point of this curve. This conjecture was proved in .
## 7 Spectral theory of two-dimensional periodic operators. KP hierarchy
A general algebraic-geometric construction of the finite-gap potentials for the Schrödinger operators and for solutions of the KP hierarchy that was presented in the previous section has been developed extensively in years. It is applicable for all soliton systems which are equivalent to various types of compatibility conditions for over-determined systems of auxiliary linear problems. In its algebraic form it is in some sense local and is a sort of inverse transform: from a set of algebraic-geometrical data to solutions of the integrable non-linear partial differential equations
$$\{algebraicgeometricaldata\}\{solutionsofNLPDE\}$$
(7.1)
In a generic case the space of algebraic-geometrical data is a union for all $`g`$ of the spaces
$$\stackrel{~}{M}_{g,N}=\{\mathrm{\Gamma }_g,P_\alpha ,k_\alpha ^1(Q),\gamma _1,\mathrm{},\gamma _g\},\alpha =1,\mathrm{},N,$$
(7.2)
where $`\mathrm{\Gamma }_g`$ is an algebraic curve of genus $`g`$ with fixed local coordinates $`k_\alpha ^1(Q),k_\alpha ^1(P_\alpha )=0,`$ in neighborhoods of $`N`$ punctures $`P_\alpha `$, and $`\gamma _1,\mathrm{},\gamma _g`$ are points of $`\mathrm{\Gamma }_g`$ in a general position. (It is to be mentioned that $`\stackrel{~}{M}_{g,N}`$ are “universal” data. For the given non-linear integrable equation the corresponding subset of data has to be specified.)
A’posteriory it can be shown that these solutions can be expressed in terms of the corresponding Riemann theta-functions and are quasi-periodic functions of all variables. Whithin this approach it is absolutely impossible to give an answer to the basic question: ”How many algebraic-geometrical solutions are there? And what is their role in the solution of the periodic Cauchy problem for two-dimensional equations of the KP type?”
The answer of the corresponding question in lower dimensions is as follows. For finite dimensional (0+1) systems a typical Lax representation has the form
$$_tU(t,\lambda )=[U(t,\lambda ),V(t,\lambda )],$$
(7.3)
where $`U(t,\lambda )`$ and $`V(t,\lambda )`$ are matrix functions that are rational (or sometimes elliptic) functions of the spectral parameter $`\lambda `$. In that case all the general solutions are algebraic-geometrical and can be represented in terms of the Riemann theta-functions.
For spatial one-dimensional evolution equations of the KdV type (1+1)-systems) the existence of direct and inverse spectral transform allows one to prove (though it is not always the rigorous mathematical statement) that algebraic-geometrical solutions are dense in the space of all periodic (in $`x`$) solutions.
It turns out that the situation for two-dimensional integrable equations is much more complicated. For one of the real forms of the KP equation that is called the KP-2 equation and coincides within (6.9), the algebraic-geometrical solutions are dense in the space of all periodic (in $`x`$ and $`y`$) solutions (). It seems, that the same statement for the KP-1 equation which can be obtained from (6.9) be replacing $`yiy`$ is wrong. One of the most important problems in the theory of two-dimensional integrable systems which are still unsolved is: ”in what sense” the KP-1 equation ahich has the operator representation (6.7) and for which a wide class of periodic solution was constructed, is an ”non-integrable” system.
The proof of the integrability of the periodic problem for the KP-2 equation is based on the spectral Floquet theory of the parabolic operator
$$M=_y_x^2+u(x,y),$$
(7.4)
with periodic potential $`u(x+l_1,y)=u(x,y+l_2)=u(x,y).`$ We are going now to present the most essential points of this theory which was developed in . It is the natural generalization of the spectral theory of the periodic Sturm-Liouville operator. We would like to mention that despite its application to the theory of non-linear equations and related topics, the structure of the Riemann surface of Bloch solutions of the corresponding linear equation that was found in has been used as a starting point for an abstract definition of the Riemann surfaces of the infinite genus ().
Solutions $`\psi (x,y,w_1,w_2)`$ of the non-stationary Shrödinger equation
$$(\sigma _y_x^2+u(x,y))\psi (x,y,w_1,w_2)=0$$
(7.5)
with a periodic potential $`u(x,y)=u(x+a_1,y)=u(x,y+a_2)`$ are called Bloch solutions, if they are eigenfunctions of the monodromy operators, i.e.
$$\psi (x+a_1,y,w_1,w_2)=w_1\psi (x,y,w_1,w_2),$$
(7.6)
$$\psi (x,y+a_2,w_1,w_2)=w_2\psi (x,y,w_1,w_2).$$
(7.7)
The Bloch functions will always be assumed to be normalized so that $`\psi (0,0,w_1,w_2)=1`$. The set of pairs $`Q=(w_1,w_2)`$, for which there exists such a solution is called the Floquet set and will be denoted by $`\mathrm{\Gamma }`$. The multivalued functions $`p(Q)`$ and $`E(Q)`$ such that
$$w_1=e^{ipa_1},w_2=e^{iEa_2}$$
are called quasi-momentum and quasi-energy, respectively.
The gauge transformation $`\psi e^{h(y)}\psi `$, where $`_yh(y)`$ is a periodic function, transfers the solutions of (7.5) into solutions of the same equation but with another potential $`\stackrel{~}{u}=u\sigma _yh`$. Consequently, the spectral sets corresponding to the potentials $`u`$ and $`\stackrel{~}{u}`$ are isomorphic. Therefore, in what follows we restrict ourselves to the case of periodic potentials such that $`_0^{a_1}u(x,y)𝑑x=0.`$
To begin with let us consider as a basic example the “free” operator $`M_0=\sigma _y_x^2`$ with zero potential $`u(x,y)=0`$. The Floquet set of this operator is parameterized by the points of the complex plane of the variable $`k`$ $`w_1^0=e^{ika_1},w_2^0=e^{\sigma ^1k^2a_2}`$, and the Bloch solutions have the form $`\psi (x,y,k)=e^{ikx\sigma ^1k^2y}.`$ The functions $`\psi ^+(x,y,k)=e^{ikx+\sigma ^1k^2y}`$ are Bloch solutions of the formal adjoint operator $`(\sigma _y+_x^2)\psi ^+=0.`$
An image of the map $`kC(w_1^0,w_2^0)`$ is the Floquet set for the free operator $`M_0`$. It is the Riemann surface with self-intersections. The self-intersections correspond to the pairs $`kk^{}`$ such that $`w_i^0(k)=w_i^0(k^{}),i=1,2.`$ The latter conditions imply the equations
$$kk^{}=\frac{2\pi N}{a_1},k^2(k^{})^2=\frac{\sigma 2\pi iM}{a_2},$$
(7.8)
where $`N`$ and $`M`$ are integers. Hence, all the resonant points have the form
$$k=k_{N,M}=\frac{\pi N}{a_1}\frac{\sigma iMa_1}{Na_2},N0,k^{}=k_{N,M}.$$
(7.9)
The basic idea of the construction of the Riemann surface of Bloch solutions of the equation (7.5) that was proposed in is to consider (7.5) as a perturbation of the free operator, assuming that the potential $`u(x,y)`$ is formally small.
For any $`k_0k_{N,M}`$ it is easy to construct a formal Bloch solution. It turns out that the corresponding formal series converges and defines a holomorphic function of $`k_0`$ for $`|k_0|>M`$ big enough and lies outside small neighborhoods of the resonant points. Moreover, it can be shown that this function can be extended on the Riemann surface that can be thought as a surface obtained from the complex plane by some kind of surgery that creates gaps in places of the resonant points.
More precisely, if $`u(x,y)`$ is a smooth real potential that has analytical continuation in some neighborhood of the real values of $`x`$ and $`y`$, then the corresponding Riemann surface of Bloch-Floquet solutions can be described in the following way.
Let us fix some finite or infinite subset $`S`$ of integer pairs $`(N>0,M)`$ . The set of pairs of complex numbers $`\pi =\{p_{s,1},p_{s,2}\}`$ where $`sS`$ would be called “admissible”, if
$$\mathrm{Re}p_{s,i}=\frac{\pi N}{a_1},|p_{s,i}k_s|=o(|k_s|^1),i=1,2,$$
(7.10)
and the intervals $`[p_{s,1},p_{s,2}]`$ do not intersect. (Here $`k_s,s=(N,M)`$ are resonant points.)
Let us define the Riemann surface $`\mathrm{\Gamma }(\pi )`$ for any admissible set $`\pi `$. It is obtained from the complex plane of the variable $`k`$ by cutting it along the intervals $`[p_{s,1},p_{s,2}]`$ and $`[\overline{p}_{s,1},\overline{p}_{s,2}]`$ and by sewing after that the left side of the first cut with the right side of the second cut and vice versa. ( After this surgery for each cut $`[p_{s,1},p_{s,2}]`$ corresponds to a nontrivial cycle $`a_s`$ on $`\mathrm{\Gamma }(\pi )`$.)
For any real periodic potential $`u(x,y)`$ which can be analytically extended into some neighborhood of the real values $`x,y`$, the Bloch solutions of the equation (7.5) are parameterized by points $`Q`$ of the Riemann surface $`\mathrm{\Gamma }(\pi )`$ corresponding to some admissible set $`\pi `$. The function $`\psi (x,y,Q)`$ which is normalized by the condition $`\psi (0,0,Q)=1`$ is meromorphic on $`\mathrm{\Gamma }`$ and has a simple pole $`\gamma _s`$ on each cycle $`a_s`$. If the admissible set $`\pi `$ contains only a finite number of pairs, then $`\mathrm{\Gamma }(\pi )`$ has finite genus and is compactified by only one point $`P_1`$ ($`k=\mathrm{}`$), in the neighborhood of which the Bloch function $`\psi `$ has the form (6.2).
The potentials $`u`$ for which $`\mathrm{\Gamma }(\pi )`$ has finite genus are called finite-gap. They coincide with the algebraic-geometrical potentials. The direct spectral transform for periodic operators (7.4) allows us to prove that like in the one-dimensional case the finite-gap potentials are dense in the space of all periodic smooth functions in two variables ().
## 8 Spectral theory of two-dimensional Schrödinger operator for fixed energy level and two-dimensional Toda lattice
In this section we discuss a spectral theory of two-dimensional periodic Schödinger operator. Unlike the one-dimensional case, spectral data for 2D linear operator are over-determined and therefore for generic operators there are no nontrivial isospectral flows. As it was noted in , deformations that preserve spectral data for one fixed energy level do exist. An analog of Lax representation for such system has the form
$$H_t=[A,H]+BH,$$
(8.1)
where $`H,A,B`$ are two dimensional operators with coefficients depending on $`x,y,t`$. Equation (8.1) is equivalent to the condition that operators $`H`$ and $`(_tA)`$ commute on the space of solutions of the equation $`H\psi =0`$. Therefore, (8.1) describes deformations preserving all the spectral data associated with the zero energy level of the operator $`H`$.
It should be mentioned that until the moment when equations (8.1) were proposed in the framework of the soliton theory the spectral problem associated with one energy level of two-dimensional periodic operators had never been considered.
For the first time an inverse algebraic-geometric spectral problem for two-dimensional Schrödinger operator in magnetic field
$$H=(i_xA_x(x,y))^2+(i_yA_y(x,y)^2+u(x,y)$$
(8.2)
was formulated and solved in .
Consider the Bloch solutions of the equation $`H\psi =ϵ\psi `$, which by definition are solutions that at the same time are eigenfuctions for the shifts operators:
$$\psi (x+T_1,y)=e^{ip_xT_1}\psi (x,y),\psi (x,y+T_2)=e^{ip_yT_2}\psi (x,y).$$
(8.3)
Here $`T_1,T_2`$ are periods of the operator $`H`$, i.e. periods of the potential $`u(x,y)`$ and periods of the magnetic field $`B(x,y)`$, which is defined by the vector potential $`(A_x,A_y)`$ by the formula $`B=_yA_x_xA_y`$.
Multivalued quantities $`p_x,p_y`$ are components of the two-dimensional quasimomentum. For fixed values of $`p_x,p_y`$ a spectrum of the operator $`H`$ restricted on the space of functions satisfying (8.3) is discrete and defines different branches of dispersion relations $`ϵ_j(p_x,p_y),j=1,\mathrm{}`$ Level lines $`ϵ_j(p_x,p_y)=ϵ_0`$ in the space of variables $`p_x,p_y`$ define the, so-called, Fermi curves. Of course, in the solid state physics all the considerations were restricted by real values of quasi-momentum.
In it was suggested to consider operators for which a complex Fermi curve does exist and is the Riemann surface of finite genus for some energy level $`ϵ_0`$. Moreover, it was assumed that this curve is compactified by two infinity points $`P_\pm `$ in neighborhoods of which the corresponding Bloch solutions have the form
$$\psi =e^{k_\pm (x\pm iy)}\left(\underset{s=0}{\overset{\mathrm{}}{}}\xi _s^\pm (x,y)k_\pm ^s\right),$$
(8.4)
where $`k_\pm ^1)`$ are local coordinates in the neighborhoods of the punctures $`P_\pm `$. It was also assumed that outside the punctures the functions $`\psi (x,y,P)`$ considered as a function of the variable $`P`$ (which is a point of the complex Fermi curve $`\mathrm{\Gamma }`$) is meromorphic function with $`g`$ poles independent of the variables $`x`$ and $`y`$.
We present here a solution to this inverse spectral problem in a more general form which is necessary for construction of exact solutions to the two-dimensional Toda lattice which has deep connections with the theory of 2D Schrödinger operators. After that we return back to the spectral problems.
Let $`\mathrm{\Gamma }`$ be a smooth algebraic curve of genus $`g`$ with fixed local coordinates $`w_\pm (P)`$ in neighborhoods of the points $`P^\pm ,w_\pm (P^\pm )=0`$. Then for each set of $`g`$ points $`\gamma _1,\mathrm{},\gamma _g`$ in general position there exists a unique function $`\psi _n(T,P),T=\{t_i^\pm ,i=1,\mathrm{},\mathrm{},\}`$ such that:
$`1^0.`$ The function $`\psi _n`$ of the variable $`P\mathrm{\Gamma }`$ is meromorphic outside the punctures and has at most simple poles at the points $`\gamma _s`$ (if all of them are distinct);
$`2^0.`$ In a neighborhood of the point $`P^\pm `$ it has the form
$$\psi _n(T,P)=w_\pm ^n\left(\underset{s=0}{\overset{\mathrm{}}{}}\xi _s^\pm (n,T)w_\pm ^s\right)\mathrm{exp}\left(\underset{i=1}{\overset{\mathrm{}}{}}w_\pm ^it_i^\pm \right),w_\pm =w_\pm (P),$$
(8.5)
$$\xi _0^+(x,T)\delta _{\alpha j}.$$
(8.6)
The proof of this statement, as well as the explicit formula for $`\psi _n`$ in terms of Riemann theta-functions, is almost identical to the way of solution of the inverse problem for finite-gap Schrödinger operators discussed in the previous sections.
Let $`d\mathrm{\Omega }_j^{(\pm )}`$ be a unique meromorphic differential holomorphic on $`\mathrm{\Gamma }`$ outside the point $`P^\pm `$, which has the form
$$d\mathrm{\Omega }_j^{(\pm )}=d(w_\pm ^j+O(w_\pm )),$$
(8.7)
near the point $`P^\pm `$, and normalized by the conditions $`_{a_k}𝑑\mathrm{\Omega }_i^{(\pm )}=0.`$ It defines a vector $`U_j^{(\pm )}`$ with coordinates
$$U_{jk}^{(\pm )}=\frac{1}{2\pi i}_{b_k}𝑑\mathrm{\Omega }_j^{(\pm )}.$$
(8.8)
Further, let us define the normalized differential $`d\mathrm{\Omega }^{(0)}`$, which is holomorphic outside the points $`P^\pm `$ where it has simple poles with residues $`\pm 1`$, respectively. From Riemann’s bilinear relations it follows that the vector of $`b`$-periods of this differential equals $`2\pi iU^{(0)}`$ , where
$$U^{(0)}=A(P^{})A(P^+).$$
(8.9)
As in the previous case, one can directly check that the function $`\psi _n(T,P)`$ given by the formula:
$$\psi _n(T,P)=\frac{\theta (A(P)+U^{(0)}n+U_i^{(\pm )}t_i^\pm +Z)\theta (A(P^+)+Z))}{\theta (A(P)+Z)\theta (A(P^+)+U^{(0)}n+U_i^{(\pm )}t_i^\pm +Z)}e^{(n\mathrm{\Omega }^{(0)}(P)+{\scriptscriptstyle t_i^\pm \mathrm{\Omega }_i^{(\pm )}(P)})},$$
(8.10)
$$\mathrm{\Omega }_i^{(\pm )}(P)=^P𝑑\mathrm{\Omega }_i^{(\pm )},$$
(8.11)
is well-defined and has all the properties of the Baker-Akhiezer function.
Note, that for $`n=0`$ and $`t_1^\pm =x\pm iy,t_i^\pm =0,i>1,`$ the analytical properties of this function coincide with the properties that were described above as analytical properties of the Bloch solutions for finite-gap two-dimensional Schrödinger operators.
¿From the uniqueness of the Baker-Akhiezer functions $`\psi _n(T,P)`$ it follows that they satisfy the linear equations
$$_+\psi _n=\psi _{n+1}+v_n\psi _n,_{}\psi _n=c_n\psi _{n1},_\pm =\frac{}{t_1^\pm },$$
(8.12)
where
$$v_n=_+\phi _n(T),c_n=e^{\phi _n(T)\phi _{n1}(T)},e^{\phi _n}=\xi _0^{}(T),$$
(8.13)
and $`\xi _0^{}(T)`$ is a leading term of the expansion of $`\psi _n`$ at the puncture $`P^{}`$. From (8.10) we get the formula
$$\phi _n=\mathrm{log}\frac{\theta (U^{(0)}(n+1)+U_i^{(\pm )}t_i^\pm +Z_0)}{\theta (U^{(0)}n+U_i^{(\pm )}t_i^\pm +Z_0)}+const,Z_0=Z+A(P^+).$$
(8.14)
for algebraic-geometric solutions of 2D Toda lattice which was obtained in .
Note that (8.12) imply that $`\psi _0`$ satisfies the equation
$$_+_{}\psi _0+v_0_{}\psi _0+(c_1+_{}v_0)\psi _0,$$
(8.15)
which is gauge equivalent to (8.2) and therefore, we do get a solution of the inverse problem that was introduced above.
The next important step was done in where algebraic-geometric spectral data corresponding to potential $`2d`$ Schrödinger operators (i.e. operators with zero magnetic field) were found.
Let $`\mathrm{\Gamma }`$ be a smooth genus $`g`$ algebraic curve with fixed local coordinates $`k_\pm ^1(Q)`$,
$`k_\pm ^1(P_\pm )=0,`$ in neighborhoods of two punctures $`P_\pm `$. Let us assume that there exists a holomorphic involution of the curve $`\sigma :\mathrm{\Gamma }\mathrm{\Gamma }`$ such that $`P_\pm `$ are its only fixed points, i.e. $`\sigma (P_\pm )=P_\pm .`$ The local parameters are to be ”odd”, i.e. $`k_\pm (\sigma (Q))=k_\pm (Q).`$ The factor-curve will be denoted by $`\mathrm{\Gamma }_0`$. The projection
$$\pi :\mathrm{\Gamma }\mathrm{\Gamma }_0=\mathrm{\Gamma }/\sigma $$
(8.16)
represents $`\mathrm{\Gamma }`$ as a two-sheet covering of $`\mathrm{\Gamma }_0`$ with the two branch points $`P_\pm `$. In this realization the involution $`\sigma `$ is a permutation of the sheets. As there are only two branching points $`g=2g_0`$, where $`g_0`$ is genus of $`\mathrm{\Gamma }_0`$. Let us consider a meromorphic differential $`d\mathrm{\Omega }(Q)`$ of the third kind on $`\mathrm{\Gamma }_0`$ with residues $`1`$ at the points $`P_\pm `$. The differential $`d\mathrm{\Omega }`$ has $`g`$ zeros that will be denoted by $`\widehat{\gamma }_i,i=1,\mathrm{},2g_0=g.`$ Let us choose for each $`i`$ a point $`\gamma _i`$ on $`\mathrm{\Gamma }`$ such that
$$\pi (\gamma _i)=\widehat{\gamma }_i,i=1,\mathrm{},g.$$
(8.17)
In it was shown that the Baker-Akhiezer function corresponding to algebraic-geometric data which have been just defined satisfies the equation
$$(_+_{}+u(x,y))\psi (x,y,Q)=0,$$
(8.18)
where
$$u=_{}\xi _1^+=_+\xi _1^{},$$
(8.19)
and $`\xi _1^\pm =\xi _1^\pm (x_+,x_{})`$ are the first coefficients in the expansion (8.4).
It should emphasized that though general formula (8.10) in terms of the Riemann $`\theta `$-functions is valid for the Baker-Akhiezer functions corresponding to the potential Shrödinger operator $`H`$, in it was found another more effective representation in terms of the, so-called, Prim theta-function.
A space of holomorphic differentials on $`\mathrm{\Gamma }`$ splits into two $`g_0`$-dimensional subspaces of even and odd (with respect to the involution $`\sigma `$) differentials. A matrix of $`b`$-periods of odd differentials defines the function $`\theta _{Pr}(z)`$ by the same formula (6.14). Then
$$\psi =\frac{\theta _{Pr}(A^{od}(Q)+U^+x_++U^{}x_{}Z)\theta _{Pr}(A^{od}(Z))}{\theta _{Pr}(A^{od}(Q)Z)\theta _{Pr}(U^+x_++U^{}x_{}Z)}e^{i(p^+(Q)x_++p^{}(Q)x_{})}.$$
(8.20)
Here $`p^\pm (Q)`$ are Abelian integrals of the second kind normalized differential $`dp^\pm `$ on $`\mathrm{\Gamma }`$ that have poles of the second order at the points $`P_\pm `$, respectively; vectors $`2\pi U^\pm `$ are the vectors of $`b`$-periods of these differentials.
As it was mentioned above the inverse algebraic-geometric spectral problem on one energy level for 2D Schrödinger operator was posed and solved at the time when no direct spectral theory was known. This theory was developed much later in , where it was shown that the Bloch solutions for (8.18) with analytical periodic potential are parameterized by points of infinite genus Riemann surface. It was proved that if this surface has finite genus, then the Bloch functions have all the analytical properties suggested in the inverse problem and therefore are just the Baker-Akhiezer-functions. Moreover, it was proved that algebraic-geometric (finite-gap) potentials are dense in the space of all periodic potentials.
## 9 Spectral theory of operators with elliptic coefficients
In this section we are going to discuss a specific spectral problem for operators with elliptic coefficients, i.e. with the coefficients that are meromorphic functions of a variable $`x`$ and have two periods $`2\omega ,2\omega ^{},\mathrm{Im}(\omega ^{}/\omega )>0.`$
Since Hermit time, it has been known that one-dimensional Schrödinger operator with potential of the form $`n(n+1)\mathrm{}(x)`$ (where $`\mathrm{}(x)=\mathrm{}(x|\omega ,\omega ^{})`$ is Weierstrass $`\mathrm{}`$-function corresponding to elliptic curve with periods $`2\omega ,2\omega ^{}`$, and $`n`$ is an integer) has only $`n`$ gaps in the spectrum. These Lame potentials had been the only known examples with the finite-gap property before the finite-gap theory was constructed in the framework of the soliton theory (see above). As we already have shown a generic algebraic-geometric potential can be expressed in terms of higher genus Riemann $`\theta `$-function. Sometimes higher genus formula can be reduced to exact expression in terms of elliptic function.
The first example of this type which is different from the Lame potentials was proposed in . Later a theory of elliptic finite-gap potentials attracted the special interest due to remarkable observation made in on their connection with the elliptic Calogero-Moser model. The most recent burst of interest is due to the unexpected connections of these systems to Seiberg-Witten solution of $`N=2`$ supersymmetric gauge theories . It turns out that the low energy effective theory for $`SU(N)`$ model with matter in the adjoint representation (identified first in with $`SU(N)`$ Hitchin system) is isomorphic to the elliptic CM system. Using this connection quantum order parameters were found in .
The elliptic Calogero-Moser (CM) system () is a system of $`N`$ identical particles on a line interacting with each other via the potential $`V(x)=\mathrm{}(x)`$. Its equations of motion have the form
$$\ddot{x}_i=4\underset{ji}{}\mathrm{}^{}(x_ix_j).$$
(9.1)
The CM system is a completely integrable Hamiltonian system, i.e. it has $`N`$ independent integrals $`H_k`$ in involution (, ). The second integral $`H_2`$ is the Hamiltonian of (9.1).
In it was shown that the elliptic solutions of the KdV equations have the form
$$u(x,t)=2\underset{i=1}{\overset{N}{}}\mathrm{}(xx_i(t))$$
(9.2)
and the poles $`x_i(t)`$ of the solutions satisfy the constraint $`_{ji}\mathrm{}^{}(x_ix_j)=0`$, which is the locus of the stationary points of the CM system. Moreover, it turns out that the dependence of the poles with respect to $`t`$ coincides with the Hamiltonian flow corresponding to the third integral $`H_3`$ of the system. In it was found that this connection becomes an isomorphism in the case of the elliptic solutions of the Kadomtsev-Petviashvilii equation. Moreover, in it was revealed that the connection of the CM systems with the KP equation is in some sense secondary and is a corollary of more fundamental connections with spectral theory of linear operators with elliptic potentials. The corresponding approach has been developed extensively in .
Let $``$ be a linear differential or difference operator in two variables $`x,t`$ with coefficients which are scalar or matrix elliptic functions of the variable $`x`$. We do not assume any special dependence of the coefficients with respect to the second variable. Then it is natural to introduce a notion of double-Bloch solutions of the equation
$$\psi =0.$$
(9.3)
We call a meromorphic vector-function $`f(x)`$, which satisfies the following monodromy properties:
$$f(x+2\omega _\alpha )=B_\alpha f(x),\alpha =1,2,$$
(9.4)
a double-Bloch function. The complex numbers $`B_\alpha `$ are called Bloch multipliers. (In other words, $`f`$ is a meromorphic section of a vector bundle over the elliptic curve.)
In the most general form a problem that we are going to address is to classify and to construct all the operators $``$ such that equation (9.3) has sufficiently enough double-Bloch solutions.
It turns out that existence of the double-Bloch solutions is so restrictive that only in exceptional cases such solutions do exist. A simple and general explanation of that is due to the Riemann-Roch theorem. Let $`D`$ be a set of points $`x_i,i=1,\mathrm{},m,`$ on the elliptic curve $`\mathrm{\Gamma }_0`$ with multiplicities $`d_i`$ and let $`V=V(D;B_1,B_2)`$ be a linear space of the double-Bloch functions with the Bloch multipliers $`B_\alpha `$ that have poles at $`x_i`$ of order less or equal to $`d_i`$ and holomorphic outside $`D`$. Then the dimension of $`D`$ is equal to:
$$\mathrm{dim}D=\mathrm{deg}D=\underset{i}{}d_i.$$
Now let $`x_i`$ depend on the variable $`t`$. Then for $`fD(t)`$ the function $`f`$ is a double-Bloch function with the same Bloch multipliers but in general with higher orders of poles because taking derivatives and multiplication by the elliptic coefficients increase these orders. Therefore, the operator $``$ defines a linear operator
$$|_D:V(D(t);B_1,B_2)V(D^{}(t);B_1,B_2),N^{}=\mathrm{deg}D^{}>N=\mathrm{deg}D,$$
and (9.3) is always equivalent to an over-determined linear system of $`N^{}`$ equations for $`N`$ unknown variables which are the coefficients $`c_i=c_i(t)`$ of expansion of $`\mathrm{\Psi }V(t)`$ with respect to a basis of functions $`f_i(t)V(t)`$. With some exaggeration one may say that in the soliton theory the representation of a system in the form of the compatibility condition of an over-determined system of the linear problems is considered equivalent to integrability.
In all of the basic examples $`N^{}=2N`$ and the over-determined system of equations has the form
$$LC=kC,_tC=MC,$$
(9.5)
where $`L`$ and $`M`$ are $`N\times N`$ matrix functions depending on a point $`z`$ of the elliptic curve as on a parameter. A compatibility condition of (9.5) has the standard Lax form $`_tL=[M,L]`$, and is equivalent to a finite-dimensional integrable system.
The basis in the space of the double-Bloch functions can be written in terms of the fundamental function $`\mathrm{\Phi }(x,z)`$ defined by the formula
$$\mathrm{\Phi }(x,z)=\frac{\sigma (zx)}{\sigma (z)\sigma (x)}e^{\zeta (z)x}.$$
(9.6)
Note, that $`\mathrm{\Phi }(x,z)`$ is a solution of the Lame equation:
$$\left(\frac{d^2}{dx^2}2\mathrm{}(x)\right)\mathrm{\Phi }(x,z)=\mathrm{}(z)\mathrm{\Phi }(x,z).$$
(9.7)
¿From the monodromy properties it follows that $`\mathrm{\Phi }`$ considered as a function of $`z`$ is double-periodic:
$$\mathrm{\Phi }(x,z+2\omega _\alpha )=\mathrm{\Phi }(x,z),$$
though it is not elliptic in the classical sense due to an essential singularity at $`z=0`$ for $`x0`$.
As a function of $`x`$ the function $`\mathrm{\Phi }(x,z)`$ is double-Bloch function, i.e.
$$\mathrm{\Phi }(x+2\omega _\alpha ,z)=T_\alpha (z)\mathrm{\Phi }(x,z),T_\alpha (z)=\mathrm{exp}\left(2\omega _\alpha \zeta (z)2\zeta (\omega _\alpha )z\right).$$
In the fundamental domain of the lattice defined by $`2\omega _\alpha `$ the function $`\mathrm{\Phi }(x,z)`$ has a unique pole at the point $`x=0`$:
$$\mathrm{\Phi }(x,z)=x^1+O(x).$$
(9.8)
The gauge transformation $`f(x)\stackrel{~}{f}(x)=f(x)e^{ax},`$ where $`a`$ is an arbitrary constant does not change poles of any function and transforms a double Bloch-function into another double-Bloch function. If $`B_\alpha `$ are Bloch multipliers for $`f`$ then Bloch multipliers for $`\stackrel{~}{f}`$ are equal to
$$\stackrel{~}{B}_1=B_1e^{2a\omega _1},\stackrel{~}{B}_2=B_2e^{2a\omega _2}.$$
(9.9)
The two pairs of Bloch multipliers that are connected with each other through the relation (9.9) for some $`a`$ are called equivalent. Note that for all equivalent pairs of Bloch multipliers the product $`B_1^{\omega _2}B_2^{\omega _1}`$ is a constant depending on the equivalence class, only.
¿From (9.8) it follows that a double-Bloch function $`f(x)`$ with simple poles $`x_i`$ in the fundamental domain and with Bloch multipliers $`B_\alpha `$ (such that at least one of them is not equal to $`1`$) may be represented in the form:
$$f(x)=\underset{i=1}{\overset{N}{}}c_i\mathrm{\Phi }(xx_i,z)e^{kx},$$
(9.10)
where $`c_i`$ is a residue of $`f`$ at $`x_i`$ and $`z`$, $`k`$ are parameters related by $`B_\alpha =T_\alpha (z)e^{2\omega _\alpha k}.`$ (Any pair of Bloch multipliers may be represented in this form with an appropriate choice of the parameters $`z`$ and $`k`$.)
Let us consider as an example the equation
$$\psi =(_t_x^2+u(x,t))\psi =0,$$
(9.11)
where $`u(x,t)`$ is an elliptic function. Then as shown in equation (9.11) has $`N`$ linear independent double-Bloch solutions with equivalent Bloch multipliers and $`N`$ simple poles at points $`x_i(t)`$ if and only if $`u(x,t)`$ has the form
$$u(x,t)=2\underset{i=1}{\overset{N}{}}\mathrm{}(xx_i(t))$$
(9.12)
and $`x_i(t)`$ satisfy the equations of motion of the elliptic CM system (9.1).
The assumption that there exist $`N`$ linear independent double-Bloch solutions with equivalent Bloch multipliers implies that they can be written in the form
$$\psi =\underset{i=1}{\overset{N}{}}c_i(t,k,z)\mathrm{\Phi }(xx_i(t),z)e^{kx+k^2t},$$
(9.13)
with the same $`z`$ but different values of the parameter $`k`$.
Let us substitute (9.13) into (9.11). Then (9.11) is satisfied if and only if we get a function holomorphic in the fundamental domain. First of all, we conclude that $`u`$ has poles at $`x_i`$, only. The vanishing of the triple poles $`(xx_i)^3`$ implies that $`u(x,t)`$ has the form (9.12). The vanishing of the double poles $`(xx_i)^2`$ gives the equalities that can be written as a matrix equation for the vector $`C=(c_i)`$:
$$(L(t,z)kI)C=0,$$
(9.14)
where $`I`$ is the unit matrix and the Lax matrix $`L(t,z)`$ is defined as follows:
$$L_{ij}(t,z)=\frac{1}{2}\delta _{ij}\dot{x}_i(1\delta _{ij})\mathrm{\Phi }(x_ix_j,z).$$
(9.15)
Finally, the vanishing of the simple poles gives the equations
$$(_tM(t,z))C=0,$$
(9.16)
where
$$M_{ij}=\left(\mathrm{}(z)2\underset{ji}{}\mathrm{}(x_ix_j)\right)\delta _{ij}2(1\delta _{ij})\mathrm{\Phi }^{}(x_ix_j,z).$$
(9.17)
The existence of $`N`$ linear independent solutions for (9.11) with equivalent Bloch multipliers implies that (9.14) and (9.16) have $`N`$ independent solutions corresponding to different values of $`k`$. Hence, as a compatibility condition we get the Lax equation $`\dot{L}=[M,L]`$ which is equivalent to (9.1). Note that the last system does not depend on $`z`$. Therefore, if (9.14) and (9.16) are compatible for some $`z`$, then they are compatible for all $`z`$. As a result we conclude that if (9.11) has $`N`$ linear independent double-Bloch solutions with equivalent Bloch multipliers then it has infinitely many of them. All the double-Bloch solutions are parameterized by points of an algebraic curve $`\mathrm{\Gamma }`$ defined by the characteristic equation
$$R(k,z)det(kIL(z))=k^N+\underset{i=1}{\overset{N}{}}r_i(z)k^{Ni}=0.$$
(9.18)
Equation (9.18) can be seen as a dispersion relation between two Bloch multipliers and defines $`\mathrm{\Gamma }`$ as $`N`$-sheet cover of $`\mathrm{\Gamma }_0`$.
As was shown in expansion of the characteristic equation (9.14) at $`z=0`$ has the form:
$$R(k,z)=\underset{i=1}{\overset{N}{}}(k+\nu _iz^1+h_i+O(z)),\nu _1=1N,\nu _i=1,i>1.$$
(9.19)
We call the sheet of $`\mathrm{\Gamma }`$ at $`z=0`$ corresponding to the branch $`k=z^1(N1)+O(1)`$, an upper sheet and mark the point $`P_1`$ on this sheet among the preimages of $`z=0`$. From (9.19) it follows that in general position when the curve $`\mathrm{\Gamma }`$ is smooth, its genus equals $`N`$.
Further consideration of analytical properties of the function $`\psi `$ given by (9.13) where $`c_i`$ are components of the eigenvector to the matrix $`L`$ shows that this function is just the Baker-Akhiezer function introduced in Section 6. Combined with expression (6.22) for $`\psi `$ in terms of $`\theta `$-function this result directly leads to the main statement of :
the coordinates of the particles $`x_i(t)`$ are roots of the equation
$$\theta (Ux+Vt+Z_0)=0,$$
(9.20)
where $`\theta (\xi )=\theta (\xi |B)`$ is the Riemann theta-function corresponding to the matrix of $`b`$-periods of holomorphic differentials on $`\mathrm{\Gamma }`$; the vectors $`U`$ and $`V`$ are the vectors of $`b`$-periods of normalized meromorphic differentials on $`\mathrm{\Gamma }`$, with poles of order 2 and 3 at the point $`P_1`$.
Among other examples of integrable systems that can be generated in the similar way are Ruijesenaars-Schneider system
$$\ddot{x}_i=\underset{si}{}\dot{x}_i\dot{x}_s(V(x_ix_s)V(x_sx_i)),V(x)=\zeta (x)\zeta (x+\eta ),$$
(9.21)
and nested Bethe ansatz equations
$$\underset{ji}{}\frac{\sigma (x_i^nx_j^{n+1})\sigma (x_i^n\eta x_j^n)\sigma (x_i^nx_j^{n1}+\eta )}{\sigma (x_i^nx_j^{n1})\sigma (x_i^n+\eta x_j^n)\sigma (x_i^nx_j^{n+1}\eta )}=1$$
(9.22)
As shown in , they are generated by spectral problems for equations
$$\psi =_t\psi (x,t)\psi (x+\eta ,t)v(x,t)\psi (x,t)=0,$$
(9.23)
and
$$\psi (x,m+1)=\psi (x+\eta )+v(x,m)\psi (x,m),$$
(9.24)
respectively. (Here $`\eta `$ is a complex number and $`v(x,t)`$ is an elliptic function.)
Strange as it is, the inverse spectral problem which is discussed here is simplier for two-dimensional operators then for one-dimensional stationary operators. For example a family of spectral curves corresponding to operators (9.11) that have double-Bloch solutions can be described explicitly. A nice formula was found in :
$$R(k,z)=f(k\zeta (z),z),f(k,z)=\frac{1}{\sigma (z)}\sigma \left(z+\frac{}{k}\right)H(k),$$
(9.25)
where H(k) is a polynomial. Note that (9.25) may be written as:
$$f(k,z)=\frac{1}{\sigma (z)}\underset{n=1}{\overset{N}{}}\frac{1}{n!}_z^n\sigma (z)\left(\frac{}{k}\right)^nH(k).$$
The coefficients of the polynomial $`H(k)`$ are free parameters of the spectral curve of the CM system.
The spectral curves corresponding to the Schrödinger operator with the same property are special case of the curves (9.25) but their explicit description is unknown. In particular, exact formula for branching points for Lame potentials is unknown. As it was mentioned above the first example of elliptic finite-gap potentials different from the Lame potentials was found in . A wide class of such potentials was found in .
In it was noted that the problem of classification of Schrödinger operators with elliptic potentials that have two double-Bloch solutions for almost all energy levels was posed by Picard though had not been solved until very recently. In using Floquet spectral theory for the Schrödinger operator it was proved that all such potentials are finite-gap. This result is an essential step in the Picard problem though its complete and effective solution is still an open problem.
|
warning/0003/physics0003052.html
|
ar5iv
|
text
|
# NATURE OF SONOLUMINESCENCE: NOBLE GAS RADIATION EXCITED BY HOT ELECTRONS IN ”COLD” WATER
## I Introduction
Sonoluminescence refers to the phenomenon of light emission during acoustic radiation of a liquid and is associated with cavitation bubbles present in the liquid. The most controllable and promising experimental data were obtained for single bubble sonoluminescence (SBSL): picosecond UV radiation of a single bubble pulsating in the field of the sound wave .
Though the case of a single pulsating bubble is, of course, much simpler than cavitation in general, it turns out that the SBSL is a very complex phenomenon which still remains not completely understood. There is a vast literature devoted to SBSL that was recently reviewed in Ref.. Interest in this phenomenon is stimulated, on the one hand, by the fact that in relatively simple and controllable experiments extraordinary conditions (ultra-high pressures, temperatures, and ultra-short light flashes) are realized at the final stage of the bubble collapse; on the other hand, it looks promising to use the SBSL to construct a source of ultraviolet ultra-short flashes that is much cheaper than lasers.
The highly involved hydrodynamics of bubble collapse has been addressed in many papers that are reviewed in Ref. . One of the most important problems to explain here is the very existence of stable pulsation of bubbles. Recently a theory has been developed which explains this regime in terms of the dissolved gas diffusion and chemical reactions in the gas within the bubble. It has been shown that an accumulation of noble gas in the bubble takes place during pulsation of the bubble and that in a stable situation the gas in the bubble consists, almost entirely, of noble gas and, of course, of water vapor. These theoretical conclusions were supported by the experimental data of Ref..
Another key problem in SBSL is the mechanism of the light emission. Many mechanisms have been proposed, criticized, and reviewed to explain UV flashes radiated by the collapsing bubble . The most popular of the suggested mechanisms are the adiabatic heating of the bubble gas , shock wave-Bremsstrahlung model , and, recently proposed proton-tunneling radiation as a result of a phase transition in water at ultra-high pressure . A general feature of these mechanisms is that extraordinary conditions are needed which can only be realized, if at all, at a very small bubble radius when the density of the bubble gas is close to the water density and both the hydrodynamics of the bubble and the properties of the bubble content and the neighboring liquid are not known from experiment but inferred from numerical simulations.
In this paper we show that the main characteristic features of SBSL can be explained even without making any assumptions about the extraordinary conditions. We are far from saying that these extraordinary conditions are not present in experiments. What we mean is that there is another mechanism of the SBSL which occurs in water near the bubble surface but not in the bubble gas as has been assumed in the most popular models of SBSL . Extraordinary conditions are not necessary for the operation of this mechanism even though it may well be the main cause of UV radiation. The mechanism under consideration is based on the idea put forward in Ref. that SBSL occurs because of electric breakdown in strong electric fields arising near the bubble surface as a result of flexoelectric effect, that is the effect of polarization of water because of gradients of pressure . Here we present a scenario and estimations that show that within the hypothesis the main features of SBSL can be explained using relatively moderate parameters, e.g., temperatures (5$`÷`$10)10<sup>3</sup>K for the bubble gas , and the natural (expected) value of the flexoelectric coefficient of water .
Let us mention that Ref. has left many questions to answer. The origin of the optical radiation of the bubble remaines unexplained. The problem is that in the visible region pure water has very low levels of absorbtion and radiation due to interband optical transition . Moreover the breakdown scenario is far from being clear. The mechanism of the breakdown in water (see, e.g.,) involves ”lucky electrons” whose acceleration in the electric field leads to development of an avalanche. At ambient temperature the concentration of free electrons in water is quite negligible (the band gap is about 6.5 eV ) and it is impossible to find a ”lucky electron” in a small volume near the bubble surface, i.e. in the region of high electric field, during the short time of the existence of this field. The spectacular synchronization of the emission pulses that in the case of breakdown appears, at least at first sight, to be hardly compatible with the fact that we are dealing with a probabilistic situation. In addition, the reference to the Penning effect to explain the role of the noble gases is not convincing because the ionization energy of the metastable state is more than the band gap of water as distinct from the case of breakdown in gases where the Penning effect is pronounced.
In this paper we consider in more detail possible processes associated with the electric field arising from the flexoelectric effect near a bubble exhibiting SBSL. A common difficulty with theories of SBSL is that little is known about the parameters of gas in the bubble when the bubble radius is near its minimum value. Consensus exists about one point: it is far from being a gas, because the lowest limit of the radius is governed by the van der Waals repulsion and the minimum volume of the bubble is close to the van der Waals hard core volume. The equations of state used for these conditions are not reliable at present. This is a challenging and fundamental problem but it is beyond the scope of this paper. What is now possible is to make order-of-magnitude estimations. That is why we shall first explain what values of relevant parameters are necessary for our scenario to be operative and then discuss whether our assumptions about these parameters are realistic. We include the value of the flexoelectric coefficient of water among these parameters as well. This coefficient has not been measured, unfortunately, and we use its estimated ”natural” value but, from the other hand, the knowledge of its precise value would hardly be of decisive importance because the bubble gas parameters are not known precisely.
Let us describe shortly the proposed scenario. At certain short time interval, $`\tau _c1`$ns, when the bubble radius is near its minimum, the acceleration of the radius and, therefore, the pressure gradient, assume gigantic values and the sign of the pressure gradient is such as to create a strong (depolarizing) electric field directed to the centre of the bubble in a thin water layer, $`1\mu m`$, near the bubble surface, because of the flexoelctric effect. What happens then is, in effect, screening of this depolarizing field by free electrons. Indeed, during the same time interval the temperature of the gas sharply increases up to at least several thousand K, which owing to thermal excitation of the bubble gas and water in a thin layer near the bubble surface makes the concentration of free electrons appreciable. The free electrons are accelerated by the strong ”flexoelectric” field up to energies sufficient to generate additional free electrons as a result of the electric breakdown of water. The hot electrons also collide with noble gas atoms dissolved in water and transfer them to higher energy excited states. Optical transitions from these states produce light radiation with a broad spectrum whose shape is determined mainly by the energy distribution of the hot electrons. The latter has the form of a Maxwell distribution with an effective electron temperature which can be very high. This is due to the fact that the noble gas atoms have a huge number of excited states with very high probabilities of optical transitions between them (see, e.g., ). Since the radiation occurs in the region of very strong and inhomogeneous electric fields the observed radiation spectrum is featureless. The characteristic time of the electric breakdown is much less than the characteristic time for the last stage of the collapse $`\tau _c`$, i.e. the polarization can continue to change after the first breakdown and the electric field can reach the breakdown threshold value more than once. As a result, several breakdowns can take place during the time interval $`\tau _c`$. During each breakdown the noble gas atoms are excited. An important specific feature of the noble gas atoms is the existence of long-living (metastable) excited states with a life time of up to several milliseconds (see, e.g., ). Therefore, once excited the noble gas atom can remain in the metastable state for the entire time interval $`\tau _c1`$ns, and, possibly, for many periods of the acoustic wave (the period is about $`30\mu `$s). That means that in the stationary state the longer is the pulse the more is its intensity.
The paper is organized as follows. In Sec.2 we discuss the flexoelectric effect in water in more detail than in and consider the sign of the flexoelectric field, which is crucial for the proposed scenario, at different stages of the bubble collapse. The scenario of SBSL is presented in Sec.3 where the values of relevant parameters necessary for realizing the scenario are estimated. In Sec.4 we discuss the energy distribution functions of hot electrons and the breakdown conditions in strong flexoelectric fields as well as the spectrum of SBSL. In Conclusions we summarize the results of our theory, discuss their relation to the experimental data, speculate about some possible modifications of the scenario and point out some problems to solve.
## II Flexoelectric effect and electric field near the bubble
We have already mentioned that the proposed mechanism of SBSL is based on the flexoelectric effect, namely the appearance of a polarization (electric field) due to gradient of density (pressure). This effect is not widely known and it was discussed in but shortly. Besides, the flexoelectric coeficients are still not determined experimentally. That is why we think it is worthwhile considering in more detail the appearance of electric field due to flexoelectric effect .
The flexoelectric effect is a particular case of a more general phenomenon: appearance of electrical polarization, $`𝐏,`$as a result of gradient of some scalar quantity, e.g., temperature, concentration of a component, mass density, $`\rho `$. This effect should take place in any substance . As usual, the material equation can be written in several forms. In particular, considering polarization as a function of dilatation, $`u=\mathrm{\Delta }\rho /\rho `$, and the electric field, $`𝐄`$, one has for an isotropic medium
$$𝐏=\alpha u+\epsilon _0\chi 𝐄,$$
(1)
where $`\chi `$ is the electric susceptibility and $`\alpha `$ is one of the flexoelectric coefficients. We shall be interested in the case of spherical symmetry and the absence of free charges. In this case $`𝐃=\epsilon _0𝐄+𝐏=0.`$ Then from Eq.1 it follows that
$$𝐄=\frac{\alpha }{\epsilon \epsilon _0}u=\frac{\alpha \beta }{\epsilon \epsilon _0}pfp,$$
(2)
since $`u=\beta p`$, where $`p`$ is the (excess) pressure, $`\beta `$ is the compressibility.
To estimate the flexoelectric coefficients it is convenient to consider separately two main mechanisms of polarization: charge displacement and dipole ordering. We shall see that in both cases the coefficient of proportionality between polarization and gradient of pressure, $`f`$, has the same characteristic value.
The charge displacement polarization is realized, e.g., in ionic crystals. When a gradient of dilatation takes place there are less and more compressed regions in each unit cell. The ions of larger radius tend to displace to the less compressed region while ions of less radius displace into the opposite direction. Since cations and anions usually have different radii their displacements produce polarization. The coefficient $`\alpha `$ can be estimated as follows . The maximum (”atomic”) gradient of dilatation is equal to $`1/d,`$ where $`d`$ is the interatomic distance.This gradient has to produce ”atomic” polarization $`P_{at}e/d^2`$, where $`e`$ is the electron charge, when the electric field is absent (compensated by some charges), i.e. $`\alpha _{ion}e/d`$ . Then from Eq.2 it follows that
$$f\frac{\alpha _{ion}\beta }{\epsilon \epsilon _0}\frac{e\beta }{d\epsilon _0}$$
(3)
where we have taken into account that $`\epsilon 1`$ for ionic crystals.
We shall argue that the same estimation is aplicable for materials with dipole ordering. The dipole ordering polarization in the absence of an external electric field but under a gradient of pressure arises because of geometric asymmetry of the molecular dipoles. For example, the geometric shape of water molecule is highly asymmetric: the negatively charged end is much more compact than the positive one: the negative charge is concentrated in the oxygen ion (O) while the positive charge is shared by two hydrogen ions (H<sup>+</sup>) located fairly far from each other. Under a pressure gradient a water molecule tries to orient itself in such a way that the oxygen ion would be located in the region of higher density while the hydrogens would be located in the region of lesser density. In other words, a gradient of pressure leads to orientation of the water molecular dipoles. It is well known that polarization due to dipole ordering is much more effective than polarization due to charge displacements. This is reflected in the fact that dielectric constants of dipolar materials are, normally, much larger than of non-polar materials. For the same reason it is natural to estimate $`\alpha _{dip}\epsilon \alpha _{ion}`$. Then from Eq.2 it follows that for water
$$f\frac{\alpha _{dip}\beta }{\epsilon \epsilon _0}\frac{e\beta }{d\epsilon _0}10^7\frac{Vm^2}{N},$$
(4)
where we have taken into accout that for water $`\beta 510^{10}\frac{m^2}{N},`$ $`d10^{10}m`$.
For what follows the sign of the flexoelectric coefficients is of importance. As it has been mentioned above the negative tip of the water molecule, the oxygen ion, tends to be located in the region of higher pressure. This means that the polarization vector is directed opposite to the gradient of pressure, i.e., the flexoelectric coefficient $`\alpha `$ is negative and the coefficient $`f`$ is positive.
Let us mention that we rather underestimated the flexoelectric coefficient of water than overestimated. Indeed, what we call $`\alpha _{ion}`$ is replaced, in fact, by $`\alpha _{at}`$ defining an ”atomic” flexoelectric coefficient that is of the same order of magnitude for all substances, i.e. it does not take into account specific features of the substance in question. It is natural to expect that a substance consisting of molecules whose electric asymmetry (existence of a dipolar moment) is accompanied by a pronounced geometric asymmetry (as it is for water) will exhibit a stronger flexoelectric effect than that estimated above. However, while this coefficient is not measured (reliable calculations seem much more problematic than measurements) we shall assume the value given by Eq.4. On the contrary, the conclusion reached about the sign of the coefficient $`f`$, which is crucial for the mechanism discussed in this paper, seems much more definite.
Above we have neglected the conductivity of water. This may seem questionable because the Debye radius of the electric field screening in water is comparable with the characteristic size of the strong field region ($`1\mu `$m). However, this neglection is justified because, as a rule, the dielectric relaxation time, $`\tau _D,`$ is much greater than $`\tau _c1`$ns, which is the maximal characteristic time of the polarization change in our case. Increasing the ionic conductivity of water by adding, for example, NaCl, one can, according to our estimations, decrease $`\tau _D`$ down to 0.1ns. In such an electrolyte the mechanism of SBSL discussed in this paper might be less effective.
## III The scenario: estimation of the main parameters
The dynamics of bubble cavitation have been studied in many papers . Here we will only discuss the short time interval when the bubble radius $`R`$ is close to its minimum value $`R_c`$ (Fig.1a). Fig.1a reflects the most essential features of Fig.4 of Ref. and Fig.12 of Ref. where experimental results were presented. Within this interval the velocity of the bubble surface $`v=dR/dt`$ reaches its maximum and reaches zero at the point $`R=R_c`$ (Fig.1b) and the acceleration $`a=d^2R/dt^2`$ reverses its sign and can achieve huge values (Fig.1c). During the negative acceleration period ($`tt_1`$, see Fig. 1c) the gradient of pressure $`p=\rho a`$ is directed from the bubble center to its periphery, thus according to Eq.4 the flexoelectric depolarizing field has the same direction (Fig.1d). That means that free electrons that could be generated in the gas or in water near the bubble surface cannot be accelerated. During the positive acceleration period ($`t_1<t<t_2`$, see Fig. 1e) the situation is opposite: the arising flexoelectric field is directed to the bubble center and the generated electrons can be accelerated by the field and produce a breakdown (Fig. 1e). According to experimental data for a bright SBSL the characteristic value of the acceleration is about (10$`{}_{}{}^{11}÷`$ $`10^{13})`$ m/s<sup>2</sup> during an interval, $`\tau _c=t_2t_1`$ 1ns , so the pressure gradient $`p=\rho a`$ can reach (10$`{}_{}{}^{14}÷`$10$`{}_{}{}^{16})`$N/m<sup>3</sup>. Note that the gradient of pressure decays with distance from the bubble surface . For purpose of estimations one can consider water as incompressible liquid where $`p=(p_s)R^2/r^2`$ where $`p_s`$ is the gradient of pressure at the bubble surface. The same estimation for the gradient of pressure can be obtained if we take into account that the pressure is about (10$`{}_{}{}^{6}÷`$10<sup>8</sup>) N/m<sup>2</sup> at the final stage of the bubble collapse and the extension of the high pressure region is about 1$`\mu `$m .
It follows from Eqs.4 and 2 that for $`p`$=(10$`{}_{}{}^{14}÷`$10$`{}_{}{}^{16})`$N/m<sup>3</sup> the value of the flexoelectric field, $`E_{\text{ ,}}`$ can reach (10$`{}_{}{}^{7}÷`$ 10$`{}_{}{}^{9})`$V/m. According to Eq.4 the electric field decays just as the gradient of pressure, i.e., the flexoelectric field in water is given by the formula $`EE_sR^2/r^2`$ where $`E_s`$ is the field at the bubble surface, i.e. the radial extension of the strong field region is about 1$`\mu `$m. The field $`E_\text{ }`$ can essentially exceed the threshold field of dynamic electric breakdown of water $`E_{th}`$ which is about (1$`÷`$3)$``$10<sup>8</sup>V/m. It does not mean, however, that electric breakdown will occur: a ”lucky electron” which capable of provoking an avalanche is needed. A similar situation takes place with the laser breakdown . At the same time the conduction electron concentration, $`n,`$ in water at room temperature is practically zero ($`n10^{90}`$cm<sup>-3</sup>): water could be considered as a wide gap amorphous semiconductor with $`E_g6.5`$eV .
The breakdown starts only the moment when the strong flexoelectric field is directed to the bubble center (Fig. 1f) and conduction electrons appear because of the sharp increase of temperature near the bubble surface . The latter takes place when the bubble radius $`RR_s`$ which is close to $`R_c`$(0.5$`÷`$1)$`\mu `$. At this moment the flexoelectric field $`E_f`$ may be much more greater than $`E_{th}`$ so the coefficient of avalanche multiplication of an electron can be practically infinite and in this case just several electrons are enough to provoke the breakdown and to screen the field.
Let us show that even relatively moderate temperatures near the bubble surface (which seem to be generally acceptable for the regimes without the shock waves) are quite sufficient to provide a sufficient number of conduction electrons to ensure the breakdown. As in Ref. we assume that the gas has the temperature of 7000K. That means that the water layer with a thickness of the thermal penetration length, $`\delta _T`$, has a temperature of about 3000K. For $`\delta _T`$ one has:
$$\delta _T=(\frac{2\kappa }{\rho c}t)^{1/2}$$
(5)
where $`\kappa `$, the thermal conductivity of water is about 0.4 J/m s K; $`c`$, the specific heat, is about 4$`10^3`$ J/m<sup>3</sup> K; $`\rho `$, the density, is 10<sup>3</sup>kg/m<sup>3</sup>. Bearing in mind that the characteristic time for the last stage of the collapse is $`\tau _c`$1ns , one finds that the thermal penetration length is $`\delta _T0.1\mu `$m. Following we will consider water as an amorphous semiconductor with a band gap $`E_g6.5`$eV and an effective density of states $`N^{}`$ $`10^{21}`$ for $`T3000`$K. The equilibrium concentration of conduction electrons in water at $`T3000`$K is $`n=N^{}\mathrm{exp}(E_g/2kT)310^{15}`$cm<sup>-3</sup>. Taking into account that the volume of the hot surface layer is approximately $`V10^{12}÷10^{13}`$cm<sup>3</sup> we find that the equilibrium number of conduction electrons in the layer is $`N>300÷3000.`$ It is important also to find the time needed for the equilibrium concentration of electrons to be established. The thermal ionization frequency (see, e.g., ) is
$$\nu _T=(N^{}\sigma v_T)\mathrm{exp}(E_g/2kT),$$
(6)
where $`\sigma `$ is the cross-section of the free carrier recombination, $`v_T`$ is the thermal velocity of conduction electrons in water. Taking into account that a reasonable value of $`\sigma `$ is about 10<sup>-15</sup>cm<sup>-2</sup> and $`v_T10^7`$cm/s we find from Eq.6 that for $`T>3000`$K the transient time to the equilibrium $`\tau _T=\nu _T^110^8`$s. From that follows that during the characteristic time for the last stage of the bubble collapse $`\tau _c`$1ns the number of conduction electrons in the layer, $`N,`$ exceeds $`30÷300`$ electrons, which is quite enough to provoke breakdown.
Note that a similar number of free electrons can be generated by thermal ionization of the gas. Indeed, the assumed gas temperature (7000K) is about two times larger and the ionization energy both of the water vapor and Ar, Kr, Xe is about 12$`÷`$16eV, i.e. also about two times larger than $`E_g`$. Authors of many works (see and the references therein) state that much higher temperatures can be reached because of shock waves forming in the bubble gas during the collapse. In principle, one can imagine a situation where high temperatures are not reached (and free electrons are not generated) before the shock waves are formed. In this case the shock wave, when it explodes, will start the electric breakdown.
As a result of breakdown the depolarizing flexoelectric field becomes screened. The total transmitted charge in the process of the screening is
$$Q_t=PS=4\pi R_s^2\epsilon _0E_s$$
(7)
where we have taken into account that in our case $`D=P+\epsilon _0E=0`$, i.e., $`P=\epsilon _0E`$. Assuming $`R_s1\mu `$m and $`E_s=(10^8÷10^9)`$ V/m we find that the maximum total number of transmitted electrons is $`N_t=Q_t/e(10^5÷10^6)`$.
Note that for considered electric fields the time of the screening (breakdown), $`\tau _s`$, is determined by the time it takes for the conduction electrons to cross the region of the strong electric field whose size is $`lR_s`$. Thus $`\tau _s=l/v_d`$ where $`v_d`$ is the electron drift velocity which in strong electric fields saturates at some value $`v_d10^5`$m/s (see Sec.4) . Assuming $`l1\mu `$m we find $`\tau _s10`$ps. Note that the value of $`\tau _s`$ is less than the observed SBSL pulse width .
The whole period of positive acceleration, $`\tau _c`$, (see Sec.3) is about two orders of magnitude larger than $`\tau _b.`$ Therefore, after the breakdown is finished and the depolarization field is screened the polarization continues to change, because of change in acceleration. The electric field arises once more and can exceed the breakdown threshold. As a result a new breakdown will take place. Such a situation can be repeated several times. Effectively, it manifests itself in an increase of the pulse duration which may achieve a fraction of a nanosecond. Therefore, within our scenario the greater is the pulse duration the greater is its energy. Analogous interrelation is observed in experiment. .
## IV Energy Distribution of Hot Electrons in the Electric Field and Spectrum of the Sonoluminescence
We have already mentioned that the electric fields under consideration are very strong, inhomogeneous and change over time. However, the characteristic space and time scales of the field change are correspondingly $`l1\mu `$m and $`\tau _c1`$ns and they are much more than the relaxation length $`\lambda `$ and relaxation time $`\tau _\epsilon `$ of the hot electron energy which are respectively about $`(10÷100)\AA `$ and $`10^{13}`$s. . Therefore one can consider the local electron energy distribution function using well known results for homogeneous static electric fields .
The considered electric fields are superstrong, i.e. the following condition is valid:
$$qE\lambda h\omega _{ph}\epsilon _{ph},$$
(8)
where $`ϵ_{ph}`$ is the characteristic energy of local oscillations in water which is practically equal to the energy of optical phonons in ice, $`ϵ_{ph}(80÷100)`$meV . Using $`\lambda `$ $`(10÷100)\AA `$ one sees that the condition 8 is satisfied for fields $`E>10^7`$V/m. The condition 8 means that in these fields an electron acquires, on average, an energy $`qE\lambda `$ which is much more than $`ϵ_{ph}`$. Since in the process of the acceleration an electron generates many phonons the electron energy distribution is nearly isotropic in the momentum space. Specifically, in this case the energy distribution of hot conduction electrons is approximated with high accuracy, up to energies of electrical breakdown $`ϵ=E_g`$, by the Maxwell function :
$$f(ϵ)\mathrm{exp}(\frac{ϵ}{kT_e})$$
(9)
with an effective electron temperature $`T_e`$ which is determined by the balance equations for energy and momentum of the hot electrons
$$\frac{dϵ}{dt}=eEv_d\frac{ϵ_{ph}v_T}{\lambda }=0,$$
(10)
$$\frac{dp}{dt}=eEm_ev_d\frac{v_T}{\lambda }=0.$$
(11)
where $`v_T`$ is the effective thermal velocity of electrons. Since for the Maxwell distribution
$$\frac{3}{2}kT_e=\frac{1}{2}m_ev_T^2,$$
(12)
from Eqs.10, 11 and 12 it follows that
$$v_d=\sqrt{\epsilon _{ph}/m_e}$$
(13)
and
$$kT_e=(eE\lambda )^2/3\epsilon _{ph.}$$
(14)
From Eqs 12 and 13 as well from Eq. 8 one sees that $`v_T/v_{dr}=\sqrt{kT_e/\epsilon _{ph}}=qE\lambda /\epsilon _{ph}>>1.`$ This is precisely the condition of validity for the thermalization of the hot electrons and use of the Maxwell distribution.
Conduction electrons in water form polarons which are conventionally called hydrated electrons . However, in the strong electric field the polaron states decay and the current carriers are the usual conduction electrons. For purpose of estimation we assume that their effective mass $`m_e`$ is close to that of free electrons $`m_0`$. Then from Eq.13 it follows that $`v_d10^5`$m/s for $`ϵ_{ph}100`$meV. It has already been mentioned that the electric fields under consideration are $`E(10^8÷10^9)`$V/m. It follows from Eq.14 that for such fields $`kT_e(1÷10)`$eV, i.e., $`T_e(10^4÷10^5)`$K for $`\lambda (10÷100)\AA `$ . One sees that the electron temperatures can be several orders higher than the gas temperature in the bubble. Recall that for our scenario (see Sec.3) it sufficient that the gas temperature be about 7000K.
Now we will discuss the role of noble gases dissolved in water. An important feature of an noble gas atom is the existence of metastable states. The life time of the metastable states can reach several milliseconds when the noble gas atoms are impurities in solids (see, e.g., ). Since the nearest order in water is essentially the same as in solids it would be natural to expect that the life-time of the metastable states for the noble gas atoms in water is not less than $`1`$ns . Hot electrons not only generate new conduction electrons (the breakdown avalanche) but also excite the noble gas atoms into metastable states. Below we will estimate from the experimental data that the characteristic value of $`kT_e(2÷5)`$eV. However, a considerable number of electrons may have energies about 10eV and can excite the noble gas atoms.
Note that the energies necessary for excitation to the metastable states, $`\epsilon _m`$, for Xe, Kr, Ar, Ne, He are about 10eV and increase monotonically from Xe to He. In other words, it is much easier to excite Xe, Kr and Ar than Ne and He. This could be the reason for the increase in SBSL intensity in the series He-Xe, which has been observed experimentally .
As we have mentioned above, the life-time of the metastable states of noble gas atoms in water is expected to be fairly long. Thus, once excited a noble gas atom can remain in a metastable state for the entire time interval of positive acceleration and multiple breakdowns, $`\tau _c1`$ns. The major part of electrons has the energies $`\epsilon <\epsilon _{m\text{ ,}}E_g;`$they collide with these metastable atoms and transfer them to higher excited states. The radiation transitions between high-energy excited states of noble gas atoms govern the SBSL spectrum.
Note that the life-time of the metastable states can exceed the period of the acoustic wave. In this case there will be an effect of accumulation of the noble gas atoms in metastable states resulting in the gradual build up of the SBSL power during several periods of the acoustic wave.
A specific feature of a noble gas is an abundance of excited states with energies higher than those of metastable states and with high probabilities of the radiation transitions between them. That is why the optical spectra of the noble gas atoms contain many lines . Because of the Stark effect in the strong electric field these lines are split. Besides, in the active region the strong electric field changes at least several times and the observed radiation spectrum is a superposition of the spectra of the atoms in different strong electric fields. In other words one can consider the density of the atomic excited states as a constant.
The probability of a hot electron having an energy $`\epsilon <E_g`$ is given by Eq.9 and at every collision it transfers the noble gas atom in a metastable state to a state with energy $`\epsilon `$, the reference point of energy being the energy of the metastable state (for our estimations we consider only one metastable state). The concentration of atoms excited during the time of a single electric breakdown, i.e. the screening time, $`\tau _s,`$ to energies within an interval $`d\epsilon `$ reads
$$dn_n^{}=n_n^{}(\sigma _{ex}v_Tn)\tau _s\mathrm{exp}(\frac{\epsilon }{kT_e})\frac{d\epsilon }{kT_e}$$
(15)
where $`n_n^{}`$ is the number of the noble gas atoms in the metastable state, $`\sigma _{ex}`$ is the cross section of the impact excitation of an atom from the metastable state to a state with energy $`\epsilon `$. Atoms excited to states with the energy $`\epsilon `$ generally do not go to the ground state directly but through intermediate excited states. It reasonable to assume that with a high probability they radiate phonons with energy $`h\nu \epsilon `$. So the spectral density of the SBSL radiation energy per pulse for unit volume can be written as
$$\stackrel{~}{P}(h\nu )d(h\nu )=h\nu w_rn_n^{}(\sigma _{ex}v_Tn)\tau _s\mathrm{exp}(\frac{h\nu }{kT_e})\frac{d(h\nu )}{kT_e}$$
(16)
where $`w_r`$ is the probability of the spontaneous radiation transition. Taking into account that $`w_r=\frac{4(2\pi )^4\nu ^3}{3c^3h}D^2`$ where $`D`$ is modulus of matrix element of the dipole moment of the transition we find
$$\stackrel{~}{P}(h\nu )d(h\nu )=\frac{4(2\pi )^4D^2}{3c^3}n_n^{}(\sigma _{ex}v_Tn)\tau _s\nu ^4\mathrm{exp}(\frac{h\nu }{kT_e})\frac{d(h\nu )}{kT_e}.$$
(17)
Recall that $`T_e`$ in Eq. 17 is a function of the coordenates because $`T_eE^2`$ (see Eq.14 ). Putting, as before, $`E=fpE_sR_s^2/r^2`$ and integrating approximately Eq.17 we find the spectral power of SBSL for a single breakdown is
$$P=\stackrel{~}{P}(h\nu ,r)𝑑V=\frac{4(2\pi )^4D^2}{3c^3h}n_n^{}(\sigma _{ex}v_TN_t)\tau _s\nu ^3\mathrm{exp}(\frac{h\nu }{kT_{es}})$$
(18)
where $`kT_{es}=(eE_s\lambda )^2/3\epsilon _{ph\text{ }}`$and $`N_t`$ is, practically, the total number of electrons participating in the breakdown. Note that the approximative Eq.18 is not sensitive to the form of decay of the field, it is important only that $`E`$ decays more steeply than $`r^1`$, i.e. it is not important that for $`p`$ we use an expression for an incompressible liquid.
The observed spectra are cut off in the shortwave region because of the absorption of water. This can be taken into account by multiplying Eq.18 by exp(-$`\alpha (h\upsilon )L`$) where $`L`$ is the size of the acoustic resonator ($`L2.5`$cm ). Because of the Urbach absorption tails the radiation maximum is located considerably below than $`h\nu E_g6.5eV`$.
It should be emphasized that the spectrum given by Eq.18 resembles the black-body spectrum but temperature is given here by the effective temperature $`T_{es}`$ of hot electrons near the bubble surface. The value of $`T_{es}`$ may be much higher than the bubble gas temperature. Experimentally observed spectra can be fitted, in the wavelength interval (200$`÷700)`$ nm, to the black-body ones with temperatures (2$`÷`$5)10<sup>4</sup>K . According to Eq.14 such electron temperatures are reached at electric fields $`E_s(2÷10)10^8`$ V/m for values of $`\epsilon _{ph}0.1`$eV and $`\lambda (10÷100)`$ Å characteristic for water . As follows from Eqs.4 and $`p=\rho a`$ such fields are realized at accelerations $`a(10^{13}÷10^{14})`$ m/s$`^2.`$ Similar accelerations are reported in .
Integrating Eq.18 multiplied by exp(-$`\alpha (h\upsilon )L`$) we find an estimate of the radiation energy for a single breakdown
$$W_r=Pw_r^1exp(\alpha (h\upsilon )L)d(h\nu )(\sigma _{ex}v_TN_t)\tau _sn_n^{}h\overline{\nu }\frac{v_T}{v_d}(\sigma _{ex}R_sn_n^{})N_th\overline{\nu }$$
(19)
where $`h\overline{\nu }`$ is the characteristic photon energy which is close to the energy of the maximum of the observed spectrum, which is about ($`5÷6)eV`$ in the case of strong electric field of our interest. Recall that in Eq.19 $`\tau _sR_s/v_d`$ is the screening (breakdown) time and $`N_t4\pi R_s^2\epsilon _0E_s/e`$ is the number of transmitted electrons in the screening process (Sec. 3) and $`n_n^{}`$ is the concentration of noble gas atoms in a metastable state near the bubble surface.
As we have mentioned before it was shown in that because of dissolved gas diffusion and chemical reactions an accumulation of noble gas in the bubble takes place while the bubble pulsates and in stationary conditions the gas in the bubble consists almost entirely of noble gas. Therefore one can assume that the concentration of the noble gas atoms, $`n_n,`$near the bubble surface is close to its saturation value. Note that this value has to be estimated for atmospheric pressure and ambient temperature rather than for the high pressures and temperatures existing at the last stage of the collapse over a very short time. This is why we assume that $`n_n`$ (10$`{}_{}{}^{18}10^{19}`$) cm<sup>-3</sup>. Now we will argue that the value of $`n_n^{}`$ can be only one order of magnitude less than $`n_n`$. Indeed, for $`kT_e(3÷5)`$ eV obtained above by fitting the SBSL spectrum to the black-body spectrum, the number of electrons having an energy higher than 10eV, $`N_h`$ may be about 1$`0^3÷10^2`$ of the total electron number $`N_t10^6`$ i.e., $`N_h`$ is about 10$`{}_{}{}^{3}÷10^4`$. Such a ”superhot” electron excites a noble gas atom to the metastable state over the time $`\tau _{ex}=(\sigma _{ex}v_Tn_n)^110^{12}`$s for $`\sigma _{ex}10^{15}`$cm$`^2,`$ $`v_T10^6`$ m/s and $`n_n`$(10$`{}_{}{}^{18}÷10^{19}`$) cm$`^3.`$ Therefore a ”superhot” electron during its participation in screening excites $`\tau _s/\tau _{ex}`$ noble gas atoms and the total number of the excited atoms is $`N_h\tau _s/\tau _{ex}10^4÷10^5`$ if one takes into account that $`\tau _s10\tau _{ex}`$ $`10^{11}`$s. The total number of the noble gas atoms in the region of the strong field ($``$1$`\mu `$m$`{}_{}{}^{3})`$ is about $`10^6÷10^7`$, i.e. the percentage of the noble gas atoms excited to the metastable state during a single breakdown can be about 10$`\%`$. Since the life time of the metastable state is much more than $`\tau _c`$ and during this time interval several breakdowns can take place (Sec.3), one may expect that the number of the noble gas atoms excited to the metastable state during the interval $`\tau _c`$ is only one order of magnitude less than the total number of these atoms i.e. $`n_n^{}`$ $``$ 10<sup>18</sup> cm<sup>-3</sup>. If the life time of the metastable states exceeds the acoustic period ($`30\mu `$s) an accumulation of noble gas atoms in the metastable states can also occur during several acoustic periods.
Now we can estimate the total radiation energy. We have mentioned above that the maximum value of $`N_t10^6`$, $`R_s`$ 1$`\mu `$m, $`\sigma _{ex}10^{15}`$cm<sup>2</sup> and $`v_T>>v_d`$. Taking into account that the positive acceleration period $`\tau _c1`$ns and the breakdown time $`\tau _s`$ is about 0.01ns one can assume that the number of breakdowns is not less than 10. Thus one finds from Eq.19 that the maximum photon number in the SBSL pulse can be more than 10<sup>7</sup>. This value corresponds to the maximum photon number observed experimentally . Note that the radiation energy is a small part of the total energy of the flexoelectric field (see Ref).
## V Conclusions
Let us emphasize once more that within the mechanism of SBSL considered the main processes occur in water near the bubble surface unlike the most widely discussed mechanisms of SBSL . Within the framework of this mechanism many experimental data about the SBSL can be explained quite naturally. Of course, since the parameters of the collapsing bubble are not reliably known when the bubble radius is close to its minimum value we can present no more than order-of-magnitude estimations.
Let us summarize the main results.
(i). The minimum duration of the SBSL flash is determined by the single breakdown (screening) time, $`\tau _s10`$ps. A larger time is possible because of the multiplicity of breakdowns. The maximum duration is limited by $`\tau _c1`$ns. The longer is the duration the greater is the energy of the flash. This is in agreement with experiments where it was found that the pulse duration changes from 30ps to 400ps when its energy increases .
(ii). According to our estimations the maximum energy of the flash corresponds to $`10^7÷10^8`$ photons with energy $`(5÷6)`$eV, which also agrees with the experiment .
(iii). Within our scenario the noble gas atoms play an important role. They do not reduce so much the breakdown threshold as they govern the radiation process. Abundance of optical transitions in these atoms and inhomogeneous broadening because of the Stark effect explain the practically continuous character of the SBSL spectrum.
(iv). In agreement with the experiment the theoretical spectrum of the SBSL resembles the black-body spectrum but the temperature is given here by the effective temperature of the hot electrons which can be about several eV, what corresponds to the observed apparent radiation temperature . At the same time the gas temperature is not directly related to the radiation temperature and can be considerably less than 1eV.
(v). The increase of influence on SBSL intensity in the noble gas series He-Xe observed in the experiment is connected with the decrease, in this series, of the energy of the lowest metastable state.
(vi). The pulse width does not depend on the spectral range of the radiated photons. This also agrees with experiment .
(vii). The effect of synchronization of the light pulses observed experimentally is explained.
Note also that within our scenario of SBSL the influence of magnetic fields on SBSL is fairly weak. Their influence becomes appreciable when they are high enough to hamper the heating of electrons. This might be the reason for the decrease of the SBSL intensity in strong magnetic fields which has been observed experimentally.
It seems that the considered mechanism of SBSL is especially effective for water because of a lucky coincidence of several conditions:
(i) the strong geometric asymmetry of the water molecule, which accompanies its electric asymmetry, provides the needed sign of the flexoelectric coefficient. This sign is such that the electric field has the ”correct direction” (accelerates electron) over the same (very short) period in which electrons are generated because of the sharp increase of the bubble gas temperature. The temporal coincidence of these two cicumstances is a possible reason of the effect of synchronization of the light pulses observed experimentally .
(ii) water is a semiconductor with a relatively narrow band gap ($`E_g=6.5`$eV) and sufficienly wide conduction band, what is necessary for the heating of the electrons in strong electric field.
(iii) high solubility of noble gases in water what makes possible high concentrations of noble gas in water near the bubble filled by the noble gas accumulated there in the process of the bubble pulsations.
Fulfillment of these conditions gives the answer to the question why water is the friendliest liquid for SBSL .
Some experimentally found characteristics of SBSL have not been explained in this paper, in particular the dependence of the energy of the pulse upon the partial pressure of noble gases and water temperature . These dependences may be not necessarily a manifestation of the radiation processes studied in this paper but they could be a consequence of a set of various factors. They may be determined by the hydrodynamics of the pulsating bubble, the solubility of noble gases in water and also by the temperature dependence of the flexoelectric coefficient of water which still remains, unfortunately, unmeasured.
Although the main features of SBSL seem to be explainable within our scenario, these explanations are qualitative rather than quantitative. We are still a long way from a quantitative theory at present. One of the main aims of the paper is to stimulate some experiments that could either support or discard the proposed mechanis of SBSL.
(i) Measurements of flexoelectric coeficients of water.
(ii) A detailed study of radiation spectra at electric breakdown of water and ice with different concentrations of noble gases.
(iii) Study of influence of water conductivity on the SBSL intensity.
After performing these experiments it makes sense to develop the theory further. In particular, it will be necessary to study in more detail the kinetics of the screening process and of the excitation of noble gas atoms while taking into account the time variation of the distribution function of the hot electrons. This problem should be considered together with that of the determination of the spatial distribution of excited noble gas atoms near the bubble surface. Diffusion and relaxation of the excited atoms should be considered together, of course, with the processes that lead to noble gas accumulation in the bubble .
In our estimations we assumed that flexoelectric coefficient of water has its ”natural” value. This was sufficient for the theoretical estimations to be in agreement with the experimental data. In fact, because of the above-mentioned strong geometric asymmetry of the water molecule, the coefficient could be even larger. If this is the case some new phenomena will take place including influence of the flexoelectric effect on hydrodynamics of the pulsating bubble and X-ray radiation induced by electrons of very high energy.
We are grateful to V.Kholodnov for useful discussions. NG and VVO thank EU ESPRIT, the Spanish CSIC and NATO for Linkage Grant Ref. OUTRLG 970308. Also, APL thanks NATO for Linkage grant HTECHLG 971213.
|
Subsets and Splits
No community queries yet
The top public SQL queries from the community will appear here once available.